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# A note on multiple Dirichlet’s 𝑞-𝐿-function
## §1. Introduction
Let $`0<q<1`$ be and for any positive integer $`k`$, define its $`q`$-analogue $`[k]_q=\frac{1q^k}{1q}`$. Let $``$ be the field of complex numbers. The $`q`$-Bernoulli numbers are usually defined as
$$\beta _{0,q}=\frac{q1}{\mathrm{log}q},\text{ }\left(q\beta _q+1\right)^n\beta _{n,q}=\delta _{n,1},$$
$`1`$
where $`\delta _{n,1}`$ is Kronecker symbol and we use the usual convention about replacing $`\beta _q^i`$ by $`\beta _{i,q},`$ cf. . Note that $`lim_{q1}\beta _{k,q}=B_k`$, where $`B_k`$ are the $`k`$-th ordinary Bernoulli numbers. In , the $`q`$-Bernoulli polynomials are also defined by
$$F_q(t,x)=\frac{q1}{\mathrm{log}q}e^{\frac{t}{1q}}t\underset{n=0}{\overset{\mathrm{}}{}}q^{n+x}e^{[n+x]_qt}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\beta _{n,q}(x)}{n!}t^n,\text{ for }x.$$
$`2`$
From (1) and (2), we can derive the below formula:
$$\beta _{n,q}(x)=\underset{i=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{i}\right)q^{xi}\beta _{i,q}[x]_q^{ni}=\left(\frac{1}{1q}\right)^n\underset{i=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{i}\right)(1)^iq^{xi}\frac{i}{[i]_q},$$
where $`\left(\genfrac{}{}{0pt}{}{n}{i}\right)`$ is the binomial coefficient. In the recent paper the multiple generalized Bernoulli numbers attached to $`\chi `$, $`B_{n,\chi }^{(r)},`$ are defined by
$$F_\chi ^r(t)=\underset{a_1,\mathrm{},a_r=1}{\overset{f}{}}\frac{\chi (_{i=1}^ra_i)t^re^{(_{i=1}^ra_i)t}}{(e^{ft}1)^r}=\underset{n=0}{\overset{\mathrm{}}{}}B_{n,\chi }^{(r)}\frac{t^n}{n!},\text{ for }|t|<2\pi /f\text{, (see [1]). }$$
$`3`$
In this paper we consider the $`q`$-analogue of the above multiple generalized Bernoulli numbers attached to $`\chi `$. Finally we construct the multiple $`qL`$-function which interpolates the $`q`$-analogue of multiple generalized Bernoulli numbers attached to $`\chi `$ at negative integers and investigate its properties.
## 2. Multiple Dirichlet’s $`qL`$-function
Let $`\chi (1)`$ be the primitive Dirichlet’s character with conductor $`f`$. Then the generalized $`q`$-Bernoulli numbers attached to $`\chi `$ are defined as
$$F_{q,\chi }(t)=t\underset{a=1}{\overset{f}{}}\chi (a)\underset{n=0}{\overset{\mathrm{}}{}}q^{fn+a}e^{[fn+a]_qt}=\underset{n=0}{\overset{\mathrm{}}{}}\beta _{n,\chi ,q}\frac{t^n}{n!},\text{ cf. [5].}$$
$`4`$
By (2)and (4), we easily see that
$$\beta _{n,\chi ,q}=[f]_q^{n1}\underset{a=1}{\overset{f}{}}\chi (a)\beta _{n,q^f}(\frac{a}{f}).$$
By the meaning of Eq.(4), we can consider the multiple generalized $`q`$-Bernoulli numbers attached to $`\chi `$ as follows:
$`F_{q,\chi }^r(t)`$ $`=(t)^r{\displaystyle \underset{a_1,\mathrm{},a_r=1}{\overset{f}{}}}\chi ({\displaystyle \underset{i=1}{\overset{r}{}}}a_i){\displaystyle \underset{n_1,\mathrm{},n_r=0}{\overset{\mathrm{}}{}}}q^{_{i=1}^r(a_i+n_if)}e^{([_{i=1}^r(a_i+n_if)]_q)t}`$ $`5`$
$`={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\beta _{n,\chi ,q}^{(r)}{\displaystyle \frac{t^n}{n!}}.`$
Note that $`lim_{q1}\beta _{n,\chi ,q}^{(r)}=B_{n,\chi }^{(r)}.`$ Let $`\mathrm{\Gamma }(s)`$ be the gamma function. Then it is easy to see that
$`{\displaystyle \frac{1}{\mathrm{\Gamma }(s)}}{\displaystyle _0^{\mathrm{}}}F_{\chi ,q}^r(t)t^{s1r}𝑑t`$ $`6`$
$`={\displaystyle \underset{n_1.\mathrm{},n_r=1}{\overset{\mathrm{}}{}}}\chi (n_1+\mathrm{}+n_r)q^{n_1+\mathrm{}+n_r}{\displaystyle \frac{1}{\mathrm{\Gamma }(s)}}{\displaystyle _0^{\mathrm{}}}t^{s1}e^{[n_1+\mathrm{}+n_r]_qt}𝑑t`$
$`={\displaystyle \underset{n_1,\mathrm{},n_r=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\chi (n_1+\mathrm{}+n_r)}{[n_1+\mathrm{}+n_r]_q^s}}q^{n_1+\mathrm{}+n_r}.`$
Thus, we can define the multiple Dirichlet’s $`q`$-$`L`$-function as follows:
$$L_q^r(s,\chi )=\underset{n_1,\mathrm{},n_r=1}{\overset{\mathrm{}}{}}q^{n_1+\mathrm{}+n_r}\frac{\chi (n_1+\mathrm{}+n_r)}{[n_1+\mathrm{}+n_r]_q^s},\text{ for }s\text{}$$
$`7`$
where $`\chi `$ is non-trivial primitive Dirichlet’s character with conductor $`f.`$ By (5), (6) and (7), we easily see that
$$L_q^r(n,\chi )=(1)^r\frac{n!}{(n+r)!}\beta _{n+r,\chi ,q}^{(r)},\text{ for }n.$$
$`8`$
Let $`s`$ be a complex variable, $`a`$ and $`F`$ be integers with $`0<a<F`$. We now consider the function $`H_{r,q}(s;a_1,\mathrm{},a_r|F)`$ as follows:
$$H_{r,q}(s;a_1,\mathrm{},a_r|F)=\underset{\frac{}{m_1,\mathrm{},m_r>0m_ia_i(modF)}}{}\frac{q^{m_1+\mathrm{}+m_r}}{[m_1+\mathrm{}+m_r]_q^s}=[F]_q^s\zeta _{r,q^F}(s,\frac{a_1+\mathrm{}+a_r}{F}),$$
$`9`$
where $`\zeta _{r,q}(s,a)=_{n_1,\mathrm{},n_r=0}^{\mathrm{}}\frac{q^{n_1+\mathrm{}+n_r+a}}{[n_1+\mathrm{}+n_r+a]_q^s},\text{ (see [5, 6]). }`$
The function $`H_{r,q}(s;a_1,\mathrm{},a_r|F)`$ is a meromorphic for $`s`$ with poles at $`s=1,\mathrm{},r.`$ In , the multiple $`q`$-Bernoulli polynomials are defined by
$$(t)^r\underset{n_1,\mathrm{},n_r=0}{\overset{\mathrm{}}{}}q^{x+_{i=1}^rn_i}e^{[x+_{i=1}^rn_i]_qt}=\underset{n=0}{\overset{\mathrm{}}{}}B_{n,q}^{(r)}(x)\frac{t^n}{n!}.$$
$`10`$
For $`n,`$ it was well known that
$$\zeta _{r,q}(n,x)=(1)^r\frac{n!}{(n+r)!}B_{n+r,q}^{(r)}(x),\text{ cf. [5]. }$$
$`11`$
Let $`F`$ be the conductor of $`\chi .`$ Then the multiple Drichlet’s $`q`$-$`L`$-function can be expressed as the sum
$$L_q^r(s,\chi )=\underset{a_1,\mathrm{},a_r=1}{\overset{F}{}}\chi (a_1+\mathrm{}+a_r)H_{r,q}(s;a_1,\mathrm{},a_r|F),\text{ for }s\text{.}$$
$`12`$
By (9) and (12), we easily see that
$$H_{r,q}(n;a_1,\mathrm{},a_r|F)=[F]_q^n(1)^r\frac{n!}{(n+r)!}B_{n+r,q^F}^{(r)}(\frac{a_1+\mathrm{}+a_r}{F}),\text{ for }n.$$
$`13`$
From (5) and (10), we can also derive the below formula:
$$\beta _{n,\chi ,q}^{(r)}=[F]_q^n\underset{a_1,\mathrm{},a_r=1}{\overset{F}{}}\chi (a_1+\mathrm{}+a_r)B_{n,q^F}^{(r)}(\frac{a_1+\mathrm{}+a_r}{F}).$$
$`14`$
By using (12), (13) and (14), we easily see that
$$L_q^r(n,\chi )=(1)^r\frac{n!}{(n+r)!}\beta _{n+r,\chi ,q}^{(r)}.$$
In Eq.(10), we note that
$$B_{n,q}^{(r)}(x)=\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)[x]_q^{nk}q^{xk}B_{k,q}^{(r)},\text{ cf. [5, 6]},$$
$`15`$
where $`B_{n,q}^{(r)}=B_{n,q}^{(r)}(0).`$ From Eq.(15), we derive the function $`H_{r,q}(s;a_1,\mathrm{},a_r|F)`$ which is modified by
$`H_{r,q}(s;a_1,\mathrm{},a_r|F)`$ $`16`$
$`={\displaystyle \frac{1}{[F]_q^r}}{\displaystyle \frac{[_{i=1}^ra_i]_q^{s+r}}{_{j=1}^r(sj)}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{s+r}{k}}\right)\left({\displaystyle \frac{[F]_q}{[_{i=1}^ra_i]_q}}\right)^kq^{(_{i=1}^ra_i)k}B_{k,q}^{(r)}.`$
By (12) and (16), we obtain the following:
$`L_q^r(s,\chi )`$ $`={\displaystyle \frac{1}{_{j=1}^r(sj)}}{\displaystyle \frac{1}{[F]_q^r}}{\displaystyle \underset{a_1,\mathrm{},a_r=1}{\overset{F}{}}}\chi (a_1+\mathrm{}+a_r)[a_1+\mathrm{}+a_r]_q^{s+r}`$ $`17`$
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{rs}{m}}\right)\left({\displaystyle \frac{[F]_q}{[a_1+\mathrm{}+a_r]_q}}\right)^mq^{(a_1+\mathrm{}+a_r)m}B_{m,q}^{(r)}.`$
Finally, we suggest the below question
###### Question
Is it possible to give the $`p`$-adic analogue of Eq.(17) which can be viewed as interpolating, in the same way that $`L_{p,q}(s,\chi )`$ interpolates $`L_q(s,\chi )\text{ in [2, 7 ] }\mathrm{?}`$
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# Odd Triplet Superconductivity and Related Phenomena in Superconductor-Ferromagnet Structures
## I Introduction
Although superconductivity has been discovered by H. Kammerlingh Onnes almost one century ago (1911), the interest in studying this phenomenon is far from declining. The great attention to superconductivity within the last 15 years is partly due to the discovery of the high temperature superconductors (HTSC) Bednorz and Müller (1986), which promises important technological applications. It is clear that issues such as the origin of the high critical temperature superconductivity, effects of external fields and impurities on HTCS, etc, will remain fields of interest for years to come.
Due to the successful investigations of the HTSC and its possible technological applications, the interest in studying properties of traditional (low $`T_c`$) superconductors was not as broad. Nevertheless this field has also undergone a tremendous development. Technologically, the traditional superconductors are often easier to manipulate than high $`T_c`$ cuprates. One of the main achievements of the last decade is the making of high quality contacts between superconductors and normal metals $`(S/N)`$, superconductors and ferromagnets $`(S/F)`$, superconductors and insulators $`(S/I)`$, etc. All these heterostructures can be very small with the characteristic sizes of submicrometers.
This has opened a new field of research. The small size of these structures provides the coherence of superconducting correlations over the full length of the $`N`$ region. The length of the condensate penetration into the $`N`$ region $`\xi _N`$ is restricted by decoherence processes (inelastic or spin-flip scattering). At low temperatures the characteristic length over which these decoherence processes occur may be quite long (a few microns). Superconducting coherent effects in $`S/N`$ nanostructures, such as conductance oscillations in an external magnetic field, were studied intensively during the last decade (see for example the review articles by Beenakker (1997); Lambert and Raimondi (1998)).
The interplay between a superconductor $`(S)`$ and a normal metal $`(N)`$ in simpler types of $`S/N`$ structures (for example, $`S/N`$ bilayers) has been under study for a long time and the main physics of this so called proximity effect is well described in the review articles by de Gennes (1964) and Deutscher and de Gennes (1969). In these works it was noticed that not only the superconductor changes the properties of the normal metal but also the normal metal has a strong effect on the superconductor. It was shown that near the $`S/N`$ interface the superconductivity is suppressed over the correlation length $`\xi _S`$, which means that the order parameter $`\mathrm{\Delta }`$ is reduced at the interface in comparison with its bulk value far away from the interface. At the same time, the superconducting condensate penetrates the normal metal over the length $`\xi _N`$, which at low temperatures may be much larger than $`\xi _S`$. Due to the penetration of the condensate into the normal metal over large distances the Josephson effect is possible in $`S/N/S`$ junctions with the thicknesses of the $`N`$ regions of the order of a few hundreds nanometers. The Josephson effects in $`S/N/S`$ junctions were studied in many papers and a good overview, both experimental and theoretical, is given by Kulik and Yanson. (1970), Likharev (1979), and Barone and Paterno (1982).
The situation described above is quite different if an insulating layer $`I`$ is placed between two superconductors. The thickness of the insulator in $`S/I/S`$ structures cannot be as large as of the normal metals because electron wave functions decay in the insulator on atomic distances. As a consequence, the Josephson current is extremely small in $`S/I/S`$ structures with a thick insulating layer.
But what about $`S/F/S`$ heterojunctions, where $`F`$ denotes a ferromagnetic metal? In principle, the electron wave function can extend in the ferromagnet over a rather large distance without a considerable decay. However, it is well known that electrons with different spins belong to different energy bands. The energy shift of the two bands can be considered as an effective exchange field acting on the spin of the electrons. The condensate of conventional superconductors is strongly influenced by this exchange field of the ferromagnets and usually this reduces drastically the superconducting correlations.
The suppression of the superconducting correlations is a consequence of the Pauli principle. In most superconductors the wave function of the Cooper pairs is singlet so that the electrons of a pair have opposite spins. In other words, both the electrons cannot be in the same state, which would happen if they had the same spin. If the exchange field of the ferromagnet is sufficiently strong, it tries to align the spins of the electrons of a Cooper pair parallel to each other, thus destroying the superconductivity. Regarding the $`S/F`$ interfaces and the penetration of the condensate into the ferromagnet, these effects mean that the superconducting condensate decays fast in the region of the ferromagnet. A rough estimate leads to the conclusion that the ratio of the condensate penetration depth in ferromagnets to the one in non-magnetic metals with a high impurity concentration is of the order of $`\sqrt{T_c/h}`$, where $`h`$ is the exchange energy and $`T_c`$ is the critical temperature of the superconducting transition. The exchange energy in conventional ferromagnets like $`Fe`$ or $`Co`$ is several orders of magnitude higher than $`T_c`$ and therefore the penetration depth in the ferromagnets is much smaller than that in the normal metals.
Study of the proximity effect in the $`S/F`$ structures started not long ago but it has already evolved into a very active field of research (for a review see Izyumov *et al.* (2002); Golubov *et al.* (2004); Buzdin (2005a); Lyuksyutov and Pokrovsky (2004)). The effect of the suppression of superconductivity by the ferromagnetism is clearly seen experimentally and this corresponds to the simple picture of the destruction of the singlet superconductivity by the exchange field as discussed above.
At first glance, it seems that due to the strong suppression of the superconductivity the proximity effect in $`S/F`$ structures is less interesting than in the $`S/N`$ systems. However, this is not so because the physics of the proximity effect in the $`S/F`$ structures is not exhausted by the suppression of the superconductivity and new very interesting effects come into play. Moreover, under some circumstances superconductivity is not necessarily suppressed by the ferromagnets because the presence of the latter may lead to a triplet superconducting pairing Bergeret *et al.* (2001a); Kadigrobov *et al.* (2001). In some cases not only the ferromagnetism tends to destroy the superconductivity but also the superconductivity may suppress the ferromagnetism Buzdin and Bulaevskii (1988); Bergeret *et al.* (2000). This may concern “real” strong ferromagnets like iron or nickel with a Curie temperature much larger than the transition temperature of the superconductor.
In all, it is becoming more and more evident from recent experimental and theoretical studies that the variety of non-trivial effects in the $`S/F`$ structures exceeds considerably what one would have expected before. Taking into account possible technological applications, there is no wonder that $`S/F`$ systems attract nowadays a lot of attention.
This review article is devoted to the study of new “exotic” phenomena in the $`S/F`$ heterojunctions. By the word “exotic” we mean phenomena that could not be expected from the simple picture of a superconductor in contact with a homogeneous ferromagnet. Indeed, the most interesting effects should occur when the exchange field is not homogeneous. These non-homogeneities can be either intrinsic for the ferromagnetic material, like e.g. domain walls, or arise as a result of experimental manipulations, such as multilayered structures with different directions of the magnetization, which can also be spoken of as a non-homogeneous alignment of the magnetic moments.
Of course, we are far from saying that there is nothing interesting to be seen when the exchange field is homogeneous. Although it is true that in this case the penetration depth of the superconducting condensate into the ferromagnet is short, the exponential decay of the condensate function into ferromagnets is accompanied by oscillations in space. These oscillations lead, for example, to oscillations of the critical superconducting temperature $`T_c`$ and the critical Josephson current $`I_c`$ in $`S/F`$ structures as a function of the thickness $`d_F`$. Being predicted by Buzdin and Kupriyanov (1990) and Radovic *et al.* (1991), the observation of such oscillatory behavior was first reported by Jiang *et al.* (1995) on $`Gd/Nb`$ structures. Indications to a non-monotonic behavior of $`T_c`$ as a function of $`d_F`$ was also reported by Wong *et al.* (1986); Strunk *et al.* (1994); Mercaldo *et al.* (1996); Obiand *et al.* (1999); Velez *et al.* (1999); Mühge *et al.* (1996).
However, in other experiments the dependence of $`T_c`$ on $`d_F`$ was monotonic. For example in Ref. Bourgeois and Dynes (2002) the critical temperature of the bilayer Pb/Ni decreased by increasing the F layer thickness $`d_F`$ in a monotonic way. In the experiments by Mühge *et al.* (1998) on $`Fe/Nb/Fe`$ structures and by Aarts *et al.* (1997) on $`V/Fe`$ systems both a monotonic and non-monotonic behavior of $`T_c`$ has been observed. This different behavior was attributed to changes of the transmittance of the $`S/F`$ interface. A comprehensive analysis taking into account the samples quality was made for different materials by Chien and Reich (1999).
More convincing results were found by measuring the Josephson critical current in a $`S/F/S`$ junction. Due to the oscillatory behavior of the superconducting condensate in the $`F`$ region the critical Josephson current should change its sign in a $`S/F/S`$ junction ($`\pi `$junction). This phenomenon predicted long ago by Bulaevskii *et al.* (1977) has been confirmed experimentally only recently Ryazanov *et al.* (2001); Kontos *et al.* (2001, 2002); Blum *et al.* (2002); Sellier *et al.* (2004); Bauer *et al.* (2004).
Experiments on transport properties of $`S/F`$ structures were also performed in the last years. For example, Petrashov *et al.* (1999) and Giroud *et al.* (1998) observed an unexpected decrease of the resistance of a ferromagnetic wire attached to a superconductor when the temperature is lowered below $`T_c`$. In both of the experiments strong ferromagnets $`Ni`$ and $`Co`$, respectively, were used. One would expect that the change of the resistance must be very small due to the destruction of the superconductivity by the ferromagnets. However, the observed drop was about $`10\%`$ and this can only be explained by a long-range proximity effect.
This raises a natural question: how can such long range superconducting effects occur in a ferromagnet with a strong exchange field? We will see in the subsequent chapters that provided the exchange field is not homogenous a long-range component of the condensate may be induced in the ferromagnet. This component is in a triplet state and can penetrate the $`F`$ region over distances comparable with $`\xi _N`$, as in the case of a normal metal.
We outline now the structure of the present review.
In Chapter II we discuss the proximity effects in $`S/N`$ structures and $`S/F`$ structures with a homogeneous magnetization. Chapter II may serve as an introduction into the field. The main results illustrated there have been presented in other reviews and we discuss them here in order to give the reader an impression about works done previously. Chapter II can also help in getting the basic knowledge about calculational methods used in subsequent chapters. One can see from this discussion that already homogeneous ferromagnets in contact to superconductors lead to new and interesting physics.
Nevertheless, the non-homogeneities bring even more. We review below several different effects arising in the non-homogeneous situation. It turns out that a non-homogeneous alignment of the exchange field leads to a complicated spin structure of the superconducting condensate. As a result, not only the singlet component of the condensate exists but also a triplet one with all possible projections of the total spin of the Cooper pair ($`S_z=0,\pm 1`$). In contrast to the singlet component, the spins of the electrons in the triplet one with $`S_z=\pm 1`$ are parallel to each other. The condensate (Gor’kov) function $`f_{tr}`$ of the triplet state is an odd function of the Matsubara frequency<sup>1</sup><sup>1</sup>1Superconductivity caused by the triplet odd in $`\omega `$ condensate is called here odd superconductivity.. The singlet part $`f_{sng}`$ is, as usual, an even function of $`\omega `$ but it changes sign when interchanging the spin indices. This is why the anticommutation relations for the equal-time functions $`f_{tr}(t,t)`$ and $`f_{sng}(t,t)`$ remain valid; in particular, $`f_{tr}(t,t)=0`$ and $`f_{sng}(t,t)0`$. Therefore the superconductivity in the $`S/F`$ structures can be very unusual: alongside with the usual BCS singlet part it may contain also the triplet part which is symmetric in the momentum space (in the diffusive case) and odd in frequency. Both components are insensitive to the scattering by non-magnetic impurities and hence survive in the $`S/F`$ structures even if the mean free path $`l`$ is short. When generated, the triplet component is not destroyed by the exchange field and can penetrate the ferromagnet over long distances of the order of $`\xi _N=\sqrt{D_F/2\pi T}`$.
In Chapter III we analyze properties of this new type of superconductivity that may arise in $`S/F`$ structures. We emphasize that this triplet superconductivity is generated by the exchange field and, in the absence of the field, one would have the conventional singlet pairing.
The superconductor-ferromagnet multilayers are a very interesting and natural object for observation of Josephson effects. The thickness of both the superconductor and ferromagnetic layers, as well as the transparency of the interface, can be varied experimentally. This makes possible a detailed study of many interesting physical quantities. As we have mentioned, an interesting manifestation of the role played by the ferromagnetism is the possibility of a $`\pi `$-junction.
However, this is not the only interesting effect and several new ones have been recently proposed theoretically. As not so much time has been passed, they have not been confirmed experimentally unambiguously but there is no doubt that proper experiments will have been performed soon. In Chapter IV we discuss new Josephson effects in multilayered $`S/F`$ structures taking into account a possible change of the mutual direction of the magnetization in the ferromagnetic layers. We discuss a simple situation when the directions of the magnetic moments in a $`SF/I/FS`$ structure are collinear and the Josephson current flows through an insulator ($`I`$) but not through the ferromagnets. Naively, one could expect that the presence of the ferromagnets leads to a reduction of the value of the critical current. However, the situation is more interesting. The critical current is larger when the magnetic moments of the $`F`$-layers are antiparallel than when they are parallel. Moreover, it turns out that the critical current for the antiparallel configuration is even larger than the one in the absence of any ferromagnetic layer. In other words, the ferromagnetism can enhance the critical current Bergeret *et al.* (2001b)
Another setup is suggested in order to observe the odd triplet superconductivity discussed in Chapter III. Here the current should flow through the ferromagnetic layers. Usually, one could think that the critical current would just decay very fast with increasing the thickness of the ferromagnetic layer. However, another effect is possible. Changing the mutual direction of the additional ferromagnetic layers one can generate the odd triplet component of the superconducting condensate. This component can penetrate the ferromagnetic layer as if it were a normal metal, leading to large values of the critical current.
Such structures can be of use for detecting and manipulating the triplet component of the condensate in experiments. In particular, we will see that in some S/F structures the type of superconductivity is different in different directions: in the longitudinal direction (in-plane superconductivity) it is caused mainly by the singlet component, whereas in the transversal direction the triplet component mainly contributes to the superconductivity. We discuss also possibilities of an experimental observation of the triplet component.
Although the most pronounced effect of the interaction between the superconductivity and ferromagnetism is the suppression of the former by the latter, the opposite is also possible and this is discussed in Chapter V. Of course, a weak ferromagnetism should be strongly affected by the superconductivity and this situation is realized in so called magnetic superconductors Bulaevskii *et al.* (1985). Less trivial is that the conventional strong ferromagnets in the $`S/F`$ systems may also be considerably affected by the superconductivity. This can happen provided the thickness of the ferromagnetic layer is small enough. Then, it can be energetically more profitable to enforce the magnetic moment to rotate in space than to destroy the superconductivity. If the period of such oscillations is smaller than the size of the Cooper pairs $`\xi _S`$, the influence of the magnetism on the superconductor becomes very small and the superconductivity is preserved. In thick layers such an oscillating structure (cryptoferromagnetic state) would cost much energy and the destruction of the superconductivity is more favorable. Results of several experiments have been interpreted in this way Garifullin *et al.* (2002); Mühge *et al.* (1998).
Another unexpected phenomenon, namely, the inverse proximity effect is also presented in Chapter V. It turns out that not only the superconducting condensate can penetrate the ferromagnets but also a magnetic moment can be induced in a superconductor that is in contact with a ferromagnet. This effect has a very simple explanation. There is a probability that some of the electrons of Cooper pairs enter the ferromagnet and its spin tends to be parallel to the magnetic moment. At the same time, the spin of the second electron of the Cooper pair should be opposite to the first one (the singlet pairing or the triplet one with $`S_z=0`$ is assumed). As a result, a magnetic moment with the direction opposite to the magnetic moment in the ferromagnet is induced in the superconductor over distances of the superconducting coherence length $`\xi _S`$.
In principle, the total magnetic moment can be completely screened by the superconductor. Formally, the appearance of the magnetic moment in the superconductor is due the triplet component of the condensate that is induced in the ferromagnet $`F`$ and penetrates into the superconductor $`S`$. It is important to notice that this effect should disappear if the superconductivity is destroyed by, e.g. heating, and this gives a possibility of an observation of the effect. In addition to the Meissner effect, this is one more mechanism of the screening of the magnetic field by superconductivity. In contrast to the Meissner effect where the screening is due to the orbital electron motion, this is a kind of spin screening.
Finally, in Chapter VI we discuss the results presented in the review and try to anticipate future directions of the research. The Appendix A contains necessary information about the quasiclassical approach in the theory of superconductivity.
We should mention that several review articles on $`S/F`$ related topics have been published recently Izyumov *et al.* (2002); Golubov *et al.* (2004); Buzdin (2005a); Lyuksyutov and Pokrovsky (2004). In these reviews various properties of the $`S/F`$ structures are discussed for the case of a homogeneous magnetization. In the review by Lyuksyutov and Pokrovsky (2004) the main attention is paid to effects caused by a magnetic interaction between the ferromagnet and superconductor (for example, a spontaneous creation of vortices in the superconductor due to the magnetic interaction between the magnetic moment of vortices and the magnetization in the ferromagnet). We emphasize that, in contrast to these reviews, we focus on the discussion of the triplet component with all possible projections of the magnetic moment ($`S_z=0,\pm 1`$) arising only in the case of a nonhomogeneous magnetization. In addition, we discuss the inverse proximity effect, that is, the influence of superconductivity on the magnetization $`M`$ of $`S/F`$ structures and some other effects. Since the experimental study of the proximity effects in the $`S/F`$ structures still remains in its infancy, we hope that this review will help in understanding the conditions under which one can observe the new type of superconductivity and other interesting effects and hereby will stimulate experimental activity in this hot area.
## II The proximity effect
In this section we will review the basic features of the proximity effect in different heterostructures. The first part is devoted to superconductors-normal metals structures, while in the second part superconductors in contact with homogeneous ferromagnets are considered.
### II.1 Superconductor-normal metal structures
If a superconductor is brought in contact with a non-superconducting material the physical properties of both materials may change. This phenomenon called the proximity effect has been studied for many years. Both experiments and theory show that the properties of superconducting layers in contact with insulating ($`I`$) materials remain almost unchanged. For example, for superconducting films evaporated on glass substrates, the critical temperature $`T_c`$ is very close to the bulk value. However, physical properties of both metals of a normal metal/superconductor ($`N/S`$, see FIG. 1) heterojunction with a high $`N/S`$ interface conductance can change drastically.
Study of the proximity effect goes back to the beginning of the 1960’s and was reviewed in many publications (see, e.g. de Gennes (1964) and Deutscher and de Gennes (1969)). It was found that the critical temperature of the superconductor in a $`S/N`$ system decreased with increasing $`N`$ layer thickness. This behavior can be interpreted as the breaking down of some Cooper pairs due to the penetration of one of the electrons of the pairs into the normal metal where they are no longer attracted by the other electrons of the pairs.
At the same time, penetrating into the normal metal the Cooper pairs induce superconducting correlations. For example, the influence of the superconductivity on the physical properties of the $`N`$ metal manifests itself in the suppression of the density of states. Experiments determining the density of states of S/N bilayers with the help of tunneling spectroscopy were performed many years ago Toplicar and Finnemore (1977); Adkins and Kington (1969). While spatially resolved density of states were later measured by Gupta *et al.* (2004); Guéron *et al.* (1996); Anthore *et al.* (2003) (see FIG. 2).
The simplest way to describe the proximity effect is to use the Ginzburg-Landau (G-L) equation for the order parameter $`\mathrm{\Delta }`$ Ginzburg and Landau (1950). This equation is valid if the temperature is close to the critical temperature of the superconducting transition $`T_c`$. In this case all quantities can be expanded in the small parameter $`\mathrm{\Delta }/T_c`$ and slow variations of the order parameter $`\mathrm{\Delta }`$ in space.
Using the G-L equation written as
$$\xi _{GL}\frac{^2\mathrm{\Delta }\left(𝐫\right)}{𝐫^2}+\mathrm{\Delta }\left(𝐫\right)(\mathrm{sgn}(T_{c,N,S}T)\mathrm{\Delta }^2\left(𝐫\right)/\mathrm{\Delta }_0^2)=0.$$
(1)
one can describe the spatial distribution of the order parameter in any $`N/S`$ structure. Here $`\xi _{GL}`$ is the coherence length in the $`N`$ and $`S`$ regions at temperatures close to the critical temperatures $`T_{cN,S}`$. In the diffusive limit this length is equal to
$$\xi _{GL}=\sqrt{\pi D_{N,S}/8|TT_{cN,S}|}$$
(2)
where $`D_{N,S}`$ is the diffusion coefficient in the $`N`$ and $`S`$ regions. The quantity $`\mathrm{\Delta }_0`$ is the bulk value of the order parameter in the superconductor $`S`$. It vanishes when T reaches the transition temperature $`T_c`$.
It should be noticed though, that the region of the applicability of Eq. (1) for the description of the $`S/N`$ contacts is rather restricted. Of course, the temperature must be close to the transition temperature $`T_c`$ but this is not sufficient. The G-L equation describes variations of the order parameters correctly only if they are slow on the scales $`v_F/T_c`$ for the clean case or $`\sqrt{D_{N,S}/T_c}`$ in the diffusive “dirty” case. This can be achieved if the normal metal is a superconducting material taken at a temperature exceeding its transition temperature $`T_{cN}`$ and the transition temperatures $`T_{cS}`$ and $`T_{cN}`$ are close to each other. If this condition is not satisfied (e.g. $`T_{cN}=0`$) one should use more complicated equations even at temperatures close to $`T_{cS}`$, as we show below.
It follows from Eq. (1) that in the $`S`$ region, far from the $`N/S`$ interface, the order parameter $`\mathrm{\Delta }\left(𝐫\right)`$ equals the bulk value $`\mathrm{\Delta }_0`$, whereas in the $`N`$ region $`\mathrm{\Delta }\left(𝐫\right)`$ decays exponentially to zero on the length $`\xi _N`$.
The order parameter $`\mathrm{\Delta }\left(𝐫\right)`$ is related to the condensate function (or Gor’kov function)
$$f(t,t^{})=<\psi _{}(t)\psi _{}(t^{})>$$
(3)
via the self-consistency equation
$$\mathrm{\Delta }_{N,S}(t)=\lambda _{N,S}f(t,t),$$
(4)
where $`\lambda _{N,S}`$ is the electron-electron coupling constant leading to the formation of the superconducting condensate.
Eq.(1) describes actually a contact between two superconductors with different critical temperatures $`T_{cN,S}`$, when the temperature is chosen between $`T_{cS}`$ and $`T_{cN}`$. In the case of a real normal metal the coupling constant $`\lambda _N`$ is equal to zero and therefore $`\mathrm{\Delta }_N=0`$. However, this does not imply that the normal metal does not possess superconducting properties in this case. The point is that many important physical quantities are related not to the order parameter $`\mathrm{\Delta }`$ but to the condensate function $`f`$, Eq. (3). For example, the non-dissipative condensate current $`j_S`$ is expressed in terms of the function $`f`$ but not of $`\mathrm{\Delta }`$. If the contact between the $`N`$ and $`S`$ regions is good, the condensate penetrates the normal metal leading to a finite value of $`j_S0`$ in this region.
In the general case of an arbitrary $`\lambda _N`$ it is convenient to describe the penetration of the condensate (Cooper pairs) into the $`N`$ region in the diffusive limit by the Usadel equation Usadel (1970) which is valid for all temperatures and for distances exceeding the mean free path $`l`$. This equation determines the so called quasiclassical Green’s functions (see Appendix A) which can be conveniently used in problems involving length scales larger than the Fermi wave length $`\lambda _F`$ and energies much smaller than the Fermi energy. Alternatively, one could try to find an exact solution (the normal and anomalous electron Green’s functions) for the Gor’kov equations, but this is in most of the cases a difficult task.
In order to illustrate the convenience of using the quasiclassical method we calculate now the change of the tunnelling density of states (DOS) in the normal metal due to the proximity effect with the help of the Usadel equation. The DOS is a very important quantity that can be measured experimentally and, at the same time, can be calculated without difficulties.
We consider the $`S/N`$ structure shown in FIG. 1 and assume that the system is diffusive (i.e. the condition $`ϵ\tau <<1`$ is assumed to be fulfilled, where $`\tau `$ is the momentum relaxation time and $`ϵ`$ is the energy) and that the transparency of the $`S/N`$ is low enough. In this case the condensate Green’s function $`f(ϵ)=𝑑tf(tt^{})\mathrm{exp}(iϵ(tt^{}))`$ is small in the $`N`$ region and the Usadel equation can be linearized (see Appendix A).
Assuming that the boundary between the superconductor and normal metal is flat and choosing the coordinate $`x`$ perpendicular to the boundary we reduce the Usadel equation in the $`N`$ region to the form
$$D_N^2f/x^2+2iϵf=0,$$
(5)
where $`D_N=v_Fl/3`$ is the classical diffusion coefficient.
The solution of this equation can be found easily and we write it as
$$f=f_0\mathrm{exp}\left(x\sqrt{2iϵ/D_N}\right),$$
(6)
where $`f_0`$ is a constant that is to be found from the boundary conditions.
We see that the solution for the condensate function $`f`$ decays in the $`N`$ region exponentially at distances inversely proportional to $`\sqrt{ϵ}`$. In many cases the main contribution to physical quantities comes from the energies $`ϵ`$ of the order of the temperature, $`ϵT`$. This means that the superconducting condensate penetrates the $`N`$ region over distances of the order of $`\xi _N=\sqrt{D_N/2\pi T}`$. At low temperatures this distance becomes very large, and if the thickness of the normal metal layer is smaller than the inelastic relaxation length, the condensate spreads throughout the entire $`N`$ region.
In order to calculate the DOS it is necessary to know the normal Green’s function $`g`$ which is related to the condensate function $`f`$ via the normalization condition (see Appendix A)
$$g^2f^2=1$$
(7)
Eqs. (5) and (7) are written for the retarded Green’s function ($`f=f^R`$, see Appendix A). They are also valid for the advanced Green’s functions provided $`(ϵ+i0)`$ is replaced by $`(ϵi0)`$. The normalized density-of-states (we normalize the DOS to the DOS of non-interacting electrons) $`\nu (ϵ)`$ is given by the expression
$$\nu (ϵ)=\mathrm{Re}g(ϵ)$$
(8)
As the condensate function $`f`$ is small, a correction $`\delta \nu `$ to the DOS due to the proximity effect is also small. In the main approximation the DOS $`\nu `$ is very close to its value in the absence of the superconductor, $`\nu 1`$. Corrections to the DOS $`\delta \nu `$ are determined by the condensate function $`f`$. From Eq. (7) one gets
$$\delta \nu f^2/2.$$
Now we consider another case when the function $`f`$ is not small and the correction $`\delta \nu `$ is of the order of unity. Then the linearized Eq. (5) may no longer be used and we should write a more general one. For a S/N system the general equation can be written as (see Appendix A)
$$iD_{S,N}(\widehat{g}\widehat{g}/x)_{S,N}/x+ϵ[\widehat{\tau }_3,\widehat{g}_{S,N}]+[\widehat{\mathrm{\Delta }}_S,\widehat{g}_{S,N}]=0.$$
(9)
This non-linear equation contains the quasiclassical matrix Green’s function $`\widehat{g}`$. Both normal $`g`$ and anomalous Green’s functions $`f`$ enter as elements of this matrix through the following relation (the phase in the superconductor is set to zero)
$$\widehat{g}_N=g_N\widehat{\tau }_3+f_Ni\widehat{\tau }_2,$$
(10)
where $`\tau _i`$, $`i=1,2,3`$ are Pauli matrices and $`[A,B]=ABBA`$ is the commutator for any matrices $`A`$ and $`B`$.
We consider a flat $`S/N`$ interface normal to the $`x`$-axis. The normal metal occupies the region $`0<x<d_N.`$ We assume that in the normal metal $`N`$ there is no electron-electron interaction ($`\lambda _N=0`$, see Eq.(4)) so that in this region the superconducting order parameter vanishes, $`\mathrm{\Delta }_N=0.`$ In the superconductor the matrix $`\widehat{\mathrm{\Delta }}_S`$ has the structure $`\widehat{\mathrm{\Delta }}_S=\mathrm{\Delta }i\widehat{\tau }_2.`$
At large distances from the $`S/N`$ interface the Green’s functions $`\widehat{g}_S`$ of the superconductor do not depend on coordinates and the first term in Eq. (9) can be neglected. Then we obtain a simpler equation
$$ϵ[\widehat{\tau }_3,\widehat{g}_S]+\mathrm{\Delta }[i\widehat{\tau }_2,\widehat{g}_S]=0.$$
(11)
The solution for this equation satisfying the normalization condition (7) is
$$g_{BCS}=ϵ/\xi _ϵ;\text{ }f_{BCS}=\mathrm{\Delta }/\xi _ϵ,$$
(12)
where $`\xi _ϵ=\sqrt{ϵ^2\mathrm{\Delta }^2}`$. Eq. (12) is just the BCS solution for a bulk superconductor.
In order to find the matrix $`\widehat{g}(x)`$ both in the $`S`$ and $`N`$ regions, Eq.(9) should be complemented by boundary conditions and this is a non-trivial problem. Starting from the initial Hamiltonian $`\widehat{H}_{tot}`$, Eq. (22), one does not need boundary conditions at the interface between the superconductor and the ferromagnet because the interface can be described by introducing a proper potential in the Hamiltonian. In this case the self-consistent Gor’kov equations can be derived.
However, deriving the Usadel equation, Eq. (178), we have simplified the initial Gor’kov equations using the quasiclassical approximation. Possible spatial variation of the interface potential on a very small scale, due to the roughness of the interface cannot be included in the quasiclassical equations. Nevertheless, this problem is avoided deriving the quasiclassical equations at distances from the interface exceeding the wavelength. In the diffusive case one should go away from the interface to distances larger than the mean free path $`l`$. In order to match the solutions in the superconducting and non-superconducting regions one should solve exact the equations near the interface and compare the asymptotic behavior of this solution at large distances with the solutions of the Usadel equation. This procedure is equivalent to solving the quasiclassical equations with some boundary conditions. These conditions were derived by Zaitsev (1984) and Kuprianov and Lukichev (1988) (see also Appendix A, where these conditions are discussed in more details). For the present case they can be written as
$$2\gamma _{S,N}(\widehat{g}\widehat{g}/x)_{S,N}=[\widehat{g}_S,\widehat{g}_N]|_{x=0}$$
(13)
where $`\gamma _{S,N}=R_b\sigma _{S,N}`$, $`R_b,`$ measured in units $`\mathrm{\Omega }cm^2`$, is the $`S/N`$ interface resistance per unit area in the normal state, and $`\sigma _{S,N}`$ are the conductivities of the $`S`$ and $`N`$ metals in the normal state.
We assume that the thickness of the normal metal $`d_N`$ is smaller than the characteristic penetration length $`\xi _N(ϵ)=\sqrt{D_N/ϵ}`$ for a given energy $`ϵ`$, that is<sup>2</sup><sup>2</sup>2The quantity $`E_{Th}=D_N/d_N^2`$ is the so called Thouless energy $`ϵ<<D_N/d_N^2=E_T`$. Then the functions $`g`$ and $`f`$ remain almost constant over the thickness of the metal, and for finding them, one can average the Usadel equation over the thickness. In other words, we assume that the thickness $`d_N`$ of the N layer satisfies the inequality
$$d_{N\text{ }}<<\sqrt{D_N/ϵ},\text{ }ϵϵ_{bN}$$
(14)
($`ϵ_{bN}`$ is a characteristic energy in the DOS of the $`N`$ layer) and average Eq. (9) over the thickness $`d_N`$ considering $`\widehat{g}_N`$ as a constant in the second term of this equation. Using the boundary condition, Eq.(13), the first term in Eq. (9) can be replaced after the integration by the commutator $`[\widehat{g}_S,\widehat{g}_N]|_{x=0}`$. At $`x=d_N`$ the product $`(\widehat{g}\widehat{g}/x)_N`$ is zero because the barrier resistance $`R_b(d_N)`$ is infinite (the current cannot flow into the vacuum). Finally, we obtain Zaitsev (1990)
$$(ϵ+iϵ_bg_S(0))[\widehat{\tau }_3,\widehat{g}_N]+ϵ_{bN}if_S(0)[i\widehat{\tau }_2,\widehat{g}_N]=0.$$
(15)
where $`ϵ_{bN}=(D_N/2\gamma _Nd_N)`$ is a new characteristic energy that is determined by the $`S/N`$ interface resistance $`R_b`$. This equation looks similar to Eq.(11) after making the replacement $`\widehat{g}_S\widehat{g}_N.`$ The solution is similar to the solution (12)
$$g_N=\stackrel{~}{ϵ}/\stackrel{~}{\xi }_ϵ;\text{ }f_N=\stackrel{~}{ϵ}_{bN}/\stackrel{~}{\xi }_ϵ,$$
(16)
where $`\stackrel{~}{ϵ}=ϵ+iϵ_{bN}g_S(0),`$ $`\stackrel{~}{\xi }_ϵ=\sqrt{\stackrel{~}{ϵ}^2\stackrel{~}{ϵ}_{bN}^2},`$ $`\stackrel{~}{ϵ}_{bN}=ϵ_{bN}if_S(0).`$ Therefore the Green’s functions in the $`N`$ layer $`g_N`$ and $`f_N`$ are determined by the Green’s functions on the $`S`$ side of the $`S/N`$ interface $`g_S(0)`$ and $`f_S(0).`$ In order to find the values of $`g_S(0)`$ and $`f_S(0),`$ one has to solve Eq. (9) on the superconducting side ($`x<0`$). However, provided the inequality
$$\gamma _N/\gamma _S=\sigma _N/\sigma _S<<1$$
(17)
is fulfilled one can easily show that, in the main approximation, the solution in the $`S`$ region coincides with the solution for bulk superconductors (12). If the transparency of the $`S/N`$ interface is not high, $`ϵ_{bN}<<\mathrm{\Delta }`$, the characteristic energies $`ϵϵ_{bN}`$ are much smaller than $`\mathrm{\Delta }`$ and the functions $`g_S(0)`$ and $`f_S(0)`$ are equal to: $`g_S(0)g_{BCS}(0)ϵ/i\mathrm{\Delta },`$ $`f_S(0)f_{BCS}(0)1/i.`$ For these energies the functions $`g_N`$ and $`f_N`$ have the same form as the BCS functions $`g_{BCS}`$ and $`f_{BCS}`$ (12) with the replacement $`\mathrm{\Delta }ϵ_{bN}`$
$$g_N=\frac{ϵ}{\sqrt{ϵ^2ϵ_{bN}^2}},\text{ }f_N=\frac{ϵ_{bN}}{\sqrt{ϵ^2ϵ_{bN}^2}},$$
(18)
where $`ϵ_{bN}=D_N/(2R_b\sigma _Nd_N)`$. The energy $`ϵ_{bN}`$ can be represented in another form
$$ϵ_{bN}=\frac{\pi ^2}{2}(\frac{R_Q}{R_bk_F^2})\mathrm{}\frac{v_F}{d_N}=\mathrm{}\frac{v_F}{d_N}(\frac{T_b}{4}).$$
(19)
$`R_Q=\mathrm{}/e^2`$ is the resistance quantum, $`v_F`$ and $`k_F`$ are the Fermi velocity and wave vector. When obtaining the latter expression, we used a relation between the barrier resistance $`R_b`$ and an effective coefficient of transmission $`T_b`$ through the S/N interface Kuprianov and Lukichev (1988); Zaitsev (1984): $`R_b\sigma _n=(2/3)(l/T_b),`$ where $`l=v_F\tau `$ is the mean free path, $`T_b=<T(\theta )\mathrm{cos}\theta /(1T(\theta )>,`$ $`\theta `$ is the angle between the momentum of an incoming electron and the vector normal to the S/N interface, and $`T(\theta )`$ is the angle dependent transmission coefficient. The angle brackets mean an averaging over $`\theta `$.
An important result follows from Eq.(18): the DOS is zero at $`|ϵ|<ϵ_{bN}`$, i.e., $`ϵ_{bN}`$ is a minigap in the excitation spectrum McMillan (1968). Remarkably, in the considered limit $`ϵ_{bN}<<\mathrm{\Delta }`$ the value of $`ϵ_{bN}`$ does not depend on $`\mathrm{\Delta }`$, but is determined by the interface transparency or, in other words, by the interface resistance $`R_b`$. The appearance of the minigap is related to Andreev reflections Andreev (1964).
Eq. (19) for the minigap is valid if the inequalities (14) and $`ϵ_{bN}<\mathrm{\Delta }`$ are fulfilled. Both inequalities can be written as
$$(D_N/\mathrm{\Delta })/d_b<d_N<d_b$$
(20)
where $`d_b=2R_b\sigma _N`$ is a characteristic length. In the case of a small interface resistance $`R_b`$ or a large thickness of the N layer, that is, if the condition $`\sqrt{D_N/\Delta },d_{b\text{ }}<d_N`$ is fulfilled, the value of the minigap in the N layer is given by Golubov and Kupriyanov (1996)
$$ϵ_{bN}=c_1\frac{D_N}{d_N^2}$$
(21)
where $`c_1`$ is a factor of the order 1. This result has been obtained from a numerical solution of the Usadel equation. The DOS for the case of arbitrary thickness $`d_N`$ and interface transparency was calculated by Pilgram *et al.* (2000).
The situation changes in the clean limit. Let us consider, for example, a normal slab of a thickness $`d_N`$ in contact with an infinite superconductor. If the Thouless energy $`E_{Th}=v_F/d_N`$ is less than $`\mathrm{\Delta }`$, then discrete energy levels $`ϵ_n`$ appear Saint-James (1964) in the N region due to Andreev reflections Andreev (1964). As a result, the DOS has sharp peaks at $`ϵ=ϵ_n`$ (for a recent review see Deutscher (2005)). If $`E_{Th}`$ is much larger than $`\mathrm{\Delta },`$ the DOS $`\nu (ϵ)`$ is zero at $`ϵ=0`$ and increases with increasing the energy $`ϵ`$ (no gap). However, this is true only for such a simple geometry. For samples of more complicated shapes the behavior of the DOS $`\nu (ϵ)`$ depends on whether the electron dynamics in the N region is chaotic or integrable Melsen *et al.* (1996); Lodder and Yu.V.Nazarov (1998); Pilgram *et al.* (2000); Taras-Semchuk and Altland (2001); Beenakker (1997).
Finally, it has been shown by Altland *et al.* (2000) and Ostrovsky *et al.* (2001) that mesoscopic fluctuations smear out the singularity in the DOS at $`|ϵ|=ϵ_{bN}`$ and the DOS in the diffusive limit is finite, although small, for $`|ϵ|<ϵ_{bN}.`$ The minigap discussed above has been observed on a Nb/Si bilayer system and on a Pb/Ag granular system by Heslinga *et al.* (1994); Kouh and Valles (2003), respectively.
From this analysis we see that the proximity effect changes the DOS of the normal metal which acquires superconducting properties. In the next section we will focus our attention on the case that the normal metal is a ferromagnet. We will see that new interesting physics will arise from the mutual interaction of superconductivity and magnetism.
### II.2 Superconductor-ferromagnet structures with an uniform magnetization
In this section we consider the proximity effect between a superconductor $`S`$ and a ferromagnet $`F`$. We assume that the ferromagnet is a metal and has a conduction band. In addition, there is an exchange field due to spins of electrons of other bands.
As has been already mentioned, the effective exchange field acts on spins of the conduction electrons in the ferromagnet, and an additional term $`\widehat{H}_{ex}`$ describing this action appears in the total Hamiltonian (for more details see Appendix A)
$$\widehat{H}_{tot}=\widehat{H}+\widehat{H}_{ex}$$
(22)
$$\widehat{H}_{ex}=d^3𝐫\psi _\alpha ^+\left(𝐫\right)\left(𝐡\left(𝐫\right)\sigma _{\alpha \beta }\right)\psi \left(𝐫\right)𝑑𝐫$$
(23)
where $`\psi ^+\left(\psi \right)`$ are creation and destruction operators, $`𝐡`$ is the exchange field, $`\sigma _{\alpha \beta }`$ are Pauli matrices, and $`\alpha ,\beta `$ are spin indices. The Hamiltonian $`\widehat{H}`$ stands for a non-magnetic part of the Hamiltonian. It includes the kinetic energy, impurities, external potentials, etc. and is sufficient to describe all properties of the system in the absence of the exchange field.
The energy of the spin-up electrons differs from the energy of the spin-down electrons by the Zeeman energy $`2h`$. Due to the presence of the term $`\widehat{H}_{ex}`$ describing the exchange interaction all functions, including the condensate Green’s function $`f`$, are generally speaking non-trivial matrices in the spin space with non-zero diagonal and off-diagonal elements.
The situation is simpler if the direction of the exchange field does not depend on coordinates. In this case, choosing the $`z`$-axis along the direction of $`𝐡`$ one can consider electrons with spin “up” and “down” separately. In this Section we concentrate on this case. This can help the reader to understand several interesting effects and get an intuition about what one can expect from the presence of the exchange field. The results of this section will also help in understanding which effects in the superconductor-ferromagnet structures can be considered as rather usual and what kind of behavior is “exotic”. We will see that the exotic phenomena occur in cases when the exchange field is not homogeneous and therefore postpone their discussion until the next chapters.
If the exchange field $`h`$ is homogeneous the matrix $`\widehat{f}`$ describing the condensate $`\widehat{f}`$ is diagonal and can be represented in the form
$$\widehat{f}=f_3\widehat{\sigma }_3+f_0\widehat{\sigma }_0$$
(24)
where $`f_3`$ is the amplitude of the singlet component and $`f_0`$ is the amplitude of the triplet component with zero projection of the magnetic moment of Cooper pairs on the $`z`$ axis ($`S_z=0`$). Note that in the case of a $`S/N`$ structure the condensate function has a singlet structure only, i.e. it is proportional to $`\widehat{\sigma }_3`$. The presence of the exchange field leads to the appearance of the triplet term proportional to $`\widehat{\sigma }_0`$
The amplitudes of the singlet and triplet components are related to the correlation functions $`\psi _\alpha \psi _\beta `$ as follows Legget (1975); Vollhardt and Wölfle (1990)
$`f_3(t)`$ $``$ $`\psi _{}(t)\psi _{}(0)\psi _{}(t)\psi _{}(0),`$
$`f_0(t)`$ $``$ $`\psi _{}(t)\psi _{}(0)+\psi _{}(t)\psi _{}(0),`$ (25)
One can see that a permutation of spins does not change the function $`f_3(0),`$ whereas such a permutation leads to a change of the sign of $`f_0\left(0\right)`$. This means that the amplitude of the triplet component taken at equal times is zero in agreement with the Pauli exclusion principle. Later we will see that in the case of a non-homogeneous magnetization all triplet components including $`\psi _{}(t)\psi _{}(0)`$ and $`\psi _{}(t)\psi _{}(0)`$ differ from zero.
Once one determines the condensate function, Eq. (24), one is able to determine physical quantities ,as DOS, the critical temperature $`T_C`$, or the Josephson critical current through a S/F/S junction.
Next paragraphs are devoted to a discussion of these physical properties in $`F/S`$ systems with homogeneous magnetization.
#### II.2.1 Density of states (DOS)
In this section we discuss the difference between the DOS in $`S/N`$ and $`S/F`$ structures. General equations for the quasiclassical Green’s functions describing the system can be written but they are rather complicated (see Appendix A). In order to simplify the problem and, at the same time, give the basic idea about the effects it is sufficient to consider some limiting cases. This will be done in the present section leaving the general equations for the Appendix A.
In the case of a weak proximity effect, the condensate function $`\widehat{f}`$ is small outside the $`S`$ region. We consider again the diffusive limit. Then, the general Eq. (178) can be linearized and one obtains an equation for the matrix $`\widehat{f}`$ similar to Eq.(5) but containing an extra term due to the exchange field $`h`$
$$D_F^2\widehat{f}_F/x^2+2i(ϵ\widehat{\sigma }_0+h\widehat{\sigma }_3)\widehat{f}_F=0.$$
(26)
The subscript $`F`$ stands for the $`F`$ region.
In the absence of the exchange field $`h`$, Eq. (26) reduces to Eq. (5). It is important to emphasize that Eq. (26) is valid for a homogeneous $`h`$ only. Any variation of $`h`$ in space makes the equation much more complicated.
Eq. (26) should be complemented by boundary conditions which take the form (see Appendix A)
$$\gamma _F\widehat{f}_F/x=\widehat{f}_S$$
(27)
where $`\gamma _F=R_b\sigma _F`$, $`R_b`$ is the boundary resistance per unit area, $`\sigma _F`$ is the conductivity of the $`F`$ region, $`\widehat{f}_{F,S}`$ are the condensate matrix functions in the $`F`$ and $`S`$ regions. Since we assume a weak proximity effect, a deviation of the $`\widehat{f}_S`$ from its $`BCS`$ value $`\widehat{f}_{BCS}=\widehat{\sigma }_3f_{BCS}`$ is small. Therefore on the right-hand side of Eq.(27) one can write $`\widehat{f}_S\widehat{\sigma }_3f_{BCS},`$ where $`f_{BCS}`$ is defined in Eq.(9). At the ferromagnet/vacuum interface the boundary condition is given by the usual expression $`_x\widehat{f}_F=0`$, which follows from the condition $`R_b\mathrm{}`$.
Using Eq. (27), one can easily solve Eq. (26). We assume, as in the previous section, that the normal metal (ferromagnet) is in a contact with the superconductor at $`x=0`$ ($`x`$ is the coordinate perpendicular to the interface). The other boundary of the ferromagnet is located at $`x=d_F`$ and the space at $`x>d_F`$ is empty.
The proper solution for the diagonal matrix elements $`f_\pm f_{11(22)}`$ can be written as
$$f_\pm (x)=\{\begin{array}{cc}\pm \frac{f_{BCS}}{\kappa _{ϵ\pm }\gamma _F}\frac{\mathrm{cosh}(\kappa _{ϵ\pm }(xd_F))}{\mathrm{sinh}(\kappa _{ϵ\pm }d_F)}& 0<x<d_F\\ \text{ }0\text{ }& x>d_F\end{array}\text{ }.$$
(28)
Here $`\kappa _{ϵ\pm }=\sqrt{2i(ϵ\pm h)/D_F}`$ is a characteristic wave vector that determines the inverse penetration depth of the condensate functions $`f_{0,3}`$ into the ferromagnet.
Usually, the exchange energy $`h`$ is much larger than the energy $`ϵ`$ ($`ϵ\mathrm{max}\{\mathrm{\Delta },T\}`$). This means that the condensate penetration depth $`\xi _F=\sqrt{D_F/h}`$ is much shorter than the penetration depth into a normal (non-magnetic) metal $`\xi _N`$. The strong suppression of the condensate in the ferromagnet is caused by the exchange interaction that tries to align the spins of electrons parallel to the magnetization. This effect destroys the Cooper pairs with zero total magnetic moment.
It is worth mentioning that the condensate function $`f_\pm `$ experiences oscillations in space. Indeed, for a thick $`F`$ layer ($`d_F>>\xi _F`$) we obtain from Eq. (28).
$$f_\pm =\pm \frac{\mathrm{\Delta }}{E_ϵ\kappa _{F\pm }\gamma _F}\mathrm{exp}(x/\xi _F)[\mathrm{cos}(x/\xi _F)\pm i\mathrm{sin}(x/\xi _F)].$$
(29)
where $`E_ϵ=\sqrt{ϵ^2\mathrm{\Delta }^2},\kappa _{F\pm }=\kappa _{ϵ\pm }(ϵ)`$ at $`ϵ=0`$. The damped oscillations of $`f_\pm `$ lead to many interesting effects and, in particular, to a non-monotonic dependence of the critical temperature on the thickness $`d_F`$ of a $`F/S`$ bilayer which will be discussed in the next section.
In order to calculate the DOS we have to use the normalization condition, Eq. (7), which is also valid for the matrix elements $`f_\pm `$ and $`g_\pm `$. Thus, for $`g_\pm `$ we obtain $`g_\pm =\sqrt{1+f_\pm ^2}`$, which can be written for small $`f_\pm `$ as $`g_\pm 1+f_\pm ^2/2`$. Then the correction to the normalized DOS in the $`F`$ region $`\delta \nu _F=\nu _F1`$ takes the form
$$\delta \nu _F(x)=\mathrm{Re}(f_+^2+f_{}^2)/4.$$
(30)
Substituting Eq. (28) into Eq.(30), we obtain finally the DOS variation at the edge of the F film
$`\delta \nu _F(d_F)`$ $`=`$ $`(1/4)\mathrm{Re}\{({\displaystyle \frac{f_{BCS}}{\gamma _F}})^2[(\kappa _{ϵ+}\mathrm{sinh}(\kappa _{ϵ+}d_F))^2+`$ (31)
$`+(\kappa _ϵ\mathrm{sinh}(\kappa _ϵd_F))^2]\}.`$
In FIG. 3 we plot the function $`\delta \nu _F(ϵ)`$ for different thicknesses $`d_F`$ and $`h/\mathrm{\Delta }=20`$. It is seen that at zero energy $`ϵ=0`$ the correction to DOS $`\delta \nu _F`$ is positive for $`F`$ films with $`d_F=0.8\xi _0`$ while it is negative for films with $`d_F=0.5\xi _0`$ where $`\xi _0=\sqrt{D_F/\mathrm{\Delta }}`$.
Such a behavior of the DOS, which is typical for S/F systems, has been observed experimentally by Kontos *et al.* (2001) in a bilayer consisting of a thin PdNi film ($`5nm<d_F<7.5nm`$) on the top of a thick superconductor. The DOS was determined by tunnelling spectroscopy. This type of dependence of $`\delta \nu _N`$ on $`d_N`$ can also be obtained in $`N/S`$ contacts but for finite energies $`ϵ`$. In the $`F/S`$ contacts the energy $`ϵ`$ is shifted, $`ϵϵ\pm h`$ (time-reversal symmetry breaking) and this leads to a non-monotonic dependence of $`\nu _F`$ on the thickness $`d_F`$ even at zero energy. On the other hand, a non oscillatory behavior of the DOS $`\nu (ϵ)`$ has been found recently in experiments on Nb/CoFe bilayers Reymond *et al.* (2000). The discrepancy between the existing theory and the experimental data may be due to the small thicknesses of the ferromagnetic layer ($`0.5nm<d_F<2.5nm`$) which is comparable with the Fermi wave length $`\lambda _F0.3nm`$. Strictly speaking, in this case the Usadel equation cannot be applied.
The DOS in $`F/S`$ structures was studied theoretically in many papers. Halterman and Valls (2002b) studied the DOS variation numerically for ballistic $`F/S`$ structures. The DOS in quasiballistic $`F/S`$ structures was investigated by Baladie and Buzdin (2001), Bergeret *et al.* (2002b) and Zareyan *et al.* (2001) and for dirty $`F/S`$ structures by Fazio and Lucheroni (1999) and Buzdin (2000). The subgap in a dirty S/F/N structure was investigated in a recent publication by Golubov *et al.* (2005).
It is interesting to note that in the ballistic case ($`\tau h>>1`$, $`\tau `$ is the momentum relaxation time) the DOS in the $`F`$ layer is constant in the main approximation in the parameter $`1/(\tau h)`$ while in the diffusive case ($`\tau h<<1`$) it experiences the damped oscillations. The reason for the constant DOS in the ballistic case is that both the parts of $`f`$ , the symmetric and antisymmetric in the momentum space, contribute to the DOS. Each of them oscillates in space. However, while in the diffusive case the antisymmetric part is small, in the ballistic case the contributions of both parts to the DOS are equal to each other, but opposite in sign, thus compensating each other.
Finally, we would like to emphasize that both, the singlet and triplet components, contribute to the DOS. As it is seen from Eq.(30), the changes of the DOS can be represented in the form $`\delta \nu _F=\mathrm{Re}(f_0^2+f_3^2)/4`$, which demonstrates explicitly this fact.
#### II.2.2 Transition temperature
As we have seen previously, the exchange field affects greatly the singlet pairing in conventional superconductors. Therefore the critical temperature of the superconducting transition $`T_c`$ is considerably reduced in $`S/F`$ structures with a high interface transparency.
The critical temperature for $`S/F`$ bilayer and multilayered structures was calculated in many works Radovic *et al.* (1991); Demler *et al.* (1997); Khusainov and Proshin (1997); Tagirov (1998); Baladie and Buzdin (2003); Fominov *et al.* (2002); Bagrets *et al.* (2003); Fominov *et al.* (2003); Proshin and Khusainov (1998, 1999); Proshin *et al.* (2001); Buzdin and Kupriyanov (1991); Tollis *et al.* (2005); You *et al.* (2004); Baladie *et al.* (2001). Experimental studies of the $`T_c`$ were also reported in many publications Jiang *et al.* (1995); Mühge *et al.* (1998); Aarts *et al.* (1997); Lazar *et al.* (2000); Gu *et al.* (2002a). A good agreement between theory and experiment has been achieved in some cases (see FIG. 4).
One has to mention that, despite of many papers published on this subject, the problem of the transition temperature $`T_c`$ in the $`S/F`$ structures is not completely clear. For example, Jiang *et al.* (1995) and Ogrin *et al.* (2000) claimed that the non-monotonic dependence of $`T_c`$ on the thickness of the ferromagnet observed on $`Gd/Nb`$ samples was due to the oscillatory behavior of the condensate function in $`F`$. However, Aarts *et al.* (1997) in an other experiment on $`V/FeV`$ have shown that the interface transparency plays a crucial role in the interpretation of the experimental data that showed both non-monotonic and monotonic dependence of $`T_c`$ on $`(d_F)`$. In other experiments Bourgeois and Dynes (2002) the critical temperature of the bilayer Pb/Ni decreases with increasing $`d_F`$ in a monotonic way.
From the theoretical point of view the $`T_c`$ problem in a general case cannot be solved exactly. In most papers it is assumed that the transition to the superconducting state is of second order, i.e. the order parameter $`\mathrm{\Delta }`$ varies continuously from zero to a finite value with decreasing the temperature $`T`$. However, generally this is not so.
Let us consider, for example, a thin $`S/F`$ bilayer with thicknesses obeying the condition: $`d_F<\xi _F`$, $`d_S<\xi _S`$, where $`d_{F,S}`$ are the thicknesses of the $`F(S)`$ layer. In this case the Usadel equation can be averaged over the thickness (see for instance, Bergeret *et al.* (2001b)) and reduced to an equation describing an uniform magnetic superconductor with an effective exchange energy $`\stackrel{~}{h}`$ and order parameter $`\stackrel{~}{\mathrm{\Delta }}`$.
This problem can easily be solved. The Green’s functions $`g_\pm `$ and $`f_\pm `$ are given by
$$g_\pm =\frac{ϵ\pm \stackrel{~}{h}}{E_ϵ},\text{ }f_\pm =\frac{\stackrel{~}{\mathrm{\Delta }}}{E_ϵ},$$
(32)
where $`E_ϵ=\sqrt{(ϵ\pm \stackrel{~}{h})^2\stackrel{~}{\mathrm{\Delta }}^2},`$ $`\stackrel{~}{h}=r_Fh`$, $`\stackrel{~}{\mathrm{\Delta }}=r_S\mathrm{\Delta }`$, $`r_F=1r_S=\nu _Fd_F/(\nu _Fd_F+\nu _Sd_S)`$. In this case the Green’s functions are uniform in space and have the same form as in a magnetic superconductors or in a superconducting film in a parallel magnetic field acting on the spins of electrons.
The difference between the $`S/F`$ bilayer system and a magnetic superconductors is that the effective exchange energy $`\stackrel{~}{h}`$ depends on the thickness of the $`F`$ layer and may be significantly reduced in comparison with its value in a bulk ferromagnet. A thin superconducting film in a strong magnetic field $`H=\stackrel{~}{h}/\mu _B`$ ($`\mu _B`$ is an effective Bohr magneton) is described by the same Green’s functions. The behavior of these systems and, in particular, the critical temperature of the superconducting transition $`T_c`$, was studied long ago by Larkin and Ovchinikov (1964); Fulde and Ferrell (1965); Sarma (1963); Maki (1968). It was established that both first and second order phase transitions may occur in these systems if $`\stackrel{~}{h}`$ is less or of the order of $`\stackrel{~}{\mathrm{\Delta }}`$. If the effective exchange field $`\stackrel{~}{h}`$ exceeds the value $`\stackrel{~}{\mathrm{\Delta }}/\sqrt{2}0.707\stackrel{~}{\mathrm{\Delta }},`$ the system remains in the normal state (the Clogston (1962) and Chandrasekhar (1962) limit). Independently from each other Larkin and Ovchinikov (1964) and Fulde and Ferrell (1965) found that in a clean system and in a narrow interval of $`\stackrel{~}{h}`$ the homogeneous state is unstable and an inhomogeneous state with the order parameter varying in space is established in the system. This state, denoted as the Fulde-Ferrel-Larkin-Ovchinnikov (LOFF) state. has not been observed yet in bulk superconductors. In bilayered $`S/F`$ systems such a state cannot be realized because of a short mean free path.
In the case of a first order phase transition from the superconducting to the normal state the order parameter $`\mathrm{\Delta }`$ drops from a finite value to zero. The study of this transition requires the use of nonlinear equations for $`\mathrm{\Delta }`$. It was shown by Tollis (2004) that under some assumptions both the first and second order phase transitions may occur in a S/F/S structure.
In the case of a second order phase transition one can linearize the corresponding equations (the Eilenberger or Usadel equation) for the order parameter and use the Ginzburg-Landau expression for the free energy assuming that the temperature $`T`$ is close to the critical temperature $`T_c`$. Just this case was considered in most papers on this topic. The critical temperature of an $`S/F`$ structure can be found from an equation which is obtained from the self-consistency condition Eq. (4). In the Matsubara representation it has the form
$$\mathrm{ln}\frac{T_c}{T_c^{}}=(\pi T_c^{})\underset{\omega }{}(\frac{1}{|\omega _n|}if_\omega /\mathrm{\Delta }),$$
(33)
where $`T_c`$ is the critical temperature in the absence of the proximity effect and $`T_c^{}`$ is the critical temperature with taking into account the proximity effect.
The function $`f_\omega `$ is the condensate (Gor’kov) function in the superconductor; it is related to the function $`f_{S3}(ϵ)`$ as $`f_{S3}(i\omega _n)=f_\omega ,`$ where $`\omega _n=\pi (2n+1)`$ is the Matsubara frequency. Strictly speaking, Eq.(33) is valid for a superconducting film with a thickness smaller than the coherence length $`\xi _S`$ because in this case $`f_\omega `$ is almost constant in space.
The quasiclassical Green’s function $`f_\omega `$ obeys the Usadel equation (in the diffusive case) or the more general Eilenberger equation. One of these equations has to be solved by using the boundary conditions at the $`S/F`$ interface (or $`S/F`$ interfaces in case of multilayered structures). This problem was solved in different situations in many works where an oscillation of $`T_c`$ as a function of the F thickness was predicted (see FIG. 4). In most of these papers it was assumed that magnetization vectors $`𝐌`$ in different $`F`$ layers are collinear. Only Fominov *et al.* (2003) considered the case of an arbitrary angle $`\alpha `$ between the $`𝐌`$ vectors in two $`F`$ layers separated by a superconducting layer.
As mentioned previously, in this case the triplet components with all projections of the spin $`S`$ of the Coopers pair arise in the $`F/S/F`$ structure. It was shown that $`T_c`$ depends on $`\alpha `$ decreasing from a maximum value $`T_{c\mathrm{max}}`$ at $`\alpha =0`$ to a minimum value $`T_{c\mathrm{min}}`$ at $`\alpha =\pi .`$ We will not discuss the problem of $`T_c`$ for $`S/F`$ structures in detail because this problem is discussed in other review articles Izyumov *et al.* (2002); Buzdin (2005a).
#### II.2.3 The Josephson effect in SFS junctions
The oscillations of the condensate function in the ferromagnet (see Eq.(29)) lead to interesting peculiarities not only in the dependence $`T_c(d_F)`$ but also in the Josephson effect in the $`S/F/S`$ junctions. Although, as it has been mentioned in the previous section, the experimental results concerning the dependence $`T_c(d_F)`$ are still controversial, there is a more clear evidence for these oscillations in experiments on the Josephson current measurements that we will discuss here.
It turns out that under certain conditions the Josephson critical current $`I_c`$ changes its sign and becomes negative. In this case the energy of the Josephson coupling $`E_J=(\mathrm{}I_c/e)[1\mathrm{cos}\phi ]`$ has a minimum in the ground state when the phase difference $`\phi `$ is equal not to $`0`$, as in ordinary Josephson junctions, but to $`\pi `$ (the so called $`\pi `$junction).
This effect was predicted for the first time by Bulaevskii *et al.* (1977). The authors considered a Josephson junction consisting of two superconductors separated by a region containing magnetic impurities. The Josephson current through a $`S/F/S`$ junction was calculated for the first time by Buzdin *et al.* (1982). Different aspects of the Josephson effect in S/F/S structures has been studied in many subsequent papers (Buzdin and Kupriyanov, 1991; Radovic *et al.*, 2003; Fogelström, 2000; Barash *et al.*, 2002; Golubov *et al.*, 2002a; Heikkilä *et al.*, 2000; Chtchelkatchev *et al.*, 2001; Zyuzin *et al.*, 2003, e.g). Recent experiments confirmed the $`0`$-$`\pi `$ transition of the critical current in S/F/S junctions Ryazanov *et al.* (2001); Kontos *et al.* (2002); Blum *et al.* (2002); Bauer *et al.* (2004); Sellier *et al.* (2004).
In the experiments of Ryazanov *et al.* (2001) and Blum *et al.* (2002), $`Nb`$ is used as a superconductor and a $`Cu_xNi_{1x}`$ alloy as a ferromagnet. Kontos *et al.* (2002) used a more complicated $`S_1/F/I/S`$ structure, where $`S_1`$ is a $`Nb/Al`$ bilayer, $`S`$ is $`Nb`$, $`I`$ is the insulating $`Al_2O_3`$ layer and $`F`$ is a thin ($`40`$Å$`<d_F<150`$Å) magnetic layer of a $`PdNi`$ alloy. All these structures exhibit oscillations of the critical current $`I_c`$.
In FIG. 5 the temperature dependence of $`I_c`$ measured by Ryazanov *et al.* (2001) is shown. It is seen that the critical current in the junction with $`d_F=27nm`$ turns to zero at $`T2`$ $`K`$, rises again with increasing temperature and reaches a maximum at $`T5.5`$ $`K`$. If temperature increases further, $`I_c`$ decreases. In FIG. 6 we also show the dependence of $`I_c`$ on the thickness $`d_F`$ measured by Blum *et al.* (2002). The measured oscillatory dependence is well fitted with the theoretical dependence calculated by Buzdin *et al.* (1982) and Bergeret *et al.* (2001c).
The $`\pi `$-state in a Josephson junction leads to some observable phenomena. As was shown by Bulaevskii *et al.* (1977), a spontaneous supercurrent may arise in a superconducting loop with a ferromagnetic $`\pi `$-junction. This current has been measured by Bauer *et al.* (2004). Note also that the fractional Shapiro steps in a ferromagnetic $`\pi `$-junction were observed by Sellier *et al.* (2004) at temperatures at which the critical current $`I_c`$ turns to zero.
Oscillations of the Josephson critical current $`I_c`$ are related to the oscillatory behavior of the condensate function $`f`$ in space (see Eq.(29)). The critical current $`I_c`$ in a $`S/F/S`$ junction can easily be obtained once the condensate function in the $`F`$ region is known. We use the following formula for the superconducting current $`I_S`$ in the diffusive limit which follows in the equilibrium case from a general expression (see Appendix A)
$$I_S=L_yL_z\sigma _F(i\pi T/4e)\underset{\omega }{}Tr(\widehat{\tau }_3[\stackrel{ˇ}{f}_+\stackrel{ˇ}{f}_+/x+\stackrel{ˇ}{f}_{}\stackrel{ˇ}{f}_{}/x]),$$
(34)
where $`L_yL_z`$ is the area of the interface and $`\sigma _F`$ is the conductivity of the $`F`$ layer.
In the considered case of a non-zero phase difference the condensate functions $`f_\pm `$ are matrices in the particle-hole space. If in Eq.(34) instead of $`f_\pm `$ we write a 4$`\times `$4 matrix for $`\stackrel{ˇ}{f}`$, then $`\mathrm{\Delta }`$ is given by $`\widehat{\mathrm{\Delta }}=\mathrm{\Delta }(i\widehat{\tau }_2\mathrm{cos}(\phi /2)i\widehat{\tau }_1\mathrm{sin}(\phi /2))\widehat{\sigma }_3.`$ We set the phase of the right (left) superconductor equal to $`\pm \phi /2`$. For simplicity we assume that the overlap between the condensate functions $`f_\pm `$ induced in the $`F`$ region by each superconductor is small. This assumption is correct in the case $`d_F>>\xi _F`$. Under this assumption the condensate function may be written in the form of two independently induced $`f`$ functions
$`\widehat{f}_\pm (x)`$ $`=`$ $`(1/\xi _ϵ\kappa _{F\pm }\gamma _F)i\widehat{\tau }_2[\widehat{\mathrm{\Delta }}_l\mathrm{exp}(\kappa _{ϵ\pm }(x+d_F/2)`$ (35)
$`+`$ $`\widehat{\mathrm{\Delta }}_r\mathrm{exp}(\kappa _{ϵ\pm }(x+d_F/2)].`$
Here $`\widehat{\mathrm{\Delta }}_{r(l)}`$ is the order parameter in the right (left) superconductor. Substituting Eq.(35) into Eq.(34), we get
$`I_SI_c\mathrm{sin}(\phi )`$ $`=`$ $`4\pi T(L_yL_z)\sigma _F/(\kappa _F\gamma _F^2)\mathrm{exp}(d_F/\xi _F)`$ (36)
$`\mathrm{cos}(d_F/\xi _F){\displaystyle \underset{\omega }{}}{\displaystyle \frac{\mathrm{\Delta }^2}{\mathrm{\Delta }^2+\omega ^2}}\mathrm{sin}\phi .`$
When deriving Eq.(36), it was assumed that the exchange energy $`h`$ is much larger than both $`T`$ and $`\mathrm{\Delta }.`$
Calculating the sum in Eq. (36), we come to the final formula for the critical current
$$I_c=\mathrm{\Delta }\mathrm{tanh}(\mathrm{\Delta }/2T)\sigma _F/(\kappa _F\gamma _F^2)\mathrm{exp}(d_F/\xi _F)\mathrm{cos}(d_F/\xi _F).$$
(37)
As expected, according to Eq.(37) the critical current oscillates with varying the thickness of the ferromagnet $`d_F`$. The period of these oscillations gives the value of $`\xi _F`$ and therefore the value of the exchange energy $`h.`$ For example, according to the experiments on Nb/CuNi performed by Blum *et al.* (2002) $`h110meV`$, which is a quite reasonable value for CuNi.
The non-monotonic dependence of the critical current on temperature observed by Ryazanov *et al.* (2001) can be obtained only in the case of an exchange energy $`h`$ comparable with $`\mathrm{\Delta }`$ (at least, the ratio $`h/\mathrm{\Delta }`$ should not be too large). If the exchange energy were not too large, the effective penetration length $`\xi _{F,eff}`$ would be temperature dependent. According to estimates presented by Ryazanov *et al.* $`h30K`$, which means that the exchange energy in this experiment was much smaller than in the one performed by Blum *et al.* and by Kontos *et al.* (in the last reference $`h35meV`$).
The conditions under which the $`\pi `$ state is realized in $`S/F/S`$ Josephson junctions of different types were studied theoretically in many papers Buzdin and Kupriyanov (1991); Buzdin and Vujicic (1992); Chtchelkatchev *et al.* (2001); Krivoruchko and Koshina (2001a); Li *et al.* (2002); Buzdin and Baladie (2003). In all these papers it was assumed that the ferromagnet consisted of a single domain with a magnetization $`M`$ fixed in space. The case of a $`S/F/S`$ Josephson junction with a two domain ferromagnet was analyzed by Blanter and Hekking (2004). The Josephson critical current $`I_c`$ was calculated for parallel and anti-parallel magnetization orientations in both ballistic and diffusive limits. It turns out that in such a junction the current $`I_c`$ is larger for the anti-parallel orientation.
A similar effect arises in a $`S/F/S`$ junction with a rotating in space magnetization, as it was shown by Bergeret *et al.* (2001c). In this case not only the singlet and triplet component with projection $`S_z=0`$, but also the triplet component with $`S_z=\pm 1`$ arises in the ferromagnet. The last component penetrates the ferromagnet over a large length of the order of $`\xi _N`$ and contributes to the Josephson current.
In FIG. 7 the temperature dependence of the critical current is presented for different values of $`Ql,`$ where $`Q=2\pi /L_m`$, $`L_m`$ is the period of the spatial rotation of the magnetization and $`l`$ is the mean free path. It is seen that at $`Q=0`$ (homogeneous ferromagnet) and low temperatures $`T`$ the critical current $`I_c`$ is negative ($`\pi `$ state$`)`$, whereas with increasing temperature, $`I_c`$ becomes positive ($`0`$ state). If $`Q`$ increases, the interval of negative $`I_c`$ gets narrower and disappears completely at $`Ql0.04`$, that is, the $`S/F/S`$ structure with a non-homogeneous $`M`$ is an ordinary Josephson junction with a positive critical current.
It is interesting to note that the $`\pi `$-type Josephson coupling may also be realized in $`S/N/S`$ junctions provided the distribution function of quasiparticles $`n(ϵ)`$ in the $`N`$ region deviates significantly from the equilibrium. This deviation may be achieved with the aid of a non-equilibrium quasiparticles injection through an additional electrode in a multiterminal $`S/N/S`$ junction. The Josephson current in such a junction is again determined by Eq.(34) in which one has to put $`h=0,`$ $`f_+=f_{}`$ and replace $`\mathrm{tanh}(ϵ\beta )=(12n(ϵ))`$ by $`(1/2)[\mathrm{tanh}((ϵ+eV)\beta )+\mathrm{tanh}((ϵeV)\beta )],`$ where $`V`$ is a voltage difference between $`N`$ and $`S`$ electrodes.
At a certain value of $`V`$ the critical current changes sign. Thus, there is some analogy between the sign reversal effect in a $`S/F/S`$ junction and the one in a multiterminal $`S/N/S`$ junction under non-equilibrium conditions.
Indeed, when calculating $`I_S`$ in a multiterminal $`S/N/S`$ junction one can shift the energy by $`eV`$ or $`eV`$. Then the function $`(1/2)[\mathrm{tanh}((ϵ+eV)\beta )+\mathrm{tanh}((ϵeV)\beta )]`$ is transformed into $`\mathrm{tanh}(ϵ\beta )`$ while in the other functions one performs the substitution $`ϵϵ\pm eV`$. So, we see that $`eV`$ is analogous to the exchange energy $`h`$ that appears in the case of a $`S/F/S`$ junction.
The sign reversal effect in a multiterminal $`S/N/S`$ junction under non-equilibrium conditions has been observed by Baselmans *et al.* (1999) and studied theoretically by Volkov (1995); Yip (1998); Wilhelm *et al.* (1998). Later Heikkilä *et al.* (2000) studied theoretically a combined effect of a non-equilibrium quasiparticle distribution on the current $`I_c`$ in a $`S/F/S`$ Josephson junction.
Concluding this Section we note that the experimental results by Ryazanov *et al.* (2001), Kontos *et al.* (2002), Blum *et al.* (2002) and Strunk *et al.* (1994) seem to confirm the theoretical prediction of an oscillating condensate function in the ferromagnet and the possibility of switching between the 0 and the $`\pi `$-state.
## III Odd triplet superconductivity in S/F structures
### III.1 Conventional and unconventional superconductivity
Since the development of the BCS theory of superconductivity by Bardeen, Cooper, and Schrieffer (1957), and over many years only one type of superconductivity was observed in experiments. This type is characterized by the $`s`$-wave pairing between the electrons with opposite spin orientations due to the electron-phonon interaction. It can be called conventional since it is observed in most superconductors with critical temperature below $`20`$ $`K`$ (the so-called low temperature superconductors).
Bednorz and Müller (1986) discovered that a $`La_{2x}Sr_xCuO_4`$ compound is a superconductor with a critical temperature of $`30K`$. This was the first known high-$`T_c`$ copper-oxide (cuprate) superconductor. Nowadays many cuprates have been discovered with critical temperatures above the temperature of liquid nitrogen. These superconductors (the so called high $`T_c`$ superconductors) show in general a $`d`$-wave symmetry and, like the conventional superconductors, are in a singlet state. That is, the order parameter $`\mathrm{\Delta }_{\alpha \beta }`$ is represented in the form: $`\mathrm{\Delta }_{\alpha \beta }=\mathrm{\Delta }(i\sigma _3)_{\alpha \beta }`$, where $`\sigma _3`$ is the Pauli matrix in the spin space. The difference between the $`s`$ and $`d`$ pairing is due to a different dependence of the order parameter $`\mathrm{\Delta }`$ on the Fermi momentum $`𝐩_F=\mathrm{}𝐤_F`$. In isotropic conventional superconductors $`\mathrm{\Delta }`$ is a $`𝐤`$-(almost) independent quantity. In anisotropic conventional superconductors $`\mathrm{\Delta }`$ depends on the $`𝐤_𝐅`$ direction but it does not change sign as a function of the momentum $`𝐤_𝐅`$ orientation in space. In high $`T_c`$ superconductors where the $`d`$-wave pairing occurs, the order parameter $`\mathrm{\Delta }(𝐤_𝐅)`$ changes sign at certain points at the Fermi surface.
On the other hand, the Pauli principle requires the function $`\mathrm{\Delta }(𝐤_𝐅)`$ to be an even function of $`𝐤_𝐅`$, which imposes certain restrictions for the dependence of the order parameter on the Fermi-momentum. For example, for $`d`$-pairing the order parameter is given by $`\mathrm{\Delta }(𝐤_𝐅)=\mathrm{\Delta }(0)(k_x^2k_y^2)`$, where $`k_{x,y}`$ are the components of the $`𝐤_𝐅`$ vector in the $`CuO`$ plane. This means that the order parameter may have either positive or negative sign depending on the direction.
The change of the sign of the order parameter leads to different physical effects. For example, if a Josephson junction consists of two high $`T_c`$ superconductors with properly chosen crystallographic orientations, the ground state of the system may correspond to the phase difference $`\phi =\pi `$ ($`\pi `$junction). In some high $`T_c`$ superconductors the order parameter may consist of a mixture of $`s`$\- and $`d`$-wave components Tsuei and Kirtley (2003).
Another type of pairing, the spin-triplet superconductivity, has been discovered in materials with strong electronic correlations. The triplet superconductivity has been found in heavy fermion intermetallic compounds and also in organic materials (for a review see Mineev and Samokhin (1999)). Recently a lot of work has been carried out to study the superconducting properties of strontium ruthenate $`Sr_2RuO_4`$. Convincing experimental data have been obtained in favor of the triplet, $`p`$-wave superconductivity. For more details we refer the reader to the review articles by Maeno *et al.* (1994) and Eremin *et al.* (2004).
Due to the fact that the condensate function $`<\psi _\alpha (r,t)\psi _\beta (r^{},t^{})>`$ must be an odd function with respect to the permutations $`\alpha \beta ,`$ $`rr^{}`$ (for equal times, $`t=t^{}`$), the wave function of a triplet Cooper pair has to be an odd function of the orbital momentum, that is, the orbital angular momentum $`L`$ is an odd number: $`L=1`$ ($`p`$-wave), $`3`$ etc. Thus, the superconducting condensate is sensitive to the presence of impurities. Only the $`s`$-wave $`(L=0)`$ singlet condensate is not sensitive to the scattering by nonmagnetic impurities (Anderson theorem). In contrast, the $`p`$-wave condensate in an impure material is suppressed by impurities and therefore the order parameter $`\mathrm{\Delta }_{\alpha \beta }=_k\mathrm{\Delta }_{\alpha \beta }(𝐤_𝐅)`$ $`_k<\psi _\alpha (r,t)\psi _\beta (r^{},t)>_k`$ is also suppressed Larkin (1965). That is why the superconductivity in impure $`Sr_2RuO_4`$ samples has not been observed. In order to observe the triplet $`p`$-wave superconductivity (or another orbital order parameter with higher odd $`L`$), one needs to use clean samples of appropriate materials.
At first glance one cannot avoid this fact and there is no hope to see a non-conventional superconductivity in impure materials. However, another nontrivial possibility for the triplet pairing exists. The Pauli principle imposes restrictions on the correlation function $`<\psi _\alpha (r,t)\psi _\beta (r^{},t)>_k`$ for equal times. In the Matsubara representation this means that the sum $`_\omega <\psi _\alpha (r,\tau )\psi _\beta (r^{},\tau ^{})>_{k,\omega }`$ must change sign under the permutation $`rr^{}`$ (for the triplet pairing the diagonal matrix elements ($`\alpha =\beta `$) of these correlation functions are not zero). This implies that the sum $`_\omega <\psi _\alpha (r,\tau )\psi _\beta (r^{},\tau ^{})>_{k,\omega }`$ has to be either an odd function of $`k`$ or just turn to zero. The latter possibility does not mean that the pairing must vanish. It can remain finite if the average $`<\psi _\alpha (r,\tau )\psi _\beta (r^{},\tau ^{})>_{k,\omega }`$ is an odd function of the Matsubara frequency $`\omega `$ (in this case it must be an even function of $`k).`$ Then the sum over all frequencies is zero and therefore the Pauli principle for the equal-time correlation functions is not violated.
This type of pairing was first suggested by Berezinskii (1975) as a possible mechanism of superfluidity in $`{}_{}{}^{3}He`$. He assumed that the order parameter $`\mathrm{\Delta }(\omega )_{\omega ,k}<\psi _\alpha (r,\tau )\psi _\beta (r^{},\tau ^{})>_{k,\omega }`$ is an odd function of $`\omega :`$ $`\mathrm{\Delta }(\omega )=\mathrm{\Delta }(\omega )`$. However, experiments on superfluid $`{}_{}{}^{3}He`$ have shown that the Berezinskii state is only a hypothetical state and the $`p`$-pairing in $`{}_{}{}^{3}He`$ has different symmetries. As it is known nowadays, the condensate in $`{}_{}{}^{3}He`$ is antisymmetric in the momentum space and symmetric (triplet) in the spin space. Thus, the Berezinskii hypothetical pairing mechanism remained unrealized for few decades.
However, in recent theoretical works it was found that a superconducting state similar to the one suggested by Berezinskii might be induced in conventional $`S/F`$ systems due to the proximity effect Bergeret *et al.* (2001a); Bergeret *et al.* (2003). In the next sections we will analyze this new type of the superconductivity with the triplet pairing that is odd in the frequency and even in the momentum. This pairing is possible not only in the clean limit but also in samples with a high impurity concentration.
It is important to note that, in spite of the similarity, there is a difference between this new superconducting state in the $`S/F`$ structures and that proposed by Berezinskii. In the $`S/F`$ structures both the singlet and triplet types of the condensate $`f`$ coexist. However, the order parameter $`\mathrm{\Delta }`$ is not equal to zero only in the $`S`$ region (we assume that the superconducting coupling in the $`F`$ region is zero) and is determined there by the singlet part of the condensate only. This contrasts the Berezinskii state where the order parameter $`\mathrm{\Delta }`$ should contains a triplet component.
Note that attempts to find conditions for the existence of the odd superconductivity were undertaken in several papers in connections with the pairing mechanism in high $`T_c`$ superconductors Kirkpatrick and Belitz (1991); Kirckpatrick and Belitz (1992); Balatsky and Abrahams (1992); Coleman *et al.* (1993a, b); Coleman *et al.* (1995); Balatsky *et al.* (1995); Abrahams *et al.* (1993); Hashimoto (2000). In these papers a singlet pairing odd in frequency and in the momentum was considered.
We would like to emphasize that, while theories of unconventional superconductivity are often based on the presence of strong correlations where one has to use a phenomenology, the triplet state induced in the $`S/F`$ structures can be studied within the framework of the BCS theory, which is valid in the weak-coupling limit. This fact drastically simplifies the problem not only from the theoretical, but also from experimental point of view since well known superconductors grown in a controlled way may be used in order to detect the triplet component.
We summarize the properties of this new type of superconductivity which we speak of as triplet odd superconductivity:
* It contains the triplet component. In particular the components with projection $`S_z=\pm 1`$ on the direction of the field are insensitive to the presence of an exchange field and therefore long-range proximity effects arise in $`S/F`$ structures.
* In the dirty limit it has a $`s`$–wave symmetry. The condensate function is even in the momentum $`𝐩`$ and therefore, contrary to other unconventional superconductors, is not destroyed by the presence of non-magnetic impurities.
* The triplet condensate function is odd in frequency.
Before we turn to a quantitative analysis let us make the last remark: we assume that in the ferromagnetic regions no attractive electron-electron interaction exists, and therefore $`\mathrm{\Delta }=0`$ in the $`F`$-regions. The superconducting condensate arises in the ferromagnet only due to the proximity effect. This will become more clear later.
Another type of triplet superconductivity in the $`S/F`$ structures that differs from the one considered in this review was analyzed by Edelstein (2001). The author assumed that spin-orbit interaction takes place at the $`S/F`$ interface due to a strong electric field which exists over interatomic distances (the so-called Rashba term in the Hamiltonian Rashba (1960)). It was also assumed that electron-electron interaction is not zero not only in the $`s`$-wave singlet channel but also in the $`p`$-wave triplet channel. The spin-orbit interaction mixes both the triplet and singlet components. Then, the triplet component can penetrate into the $`F`$ region over a large distance.
However, in contrast to odd superconductivity, the triplet component analyzed by Edelstein is odd in the momentum and therefore must be destroyed by scattering on ordinary nonmagnetic impurities. This type of triplet component was also studied in $`2`$-dimensional systems and in $`S/N`$ structures in the presence of the Rashba-type spin-orbit interaction Edelstein (1989); Gor’kov and Rashba (2001); Edelstein (2001).
### III.2 Odd triplet component (homogeneous magnetization)
As we have mentioned in section II.2, even in the case of a homogeneous magnetization the triplet component with the zero projection $`S_z=0`$ of the total spin on the direction of the magnetic field appears in the $`S/F`$ structure. Unlike the singlet component it is an odd function of the Matsubara frequency $`\omega `$. In order to see this, we look for a solution of the Usadel equation in the Matsubara representation. In this representation the linearized Usadel equation for the ferromagnet takes the form
$$D_F^2\widehat{f}_F/x^22(|\omega |\widehat{\sigma }_0ih_\omega \widehat{\sigma }_3)\widehat{f}_F=0,$$
(38)
where $`\omega =\pi T(2n+1)`$ is the Matsubara frequency and $`h_\omega =\mathrm{sgn}(\omega )h`$.
The solution of Eq. (38) corresponding to Eq.(29) can be written as
$$f_\pm (\omega )=\pm (\mathrm{\Delta }/i\xi _\omega \kappa _\pm (\omega )\gamma _F)\mathrm{exp}(\kappa _\pm (\omega )x).$$
(39)
where
$$\kappa _\pm (\omega )=\sqrt{2(|\omega |ih_\omega )/D_F}$$
(40)
and $`\xi _\omega =\sqrt{\omega ^2+\mathrm{\Delta }^2}`$.
For the amplitudes of the triplet $`(f_0=(f_++f_{})/2)`$ and singlet $`(f_3=(f_+f_{})/2)`$ components we get in the ferromagnet
$$f_{3,0}(\omega ,x)=(\mathrm{\Delta }/2i\xi _\omega \gamma _F)[\frac{\mathrm{exp}(\kappa _+(\omega )x)}{\kappa _+(\omega )}\pm \frac{\mathrm{exp}(\kappa _{}(\omega )x)}{\kappa _{}(\omega )}].$$
(41)
Eqs. (39) and (41) show that both the singlet and the triplet component with $`S_z=0`$ of the condensate functions decay in the ferromagnet on the scale of $`\mathrm{Re}\kappa _+`$ having oscillations with $`\mathrm{Im}\kappa _+`$. Taking into account that $`\kappa _+(\omega )=\kappa _{}(\omega )`$, we see that $`f_3(\omega )`$ is an even function of $`\omega ,`$ whereas the amplitude of the triplet component, $`f_0(\omega )`$, is an odd function of $`\omega `$. The mixing between the triplet and singlet components is due to the term proportional to $`h_\omega \widehat{\sigma }_3`$ in Eq.(38). This term breaks the time-reversal symmetry.
Due to the proximity effect the triplet component $`f_0`$ penetrates also into the superconductor and the characteristic length of the decay is the coherence length $`\xi _S`$. The spatial dependence of this component inside the superconductor can be found provided the Usadel equation is linearized with respect to a deviation of the $`\widehat{f}_S`$ matrix from its bulk BCS form $`\widehat{f}_{BCS}`$. In the presence of an exchange field the Green’s functions $`\stackrel{ˇ}{g}`$ are $`4\times 4`$ matrices in the particle-hole and spin space. In the case of a homogeneous magnetization they can be represented as a sum of two terms (the $`\widehat{\tau }`$ matrices operate in the particle-hole space)
$$\stackrel{ˇ}{g}=\widehat{g}\widehat{\tau }_3+\widehat{f}i\widehat{\tau }_2,$$
(42)
where $`\widehat{g}`$ and $`\widehat{f}`$ are matrices in the spin space.
In a bulk superconductor these matrices are equal to
$$\widehat{g}_{BCS}=g_{BCS}(\omega )\widehat{\sigma }_0;\text{ }\widehat{f}_{BCS}=f_{BCS}(\omega )\widehat{\sigma }_3,$$
(43)
where
$$g_{BCS}(\omega )=\omega /\xi _\omega ;\text{ }f_{BCS}(\omega )=\mathrm{\Delta }/i\xi _\omega .$$
(44)
and $`\xi _\omega =\sqrt{\omega ^2+\mathrm{\Delta }^2}`$.
We linearize now the Usadel equation with respect to a small deviation $`\delta \stackrel{ˇ}{g}_S\delta \widehat{g}_S\widehat{\tau }_3+\delta \widehat{f}_Si\widehat{\tau }_2=\stackrel{ˇ}{g}_S\stackrel{ˇ}{g}_{BCS}`$ and obtain for the condensate function $`\delta \widehat{f}_S`$ in the superconductor the following equation
$$(^2/x^2)\delta \widehat{f}_S\kappa _S^2\delta \widehat{f}_S=2i(\delta \mathrm{\Delta }/D_S)g_{BCS}^2\widehat{\sigma }_3,$$
(45)
where $`\kappa _S^2=2\sqrt{(\omega ^2+\mathrm{\Delta }^2)}/D_S`$ and $`\delta \mathrm{\Delta }\left(x\right)`$ is a deviation of the superconducting order parameter from its BCS value in the bulk.
A solution for Eq.(45) determines the triplet component $`\delta f_{S0}`$ and a correction $`\delta f_{S3}`$ to the singlet component. To find the component $`\delta f_{S3}`$ is a much more difficult task than to find $`\delta f_{S0}`$ because $`\delta \mathrm{\Delta }(x)`$ is a function of $`x`$ and, in its turn, is determined by the amplitude $`\delta f_{S3}.`$ Therefore, the singlet component $`\delta f_{S3}`$ obeys a non-linear integro-differential equation. That is why the critical temperature $`T_c`$ can be calculated only approximately Buzdin and Kupriyanov (1990); Radovic *et al.* (1991); Demler *et al.* (1997); Izyumov *et al.* (2002); Tagirov (1998); Baladie and Buzdin (2003); Bagrets *et al.* (2003). Fominov *et al.* (2002) proposed an analytical trick that reduces the $`T_c`$ problem to a form allowing a simple numerical solution.
On the contrary, the triplet component $`\delta f_{S0}`$ proportional to $`\widehat{\sigma }_0`$ can be found exactly (in the linear approximation). The solution for $`\delta f_{S0}\left(0\right)`$ takes the form
$$\delta f_{S0}(x)=\delta f_{S0}(0)\mathrm{exp}(\kappa _S(\omega )x).$$
(46)
The constant $`\delta f_{S0}(0)`$ can be found from the boundary condition (see Appendix A)
$$\delta f_{S0}(x)/x|_{x=0}=f_{F0}(0)/\gamma _S.$$
(47)
As follows from this equation, the triplet component in the superconductor $`\delta f_{S0}`$ has the same symmetry as the component $`f_{F0},`$ that is, it is odd in frequency. So, the triplet component of the condensate is inevitably generated by the exchange field both in the ferromagnet and superconductor. Both the singlet component and the triplet component with $`S=0`$ decay fast in the ferromagnet because the exchange field $`h`$ is usually very large (see Eq. (40)). At the same time, the triplet component decays much slower in the superconductor because the inverse characteristic length of the decay $`k_S`$ is much smaller.
To illustrate some consequences of the presence of the triplet component in the superconductor, we use the fact that the normalization condition $`\stackrel{ˇ}{g}^2=1`$ results in the relation
$$g_0g_3=f_3f_0$$
(48)
The function $`g_0`$ entering Eq. (48) determines the change of the local DOS
$$\nu (ϵ)=\mathrm{Re}g_0(ϵ)$$
(49)
while the function $`g_3`$ determines the magnetic moment $`M_z`$ of the itinerant electrons (see Appendix A)
$$M_z=\mu _B\nu i\pi T\underset{\omega }{}g_3(\omega )$$
(50)
We see that the appearance of the triplet component in the superconductor leads to a finite magnetic moment in the $`S`$-region, which can be spoken of as an inverse proximity effect. This problem will be discussed in more detail in section V.2.
Thus, even in the case of a homogeneous magnetization, the triplet component with $`S_z=0`$ arises in the $`S/F`$ structure. This fact was overlooked in many papers and has been noticed for the first time by Bergeret *et al.* (2003). This component, as well as the singlet one, penetrates the ferromagnet over a short length $`\xi _F`$ because it consists of averages of two operators with opposite spins $`<\psi _{}\psi _{}>`$ and is strongly suppressed by the exchange field. The triplet component with projections $`S_z=\pm 1`$ on the direction of the field results in more interesting properties of the system since it is not suppressed by the exchange interaction. It can be generated by a non-homogeneous magnetization as we will discuss in the next section.
### III.3 Triplet odd superconductivity ( inhomogeneous magnetization)
According to the results of the last section the presence of an exchange field leads to the formation of the triplet component of the condensate function. In a homogeneous exchange field, only the component with the projection $`S_z=0`$ is induced.
A natural question arises: Can the other components with $`S_z=\pm 1`$ be induced? If they could, this would lead to a long range penetration of the superconducting correlations into the ferromagnet because these components correspond to the correlations of the type $`<\psi _{}\psi _{}>`$ with parallel spins and are not as sensitive to the exchange field as the other ones.
In what follows we analyze some examples of $`S/F`$ structures in which all projections of the triplet component are induced. The common feature of these structures is that the magnetization is nonhomogeneous.
In order to determine the structure of the condensate we will use as before the method of quasiclassical Green’s functions. This allows us to investigate all interesting phenomena except those that are related to quantum interference effects.
The method of the quasiclassical Green’s functions can be used at spatial scales much longer than the Fermi wave length <sup>3</sup><sup>3</sup>3Note that, as was shown by Galaktionov and Zaikin (2002); Shelankov and Ozana (2000), in the ballistic case and in the presence of several potential barriers some effects similar to the quantum interference effects may be important. We do not consider purely ballistic systems assuming that the impurity scattering is important. In this case the quasiclassical approach is applicable. The applicability of the quasiclassical approximation was discussed long ago by Larkin and Ovchinnikov Larkin and Ovchinnikov (1968).. As we have mentioned already, in order to describe the $`S/F`$ structures the Green’s functions have to be $`4\times 4`$ matrices in the particle-hole and spin space. Such $`4\times 4`$ matrix Green’s functions (not necessarily in the quasiclassical form) have been used long ago by Vaks *et al.* (1962); Maki (1969). Equations for the quasiclassical Green’s functions in the presence of the exchange field similar to the Eilenberger and Usadel equations can be derived in the same way as the one used in the non-magnetic case (see Appendix A). For example, a generalization of the Eilenberger equation was presented by Bergeret, Efetov, and Larkin (2000) and applied to the study of cryptoferromagnetism.
#### III.3.1 F/S/F trilayer structure
We start the analysis of the non-homogeneous case by considering the $`F/S/F`$ system shown in FIG. 8. The structure consists of one $`S`$ layer and two $`F`$ layers with magnetizations inclined at the angle $`\pm \alpha `$ with respect to the $`z`$-axis (in the $`yz`$ plane).
We want to demonstrate now that the triplet component with $`S_z=\pm 1`$ inevitably arises due to the overlap of the triplet components generated by the ferromagnetic layers in the $`S`$ layer. It is not difficult to understand why it should be so.
As we have seen in the previous section, each of the layers generates the triplet component with the zero total projection of the spin, $`S_z=0,`$ on the direction of the exchange field. If the magnetic moments of the layers are collinear to each other (parallel or antiparallel), the total projection remains zero. However, if the moments of the ferromagnetic layers are not collinear the superposition of the triplet components coming from the different layers should have all possible projections of the total spin.
From this argument we can expect the generation of the triplet component with all projections of the total spin provided the thickness of the $`S`$ layer is not too large. Since the only relevant length in the superconductors is $`\xi _S\sqrt{D_S/\pi T_c}`$, we assume that the thickness of the superconducting layer $`S`$ does not exceed this length.
Now we perform explicit calculations that support the qualitative conclusion about the generation of the triplet component with all projections of the total spin. We consider the diffusive case when the Usadel equation is applicable. This means that the condition
$$h\tau <<1$$
(51)
is assumed to be fulfilled ($`\tau `$ is the elastic scattering time).
The linearized Usadel equation in the $`F`$ region takes the form (see Appendix A)
$$_{xx}^2\stackrel{ˇ}{f}\kappa _\omega ^2\stackrel{ˇ}{f}+\frac{i\kappa _h^2}{2}\left\{\widehat{\tau }_0[\widehat{\sigma }_3,\stackrel{ˇ}{f}]_+\mathrm{cos}\alpha \pm \widehat{\tau }_3[\widehat{\sigma }_2,\stackrel{ˇ}{f}]\mathrm{sin}\alpha \right\}=0,$$
(52)
where $`\stackrel{ˇ}{f}`$ is a $`4\times 4`$ matrix (condensate function) which is assumed to be small and $`[\widehat{\sigma }_3,\stackrel{ˇ}{f}]_+=\widehat{\sigma }_3.\stackrel{ˇ}{f}+\stackrel{ˇ}{f}.\widehat{\sigma }_3`$. The wave vectors $`\kappa _\omega `$ and $`\kappa _h`$ entering Eq. (52) have the form
$$\kappa _\omega ^2=2|\omega |/D_F$$
(53)
and
$$\kappa _h^2=2h\mathrm{sgn}\omega /D_F$$
(54)
The magnetization vector $`𝐌`$ lies in the ($`y,z`$)-plane and has the components: $`𝐌=M\{0,\pm \mathrm{sin}\alpha ,\mathrm{cos}\alpha \}`$. The sign “$`+`$” (“$``$”) corresponds to the right (left) $`F`$ film. We consider here the simplest case of a highly transparent $`S/F`$ interface and temperatures close to the critical temperature of the superconducting transition $`T_c`$. In this case the function $`\stackrel{ˇ}{f}`$, being small, obeys a linear equation similar to Eq.(45)
$$(^2\stackrel{ˇ}{f}/x^2)\kappa _S^2\stackrel{ˇ}{f}=2i(\delta \stackrel{˘}{\mathrm{\Delta }}/D_S)g_{BCS}^2,$$
(55)
where $`\kappa _S^2=2|\omega |/D_S`$.
The boundary conditions at the $`S/F`$ interfaces are
$`\stackrel{ˇ}{f}_{x=d_S+0}`$ $`=`$ $`\stackrel{ˇ}{f}_{x=d_S0}`$ (56)
$`\gamma (\stackrel{ˇ}{f}/x)|_F`$ $`=`$ $`(\stackrel{ˇ}{f}/x)|_S.`$ (57)
where $`\gamma =\sigma _F/\sigma _S`$ and $`\sigma _F\left(\sigma _S\right)`$ is the conductivity in the ferromagnet (superconductor).
The first condition, Eq. (56), corresponds to the continuity of the condensate function at the $`S/F`$ interface with a high transparency, whereas Eq (57) ensures the continuity of the current across the $`S/F`$ interface Volkov *et al.* (2003).
A solution for Eqs.(52-55) with the boundary conditions (56-57) can easily be found. The matrix $`\stackrel{ˇ}{f}`$ can be represented as
$$\stackrel{ˇ}{f}=i\widehat{\tau }_2\widehat{f}_2+i\widehat{\tau }_1\widehat{f}_1,$$
(58)
where
$$\widehat{f}_1=b_1(x)\widehat{\sigma }_1,\text{ }\widehat{f}_2=b_3(x)\widehat{\sigma }_3+b_0(x)\widehat{\sigma }_0,$$
(59)
In the left $`F`$ layer the functions $`b_k(x)`$ are to be replaced by $`\overline{b}_k(x).`$ For simplicity we assume that the thickness of the $`F`$ films $`d_F`$ exceeds $`\xi _F`$ (the case of an arbitrary $`d_F`$ was analyzed by Bergeret *et al.* (2003)). Using the representation, Eqs. (58-59), we find the functions $`b_i(x)`$ and $`\overline{b}_i(x)`$. They are decaying exponential functions and can be written as
$$b_k(x)=b_k\mathrm{exp}(\kappa (xd_S)),\text{ }\overline{b}_k(x)\text{ }=\overline{b}_k\mathrm{exp}(\kappa (x+d_S))$$
(60)
Substituting Eq.(60) into Eqs.(52)-(55), we obtain a set of linear equations for the coefficients $`b_k.`$ The condition for the existence of non-trivial solutions yields an equation for the eigenvalues $`\kappa .`$ This equation reads
$$(\kappa ^2\kappa _\omega ^2)[(\kappa ^2\kappa _\omega ^2)^2+\kappa _h^4]=0$$
(61)
Eq. (61) is of the sixth order and therefore has six solutions. Three of these solutions should be discarded because the corresponding to $`b_k\left(x\right)`$ grow when going away from the interface. The remaining three solutions of Eq. (61) give three different physical values of $`\kappa `$.
If the exchange energy $`h`$ is sufficiently large ($`h>>\{T,\mathrm{\Delta }\}`$), the eigenvalues are
$`\kappa `$ $`=`$ $`\kappa _\omega `$ (62)
$`\kappa _\pm `$ $``$ $`(1\pm i)\kappa _h`$ (63)
We see that these solutions are completely different. The roots $`\kappa _\pm `$ proportional to $`\kappa _h`$ (cf. Eq. (54)), are very large and therefore the corresponding solutions $`b_k\left(s\right)`$ decay very fast (similar to the singlet component). This is the solution that exists for a homogeneous magnetization (collinear magnetization vectors).
In contrast, the value for $`\kappa `$ given by Eq. (62), is much smaller (see Eq. (53)) and corresponds to a slow decay of the superconducting correlations. The solutions corresponding to the root given by Eq.(62) describe a long-range penetration of the triplet component into the ferromagnetic region. For each root one can easily obtain relations between the coefficients $`b_k(x)`$. As a result, we obtain
$`b_1(x)`$ $`=`$ $`b_\omega e^{\kappa _\omega (xd_S)}`$ (64)
$`\mathrm{sin}\alpha \left[b_{3+}e^{\kappa _+(xd_S)}b_3e^{\kappa _{}(xd_S)}\right]`$
$`b_0(x)`$ $`=`$ $`\mathrm{tan}\alpha b_\omega e^{\kappa _\omega (xd_S)}`$ (65)
$`\mathrm{cos}\alpha \left[b_{3+}e^{\kappa _+(xd_S)}b_3e^{\kappa _{}(xd_S)}\right]`$
and
$$b_3(x)=b_{3+}\mathrm{exp}(\kappa _+(xd_S))+b_{3+}\mathrm{exp}(\kappa _{}(xd_S))$$
(66)
The function $`b_1(x)`$ is the amplitude of the triplet component penetrating into the $`F`$ region over a long distance of the order of $`\kappa _\omega ^1\xi _N.`$ Its value as well as the values of the other functions $`b_k(x)`$ is found from the boundary conditions (56-57) at the $`S/F`$ interfaces.
What remains to be done is to match the solutions for the superconductor and the ferromagnets at the interfaces between them. The solution for the superconductor satisfies Eq.(55) and can be written as
$`f_3(x)`$ $`=`$ $`\mathrm{\Delta }/iE_\omega +a_3\mathrm{cosh}(\kappa _Sx)`$ (67)
$`f_0(x)`$ $`=`$ $`a_0\mathrm{cosh}(\kappa _Sx)`$ (68)
$`f_1(x)`$ $`=`$ $`a_1\mathrm{sinh}(\kappa _Sx),`$ (69)
where $`E_\omega =\sqrt{\omega ^2+\mathrm{\Delta }^2}`$.
Matching these solutions with Eqs. (64-66) at the $`S/F`$ interfaces we obtain the coefficients $`b_k`$ and $`\overline{b}_k`$ as well as $`a_k`$. Note that $`b_{3\pm }=`$ $`\overline{b}_{3\pm }`$ and $`b_\omega =`$ $`\overline{b}_\omega .`$ Although the solution can be found for arbitrary parameters entering the equations, we present here for brevity the expressions for $`b_{3\pm }`$ and $`b_\omega `$ in some limiting cases only.
Let us consider first the case when the parameter $`\gamma \kappa _h/\kappa _S`$ is small, that is, we assume the condition
$$\gamma \kappa _h/\kappa _S\frac{\nu _F}{\nu _S}\sqrt{\frac{D_Fh}{D_S\pi T_c}}<<1$$
(70)
to be fulfilled.
Here $`\nu _{F,S}`$ is the density of states in the ferromagnet and superconductor, respectively (in the quasiclassical approximation the DOS for electrons with spin up and spin down is nearly the same: $`h<<ϵ_F`$). The condition, Eq. (70), can be fulfilled in the limit $`D_F<<D_S`$. Taking, for example, $`\nu _F\nu _S,`$ $`l_F30`$ Å and $`l_S300`$ Å, we find that $`h`$ should be smaller than $`30T_c`$.
In this limit the coefficients $`b_{1,3\pm }`$ and $`a_1`$ can be written in a rather simple form
$$b_\omega \frac{2\mathrm{\Delta }}{E_\omega }(\frac{\gamma \kappa _h}{\kappa _S})\frac{\mathrm{sin}\alpha \mathrm{cos}^2\alpha }{\mathrm{sinh}(2\mathrm{\Theta }_S)},$$
(71)
$$b_{3+}b_3\frac{\mathrm{\Delta }}{2i\xi _\omega },$$
(72)
$$a_3=\frac{\mathrm{\Delta }}{iE_\omega }\frac{\gamma \kappa _h}{\kappa _S}\frac{1}{\mathrm{sinh}(2\mathrm{\Theta }_S)},$$
(73)
where $`\mathrm{\Theta }=\kappa _Sd_S.`$
As follows from the first of these equations, Eq. (72), the correction to the bulk BCS solution for the singlet component is small in this approximation and this justifies our approach.
At the $`S/F`$ interface the amplitude of the triplet component $`b_\omega `$ is small in comparison with the magnitude of the singlet one $`b_{3+}`$ . However the triplet component decays over a long distance $`\xi _N`$ while the singlet one vanishes at distances exceeding the short length $`\xi _F`$. The amplitudes $`b_\omega `$ and $`b_{3\pm }`$ become comparable if the parameter $`\gamma \kappa _h/\kappa _S`$ is of the order of unity.
It follows also from Eq.(71) that the amplitude of the triplet component $`b_\omega `$ is zero in the case of collinear vectors of magnetization, i.e. at $`\alpha =0`$ or $`\alpha =\pi /2`$. It reaches the maximum at the angle $`\alpha _m`$ for which $`\mathrm{sin}\alpha _m=1/\sqrt{3}`$. Therefore the maximum angle-dependent factor in Eq.(71) is $`\mathrm{sin}\alpha _m\mathrm{cos}^2\alpha _m=2/3\sqrt{3}0.385.`$
One can see from Eq.(71) that $`b_\omega `$ becomes exponentially small if the thickness $`d_{S\text{ }}`$of the $`S`$ films significantly exceeds the coherence length $`\xi _S\sqrt{D_S/\pi T_c}`$. This means that in order to have a considerable penetration of the superconducting condensate into the ferromagnet one should not make the superconducting layer too thick.
On the other hand, if the thickness $`d_S`$ is too small the critical temperature $`T_c`$ is suppressed. In order to avoid this suppression one has to use, for instance, an $`F/S/F`$ structure with a small width of the $`F`$ films. Similar systems were considered by Beckmann *et al.* (2004), where non-local effects of Andreev reflections in a $`S/F`$ nanostructure were studied .
Another limiting case that allows a comparatively simple solution is the limit of small angles $`\alpha `$ Volkov *et al.* (2003) but an arbitrary parameter $`\gamma \kappa _h/\kappa _S`$, Eq. (70). At small angles $`\alpha `$ the amplitudes of the triplet and singlet component are given by the following formulae
$$b_\omega \frac{\mathrm{\Delta }}{E_\omega }\frac{\mathrm{sin}\alpha (\gamma \kappa _h/\kappa _S)\mathrm{tanh}\mathrm{\Theta }_S}{\mathrm{cosh}^2\mathrm{\Theta }_S|\mathrm{tanh}\mathrm{\Theta }_S+(\gamma \kappa _h/\kappa _S)|^2(1+(\gamma \kappa _h/\kappa _S)\mathrm{tanh}\mathrm{\Theta }_S)},$$
(74)
$$b_{3\pm }\frac{\mathrm{\Delta }}{2iE_\omega }\frac{1}{1+(\gamma \kappa _\pm /\kappa _S)\mathrm{tanh}\mathrm{\Theta }_S},$$
(75)
One can see from Eqs. (74-75) that, provided the parameter given by Eq. (70) is not small and $`\alpha ,|\mathrm{\Theta }_S|1`$, the amplitudes $`b_\omega `$ and $`b_{3\pm }`$ are again comparable with each other.
The amplitudes of the triplet and singlet components were calculated by Bergeret *et al.* (2003) in a more general case of an arbitrary $`S/F`$ interface transparency and a finite thickness of the $`F`$ films.
In FIG. 9 we plot the spatial dependence of the triplet (TC) and singlet (SC) components in a $`F/S/F`$ structure. It is seen from this figure that, as expected, the triplet component decays slowly, whereas the singlet component decays fast over the short length $`\xi _h`$. For this reason, in a multilayered $`S/F`$ structure with a varying direction of the magnetization vector $`𝐌`$ and thick $`F`$ layers ($`\xi _h<<d_F`$) a Josephson-like coupling between neighboring layers can be realized via the odd triplet component. In this case the in-plane superconductivity is caused by both triplet and singlet components. Properties of such S/F multilayered structures will be discussed in the next chapter.
Let us mention an important fact. The quasiclassical Green’s function $`\stackrel{ˇ}{g}(\vartheta )`$ in the diffusive case can be expanded in spherical harmonics. In the present approach, only the first two terms of this expansion are taken into account such that
$$\stackrel{ˇ}{g}=\stackrel{ˇ}{g}_{sym}+\stackrel{ˇ}{g}_{as}\mathrm{cos}\vartheta $$
(76)
where $`\vartheta `$ is the angle between the momentum $`p`$ and the $`x`$-axis, and $`\stackrel{ˇ}{g}_{as}=l\stackrel{ˇ}{g}_{sym}\stackrel{ˇ}{g}_{sym}/x`$ is the antisymmetric part of $`\stackrel{ˇ}{g}(\vartheta )`$ and $`\stackrel{ˇ}{g}_{sym}`$ is the isotropic part of $`\stackrel{ˇ}{g}(\vartheta )`$, which does not depend on $`\vartheta `$. The antisymmetric part of $`\stackrel{ˇ}{g}`$ determines the electric current in the system.
Higher order terms in the expansion of $`\stackrel{ˇ}{g}`$ are small in the diffusive limit and can be neglected. In the case of a weak proximity effect the antisymmetric part of the condensate function in the $`F`$ region can be written as
$$\stackrel{ˇ}{f}_{as}\mathrm{cos}\vartheta l\widehat{\tau }_3\widehat{\sigma }_0\mathrm{sgn}\omega \stackrel{ˇ}{f}_{sym}/x\mathrm{cos}\vartheta .$$
(77)
This expression follows from the fact that $`\stackrel{ˇ}{g}_0\widehat{\tau }_3\widehat{\sigma }_0\mathrm{sgn}\omega `$ (corrections to $`\stackrel{ˇ}{g}_0`$ are proportional to $`\stackrel{ˇ}{f}_0^2`$). Eq. (77) holds for both the singlet and triplet components.
As we have clarified previously, the symmetric part $`\stackrel{ˇ}{f}_0`$ is an odd function of $`\omega `$. Thus, according to Eq.(75) the antisymmetric part is an even function of $`\omega `$ so that the total condensate function $`\stackrel{ˇ}{f}=\stackrel{ˇ}{f}_0+\stackrel{ˇ}{f}_1\mathrm{cos}\vartheta `$ is neither odd nor even function of $`\omega `$. However, in the diffusive limit it is still legitimate to speak about the odd superconductivity since the symmetric part is much larger than antisymmetric part of $`\stackrel{ˇ}{f}`$.
If the parameter $`h\tau `$ is not small, i.e. the system is not diffusive, the symmetric and antisymmetric parts are comparable, and one cannot speak of the odd superconductivity. All this distinguishes the superconductivity in $`S/F`$ structures from the odd superconductivity suggested by Berezinskii (1975) who assumed that the order parameter $`\mathrm{\Delta }(\omega )`$ was an odd function of $`\omega `$. In our discussion it is assumed that the order parameter $`\mathrm{\Delta }`$ is an $`\omega `$ independent quantity and it is determined by the singlet component of the condensate function $`\stackrel{ˇ}{f}_0`$.
#### III.3.2 Domain wall at the S/F interface
In the previous section we have seen how the generation of the triplet component takes place. The appearance of this component leads to long range effects in a structure where the angle between the directions of magnetization in the different layers can be changed experimentally. This is an example of a situation when the long range triplet component of the superconducting condensate can be produced under artificial experimental conditions.
In this section we show that the conditions under which the triplet long range superconducting correlations occur are considerably more general. It is well known that the magnetization of any ferromagnet can be quite inhomogeneous due to the presence of domain walls. They are especially probable near interfaces between the ferromagnets and other materials. Therefore, making an interface between the ferromagnets and superconductors one produces almost inevitably domain walls, and one should take special care to get rid off them.
In this section we consider a domain wall like structure and show that it will also lead to the triplet long range correlations. This structure is shown schematically in FIG. 10. It consists of a $`S/F`$ bilayer with a non-homogeneous magnetization in the $`F`$ layer. We assume for simplicity that the magnetization vector $`𝐌=M(0,\mathrm{sin}\alpha (x),\mathrm{cos}\alpha (x))`$ rotates in the $`F`$ film starting from the $`S/F`$ interface ($`x=0`$) and the rotation angle has a simple piece-wise $`x`$-dependence
$$\alpha (x)=\{\begin{array}{c}Qx,\text{ }0<x<w\hfill \\ \alpha _w=Qw,\text{ }w<x\hfill \end{array}$$
(78)
This form means that the $`𝐌`$ vector is aligned parallel to the $`z`$-axis at the $`S/F`$ interface and rotates by the angle $`\alpha _w`$($`=\alpha (w)`$) over the length $`w`$ ($`w`$ may be the width of a domain wall). At $`x>w`$ the orientation of the vector $`𝐌`$ is fixed.
We calculate the condensate function in the $`F`$ region and show that it contains the long range triplet component (LRTC). As in the preceding section, we assume that the condensate function in the $`F`$ region is small. The smallness of $`\stackrel{ˇ}{f}`$ in this case is either due to a mismatch of the Fermi velocities in the superconductor and ferromagnet or due to a possible potential barrier at the $`S/F`$ interface. In such cases the transparency of the interface is small and only a small portion of the superconducting electrons penetrates the ferromagnet.
Due to the smallness of the transparency of the interface the function $`\stackrel{ˇ}{f}`$ can experience a jump at the $`S/F`$ interface, which contrasts the preceding case. The boundary condition for the $`4\times 4`$ matrix $`\stackrel{ˇ}{f}`$ has the same form as in Eq.(27)
$$\gamma _F_x\stackrel{ˇ}{f}=\stackrel{ˇ}{f}_S$$
(79)
The function $`\stackrel{ˇ}{f}_S`$ on the right-had side is the condensate matrix Green’s function in the superconductor that, in the limit considered here, should be close to the bulk solution
$$\stackrel{ˇ}{f}_S=f_{BCS}i\widehat{\tau }_2\widehat{\sigma }_3$$
(80)
We have to solve again Eq.(52) with the boundary conditions (79). Therefore we assume that the domain wall thickness $`w`$ is larger than the mean free path $`l`$and the condition, Eq. (51) is fulfilled (dirty limit). This case was analyzed by Bergeret *et al.* (2001b). Another case of a thin domain wall ($`w<l`$) was considered by Kadigrobov *et al.* (2001).
The problem of finding the condensate functions in the case of the magnetization varying continuously in space is more difficult than the previous one because the angle $`\alpha `$ depends now on $`x`$. However, one can use a trick that helps to solve the problem, namely, we exclude the dependence $`\alpha (x)`$ introducing a new matrix $`\stackrel{ˇ}{f}_n`$ related to $`\stackrel{ˇ}{f}`$ via an unitary transformation (a rotation in the particle-hole and spin-space)
$$\stackrel{ˇ}{f}=\stackrel{ˇ}{U}.\stackrel{ˇ}{f}_n.\stackrel{ˇ}{U}^+$$
(81)
where $`\stackrel{ˇ}{U}=\mathrm{exp}(i\widehat{\tau }_3\widehat{\sigma }_1\alpha (x)/2).`$
Performing this transformation we obtain instead of Eq.(52) a new equation
$`(_{xx}^2Q^2/2)\stackrel{ˇ}{f}_n`$ $``$ $`\kappa _\omega ^2\stackrel{ˇ}{f}_n+i\kappa _h^2[\widehat{\sigma }_3,\stackrel{ˇ}{f}_n]_+`$
$`{\displaystyle \frac{Q^2}{2}}(\widehat{\sigma }_1\stackrel{ˇ}{f}_n\widehat{\sigma }_1)`$ $`+`$ $`iQ\widehat{\tau }_3[\widehat{\sigma }_1,_x\stackrel{ˇ}{f}_n]_+=0`$ (82)
Correspondingly, the boundary condition, Eq. (79), takes the form
$$\gamma _F\{(Q/2)i\widehat{\tau }_3[\widehat{\sigma }_{1,}\stackrel{ˇ}{f}_n]_++\stackrel{ˇ}{f}_n/x\}=\stackrel{ˇ}{f}_s$$
(83)
Eq.(82) complemented by this boundary condition has to be solved in the region $`0<x<w`$. In the region $`w<x`$ one needs to solve Eq.(52) with $`Q=0`$. Both the solutions should be matched at $`x=w`$ under the assumption that there is no barrier at this point. Therefore, the matrix $`\stackrel{ˇ}{f}_n`$ and its ”generalized” derivative should be continuous at $`x=w`$
$`\stackrel{ˇ}{f}_n_{x=w0}`$ $`=`$ $`\stackrel{ˇ}{f}_n_{x=w+0}`$ (84)
$`{\displaystyle \frac{Q}{2}}i\widehat{\tau }_3[\widehat{\sigma }_{1,}\stackrel{ˇ}{f}_n]_+`$ $`+`$ $`_x\stackrel{ˇ}{f}_n_{x=w0}=_x\stackrel{ˇ}{f}_n_{x=w+0}`$ (85)
In this case the solution has the same structure as Eq.(58) but small changes should be done. The eigenvalues $`\kappa `$ obey the equation
$$[(\kappa ^2Q^2\kappa _\omega ^2)^2+4Q^2\kappa _\omega ^2](\kappa ^2\kappa _\omega ^2)+\kappa _h^4[\kappa ^2Q^2\kappa _\omega ^2]=0$$
(86)
where $`\kappa _{\omega ,h}^2`$ are determined in Eqs.(53,54). The eigenvalue given by Eq. (62) changes. Now it is equal to
$$\kappa _Q^2=Q^2+\kappa _\omega ^2,$$
(87)
while the eigenvalues $`\kappa _\pm ,`$ Eq. (63), remain unchanged provided the condition
$$Q,\kappa _\omega <<\kappa _h$$
(88)
is fulfilled.
In the opposite limit of large $`Q>>\kappa _h`$, the eigenvalues $`\kappa _\pm `$ take the form
$$\kappa _\pm =\pm iQ[1i\kappa _h^2/\sqrt{2}Q^2],$$
(89)
Thus, in this limit $`\kappa _\pm `$ is imaginary in the main approximation, which means that the function $`\stackrel{ˇ}{f}_n(x)`$ oscillates fast in space with the period $`2\pi /Q`$. In this case the eigenvalues (87) change also and have the form
$$\kappa ^2=\kappa _\omega ^2+\frac{\kappa _h^4}{Q^2}$$
(90)
Therefore the limit of a very fast rotating magnetization ($`\kappa _h/Q0`$) is analogous to the case of a normal metal, i.e. when the condensate penetrates the ferromagnet over the length $`\kappa _\omega ^1\sqrt{D_F/2\pi T}`$ which is the characteristic penetration length of the condensate in a S/N system.
More interesting and realistic is the opposite limit when the condition (88) is fulfilled and the long-range penetration of the triplet component into the ferromagnet becomes possible.
In the limit of large $`\kappa _h`$, (Eq. (88)), the singlet component penetrates the ferromagnet over a short length of the order $`\xi _h=1/\kappa _h`$ while the LRTC penetrates over the length $`1/\kappa _Q`$. As follows from Eq. (87), this penetration length is about $`1/Q`$ (provided $`w/\alpha _w`$ is smaller than the length $`\xi _N`$).
Now let us find the amplitude of the LRTC. The solution for Eq.(82) in the interval $`0<x<w`$ is determined by Eqs.(58, 59) with the functions $`b_i\left(x\right)`$, $`i=0,1,3`$ given by the following formulae
$$b_1(x)=b_Q\mathrm{exp}(\kappa _Qx)+\overline{b}_Q\mathrm{exp}(\kappa _Qx)$$
(91)
$$b_0(x)=b_{3+}\mathrm{exp}(\kappa _+x)+b_3\mathrm{exp}(\kappa _{}x)$$
(92)
and
$$b_3(x)=b_{3+}\mathrm{exp}(\kappa _+x)+b_3\mathrm{exp}(\kappa _{}x)$$
(93)
In the region $`w<x`$ the solution for the condensate function $`\stackrel{ˇ}{f}_n`$ takes the form
$$\stackrel{ˇ}{f}_n=i\widehat{\tau }_1\widehat{\sigma }_1c_\omega \mathrm{exp}(\kappa _\omega (xw))$$
(94)
where $`c_\omega `$ is a coefficient that has to be found by matching the solutions at $`x=w`$.
Terms of the order of $`Q/\kappa _h`$ are small and they are omitted now.
Then we find from the matching conditions at the $`S/F`$ interface, Eq. (85), the following relations for the coefficients
$$b_{3\pm }=\frac{f_{BCS}}{2\gamma _F\kappa _\pm }$$
(95)
and
$$b_Q=\overline{b}_Q=(Q/\kappa _Q)(b_{3+}b_3)$$
(96)
(the parameter $`\gamma _F`$ given by Eq. (79))
One can see from the above equations that the condensate function $`|\stackrel{ˇ}{f}|`$ is small provided parameter $`|\gamma _F\kappa _\pm |`$ is large. It follows from Eq.(96) that the amplitude of the LRTC, $`b_Q`$, is not zero only if the magnetization is nonhomogeneous, i.e., $`Q0.`$
Matching the solutions (91-94) at $`x=w`$, we find for the amplitude of the LRTC
$$c_\omega =\frac{if_{BCS}}{2\gamma _F}[\frac{Q}{\kappa _Q\mathrm{sinh}\alpha _w+\kappa _\omega \mathrm{cosh}\alpha _w}](\frac{h\mathrm{sgn}\omega /D_F}{|\kappa _+|^2\mathrm{Re}\kappa _+})$$
(97)
where $`\alpha _w=Qw`$ is the total angle of the magnetization rotation. As it has been mentioned, the amplitude of the LRTC is an odd function of $`\omega `$.
As one can see from the last expression the amplitude $`c_\omega `$ increases from zero when increasing $`Q`$, reaches a maximum at $`Q_{max}`$ corresponding a certain angle $`\alpha _{max}`$ and then exponentially decreases at $`\alpha _w>>`$ $`\alpha _{max}.`$
The maximum of $`c_\omega `$ is achieved at
$$\alpha _{max}=(w\kappa _\omega )\sqrt{\sqrt{5}1}/\sqrt{2}0.786(w\kappa _\omega ),$$
(98)
At $`\alpha _w=\alpha _{max}`$ the ratio in the square brackets in Eq.(97) is equal to $`0.68`$ . This means that the amplitude of the LRTC is of the order of the singlet component at the $`S/F`$ interface. The width $`w`$ should not be too small because in deriving the expression for $`c_Q`$ we assumed the condition $`w>>\xi _h.`$
In FIG. 11 we represented the dependence of $`\left|c_\omega \right|`$ on $`\alpha _w`$ for a fixed $`w`$. The spatial dependence of the LRTC and the singlet component is shown in FIG. 12. It is seen that for the parameters chosen the LRTC is larger than the singlet component and decays much slower with increasing the distance $`x`$.
If the magnetization vector $`𝐌`$ rotates by the angle $`\pi `$ (a domain wall) over a small length $`w`$ so that $`Q\pi /w>>\kappa _w`$, then the ratio in brackets in Eq.(97) is equal to
$$(\frac{Q}{\kappa _Q\mathrm{sinh}\alpha _w+\kappa _\omega \mathrm{cosh}\alpha _w})Q/(Q\mathrm{sinh}\pi )0.087$$
(99)
which shows that the amplitude of the LRTC in this case is smaller than the amplitude of the singlet component.
We can conclude from this analysis that in order to get a large LRTC, a small total angle of the rotation of the magnetization vector is more preferable.
The amplitude of the condensate function calculated here enters different physical quantities. In section III.4 we discuss how the long-range penetration of the triplet component into the ferromagnet affects transport properties of $`F/S`$ structures.
It is interesting to note that the type of magnetic structure discussed in this section differs drastically from the one in the case of an in-plane rotating magnetization. The latter was considered recently by Champel and Eschrig (2005a, b). It was assumed that the magnetization vector $`M_F`$ was parallel to the S/F interface and rotates; that is, it has the form $`M_F=M_0\{0,\mathrm{sin}(Qy),\mathrm{cos}(Qy)\}`$ (the x-axis is normal to the S/F interface plane). As shown by Champel and Eschrig (2005b), the odd triplet component arises also in this case but it penetrates into the ferromagnetic region over a short distance of the order of $`\xi _h`$.
#### III.3.3 Spin-active Interfaces
In almost all papers containing discussions of the $`S/F`$ structures it is assumed that the $`S/F`$ interface is spin-inactive, i.e. the spin of an electron does not change when the electron goes through the interface.
Although in many cases it is really so , one can imagine another situation when the spin of an electron passing through the interface changes. One can consider a region with a domain wall at the interface also as a “ spin-active interface” provided the width $`w`$ of the domain wall is very small but the product $`Qw`$ is of the order unity. As we have seen in section III.3.2, at such type of interfaces the triplet condensate arises.
Boundary conditions at spin-active $`S/F`$ interfaces for the quasiclassical Green’s functions were derived in a number of publications Millis *et al.* (1988); Kopu *et al.* (2004) and were used in studying different problems. Kulic and Endres (2000) employed these boundary conditions in the study of a system similar to the one shown in FIG. 8. Contrary to Bergeret *et al.* (2003), they assumed that the ferromagnets $`F`$ are insulators so that the condensate does not penetrate them. Nevertheless, the calculated critical temperature $`T_c`$ of the superconducting transition depends on the mutual orientation of the magnetization $`M_F`$ in the ferromagnets. In accordance with Baladie and Buzdin (2003); Tagirov (1998); Fominov *et al.* (2002) where metallic ferromagnets were considered in a $`F/S/F`$ structure, Kulic and Endres found that the critical temperature $`T_c`$ was maximal for the antiparallel magnetization orientation. If the directions of magnetization vector $`M_F`$ are perpendicular to each other, a triplet component also arises in the superconductor. The authors considered a clean case only, so that the influence of impurity scattering on the triplet component remained unclear.
According to Huertas-Hernando *et al.* (2002) a spin-active N/F interface plays an important role in the absolute spin-valve effect which can take place in a S/N/F mesoscopic structure. The authors considered a structure with a thin normal metal layer (N) and a ferromagnetic insulator F. The DOS variation in a conventional superconductor which is in contact with a ferromagnetic insulator was analyzed by Tokuyasu *et al.* (1988).
Eschrig *et al.* (2003) considered a clean $`S/F/S`$ Josephson junction in which the ferromagnet $`F`$ was a half metal so that the electrons with only one spin orientation (say the spin-up $``$ electrons) existed in the ferromagnet. In this case only the triplet component corresponding to the condensate function $`<\psi _{}\psi _{}>`$ may penetrate the ferromagnet. Assuming the p-wave triplet condensate function, the authors have calculated the critical Josephson current $`I_c`$ . They showed that the $`\pi `$ state (negative critical current $`I_c`$ ) is possible in this junction. The dc Josephson effect in a junction consisting of two superconductors and a spin-active interface between them was analyzed by Fogelström (2000).
It would be of interest to analyze the influence of impurities on the critical current in such type of Josephson junctions because, as we noted, in a clean case the singlet component can penetrate the ferromagnet (not a half metal) over a large distance.
### III.4 Long-range proximity effect
In the last decade transport properties of mesoscopic superconductor/normal metal $`S/N`$ structures were intensively studied (see for example the review articles by Beenakker (1997); Lambert and Raimondi (1998) and references therein). In the course of these studies many interesting phenomena have been discovered. Among them is a non-monotonic voltage and temperature dependence of the conductance in $`S/N`$ mesoscopic structures, i.e. structures whose dimensions are less than the phase coherence length $`L_\phi `$ and the inelastic scattering length $`l_\epsilon `$. This means that the resistance $`R`$ of a $`S/N`$ structure changes non-monotonically when the temperature decreases below the critical temperature $`T_c`$.
This complicated behavior is due to the fact that there are two contributions to the resistance in such systems: the one coming from the $`S/N`$ interface resistance and the resistance of the normal wire itself. The experimentally observed changes of the resistance can be both positive ($`\delta R>0`$) and negative ($`\delta R<0`$) Quirion *et al.* (2002); Shapira *et al.* (2000). The increase or decrease of the resistance $`R`$ depends, in particular, on the interface resistance $`R_{S/N}`$. If the latter is very small, the resistance of the $`S/N`$ structure is determined mainly by the resistance of the $`N`$ wire $`R_N`$. This resistance decreases with decreasing the temperature $`T`$, reaches a minimum at a temperature of the order of the Thouless energy $`D_N/L_N^2`$, and increases again returning to the value in the normal state $`R_N(T_c)`$ at low $`T`$, where $`D_N`$ is the diffusion coefficient and $`L_N`$ is the length of the $`N`$ wire. This is the so called re-entrant behavior observed in many experiments Gubankov and Margolin (1979); Pothier *et al.* (1994); Dimoulas *et al.* (1995); Petrashov *et al.* (1995); Charlat *et al.* (1996); Chien and Chandrasekhar (1999); Shapira *et al.* (2000).
Theoretical explanations for the non-monotonic behavior of the resistance variation as a function of the temperature $`T`$ or voltage $`V`$ in $`S/N`$ structures have been presented by Artemenko *et al.* (1979); Volkov *et al.* (1993); Nazarov and Stoof (1996); Volkov *et al.* (1996); Golubov *et al.* (1997); Shapira *et al.* (2000). Such a variation of the resistance of the normal metal wire can be explained in terms of the proximity effect that leads to the penetration of the condensate into the $`N`$ wire. Due to this penetration there are two types of contributions to the conductance $`G_N`$ Volkov and Pavlovskii (1996); Golubov *et al.* (1997). One of them reduces the DOS in the $`N`$ wire and therefore reduces the conductance $`G_N.`$ The other term, similar to the Maki-Thompson term Golubov *et al.* (1997); Volkov and Pavlovskii (1996), leads to an increase of the conductance of the $`N`$ wire.
In principle, the magnitude of the conductance variation $`\delta G_N`$ may be comparable with the conductance $`G_N`$. So, there are no doubts that the proximity effect plays a very important role in many experiments on $`S/N`$ structures.
Recently, similar investigations have been carried out also on mesoscopic $`F/S`$ structures in which ferromagnets ($`F`$) were used instead of normal (nonmagnetic) metals. According to our previous discussion, the depth of the condensate penetration into an impure ferromagnet equals $`\xi _F=\sqrt{\mathrm{}D/h}`$. This length is extremely short ($`550`$Å) for strong ferromagnets like $`Fe`$ or $`Ni`$. Therefore one might expect that the influence of the proximity effect on the transport properties of such structures should be negligibly small.
It was a great surprise that experiments carried out recently on $`F/S`$ structures showed that the resistance variation $`\delta R`$ were quite visible (varying from about $`1`$ to $`10\%`$) when decreasing the temperature below $`T_c`$ Lawrence and Giordano (1996a, b); Petrashov *et al.* (1999); Aumentado and Chandrasekhar (2001); Giroud *et al.* (1998). For example in the experiments by Lawrence and Giordano (1996a, b), where an $`Sn/Ni`$ structure was studied, the effective condensate penetration length estimated from the measured resistance was about $`400`$Å. This quantity exceeds $`\xi _F`$ by order of magnitude. Similar results have been obtained by Giroud *et al.* (1998) on $`Co/Al`$ structures, by Petrashov *et al.* (1999) on a $`Ni/Al`$ structures and by Aumentado and Chandrasekhar (2001) on $`Ni/Al`$ structures.
It is worth mentioning that the change of the resistance was both positive and negative. In some experiments the variation $`\delta R_F`$ was related to a change of the interface resistance Aumentado and Chandrasekhar (2001), whereas in others Lawrence and Giordano (1996a, b); Petrashov *et al.* (1999); Giroud *et al.* (1998) to the resistance variation of the ferromagnetic wire $`\delta R_F`$.
In FIG. 13 we show the temperature dependence of the resistance of a $`Ni`$ wire attached to an $`Al`$ bank measured by Petrashov *et al.* (1999). According to estimates of $`\xi _F`$ performed in this experiments, the observed $`\delta R_F`$ is by two orders of magnitude larger than it might be expected from the conventional theory of $`S/F`$ the contacts. Therefore these results cannot be explained in terms of the penetration of the singlet component.
In FIG. 14 we show similar data from the experiment on $`Co/Al`$ structures performed by Giroud *et al.* (1998). In this experiments a reentrance behavior of $`\delta R`$ was observed. In the limit of very low temperatures $`T0`$ the resistance was even larger than in the normal state.
The final explanation of this effect remains until now unclear. However, the long range proximity effects considered in the previous sections may definitely contribute to the conductance variation. In order to support this point of view we analyze qualitatively the changes of the conductance due to the LRTC penetration into the ferromagnet and demonstrate that the LRTC may lead to the conductance variation comparable with that observed in the experiments.
However, before presenting these calculations it is reasonable to understand if one can explain the experiments in a more simple way. Actually, the resistance of the $`S/F`$ structures has been analyzed in many theoretical works. For example, de Jong and Beenakker (1994); Golubov (1999); Belzig *et al.* (2000) analyzed a ballistic $`S/F`$ contact. It was shown that at zero exchange field ($`h=0`$), the contact conductance $`G_{F/S}`$ is twice as large as its conductance $`G_{F/N}`$ in the normal state (above $`T_c`$), as it should be. This agrees with a conductance in a $`N/S`$ ballistic contact according to theoretical predictions. At the same time, it drops to zero at $`h=E_F,`$ where $`E_F`$ is the Fermi energy.
The conductance of a diffusive point contact $`G_{F/S}`$ has been calculated by Golubov (1999) who showed that $`G_{F/S}`$ was always smaller than the conductance $`G_{F/N}`$ in the normal state. In the case of a mixed conductivity mechanism (partly diffusive and partly ballistic) the conductance $`G_{F/S}`$ has been calculated by Belzig *et al.* (2000). According to their calculations it can be both larger or smaller than the conductance in the normal state $`G_{F/N}`$.
The resistance $`R_F`$ of a ferromagnetic wire attached to a superconductor was calculated by Falko *et al.* (1999); Jedema *et al.* (1999); Bergeret *et al.* (2002a) and let us shortly describe what happens in such a system.
The proximity effect was neglected in these works but a difference in the conductivities $`\sigma _{}`$ for spin-up and down electrons was taken into account. The change of the conductance (or resistance) $`\delta G_F`$ is caused by a different form of the distribution functions below and above $`T_c`$ because of Andreev reflections.
The conductance $`G_F(T_c)`$ of the $`F`$ wire in the normal state ($`T>T_c`$) is given by the simple expression
$$G_F(T_c)=G_{}+G_{},$$
(100)
where $`G_{}=\sigma _{}L_FA,`$ $`L_F`$ and $`A`$ are the length and cross-section area of the $`F`$ wire.
This means that the total conductance is the sum of the conductances of the spin-up and down channels. In this case not only the electric current but also the spin current is not zero. It turns out that below $`T_c`$ ($`T<T_c`$) the conductance decreases and at zero temperature it is equal to
$$G_F(0)=4G_{}G_{}/(G_{}+G_{})$$
(101)
Eq. (101) shows that the zero-temperature conductance $`G_F(0)`$ for the system considered is smaller than the normal state conductance $`G_F(T_c).`$
It is possible to obtain the explicit formulae not only in the limiting cases, Eq. (100, 101), but also to describe the system at arbitrary temperatures. The general formula for the conductance of the $`F`$ wire can be written as
$$G_F(T)=G_F(0)\mathrm{tanh}(\mathrm{\Delta }/2T)+G_F(T_c)(1\mathrm{tanh}(\mathrm{\Delta }/2T))$$
(102)
Eqs. (100) and (101) can be obtained from Eq. (102) by putting $`\mathrm{\Delta }`$ or $`T`$ to zero. Eqs.(100-102) are valid provided the length $`L_F`$ satisfies the condition
$$l_{}<L_F<L_{SO},L_{in},$$
(103)
where $`l_{}`$ is the mean free path of spin-up and spin-down electrons,while $`L_{SO}`$ and $`L_{in}`$ are the spin-orbit and inelastic relaxation length, respectively.
The resistance of multiterminal $`S/F`$ structures was calculated by Mélin (2001); Mélin and Peysson (2003); Mélin and Feinberg (2004) on the basis of the tunnel Hamiltonian method. The influence of superconducting contacts on giant magnetoresistance in multilayered structures was studied by Taddei *et al.* (2001). Tkachov *et al.* (2002) studied an enhancement of Andreev reflection at the S/F interface due to inelastic magnon-assisted scattering.
One can conclude from the works listed above that neglecting the penetration of the LRTC into the $`F`$ wire an increase in the conductance $`G_F`$ cannot be explained. Therefore, let us discuss the consequences of the LRTC penetration into the ferromagnetic wire. In order to avoid the consideration of the $`S/F`$ interface contribution to the total resistance, we consider a cross geometry (see FIG. 15) and assume that the resistance of the interface between the $`F`$ wire and $`F`$ reservoirs is negligible. Such a geometry was used, for example, in the experiments by Petrashov *et al.* (1995). The structure under consideration consists of two $`F`$ wires attached to the $`F`$ and $`S`$ reservoirs. We assume that there is a significant mismatch between parameters of the superconductor and ferromagnet so that the condensate amplitude induced in $`F`$ is small and is determined by Eqs.(74) or (97).
According to our results obtained previously the long range proximity effect is possible provided there is a domain wall near the interface between the superconductor and ferromagnet and we assume that this is the case for the setup shown in FIG. 15. Another possibility to generate the triplet condensate would be to attach to the superconductor an additional ferromagnet with a non-collinear magnetization.
The conductance can be found on the basis of a general formula for the current (see for example the book by Kopnin (2001) and Appendix A)
$$I=(1/16e)(L_yL_z)\sigma _F\mathrm{Tr}\widehat{\sigma }_0\widehat{\tau }_3dϵ[\stackrel{ˇ}{g}^R_x\stackrel{ˇ}{g}^K+\stackrel{ˇ}{g}^K_x\stackrel{ˇ}{g}^A]$$
(104)
where $`\sigma _F`$ is the conductivity of the $`F`$ wire in the normal state.
The matrix Green’s function $`\stackrel{ˇ}{g}^K=\stackrel{ˇ}{g}^R\stackrel{ˇ}{F}\stackrel{ˇ}{F}`$ $`\stackrel{ˇ}{g}^A`$ is the Keldysh function related to a matrix distribution function $`\stackrel{ˇ}{F}`$. The distribution function consists of two parts, namely, one of them is symmetric with respect to the energy $`ϵ`$, the other one is antisymmetric in $`ϵ`$ and determines the dissipative current.
In the limit of a weak proximity effect the retarded (advanced) Green’s function has the form
$$\stackrel{ˇ}{g}^{R(A)}\pm \widehat{\tau }_3\widehat{\sigma }_0+\stackrel{ˇ}{f}^{R(A)},$$
(105)
where $`\stackrel{ˇ}{f}^{R(A)}`$ is given by Eqs.(74) or (97).
We have to find the conductance of the vertical $`F`$ wire in FIG. 15. In the main approximation the distribution function in this $`F`$ wire is equal to
$$\stackrel{ˇ}{F}=F_0\widehat{\tau }_0\widehat{\sigma }_0+F_3\widehat{\tau }_3\widehat{\sigma }_0,$$
(106)
where $`F_{0,3}=[\mathrm{tanh}((ϵ+V)/2T)\pm \mathrm{tanh}((ϵV)/2T)].`$
The distribution function $`F_3`$ symmetric in $`ϵ`$ determines the current $`I.`$ The differential conductance $`G_d=dI/dV`$ can be represented as
$$G_d=G_0+\delta G,$$
(107)
where $`G_0=\sigma _FL_FA`$ is the conductance in the normal state (here we neglect for simplicity the difference between $`\sigma _{}`$ and $`\sigma _{}`$).
The normalized correction to the conductance due to the proximity effect $`\delta S(T)\delta G/G_0`$ can be found using a general formula Bergeret *et al.* (2001a)
$$\delta S(T)=(32T)^1\mathrm{Tr}\widehat{\sigma }_0𝑑ϵ<(\widehat{f}^R\widehat{f}^A)^2>F_V^{}(ϵ)$$
(108)
where
$`F_V^{}(ϵ)=[`$ $`\mathrm{cosh}^2((ϵ+eV)/2T)+\mathrm{cosh}^2((ϵeV)/2T)]/2`$.
The angle brackets $`<\mathrm{}>`$ denote the average over the length of the ferromagnetic wire between the F (or N) reservoirs. The functions $`\widehat{f}^{R(A)}`$ are given by expressions similar to Eq.(97). This formula shows that if $`T<D_F/L^2`$, on the order of magnitude $`\delta S(T)|f_{tr}|^2`$,where $`L`$ is the length of the ferromagnetic wire and $`|f_{tr}|`$ is the amplitude of the triplet component at the S/F interface at a characteristic energy $`ϵ_{ch}\mathrm{min}\{T,D_F/L\}.`$ According to Eq.(97) the amplitude of the triplet component is of the order of $`c_1(\rho \xi _h/R_b)`$ where $`\rho `$ is the resistivity of the ferromagnet and $`c_1`$ is determined by the factor in the square brackets, that is, by the characteristics of the domain wall. In principle the amplitude $`|f_{tr}|`$ may be of the order of 1.
Strictly speaking, both the singlet and triplet components contribute to the conductance. However if the length $`L_F`$ much exceeds the short length $`\xi _F`$ only the contribution of the LRTC is essential.
In FIG. 16 we present the temperature dependence of the correction to the conductance $`\delta G(T).`$ It is seen that with increasing the temperature $`\delta G_F(T)`$ decreases in a monotonous way. This dependence differs from the re-entrant behavior discussed above that occurs in the $`S/N`$ structures. The reason for this difference is that the time-reversal symmetry in $`S/F`$ structures is broken and this leads to a difference in transport properties. In a $`S/N`$ system, a relation$`\widehat{f}^R(ϵ)=\widehat{f}^A(ϵ)|_{ϵ=0}`$ holds and this equality is a consequence the time-reversal symmetry. That is why $`\delta G(0)=\delta G(T_c)=0`$ in $`S/N`$ structures, whereas in a $`S/F`$ structure $`\widehat{f}^R(ϵ)\widehat{f}^A(ϵ)|_{ϵ=0}`$ and that is why $`\delta S(T)_{T=0}0.`$
Although the LRTC may be the reason for the enhancement of the conductivity in the $`S/F`$ structures (this possibility was also pointed out in the work by Giroud *et al.* (2003)), our understanding is based on the assumption that the magnetic moment is fixed and does not change with the temperature. Dubonos *et al.* (2002) suggested another mechanism based on an assumption about a domain redistribution when the temperature drops below $`T_c.`$ The ferromagnetic wires (or strips) used in different experiments may consist of many domains. Their form and number depend on the sample geometry and parameters of the system. When the temperature decreases below $`T_c`$, stray magnetic fields excite the Meissner currents in the superconductor attached to the $`F`$ wire. Therefore the demagnetizing factors change, which may lead to a new domain structure. At the same time, the total conductance (or resistance) $`G_F`$ depends on the form and the number of domains. So, one might expect that the conductance $`G_F(T)`$ below $`T_c`$ would differ from $`G_F`$ in the normal state. This idea was supported by measurements carried out by Dubonos *et al.* (2002). In this work a structure consisting of a two-dimensional electron gas and five Hall probes was used. An $`F/S`$ system ($`Ni`$+$`Al`$ disks) was placed on top of this structure. Measuring the Hall voltage, the authors were able to probe local magnetic fields around the ferromagnetic disks. They found that these fields really changed when the temperature dropped below $`T_c`$.
On the other hand, the Meissner currents and, hence, the effect of the redistribution of the domain walls may be considerably reduced in wires, as discussed previously. Changing the thickness of the superconducting wires in a controlled way and measuring the conductance could help to distinguish experimentally between the contribution to the conductivity of the triplet condensate and the effects of the redistribution of the domain walls.
An experiment in which the domain redistribution was excluded has been performed by Nugent *et al.* (2004). The authors measured the resistance variation of a ferromagnetic wire (Ni<sub>1-x</sub>Cu<sub>x</sub>) lowering the temperature $`T`$ below the critical temperature $`T_c`$ of the superconductor (Al or Pb), which was attached to the middle part of the ferromagnetic wire. A magnetic field, strong enough to align all domains in the ferromagnet in one direction but not too strong to suppress the superconductivity, was applied to the system. Under these conditions a small increase in the resistance ($`\delta R/R310^3`$) was observed when the temperature $`T`$ drops below $`T_c`$. The analysis presented above shows that the triplet component leads to an increase of the conductance but not in the resistance of the ferromagnetic wire. Therefore this particular experiment can hardly be explained in terms of the long-range proximity effect. Perhaps the small increase in the resistance of the ferromagnetic wire observed in Nugent *et al.* (2004) was related to the ”kinetic” mechanism discussed above (see Eq. (102)) or to weak localization corrections caused by the triplet Cooperons McCann *et al.* (2000). According to McCann *et al.* (2000) the change of the resistance of the ferromagnetic wire is positive (contrary to the contribution of the LRTC) and its order of magnitude is $`(e^2/\mathrm{})R_F`$, where $`R_F`$ is the resistance of the F wire in the normal state. In order to clarify the role of the LRTC in the transport properties of S/F structures, further theoretical and experimental investigations are needed. Note that strong ferromagnets like Fe are not suitable materials for observing the contribution of the LRTC into the conductance variation because of the strong exchange field $`h`$. In this case, according to Eq.(71) and Eq.(97), the amplitude of the LRTC is small because it contains $`h`$ in the denominator.
## IV Josephson effect in S/F systems (inhomogeneous magnetization)
As we have mentioned above, one of the most interesting issues in the $`S/F`$ structures is the possibility of switching between the so called 0- and $`\pi `$-states in Josephson $`S/F/S`$ junctions. The $`\pi `$-state denotes the state for which the Josephson critical current $`I_c`$ becomes negative. This occurs for a certain thickness $`d_F`$ and temperature $`T`$. In this state the minimum of the Josephson coupling energy $`E_J=(\mathrm{}I_c/e)(1\mathrm{cos}\varphi )`$ corresponds to a phase difference of $`\varphi =\pi `$ but not to $`\varphi =0`$ as in conventional Josephson junctions.
The reason for the sign reversal of $`I_c`$ is the oscillatory dependence of the condensate functions $`\widehat{f}`$ on the thickness $`d_F`$ (see Eq.(37)). Since the critical current $`I_c`$ is sensitive to the phase of the condensate function at the boundary, the $`\pi `$-state is a rather natural consequence of the oscillations.
The possibility of the $`\pi `$ state was predicted by Bulaevskii *et al.* (1977) and Buzdin *et al.* (1982), and studied later in many other works (e.g. Radovic *et al.*, 1991; Buzdin and Vujicic, 1992). Experimentally, this phenomenon manifests in a non-monotonic dependence of the critical temperature on the thickness of the junction observed in many experiments and discussed in Section II.2.2. Another manifestation of the transition from the $`0`$\- state to the $`\pi `$\- state is the sign reversal of the critical current observed in the experiment by Ryazanov *et al.* (2001) on $`Nb/Cu_xNi_{1x}/Nb`$ Josephson junctions (see FIG. 5). The proper choice of an alloy with a weak ferromagnetic coupling was crucial for the observation of the effect.
Subsequent experiments, Kontos *et al.* (2002), Blum *et al.* (2002) and Guichard *et al.* (2003), corroborated the observed change of the sign of the Josephson coupling varying the thickness of the intermediate $`F`$ layer. Qualitatively, the experimental data on the Josephson effect in the $`S/F/S`$ structures are in agreement with the theoretical works above mentioned. However, a more accurate control and understanding of the $`0`$-$`\pi `$ transition demands knowledge of the magnetic structure of the ferromagnetic materials.
Almost in all theoretical papers very simplified models of the $`S/F/S`$ junction were analyzed. For example, Blanter and Hekking (2004) assumed that the $`F`$ layer consisted either of one domain or two domains with the collinear orientation of the magnetization. In this case and according to the discussion of section III.3 the LRTC is absent in the system.
If the $`F`$ layer is a single domain layer, the critical current $`I_c`$ is maximal at a non-zero external magnetic field $`H_{ext}`$ equal to $`4\pi M_F`$, where $`M_F`$ is the magnetization of the $`F`$ layer. At the same time, in experiments Ryazanov *et al.* (2001); Kontos *et al.* (2001, 2002); Blum *et al.* (2002); Strunk *et al.* (1994); Sellier *et al.* (2004) a decrease of the current $`I_c`$ with increasing field $`H_{ext}`$ was observed and it was maximal at $`H_{ext}=0`$. This means, as it was assumed in these experimental works, that the F layer in real junctions contains many magnetic domains. In this case the Josephson critical current $`I_c`$ may change sign in the $`S/F/S`$ junctions with a multidomain magnetic structure even if the local Josephson current density $`j_c`$ is always positive. The reason for the sign reversal of $`I_c`$ in this case is a spatial modulation of the phase difference $`\varphi (x)`$ due to an alternating magnetization $`M\left(x\right)`$ in the domains Volkov and Anishchanka (2004). In order to determine the mechanism that leads to the sign reversal of the critical current further experiments are needed.
In this chapter we discuss a new phenomenon, namely, how the Josephson coupling between the $`F`$ layers in the $`S/F`$ structures is affected by the LRTC.
First, we consider a planar $`S/F/S`$ Josephson junction with a ferromagnet magnetization $`𝐌_𝐅`$ rotating in the direction normal to the junction plane. This model is an idealization of a real multidomain structure with different magnetization orientations. In this case, as we discussed in Section III.4, the LRTC arising in the structure affects strongly the critical current $`I_c.`$
Next, we will analyze a multilayered $`S/F/S/`$… structure in which the vector $`𝐌_F`$ has a different direction in different $`F`$ layers. Again, in this case the LRTC arises in the system. Interestingly, if the thickness of the $`F`$ layers $`d_F`$ is much larger than the penetration length $`\xi _F`$ of the singlet component but less or of the order of $`\xi _N`$, then the Josephson coupling between the $`F`$ layers is realized due to the LRTC (odd triplet superconductivity in the transverse direction). At the same time, the in-plane superconductivity is due to the conventional singlet superconducting pairing.
Finally we will discuss the dc Josephson effect in a $`SF/I/FS`$ junction (here $`SF`$ is a superconductor-ferromagnet bilayer and $`I`$ is a thin insulating layer). In this structure, the exchange field may lead not only to a suppression of the Josephson coupling as one could naively expect but, under a certain condition, to its enhancement.
Let us consider first a planar $`S/F/S`$ Josephson junction. We assume the following spacial dependence of the magnetization vector in the F layer: $`𝐌_F=M_F(0,\mathrm{sin}(Qx),\mathrm{cos}(Qx))`$, where the $`x`$-axis is normal to the junction plane.
In this case, as we have seen in Section III.3.2, the LRTC arises. Due to the long range penetration into the ferromagnet the triplet component can give a very important contribution to the Josephson current. A general expression for the Josephson current can be written in the form
$$I_J=(L_yL_z/4e)\sigma _F(\pi T)\mathrm{Tr}(\widehat{\sigma }_0\widehat{\tau }_3.\underset{\omega }{}\stackrel{ˇ}{f}_\omega _x\stackrel{ˇ}{f}_\omega ).$$
(109)
We assume that the impurity concentration is sufficiently high and therefore the condensate function $`\stackrel{ˇ}{f}_\omega `$ should be found from the Usadel equation. In the limit of a weak proximity effect (the $`S/F`$ interface transparency is not too high) this equation can be linearized and solved exactly. The solution for the $`\stackrel{ˇ}{f}_\omega `$ matrix in the $`F`$ region can be found in a similar way as it was done in Section III.3.2. Due to the rotation of the magnetization the condensate function contains the LRTC. We obtain for the Josephson current Bergeret *et al.* (2001c) the following expression
$$I_J=I_c\mathrm{sin}\varphi $$
(110)
where the critical current $`I_c`$ is equal to
$$I_c=(L_yL_z\sigma _F/l)\stackrel{~}{\gamma }_F^2\mathrm{Re}\underset{\omega >0}{}f_s^2\left[\frac{e^{\kappa _+d_F}}{\kappa _+l}+\frac{(Ql)^2e^{\kappa _ld_F}}{2(3h\tau )^{3/2}}\right],$$
(111)
and $`\kappa _l^2=2|\omega _n|/\mathrm{\Delta }+Q^2`$. The parameter $`\stackrel{~}{\gamma }_F=(3/4)<\mu T(\mu )>`$ is an effective, averaged over angles, transmittance coefficient which characterizes the S/F interface transparency and $`\kappa _+`$ is defined in Eq. (89)
The first term in the brackets containing the parameter $`\kappa _+`$corresponds to Eq. (36). It decays by increasing the thickness $`d_F`$ over the short characteristic length $`\xi _F=\sqrt{D_F/h}`$ and can change the sign. The second term in Eq. (111) originates from the rotation of $`h`$ along the $`x`$-axis. It decays with the thickness $`d_F`$ over another characteristic length $`\kappa _l^1`$ that can be much larger than $`\xi _F`$. Therefore this term results in a drastic change of the critical current.
The presence of the second term in Eq. (111) is especially interesting in the case when the thickness $`d_F`$ of the ferromagnetic spacer between the superconductors obeys $`\xi _F<d_F<\kappa _l^1`$. Then the main contribution to the Josephson coupling comes from the long-range triplet component of the condensate. Another important feature of this limit is that for sufficiently large values of $`Ql`$ the critical current is always positive (no possibility for the $`\pi `$-contact). This can be seen from FIG. 7.
The fact that the superconductivity looses its “exotic properties” at large $`Q`$ is quite understandable. The superconductivity is sensitive not to the local values of the exchange field but to its average on the scales of the order of the superconducting coherence length. If the exchange field oscillates very fast such that the period of the oscillations is much smaller than the superconducting coherence length, its average on this scale vanishes and therefore all new properties of the superconductivity originating from the presence of the exchange field become negligible.
To conclude this introduction we summarize the results known for $`S/F/S`$ Josephson junctions.
When the magnetization in the ferromagnetic $`F`$ is homogeneous, we have to distinguish between two different cases.
In the dirty limit ($`h\tau 1`$) the change of the sign of the critical current occurs if the thickness of the $`F`$ layer $`d_F`$ is of the order of $`\xi _F`$. The condensate function in the $`F`$ layer decays exponentially over this $`\xi _h`$ and oscillates with the same period.
In the opposite clean limit, $`h\tau 1`$, the condensate function oscillates in space with the period $`v_F/h`$ and decays exponentially over the mean free path $`l`$.
Finally, if the ferromagnetic region contains a domain wall described by a vector $`Q`$, a long-range component of the condensate appears. It decays in the $`F`$ film over a considerably larger length of the order $`\xi _N=\sqrt{D/2\pi T}`$ that can greater exceed the characteristic length ($`\sqrt{D/h}`$) in a homogeneous $`F`$ layer ($`Q=0`$). In this case the coupling between the superconductors survives even if the thickness of $`F`$ is larger than $`\xi _F`$.
It is clear that the presence of a domain wall between the superconductors is something that cannot be controlled very well experimentally. Therefore in the next section we discuss a possible experiment on $`S/F`$ multilayered structures that may help in detecting the LRTC by measuring the Josephson critical current.
### IV.1 Josephson coupling between S layers via the triplet component
In this subsection we analyze another type of multilayered $`S/F`$ structure in which the LRTC arises. This is a multilayered periodic $`\mathrm{}S/F_{n1}/S/F_n/S/F_{n+1}/S\mathrm{}`$ structure with alternating magnetization vector $`𝐌_{F,n}`$ in different $`F`$ layers. We assume that the vector $`𝐌_{F,n}`$ is turned with respect to the vector $`𝐌_{F,n1}`$ by an angle $`2\alpha `$, such that the angle increases monotonously with increasing $`n`$. We call this arrangement of the magnetization the one with a positive chirality.
In an infinite system the magnetization vector $`𝐌_F`$ averaged over $`n`$ is equal to zero (it rotates when one moves from the $`nth`$ to the ($`n+1)th,`$ layer etc.). Another type of chirality (negative chirality) is the arrangement when the angle between vectors $`𝐌_{F,n}`$ and $`𝐌_{F,n1}`$ is equal to $`2\alpha (1)^n`$. In this case the averaged vector $`𝐌_F`$ is not zero.
In Section III.3.1 we have seen that in a $`F/S/F`$ structure with a non-collinear orientation of the magnetization vectors in the $`F`$ layers the LRTC arises. If one assumes that the thickness of the $`F`$ layers $`d_F`$ is larger than the coherence length in the normal metal $`\xi _N,`$ the overlap of the condensate functions created in a $`F`$ layer by neighboring $`S`$ layers is weak, and the solutions given by Eqs.(58-66) remain valid for the multilayered $`S/F`$ structure.
Using these solutions one can calculate the Josephson current between neighboring $`S`$ layers. As the thickness $`d_F`$ is assumed to be much larger than $`\xi _F`$ (as usual, we assume that $`\xi _F<<\xi _N`$), the Josephson coupling between the $`S`$ layers is solely due to the LRTC. So, in such systems we come to a new type of the superconductivity: an odd triplet out-of-plane superconductivity and the conventional singlet in-plane superconductivity (the triplet component gives only a small contribution to the in-plane superconductivity).
Using the general Eqs. (58-66) one can perform explicit calculations for this case without considerable difficulties. As a result, the Josephson critical current $`I_c`$ can be written as follows Bergeret *et al.* (2003)
$$eR_FI_c=\pm 2\pi T\underset{\omega }{}\kappa _\omega d_Fb_1^2(\alpha )\left(1+\mathrm{tan}^2\alpha \right)e^{d_F\kappa _\omega },$$
(112)
where
$$b_1(\alpha )=f_{BCS}\mathrm{sin}\alpha \frac{\stackrel{~}{\kappa }_S^2(\stackrel{~}{\kappa }_+\stackrel{~}{\kappa }_{})\mathrm{sgn}\omega }{\mathrm{cosh}^2\mathrm{\Theta }_S\left(M_+T_{}+M_{}T_+\right)(g_{BCS}+\gamma _F\kappa _\omega \mathrm{tanh}\mathrm{\Theta }_F)},$$
$`\mathrm{\Theta }_S=\kappa _sd_S`$, $`\mathrm{\Theta }_F=\kappa _\omega d_F`$, $`\stackrel{~}{\kappa }_\pm =\kappa _\pm /(g_{BCS}+\gamma _F\kappa _\pm )`$, $`\stackrel{~}{\kappa }=\kappa _\omega /(g_{BCS}+\gamma _F\kappa _\omega \mathrm{tanh}\mathrm{\Theta }_F)`$, $`\stackrel{~}{\kappa }_S=\kappa _S/(g_{BCS}\gamma )`$ and
$`M_\pm `$ $`=`$ $`T_\pm (\stackrel{~}{\kappa }_S\mathrm{coth}\mathrm{\Theta }_S+\stackrel{~}{\kappa }\mathrm{tanh}\mathrm{\Theta }_F)+\mathrm{tan}^2\alpha C_\pm (\stackrel{~}{\kappa }_S\mathrm{tanh}\mathrm{\Theta }_S+\stackrel{~}{\kappa }\mathrm{tanh}\mathrm{\Theta }_F)`$
$`T_\pm `$ $`=`$ $`\stackrel{~}{\kappa }_S\mathrm{tanh}\mathrm{\Theta }_S+\stackrel{~}{\kappa }_\pm `$
$`C_\pm `$ $`=`$ $`\stackrel{~}{\kappa }_S\mathrm{coth}\mathrm{\Theta }_S+\stackrel{~}{\kappa }_\pm .`$
$`R_F`$ is defined as $`R_F=2d_F/(L_yL_z\sigma _F)`$. Eq. (112) describes the layered systems with both the positive (“+” sign) and negative (“-” sign) chirality.
One can see from Eq. (112) that in the case of positive chirality the critical current is positive, while if the chirality is negative the system is in the $`\pi `$-state (negative current). This means that changing the configuration of the magnetization, one can switch between the $`0`$ and $`\pi `$ state.
It is important to emphasize that the nature of the $`\pi `$-contact obtained here differs from that predicted by Bulaevskii *et al.* (1977) and observed by Ryazanov *et al.* (2001). In the latter cases the transition is due to the change of the values of either the exchange field, the temperature or the thickness of the $`F`$ film. In the case considered in this section, the negative Josephson coupling originates from the presence of the triplet component and can be realized in $`S/F`$ structures with negative chirality. Since for the positive chirality the Josephson current is positive, the result obtained gives an unique opportunity to switch experimentally between the $`0`$ and $`\pi `$-states by changing the angles of the mutual magnetization of the layers.
A similar dependence of the Josephson current $`I_c`$ on the chirality was predicted in a Josephson junction $`S_mIS_m`$ ($`I`$ is an insulator) between two magnetic superconductors $`S_m`$ by Kulic and Kulic (2001). For the magnetic superconductors considered in that work, the magnetization vector $`𝐌`$ rotated with the angle of rotation equal to $`\alpha =x𝐐𝐧_x,`$where $`𝐐`$ is the wave vector of the x- dependence of the angle $`\alpha `$, $`𝐧_x`$ is the unit vector normal to the insulating layer $`I.`$ Therefore the chirality (or spiral helicity, in terms of Kulic and Kulic) in this case is determined by the sign of the product $`𝐐_R𝐐_L,`$ where $`𝐐_{L,R}`$ is the wave vector in the left (right) magnetic superconductor.
However, there is an essential difference between the multilayered $`S/F`$ structure discussed here and the magnetic superconductors. In the magnetic superconductors with a spiral magnetization the triplet component also exists but, in contrast to the $`S/F`$ structures, the singlet and triplet components cannot be separated. In particular, in the case of a collinear alignment of $`𝐌,`$ the Josephson coupling in the $`S/F`$ structures with thick $`F`$ layers disappears, whereas it remains finite in the $`S_mIS_m`$ system.
FIG. 17 shows the dependence of the Josephson current $`I_c`$ on the angle $`\alpha `$ given by Eq. (112). If the mutual orientation of $`𝐌`$ is parallel ($`\alpha =0`$) or antiparallel ($`2\alpha =\pi `$) the amplitude of the triplet component is zero and therefore there is no coupling between the neighboring $`S`$ layers, i.e. $`I_c=0`$. For any other angles between the magnetizations the amplitude of the triplet component is finite and this leads to a non-zero critical current. At $`2\alpha =\pi /2`$ ( perpendicular orientation of $`𝐌`$) the Josephson current $`I_c`$ reaches its maximum value.
Another possible experiment for detecting the long range triplet component is the measurement of the density of states in the $`F/S/F`$ system as it is shown in FIG. 18. Kontos *et al.* (2001) determined the spatial changes of the DOS in a $`PdNi/Al`$ structure with the help of planar tunnelling spectroscopy. This method could also be used in order to detect the LRTC. It is clear that if the thickness of the $`F`$ layer in FIG. 18 is larger than the penetration of the short-range components, then any change of the DOS at the outer boundary of the $`F`$ layer may occur only due to the long range penetration of the triplet component. If the magnetizations of both $`F`$ layer are collinear no effect is expected to be observed, while for a non-collinear magnetization a change of the DOS should be seen.
### IV.2 Enhancement of the critical Josephson current
Another interesting effect in the $`S/F`$ structures that we would like to discuss is the enhancement of the Josephson critical current by the exchange field. The common wisdom is that any exchange field should reduce or destroy the singlet superconductivity. In the previous sections we have seen that this is not always so and the superconductivity can survive in the presence of a strong exchange field. But still, it is not so simple to imagine that the superconducting properties can be enhanced by the exchange field.
Surprisingly, this possibility exists and we will demonstrate now how this unusual phenomenon occurs. Although the LRTC is not essential to get the critical current enhancement, the short-range triplet component arises in this case and it plays a certain role in this effect. The enhancement of the Josephson current in the $`S/F/I/F/S`$ tunnel structures ($`I`$ stands for an insulating layer, see FIG. 19) was predicted by Bergeret *et al.* (2001b) and further considered in a subsequent work by Golubov *et al.* (2002b). As we will see below, if the temperature is low enough and the $`S/F`$ interface transparency is good, one can expect an enhancement of the critical current with increasing the exchange field provided the magnetizations of the $`F`$ layers are antiparallel to each other.
This surprising result can be obtained in a quite simple way in the limit when the thicknesses $`d_S`$ and $`d_F`$ of the $`S`$ and $`F`$ layers are smaller than the superconducting coherence length $`\xi _S\sqrt{D/2\pi T_c}`$ and the penetration length of the condensate into the ferromagnet $`\xi _F\sqrt{D/h}`$, respectively. In this case one can assume that the quasiclassical Green’s functions does not depend on the space coordinates and, in particular, the superconducting order parameter $`\mathrm{\Delta }`$ is a constant in space. Moreover, instead of considering the dependence of the exchange field $`h`$ on the coordinates one can replace it by a homogeneous effective exchange field $`h_{eff}`$ with a reduced value. Therefore we use in our calculations effective fields $`\mathrm{\Delta }_{eff}`$ and $`h_{eff}`$ defined as
$`\mathrm{\Delta }_{eff}/\mathrm{\Delta }`$ $`=`$ $`\nu _Sd_S\left(\nu _Sd_S+\nu _Fd_F\right)^1,`$ (113)
$`h_{eff}/h`$ $`=`$ $`\nu _Fd_F\left(\nu _Sd_S+\nu _Fd_F\right)^1,`$ (114)
where $`\nu _S`$ and $`\nu _F`$ are the densities of states in the superconductor and ferromagnet, respectively.
With this simplification, we can write the Gor’kov equations for the normal and anomalous Green’s functions in the spin-space as
$`\left(i\omega +\xi \sigma 𝐡\right)\widehat{G}_\omega +\widehat{\mathrm{\Delta }}\widehat{F}_\omega ^+`$ $`=`$ $`1`$ (115)
$`\left(i\omega +\xi \sigma 𝐡\right)\widehat{F}_\omega +\widehat{\mathrm{\Delta }}\widehat{G}_\omega `$ $`=`$ $`0,`$ (116)
where $`\sigma =(\widehat{\sigma }_1,\widehat{\sigma }_2,\widehat{\sigma }_3)`$ are the Pauli matrices and $`\xi =ϵ\left(𝐩\right)ϵ_F,`$ $`\epsilon _F`$ is the Fermi energy, $`\epsilon \left(𝐩\right)`$ is the spectrum, and $`\omega =\left(2n+1\right)\pi T`$ are Matsubara frequencies. (We omit the subscript $`eff`$ in Eqs. (115-116) and below).
In order to calculate the Josephson current $`I_J`$ through the tunnel junction represented in FIG. 19 we use the well known standard formula
$$I_J=\left(2\pi T/eR\right)\mathrm{Tr}\underset{n}{}\widehat{f}_\omega (h_1)\widehat{f}_\omega (h_2)\mathrm{sin}\phi ,$$
(117)
where
$$\widehat{f}_\omega =\frac{i}{\pi }\widehat{F}_\omega d\xi $$
(118)
is the quasiclassical anomalous Green’s function, $`\phi `$ is the phase difference between both the superconductors, $`R`$ is the barrier resistance and $`h_{1,2}`$ are the exchange fields of the left and the right $`F`$-layers.
The only difference between Eqs. (117, 118) and the corresponding equations in the absence of the exchange field is the dependence of the condensate function $`\widehat{f}_\omega `$ on $`h`$. This dependence can be found immediately from Eqs. (115, 116).
$$\widehat{f}_\omega =\widehat{\mathrm{\Delta }}\left(\left(\omega +i\sigma 𝐡\right)^2+\mathrm{\Delta }^2\right)^{1/2}$$
(119)
What remains to be done is to insert the condensate function $`\widehat{f}`$ into Eq. (117) for certain exchange fields $`h_1`$ and $`h_2`$ and to calculate the sum over the Matsubara frequencies $`\omega `$. Although it is possible to carry out these calculations for arbitrary vectors $`h_1`$ and $`h_2,`$ we restrict our consideration by the cases when the absolute values the magnetizations $`h_1`$ and $`h_2`$ are equal but the magnetizations are either parallel or antiparallel to each other. This simplifies the computation of the Josephson current but, at the same time, captures the essential physics of the phenomenon.
Using Eqs. (117-119) and assuming first that $`h_1`$ and $`h_2`$ are parallel to each other we write the expression for the critical current as Bergeret *et al.* (2001b)
$$I_J=I_c\mathrm{sin}\phi $$
(120)
$$I_c^{\left(p\right)}=\frac{\mathrm{\Delta }^2\left(T\right)4\pi T}{eR}\underset{\omega }{}\frac{\omega ^2+\mathrm{\Delta }^2(T,h)h^2}{\left(\omega ^2+\mathrm{\Delta }^2(T,h)h^2\right)^2+4\omega ^2h^2},$$
(121)
The corresponding equation for the antiparallel configuration is different from Eq. (121) and can be written as
$$I_c^{\left(a\right)}=\frac{\mathrm{\Delta }^2\left(T\right)4\pi T}{eR}\underset{\omega }{}\frac{1}{\sqrt{\left(\omega ^2+\mathrm{\Delta }^2(T,h)h^2\right)^2+4\omega ^2h^2}}.$$
(122)
One can easily check that the critical current $`I_c^{\left(p\right)}`$ for the parallel configuration, Eq. (121), is always smaller than the current $`I_c^{\left(a\right)}`$ for the antiparallel case. These two expressions are equal to each other only in the absence of any magnetization.
In FIGS. 20 and 21 we represent the dependence of the critical current on the strength of the exchange field. We see from FIG. 20 that for the parallel configuration the exchange field reduces the value of the Josephson current and this is exactly what one could expect. At the same time, the critical current grows with the exchange field for the antiparallel configuration at low temperatures, which is a new intriguing effect (see FIG. 21).
This unexpected result can be understood from Eq. (122) rather easily without making numerics. In the limit $`T0`$ the sum over the Matsubara frequencies can be replaced by an integral and one can take for the superconducting order parameter $`\mathrm{\Delta }`$ the values $`\mathrm{\Delta }=\mathrm{\Delta }_0`$ if $`h<\mathrm{\Delta }_0,`$ and $`\mathrm{\Delta }=0`$ if $`h>\mathrm{\Delta }_0`$, where $`\mathrm{\Delta }_0`$ is the BCS order parameter in the absence of an exchange field (see, e.g. Larkin and Ovchinikov (1964)).
Inserting this solution into Eq. (122) one can see that the Josephson critical current $`I_c^{\left(a\right)}`$ grows with increasing exchange field. Moreover, formally it diverges logarithmical when $`h\mathrm{\Delta }_0`$
$$I_c^{\left(a\right)}\left(h\mathrm{\Delta }_0\right)\frac{I_c\left(0\right)}{\pi }\mathrm{ln}\left(\mathrm{\Delta }_0/\omega _0\right),$$
(123)
where $`I_c\left(0\right)`$ is the critical current in the absence of the magnetic moment at $`T=0`$, and $`\omega _0`$ is a parameter needed to cut the logarithm at low energies.
When deriving Eqs. (121, 122) the conventional singlet superconducting pairing was assumed. The electrons of a Cooper pair have the opposite spins. This picture of a superconducting pairs with the opposite spins of the electrons helps in the understanding of the effect.
If the magnetic moments in both the magnetic layers are parallel to each other, they serve as an obstacle for the Cooper pair because the pairs located in the region of the ferromagnet demand more energy. This leads to a reduction of the Josephson current. However, if the exchange fields of the different layers are antiparallel, they may favor the location of the Cooper pairs in the vicinity of the Josephson junction. A certain probability exists that one of the electrons of the pair is located in one layer, whereas the other is in the second layer. Such a possibility is energetically favorable because the spins of the electrons of the pair can now have the same direction as the exchange fields of the layers. Then it is more probable for the pairs to be near the junction even in comparison with a junction without exchange fields and, as a result, the critical current may increase.
The results presented above have been obtained for the $`SF/I/FS`$ structure by Bergeret *et al.* (2001b). Earlier a formula for the Josephson critical current in the $`S_mIS_m`$ ($`S_m`$ is the magnetic superconductor) junction was presented by Kulic and Kulic (2001). From that formulae one could, in principle, derive an enhancement of the critical current for the antiparallel $`M`$ orientation in magnetic superconductors $`S_m`$. Unfortunately, the authors seem to have missed this interesting effect.
Some remarks should be made at this point:
1) The results presented above are valid in the tunnelling regime, i.e. when the transparency of the insulating barrier $`I`$ is low enough. Golubov *et al.* (2002b) have shown that a smearing of the singularity of $`I_c^{(a)}`$ is provided by a finite temperature or a not very low barrier transparency. The maximum of the critical current for the antiparallel configuration $`I_c^{(a)}`$ decreases with decreasing resistance of the $`I`$ layer. The effect becomes weaker as the thickness of the $`F`$ layer grows.
2) We assumed that the $`S/F`$ interface was perfect. In a structure with a large $`S/F`$ interface resistance $`R_{S/F}`$ the bulk properties of the $`S`$ film are not considerably influenced by the proximity of the $`F`$ film (to be more precise, the condition $`R_{S/F}>(\nu _Fd_F/\nu _Sd_S)\rho _F\xi _F`$ must be satisfied, where $`\rho _F`$ is the specific resistance of the $`F`$ film). Then, as one can readily show (see section II.2), a minigap $`ϵ_{bF}=\left(D\rho \right)_F/\left(2R_{S/F}d_F\right)`$ arises in the $`F`$ layer. The Green’s functions in the $`F`$ layers have the same form as before with $`\mathrm{\Delta }`$ replaced by $`ϵ_{bF}`$. The singularity in $`I_c(h)`$ first occurs at $`h`$ equal to $`ϵ_{bF}`$.
A physical explanation for the singular behavior of the critical current $`I_c^{(a)}`$ was given by Golubov *et al.* (2002b) These authors noticed that the density of states in the $`F`$ layer has a singularity when $`h=ϵ_{bF}`$. At this value of $`h`$ the maximum of $`I_c^{(a)}`$ is achieved due to an overlap of two $`ϵ^{1/2}`$ singularities. This leads to the logarithmic divergency of the critical current in the limit $`T0`$ in analogy with the well known Riedel peak in $`SIS`$ tunnel junctions for the voltage difference $`2\mathrm{\Delta }`$. In the latter case the shift of the energy is due to the electric potential.
Golubov *et al.* (2002b) have also shown that for the parallel configuration, at $`h=ϵ_{bF}`$ the critical current changes its signs, i.e. there is a transition from $`0`$ to a $`\pi `$ junction. Similar results were obtained by Krivoruchko and Koshina (2001a, b). The case of an arbitrary $`S/F`$ transparency was also studied by Li *et al.* (2002); Chtchelkatchev *et al.* (2002); Barash *et al.* (2002). In the paper by Barash *et al.* (2002) the authors calculated the Josephson current as a function of the angle between the magnetizations in the $`F`$ film.
## V Reduction of the Magnetization due to Superconductivity: Inverse Proximity Effect
Until now we have been studying the superconducting properties of different $`S/F`$ structures for a fixed magnetization. This means that we assumed a certain value for this quantity and its dependence on coordinates. The implied justification of this assumption was that the ferromagnetism is a stronger phenomenon than the superconductivity and the magnetic moment of conventional ferromagnets can hardly be affected by the superconductivity.
This assumption is certainly correct in many cases but not always. Often the presence of the superconductivity can drastically change magnetic properties of the ferromagnets even if they are strong.
Experiments performed by Mühge *et al.* (1998) and Garifullin *et al.* (2002) showed that the total magnetization of certain $`S/F`$ bilayers with strong ferromagnets decreased with lowering the temperature below the critical superconducting transition temperature $`T_c`$. As an explanation, it was suggested that due to the proximity effect domains with different magnetization appeared in the magnetic materials and this could reduce the total magnetization. At the same time, quantitative estimates based on an existing theory Buzdin and Bulaevskii (1988) led to a conclusion that this mechanism was not very probable.
In this Chapter we address the problem of the reduction of the magnetic moment by the presence of a superconductor assuming again that, in the absence of the ferromagnet, we would have the conventional singlet superconducting pairing. It turns out that two different and independent mechanisms that lead to a decrease of the magnetization in $`S/F`$ heterostructures due to the proximity effect exist and we give a detailed account of them.
In order to study the magnetic properties we have to choose a model. One can distinguish two different types of the ferromagnetism: a) itinerant ferromagnetism due to the spin ordering of free electrons and b) ferromagnetism caused by localized spins. Most of ferromagnetic metals show both of the types of ferromagnetism simultaneously, i.e. their magnetization consists of both contributions.
We consider a model in which the conducting electrons interact with the localized moments via an effective exchange interaction. The corresponding term in the Hamiltonian is taken in the form (see Appendix A):
$$d^3r\psi ^{}(𝐫)_\alpha (J𝐒(𝐫)\sigma ))_{\alpha \beta }\psi (𝐫)_\beta .$$
(124)
This term is suitable to describe $`sd`$ or $`sf`$ interaction between the $`s`$ and localized $`d`$ and $`f`$ electrons. We also consider the ferromagnetic interaction between the localized moments. This interaction can be very complicated and to determine it, one should know the detailed band structure of the metal as well as different parameters. However, all these details are not important for us and we write the interaction between the localized spins phenomenologically as
$$\underset{ij}{}𝒥_{ij}𝐒_i𝐒_j.$$
(125)
It is assumed that $`𝒥`$ is positive. This interaction, Eq. (125) , is responsible for the ferromagnetic alignment of the localized moments and is known as the Heisenberg Hamiltonian.
So , we consider a metallic ferromagnet in which the conduction electrons interact with localized magnetic moments. The ferromagnetic interaction (125) assures a finite magnetic moment of the background. The total magnetization is the sum of the background magnetization (localized moments) and the magnetization of the polarized free electrons.
In the next two sections we discuss the two different mechanisms that lead to a decrease of the magnetization at low temperatures. In Section V.1 we consider a possibility of changing the magnetic order of the localized magnetic spins in a $`F`$ film deposited on top of a bulk superconductor. The contribution from free electrons to the magnetization is first assumed to be small. We will see that for not too strong ferromagnetic coupling $`𝒥`$ the proximity effect may lead to an inhomogeneous magnetic state. Contrary to this case, we consider in Section V.2 an itinerant ferromagnet in which the main contribution to the magnetization is due to free electrons. We will show that the magnetization of free electrons may decrease at low temperatures due to a some kind of spin screening. Thus, both effects may lead to the decrease in the magnetization observed in experiments Mühge *et al.* (1998); Garifullin *et al.* (2002).
### V.1 Cryptoferromagnetic state
In 1959 Anderson and Suhl suggested that superconductivity could coexist with a nonhomogeneous magnetic order in some type of materials. Anderson and Suhl called this state cryptoferromagnetic state.
The reason for this coexistence is that, if the magnetization direction varies over a scale smaller than the superconducting coherence length, the superconductivity may survive despite the ferromagnetic background. This is due to the fact that the superconductivity is sensitive to the ferromagnetic moment averaged on the scale of the size of Cooper pairs rather than to its local values.
In 1988 Buzdin and Bulaevskii discussed properties of a bilayer system consisting of a conventional superconductor in contact with a ferromagnet. They have shown that the magnetic ordering in the magnet might take the form of a structure consisting of small size domains, such that the superconductivity is not destroyed. Of course, as follows from Eq. (125), the formation of a domain-like structure costs a magnetic energy but this is compensated by the energy of the superconductor that would have been lost if the magnetic order remained ferromagnetic.
This is only possible if the stiffness of the magnetic order parameter ($`𝒥`$) is not too large. For instance this nonhomogeneous magnetization occurs in magnetic superconductors as those studied by Bulaevskii *et al.* (1985). But can one see it in the heterostructures containing strong ferromagnets like $`Fe`$ or $`Ni`$ in contact with conventional superconductors?
At first glance, it seems impossible, since the Curie temperature of, for example, iron is hundred times or more larger than the critical temperature of a conventional superconductor. Therefore any change of the ferromagnetic order look much less favorable energetically than the destruction of the superconductivity in the vicinity of the $`S/F`$ interface.
This simple argument was however questioned in the experiments performed by Mühge *et al.* (1998) on $`Fe/Nb`$ bilayers and by Garifullin *et al.* (2002) on $`V/Pd_{1x}Fe_x`$ structures. Direct measurements of the ferromagnetic resonance has shown that in several samples with thin ferromagnetic layers the average magnetic moment started to decrease below the superconducting transition temperature $`T_c`$.
Of course, one can reduce the influence of the ferromagnet on the superconductor by diminishing the thickness of the ferromagnet. Using the formulae obtained by Buzdin and Bulaevskii (1988), Mühge *et al.* (1998) estimated the thickness of the ferromagnet for which the superconductivity was still possible and got a value of the order of $`1`$Å, which created a doubt on the explanation of the experiment in this way.
At the same time, the use of the formulae derived by Buzdin and Bulaevskii was not really justified because the calculations were done for thick but weak ferromagnets assuming a strong anisotropy of the ferromagnet that was necessary for a formation of the domain walls with the magnetization vector changing its sign but not its axis.
Bergeret, Efetov, and Larkin (2000) investigated theoretically the possibility of a cryptoferromagnetic-like (CF) state in $`S/F`$ bilayers with parameters corresponding to the experiments by Mühge *et al.* and Garifullin *et al.*. In that work a CF state with a magnetic moment that rotates in space was considered. This corresponds to a weak anisotropy of the ferromagnet, which was the case in the samples studied in Mühge *et al.* (1998). In particular, Bergeret *et al.* (2000) studied a phase transition between the CF and the ferromagnetic (FM) phases. The calculations were carried out in the limit $`d_F\xi _h=v_0/h,`$ $`T_chϵ_0`$, $`v_0`$ and $`ϵ_0`$ are the Fermi velocity and Fermi energy, respectively. This limit is consistent with the parameters of the experiment of Mühge *et al.* (1998), Garifullin *et al.* (2002). We present here the main ideas of this work.
The Hamiltonian describing the bilayer structure in FIG. 22 can be written as
$$H\left(\gamma \right)=H_0+H_{BCS}\gamma 𝑑𝐫\mathrm{\Psi }_\alpha ^+(𝐫)\left[𝐡(𝐫)\sigma \right]_{\alpha \beta }\mathrm{\Psi }_\beta (𝐫)+H_M,$$
(126)
where the integration must be taken in the region $`d<x<0`$. Here $`H_0`$ is the one-particle electron energy (including an interaction with impurities), $`H_{BCS}`$ is the usual term describing the conventional BCS superconductivity in the superconductor $`S`$ and the third term describes the interaction between localized moments and conduction electrons, where $`\gamma `$ is a constant that will be put to $`1`$ at the end (see Appendix A).
The term $`H_M`$ describes the interaction between the localized moments in the ferromagnet (cf. Eq.(125)). We assume that the magnetization of the localized spins is described by classical vectors and take into account the interaction between neighboring spins only. In the limit of slow variations of the magnetic moment in space with account of Eq. (125), the Hamiltonian $`H_M`$ can be written in the form
$$H_M=𝒥\left[\left(S_x\right)^2+\left(S_y\right)^2+\left(S_z\right)^2\right]𝑑V,$$
(127)
where the magnetic stiffness $`𝒥`$ characterizes the strength of the coupling between the localized moments in the $`F`$ layer and the $`S_i`$ are the components of a unit vector that are parallel to the local direction of the magnetization.
We assume that the magnetic moments are directed parallel to the $`S/F`$ interface and write the spin vector $`𝐒`$ as $`𝐒=(0,\mathrm{sin}\theta ,\mathrm{cos}\theta )`$. A perpendicular component of the magnetization would induce strong Meissner currents in the superconductor, which would require a greater additional energy.
The condition for an extremum of the energy $`H_M`$, Eq. (127) can be written as
$$\mathrm{\Delta }\theta =0$$
(128)
Solutions of Eq. (128) can be written in the form $`\theta =Qy`$, where $`Q`$ is the wave vector characterizing the rotation in space (see FIG. 22). The value $`Q=0`$ corresponds to the ferromagnetic state.
What we want to do now is to compare the energies of the ferromagnetic and cryptoferromagnetic states. The latter will be considered for the case with a rotating in space magnetic moment $`\theta =Qy`$. This should be energetically more favorable than the domain-like structure one provided the magnetic anisotropy of $`F`$ is low. Such a CF state corresponds to a so called Neel wall (see for example Aharoni (1996)).
Strictly speaking, one has to take into account also a magnetostatic energy due to a purely magnetic interaction of the magnetic moments. However, if the condition
$$\frac{𝒥}{M_s^2}d^2$$
(129)
where $`M_s`$ is the magnetic moment per volume, is fulfilled one can neglect its contribution with respect to the one of the exchange energy Aharoni (1996).
Taking typical values of the parameters for $`Fe`$: $`M_s=800`$emu/cm<sup>3</sup> and $`J=2.10^6`$erg/cm one can see that the condition (129) requires that the thickness $`d`$ of the ferromagnet is smaller than $`10nm`$, which corresponds to comparatively thick layers. Throughout this section this condition is assumed to be fulfilled.
In this case the magnetic energy $`\mathrm{\Omega }_M`$ (per unit surface area) is given by the simple expression
$$\mathrm{\Omega }_M=JdQ^2.$$
(130)
In order to calculate the superconducting energy $`\mathrm{\Omega }_S`$ one has to take into account the fact that the order parameter should be destroyed, at least partially, near the contact with the ferromagnet. This means that the order parameter $`\mathrm{\Delta }`$ is a function of the coordinate $`x`$ perpendicular to the interface. As we want to minimize the energy we should look for a non-homogeneous solution for $`\mathrm{\Delta }\left(x\right)`$ of non-linear equations describing the superconductivity. Near the critical temperature $`T_c`$ one can use Ginzburg-Landau equations. The proper solution of these equations can be written in the form
$$\mathrm{\Delta }\left(x\right)=\mathrm{\Delta }_0\mathrm{tanh}\left(\frac{x}{\sqrt{2}\xi _{GL}\left(T\right)}+C\right)$$
(131)
where $`\mathrm{\Delta }_0`$ the value of the order parameter in the bulk, and $`\xi _{GL}`$ is the correlation length of the superconductor defined in Eq. (2). Near $`T_c`$ this length can be much larger than the length $`\xi _S`$. The parameter $`C`$ in Eq. (131) is a number that has to be found from boundary conditions.
The solution for $`\mathrm{\Delta }\left(x\right)`$, Eq. (131) is applicable at distances exceeding the length $`\xi _S`$ and therefore we cannot use it near the interface.
Having fixed the constant $`C`$ one can compute the decrease of the superconducting energy due to the suppression of superconductivity in the $`S`$ layer using the Ginzburg-Landau free energy functional (e.g. de Gennes, 1966). The decrease of the superconducting energy $`\mathrm{\Omega }_S`$ per unit area at the $`F/S`$ interface is a function of $`C`$ and can be written as
$$\mathrm{\Omega }_S=\frac{\sqrt{\pi }}{6\sqrt{2}}|\tau |^{3/2}\left(2+K\right)(1K)^2,$$
(132)
where $`K=\mathrm{tanh}C`$, and $`\tau =(TT_c)/T_c.`$
It remains only to determine the contribution from the third term of the Hamiltonian (126). The corresponding free energy $`\mathrm{\Omega }_{M/S}`$ is given by the expression:
$$\mathrm{\Omega }_{M/S}=i\pi T\nu _0\frac{\mathrm{Tr}}{2}\underset{\omega }{}_0^1𝑑\gamma d^3𝐫(𝐡\sigma )\widehat{g}_0,$$
(133)
where $`\nu _0`$ is the density of states and $`\widehat{g}_0`$ is the quasiclassical Green’s function averaged over all directions of the Fermi velocity.
Since the exchange field $`h`$ in a strong ferromagnet may be much higher than the value of $`\tau ^1`$ (here $`\tau `$ is the momentum relaxation time), one has to solve the Eilenberger equation in the $`F`$ region and the Usadel equation in the $`S`$ region. Solutions for these equations in both the superconductor and ferromagnet were obtained by Bergeret *et al.* (2000).
Thus, the total energy is given by $`\mathrm{\Omega }=\mathrm{\Omega }_M+\mathrm{\Omega }_S+\mathrm{\Omega }_{M/S}`$, Eqs. (130, 132, 133). As a result, one can express the free energy as a function of two unknown parameters, $`K`$ and $`Q`$. One can find these parameters from the condition that the free energy must be minimal, which leads to the equations
$$\mathrm{\Omega }/K=\mathrm{\Omega }/Q=0$$
(134)
One can show that the CF-F transition is of second order, which means that near the transition the parameter $`Q`$ is small. At the transition it vanishes and this gives an equation binding the parameters. Solving the equation numerically we come to the phase diagram of FIG. 23 determining the boundary between the ferromagnetic and cryptoferromagnetic states. The parameters $`a`$ and $`\lambda `$ used in FIG. 23 are defined as
$$a^2\frac{2h^2d_f^2}{DT_c\eta ^2},\lambda \frac{𝒥d_F}{\nu _F\sqrt{2T_cD^3}}\frac{7\zeta (3)}{2\pi ^2},$$
(135)
where $`\eta `$ is the ratio between the Fermi velocities $`v_0^F/v_0^S`$. It is clear from Eqs. (135) that the parameter $`a`$ is related to the exchange energy $`h`$, while $`\lambda `$ is the related to the magnetic stiffness $`𝒥`$.
The conclusion that the phase transition between F and CF states should be of the second order was drawn neglecting the magnetostatic interaction. The direct magnetic interaction can change this transition to a first order one Buzdin (2005b). However, in the limit of Eq. (129), this first order transition will be inevitably close to the second order one. Such a modification of the type of the phase transition is out of the focus of this review.
Let us make estimates for the materials used in the experiments. Performing ferromagnetic resonance measurements, Mühge *et al.* (1998) have observed a decrease of the effective magnetization of a $`Nb/Fe`$ bilayer. The stiffness $`𝒥`$ for materials like $`Fe`$ and $`Ni`$ is $`60K/`$Å. The parameters characterizing $`Nb`$ can be estimated as follows: $`T_c=10K`$ , $`v_F10^8`$cm/s, and $`l100`$Å. The thickness of the magnetic layer is of order $`d=10`$Å, and the exchange field $`h10^4K`$ which is proper for iron.
Assuming that the Fermi velocities and energies of the ferromagnet and superconductor are close to each other, we obtain $`a25`$ and $`\lambda 6.10^3`$. It is clear from FIG. 23 that the cryptoferromagnetic state is hardly possible in the $`Fe/Nb`$ samples used in the experiment Mühge *et al.* (1998).
However, one can in principle explain the observed, decrease of the magnetization taking a closer look at the structure of the $`S/F`$ interface. In the samples analyzed by Mühge *et al.* the interface between the $`Nb`$ and $`Fe`$ layers is rather rough. So, one can expect that in the magnetic layers there were “islands” with smaller values of $`𝒥`$ and/or $`h`$. A reduction of these parameters in the $`Fe/Nb`$ bilayers is not unrealistic because of the formation of non-magnetic “dead” layers that can also affect the parameters of the ferromagnetic layers. If the cryptoferromagnetic state were realized only on the islands, the average magnetic moment would be reduced but would remain finite. Such a conclusion correlates with what one observes experimentally. One can also imagine islands very weakly connected to the rest of the layer, which would lead to smaller energies of a non-homogeneous state.
Let us now consider the experiment by Garifullin *et al.* (2002) on $`Pd_{0.97}Fe_{0.03}/V`$. Due to the low concentration of iron, the magnetic stiffness and the exchange field of the $`F`$-layers is much lower than the one in the case of a pure iron. For this system, one estimates the parameters as(see Garifullin *et al.* (2002)) $`J60K/nm`$, $`h100K`$. Assuming again that the Fermi velocities of $`V`$ and $`Pd_{1x}Fe_x`$ are close to each other, Garifullin *et al.* (2002) obtained for the sample with $`d_F=1.2`$nm the following values of the parameters $`a1.2`$ and $`\lambda \mathrm{1.3.10}^3`$.
Using these values for $`a`$ and $`\lambda `$ one can see from the phase diagram in FIG. 23 that there can be a transition from the F to the CF state at $`|\tau |0.2`$, which corresponds to $`T2.4K`$. The decrease of the effective magnetization $`M_{eff}`$ with decreasing temperature was not observed in samples with larger $`F`$ thickness $`d_F`$: $`M_{eff}`$ was a temperature-independent constant for the sample with $`d_F=4.4nm`$ and $`d_S=37.2nm`$. In the sample with $`d_F=1.2nm`$ and $`d_S=40nm`$ the effective magnetization $`M_{eff}`$ decreased by $`50\%`$ with cooling from $`T4K`$ to $`T1.5K`$. This fact is again in accordance with the predictions of Bergeret *et al.* (2000).
The results of this section demonstrate that not only ferromagnets change superconducting properties but also superconductivity can affect ferromagnetism. This result is valid, in particular, for strong ferromagnets, although the thickness of the ferromagnetic layers must be small in this case.
The exchange interaction between the superconducting condensate and the magnetic order parameter reduces the energy of the system if the direction of the magnetization vector $`M_F`$ is not constant in space but oscillates. Provided the energy of the anisotropy is small, this interaction leads to the formation of a spiral magnetic structure in the $`F`$ film.
As we will see in the next section the appearance of the CF-state is not the only effect that leads to a reduction of the effective magnetization in S/F structures. We will show that the proximity effect may also lead to a change of the absolute value of the magnetic moment $`M_F`$ in the ferromagnet and to an induced magnetization $`M_S`$ in the superconductor.
### V.2 Ferromagnetism induced in a superconductor
In the previous section we have seen that the superconductivity can affect the magnetic ordering changing the orientation of magnetic moments in the ferromagnetic film. In this section we want to demonstrate that another mechanism for a change of the total magnetization of a $`S/F`$ system exists. In contrast to the phenomenon discussed in the previous Section, the orientation of the magnetic moments in the $`F`$ film does not change but the magnitude of the magnetization both in the $`F`$ and $`S`$ films does.
This change is related to the contribution of free electrons both in the ferromagnet ($`\delta M_F`$) and in the superconductor ($`M_S`$) to the total magnetization. On one hand, the DOS in the $`F`$ film is reduced due to the proximity effect and therefore $`\delta M_F`$ is reduced. On the other hand, the Cooper pairs in the $`S`$ film are polarized in the direction opposite to $`M_F`$, where $`M_F`$ is the magnetization of free electrons in the ferromagnet.
Let us consider first a bulk ferromagnet and derive a relation between the exchange field and the magnetization of the free electrons. The exchange field $`h=JS`$ in the ferromagnet can be due to the localized moments (see Eq. (124)) or due to the free electrons in the case of an itinerant ferromagnet<sup>4</sup><sup>4</sup>4In many papers the exchange ”field” $`h`$ is defined in another way ($`h=JS`$) so that the energy minimum corresponds to orientation of the vector $`<\sigma >`$ antiparallel to the vector $`h.`$ In this case the magnetic moment $`m=\mu _B<\sigma >`$ is parallel to $`h`$. Both definitions lead to the same results. In some ferromagnets both the localized and itinerant moments contribute to the magnetization.
The magnetization of the free electrons is given by
$$M=\frac{i}{4}\mu _B\frac{d\omega }{2\pi }\frac{d^3p}{(2\pi )^3}\mathrm{Tr}\widehat{\tau }_3\widehat{\sigma }_3\left(\stackrel{ˇ}{G}^R\stackrel{ˇ}{G}^A\right)n_p,$$
(136)
where $`\mu _B`$ is an effective Bohr magneton and $`n_P`$ is the Fermi distribution function of the free electrons. The expression in front of $`n_P`$ in Eq. (136) determines the DOS that depends on the exchange field $`𝐡`$. We assume that the magnetization is oriented along the $`z`$-axis.
Using Eq. (136) one can easily compute the contribution of the free electrons to the magnetization in a bulk ferromagnet. In the simplest case of a normal metal with a quadratic energy spectrum we have
$$M_F=\frac{\mu _B}{(2\pi )^2}p^2𝑑p\left[n(\xi _ph)n(\xi _p+h)\right],$$
(137)
where $`\xi _p=p^2/2mϵ_F`$. At $`T=0`$ the magnetization is given by:
$$M_{F0}=\frac{\mu _B}{2(3\pi ^2)}\left(p_+^3p_{}^3\right)$$
(138)
where $`p_\pm =\sqrt{2m(ϵ_F\pm h)}`$ are the Fermi momenta for spin up and spin down electrons. In the quasiclassical limit it is assumed that $`hϵ_F`$, and therefore
$$M_{F0}\mu _B\nu h,$$
(139)
where $`\nu =p_{F0}m/\pi ^2`$ is the density of states at the Fermi level, and $`p_{F0}=\sqrt{2mϵ_F}`$ is the Fermi momentum in the absence of the exchange field<sup>5</sup><sup>5</sup>5Actually Eq.(139) is valid not only in the case of a quadratic spectrum but also in a more general case.. For the temperature range $`Th`$ we are interested in, one can assume that the magnetization of the ferromagnet does not depend on $`T`$ and is given by Eq. (139).
Now let us consider a $`S/F`$ system with a thin F layer (see FIG. 24) and ask a question: Is the magnetization of the itinerant electrons modified by the proximity effect? We assume that the exchange field of the ferromagnet $`F`$ is homogeneous and aligned in the $`z`$\- direction, which is the simplest situation.
At first glance, it is difficult to expect anything interesting in this situation and, to the best of our knowledge, such a system has not been discussed until recently.
However, physics of this heterostructure is actually very interesting and is general for any shape of the $`S`$ and $`F`$ regions. It turns out that the proximity effect reduces the total magnetization of the system and this effect can be seen as a certain kind of “spin screening”.
Before doing explicit calculations we would like to explain the phenomenon in simple words. If the temperature is above $`T_c`$, the total magnetization of the system $`M_{tot}`$ equals $`M_{0F}d_F`$, where $`d_F`$ is the thickness of the $`F`$-layer. When the temperature is lowered below $`T_c,`$ the $`S`$ layer becomes superconducting and the Cooper pairs with the size of the order of $`\xi _S\sqrt{D_S/2\pi T_c}`$ arise in the superconductor. Due to the proximity effect the Cooper pairs penetrate the ferromagnet. In the case of a homogeneous magnetization the Cooper pairs consist, as usual, of electrons with the opposite spins, such that the total magnetic moment of a pair is equal to zero. The exchange field is assumed to be not too strong, otherwise the pairs would break down.
It is clear from this simple picture that pairs located entirely in the superconductor cannot contribute to the magnetic moment of the superconductor because their magnetic moment is simply zero, which is what one could expect. Nevertheless, some pairs are located in space in a more complicated manner: one of the electrons of the pair is in the superconductor, while the other moves in the ferromagnet. These are the pairs that create the magnetic moment in the superconductor. This follows from the simple fact that the direction along the magnetic moment $`𝐌`$ in the ferromagnet is preferable for the electron located in the ferromagnet (we assume a ferromagnetic type of exchange field) and this makes the spin of the other electron of the pair be antiparallel to $`𝐌`$. So, all such pairs with one electron in the ferromagnet and one in the superconductor equally contribute to the magnetic moment in the bulk of the superconductor. As a result, a ferromagnetic order is created in the superconductor, the direction of the magnetic moment in this region being opposite to the direction of the magnetic moment $`𝐌`$ in the ferromagnet. Moreover, the induced magnetic moment penetrates the superconductor over the size of the Cooper pairs $`\xi _S`$ that can be much larger than $`d_F`$.
This means that although the magnetization $`M_S`$ induced in the superconductor is less than the magnetization in the ferromagnet $`M_{F0}`$, the total magnetic moment in the superconductor $`\overline{M}_S=_Sd^3rM_S(r)`$ may be comparable with the magnetic moment of the ferromagnet in the normal state $`\overline{M}_{F0}=M_{F0}V_F`$, where $`V_F=d_F`$ in the case of a flat geometry ($`\overline{M}_{F0}`$ is the magnetic moment per unit square) and $`V_F=4\pi a_F^3/3`$ is the volume of the spherical ferromagnetic grain. It turns out that the total magnetic moment of the ferromagnetic region ( film or grain) $`\overline{M}_{F0}=\mu _B\nu _FhV_F`$ due to free electrons is compensated at zero temperature by the total magnetic moment $`\overline{M}_S`$ induced in the superconductor. This statement is valid if the condition
$$\mathrm{\Delta }<<h<<E_{Th}=D_F/d_F^2$$
(140)
is fulfilled. If the thickness of the F film (or radius of the F grain) is not small in comparison with the correlation length $`\xi _S`$, the situation changes: the induced magnetic moment $`\overline{M}_S`$ is much smaller than $`\overline{M}_{F0}`$ but a variation of the magnetic moment of the ferromagnetic film (or grain) $`\delta M_F`$ becomes comparable with $`\overline{M}_{F0}`$. The latter is caused by a change in the density of states of the ferromagnet due to the proximity effect. However, the case of a large ferromagnet size is less interesting because the exchange field $`h`$ should be smaller than $`\mathrm{\Delta }`$ (the full screening of $`\overline{M}_{F0}`$ occurs only if the second condition in Eq.(140) is fulfilled).
Using similar arguments we can come to a related effect: the magnetic moment in the ferromagnet should be reduced in the presence of the superconductivity because some of the electrons located in the ferromagnet condensate into Cooper pairs and do not contribute to the magnetization.
¿From this qualitative and somewhat oversimplified picture one can expect that the total magnetization of the $`S/F`$ system will be reduced for temperatures below $`T_c`$. Both the mechanism studied here and that of the last section lead to a negative change of the total magnetization. Thus, independently of the origin of ferromagnetism, they can explain, at least qualitatively, the experimental data of Mühge *et al.* (1998) and Garifullin *et al.* (2002).
The ideas presented above can be confirmed by calculations based on the Usadel equation. In order to determine the change of the magnetization it is enough to compute the quasiclassical Green’s functions $`\stackrel{ˇ}{g}^{R(A)}=(i/\pi )𝑑\xi \stackrel{ˇ}{G}^{R(A)}`$ and, in particular, the component proportional to $`\widehat{\tau }_3`$$`\widehat{\sigma }_3`$.
The matrix Green’s function has the form (we write $`\stackrel{ˇ}{g}`$ in Matsubara representation: $`\stackrel{ˇ}{g}(\omega )=\stackrel{ˇ}{g}^R(i\omega )`$ for positive $`\omega `$)
$$\stackrel{ˇ}{g}=\widehat{\tau }_3\widehat{g}+i\widehat{\tau }_2\widehat{f}.$$
(141)
In the ferromagnet we represent, for convenience, the matrix $`\widehat{f}`$ in the spin-space as
$$\left(\begin{array}{cc}f_+& 0\\ 0& f_{}\end{array}\right)$$
(142)
The diagonal form of the matrix is a consequence of the uniformity of the exchange field $`h`$. The matrix $`\widehat{g}`$ has the same form.
In order to find the function $`g_3`$ that determines the magnetization, we have to solve the Usadel equation (178) in the $`F`$ and $`S`$ region and to match the corresponding solutions with the help of the boundary conditions (181).
The simplest case when the Usadel equation can be solved analytically is the case of a thin $`F`$ layer. We suppose that the thickness $`d_F`$ of the $`F`$ layer is small compared with the characteristic length $`\xi _F`$ of the condensate penetration into the ferromagnet (this condition is fulfilled in the experiments by Garifullin *et al.* (2002)). In this case we can average the exact Usadel equation (178) over $`x`$ in the $`F`$ layer assuming that the Green’s functions are almost constant in space. In addition, provided the ratio $`\sigma _F/\sigma _S`$ is small enough, the Green’s functions in the superconductor are close to the bulk values $`f_{BCS}`$ and $`g_{BCS}`$. This allows us to linearize the Usadel equation in the superconductor. The component of the Green’s function in $`S`$ that enters the expression for the magnetization can be obtained from the boundary condition (181) and is given by
$$g_{S3}(x)=\frac{1}{\gamma _S\kappa _S}\left(g_{BCS}f_{F0}+f_{BCS}g_{F3}\right)e^{\kappa _sx},$$
(143)
where $`\kappa _S^2=2\sqrt{\omega ^2+\mathrm{\Delta }^2}/D_S`$, $`f_{F0}=(f_++f_{})/2`$, $`g_{F3}=(g_+g_{})/2`$ and $`g_\pm `$ and $`f_\pm `$ are the components of the matrices $`\widehat{g}`$ and $`\widehat{f}`$. They are defined as
$$g_{F\pm }=\stackrel{~}{\omega }_\pm /\zeta _{\omega \pm },\text{ }f_{F\pm }=\pm ϵ_{bF}f_{BCS}/\zeta _{\omega \pm },$$
(144)
where $`\stackrel{~}{\omega }_\pm =\omega +ϵ_{bF}g_{BCS}ih`$, $`\zeta _{\omega \pm }=\sqrt{\stackrel{~}{\omega }_\pm ^2(ϵ_{bF}f_{BCS})^2}`$, $`ϵ_{bF}=D_F/(2\gamma _Fd_F)`$. The magnetization variation is determined by the expression
$$\delta M=i\pi \nu T\underset{\omega =\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{Tr}\text{ }(\widehat{g}\widehat{\sigma }_3),$$
(145)
Using Eqs. (143-145) for $`\mathrm{Tr}(\widehat{g}\widehat{\sigma }_3)/2g_3=(g_+g_{})/2,`$ one can easily calculate $`\delta M`$. In FIG. 25 we show the change of the magnetization $`\delta M`$ induced in the superconductor as a function of the temperature. We see that for low enough temperatures the decrease of the magnetization can be very large. At the same time, the change of the magnetization in the ferromagnet is small Bergeret *et al.* (2004a).
It is interesting to calculate the total magnetic moment $`\delta \overline{M}_S`$ induced in the superconducting film and compare it with the total magnetization of the ferromagnet $`M_{F0}d_F`$ (as we have mentioned, the magnetization variation $`\delta M_F`$ in the ferromagnet is small and can be neglected).
The total magnetization of the superconductor is given by
$$\delta \overline{M}_S=_{d_s}^0𝑑x\delta M_S(x).$$
Assuming that $`hϵ_{bF}=D_F/(2\gamma _Fd_F)`$ or $`h[D_F/(2d_F^2)](\rho _Fd_F/R_b)`$, we can easily compute the ratio
$$\frac{\delta \overline{M}_S}{M_{F0}d_F}\pi \frac{D_S\nu _S\mathrm{\Delta }^2T}{d_F\gamma _S\nu _Fϵ_{bF}}\underset{\omega }{}\frac{1}{(\omega ^2+\mathrm{\Delta }^2)^{3/2}}=1,$$
(146)
where $`\rho _F`$ is the resistivity of the F region.
We see that in the case of a thin ferromagnet at low temperatures and a not too strong exchange field the magnetization induced in the superconductor compensates completely the magnetization in the ferromagnet. This result follows from the fact that the magnetization induced in the superconductor (it is proportional to $`g_{S3}`$) spreads over distances of the order of $`\xi _S`$. In view of this result one can expect that the magnetic moment of a small ferromagnetic particle embedded in a superconductor should be completely screened by the Cooper pairs. We discuss the screening of a ferromagnet particle by the Cooper pairs in the next subsection.
It is worth mentioning that the problem of finding the magnetization in a $`S/F`$ structure consisting of thin $`S`$ ($`d_S<\xi _S`$) and $`F`$ ($`d_F<\xi _F`$) layers is equivalent to the problem of magnetic superconductors where ferromagnetic (exchange) interaction and superconducting correlations coexist. If we assume a strong coupling between the thin $`S`$ and $`F`$ layers, we can again average the equations over the thickness of the structure and arrive at the Usadel equation for the averaged Green’s function with an effective exchange field $`\stackrel{~}{h}=hd_F/d`$ and an effective order parameter $`\stackrel{~}{\mathrm{\Delta }}=\mathrm{\Delta }d_S/d,`$ where $`d=d_S+d_F`$. In this case the magnetization is given by $`M=g\mu _B\nu \sqrt{\stackrel{~}{h}^2\stackrel{~}{\mathrm{\Delta }}^2}\mathrm{\Theta }(\stackrel{~}{h}\stackrel{~}{\mathrm{\Delta }})`$, where $`\mathrm{\Theta }\left(x\right)`$ is the step function. This means that the total magnetization $`M`$ is zero for $`\stackrel{~}{h}<\stackrel{~}{\mathrm{\Delta }}`$. This result agrees with those obtained by Karchev *et al.* (2001); Shen *et al.* (2003) who studied the problem of the coexistence of superconductivity and itinerant ferromagnetism in magnetic superconductors.
One of the assumptions made for obtaining the previous results is the quasiclassical condition $`h/ϵ_F1`$. For some materials the latter is not fulfilled and one has to go beyond the quasiclassical approach. Halterman and Valls (2002a) studied the imbalance of spin up and spin down electrons in pure $`S/F`$ structures (i.e. without impurities) in the case of strong exchange fields ($`h/ϵ_F1`$). In that case superconductivity is strongly suppressed at the $`S/F`$ interface. Solving the Bogoliubov-de Gennes equations numerically the authors showed that there was a magnetic “leakage” from the ferromagnet into the superconductor, which lead to a polarization of the electrons in $`S`$ over the short length scale $`\lambda _F`$. The direction of the induced magnetic moment in the superconductor was parallel to that in the ferromagnet, which contrasts our finding.
At the same time, the limit of a very strong exchange field considered by Halterman and Valls (2002a) differs completely from ours. It is clear that due to the strong suppression of the superconductivity at the $`S/F`$ interface, the magnetic moment cannot be influenced by the superconductivity and therefore thick ferromagnetic layers with exchange energies of the order of the Fermi energy are not suitable for observing the reduction of the magnetization described above.
The DOS for states with spin-up and spin-down electrons in a $`S/F`$ structure has been calculated on the basis of the Usadel equation by Fazio and Lucheroni (1999). The authors have found that the DOS of these states was different in the superconductor over the length of the order $`\xi _S`$. However, the change of the magnetization has not been calculated in this work.
This has been done later by Krivoruchko and Koshina (2002) for a $`S/F`$ structure. Using the Usadel equation, the authors numerically calculated the magnetization induced in the superconductor. They found that the magnetic moment leaked from the $`F`$ layer into the $`S`$ layer and changed the sign at some distance of the order of $`\xi _S`$, thus becoming negative at sufficiently large distances only. In our opinion, the “leakage” of the magnetic moment $`M_S`$ obtained in that paper is a consequence of the use by the authors of a wrong expression for the magnetic moment. They did not add to the formula obtained in the quasiclassical approximation a contribution from the energies levels located far from the Fermi energy. The latter contribution is not captured by the quasiclassical approach and should be written additionally.
We have seen that under certain conditions a finite magnetic moment is induced inside the superconductor. Does this magnetic moment affect the superconductivity? The magnetic field $`B_S`$ in the superconductor equals the magnetization $`4\pi M_s`$. The induced magnetization in the superconductor $`M_S`$ is smaller than the magnetization in the ferromagnet: $`M_S=M_F\mathrm{max}(d_F/\{\xi _S,d_S\})`$. The critical field for superconducting thin films is given by the expression $`H_c(\lambda _L/d_S)H_{bulk}`$, where $`\lambda _L`$ is the London penetration depth, and $`H_{bulk}`$ is the critical field of the bulk material. The superconductivity is not affected by the induced field $`B_S`$ if the field $`B_S4\pi M_F(d_F/\xi _S)`$ (we set $`d_S\xi _S`$) is smaller than $`H_c`$. Therefore the condition $`4\pi M_F<(\lambda _L/d_F)H_{bulk}`$ should be satisfied. If we take $`\lambda _L1\mu m`$ and $`d_F50`$Å, we arrive at the condition $`4\pi M_F<200H_{bulk}`$. This condition is fulfilled easily for the case of not too strong ferromagnets. Due to the presence of the magnetization in the ferromagnet and superconductor spontaneous currents arise in the system. The spontaneous Meissner currents induced by the magnetization in S/F structures were studied by Bergeret *et al.* (2001c); Krawiec *et al.* (2004).
The phenomenon discussed in this section can be considered as an alternative mechanism of the decrease of the total magnetic moment observed by Garifullin *et al.* (2002). In order to clarify which of these two effects is more important for the experimental observations one needs more information.
The most direct check for the cryptoferromagnetic phase would be measurements with polarized neutrons. In a recent work by Stahn *et al.* (2005), in which a multilayered $`S/F/S/F\mathrm{}`$ structure was studied. This structure consists of the high $`T_c`$ superconductor $`YBa_2Cu_3O_7`$ ($`S`$ layer) and of the ferromagnet $`La_{2/3}Ca_{1/3}MnO_3`$ ($`F`$ layer). Two samples with the $`S`$ and $`F`$ layers of the same thickness were used. Layers of sample 1(2) are $`98`$Å($`160`$Å) thick. The Curie temperature of the ferromagnet and the temperature of the superconducting transition are equal to $`165K`$ and $`75K`$ respectively. By using neutron reflectometry the authors obtained an information about the spatial distribution of the magnetic moment in the structure. Analyzing the temperature dependence of the Bragg peaks intensity they came to the conclusion that the most probable scenario to explain important features of this dependence observed was the assumption that an induced magnetization arises in the $`S`$ layers. If this explanation was correct, the sign of the induced magnetization had to be opposite to the sign of the magnetization in the $`F`$ layers. It is quite reasonable to think that the mechanism discussed above for conventional superconductors should be present also in high $`T_c`$ superconductors and then the theoretic scenario analyzed in this section can serve as an explanation of the experiment.
### V.3 Spin screening of the magnetic moment of a ferromagnetic particle in a superconductor
Let us consider now a ferromagnetic particle (grain) embedded into a superconductor (see FIG. 26). As in the previous subsection, we analyze the magnetic moment induced in the superconductor around the particle and compare it with the magnetic moment of the $`F`$ particle $`(4\pi a^3/3)M_{F0}`$ (we assume that the particle has a spherical form and radius $`a`$).
It is well known that the superconducting currents (Meissner currents) in a superconductor screen a magnetic field that decays from the surface over the London penetration length $`\lambda _L`$ and vanishes in the bulk of the superconductor. The same length characterizes the decay of the magnetic field created by a ferromagnetic $`(F)`$ grain embedded in a superconductor if the radius of the grain $`a`$ is larger than $`\lambda _L`$. However, if the radius $`a`$ is small, the Meissner effect can be neglected and a stray magnetic field around the grain should decay, as in a normal metal, over a length of the order $`a`$. We consider now just this case.
Above the critical temperature $`T_c`$ the stray magnetic field polarizes the spins of free electrons and induces a magnetic moment. This magnetic moment is very small because the Pauli paramagnetism is weak ($`\mu _B^2\nu 10^6`$). In addition, the total magnetic moment induced by the stray magnetic field is zero. The penetration depth $`\lambda _L`$ can be of the order of hundreds of interatomic distances or larger, so that if $`a`$ is smaller or of the order of $`10nm`$, the Meissner effect can be neglected.
The screening of the magnetic moment is a phenomenon specific for superconductors. It is usually believed that in a situation, when the screening due to the orbital electron motion can be neglected (small grains and thin films), the total magnetic moment is just the magnetic moment of the ferromagnetic particle and no additional magnetization is induced by the electrons of the superconductor.
This common wisdom is quite natural because in conventional superconductors the total spin of a Cooper pair is equal to zero and the polarization of the conduction electrons is even smaller than in the normal metal. Spin-orbit interactions may lead to a finite magnetic susceptibility of the superconductor but it is positive and smaller anyway than the one in the normal state Abrikosov and Gor’kov (1962); Abrikosov (1988).
Let us now take a closer look at the results of the last subsection. We have seen that the proximity effect induces in the superconductor a magnetic moment with the sign opposite to the one in the ferromagnet. In view of this result it is quite natural to expect that the magnetic moment of a small ferromagnetic particle embedded in a superconductor may be screened by the Cooper pairs as it is sketched in FIG. 26. So, let us consider this situation in more detail.
We consider a ferromagnetic grain of radius $`a`$ embedded in a bulk superconductor. If the size of the particle is smaller than the length $`\xi _F`$ we can again assume that the quasiclassical Green’s functions in the $`F`$ region are almost constant and given by Eq. (144), where now $`ϵ_{bF}=3D_F/(2\gamma _Fa)`$. In the superconductor we have to solve the linearized Usadel equation for the component $`g_{S3}`$ determining the magnetization
$$^2g_{S3}\kappa _S^2g_{S3}=0,$$
(147)
where $`^2=_{rr}+(2/r)_r`$ is the Laplace operator in spherical coordinates.
Using the boundary conditions Eq. (181) we write the solution of this equation as
$$g_{S3}=\frac{f_{BCS}}{\gamma _S}\left(g_{BCS}f_{F0}f_{BCS}g_{F3}\right)\frac{a^2}{1+\kappa _Sa}\frac{e^{\kappa _S(ra)}}{r},$$
(148)
where $`f_{F0}=(f_{F+}+f_F)/2`$ and $`g_{F3}=(g_{F+}g_F)/2`$.
We assume again that the transmission coefficient through the $`S/F`$ interface is not small and the condition $`\mathrm{\Delta }<<h(D_F/a^2)`$ is fulfilled. In this case the expression for $`g_{S3}`$ drastically simplifies. Indeed, in this limit $`g_{F3}=f_{F0}f_{BCS}/g_{BCS}`$ and $`f_{F0}ihf_{BCS}g_{BCS}/ϵ_{bF}`$. Therefore Eq.(148) acquires the form
$$𝐠_{S3}=\frac{f_{BCS}^2}{\gamma _S}\frac{a^2}{r}\frac{ih}{ϵ_{bF}}𝐞^{\kappa _S(ra)},$$
(149)
This solution can be obtained from Eq.(147) if one writes down the term $`4\pi A\delta (r)`$ on the right-hand side of this equation with $`A=f_{BCS}^2a^2ih/(\gamma _Sϵ_{bF}).`$ This means that the ferromagnetic grain acts on Cooper pairs as a magnetic impurity embedded into a dirty superconductor. It induces a ferromagnetic cloud of the size of the order $`\xi _S`$ with a magnetic moment $`\mu _B\nu hV_F.`$
In order to justify the assumptions made above we estimate the energy $`D_F/a^2`$ assuming that the mean free path is of the order of $`a`$. For $`a=30`$Å and $`v_F=10^8cm/\mathrm{sec}`$ we get $`D_F/a^21000K;`$. This condition is fulfilled for ferromagnets with the exchange energy of the order of several hundreds $`K`$.
In the limit of low temperatures the calculation of the magnetic moment becomes very easy and we obtain for the magnetic moment $`\overline{M}_S`$ induced in the superconductor the following expression
$$\frac{\overline{M}_S}{M_{F0}(4\pi a^3/3)}=1$$
(150)
This is a remarkable result which shows that the induced magnetic moment is opposite in sign to the moment of the ferromagnetic particle and their absolute values are equal to each other. In other words, the magnetic moment of the ferromagnet is completely screened by the superconductor Bergeret *et al.* (2004b). The characteristic radius of the screening is the coherence length $`\xi _S`$, which contrasts the orbital screening due to the Meissner effect characterized by the London penetration depth $`\lambda _L`$.
To avoid misunderstanding we emphasize once again that the full screening occurs only if the magnetization (per unit volume) of the ferromagnetic grain $`M_{F0}`$ is given by Eq.(139), which means that the ferromagnetic grain is an itinerant ferromagnet. If the magnetization of the ferromagnet is caused by both localized moments ($`M_{loc})`$ and itinerant electrons ($`M_{itin})`$, the full screening is not achieved. Moreover, the magnetization $`M_{loc}`$ may be larger than $`M_{itin}`$ and have opposite direction. In this case we would have an anti-screening Bergeret and García (2004).
Actually, we have discussed the diffusive case only. However, it turns out that the spin screening occurs also in the clean case provided the exchange field is not too high: $`h<<v_F/d_F`$, where $`v_F`$ and $`d_F`$ are the Fermi velocity and the thickness (radius) of the ferromagnetic film or grain Kharitonov *et al.* (2005); Bergeret *et al.* (2005).
The energy spectrum of a superconductor with a point-like classical magnetic moment was studied many years ago by Shiba (1968), Sakurai (1970) and Rusinov (1969), and more recently by Salkola *et al.* (1997). The magnetic impurity leads to a bound state $`\beta _0`$ inside the superconducting energy gap. There is some critical strength $`h_cϵ_F`$ of the exchange coupling $`h`$ that separates two different ground states of the system denoted by $`\psi `$ if $`h<h_c`$ and $`\psi ^{}`$ if $`h>h_c`$. The bound state $`\beta _0`$ corresponds to a localized quasiparticle with spin “up”<sup>6</sup><sup>6</sup>6One assumes that the magnetic impurity has spin up.. Since the total electronic spin in the state $`\psi `$ is zero one says that the continuum localizes a spin “up”. The energy needed to create a quasiparticle excitation decreases when increasing $`h`$. At $`h=h_c`$ the state $`\psi `$ becomes unstable against a spontaneous creation of an excitation with spin “up” and the transition to the state $`\psi ^{}`$ occurs. In this state the electronic spin at the impurity site is now equal to $`1/2`$. All the works considering this problem focused the attention on the subgap structure of the spectrum and did not addressed the problem of the screening of the magnetic moment by the continuum spectrum. This is of no surprise because a sufficiently large magnetic moment of the impurity ($`S1`$) cannot be screened by the quasiparticles.
### V.4 Spin-orbit interaction and its effect on the proximity effect
In this section we discuss the influence of the spin-orbit (SO) interaction on the proximity effect. Although in general its characteristic energy scale is much smaller than the exchange energy $`h`$, it can be comparable with the superconducting gap $`\mathrm{\Delta }`$ and therefore this effect can be very important. Since the SO scattering leads to a mixing of the spin channels, we expect that it will affect not only the singlet component of the condensate but also the triplet one in the ferromagnet.
In conventional superconductors the SO interaction does not affect thermodynamic properties. However, a non-vanishing magnetic susceptibility at zero temperature (Knight shift) observed in small superconducting samples and films can be explained only if the SO interaction is taken into account Abrikosov and Gor’kov (1962). In the $`F/S`$ structures considered here the exchange field $`h`$ breaks the time-reversal symmetry in analogy to the external magnetic field in the Knight shift problem. Therefore the SO interaction in the superconductor is expected to influence the inverse proximity effect studied in this Chapter.
In this Section we will generalize the analysis of the long-range proximity effect and the inverse proximity effect presented above taking the SO interaction into account. The quasiclassical equations in the presence of the SO interaction were derived by Alexander *et al.* (1985) and used for the first time for the $`F/S`$ systems by Demler *et al.* (1997).
The derivation of these equations is presented in the Appendix A. The resulting Usadel equation takes the form
$`iD_𝐫(\stackrel{ˇ}{g}_𝐫\stackrel{ˇ}{g})`$ $`+`$ $`i\left(\widehat{\tau }_3_t\stackrel{ˇ}{g}+_t^{}\stackrel{ˇ}{g}\widehat{\tau }_3\right)+[\stackrel{ˇ}{\mathrm{\Delta }},\stackrel{ˇ}{g}]+[𝐡\stackrel{ˇ}{𝐒},\stackrel{ˇ}{g}]`$ (151)
$`+`$ $`{\displaystyle \frac{i}{\tau _{s.o.}}}[\stackrel{ˇ}{𝐒}\widehat{\tau }_3\stackrel{ˇ}{g}\widehat{\tau }_3\stackrel{ˇ}{𝐒},\stackrel{ˇ}{g}]=0.`$
All symbols are defined in the Appendix A. The spin-orbit relaxation time $`\tau _{s.o.}`$ takes very different values depending on the material used in the experiments. Some estimates for the values of $`1/h\tau _{s.o.}`$can be found in Oh *et al.* (2000). For example, for transition metals like $`Fe`$ one obtains $`1/h\tau _{s.o.}10^2`$, while for a typical magnetic rare earth the value $`1/h\tau _{s.o.}0.3`$ is more typical. In the latter case the SO interaction should clearly affect the penetration of the condensate into the ferromagnet.
In order to study the influence of the SO interaction on both the long-range and the inverse proximity effect we will use Eq.(151). We consider first the well known problem of the Knight shift. This example will show the convenience of using the quasiclassical approach.
#### The Knight shift in superconductors
Since the pioneering work of Abrikosov and Gor’kov (1962) it is well established that the magnetic susceptibility of small superconducting samples is not zero due to the spin-orbit interaction. This explains the experiments performed for the first time many years ago by Androes and Knight (1961) who used the nuclear magnetic resonance technique.
Let us consider a superconductor in an external magnetic field $`H`$. In the Usadel equation, Eq. (151), the field $`H`$ plays the role of the exchange energy $`h`$. We are interested in the linear response to this filed, i.e. in the magnetic susceptibility $`\chi _S`$ of the superconductor. We assume that the superconductor is homogeneous and therefore we drop the gradient term in Eq. (151):
$`\omega [\widehat{\tau }_3,\stackrel{ˇ}{g}]+i[\stackrel{ˇ}{\mathrm{\Delta }},\stackrel{ˇ}{g}]+iH[\stackrel{ˇ}{n},\stackrel{ˇ}{g}](1/\tau _{s.o.})[\stackrel{ˇ}{𝐒}\widehat{\tau }_3\stackrel{ˇ}{g}\widehat{\tau }_3\stackrel{ˇ}{𝐒},\stackrel{ˇ}{g}]`$ $`=`$ $`0,`$ (152)
$`\stackrel{ˇ}{g}^2`$ $`=`$ $`1.`$ (153)
The solution of Eq. (152) has the form
$$\stackrel{ˇ}{g}=\left(g_{BCS}+g_3\widehat{\sigma }_3\right)\widehat{\tau }_3+\left(f_{BCS}\widehat{\sigma }_3+f_0\right)i\widehat{\tau }_2,$$
(154)
where the functions $`g_3`$ and $`f_0`$ are corrections to the normal $`g_{BCS}`$ and anomalous $`f_{BCS}`$ Green’s functions. In the particle-hole space the matrix $`\stackrel{ˇ}{g}`$ has the usual form, i.e. it is expanded in matrices $`\widehat{\tau }_3`$ and $`i\widehat{\tau }_2`$. In the spin space the triplet component (the $`g_3`$ and $`f_0`$ terms) appears due to the magnetic field acting on the spins. Using Eqs. (152-154) one can readily obtain
$$g_3=i\frac{\mathrm{\Delta }^2H}{E_\omega ^2(E_\omega +4/\tau _{s.o.})}.$$
(155)
where $`E_\omega =\sqrt{\mathrm{\Delta }^2+\omega ^2}`$.
Substituting Eq. (155) into Eq. (145) we can write the magnetization $`M`$ as follows
$$M=M_0\mu _B\nu \left(2\pi T\mathrm{\Delta }^2\underset{\omega }{}\frac{1}{E_\omega ^2(E_\omega +4/\tau _{s.o.})}\right)H$$
(156)
The first term in Eq. (156) cannot be calculated in the framework of the quasiclassical theory and one should use exact Green’s functions. It corresponds to the Pauli paramagnetic term given by $`M_0=\mu _B\nu H`$. In the quasiclassical approach this term is absent. This term does not depend on temperature on the energy scale of the order of $`T_c`$ and originates from a contribution of short distances where the quasiclassical approximation fails.
This situation is rather typical for the quasiclassical approach and one usually adds to formulae obtained within this approach contributions coming from short distances or times by hand (see, for example, Artemenko and Volkov (1980); Kopnin (2001); Rammer and Smith (1986)). Eq. (156) was first obtained by Abrikosov and Gor’kov (1962).
In the absence of the spin orbit interaction the magnetization at $`T=0`$ is, as expected, equal to zero. However, if the SO interaction is finite the spin susceptibility $`\chi _S`$ does not vanish at $`T=0`$. It is interesting that, as follows from Eq. (152), the singlet component of the condensate is not affected by the SO interaction. The origin of the finite susceptibility is the existence of the triplet component $`f_0`$ of the condensate.
In the $`S/F`$ structures there is no exchange field in the superconductor and therefore the situation is in principle different. However, we have seen that due to the proximity effect the triplet component $`f_0`$ is induced in the superconductor.
From the above analysis one expects that the SO interaction may affect the penetration length of such component in the superconductor. In the next sections we consider the influence of the SO the superconducting condensate in both the ferromagnet and the superconductor.
#### V.4.1 Influence of the Spin-Orbit interaction on the long-range Proximity Effect
Now we consider again the $`S/F/S/F/S`$ structure of Section IV.1 and assume that the long-range triplet component is created, which is possible provided the angle $`\alpha `$ between the magnetizations differs from $`0`$ and $`\pi `$. In order to understand how the SO interaction affects the triplet component it is convenient to linearize Eq. (151) in the $`F`$-layer assuming, for example, that the proximity effect is weak. One can easily obtain a linearized equation similar to Eq. (52) for the condensate function $`\stackrel{ˇ}{f}`$. The solution of this equation is represented again in the form
$$\stackrel{ˇ}{f}(x)=i\widehat{\tau }_2(f_0(x)\widehat{\sigma }_0+f_3(x)\widehat{\sigma }_3)+i\widehat{\tau }_1f_1(x)\widehat{\sigma }_1.$$
(157)
The functions $`f_i(x)`$ are given as before by $`f_i(x)=_jb_j\mathrm{exp}[\kappa _jx]`$ but now the new eigenvalues $`\kappa _j`$ are written as
$`\kappa _\pm ^2`$ $`=`$ $`\pm {\displaystyle \frac{2i}{D_F}}\sqrt{h^2\left({\displaystyle \frac{4}{\tau _{so}}}\right)^2}+{\displaystyle \frac{4}{\tau _{so}D_F}}`$ (158)
$`\kappa _0^2`$ $`=`$ $`\kappa _\omega ^2+2\left({\displaystyle \frac{4}{\tau _{s.o.}D_F}}\right).`$ (159)
We see from these equations that both the singlet and triplet components are affected by the spin-orbit interaction making the decay of the condensate in the ferromagnet faster. In the limiting case, when $`4/\tau _{so}>h,T_c`$, both the components penetrate over the same distance $`\xi _{s.o.}=\sqrt{\tau _{so}D_F}`$ and therefore the long-range effect is suppressed. In this case the characteristic oscillations of the singlet component are destroyed Demler *et al.* (1997). In the more interesting case $`4/\tau _{so}T_c<h`$, the singlet component does not change and penetrates over the short distance $`\xi _F`$ . At the same time, the triplet component is more sensitive to the spin-orbit interaction and the penetration length equals $`\mathrm{min}`$($`\xi _{so},\xi _T`$)$`>\xi _F`$.
Therefore, if the spin-orbit interaction is not very strong, the penetration of triplet condensate over the long distances discussed in the preceding sections is still possible, although the penetration length is reduced.
#### V.4.2 Spin-Orbit Interaction and the Inverse Proximity Effect
Studying a $`S/F`$ bilayer we have seen that the induced magnetic moment in the superconductor $`S`$ is related to the appearance of the triplet component $`f_0`$. Moreover, we have shown that this component is affected by the SO interaction, while the singlet one $`f_3`$ is not. So, one should expect that the SO interaction may change the scale over which the magnetic moment is induced in the superconductor and one can estimate easily this length.
Assuming that the Green’s functions in the superconductor take values close to the bulk values we linearize the Usadel equation (151) in the superconductor. The solution has the same form as before, Eq. (143), but $`\kappa _S`$ should be replaced by
$$\kappa _S^2\kappa _S^2+\kappa _{so}^2,$$
(160)
where $`\kappa _{so}^2=8D_S/\tau _{so}`$. Therefore, the length of the penetration of $`g_{S3}`$ and, in its turn, of $`M_S`$ into the $`S`$ region decreases if $`\kappa _S^2\xi _S^2<\kappa _{so}^2`$.
In principle, one can measure the spatial distribution of the magnetic moment in the $`S`$ region as it was done by Luetkens *et al.* (2003) by means of muon spin rotation and get an information about the SO interaction in superconductors. As Eq. (160) shows, this would be an alternative method to measure the strength of the SO interaction in superconductors, complementary to the measurement of the Knight shift Androes and Knight (1961).
## VI Discussion of the results and outlook
In this review we have discussed new unusual properties of structures consisting of conventional superconductors in a contact with ferromagnets. It has been known that such systems might exhibit very interesting properties like a non-monotonous reduction of the superconducting temperature as a function of the thickness of the superconductor, possibility of a $`\pi `$-contact in Josephson junctions with ferromagnetic layers, etc.
However, as we have seen, everything is even more interesting and some spectacular phenomena are possible that even might look at first glance as a paradox. The common feature of the effects discussed in this review is that almost all of them originate in situations when the exchange field is not homogeneous. As a consequence of the inhomogeneity, the spin structure of the superconducting condensate function becomes very non-trivial and, in particular, the triplet components are generated. In the presence of the inhomogeneous exchange field, the total spin of a Cooper pair is not necessarily equal to zero and the total spin equal to unity with all projections onto the direction of the exchange field is possible.
We have discussed the main properties of the odd triplet superconductivity in the $`S/F`$ structures. This superconductivity differs from the well known types of superconductivity: a) singlet superconductivity with the s-wave (conventional $`T_c`$ superconductors) and d-wave (high $`T_c`$ superconductors) types of pairing; b) odd in momentum $`p`$ and even in frequency $`\omega `$ triplet superconductivity observed, e.g., in $`Sr_2RuO_4`$.
The odd triplet superconductivity discussed in this Review has a condensate (Gor’kov) function that is an odd function of the Matsubara frequency $`\omega `$ and an even function (in the main approximation) of the momentum $`p`$ in the diffusive limit. It is insensitive to the scattering on nonmagnetic impurities and therefore may be realized in thin film $`S/F`$ structures where the mean free path is very short.
For the first time, the condensate function of this type has been suggested by Berezinskii (1975) many years ago as a possible candidate to describe superfluidity in $`He^3`$. Later, it has been established that the superfluid condensate in $`He^3`$ had a different structure - it was odd in $`p`$ and even in $`\omega `$. In principle, there is an important difference between the triplet superconductivity discussed here and that predicted by Berezinskii who assumed that the order parameter $`\mathrm{\Delta }`$ was also an odd function of $`\omega `$. In our case the order parameter $`\mathrm{\Delta }`$ is determined by the singlet, s-wave condensate function and has the ordinary BCS structure (i.e., it does not depend on the momentum $`p`$ and frequency $`\omega `$). On the other hand the structure of the triplet condensate function $`\stackrel{ˇ}{f}`$ in the diffusive case considered here coincides with that suggested by Berezinskii: it is an odd function of the Matsubara frequency $`\omega `$ and, in the main approximation, is constant in the momentum space. The antisymmetric part of $`\stackrel{ˇ}{f}`$ is small compared with the symmetric part, being odd in $`p`$ and even in $`\omega .`$
The triplet component with the projection of the total spin $`S_z=\pm 1`$ penetrates the ferromagnet over a long distance of the order of $`\xi _N\sqrt{D_F/2\pi T},`$ which shows that the exchange field does not affect the triplet part of the condensate. At the same time, the exchange field suppresses the amplitude of the singlet component at the $`S/F`$ interface that determines the amplitude of the triplet component. The long-range triplet component arises only in the case of a nonhomogeneous magnetization. The triplet component appears also in a system with a homogeneous magnetization but in this case it corresponds to the projection $`S_z=0`$ and penetrates the ferromagnet over a short length $`\xi _F=\sqrt{D_F/h}<<\xi _N`$.
The triplet component exists also in magnetic superconductors Bulaevskii *et al.* (1985); Kulic and Kulic (2001) with a spiral magnetic structure. However, it always coexists with the singlet component and cannot be separated from it. In contrast, in the multilayered $`S/F`$ structures with a nonhomogeneous magnetization and with the thickness of the $`F`$ layers $`d_F`$ exceeding $`\xi _F`$, the Josephson coupling between $`S`$ layers is realized only through the long-range triplet component and this separates the singlet and triplet components from each other. As a result, the “real” odd triplet superconductivity may be realized in the transverse direction in such structures.
Another interesting peculiarity of the $`S/F`$ structures is the inverse proximity effect, namely, the penetration of the magnetic order parameter (spontaneous magnetic moment $`M`$) into the superconductor and a spatial variation of the magnetization direction in the ferromagnet under the influence of the superconductivity. It turns out that both effects are possible. A homogeneous distribution of the magnetization $`M_F`$ in the $`S/F`$ bilayer structures may be energetically unfavorable in $`F`$ even in a one-domain case resulting in a nonhomogeneous distribution of $`𝐌_F`$ in the ferromagnet.
Moreover, the magnetic moment penetrates the superconductor (induced ferromagnetism) changing sign at the $`S/F`$ interface. Therefore the total magnetic moment of the system is reduced. Under some condition the full spin screening of $`M_F`$ occurs. For example, at zero temperature the itinerant magnetic moment of a ferromagnetic grain embedded into a superconductor is completely screened by spins of the Cooper pairs in $`S`$. The radius of the screening cloud is of the order of the superconducting coherence length $`\xi _S`$. If the magnetization vector $`𝐌_F`$ is oriented in the opposite direction to the ferromagnetic exchange field $`𝐡`$, the anti-screening is possible.
As concerns the experimental situation, certainly there are indications in favor of the long-range triplet component, although an unambiguous evidence does not exist so far. For example, the resistance of ferromagnetic films or wires in the $`S/F`$ structures changes on distances that exceed the length of the decay of the singlet component $`\xi _h`$ Aumentado and Chandrasekhar (2001); Petrashov *et al.* (1999); Giroud *et al.* (1998). A possible reason for this long-range proximity effect in the $`S/F`$ systems is the long-range penetration of the triplet component. However a simpler effect might also be the reason for this long-range proximity effect. It is related to a rearrangement of a domain structure in the ferromagnet when the temperature lowers below $`T_c`$. The Meissner currents induced in the superconductor by a stray magnetic field affect the domain structure, and the resistance of the ferromagnet may change Dubonos *et al.* (2002). At the same time, the Meissner currents should be considerably reduced in an one dimensional geometry for the ferromagnet like that used in Giroud *et al.* (1998) and the explanation in terms of the long range penetration of the triplet component are more probable here.
Sefrioui *et al.* (2003) also obtained some indications on the existence of a triplet component in a multilayered $`S/F/S/F`$… structure. The samples used by Sefrioui *et al.* contained the high $`T_c`$ material $`YBa_2Cu_3O_7`$ (as a superconductor) and the half-metallic ferromagnet $`La_{0.7}Ca_{0.3}MnO_3`$ (as a ferromagnet). They found that superconductivity persisted even in the case when the thickness of the $`F`$ layers $`d_F`$ essentially exceeded $`\xi _F`$ ($`d_F10nm`$ and $`\xi _F5nm`$). In a half-metal ferromagnet with spins of free electrons aligned in one direction the singlet Cooper pairs cannot exist. Therefore it is reasonable to assume that the superconducting coupling between neighboring $`S`$ layers is realized via the triplet component Volkov *et al.* (2003); Eschrig *et al.* (2003).
A reduction of the magnetic moment of the $`S/F`$ structures due to superconducting correlations has been observed already Garifullin *et al.* (2002). This reduction may be caused both by the spin screening of the magnetic moment $`M_F`$ and by the rotation of $`M_F`$ in space Bergeret *et al.* (2000); Bergeret *et al.* (2004a). Perhaps, the spin screening can be observed directly by probing the spatial distribution of the magnetic field (or magnetic moment $`M`$) with the aid of the muon spin rotation technique Luetkens *et al.* (2003). The variation of the magnetic moment $`M`$ occurs on a macroscopic length $`\xi _S`$ and therefore can be detected.
An evidence in favor of the inverse proximity effect has also been obtained in another experimental work Stahn *et al.* (2005). Analyzing data of neutron reflectometry on a multilayered $`YBa_2Cu_3O_7/La_{2/3}Ca_{1/3}MnO_3`$ structure, the authors concluded that a magnetic moment was induced in the superconducting $`YBa_2Cu_3O_7`$ layers. The sign of this induced moment was opposite to the sign of the magnetic moment in the ferromagnetic $`La_{2/3}Ca_{1/3}MnO_3`$ layers, which correlates with our prediction.
In spite of these experimental results that may be considered as, at least preliminary, confirmation of the existence of the triplet component in the $`S/F`$ structures, there is a need in additional experimental studies of the unconventional superconductivity discussed in this review. One of the important issues would be to understand whether the long range proximity effects already observed experimentally are due to the triplet pairing or to a simple redistribution of the domain walls by the Meissner currents. We believe that measurements on thin ferromagnetic wires where the Meissner currents are reduced may clarify the situation.
It is very interesting to distinguish between the two possible inverse proximity effects experimentally. Although both the formation of the cryptoferromagnetic state and the induction of the magnetic moments in the superconductors are very interesting effects, it is not clear yet which of these effects causes the magnetization reduction observed by Mühge *et al.* (1998); Garifullin *et al.* (2002).
The enhancement of the Josephson current by the presence of the ferromagnet near the junction is one more theoretical prediction that has not been observed yet but, certainly, this effect deserves an attention. An overview for experimentalists interested in all these subjects is presented in Appendix B, where we discuss briefly different experiments on S/F structures, focusing our attention on the materials for which, we expect, the main effects discussed in this review may be observed.
In addition, further theoretical investigations are needed. The odd triplet component has been studied mainly in the diffusive limit ($`h\tau <<1`$). It would be interesting to investigate the properties of the triplet component for an arbitrary impurity concentration ($`h\tau 1`$). No theoretical work on dynamics of magnetic moments in the $`S/F`$ structures has been performed yet, although the triplet component may play a very important role in the dynamics of the $`S/F`$ structures. Transport properties of the $`S/F`$ structures require also further theoretical considerations. It would be useful to study the influence of domain structures on properties of the $`S/F`$ structures, etc. In other words, physics of the proximity effects in the superconductor-ferromagnet structures is evolving into a very popular field of research, both experimentally and theoretically.
The study of the proximity effect in $`S/F`$ structures may be extended to include ferromagnets in contact with high temperature superconductors. Although some experiments have been done already Sefrioui *et al.* (2003); Stahn *et al.* (2005), one can expect much more broad experimental investigations in the future. The modern technique allows the preparation of multilayered $`S/F/S/F`$.. structures consisting of thin ferromagnetic layers (as $`La_{2/3}Ca_{1/3}MnO_3`$) and thin layers of high $`T_c`$ superconductor (as $`YBa_2Cu_3O_7`$) with variable thicknesses. It would be very interesting to study, both experimentally and theoretically, such a system with non-collinear magnetization orientations. In this case d-wave singlet and odd triplet superconductivity should coexist in the system. It is known that many properties of the ordinary BCS superconductivity remain unchanged in the high $`T_c`$ superconductors. This means that many effects considered in this review can also occur in S/F structures containing high $`T_c`$ materials, but there will certainly be differences with respect to the conventional superconductors with the s-pairing.
We hope that this review will encourage experimentalists and theoreticians to make further investigations in this fascinating field of research.
## Acknowledgements
We appreciate fruitful discussions with A. I. Buzdin, Y. V. Fominov, A. A. Golubov, I. A. Garifullin, A. Gerber, A. Palevski, L. R. Tagirov, K. Westerholt and H. Zabel. We would like to thank SFB 491 for financial support. F.S.B would like to thank the E.U. network DIENOW for financial support.
## Appendix A Basic equations
Throughout this review we use mainly the well established method of quasiclassical Green’s functions. Within this method the Gor’kov equations can be drastically simplified by integrating the Green’s function over the momentum. This method was first introduced by Larkin and Ovchinnikov (1968) and Eilenberger (1968), and then extended by Usadel (1970) for a dirty case and by Eliashberg (1971) for a non-equilibrium case . The method of the quasiclassical Green’s functions is discussed in many reviews Serene and Reiner (1983); Rammer and Smith (1986); Larkin and Ovchinnikov (1984); Belzig *et al.* (1999) and in the book by Kopnin (2001). In this Appendix we present a brief derivation of equations for the quasiclassical Green’s functions and write formulae for the main observable quantities in terms of these functions. A special attention will be paid to the dependence of these functions on the spin variables that play a crucial role in the S/F structures. In particular, we take into account the spin-orbit interaction alongside with the exchange interaction in the ferromagnet.
We start with a general Hamiltonian describing a conventional BCS-superconductor/ferromagnet structure:
$`\widehat{H}`$ $`=`$ $`{\displaystyle \underset{\{p,s\}}{}}\{a_{sp}^+[((\xi _p\delta _{pp^{}}+eV)+U_{imp})\delta _{ss^{}}+U_{s.o.}`$ (161)
$`(𝐡.\sigma )]a_{s^{}p^{}}(\mathrm{\Delta }a_{\overline{s}\overline{p}}^{}a_{s^{}p^{}}^{}+c.c.)\}.`$
The summation is carried out over all momenta $`(p,p^{})`$ and spins $`(s,s^{})`$ (the notation $`\overline{s}`$, $`\overline{p}`$ means inversion of both spin and momentum), $`\xi _p=p^2/2mϵ_F`$ is the kinetic energy counted from the Fermi energy $`ϵ_F`$, $`V`$ is a smoothly varying electric potential. The superconducting order parameter $`\mathrm{\Delta }`$ must be determined self-consistently. It vanishes in the ferromagnetic regions. The potential $`U_{imp}=U(pp^{})`$ describes the interaction of the electrons with nonmagnetic impurities, and $`U_{s.o.}`$ describes a possible spin-orbit interaction Abrikosov and Gor’kov (1962):
$$U_{s.o.}=\underset{i}{}\frac{u_{s.o.}^{(i)}}{p_F^2}\left(𝐩\times 𝐩^{}\right)\sigma .$$
Here the summation is performed over all impurities.
The representation of the Hamiltonian in the form (161) implies that we use the mean-field approximation for the superconducting ($`\mathrm{\Delta }`$) and magnetic ($`𝐡`$) order parameter. The exchange field $`𝐡`$ is parallel to the magnetization $`𝐌_F`$ in $`F`$<sup>7</sup><sup>7</sup>7(we remind that the exchange field $`h`$ is measured in energy units, see also the Footnote on page 78). In strong ferromagnets the magnitude of $`𝐡`$ is much higher than $`\mathrm{\Delta }`$ and corresponds to an effective magnetic field $`H_{exc}=h/\mu _B`$ of the order $`10^6Oe`$ (where $`\mu _B=g\mu _{Bohr},`$ $`g`$ is the $`g`$-factor and $`\mu _{Bohr}`$ is the Bohr magneton).
In order to describe the ferromagnetic region we use a simplified model that catches all physics we are interested in. Ferromagnetism in metals is caused by the electron-electron interaction between electrons belonging to different bands that can correspond to localized and conducting states. Only the latter participate in the proximity effect. If the contribution of free electrons strongly dominates (an itinerant ferromagnet), one has $`M_FM_e`$ and the exchange energy is caused mainly by free electrons.
If the polarization of the conduction electrons is due to the interaction with localized magnetic moments, the Hamiltonian $`\widehat{H}_F`$ can be written in the form
$$\widehat{H}_F=h_1\underset{\{p,s\}}{}\left\{a_{sp}^+𝐒\sigma _{ss^{}}a_{s^{}p^{}}\right\}$$
(162)
where $`𝐒=_a𝐒_a\delta (rr_a)`$, $`𝐒_a`$ is the spin of a particular ion. A constant $`h_1`$ is related to $`h`$ via the equation: $`h=h_1n_MS_0`$ , where $`n_M`$ is the concentration of magnetic ions and $`S_0`$ is a maximum value of $`S_a`$ (we consider these spins as classical vectors; see Ref. Gor’kov and Rusinov (1963)). In this case the magnetization is a sum: $`M=M_{loc}+M_e`$, and the magnetization $`M_e`$ can be aligned parallel ($`h_1>0`$, the ferromagnetic type of the exchange field) to $`M`$ or antiparallel ($`h_1<0`$, the antiferromagnetic type of the exchange field). In the following we will assume a ferromagnetic exchange interaction ($`M_e`$ and $`M`$ are oriented in the same direction$`)`$. In principle, one can add to Eq. (162) the term $`_{\{a,b\}}\left\{𝐒_a𝐒_b\right\}`$ describing a direct interaction between localized magnetic moments but in the most part of the review this term is not important except Section V.A. where the cryptoferromagnetic state is discussed.
Starting from the Hamiltonian (161) and using a standard approach Larkin and Ovchinnikov (1984), one can derive the Eilenberger and Usadel equations. Initially these equations have been derived for $`2\times 2`$ matrix Green’s functions $`g_{n,n^{}}`$, where indices $`n,n^{}`$ relate to normal $`(g_{11},g_{22})`$ and anomalous or condensate $`(f_{12},f_{21})`$ Green’s functions. These functions describe the singlet component. In the case of a non-homogenous magnetization considered in this review one has to introduce additional Green’s functions depending on spins and describe not only the singlet but also the triplet component. These matrices depend not only on $`n,n^{}`$ indices but also on the spin indices $`s,s^{},`$ and are $`4\times 4`$ matrices in the spin and Gor’kov space (sometimes the $`n,n^{}`$ space is called the Nambu or Nambu-Gor’kov space).
In order to define the Green’s functions in a customary way it is convenient to write the Hamiltonian (161) in terms of new operators $`c_{nsp}^{}`$ and $`c_{nsp}`$ that are related to the creation and anhilation operators $`a_s^+`$ and $`a_s`$ by the relation (we drop the index $`p`$ related to the momentum)
$$c_{ns}=\{\begin{array}{c}a_s,n=1\\ a_{\overline{s}}^{},n=2.\end{array}$$
(163)
These operators (for $`s=1`$) were introduced by Nambu Nambu (1960) . The new operators allow one to express the anomalous averages $`<a_{}a_{}>`$ introduced by Gor’kov as the conventional averages $`<c_1c_2^+>`$ and therefore the theory of superconductivity can be constructed by analogy with a theory of normal systems. Thus, the index $`n`$ operates in the particle-hole (Numbs-Gor’kov) space, while the index $`s`$ operates in the spin space. In terms of the $`c_{ns}`$ operators the Hamiltonian can be written in the form
$$H=\underset{\{p,n,s\}}{}c_{ns}^+_{(nn^{})(ss^{})}c_{n^{}s^{}},$$
(164)
where the summation is performed over all momenta, particle-hole and spin indices. The matrix $`\stackrel{ˇ}{}`$ is given by
$`\stackrel{ˇ}{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{[(\xi _p\delta _{pp^{}}+eV)+U_{imp}]\widehat{\tau }_3\widehat{\sigma }_0+\stackrel{~}{\widehat{\mathrm{\Delta }}}\widehat{\sigma }_3𝐡\widehat{\tau }_3\stackrel{ˇ}{𝐒}`$ (165)
$`+{\displaystyle \underset{i}{}}{\displaystyle \frac{u_{s.o.}^{(i)}}{p_F^2}}(𝐩\times 𝐩^{})\stackrel{ˇ}{𝐒}\}.`$
The matrices $`\widehat{\tau }_i`$ and $`\widehat{\sigma }_i`$ are the Pauli matrices in the particle-hole and spin space respectively; $`i=0,1,2,3`$, where $`\widehat{\tau }_0`$ and $`\sigma _0`$ are the corresponding unit matrices. The matrix vector $`\stackrel{ˇ}{𝐒}`$ is defined as
$$\stackrel{ˇ}{𝐒}=(\widehat{\sigma }_1,\widehat{\sigma }_2,\widehat{\tau }_3\widehat{\sigma }_3),$$
and the matrix order parameter equals $`\stackrel{~}{\widehat{\mathrm{\Delta }}}=\widehat{\tau }_1Re\mathrm{\Delta }\widehat{\tau }_2Im\mathrm{\Delta }`$. Now we can define the matrix Green’s functions (in the particle-hole$``$spin space) in the Keldysh representation in a standard way
$$\stackrel{ˇ}{G}(t_i,t_k^{})=\frac{1}{i}T_C\left(c_{ns}(t_i)c_{n^{}s^{}}^{}(t_k^{})\right),$$
(166)
where the temporal indices take the values $`1`$ and $`2`$, which correspond to the upper and lower branch of the contour $`C`$, running from $`\mathrm{}`$ to $`+\mathrm{}`$ and back to $`\mathrm{}`$.
One can introduce a matrix in the Keldysh space of the form
$$\stackrel{ˇ}{𝐆}(t,t^{})=\left(\begin{array}{cc}\stackrel{ˇ}{G}(t,t^{})^R& \stackrel{ˇ}{G}(t,t^{})^K\\ 0& \stackrel{ˇ}{G}(t,t^{})^A\end{array}\right)$$
(167)
where the retarded (advanced) Green’s functions $`\stackrel{ˇ}{G}(t,t^{})^{R(A)}`$ are related to the matrices $`\stackrel{ˇ}{G}(t_i,t_k^{}):`$ $`\stackrel{ˇ}{G}(t,t^{})^{R(A)}=\stackrel{ˇ}{G}(t_1,t_1^{})\stackrel{ˇ}{G}(t_{1(2)},t_{2(1)}^{})`$. All these elements are $`4\times 4`$ matrices. These functions determine thermodynamic properties of the system (density of states, the Josephson current etc). The matrix $`\stackrel{ˇ}{G}(t,t^{})^K=`$ $`\stackrel{ˇ}{G}(t_1,t_2^{})+\stackrel{ˇ}{G}(t_2,t_1^{})`$ is related to the distribution function and has a nontrivial structure only in a nonequilibrium case. In the equilibrium case it is equal to: $`\stackrel{ˇ}{G}(ϵ)^K=d(tt^{})\stackrel{ˇ}{G}(tt^{})^K\mathrm{exp}(iϵ(tt^{}))=[\stackrel{ˇ}{G}(ϵ)^R\stackrel{ˇ}{G}(ϵ)^A]\mathrm{tanh}(ϵ/2T).`$
In order to obtain the equations for the quasiclassical Green’s functions, we follow the procedure introduced by Larkin and Ovchinnikov (1984). The equation of motion for the Green’s functions is
$$\left(i_t\stackrel{ˇ}{H}\stackrel{ˇ}{\mathrm{\Sigma }}_{imp}\stackrel{ˇ}{\mathrm{\Sigma }}_{s.o.}\right)\stackrel{ˇ}{𝐆}=\stackrel{ˇ}{1},$$
(168)
where
$$\stackrel{ˇ}{H}=\widehat{\tau }_3\frac{_𝐫^2}{2m}ϵ_F𝐡\widehat{\tau }_3\stackrel{ˇ}{𝐒}+\stackrel{~}{\widehat{\mathrm{\Delta }}}\widehat{\sigma }_3$$
and $`\stackrel{ˇ}{\mathrm{\Sigma }}_{imp}`$ and $`\stackrel{ˇ}{\mathrm{\Sigma }}_{s.o.}`$ are the self-energies given in the Born approximation by
$`\stackrel{ˇ}{\mathrm{\Sigma }}_{imp}`$ $`=`$ $`N_{imp}u_{imp}^2\widehat{\tau }_3\stackrel{ˇ}{𝐆}\widehat{\tau }_3,\stackrel{ˇ}{𝐆}=\nu {\displaystyle 𝑑\xi _p\frac{d\mathrm{\Omega }}{4\pi }\stackrel{ˇ}{𝐆}}`$
$`\stackrel{ˇ}{\mathrm{\Sigma }}_{s.o.}`$ $`=`$ $`N_{imp}u_{s.o.}^2\stackrel{ˇ}{𝐆}_{s.o.},`$
$`\stackrel{ˇ}{𝐆}_{s.o.}`$ $`=`$ $`\nu {\displaystyle 𝑑\xi _p\frac{d\mathrm{\Omega }}{4\pi }(𝐧\times 𝐧^{})\stackrel{ˇ}{𝐒}\stackrel{ˇ}{𝐆}\stackrel{ˇ}{𝐒}(𝐧\times 𝐧^{})}.`$ (169)
Here $`N_{imp}`$ is the impurity concentration, $`\nu `$ is the density of states at the Fermi level and $`𝐧`$ is a unit vector parallel to the momentum.
Next step is to subtract from Eq.(168), multiplied by $`\widehat{\tau }_3`$ from the left, its conjugate equation multiplied by $`\widehat{\tau }_3`$ from the right. Then one has to go from the variables ($`𝐫,𝐫^{}`$) to ($`(𝐫+𝐫^{})/\mathrm{𝟐},𝐫𝐫^{}`$) and to perform a Fourier transformation with respect to the relative coordinate. By making use of the fact that the Green’s functions are peaked at the Fermi surface, one can integrate the resulting equation over $`\xi _p,`$ and finally one obtains
$`\widehat{\tau }_3_t\stackrel{ˇ}{g}+_t^{}\stackrel{ˇ}{g}\widehat{\tau }_3`$ $`+`$ $`𝐯_𝐅\stackrel{ˇ}{g}i[𝐡\stackrel{ˇ}{𝐒},\stackrel{ˇ}{g}]i[\stackrel{ˇ}{\mathrm{\Delta }},\stackrel{ˇ}{g}]`$
$`+{\displaystyle \frac{1}{2\tau }}[\stackrel{ˇ}{g},\stackrel{ˇ}{g}]`$ $`+`$ $`{\displaystyle \frac{1}{2\tau _{s.o.}}}[\widehat{\tau }_3\stackrel{ˇ}{g}_{s.o.}\widehat{\tau }_3,\stackrel{ˇ}{g}]=0`$ (170)
where $`\stackrel{ˇ}{\mathrm{\Delta }}=\widehat{\tau }_3\stackrel{~}{\widehat{\mathrm{\Delta }}}`$ and the quasiclassical Green’s functions $`\stackrel{ˇ}{g}(t_i,t_k^{})`$ are defined as
$$\stackrel{ˇ}{g}(𝐩_𝐅,𝐫)=\frac{i}{\pi }\left(\widehat{\tau }_3\widehat{\sigma }_0\right)𝑑\xi _p\stackrel{ˇ}{𝐆}(t_i,t_k^{};𝐩,𝐫),$$
(171)
and $`𝐯_𝐅`$ is the Fermi velocity. The scattering times appearing in Eq. (170) are defined as
$`\tau ^1`$ $`=`$ $`2\pi \nu N_{imp}u_{imp}^2`$ (172)
$`\tau _{s.o.}^1`$ $`=`$ $`{\displaystyle \frac{1}{3}}\pi \nu N_{imp}{\displaystyle \frac{d\mathrm{\Omega }}{4\pi }u_{s.o.}^2\mathrm{sin}^2\theta }`$ (173)
Eq. (170) is a generalization of an equation derived by Larkin and Ovchinnikov (1968), and Eilenberger (1968) for a general nonequilibrium case. This generalization (in the absence of spin-dependent interactions) has been done by Eliashberg (1971) and Larkin and Ovchinnikov Larkin and Ovchinnikov (1984). A solution for Eq.(170) is not unique. The proper solutions must obey the so called normalization condition
$$(dϵ_1/2\pi )\stackrel{ˇ}{g}(𝐩_𝐅,𝐫;ϵ,ϵ_1).\stackrel{ˇ}{g}(𝐩_𝐅,𝐫;ϵ_1,ϵ^{})=1$$
(174)
Generalization for the case of exchange and spin-orbit interaction was presented in Bergeret *et al.* (2000) and Bergeret *et al.* (2001c). The solution for Eq.(170) can be obtained in some limiting cases, for example, in a homogeneous case. However finding its solution for nonhomogeneous structures with an arbitrary impurity concentration may be a quite difficult task. Further simplifications can be made in the case of a dirty superconductor when the energy $`\tau ^1`$ related to the elastic scattering by nonmagnetic impurities is larger than all other energies involved in the problem, and the mean free path $`l`$ is smaller than all characteristic lengths (except the Fermi wave length that is set in the quasiclassical theory to zero). In this case one can expand the solution of Eq. (170) in terms of spherical harmonics and retain only the first two of them, i.e.
$$\stackrel{ˇ}{g}(𝐩_𝐅,𝐫;)=\stackrel{ˇ}{g}_s(𝐫)+(𝐩_F/p_F)\stackrel{ˇ}{𝐠}_a(𝐫),$$
(175)
where $`\stackrel{ˇ}{g}_s(𝐫)`$ is a matrix that depends only on coordinates. The second term is the antisymmetric part (the first Legendre polynomial) that determines the current. It is assumed that the second term is smaller than the first one. The parameter $`l/x_0`$ determines it’s smallness, where $`l`$ is the mean free path and $`x_0`$ is a characteristic length of the problem. In $`S/F`$ structures $`x_0\sqrt{D_F/h\text{ }}`$ is the shortest length because usually $`h>\mathrm{\Delta }.`$ In the limit $`l/x_0<<1,`$ that is, if the product $`h\tau `$ is small, one can express $`\stackrel{ˇ}{𝐠}_a(𝐫)`$ from Eq.(170) in terms of $`\stackrel{ˇ}{g}_s(𝐫)`$
$$\stackrel{ˇ}{𝐠}_a(𝐫;ϵ,ϵ^{})=l\stackrel{ˇ}{g}_s(𝐫;ϵ,ϵ_1)\stackrel{ˇ}{g}_s(𝐫;ϵ_1,ϵ^{}),$$
(176)
When obtaining Eq.(176), we used the relations
$`\stackrel{ˇ}{g}_s(𝐫;ϵ,ϵ_1)\stackrel{ˇ}{g}_s(𝐫;ϵ_1,ϵ^{})=1,`$ (177)
$`\stackrel{ˇ}{𝐠}_{as}(𝐫;ϵ,ϵ_1)\stackrel{ˇ}{g}_s(𝐫;ϵ_1,ϵ^{})+\stackrel{ˇ}{g}_s(𝐫;ϵ,ϵ_1)\stackrel{ˇ}{𝐠}_a(𝐫;ϵ_1,ϵ^{})=0`$
The symbolically written products in Eqs.(176),(177) imply an integration over the internal energy $`ϵ_1`$ as it is shown in Eq.(174).
The equation for the isotropic component of the Green’s function after averaging over the direction of the Fermi velocity $`𝐯_𝐅`$ reads
$`iD(\stackrel{ˇ}{g}\stackrel{ˇ}{g})`$ $`+`$ $`i\left(\widehat{\tau }_3_t\stackrel{ˇ}{g}+_t^{}\stackrel{ˇ}{g}\widehat{\tau }_3\right)+[\stackrel{ˇ}{\mathrm{\Delta }},\stackrel{ˇ}{g}]+[𝐡\stackrel{ˇ}{𝐒},\stackrel{ˇ}{g}]`$ (178)
$`+`$ $`{\displaystyle \frac{i}{\tau _{s.o.}}}[\stackrel{ˇ}{𝐒}\widehat{\tau }_3\stackrel{ˇ}{g}\widehat{\tau }_3\stackrel{ˇ}{𝐒},\stackrel{ˇ}{g}]=0,`$
where $`D`$ is the diffusion coefficient.
If we take the elements $`(11)`$ or $`(22)`$ of the supermatrix $`\stackrel{ˇ}{g}`$, we obtain the Usadel equation for the retarded and advanced Green’s functions $`\stackrel{ˇ}{g}^{R(A)}(t,t^{})`$ generalized for the case of the exchange field acting on the spins of electrons. In this review we are interested mainly in stationary processes, when the matrices $`\stackrel{ˇ}{g}^{R(A)}(t,t^{})`$ depend only on the time difference $`(tt^{}).`$ Performing the Fourier transformation $`\stackrel{ˇ}{g}^{R(A)}(ϵ)=d(tt^{})\stackrel{ˇ}{g}^{R(A)}(tt^{})\mathrm{exp}(iϵ(tt^{}))`$, we obtain for $`\stackrel{ˇ}{g}^{R(A)}(ϵ)`$ the following equation (we drop the indices $`R(A)`$)
$$D_x\left(\stackrel{ˇ}{g}_x\stackrel{ˇ}{g}\right)+iϵ[\widehat{\tau }_3\widehat{\sigma }_0,\stackrel{ˇ}{g}]+ih\left\{[\widehat{\tau }_3\widehat{\sigma }_3,\stackrel{ˇ}{g}]\mathrm{cos}\alpha (x)+[\widehat{\tau }_0\widehat{\sigma }_2,\stackrel{ˇ}{g}]\mathrm{sin}\alpha (x)\right\}+i[\stackrel{ˇ}{\mathrm{\Delta }},\stackrel{ˇ}{g}]+\frac{i}{\tau _{s.o.}}[\stackrel{ˇ}{𝐒}\widehat{\tau }_3\stackrel{ˇ}{g}\widehat{\tau }_3\stackrel{ˇ}{𝐒},\stackrel{ˇ}{g}]=0.$$
(179)
It is assumed here that $`𝐡`$ has the components $`h(0,\mathrm{sin}\alpha ,\mathrm{cos}\alpha )`$. This equation was first obtained by Usadel (1970) and it is known as the Usadel equation. An inclusion of the exchange and spin-orbit interaction was made in Alexander *et al.* (1985); Demler *et al.* (1997).
Eq. (178) can be solved analytically in many cases and it is used in most of previous sections in order to describe different $`S/F`$ structures. Solutions for the Usadel equation must obey the normalization condition
$$\stackrel{ˇ}{g}(𝐩_𝐅,𝐫;ϵ).\stackrel{ˇ}{g}(𝐩_𝐅,𝐫;ϵ)=1$$
(180)
The Usadel equation is complemented by the boundary conditions presented by Kuprianov and Lukichev (1988) on the basis of the Zaitsev’s boundary conditions Zaitsev (1984). Various aspects of the boundary conditions have been discussed by Lambert *et al.* (1997); Nazarov (1999); Kopu *et al.* (2004); Xia *et al.* (2002). In the absence of spin-flip processes at the interface they take the form:
$$\stackrel{ˇ}{g}_1_x\stackrel{ˇ}{g}_1=\frac{1}{2\gamma _a}[\stackrel{ˇ}{g}_1,\stackrel{ˇ}{g}_2],$$
(181)
where $`\gamma _1=R_b\sigma _1`$, $`\sigma _1`$ is the conductivity of the conductor $`1`$ and $`R_b`$ is the interface resistance per unit area, the $`x`$-coordinate is assumed to be normal to the plane of the interface.
The boundary condition (181) implies that we accept the simplest model of the S/F interface which is used in most papers on S/F structures. We assume that the interface separates two dirty regions: a singlet superconductor and a ferromagnet. The superconductor and the ferromagnet are described in the mean field approximation with different order parameters: the off-diagonal order parameter $`\mathrm{\Delta }`$ in the superconductor (in the weak coupling limit) and the exchange field $`h`$ in the ferromagnet acting on the spins of free electrons. No spin-flip scattering processes are assumed at the S/F interface. A generalization of the boundary conditions we use to the case of a spin-active S/F interface was carried out in the papers by Millis *et al.* (1988); Kopu *et al.* (2004); Fogelström (2000); Eschrig (2000).
Eqs. (178) and (181) together with the self-consistency equation that determines the superconducting order parameter $`\mathrm{\Delta }`$, constitute a complete set of equations from which one can obtain the Green’s functions.
The Usadel equation can be solved in some particular cases. We often use the linearized Usadel equation. In order to obtain the linearized Usadel equation we represent the Green’s functions $`\stackrel{ˇ}{g}`$ in the superconductor in the form
$$\stackrel{ˇ}{g}(𝐩_𝐅,𝐫;\omega )=\stackrel{ˇ}{g}_{BCS}(\omega )+\delta \stackrel{ˇ}{g}_S+\delta \stackrel{ˇ}{f}_S,$$
(182)
where $`\stackrel{ˇ}{g}_{BCS}(\omega )=\widehat{\tau }_3g_{BCS}(\omega )+i\widehat{\tau }_2f_{BCS}`$, $`g_{BCS}(\omega )=(i\omega /\mathrm{\Delta })f_{BCS},`$ $`f_{BCS}=\mathrm{\Delta }/i\sqrt{\omega ^2+\mathrm{\Delta }^2}`$. We have written the matrix $`\stackrel{ˇ}{g}`$ in the so called Matsubara representation. This means that a substitution $`ϵi\omega `$ ($`\omega =\pi T(2n+1),n=0,\pm 1,\pm 2,\mathrm{}`$..) is done and $`\stackrel{ˇ}{g}(\omega )`$ coincides with $`\stackrel{ˇ}{g}^R(ϵ)`$ for positive $`\omega `$ and with $`\stackrel{ˇ}{g}^A(ϵ)`$ for negative $`\omega `$. The linearized Usadel equation has the form
$$_{xx}^2\delta \stackrel{ˇ}{f}_S\kappa _S^2\delta \stackrel{ˇ}{f}_S=2i(\delta \stackrel{˘}{\mathrm{\Delta }}/D_S)g_{BCS}^2$$
(183)
in the $`S`$ region and
$$_{xx}^2\delta \stackrel{ˇ}{f}\kappa _\omega ^2\delta \stackrel{ˇ}{f}+i\kappa _h^2\left\{[\widehat{\sigma }_3,\delta \stackrel{ˇ}{f}]_+\mathrm{cos}\alpha \pm \widehat{\tau }_3[\widehat{\sigma }_2,\delta \stackrel{ˇ}{f}]_{}\mathrm{sin}\alpha \right\}=0$$
(184)
in the $`F`$ region. Here $`\kappa _S^2=2E_\omega /D_S`$, $`\kappa _\omega ^2=2|\omega |/D_F`$, $`\kappa _h^2=h`$ $`\mathrm{sgn}\omega /D_F`$ and $`[A,B]_\pm =AB\pm BA,\delta \stackrel{˘}{\mathrm{\Delta }}=i\widehat{\tau }_2\widehat{\sigma }_3\delta \mathrm{\Delta }`$. The signs $`\pm `$ in Eq. (184) correspond to the right and left layer respectively.
The boundary conditions for $`\delta \stackrel{ˇ}{f}_S`$ and $`\delta \stackrel{ˇ}{f}_Ff`$ (in zero-order approximation $`\delta \stackrel{ˇ}{f}_F=0`$) are obtained from Eq.(181). They have the form
$$_x\delta \stackrel{ˇ}{f}_S=(1/\gamma _S)[g_{BCS}^2\delta \stackrel{ˇ}{f}g_{BCS}\widehat{f}_{BCS}\widehat{\sigma }_3g_{F3}g_{BCS}\stackrel{ˇ}{f}_S]$$
(185)
$$_x\stackrel{ˇ}{f}_F=(1/\gamma _F)[g_{BCS}\delta \stackrel{ˇ}{f}\stackrel{ˇ}{f}_S]$$
(186)
where $`\gamma _{F,S}=R_b\sigma _{F,S}`$.
If the Green’s functions are known, one can calculate macroscopic quantities such as the current, magnetic moment etc. For example, the current is given by Larkin and Ovchinnikov (1984)
$$I_S=(L_yL_z/16)\sigma _F\mathrm{Tr}\left(\widehat{\tau }_3\widehat{\sigma }_0\right)𝑑ϵ(\stackrel{ˇ}{g}_s\stackrel{ˇ}{g}_s/x)_{12}$$
(187)
where $`L_{y,z}`$ are the widths of the films in $`y`$ and $`z`$ direction (the current flows in the transverse $`x`$-direction) and subscript $`(12)`$ shows that one has to take the Keldysh component of the supermatrix $`\stackrel{ˇ}{g}_s\stackrel{ˇ}{g}_s/x`$. A variation of the magnetic moment due to proximity effect is determined by formulae
$$\delta M_z=\mu _B\nu (1/2)i\pi T\underset{\omega }{}Tr(\widehat{\tau }_3\widehat{\sigma }_3\delta \stackrel{ˇ}{g})$$
(188)
$$\delta M_{x,y}=\mu _B\nu (1/2)i\pi T\underset{\omega }{}Tr(\widehat{\tau }_0\widehat{\sigma }_{1,2}\delta \stackrel{ˇ}{g})$$
(189)
where $`\nu `$ is the density-of-states at the Fermi level in the normal state and $`\mu _B=g\mu _{Bohr}`$ is an effective Bohr magneton.
Finally, it is important to make remarks concerning the notations used in this review. In most works where the $`S/F`$ structures with homogeneous magnetization are studied, the Green’s function $`\stackrel{ˇ}{g}`$ is a $`2\times 2`$ matrix with the usual normal and Gor’kov’s components. Of course, this simplification can be made provided the magnetizations of the $`F`$ layers involved in the problem are aligned in one direction. However, this simple form leads to erroneous results if the magnetizations are arbitrarily oriented with respect to each other. The $`4\times 4`$ form of the Green’s function is unavoidable if one studies structures with a non-homogeneous magnetization. Of course, the $`c`$\- operators in Eq. (163) can be defined in different ways. For example, Maki (1969) introduced a spinor representation of the field operators, which is equivalent to letting in Eq. (163) the spin index of the operator $`a`$ unchanged when $`n=2`$. This notation was used in later works (Demler *et al.*, 1997; Alexander *et al.*, 1985, e.g) in which the Green’s functions have a $`2\times 2`$ block matrix form. The diagonal blocks represent the normal Green’s functions, while the off-diagonal blocks represent the anomalous one. With this notation the matrix, Eq. (165), changes its form. For example, the term containing $`\mathrm{\Delta }`$ is proportional to $`i\widehat{\sigma }_2`$ and not to $`\widehat{\sigma }_3`$. The choice of the notation depends on the problem to solve. In order to study the triplet superconductivity induced in $`S/F`$ systems and to see explicitly the three projections ($`S_z=0,\pm 1`$) of the condensate function it is more convenient to use the operators defined in Eq.(163) (see for example Fominov *et al.* (2003) and Bergeret *et al.* (2001c)).
## Appendix B Future direction of the experimental research
As we have seen throughout the paper there is a great number of experiments on S/F structures. The variety of superconducting and ferromagnetic materials is very large. In this section we review briefly some of these experiments. We will not dwell on specific fabrication techniques but rather focus on the discussion: which pairs of material (S and F) are more appropriate for the observation of the effects studied in this review.
First experiments on S/F structures used strong ferromagnets (large exchange fields) as Fe, Ni, Co or Gd and conventional superconductors like Nb, Pb, V, etc. Hauser *et al.* (1963). In these experiments the dependence of the superconducting transition temperature on the thicknesses of the S and F layers has been measured. In other words the suppression of the superconductivity due to the strong exchange field of the ferromagnet was analyzed. It is clear that for such strong ferromagnets the spin splitting is large and therefore a mismatch in electronic parameters of the S and F regions is large. This leads to a low interface transparency and a weak proximity effect. This was confirmed by Aarts *et al.* (1997) in experiments on V/V<sub>1-x</sub>Fe<sub>x</sub> multilayers. Varying the concentration of Fe in the VFe alloys they could change the values of the exchange field and indirectly the transparency of the interface. Such systems consisting of a conventional superconductor and a ferromagnetic alloy, both with similar band structure (in the above experiment the mismatch was $`<5\%`$) , are good candidates for observing the effects discussed in sections IV.1., V.2. and V.3.
Weak ferromagnets have been used in the last years in many experiments on S/F structures. Before we turn our attention to ferromagnets with small exchange fields it is worth mentioning the experiment by Rusanov *et al.* (2004). They analyzed the so-called spin switch effect. In particular they studied the transport properties of Permalloy(Py)/Nb bilayers. They observed an enhancement of superconductivity in the resistive transition in the field range where the magnetization of the Py switches and many domains are present. Interesting for us is that Py shows a well defined magnetization switching at low fields and therefore it could be used in order to detect the long-range triplet component that appears when the magnetization of the ferromagnet is not homogeneous (see section III.3). Finally, a magnetic configuration analysis of the strong-ferromagnetic structures used in transport experiments as those performed by Petrashov *et al.* (1999); Giroud *et al.* (1998) may also serve to confirm the predictions of section III.3. As it was discussed before the increase in the conductance of the ferromagnet for temperatures below the superconducting $`T_c`$ may be explained assuming a long-range proximity effect.
The proximity effect in S/F is stronger if one uses dilute ferromagnetic alloys. Thus, such materials are the best candidates in order to observed most of the effects discussed in this review. The idea of using ferromagnetic alloys with small exchanges field was used by Ryazanov *et al.* (2001). They were the first in observing the sign reversal of the critical current in a S/F/S Josephson junction. Nb was used as superconductor while Cu<sub>0.48</sub>Ni<sub>0.52</sub> alloy as a ferromagnet (exchange field $`25K`$). (Later on similar results were obtained by Kontos *et al.* (2002) on Nb/Al/Al<sub>2</sub>O<sub>3</sub>/PdNi/Nb structures). The CuNi alloy was also used in the experiment by Gu *et al.* (2002b) on F/S/F structures. In this experiment the authors determined the dependence of the superconducting transition temperature on the relative magnetization-orientation of the two F layers. In order to get different alignments between the two CuNi layers an exchange-biased spin-valve stack of CuNi/Nb/CuNi/Fe<sub>50</sub>Mn<sub>50</sub> was employed. With a small magnetic field the authors could switch the magnetization direction of the free NiCu layer. This technique could be very useful in order to observe the Josephson coupling via the triplet component as described in section IV.1.
Finally, it is worth mentioning the experiment by Stahn *et al.* (2005) on YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub>/La<sub>2/3</sub>Ca<sub>1/3</sub>MnO<sub>3</sub>. Using the neutron reflectometry technique they observed a induced magnetic moment in the superconductor. Although the materials employed in this experiment cannot be quantitatively described with the methods presented in this review (the ferromagnet used by Stahn *et al.* (2005) is a half-metal with a exchange field comparable to the Fermi-energy and the superconductor is unconventional), the experimental technique may be used in other experiments in order to detected the induced magnetization predicted in sections V.2. and V.3.
## List of Symbols and Abbreviations
| S | | superconductor |
| --- | --- | --- |
| N | | nonmagnetic normal metal |
| F | | ferromagnetic metal |
| I | | insulator |
| LRTC | | long-range triplet component |
| $`\widehat{\tau }_i`$, $`i=1,2,3`$ | | pauli matrices in particle-hole space |
| $`\widehat{\sigma }_i`$, $`i=1,2,3`$ | | Pauli matrices in spin space |
| $`\widehat{\tau }_0`$, $`\widehat{\sigma }_0`$ | | unit matrices. |
| $`D`$ | | diffusion coefficient |
| $`\nu `$ | | density of states |
| $`\omega =\pi T(2n+1)`$ | | Matsubara Frequency |
| $`ϵ`$ | | real frequency (energy) |
| $`g_{BCS}`$ | | quasiclassical normal Green’s function for a bulk superconductor |
| $`f_{BCS}`$ | | quasiclassical anomalous Green’s function for a bulk superconductor |
| $`T_c`$ | | superconducting critical temperature |
| $`I_c`$ | | Josephson critical current |
| $`R_b`$ | | interface resistance per unit area |
| $`ϵ_{bN}=D_N/2R_b\sigma _Nd_N`$ | | minigap induced in a normal metal |
| $`\sigma _{S,F}`$ | | conductivity in the normal state |
| $`\gamma _{S,F}`$ | | $`R_b\sigma _{S,F}`$ |
| $`\gamma `$ | | ratio $`\sigma _F/\sigma _S`$ |
| $`𝒥`$ | | magnetic coupling between localized magnetic moments. |
| $`h`$ | | exchange field acting on the spin of conducting electrons |
| $`\xi _N=\sqrt{\frac{D_N}{2\pi T}}`$ | | characteristic penetration length of the condensate into a dirty normal metal |
| $`\xi _F=\sqrt{\frac{D_F}{h}}`$ | | characteristic penetration length of the condensate into a dirty ferromagnet |
| $`\xi _S=\sqrt{\frac{D_S}{2\pi T_c}}`$ | | superconducting coherence length for a dirty superconductor |
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# 1 Introduction and summary
## 1 Introduction and summary
To quantize an electromagnetic field one could fix a gauge by imposing an operator constraint on physical states of the theory. For instance, in the Gupta-Bleuler approach one postulates
$`G(A)|\mathrm{\Psi }=0,`$ (1)
where $`G(A)`$ denotes a function of the gauge field operator $`A`$ or its derivative, and $`|\mathrm{\Psi }`$ is an arbitrary physical state of the theory.
Alternatively, one could choose to impose the constraint on a gauge field already in a classical theory in such a way that (1) is automatic upon quantization . This can be achieved, e.g., by introducing into the classical action of the theory a Lagrange multiplier $`\lambda `$ times the function $`G(A)`$
$`{\displaystyle d^4x\lambda (x)G(A(x))}.`$ (2)
The field $`\lambda `$ has no kinetic or potential terms. Variation of the action w.r.t. $`\lambda `$ gives a constraint $`G(A(x))=0`$, which is just a classical counterpart of (1). Small fluctuations of the fields in this theory can consistently be quantized for Abelian as well as non-Abelian and gravitational fields . The resulting theories can be completed to be invariant under the Becchi-Rouet-Stora-Tyutin (BRTS) transformations , and a Hilbert space of physical states can be defined by requiring that the states carry zero BRST and ghost charges . Quantum effects in the resulting theory are identical to those of the Gupta-Bleuler approach.
Nevertheless, there is one difference in using the classically constrained theory that has not been explored. This difference could be seen in classical equations of motion. Variation of the action w.r.t. the gauge fields gives an equation of motion in which there are new terms proportional to $`\lambda `$ and/or its derivatives. Hence, classical field equations are modified, and, depending on allowed boundary conditions, new solutions could emerge.
This becomes especially important for gravity. General Relativity (GR) with a non-zero cosmological constant does not admit Minkowski space as a solution of equations of motion. We will discuss classically constrained General Relativity (CGR) in section 3 and show that the latter does admit flat space as a solution even if the cosmological constant is not zero. This difference is clearly very important.
Unimodular gravity (UGR) is an interesting example of a partially constrained theory. In UGR the full reparametrization invariance of GR is restricted to a subgroup of volume preserving transformations. One practical difference between UGR and GR is that the cosmological constant problem in UGR is somewhat relaxed. This is because UGR with a cosmological constant admits an infinite number of maximally symmetric solutions labeled by the value of the space-time curvature. However, this does not explain why one should choose the desirable (almost) flat solution among a continuum of maximally symmetric ones. Another issue in UGR is related to quantum loops of matter and gravity. It is likely that in this theory the Lagrange multiplier $`\lambda `$ acquires quadratic terms via the quantum loops (see, section 3); if so, new interactions would be needed to maintain the classical solutions of UGR in a quantum theory.
Constrained GR improves on the above-mentioned aspects: (i) it admits only two maximally symmetric solutions – one with zero curvature, and another one with the curvature obtained in GR; (ii) its classical properties, due to the BRST invariance of that model, are not modified by the quantum loop effects.
Reparametrization invariance in CGR is completely constrained. We will show in section 3, that one can still define a locally inertial reference frame in a small region around an arbitrary space-time point. This is because the constraint that fixes gauge in the whole space, allows, in the neighborhood of a given point, for the point-dependent gauge transformations that locally eliminate effects of gravity. Thus the equivalence principle is preserved.
Do the constrained theories solve the “old” cosmological constant problem ? In UGR the answer is negative because the theory admits an infinite number of maximally symmetric solutions. Clearly, in CGR, there still exists a conventional de Sitter solution of GR which can be used for inflation in the early universe. But the question is whether there exists an infinite number of other non-maximally symmetric solutions in the theory with a cosmological constant. If some of these solutions are physical, one should understand why in our Universe the (almost) flat solution is preferred. If, on the other hand, the non-maximally symmetric solutions can be disregarded for one reason or other, then CGR could be a good starting point for trying to accommodate inflation in the early universe and still solve the “old” cosmological constant problem. These important issues, including the question of stability of the new solutions, are being studied in . The purpose of the present work is to investigate CGR to determine whether it is a consistent low-energy quantum field theory, which by itself is a legitimate theoretical question. Concrete applications of this model will be discussed in .
Even in the most optimistic scenario for CGR, one should explain why the space-time curvature is not exactly zero but $`H_0^2(10^{42}\mathrm{GeV})^2`$, as suggested by recent observations. Where could this scale come from? One possibility is to introduce a pseudo Nambu-Goldstone boson potential that gives rise to required “dark energy” <sup>1</sup><sup>1</sup>1In this case though, one needs a VEV of a scalar to be somewhat higher than the Planck mass.. Another way is to try introducing a graviton mass $`m_gH_0`$. Although both of the above approaches postulate the existence of a new small scale, this scale is stable w.r.t. quantum corrections (i.e., it is technically natural). In this regard, we briefly discuss in the present work massive deformations of classically constrained gauge and gravitational theories. We will find that the UV behavior of the propagators of massive gauge and gravitational quanta are softened. Although these results are encouraging, at this stage we still lack an understanding of whether the non-linear unitarity of the S-matrix on a Hilbert space of physical states can be preserved using the BRST and ghost charges of these models, or whether some further modifications may be needed. Studies on these issues will be reported elsewhere.
The work is organized as follows: in section 2 we discuss a constrained theory of a photon and its classical equations and solutions. In section 3 we discuss a classically constrained theory of gravity (CGR). We find that CGR with a cosmological constant has a remarkable property – it admits two maximally symmetric solutions one of which is a flat space. We find general cosmological solutions in CGR as well as the expression for the Schwarzschild metric. An important question that is also addressed in section 3 is that of radiative stability. Using the BRST invariant version of the theory we argue that quantum corrections do not ruin the obtained classical results. Massive deformations of CGR are also briefly discussed in section 3. In Appendix A we discuss the spectrum of the constrained Abelian gauge theory in various approaches, including Stückelberg’s method. We also look at the massive deformation of this model pointing out that the massive propagator, unlike in the Proca theory, has smooth UV behavior and a nonsingular massless limit. Appendix B deals with constrained non-Abelian gauge fields. After briefly discussing classical equations, we study the spectrum of the theory. Comments on massive deformations of the non-Abelian theories are also included. The results are similar to those of the Abelian case. Some parts of Appendix A and B are of a review character, but we felt that including these discussions would make our presentation more complete.
## 2 Constrained photon
As an instructive example we consider electrodynamics with an imposed classical constraint. We call this model constrained QED (CQED) even though we will only quantize it later. We start with the Lagrangian density
$`={\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+A_\mu J^\mu +\lambda (_\mu A^\mu ).`$ (3)
Here $`\mu ,\nu =0,1,2,3;`$ $`J^\mu `$ is a current, and $`\lambda `$ is a Lagrange multiplier (our choice of the Lorentzian signature is “mostly negative”).
We start by studying classical properties of this theory. The equations of motion that follow from the above Lagrangian are
$`^\nu F_{\nu \mu }+J_\mu _\mu \lambda =0,`$ (4)
$`_\mu A^\mu =\mathrm{\hspace{0.17em}0}.`$ (5)
Taking a derivative of Eq. (4) one obtains the following relation
$`\mathrm{}\lambda =^\mu J_\mu .`$ (6)
If the current $`J_\mu `$ is conserved, $`\lambda `$ is a harmonic function, $`\mathrm{}\lambda =0`$. One particular solution of this equation is $`_\mu \lambda =C_\mu `$, where $`C_\mu `$ denotes an arbitrary space-time constant four-vector. Physically, the effective conserved current to which the gauge field is coupled in (3) is
$`J_\mu ^{\mathrm{eff}}=J_\mu _\mu \lambda =J_\mu C_\mu .`$ (7)
The last term on the r.h.s. of (7) acts as a constant background current density determined by a vector $`C_\mu `$. Since $`C_\mu `$ is an integration constant, there is a continuous set of $`C_\mu `$’s that one could choose from. Setting the value of $`C_\mu `$ is equivalent of choosing the corresponding boundary conditions. The equations of motion (4) and (5) were derived by varying the action corresponding to (3) with the following boundary conditions
$`\delta A_\mu |_{\mathrm{boundary}}=0,\delta \lambda |_{\mathrm{boundary}}=\mathrm{finite}\mathrm{function}.`$ (8)
Therefore, all the solutions should obey (8). To demonstrate that such solutions exist let us consider a simple example of a spherically symmetric localized charge density for which $`J_0=\rho \theta (r_0r)`$, $`J_i=0`$. Here $`\rho `$ is a constant charge density, $`\theta (r)`$ denotes the step function, $`r`$ is the radial coordinate, and $`r_0`$ is the radius of the charge distribution. We substitute this source into the RHS of (4). In addition we chose a solution for the Lagrange multiplier to be $`_0\lambda C_0=\rho `$, and $`_j\lambda C_j=0`$. For this source a solution of Eqs. (4) and (5), and the corresponding electric field read
$`A_0`$ $`={\displaystyle \frac{\rho r_0^3}{3r}}{\displaystyle \frac{\rho r^2}{6}}\mathrm{for}rr_0,A_0={\displaystyle \frac{\rho r_0^2}{2}}\mathrm{for}rr_0,`$ (9)
$`\stackrel{}{E}`$ $`={\displaystyle \frac{\rho }{3}}(1{\displaystyle \frac{r_0^3}{r^3}})\theta (rr_0)\stackrel{}{r}.`$
The above solution satisfies the boundary conditions (8) at $`r=r_0`$. In the conventional electrodynamics the source $`J_0=\rho \theta (r_0r)`$, $`J_i=0`$, yields an electric field that is well-known and differs from (9). The origin of this difference is clear – in CQED the quantity $`J_\mu `$ that is specified in the action is not the whole source producing the gauge field! An additional integration constant appears in the equations of motion and the total source is $`J_\mu ^{\mathrm{eff}}`$ (7). The above theory reduces to conventional electrodynamics when we choose $`\lambda =0`$. This corresponds to what we measure in ordinary experiments.
Similar solutions with a nonzero value of the Lagrange multiplier will play an important role for gravity with a nonzero cosmological constant. These will be discussed in section 3 (non-Abelian gauge fields are discussed in Appendix B).
So far we have not emphasized the fact that the Lagrangian (3) is not gauge invariant. In fact, variation of (3) under the gauge transformation $`\delta A_\mu =_\mu \alpha (x)`$ vanishes (up to a surface term) for configurations satisfying $`\mathrm{}\lambda =0`$, therefore (3) has gauge invariance when the $`\lambda `$ field is on-shell, even though the gauge field could be off-shell. Given that the $`\lambda `$ is not propagating, and remains such at the quantum level (see below), this suggests two physical degrees of freedom for a photon, while off-shell there are four degrees of freedom. This will be established more rigorously below. As it is shown in Appendix A, the two extra off-shell degree of freedom are decoupled from conserved sources and have no relevance in the Abelian case.
The quantum theory of CQED could be approached in a number of different ways. One could restore first a manifest gauge invariance of (3) using the Stückelberg method, and then quantize the resulting gauge invariant theory. We will discuss this in Appendix A and B for Abelian and non-Abelian gauge fields respectively.
In the Feynman integral formulation one could think of the constrained approach as follows: the gauge and auxiliary fields can be decomposed in their classical and quantum parts, $`A=A_{cl}+\delta A`$, and $`\lambda =\lambda _{cl}+\delta \lambda `$. For classical solutions, $`A_{cl}`$ and $`\lambda _{cl}`$, we allow boundary conditions that are different from the conventional ones. In particular, we allow for a nonzero solution of the $`\mathrm{}\lambda =0`$ equation, a zero-mode. This is an unconventional step, since the usual FP term that appears in the path integral is the determinant of the operator acting on $`\lambda `$, that is $`det(\mathrm{})`$ in this particular case. Because of the zero-mode, the determinant, $`det(\mathrm{})`$, would have been zero, leading to an ill-defined partition function. However, a right way to formulate the path integral is to separate the zero-mode, and integrate only w.r.t. the small fluctuations for which $`det(\mathrm{})`$ is non-zero, and for which the conventional radiation boundary conditions are imposed.
This is the procedure that we will be assuming throughout the text. Then, quantization of the fluctuations $`\delta A`$ and $`\delta \lambda `$ over the classical background $`A_{cl}`$, $`\lambda _{cl}`$ can be performed in a BRST invariant way. That is, one could start by postulating BRST invariance of the Lagrangian for the fluctuations by adding in the Faddeev-Popov (FP) ghosts without any reference to a local gauge symmetry, but elevating BRST invariance to a fundamental guiding principle in constructing the Lagrangian. The resulting BRST symmetric theory of small fluctuations could be quantized in a conventional way. This is what we summarize below. For simplicity of presentation we will replace $`\delta A_\mu `$ and $`\delta \lambda `$ by $`A_\mu `$ and $`\lambda `$ keeping in mind that these are fluctuations over a classical background (the same replacement will be assumed for non-Abelian and gravity fields considered in the next sections and Appendices).
In a conventional Feynman integral approach the measure in the path integral should be modded out by the gauge equivalent classes. The FP trick does this job by introducing the gauge-fixing term along with the FP ghosts. In this regard the following natural question arises – since (3) (or its non-Abelian counterpart) is not gauge invariant why do we need to introduce the FP ghosts in the theory? Naively, it would seem that we should use the path integral
$`{\displaystyle DA\delta (G(A))\mathrm{exp}\left(i_{\mathrm{Gauge}\mathrm{fields}}\right)},`$ (10)
instead of the one with the FP ghosts
$`{\displaystyle DA\delta (G(A))\mathrm{det}\left|\frac{\delta G^\omega }{\delta \omega }\right|\mathrm{exp}\left(i_{\mathrm{Gauge}\mathrm{fields}}\right)},`$ (11)
where $`G`$ is a gauge-fixing condition (for instance it could be that $`G=_\mu A^\mu `$), and $`G^\omega `$ is gauge transformed $`G`$ with $`\omega `$ being the transformation parameter (the above definition is consistent as far as we do not include the zero modes in the integration and $`det(\mathrm{})`$ is nonzero, as we discussed above). In QED the difference is irrelevant because the FP ghost are decoupled from the rest of the physics, however, in the case of non-Abelian fields (10) would lead to a non-unitary theory. Because of the absence of gauge invariance in (3), the argument to introduce the FP ghosts in the path integral of (3) cannot be the same as in the conventional case. Nevertheless, the FP ghosts can be motivated by a symmetry argument, namely by requiring the BRST invariance of the constrained action.
Let us modify (3) by adding the FP ghost $`c`$ and anti-ghost $`\overline{c}`$ which are Grassmann variables, $`c^2=(\overline{c})^2=0`$ and $`\overline{c}^+=\overline{c},c^+=c`$. The Lagrangian reads as follows
$`={\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+A_\mu J^\mu +\lambda (_\mu A^\mu )i\overline{c}\mathrm{}c.`$ (12)
Since on the classical backgrounds considered here the FP ghost fields vanish, the classical properties of (3) and (12) are identical<sup>2</sup><sup>2</sup>2In general, a ghost-antighost bilinear could have nonzero expectation values on certain states, however, these stats do not satisfy the zero BRST and Ghost charge conditions, see below.. The Lagrangian (12), however, is invariant under the following continuous BRST transformations:
$`\delta A_\mu =i\zeta _\mu c,\delta c=0,\delta \overline{c}=\lambda \zeta ,\delta \lambda =0,`$ (13)
where $`\zeta `$ is a coordinate independent Grassmann transformation parameter such that $`(\zeta c)^+=c^+\zeta ^+=c\zeta `$. The Lagrangian (12) gives the path integral that is identical to that of the conventional approach (11) which takes the form
$`{\displaystyle 𝒟A𝒟\lambda 𝒟\overline{c}𝒟c\mathrm{exp}\left[i\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+A^\mu J_\mu +\lambda (_\mu A^\mu )i\overline{c}\mathrm{}c\right]}.`$ (14)
Notice that the rationale for writing down (14) and the Lagrangian (12) is different from the motivation that led to the path integral (11). In the constrained approach the rules for constructing a Lagrangian and path integral are:
(1) Impose classical constraints on fields using the Lagrange multiplier technique. Find the corresponding zero-modes and treat them separately from small fluctuations of the fields.
(2) Introduce the FP ghosts to obtain the BRST invariance of the gauge non-invariant theory.
(3) Use this Lagrangian to set up the path integral in a straightforward way.
The presence of the BRST symmetry guarantees that all the Ward-Takahashi identities (the Slavnov-Taylor identities in the non-Abelian case) of the conventional theory are preserved in the constrained approach, even though the classical equations of motion in this approach are different as discussed above.
Let us now turn to the loop corrections that emerge in (12). In particular we would like to make sure that no kinetic or potential terms are generated for $`\lambda `$. We do have a symmetry $`\lambda (x)\lambda (x)+\beta (x)`$, where $`\beta `$ is an arbitrary function, w.r.t. which (12) is invariant. Kinetic or potential terms of $`\lambda `$ would break it. The question is whether this symmetry is preserved by the loop corrections. To address this issue we calculate the propagator of the gauge fields. This can be done in a few different ways. The easiest one is to add a fictitious term $`\frac{1}{2}\gamma \lambda ^2`$ to the Lagrangian density and then take the limit $`\gamma 0`$
$$\begin{array}{cc}\hfill Z[J_\mu ]& \underset{\gamma 0}{lim}𝒟A𝒟\lambda 𝒟\overline{c}𝒟c\mathrm{exp}\left[i_g+A^\mu J_\mu +\lambda (_\mu A^\mu )\frac{1}{2}\gamma \lambda ^2i\overline{c}\mathrm{}c\right]\hfill \\ \hfill & \underset{\gamma 0}{lim}𝒟A𝒟\overline{c}𝒟c\mathrm{exp}\left[i_g+A^\mu J_\mu +\frac{1}{2\gamma }(_\mu A^\mu )^2i\overline{c}\mathrm{}c\right],\hfill \end{array}$$
(15)
where $`_g\frac{1}{4}F_{\mu \nu }F^{\mu \nu }`$. The propagator is obtained following the standard procedure and the result is
$$\mathrm{\Delta }^{\mu \nu }=\underset{\gamma 0}{lim}\frac{\eta ^{\mu \nu }(1+\gamma )^\mu ^\nu /\mathrm{}}{\mathrm{}}=\frac{\eta ^{\mu \nu }^\mu ^\nu /\mathrm{}}{\mathrm{}},$$
(16)
which coincides with the standard transverse QED propagator in Landau gauge<sup>3</sup><sup>3</sup>3 Here and below we always assume the Feynman (causal) prescription for the poles in the propagators.. It is straightforward to check that loop corrections preserve the transversality of the gauge field propagator. All the diagrams that renormalize the propagator consist of the standard bubble diagrams produced by the contractions between two currents. Because the current to which the gauge field is coupled is conserved, the propagator remains transverse to all loops
$`^\mu \lambda A_\mu A_\nu ^\nu \lambda ^\mu \lambda (\eta _{\mu \nu }\mathrm{}_\mu _\nu )^\nu \lambda =0.`$ (17)
As a result, loops cannot generate the kinetic term for $`\lambda `$.
Last but not least, the BRST symmetry of (12) can be used to define the Hilbert space of physical states by imposing the standard zero BRST- and Ghost- charge conditions $`Q^{\mathrm{BRST}}|\mathrm{Phys}>=0`$, $`Q^{\mathrm{Ghost}}|\mathrm{Phys}>=0`$. In the Abelian case this reduces to the Nakanishi-Lautrup condition $`\lambda ^{()}|\mathrm{Phys}>=0`$ (here $`\lambda ^{()}`$ denotes a negative frequency part of a fluctuation of $`\lambda `$ over a classical background) which ensures that the quantum of $`\lambda `$ is not in a set of physical in and out states of the theory.
## 3 Constrained models of gravity
There exists a modification of GR, unimodular gravity (UGR) , that partially restricts its gauge freedom and yet reproduces the observable results of the Einstein theory. In UGR reparametrization invariance is restricted to the volume-preserving diffeomorphisms that keep the value of $`\sqrt{g}`$ intact ($`g|detg_{\mu \nu }|`$). Such a theory can be formulated by using the Lagrange multiplier $`\lambda `$ (we set $`8\pi G_N=1`$, unless indicated otherwise):
$$=\frac{\sqrt{g}}{2}R\lambda (\sqrt{g}1)+_M,$$
(18)
where $`_M`$ is the Lagrangian of the matter fields. The equations of motion and Bianchi identities of this theory require $`\lambda `$ to be an arbitrary space-time constant . As a result, UGR is equivalent at the classical level to GR except that cosmological term becomes an arbitrary integration constant. Since the latter can take any value, there are an infinite number of solutions parametrized by a constant scalar curvature. Such a theory, at least at the classical level, seems to be more favorable than GR – for an arbitrary large value of the vacuum energy in the Lagrangian (arising, e.g., from particle physics) one is always able to find an almost-flat solution. Of course this still does not explain why one should choose the desirable almost-flat solution among a continuum of maximally symmetric ones. This is one aspect of UGR on which one would like to improve.
Another, perhaps more pertinent question in UGR emerges when one considers quantum loops of matter and gravity. Inspection of loop diagrams (see section 4.2) suggest that the Lagrange multiplier $`\lambda `$ in (18) would acquire the mass ($`\lambda ^2`$) as well as kinetic ($`(_\mu \lambda )^2`$) terms due to the quantum effects. As a result, $`\lambda `$ would cease to be an auxiliary field, and all the classical results of UGR would have to be reconsidered.
We will discuss a model that completely constrains reparametrization invariance of GR. This theory has the following two important properties: (I) It allows, like UGR does, an adjustment of the cosmological constant via the integration constant mechanism, but only admits two maximally symmetric solutions, one of which is a flat space. (II) Unlike UGR its classical properties are stable w.r.t. quantum corrections, i.e., the Lagrange multiplier remains an auxiliary field even in the quantum theory. This theory also preserves the equivalence principle.
### 3.1 Constrained gravity and the cosmological constant
Consider the following Lagrangian
$$=\frac{\sqrt{g}}{2}R+\sqrt{g}g^{\mu \nu }_\mu \lambda _\nu +_M+\mathrm{surface}\mathrm{terms}.$$
(19)
Here, $`\lambda _\nu `$ is a vector that serves as a Lagrange multiplier and $`_M`$ denotes the Lagrangian of other fields which can also include a vacuum energy term (the cosmological constant) produced by classical and/or quantum effects. Versions of this model have been discussed in the literature previously (see, e.g., the last reference in ) with the purpose of introducing the de Donder gauge-fixing condition in the context of quantization of gravity. Here, instead, we regard this model as a classical theory, which is subsequently quantized, but the classical equations of which could admit rather interesting solutions that are absent in GR. Thus, we first concentrate on classical effects, leaving the discussion of quantum corrections for the next subsection<sup>4</sup><sup>4</sup>4One could choose to impose a different constraint using the Lagrange multiplier in (19) (for instance, an axial-gauge constraint). As long as it is an acceptable gauge-fixing condition for small fluctuations, and the corresponding FP ghosts are taken care of consistently (in the axial-gauge the FP ghost are not needed) the quantum effects of the fluctuations won’t depend on the choice of this constraint. However, the classical solutions could differ for different constraints. Here we choose the constraint that is Lorentz-invariant..
Among others, we will be considering below solutions with fixed boundaries. For such solutions one should add to the Einstein-Hilbert action the Gibbons-Hawking boundary term. Moreover, the boundary conditions that we allow for are:
$`\delta g_{\mu \nu }|_{\mathrm{boundary}}=0,\mathrm{and}\delta \lambda _\mu |_{\mathrm{boundary}}=0.`$ (20)
Under these conditions the variation w.r.t. the metric gives
$`G_{\mu \nu }\left(_\mu \lambda _\nu +_\nu \lambda _\mu \right)+g_{\mu \nu }^\sigma \lambda _\sigma =T_{\mu \nu }.`$ (21)
In addition to this, we get a constraining equation by varying the action w.r.t. the Lagrange multiplier
$`_\mu (\sqrt{g}g^{\mu \nu })=0.`$ (22)
The above model is similar in spirit to a constrained gauge theory discussed in the previous section. It is easily checked that (22) is equivalent to the condition
$`\sqrt{g}\mathrm{\Gamma }_{\mu \nu }^\alpha g^{\mu \nu }=0,`$ (23)
and, in the linearized approximation, gives rise to de Donder (harmonic) gauge fixing of linearized GR. Due to this condition, $`_\mu A^\mu =g^{\mu \nu }_\mu A_\nu ^\mu A_\mu `$, where $`_\mu `$ is a covariant derivative acting on an arbitrary four-vector $`A_\mu `$.
The constraint (22) (or (23)) fixes completely reparametrization invariance of the theory. How is then the equivalence principle recovered? A relevant property of (22) (or (23)) is this: for any given point $`x^\mu =x_0^\mu `$ in the coordinate system $`\{x^\mu \}`$, it allows for the $`x_0^\mu `$-dependent coordinate transformations that eliminate the connection in a small neighborhood of this point. These transformations can be written as
$`x^\mu =x^\mu +{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{\alpha \beta }^\mu (x_0)(xx_0)^\alpha (xx_0)^\beta .`$ (24)
It is straightforward to check that for (24) $`g^{\alpha \beta }(x)|_{x=x_0}=g^{\alpha \beta }(x^{})|_{x^{}=x_0}`$, and
$`\mathrm{\Gamma }_{\alpha \beta }^\mu (x)|_{x=x_0}=\mathrm{\Gamma }_{\alpha \beta }^\mu (x^{})|_{x^{}=x_0}+\mathrm{\Gamma }_{\alpha \beta }^\mu (x)|_{x=x_0}.`$ (25)
As a result, $`\mathrm{\Gamma }_{\alpha \beta }^\mu (x^{})|_{x^{}=x_0}=0`$, and at this point the metric can simultaneously be brought to the Minkowski form. The transformation (24), on the surface of (22) (or (23)), is trivially consistent with (22) (or (23))
$`g^{\alpha \beta }\mathrm{\Gamma }_{\alpha \beta }^\mu (x)|_{x=x_0}=g^{\alpha \beta }(x^{})\mathrm{\Gamma }_{\alpha \beta }^\mu (x^{})|_{x^{}=x_0}=0.`$ (26)
Since $`x_0`$ was arbitrary, the above arguments can be repeated for any other point in space-time. One should emphasize again that the coordinate transformation (24) is point-specific, i.e., at different points of space-time one should perform transformations that depend parametrically on that very point. This is why they are not gauge transformations in the entire space-time. Summarizing, although (22) (or (23)) picks a global coordinate frame, it allows for the point-dependent coordinate transformations that can eliminate gravity locally. Hence, the equivalence principle. The local Lorentz transformations are also preserved.
To obtain the equation which the Lagrange multiplier has to satisfy we apply a covariant derivative to both sides of (21). This is subtle since the second term on the l.h.s. of (21) is not a tensor and we should define the action of a covariant derivative on this object. We adopt the following straightforward procedure: apply to both sides of (21) the operator
$`^\alpha g^{\alpha \mu }\left(_\mu \mathrm{\Gamma }_\mu ^{}\mathrm{\Gamma }_\mu ^{}\right),`$ (27)
where the standard index arrangement should be used in place of the asterisks even if the two-index object on which this operator is acting does not transform as a tensor. Then, using the Bianchi identities and covariant conservation of the stress-tensor we obtain:
$$g^{\mu \alpha }_\mu (_\alpha \lambda _\nu +_\nu \lambda _\alpha )=_\nu ^\sigma \lambda _\sigma .$$
(28)
The above equation can be simplified substantially due to (22). The left hand side of (28) can be reduced to
$`g^{\mu \alpha }_\mu \left(_\alpha \lambda _\nu +_\nu \lambda _\alpha \right)g^{\mu \alpha }\left(_\alpha \lambda _\beta +_\beta \lambda _\alpha \right)\mathrm{\Gamma }_{\nu \mu }^\beta ,`$
while the right hand side simplifies to give
$`(_\nu g^{\alpha \beta })_\alpha \lambda _\beta +g^{\alpha \beta }_\nu _\alpha \lambda _\beta .`$
Combining the above two expressions together we find from (28)
$$g^{\mu \alpha }_\mu _\alpha \lambda _\nu =0.$$
(29)
As long as $`g^{\mu \nu }`$ is non-singular general solutions for $`\lambda _\nu `$ could be found.
Clearly, the system of equations (21), (22) and (28) (or (29)), could admit new solutions that are absent in the Einstein theory. For instance, the Einstein equations with a nonzero cosmological constant (i.e., $`T_{\mu \nu }=\mathrm{\Lambda }g_{\mu \nu }`$), do not admit Minkowski space as a solution. In contrast with this, the system (21), (22) and (28) is satisfied by the flat space metric
$`g_{\mu \nu }=\eta _{\mu \nu },_\mu \lambda _\nu +_\nu \lambda _\mu =\mathrm{\Lambda }\eta _{\mu \nu }.`$ (30)
It is remarkable that one can obtain a flat solution even though the vacuum energy in the Lagrangian is not zero! Similar property exists in unimodular gravity when $`T_{\mu \nu }=\mathrm{\Lambda }g_{\mu \nu }`$, but there are an infinite number of maximally symmetric solutions labeled by the value of a constant curvature. Is this also true in the model (19)? We analyze this issue below. First we notice that the system (21), (22) and (28) admits a de Sitter solution of conventional general relativity (we will focus on the case of a positive cosmological constant only from now on and interesting results can be obtained for both positive and negative $`\mathrm{\Lambda }`$, ). The conventional dS metric solves (21) with $`\lambda _\mu =const.`$. What is less obvious is how this solution satisfies (22). To understand this we start with the dS solution in the co-moving coordinate system (this is also applicable to any FRW cosmology)
$$\mathrm{d}s^2=\mathrm{d}t^2a^2(t)\left(\mathrm{d}x^2+\mathrm{d}y^2+\mathrm{d}z^2\right).$$
(31)
As it can be checked directly, (31) does not satisfy (22). However, we can define a new time variable
$$\tau \frac{\mathrm{d}t}{a^3(t)},$$
(32)
for which the interval becomes
$$\mathrm{d}s^2=a^6(\tau )\mathrm{d}\tau ^2a^2(\tau )(\mathrm{d}x^2+\mathrm{d}y^2+\mathrm{d}z^2).$$
(33)
This metric satisfies (22) identically. Therefore, a conventional dS metric, or any other spatially-flat FRW cosmology, is a solution of the system (21), (22), (28).
Having established that the flat and conventional dS spaces are two solutions of the theory, let us now look at other possible maximally symmetric solutions. In general, the following ansatz
$`_\mu \lambda _\nu +_\nu \lambda _\mu =cg_{\mu \nu }^{dS},`$ (34)
where $`c`$ is an arbitrary constant, $`g_{\mu \nu }^{dS}`$ is a dS solution with $`R=4(\mathrm{\Lambda }c)`$, does satisfy the equations (21), (22) and (28). However, this is not enough to claim that (34) is a legitimate solution of the theory. This is because equation (34) itself might not be solvable in terms of $`\lambda _\mu `$ given that $`g_{\mu \nu }`$ is a dS metric obeying (22); solvability for $`\lambda _\mu `$ is a necessary condition since it is w.r.t. $`\lambda _\mu `$ that we varied the action. It is straightforward to check that there is no solution for $`\lambda _\mu `$ that would satisfy equation (34) if the metric is given by (33). Could there be other forms of dS space that are non-trivially different from (33) and yet satisfy (22)? The answer is no. To see this consider a dS solution in GR in two different coordinate systems. Let us assume the opposite, that both of these coordinate systems can be gauge transformed in GR to two different coordinate systems for which (22) is valid. If true, this would mean that the condition (22) does not completely fix the gauge freedom of GR. On the other hand, we find by performing gauge transformation of (22) that this is only possible if the gauge transformation itself is trivial. Therefore, the form (33) is the unique dS solution which satisfies the constraints (22).
The fact that (34) cannot be a maximally symmetric solution with nonzero $`R(g)`$ could also be established just by looking at a general expression of the Ricci scalar in terms of the metric $`g_{\mu \nu }`$ in which the substitution $`cg_{\mu \nu }=_\mu \lambda _\nu +_\nu \lambda _\mu `$ is made. If this ansatz could describe a dS space, one could always go to a weak-field regime where the Ricci curvature, as a function of the metric, has to be nonzero. However, the above ansatz for $`g_{\mu \nu }`$ gives zero $`R`$ in the leading order of the weak-field approximation.
Summarizing, we conclude that in the class of maximally symmetric spaces there are only two solutions of the theory: (I) The flat space defined in (30); (II) The (A)dS solution as it would appear in conventional GR transformed to a new coordinate system. There could in principle exist other, non-maximally symmetric solutions. It would be interesting to study whether those solutions are physical. If they are one should look for arguments why the maximally symmetric solutions could be preferred in our Universe. On the other hand, if the non-maximally-symmetric solutions are not there, then the model (19) could be a playground for studying the fate of the “old” cosmological constant problem . It is interesting to point out that by integrating out the Lagrange multiplier field one gets a non-local action, with some similarities but also differences of the resulting non-local theory with that of Ref. .
We would also point out that the Schwarzschild metric of conventional GR is also a legitimate solution of the CGR. A simplest way to address this is to choose $`\lambda _\nu =0`$ and make sure that the known GR solution itself satisfies (22) in a particular coordinate system. Usual solutions of GR can be transformed to satisfy Eq. (22). This is in particular true when $`g_{\mu \nu }`$ is diagonal and each element of $`\sqrt{g}g_{\mu \mu }`$ (there is no summation w.r.t. $`\mu `$ here) factorizes into the products of the form $`\sqrt{g}g_{\mu \mu }=h(x^\alpha )j(x^\nu )\mathrm{}f(x^\lambda )`$, where each function depends on one coordinate only. In this case, the constraints (22) turn into four separate partial differential equations
$$_\mu (\sqrt{g}g_{\mu \mu }^1)=0,\mu =0,1,2,3,\text{no summation w.r.t. }\mu .$$
(35)
Suppose we introduce one new coordinate $`\stackrel{~}{x}^\alpha =\stackrel{~}{x}^\alpha (x^\alpha )`$ such that it depends only on $`x^\alpha `$, and leave all the other coordinates intact. In the new coordinate system
$$\sqrt{\stackrel{~}{g}}=x^\alpha \sqrt{g},$$
(36)
while
$$\stackrel{~}{g}^{\alpha \alpha }=(x^\alpha )^2g^{\alpha \alpha },\text{and}\stackrel{~}{g}^{\mu \mu }=g^{\mu \mu }\text{for}\mu \alpha ,$$
(37)
where we have defined $`x^\alpha \frac{\mathrm{d}x^\alpha }{\mathrm{d}\stackrel{~}{x}^\alpha }`$ and chosen it to be positive.
The $`\alpha `$-th equation of (35) in the new coordinate system takes the form:
$$\frac{(\sqrt{\stackrel{~}{g}}\stackrel{~}{g}^{\alpha \alpha })}{\stackrel{~}{x}^\alpha }=\frac{[\sqrt{g}g^{\alpha \alpha }(x^\alpha )^1]}{\stackrel{~}{x}^\alpha }=0.$$
(38)
If $`\sqrt{g}g^{\alpha \alpha }`$ can be factorized, say as, $`\sqrt{g}g^{\alpha \alpha }=h(x^\alpha )\psi `$ where $`\psi `$ depends only on coordinates other than $`x^\alpha `$, we can find the desired $`\stackrel{~}{x}^\alpha `$ by simply demanding
$$\frac{h(x^\alpha )}{x^\alpha (\stackrel{~}{x}^\alpha )}=1,\text{or any constant if more convenient,}$$
(39)
and solving this ordinary differential equation. It is not difficult to see that one can carry on the same procedure for each $`x^\mu `$’s without invalidating the previous results, and, therefore, eventually find the new coordinate system that satisfies (22).
The above procedure is directly applicable to the Schwarzschild and FRW solutions. For the latter the result was already given above (see (33)). Here we perform the change of coordinates for the Schwarzschild metric. In a spherically symmetric coordinates
$$\mathrm{d}s^2=\left(1\frac{r_g}{r}\right)\mathrm{d}t^2\left(1\frac{r_g}{r}\right)^1\mathrm{d}r^2r^2\left(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\phi ^2\right).$$
(40)
Here, $`r_g=2G_NM`$ is the horizon radius of an object. The above described procedure leads to the new coordinate system
$$\begin{array}{cc}\hfill \stackrel{~}{r}=r_g\mathrm{ln}\frac{rr_g}{r}& r=\frac{r_g}{1e^{\stackrel{~}{r}/r_g}},\hfill \\ \hfill \stackrel{~}{\theta }=\mathrm{ln}\mathrm{tan}\frac{\theta }{2}& \theta =2\mathrm{tan}^1e^{\stackrel{~}{\theta }},\hfill \end{array}$$
(41)
in which the Schwarzschild metric becomes
$$\mathrm{d}s^2=e^{\stackrel{~}{r}/r_g}\mathrm{d}t^2\frac{e^{\stackrel{~}{r}/r_g}}{(e^{\stackrel{~}{r}/r_g}1)^4}\mathrm{d}\stackrel{~}{r}^2\frac{r_g^2}{(e^{\stackrel{~}{r}/r_g}1)^2\mathrm{cosh}^2\stackrel{~}{\theta }}\left(\mathrm{d}\stackrel{~}{\theta }^2+\mathrm{d}\phi ^2\right).$$
(42)
This metric satisfies (22). The new variable $`\stackrel{~}{r}`$ spans the interval $`(\mathrm{},0]`$ as the coordinate $`r`$ increases from $`r_g`$ to $`+\mathrm{}`$, while the new angular variable $`\stackrel{~}{\theta }`$ covers the interval $`(\mathrm{},+\mathrm{})`$. One can also easily describe the interior of the Schwarzschild solution by flipping the sign of the argument of the log in (41).
### 3.2 Radiative stability
In this section we discuss the issue of quantum loop corrections to (19). In particular, we would like to understand whether this theory is stable w.r.t. the loops. It is clear that quantum-gravitational and matter loops will generate higher dimensional operators entering the action suppressed by the UV cutoff of this theory. This is similar to any theory that is not renormalizable and should be regarded as an effective field theory below its UV cutoff (for an introduction to an effective field theory treatment of gravity, see, e.g., ). However, there is another question of a vital importance in the present context. This is whether loop corrections can generate the potential and/or kinetic terms for the Lagrange multiplier. If this happens, $`\lambda _\mu `$ cannot be regarded as an auxiliary field and all the results of the previous subsection would be ruined.
We will argue below that this problem is avoided in CGR because of the specific form of (19) which can be completed to a BRST invariant theory. Following Ref. , we introduce the variables (as in the previous sections, below we are discussing small fluctuations of the fields)
$`\gamma ^{\mu \nu }=\sqrt{g}g^{\mu \nu },\gamma _{\mu \nu }={\displaystyle \frac{g_{\mu \nu }}{\sqrt{g}}}.`$ (43)
It is straightforward to rewrite the Lagrangian (19) in terms of these variables and include the FP ghost term:
$`={\displaystyle \frac{1}{2}}\gamma ^{\mu \nu }\left(R_{\mu \nu }(\gamma )_\mu \lambda _\nu _\nu \lambda _\mu \right)+{\displaystyle \frac{i}{2}}\left(_\mu \overline{c}_\nu +_\nu \overline{c}_\mu \right)_\alpha ^{\mu \nu }c^\alpha .`$ (44)
Here the terms in the first parenthesis represent the gravitational part of (19), while the last term introduces the vector-like FP ghost and anti-ghost fields for which the operator $`_\alpha ^{\mu \nu }`$ is defined as follows:
$`_\alpha ^{\mu \nu }\gamma ^{\mu \tau }\delta _\alpha ^\nu _\tau +\gamma ^{\nu \tau }\delta _\alpha ^\mu _\tau _\alpha (\gamma ^{\mu \nu }).`$
The important point is that the standard Einstein-Hilbert action can be rewritten in terms of (43) and their first derivatives
$`{\displaystyle \frac{1}{2}}\gamma ^{\mu \nu }R_{\mu \nu }(\gamma )={\displaystyle \frac{1}{8}}_\rho \gamma ^{\mu \tau }_\sigma \gamma ^{\lambda \nu }\left(\gamma ^{\rho \sigma }\gamma _{\lambda \mu }\gamma _{\tau \nu }{\displaystyle \frac{1}{2}}\gamma ^{\rho \sigma }\gamma _{\mu \tau }\gamma _{\lambda \nu }2\delta _\tau ^\sigma \delta _\lambda ^\rho \gamma _{\mu \nu }\right).`$ (45)
The FP ghost term in (44) ensures the BRST invariance of this theory . The respective BRST transformations with a continuous Grassmann variable $`\zeta `$ are:
$`\delta \gamma ^{\mu \nu }=i_\alpha ^{\mu \nu }c^\alpha \zeta ,`$
$`\delta c^\mu =ic^\tau _\tau c^\mu \zeta ,`$
$`\delta \overline{c}^\mu =\lambda _\mu \zeta ,\delta \lambda _\mu =0.`$
Presence of this symmetry ensures that the Lagrange multiplier in (44) does not acquire the kinetic term through the loop corrections. This is because the potentially dangerous term in the effective Lagrangian
$`_\mu \lambda _\nu \gamma ^{\mu \nu }(x)\gamma ^{\alpha \beta }(0)_\alpha \lambda _\beta ,`$ (46)
is zero (up to a total derivatives) due to the transversality of a two point graviton correlation function. The latter can be seen order by order in perturbation theory. Let us for simplicity consider this for an expansion about a flat space which is a consistent solution of the theory even if the cosmological constant is present. We introduce the notations $`\gamma ^{\mu \nu }=\eta ^{\mu \nu }\phi ^{\mu \nu }`$, and look at the two-point correlation function of the $`\phi ^{\mu \nu }`$ field. This correlator has been studied in detail in the conventional approach, in which there is no Lagrange multiplier term in (44), but instead, a standard quadratic gauge fixing term $`\frac{(_\mu \gamma ^{\mu \nu })^2}{2\zeta }`$ with the gauge parameter $`\zeta `$ is introduced. Our theory, on a fixed classical background, corresponds to the limit $`\zeta 0`$. Hence, all the results concerning the quantum loops derived in the conventional approach on a fixed background are also applicable here with the condition that $`\zeta 0`$. The BRST invariance of the theory can be used to deduce the Slavnov-Taylor identities (see, e.g., ). The latter guarantee that order by order in perturbation theory the two point correlation function of the $`\phi ^{\mu \nu }`$ field is transverse (this corresponds to the $`\zeta 0`$ gauge results of Refs. )
$`\phi ^{\mu \nu }(x)\phi ^{\alpha \beta }(0)\mathrm{\Pi }^{\mu \alpha }\mathrm{\Pi }^{\nu \beta }+\mathrm{\Pi }^{\mu \beta }\mathrm{\Pi }^{\nu \alpha }\mathrm{\Pi }^{\mu \nu }\mathrm{\Pi }^{\alpha \beta },`$
$`\mathrm{where}\mathrm{\Pi }^{\mu \nu }\eta ^{\mu \nu }{\displaystyle \frac{^\mu ^\nu }{\mathrm{}}}.`$
In conclusion, the presence or absence of the FP ghosts does not affect the classical equations of motion, and for the classical analysis it is acceptable to ignore them and study the Lagrangian (19). The only quantum mechanically consistent theory is that with the FP ghosts described by (44) and all the classical results obtained above hold in this theory. Furthermore those results are stable w.r.t. the quantum loop corrections.
We would like to comment on a similar issue in the context of UGR . The Lagrangian of this theory, as given in (18), is likely to generate quadratic terms for $`\lambda `$ via loops. To see this we start with a one-loop diagram of Fig. 1.
This diagram is logarithmically divergent and will generate an additional term proportional to
$`{\displaystyle \frac{\lambda ^2}{M_{\mathrm{Pl}}^4}}\mathrm{log}\left({\displaystyle \frac{\mu _{UV}}{\mu _{IR}}}\right),`$ (47)
where $`\mu _{UV}`$ and $`\mu _{IR}`$ denote the UV and IR scales respectively, and we have restored $`M_{\mathrm{Pl}}^2`$ in front of the $`\sqrt{g}R`$ which resulted in the $`1/M_{\mathrm{Pl}}^4`$ coefficient in (47). Naively, the above term may seem to be irrelevant because of the $`1/M_{\mathrm{Pl}}^4`$ suppression. However, to get a right scaling we should restore the canonical dimensionality of $`\lambda `$. This is achieved by substitution $`\lambda M_{\mathrm{Pl}}^3\alpha `$, after which the $`\alpha `$-dependent terms in the effective Lagrangian take the form
$`M_{\mathrm{Pl}}^3\alpha (\sqrt{g}1)+M_{\mathrm{Pl}}^2\alpha ^2\mathrm{log}\left({\displaystyle \frac{\mu _{UV}}{\mu _{IR}}}\right).`$ (48)
Because of the new induced term, the field $`\alpha `$ acquires a Planckian mass and ceases to be an auxiliary field. Furthermore, we could also look at a two-loop diagram of Fig 2. A simple power-counting of the momenta running in the loops shows that this diagram generates not only a mass term for $`\alpha `$ but also its kinetic term:
$`{\displaystyle \frac{(_\mu \lambda )^2}{M_{\mathrm{Pl}}^6}}\mathrm{log}\left({\displaystyle \frac{\mu _{UV}}{\mu _{IR}}}\right)(_\mu \alpha )^2\mathrm{log}\left({\displaystyle \frac{\mu _{UV}}{\mu _{IR}}}\right).`$
Higher loops are also expected to generate similar terms, and the above arguments are hard to avoid unless the theory (18) is amended by new interactions. Perhaps the BRST invariant completion of UGR proposed in Ref. can cure this; it would be interesting to perform explicit calculations in the framework of to see whether the radiative stability is restored. Since the model of is BRST invariant one would expect a positive outcome. However, we should point out that the bosonic part of contains additional gauge fixing terms needed to completely restrict parametrization invariance of the theory, and, from this perspective, it differs from UGR.
Finally, we would like to calculate a response of the graviton field to a source. In the linearized approximation the number of physical propagating degrees of freedom of CGR should be the same as in GR. In the linearized approximation $`_\mu \phi ^{\mu \nu }=0`$ due to (22), and $`^2\lambda _\nu =0`$ due to (28). Then, the equation (21) simplifies to
$`\mathrm{}\phi _{\mu \nu }=T_{\mu \nu }.`$ (49)
The above equation is identical to an expression for the response in the Einstein theory. This of course is a consequence of the fact that the free propagator of (44), coincides with the graviton propagator of GR in the harmonic gauge $`_\mu \phi ^{\mu \nu }=_\mu h^{\mu \nu }+\frac{1}{2}^\nu h=0`$, where $`g^{\mu \nu }\eta ^{\mu \nu }h^{\mu \nu }`$.
### 3.3 Comments on massive theories
It is difficult to construct a consistent Lorentz-invariant nonlinear model of a massive graviton propagating on a Minkowski background. In the linearized approximation the only consistent massive deformation of GR is the Fierz-Pauli (F-P) model
$`={\displaystyle \frac{\sqrt{g}}{2}}R+{\displaystyle \frac{\sqrt{g}}{4}}m_g^2(h_{\mu \nu }h^{\mu \nu }h^2),`$ (50)
were $`h_{\mu \nu }g_{\mu \nu }\eta _{\mu \nu }`$. In the quadratic approximation this model describes a massive spin-2 state with 5 degrees of freedom. However, a nonlinear completion of this theory is not unique, and so far there is no known non-linear theory in four-dimensions that would be consistent. A rather general class of nonlinear completions of the F-P theory give rise to unbounded from below Hamiltonian . This manifests itself in classical instabilities for which the time scale can be substantially shorter than the scale of the inverse graviton mass .
The reason for this instability is that at the nonlinear level a ghost-like sixth “degree of freedom” shows up. This should have been expected because of the following. Ten degrees of freedom of $`g_{\mu \nu }`$ in (50) are restricted only by four independent Bianchi identities. Hence, six degrees of freedom should remain. The absence of the sixth degree of freedom in the linearized theory was just an artifact of the linearized approximation itself .
It is interesting to ask the following question: what if we start with a mass deformation of a constrained gravitational theory instead of modifying GR as in (50)? Straightforward calculations show that the mass deformation of unimodular gravity (18) or the constrained gravity (19) leads to a theory with a ghost in the linearized approximation. This ghost can be removed, at least in the linearized theory, if one considers a mass deformation of a model with both constraints $`\sqrt{g}=1`$ and $`_\mu \sqrt{g}g^{\mu \nu }=0`$. Because this set of equations imposes 5 conditions on ten components of $`g_{\mu \nu }`$, one should expect this theory to propagate 5 physical degrees of freedom. Below we will discuss the advantages as well as difficulties of this approach.
The Lagrangian of the massive “hybrid model” that combines the two constraints mentioned above takes the form:
$$=\frac{\sqrt{g}}{2}R+\frac{\sqrt{g}}{4}m_g^2(h_{\mu \nu }h^{\mu \nu }h^2)+\sqrt{g}g^{\mu \nu }_\mu \lambda _\nu \lambda (\sqrt{g}1).$$
(51)
Variation of the action w.r.t.the Lagrange multipliers $`\lambda `$ and $`\lambda _\mu `$ yields the constraints:
$$\sqrt{g}=1;_\mu \sqrt{g}g^{\mu \nu }=0.$$
(52)
We now turn to the linearized approximation about a flat space to study a graviton propagator. In this approximation the constraints (52) reduce to $`h=0`$, and $`^\mu h_{\mu \nu }=0`$ respectively. The equation of motion becomes
$$G_{\mu \nu }m_g^2h_{\mu \nu }_{(\mu }\lambda _{\nu )}+(^\alpha \lambda _\alpha )\eta _{\mu \nu }\lambda \eta _{\mu \nu }=T_{\mu \nu }.$$
(53)
The trace equation and the Bianchi identities give respectively
$$4\lambda +2^\mu \lambda _\mu =T,_\mu \lambda +^2\lambda _\mu =0,$$
which can be solved to obtain
$$\lambda _\mu =\frac{_\mu T}{6\mathrm{}},\lambda =\frac{T}{6}.$$
(54)
Substituting these solutions into (53) we find
$$\begin{array}{c}\hfill h_{\mu \nu }=\frac{T_{\mu \nu }\frac{1}{3}T\eta _{\mu \nu }}{\mathrm{}+m_g^2}+\frac{_\mu _\nu T}{3\mathrm{}(\mathrm{}+m_g^2)}.\end{array}$$
(55)
From the tensorial structure of (55) we conclude that the theory propagates five physical polarizations, as it should. Moreover, unlike F-P gravity, the propagator (55) has a well-defined $`m_g0`$ limit. This is similar to the soft behavior of massive gauge field propagator discussed in Appendix A and B, and to “softly massive gravity” emerging in higher dimensional constructions .
Could (51) be a consistent model of a massive graviton in 4D? There still is a long way to go in order to find out whether (51) is a theoretically sound theory. There are three major checks one should perform.
(i) The main problem of the F-P gravity stems from the fact that the Hamiltonian of the nonlinear theory is unbounded below. Hence, one should understand whether the same problem is evaded by the hybrid model (51). We studied this question partially and have shown that the unbounded terms that appear in the F-P massive gravity do not arise in (51). To understand this we look at the ADM decomposition of the metric
$$\begin{array}{cc}g_{\mu \nu }=\left(\begin{array}{cc}N^2\stackrel{~}{\gamma }_{ij}N^iN^j& N^j\stackrel{~}{\gamma }_{ij}\\ N^j\stackrel{~}{\gamma }_{ij}& \stackrel{~}{\gamma }_{ij}\end{array}\right),& g^{\mu \nu }=\left(\begin{array}{cc}\frac{1}{N^2}& \frac{N^i}{N^2}\\ \frac{N^j}{N^2}& \stackrel{~}{\gamma }^{ij}+\frac{N^iN^j}{N^2}\end{array}\right).\end{array}$$
In terms of the new variables the Einstein-Hilbert Lagrangian takes the form
$$\frac{\sqrt{\stackrel{~}{\gamma }}}{2}N(R^{(3)}+K_{ij}K^{ij}K^2),$$
(56)
where $`\stackrel{~}{\gamma }=det\stackrel{~}{\gamma }_{ij}`$, $`R^{(3)}`$ is the 3-dimension Ricci curvature calculated with the metric $`\stackrel{~}{\gamma }_{ij}`$ and the extrinsic curvature tensor
$$K_{ij}=\frac{1}{2N}(\dot{\stackrel{~}{\gamma }}_{ij}D_iN_jD_jN_i),$$
contains a covariant derivative $`D_i`$ with respect to the metric $`\stackrel{~}{\gamma }_{ij}`$. The problem of the F-P gravity arises because the lapse $`N`$ acquires a quadratic term in non-linear realizations of F-P gravity. Hence, it ceases to be a Lagrange multiplier and does not restrict the propagation of an extra sixth “degree of freedom” which is ghost-like. In the hybrid model, however, because of $`\sqrt{g}=1`$ we get that $`N=1/\sqrt{|det\stackrel{~}{\gamma }_{\mu \nu }|}`$. This constraint enables one to remove the sixth degree of freedom and the dangerous terms that previously led to unboundedness of the F-P Hamiltonian . Regretfully, the expression for the Hamiltonian of the hybrid model is rather complicated and it is difficult to see that there are no other sources of rapid classical instability there.
(ii) Even if the rapid classical instabilities are removed, the main question is whether the BRST invariant completion of (51) exists. The BRST symmetry, could be a guiding principle determining a unique nonlinear completion of (51)(or any other massive theory), which at this stage is completely arbitrary. One could hope that the BRST and ghost charges can be used to define the Hilbert space of physical states of the theory so that even if the Hamiltonian is not bounded below, the states of negative energy do not appear in the final states (i.e., they are projected out by the conditions $`Q_{\mathrm{BRST}}|\mathrm{Phys}=Q_{\mathrm{Ghost}}|\mathrm{Phys}=0`$ and cannot be emitted in any process). At the moment it is not clear whether such a construction is possible, but we plan to return to this set of questions in future.
(iii) The question of radiative stability of (51) is something one should worry about. In general the Lagrange multipliers of the massive theory will acquire the mass and kinetic terms, and this would lead to propagation of a new degree of freedom. Only hope here could be to complete (51) in a BRST invariant way so that the resulting theory does not generate the quadratic and higher terms for the Lagrange multipliers.
Even if all the above three issues (i-iii) are positively resolved, one needs to amend the massive model to make it consistent with the data. The point is that the scalar polarization of a massive graviton couples to sources and gives rise to contradictions with the Solar system data. In the model of Ref. the similar problems is solved due to nonlinear effects that screen the undesirable scalar polarization at observable distances . In the present model, at least naively, such a mechanism does not seem to be operative, and some new ideas are needed.
Acknowledgments
We thank Savas Dimopoulos, Gia Dvali, Andrei Gruzinov, Nemanja Kaloper, Matthew Kleban, Massimo Porrati, Adam Schwimmer, Tom Taylor, Jay Wacker and Dan Zwanziger for useful discussions and comments. The work was supported by NASA grant NNG05GH34G. GG is also partly supported by NSF grant PHY-0403005, and YS by graduate student funds provided by New York University. YS would also like to thank Andrei Gruzinov for his support through David and Lucile Packard Foundation Fellowship.
Appendix A: Qantization of CQED: Stückelberg formalism
There is another way of quantizing (3). We can restore the gauge invariance of (3) using the Stückelberg method and then follow the standard FP procedure of fixing the gauge and introducing the FP ghosts. Let us discuss this in some detail.
We start by rewriting (3) as follows:
$$=\frac{1}{4}F_{\mu \nu }^2+\lambda (^\mu B_\mu \mathrm{}\phi )+J^\mu (B_\mu _\mu \phi ),$$
(57)
here we have introduced the notations:
$$B_\mu =A_\mu +_\mu \phi \text{and}F_{\mu \nu }_\mu B_\nu _\nu B_\mu =_\mu A_\nu _\nu A_\mu .$$
(58)
A Lagrangian, similar to (57), in a context of a theory with a non-conserved current was recently discussed in . The Lagrangian (57) is invariant under the following gauge transformation
$$\delta B_\mu =_\mu \alpha \text{and}\delta \phi =\alpha ,$$
(59)
where $`\alpha `$ is an arbitrary function. To fix this freedom we choose a gauge similar to the $`R_\xi `$-gauge that eliminates a mixing terms between $`\lambda `$ and $`B`$. This can be achieved by adding into the Lagrangian the following gauge fixing term:
$$\mathrm{\Delta }_{GF}=\frac{1}{2\xi }(^\mu B_\mu \xi \lambda )^2.$$
(60)
Here, $`\xi `$ is an arbitrary gauge parameter. Furthermore, it is easy to find the FP determinant and introduce the FP ghosts into the theory. The total Lagrangian that includes the gauge fixing and FP ghost terms reads:
$$_{\mathrm{tot}}=\frac{1}{4}F_{\mu \nu }^2+\frac{1}{2\xi }(^\mu B_\mu )^2i\overline{c}\mathrm{}c\lambda \mathrm{}\phi +\frac{\xi }{2}\lambda ^2+J^\mu (B_\mu _\mu \phi ).$$
(61)
The first three terms on the right hand side of the above expression constitute a free Lagrangian of gauge-fixed QED. In addition there are other states in (61). To understand their nature, we integrate out from (61) the auxiliary field $`\lambda `$. As a result, the terms of (61) containing $`\lambda `$ get replaced as
$`\lambda \mathrm{}\phi +{\displaystyle \frac{\xi }{2}}\lambda ^2{\displaystyle \frac{1}{2\xi }}(\mathrm{}\phi )^2.`$ (62)
Because the propagator of $`\phi `$ is $`\frac{\xi }{\mathrm{}^2}=lim_{\gamma 0}\frac{1}{\gamma }(\frac{\xi }{\mathrm{}}\frac{\xi }{\mathrm{}+\gamma })`$, it is more appropriate to think of two states, described by (62), one of which has a positive-sign kinetic term and the other one is ghost-like <sup>5</sup><sup>5</sup>5This could also be understood directly from (61) by making substitutions: $`\lambda =a+b,\phi =ab`$, which generate two kinetic terms one for $`a`$ with a right sign and another one for $`b`$ with a wrong sign. These two fields will also have $`\xi `$-dependent masses and mass mixing that originate from the $`\xi \lambda ^2`$ term.. However, these field are not present in the physical on-shell spectrum of the theory. In the limit $`\xi 0`$ the field $`\phi `$ is frozen, the Lorentz-gauge fixing condition, $`^\mu B_\mu =0`$, is enforced, and the on-shell spectrum consists of two physical polarizations of a photon. Since physics cannot depend on a choice of the value of the gauge parameter $`\xi `$, the model (61) propagates two on-shell polarizations. Off-shell, however, there are four propagating degrees of freedom. The two non-physical degrees of freedom are longitudinal and time-like components of a photon, while $`\lambda `$ plays a role of the canonically conjugate momentum to $`A_0`$ (alternatively, if $`\lambda `$ is adopted as a canonical coordinate $`A_0`$ becomes its conjugate momentum).
It is straightforward to read off the propagators from (61)
$$\mathrm{\Delta }_B^{\mu \nu }=\frac{\eta ^{\mu \nu }(1+\xi )^\mu ^\nu /\mathrm{}}{\mathrm{}},\mathrm{\Delta }_\phi ^\mu ^\nu =\frac{\xi ^\mu ^\nu }{\mathrm{}^2}.$$
(63)
Adding these two contributions together we find that the dependence on the gauge parameter $`\xi `$ cancels out and we obtain (16). This confirms our previous conclusion that the $`\phi `$ field is a gauge artifact. It is also instructive to rewrite the Lagrangian (61) in terms of the original field $`A_\mu `$. The latter looks as follows:
$`_{\mathrm{tot}}={\displaystyle \frac{1}{4}}F_{\mu \nu }^2+{\displaystyle \frac{1}{2\xi }}(^\mu A_\mu )^2{\displaystyle \frac{1}{\xi }}(^\mu A_\mu )\mathrm{}\phi i\overline{c}\mathrm{}c+A_\mu J^\mu ,`$ (64)
where we have integrated out the $`\lambda `$ field. Here we see a difference from conventional QED. The additional $`\xi `$ dependent term ensures that the resulting free propagator, for arbitrary values of $`\xi `$, coincides with the Landau gauge propagator of QED. This is due to an extra field $`\phi `$ which has no kinetic term but acquires one through the kinetic mixing with the gauge field (64). The mixing term itself is gauge-parameter dependent and this is what cancels the $`\xi `$ dependence of the QED propagator.
### A.1. On massive deformation of CQED
In this subsection we study the massive deformation of the Lagrangian (3). We expect, because of the classical constraint, the massive theory to be different off-shell from the conventional massive electrodynamics (the Proca theory).
To proceed, we add to the Lagrangian (3) the following mass term:
$`\mathrm{\Delta }={\displaystyle \frac{1}{2}}m_\gamma ^2A_\mu ^2.`$ (65)
One can think of this term as arising from some higher-dimensional operator in which certain fields acquiring VEV’s generate (65) while these fields themselves become heavy and decouple from the low-energy theory.
It is straightforward to see that Eq. (4) gets modified by the mass term
$$^\nu F_{\nu \mu }+m_\gamma A_\mu +J_\mu _\mu \lambda =0,$$
(66)
while the constraint $`_\mu A^\mu =0`$ remains intact. Taking a derivative of (66) one obtains the equation of motion for $`\lambda `$
$$\mathrm{}\lambda =^\mu J_\mu .$$
(67)
From the above we find a solution with a nonzero constant background current $`_\mu \lambda =C_\mu `$ which is identical to that of the constrained massless theory.
It is interesting that the constrained massive theory also has a continuous BRST invariance if the FP ghost fields are introduced. Indeed, consider the Lagrangian
$`={\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+\lambda (_\mu A^\mu )+{\displaystyle \frac{1}{2}}m_\gamma ^2A_\mu ^2i\overline{c}\mathrm{}c.`$ (68)
As before, the presence of the FP ghost in (68) does not affect the discussions of the classical equations. On the other hand, due to these fields the Lagrangian (68) becomes invariant under the following continuous BRST transformations:
$`\delta A_\mu =i\zeta _\mu c,\delta c=0,\delta \overline{c}=\lambda \zeta ,\delta \lambda =im_\gamma ^2\zeta c,`$ (69)
where $`\zeta `$ is a continuous Grassmann parameter. Notice that $`\lambda `$ transforms to compensate for the non-invariance of the mass term<sup>6</sup><sup>6</sup>6The above transformations (69) differ from the standard BRST transformations of a massive Abelian theory in which massive FP ghost need to be introduced. We emphasize here the transformations (69) because they will be straightforwardly generalized to the non-Abelian massive case in the next section.. Moreover, $`\delta ^2A_\mu =\delta ^2c=\delta ^2\lambda =0`$, while $`\delta ^2\overline{c}0`$, but $`\delta ^3\overline{c}=0`$.
Let us now turn to the discussion of the spectrum of this theory. First we evaluate the propagator of the gauge field. For this we follow the method used in section 2. The result is:
$$\mathrm{\Delta }^{\mu \nu }=\underset{\gamma 0}{lim}\left[\eta ^{\mu \nu }\frac{1+\gamma }{\mathrm{}\gamma m_\gamma ^2}^\mu ^\nu \right]\frac{1}{\mathrm{}+m_\gamma ^2}=\frac{\eta ^{\mu \nu }^\mu ^\nu /\mathrm{}}{\mathrm{}+m_\gamma ^2}.$$
(70)
This should be contrasted with the propagator of conventional massive QED (the Proca theory):
$$\mathrm{\Delta }_{\mathrm{Proca}}^{\mu \nu }=\frac{\eta ^{\mu \nu }+^\mu ^\nu /m_\gamma ^2}{\mathrm{}+m_\gamma ^2}.$$
(71)
The key difference of (70) from (71) is the absence in (70) of the longitudinal term that is inversely proportional to the mass square. Because of this, the propagator (70) does have a good UV behavior, while the propagator (71) does not. In an Abelian theory with a conserved current this hardly matters since the longitudinal parts of the propagators do not contribute to physical amplitudes. However, this could become important for non-Abelian and gravitational theories where the matter currents are only covariantly conserved, so we carry on with this discussion.
A natural question that arises is what is the mechanism that softens the UV behavior of (70) as compared to (71)? Usually in massive gauge theories this is achieved by introducing a Higgs field that regulates the UV behavior of (71). Therefore, on top of the three physical polarizations of a massive gauge field, there should be a new state that replaces the role of the Higgs. To see this state manifestly we rewrite (70) as follows:
$$\mathrm{\Delta }^{\mu \nu }=\frac{\eta ^{\mu \nu }+^\mu ^\nu /m_\gamma ^2}{\mathrm{}+m_\gamma ^2}+\frac{^\mu ^\nu }{m_\gamma ^2\mathrm{}}.$$
(72)
The first term on the r.h.s. is just the Proca propagator of three massive polarizations; the additional term represents a new massless derivatively-coupled degree of freedom. To uncover the nature of this extra state we rewrite the Lagrangian of the constrained massive theory as follows:
$$=\frac{1}{4}F_{\mu \nu }^2+\frac{1}{2}m_\gamma ^2\left(A_\mu \frac{_\mu \lambda }{m_\gamma ^2}\right)^2\frac{(_\mu \lambda )^2}{2m_\gamma ^2}i\overline{c}\mathrm{}c+J_\mu A^\mu .$$
(73)
Notice that the quantity in the parenthesis is invariant under the BRST transformations (69), while the non-invariance of the $`\lambda `$-kinetic term under (69) is compensated by the terms coming from the FP ghost kinetic term. Defining a field
$$B_\mu A_\mu \frac{_\mu \lambda }{m_\gamma ^2},$$
(74)
we end up with the following theory:
$$=\frac{1}{4}F_{\mu \nu }^2+\frac{1}{2}m_\gamma ^2B_\mu ^2+B_\mu J^\mu \frac{1}{2m_\gamma ^2}(_\mu \lambda )^2i\overline{c}\mathrm{}c+\frac{_\mu \lambda }{m_\gamma ^2}J^\mu .$$
(75)
The first three terms on the r.h.s. of the above expression represent the Proca Lagrangian of massive electrodynamics. There are also additional terms in (75). These are the $`\lambda `$ kinetic term and FP ghost kinetic term. These two terms form a sector of the theory that is invariant under the continuous BRST transformations $`\delta \lambda =im_\gamma ^2\zeta c`$, $`\delta \overline{c}=\zeta \lambda `$ and $`\delta c=0`$. This symmetry is an exact one if the current $`J_\mu `$ is either conserved, as it is in an Abelian theory, or transforms w.r.t. BRST in an appropriate way as it will for a non-Abelian theory discussed in the next section. The $`\lambda `$ kinetic term in (75) has a wrong sign, and this state is ghost-like. While $`J_\mu `$ is conserved, $`\lambda `$ is decoupled from the rest of the physics and can be ignored for all the practical purposes. Nevertheless, it is interesting to understand whether the state $`\lambda `$ could belong to a Hilbert space of physical states. This space could be defined by introducing the BRST and Ghost charges and postulating $`Q^{\mathrm{BRST}}|\mathrm{Phys}>=0`$ and $`Q^{\mathrm{Ghost}}|\mathrm{Phys}>=0`$. The above conditions could ensure that quanta of the $`\lambda `$ field do not belong to the set of physical in or out states of the theory. However, as we discussed previously, the BRST transformations (69) are rather peculiar and the construction of the Hilbert space of physical states could differ from the conventional one.
Here we discuss the properties of the physical degrees of freedom in the Hamiltonian formalism. Let us ignore the external current $`J_\mu `$ for the time being. The canonical momenta conjugate to $`A_\mu `$ read
$`\pi _\mu {\displaystyle \frac{}{\dot{A}^\mu }}=F_{0\mu }+\lambda \delta _{0\mu }.`$ (76)
The key difference of this from QED is that (3) contains time derivative of $`A_0`$, and, as a result, the primary constraint of QED, $`\pi _0=0`$, is replaced by the relation
$`\pi _0=\lambda .`$ (77)
Thus, $`\lambda `$ is just a conjugate momentum for $`A^0`$. The above expression can be used to determine $`\lambda `$, and if so, it does not constrain $`\pi _0`$. Excluding $`\lambda `$ by (77), the extended Hamiltonian takes the form
$`={\displaystyle \frac{1}{2}}\pi _j^2+{\displaystyle \frac{1}{2}}(ϵ_{ijk}_jA_k)^2+\pi _0_jA_j+A_0\left(_j\pi _j\right).`$ (78)
To study this system further we look at the two equations governing the time evolutions of $`\pi _0`$ and $`A_0`$ as
$`\dot{\pi }_0=\{\pi _0\},\dot{A}_0=\{A_0\},`$ (79)
where $`\{..\}`$ denotes the canonical Poisson brackets. This reduces to the following relations
$`\dot{\pi }_0+_i\pi _i=0,_\mu A^\mu =0,`$ (80)
both of which were already given in a different form by equations of motion derived in section 2. Furthermore, requiring that the time derivatives of these two expressions are identically zero
$`\{\dot{\pi }_0+_i\pi _i,\}=0,\{_\mu A^\mu ,\}=0,`$ (81)
we find the following two additional equations
$`\mathrm{}\pi _0=0,\mathrm{}A^0+_i\pi _i=0.`$ (82)
Further time derivatives are identically satisfied.
From the first equation of (82) we find that $`\pi _0`$ can either be a plane wave or a trivial harmonic function (or a superposition of the two). The second equation of (82), however, dictates that $`\pi _0`$ can only assume the trivial solutions. To show this suppose that the solution for $`\pi _0`$ is a plane wave. Then, according to
$`\mathrm{}A^0=_i\pi _i=\dot{\pi }_0,`$ (83)
$`A^0`$ cannot have a nonsingular solution because of the mass-shell condition. Therefore, to avoid non-physical solutions a trivial solution for $`\pi _0`$ should be taken. As a result, equation (80) and (82) together remove $`\pi _0`$ and $`A^0`$ from the list of free variables and impose one condition on $`\pi _i`$’s and one on $`A^i`$’s, thus, reducing the number of physical degrees of freedom to two. This is consistent with the Nakanishi-Lautrup condition that prescribes to physical states to satisfy $`\lambda ^{()}|\mathrm{Phys}=\pi _0^{()}|\mathrm{Phys}=0`$.
In the massive case the Hamiltonian density is modified to be
$`={\displaystyle \frac{1}{2}}\pi _j^2+{\displaystyle \frac{1}{2}}(ϵ_{ijk}_jA_k)^2+\pi _0_jA_j+A_0\left(_j\pi _j\right)+{\displaystyle \frac{1}{2}}m_\gamma ^2(A_i^2A_0^2).`$ (84)
Time derivatives of $`\pi _0`$ and $`A_0`$ give rise to the following relations
$`\dot{\pi }_0+_i\pi _i=m_\gamma ^2A_0,_\mu A^\mu =0.`$ (85)
Their time derivatives lead to Eqs. (82). The above four equations, however, are no longer enough to remove two extra degrees of freedom. Indeed,
$`\mathrm{}A^0+m_\gamma ^2A^0=_i\pi _i=\dot{\pi }_0,`$ (86)
and since $`m_\gamma 0`$, even if $`\pi _0`$ were a plane wave solution of $`\mathrm{}\pi _0=0`$, non-singular solutions $`A^0`$ can be obtained. Therefore $`_i\pi _i`$ is no longer constrained to zero, as it was in the massless case. Hence, such a theory propagates three degrees of freedom.
Furthermore, using the relations (85) one can rewrite the Hamiltonian as follows
$$=\frac{1}{2}\pi _j^2+\frac{1}{2}(ϵ_{ijk}_jA_k)^2+\frac{1}{2}m_\gamma ^2\left(A_i\frac{_i\pi _0}{m_\gamma ^2}\right)^2\frac{(_i\pi _0)^2}{2m_\gamma ^2}\frac{\dot{\pi }_0^2}{2m_\gamma ^2}+\frac{(_i\pi _i)^2}{2m_\gamma ^2}.$$
(87)
Defining a new field
$$B_i=A_i\frac{_i\pi _0}{m_\gamma ^2},$$
(88)
we find the Hamiltonian
$$=\frac{1}{2}\pi _j^2+\frac{(_i\pi _i)^2}{2m_\gamma ^2}+\frac{1}{2}(ϵ_{ijk}_jB_k)^2+\frac{1}{2}m_\gamma ^2B_i^2\frac{1}{2m_\gamma ^2}[\pi _\lambda ^2+(_i\lambda )^2],$$
(89)
where the pairs of canonical coordinates and their conjugate momenta are $`B_i`$ and $`\pi _i`$ and $`\lambda `$ and $`\pi _\lambda `$. The $`\lambda `$ field makes a negative contribution to the energy density. However, this field is decoupled from all the sources and, thus, cannot be produced to grow the negative energy.
Appendix B: Constrained non-Abelian gauge fields
Here we generalize the constrained approach to the massless and massive non-Abelian fields. To understand the essence of these models we start with the discussion of the massless case. The Lagrangian reads as follows:
$`={\displaystyle \frac{1}{4}}F_{\mu \nu }^aF^{a\mu \nu }+A_\mu ^aJ^{a\mu }+\lambda ^a(_\mu A^\mu )^a.`$ (90)
Equations of motion for the gauge fields and Lagrange multiplier are respectively:
$`D^\nu F_{\nu \mu }^a+J_\mu ^a_\mu \lambda ^a=0,`$ (91)
$`_\mu A^{a\mu }=\mathrm{\hspace{0.17em}0}.`$ (92)
An equation for $`\lambda `$ follows by taking a covariant derivative of (91)
$`D^\mu _\mu \lambda ^a=D^\mu J_\mu ^a.`$ (93)
In a theory with a covariantly conserved source there could exist new non-trivial classical solutions for $`\lambda `$ such that $`_\mu \lambda `$ is also covariantly conserved. One particular solution is a path-ordered Wilson line
$`_\mu \lambda (x)=𝒫\mathrm{exp}\left(ig{\displaystyle _y^x}A_\mu (z)𝑑z^\mu \right)_\mu \lambda (y).`$ (94)
Thus, if there is a non-zero external colored current solution $`_\mu \lambda `$ at some point, then its value at any other point can be calculated according to (94).
The next issue to be address is that of quantum consistency of this approach. Again an important fact is that the Lagrangian (90) can be completed to a BRST invariant form:
$`={\displaystyle \frac{1}{4}}F_{\mu \nu }^aF^{a\mu \nu }+A_\mu ^aJ^{a\mu }+\lambda ^a(^\mu A_\mu ^a)i\overline{c}^a^\mu D_\mu ^{ab}c^b.`$ (95)
Since at the classical solutions that we consider the FP ghost fields vanish, (95) recovers all the classical results of (90). The explicit BRST transformations leaving (95) invariant are:
$`\delta A_\mu ^a=i\zeta D_\mu ^{ab}c^b,\delta c^a={\displaystyle \frac{i}{2}}g\zeta f^{abd}c^bc^d,\delta \overline{c}^a=\lambda ^a\zeta ,\delta \lambda ^a=0.`$ (96)
The expression for the BRST and Ghost currents, as well as the subsidiary conditions that guarantee the unitarity and completeness of the physical Hilbert space of states are the standard ones .
As in the case of a photon, the massless free propagator reads:
$`\mathrm{\Delta }_{\mu \nu }^{ab}=\delta ^{ab}{\displaystyle \frac{\eta _{\mu \nu }_\mu _\nu /\mathrm{}}{\mathrm{}}}.`$ (97)
Finally, due to the presence of the BRST symmetry the Lagrange multiplier does not acquire kinetic term via the loop corrections. This is because the two-point correlation functions of the $`A_\mu `$ field is guarantied to be transverse due to the presence of the FP ghost and the transverse structure of the tree-level propagator (97)
$`^\mu \lambda ^aA_\mu ^aA_\nu ^b^\nu \lambda ^b^\mu \lambda ^a\delta ^{ab}(\eta _{\mu \nu }\mathrm{}_\mu _\nu )^\nu \lambda ^b=0,`$ (98)
to all orders of perturbation theory.
### B.1. Stückelberg formalism
Instead of using BRST symmetry to arrive from (90) to (95), here we follow a conventional route. First we restore gauge symmetry of (90) and then fix that restored gauge invariance and introduce the appropriate FP ghosts. Let us start by defining new variables:
$`igA_\mu =U^+D_\mu U,\mathrm{where}U=e^{it^a\pi ^a},D_\mu =_\mu +igB_\mu ,`$ (99)
and rewrite the Lagrangian (90) in the following form
$`={\displaystyle \frac{1}{2}}\mathrm{Tr}\left({\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+{\displaystyle \frac{\lambda ^\mu [U^+D_\mu U]}{ig}}+{\displaystyle \frac{U^+D_\mu U}{ig}}J^\mu \right).`$ (100)
The above Lagrangian is gauge invariant under the local transformations of $`Ue^{i\alpha ^at^a}U`$ and $`B_\mu e^{i\alpha ^at^a}B_\mu e^{i\alpha ^at^a}+\frac{i}{g}[_\mu e^{i\alpha ^at^a}]e^{i\alpha ^at^a}`$, where $`t^a`$ denote the generators of a local gauge group, and $`\alpha ^a`$’s are gauge transformation functions<sup>7</sup><sup>7</sup>7 Notice that the current $`J_\mu `$ and the Lagrange multiplier $`\lambda `$ are not supposed to transform under these gauge transformations. This can be achieved by rewriting the fundamental fields out of which $`J_\mu `$ is constructed (as well as rewriting $`\lambda `$) in terms of new fields rescaled by $`U`$’s. Under the gauge transformations the new field transform in a conventional way but their variance is compensated by transformations of $`U`$’s, so that $`J_\mu `$ and $`\lambda `$ stay invariant.. We can proceed further and fix this gauge freedom by introducing into the Lagrangian the following term
$`{\displaystyle \frac{1}{4\xi }}\mathrm{Tr}(_\mu B^\mu \xi \lambda )^2.`$ (101)
This term removes quadratic mixing between $`B_\mu `$ and $`\lambda `$, and after integrating out $`\lambda `$ and rewriting it back in terms of the $`A_\mu `$ field, the result reads as follows:
$`={\displaystyle \frac{1}{2}}\mathrm{Tr}\left({\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+{\displaystyle \frac{1}{2\xi }}(_\mu A^\mu )^2+{\displaystyle \frac{1}{\xi }}(_\mu A^\mu )(^\mu \mathrm{\Pi }_\mu )i\overline{c}_\mu D_\mu c+A_\mu J^\mu \right),`$ (102)
where $`\mathrm{\Pi }_\mu =UA_\mu U^++\frac{i}{g}[_\mu U]U^+A_\mu `$. All the terms of the Lagrangian (102), except the third one, is what one gets in the conventional approach. The third, gauge dependent term, has a structure that guarantees that the free propagator for an arbitrary value of $`\xi `$ coincides with the Landau gauge propagator of the conventional approach. This is due to the extra fields $`\pi ^a`$’s which by themselves have no kinetic term but acquire one through the gauge-parameter-dependent kinetic mixing term with $`A_\mu `$ in (102).
### B.2. Comments on massive constrained non-Abelian fields
Let us start with a component form of the constrained massive Lagrangian
$`={\displaystyle \frac{1}{4}}F_{\mu \nu }^aF^{a\mu \nu }+A_\mu ^aJ^{a\mu }+\lambda ^a(_\mu A^\mu )^a+{\displaystyle \frac{1}{2}}M^2A_\mu ^aA_\mu ^ai\overline{c}^a^\mu D_\mu ^{ab}c^b.`$ (103)
It is interesting that the action of this theory is BRST invariant under the following transformations (the Lagrangian transforms as a total derivative)
$`\delta A_\mu ^a=i\zeta D_\mu ^{ab}c^b,\delta c^a={\displaystyle \frac{i}{2}}g\zeta f^{abd}c^bc^d,\delta \overline{c}^a=\lambda ^a\zeta ,\delta \lambda ^a=iM^2\zeta c^a.`$ (104)
Notice that the Lagrange multiplier does transforms w.r.t. the BRST and compensates for the non-invariance of the mass term. The above BRST transformations have the following peculiar properties:
$`\delta ^2A_\mu ^a=\delta ^2c^a=0;\delta ^2\overline{c}^a0\delta ^2\lambda ^a;\delta ^4\overline{c}^a=\delta ^3\lambda ^a=0.`$ (105)
The respective BRST and Ghost currents (the latter is due to the ghost-rescaling symmetry) $`\overline{c}e^\alpha \overline{c},ce^\alpha c`$) can be derived
$`J_\mu ^{\mathrm{BRST}}=F_{\mu \nu }D_\nu c+\lambda D_\mu c{\displaystyle \frac{ig}{2}}_\mu (\overline{c}^a)f^{abd}c^bc^d,J_\mu ^{\mathrm{Ghost}}=i(\overline{c}^aD_\mu c^a(_\mu \overline{c}^a)c^a).`$ (106)
The physical states could be defined in analogy with the standard approach
$`Q_{\mathrm{BRST}}|\mathrm{Phys}>=0,Q_{\mathrm{Ghost}}|\mathrm{Phys}>=0,`$ (107)
where the BRST and Ghost charges are defined as $`Q_{\mathrm{BRST}}=d^3xJ_0^{\mathrm{BRST}}(t,x)`$ and $`Q_{\mathrm{Ghost}}=d^3xJ_0^{\mathrm{Ghost}}(t,x)`$. This suggest that the $`\lambda `$ state, which upon the diagonalization of the Lagrangian acquires a ghost-like kinetic term, should not be a part of the physical Hilbert space of in and out states of the theory. The fictitious $`\lambda `$ particle should be allowed to propagate as an intermediate state in Feynman diagrams softening the UV behavior of the theory, however, it cannot be emitted as a final in or out state of the theory. In this respect it should be similar to a FP ghost. However, because the peculiar properties of the above BRST transformations (105), it still remains to be seen that for a rigorous construction of a positive semi-definite norm Hilbert space of states with unitary S-matrix elements the conditions (107) are enough. Besides the quantum effects, one should make sure that rapid classical instabilities are also removed. These may require further modification of the model. Detailed studies of this issue will be presented elsewhere.
Having the formalism of the previous subsection developed it is easy now to discuss the spectrum of massive theory. First we restore the gauge invariance of the massive Lagrangian (103) by using the variables (99), and then gauge fix it by (101). The resulting Lagrangian is
$`={\displaystyle \frac{1}{2}}\mathrm{Tr}\left({\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+\lambda ^\mu \left({\displaystyle \frac{U^+D_\mu U}{ig}}\right)+{\displaystyle \frac{M^2}{2}}\left({\displaystyle \frac{U^+D_\mu U}{ig}}\right)^2+{\displaystyle \frac{1}{2\xi }}(_\mu B^\mu \xi \lambda )^2\right).`$ (108)
The mixing term between the gauge field and Goldstones is canceled and we integrate out the $`\lambda `$ field. The resulting Lagrangian reads:
$`={\displaystyle \frac{1}{2}}\mathrm{Tr}\left({\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }(A)+{\displaystyle \frac{1}{2}}M^2A_\mu ^2+{\displaystyle \frac{1}{2\xi }}(_\mu A^\mu )^2+{\displaystyle \frac{1}{\xi }}(_\mu A^\mu )(^\mu \mathrm{\Pi }_\mu )+A_\mu J^\mu \right),`$ (109)
where, as before, we defined $`\mathrm{\Pi }_\mu =UA_\mu U^++\frac{i}{g}[_\mu U]U^+A_\mu `$.
A remarkable feature of this model is that the propagator takes the form:
$`\mathrm{\Delta }_{\mu \nu }^{ab}=\delta ^{ab}\left(\eta _{\mu \nu }{\displaystyle \frac{(1+\xi )}{\mathrm{}\xi M^2}}_\mu _\nu \right){\displaystyle \frac{1}{\mathrm{}+M^2}}\delta ^{ab}{\displaystyle \frac{\xi _\mu _\nu }{(\mathrm{}\xi M^2)\mathrm{}}},`$ (110)
which has a smooth UV behavior and non-singular massless limit.
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# Weak curvature conditions and functional inequalities
## 1. Democratic couplings
We recall some notation from \[10, Section 2\]. Let $`(X,d)`$ be a compact length space; see for background material on such spaces. (Many results of the paper extend to the locally compact case, but for simplicity we will assume compactness.) Let $`\mathrm{\Gamma }`$ denote the set of minimizing constant-speed geodesics $`\gamma :[0,1]X`$, with the time-$`t`$ evaluation map denoted by $`e_t:\mathrm{\Gamma }X`$. The endpoint map $`E:\mathrm{\Gamma }X\times X`$ is $`E=(e_0,e_1)`$.
We let $`P(X)`$ denote the Borel probability measures on $`X`$. A transference plan $`\pi P(X\times X)`$ between $`\mu _0,\mu _1P(X)`$ is a probability measure whose marginals are $`\mu _0`$ and $`\mu _1`$. The 2-Wasserstein space $`P_2(X)`$ is $`P(X)`$ equipped with the metric of optimal transport, $`W_2(\mu _0,\mu _1)=[inf_{X\times X}d(x_0,x_1)^2𝑑\pi (x_0,x_1)]^{1/2}`$. Here the infimum is over transference plans between $`\mu _0`$ and $`\mu _1`$. A transference plan is said to be optimal if it achieves the infimum in the above variational problem. When such a $`\pi `$ is given, we can disintegrate it with respect to its first marginal $`\mu _0`$ or its second marginal $`\mu _1`$. We write this in a slightly informal way:
(1.1)
$$d\pi (x_0,x_1)=d\pi (x_1|x_0)d\mu _0(x_0)=d\pi (x_0|x_1)d\mu _1(x_1).$$
A dynamical transference plan consists of a transference plan $`\pi `$ and a Borel measure $`\mathrm{\Pi }`$ on $`\mathrm{\Gamma }`$ such that $`E_{}\mathrm{\Pi }=\pi `$; it is said to be optimal if $`\pi `$ itself is. If $`\mathrm{\Pi }`$ is a dynamical transference plan then for $`t[0,1]`$, we put $`\mu _t=(e_t)_{}\mathrm{\Pi }`$. Then $`\mathrm{\Pi }`$ is optimal if and only if $`\{\mu _t\}_{t[0,1]}`$ is a Wasserstein geodesic \[10, Lemma 2.4\]. Any Wasserstein geodesic arises from some optimal dynamical transference plan in this way \[10, Proposition 2.10\].
###### Definition 1.2 (democratic coupling).
Given $`\mu P(X)`$, the democratic transference plan between $`\mu `$ and itself is the tensor product $`\mu \mu P(X\times X)`$. A probability measure $`\mathrm{\Pi }P(\mathrm{\Gamma })`$ is said to be a dynamical democratic transference plan between $`\mu `$ and itself if $`E_{}\mathrm{\Pi }=\mu \mu `$.
###### Example 1.3.
Let $`(X,d)`$ be equipped with a reference measure $`\nu P(X)`$. Suppose that one has almost-everywhere uniqueness of geodesics in the following sense :
(1.4)
$$\{\begin{array}{cc}\text{For }\nu \nu \text{-almost all }(x_0,x_1)X\times X\text{,}\hfill & \\ \text{there is a unique geodesic }\gamma =\gamma _{x_0,x_1}\mathrm{\Gamma }\text{ with }\gamma (0)=x_0\text{ and }\gamma (1)=x_1.\hfill & \end{array}$$
Define $`S:X\times X\mathrm{\Gamma }`$ measurably by $`S(x_0,x_1)=\gamma _{x_0,x_1}`$. If $`\mu `$ is absolutely continuous with respect to $`\nu `$ then there is a unique dynamical democratic transference plan between $`\mu `$ and itself given by
(1.5)
$$\mathrm{\Pi }=S_{}(\mu \mu ).$$
###### Definition 1.6.
A compact measured length space $`(X,d,\nu )`$ is a compact length space $`(X,d)`$ equipped with a Borel probability measure $`\nu P(X)`$. Given $`C>0`$, the triple $`(X,d,\nu )`$ is said to satisfy $`\mathrm{DM}(\mathrm{C})`$ if for each ball $`B=B_r(x)X`$ with $`\nu [B]>0`$, there is a dynamical democratic transference plan $`\mathrm{\Pi }`$ from $`\mu =\frac{1_B}{\nu [B]}\nu `$ to itself with the property that if we put $`\mu _t=(e_t)_{}\mathrm{\Pi }`$ then
(1.7)
$$_0^1\mu _t𝑑t\frac{C}{\nu [B]}\nu .$$
We recall that a sequence $`\{(X_i,d_i)\}_{i=1}^{\mathrm{}}`$ of compact metric spaces converges to a compact metric space $`(X,d)`$ in the Gromov-Hausdorff topology if there is a sequence of Borel maps $`f_i:X_iX`$ and a sequence of positive numbers $`ϵ_i0`$ so that
1. For all $`x_i,x_i^{}X_i`$, $`|d_X(f_i(x_i),f_i(x_i^{}))d_{X_i}(x_i,x_i^{})|ϵ_i`$.
2. For all $`xX`$ and all $`i`$, there is some $`x_iX_i`$ such that $`d_X(f_i(x_i),x)ϵ_i`$.
The maps $`f_i`$ are called $`ϵ_i`$-approximations. If each $`(X_i,d_i)`$ is a length space then so is $`(X,d)`$. A sequence $`\{(X_i,d_i,\nu _i)\}_{i=1}^{\mathrm{}}`$ of compact measured length spaces converges to $`(X,d,\nu )`$ in the measured Gromov-Hausdorff topology if in addition, one can choose the $`f_i`$’s so that $`lim_i\mathrm{}(f_i)_{}\nu _i=\nu `$ in the weak-$``$ topology on $`P(X)`$.
###### Theorem 1.8.
Suppose that $`\{(X_i,d_i,\nu _i)\}_{i=1}^{\mathrm{}}`$ is a sequence of compact measured length spaces that converges to $`(X,d,\nu )`$ in the measured Gromov-Hausdorff topology. Suppose that each ball $`B`$ in $`X`$ has $`\nu [B]=\nu [\overline{B}]`$. If each $`(X_i,d_i,\nu _i)`$ satisfies $`\mathrm{DM}(\mathrm{C})`$ then so does $`(X,d,\nu )`$.
###### Proof.
Let $`f_i:X_iX`$ be a sequence of $`\epsilon _i`$-approximations, with $`\epsilon _i0`$, that realizes the Gromov-Hausdorff convergence. Let $`B=B_r(x)`$ be a ball in $`X`$ with $`\nu [B]>0`$. For each $`i`$, choose a point $`x_iX_i`$ so that $`d_X(f_i(x_i),x)\epsilon _i`$ and put $`B_i=B_r(x_i)`$. By elementary estimates,
(1.9)
$$f_i^1(B_{r2\epsilon _i}(x))B_if_i^1(f_i(B_i))f_i^1(B_{r+2\epsilon _i}(x)).$$
Combining this with the convergence of $`(f_i)_{}\nu _i`$ to $`\nu `$, and the fact that $`\nu [B]=\nu [\overline{B}]`$, it is easy to deduce that $`\nu _i[B_i]\nu [B]`$ (and in particular $`\nu _i[B_i]>0`$ for $`i`$ large enough). A similar “squeezing” argument shows that $`_{B_i}\phi f_i𝑑\nu _i`$ converges to $`_B\phi 𝑑\nu `$ for all nonnegative continuous functions $`\phi `$. As a consequence, if we put $`\mu =\frac{1_B}{\nu [B]}\nu `$ and (for $`i`$ large enough) $`\mu _i=\frac{1_{B_i}}{\nu _i[B_i]}\nu _i`$ then $`lim_i\mathrm{}(f_i)_{}\mu _i=\mu `$ in the weak-$``$ topology.
For each $`i`$, we can introduce a dynamical democratic transference plan $`\mathrm{\Pi }_i`$ as in Definition 1.6, relative to the ball $`B_i`$. We write $`\mu _{i,t}=(e_t)_{}\mathrm{\Pi }_i`$. From Theorem A.45, there is a dynamical transference plan $`\mathrm{\Pi }P(\mathrm{\Gamma }(X))`$, with associated transference plan $`\pi =E_{}\mathrm{\Gamma }`$ and measures $`\mu _t=(e_t)_{}\mathrm{\Pi }`$, so that $`lim_i\mathrm{}(f_i,f_i)_{}\pi _i=\pi `$ and $`lim_i\mathrm{}(f_i)_{}\mu _{i,t}=\mu _t`$.
For any $`F_1,F_2C(X)`$, we have
(1.10) $`{\displaystyle _{X\times X}}F_1(x)F_2(y)𝑑\pi (x,y)`$ $`=\underset{i\mathrm{}}{lim}{\displaystyle _{X\times X}}F_1(x)F_2(y)d((f_i,f_i)_{}\pi _i)(x,y)`$
$`=\underset{i\mathrm{}}{lim}{\displaystyle _{X_i\times X_i}}(f_i^{}F_1)(x_i)(f_i^{}F_2)(y_i)𝑑\pi _i(x_i,y_i)`$
$`=\underset{i\mathrm{}}{lim}{\displaystyle _{X_i}}(f_i^{}F_1)𝑑\mu _i{\displaystyle _{X_i}}(f_i^{}F_2)𝑑\mu _i`$
$`=\underset{i\mathrm{}}{lim}{\displaystyle _X}F_1d(f_i)_{}\mu _i{\displaystyle _X}F_2d(f_i)_{}\mu _i`$
$`={\displaystyle _X}F_1𝑑\mu {\displaystyle _X}F_2𝑑\mu .`$
Thus $`E_{}\mathrm{\Pi }=\pi =\mu \mu `$, so $`\mathrm{\Pi }`$ is still a dynamical democratic transference plan.
It remains to check (1.7). Let $`\phi `$ be a nonnegative continuous function on $`X`$. For large $`i`$, we can write
(1.11)
$$_0^1_{X_i}(f_i)^{}\phi 𝑑\mu _{i,t}\frac{C}{\nu _i[B_i]}_{X_i}(f_i)^{}\phi 𝑑\nu _i.$$
In other words,
(1.12)
$$_0^1_X\phi d(f_i)_{}\mu _{i,t}\frac{C}{\nu _i[B_i]}_X\phi d(f_i)_{}\nu _i.$$
Passing to the limit as $`i\mathrm{}`$ gives
(1.13)
$$_0^1_X\phi 𝑑\mu _t\frac{C}{\nu [B]}_X\phi 𝑑\nu .$$
Since $`\phi `$ is arbitrary, this proves (1.7). ∎
## 2. From $`\mathrm{DM}`$ to a scale-invariant local Poincaré inequality
We first recall some notation and definitions about metric-measure spaces $`(X,d,\nu )`$. If $`B=B_r(x)`$ is a ball in $`X`$ then we write $`\lambda B`$ for $`B_{\lambda r}(x)`$. The measure $`\nu `$ is said to be doubling if there is some $`D>0`$ so that for all balls $`B`$, $`\nu [2B]D\nu [B]`$. The constant $`D`$ is called the doubling constant. An upper gradient for a function $`uC(X)`$ is a Borel function $`g:X[0,\mathrm{}]`$ such that for each curve $`\gamma :[0,1]X`$ with finite length $`L(\gamma )`$ and constant speed,
(2.1)
$$\left|u(\gamma (1))u(\gamma (0))\right|L(\gamma )_0^1g(\gamma (t))𝑑t.$$
If $`u`$ is Lipschitz then an example of an upper gradient is obtained by defining
(2.2)
$$g(x)=\{\begin{array}{cc}\underset{yx}{lim\; sup}\frac{|u(y)u(x)|}{d(x,y)}\text{if }x\text{ is not isolated}\hfill & \\ & \\ g(x)=0\text{if }x\text{ is isolated}.\hfill & \end{array}$$
There are many forms of local Poincaré inequalities. The strongest one, in a certain sense, is as follows :
###### Definition 2.3.
A metric-measure space $`(X,d,\nu )`$ admits a local Poincaré inequality if there are constants $`\lambda 1`$ and $`P<\mathrm{}`$ such that for all $`uC(X)`$ and $`B=B_r(x)`$ with $`\nu [B]>0`$, each upper gradient $`g`$ of $`u`$ satisfies
(2.4)
$$\underset{B}{}|uu_B|d\nu Pr\underset{\lambda B}{}gd\nu .$$
Here the barred integral is the average (with respect to $`\nu `$) and $`u_B`$ is the average of $`u`$ over the ball $`B`$. In the case of a length space, the local Poincaré inequality as formulated in Definition 2.3 actually implies stronger inequalities, for which we refer to \[6, Chapters 4 and 9\]. It is known that the property of admitting a local Poincaré inequality is preserved under measured Gromov-Hausdorff limits . (This is an extension of the earlier result \[2, Theorem 9.6\]; Cheeger informs us that in unpublished work he also proved the extension.)
The use of a condition like $`\mathrm{DM}`$ to prove a local Poincaré inequality is implicit in the work of Cheeger and Colding \[3, Proof of Theorem 2.11\]. The next theorem makes the link explicit.
###### Theorem 2.5.
If the compact measured length space $`(X,d,\nu )`$ satisfies $`\mathrm{DM}(\mathrm{C})`$ and $`\nu `$ is doubling, with doubling constant $`D`$, then $`(X,d,\nu )`$ admits a local Poincaré inequality (2.4) with $`\lambda =2`$ and $`P=\mathrm{\hspace{0.25em}2}CD`$.
###### Proof.
Let $`x_0`$ be a given point in $`X`$. Given $`r>0`$, write $`B=B_r(x_0)`$. Note that from the doubling condition, $`\nu [B]>0`$. Put $`\mu =\frac{1_B}{\nu [B]}\nu `$.
For $`y_0X`$, we have
(2.6)
$$u(y_0)u_B=_X\left(u(y_0)u(y_1)\right)𝑑\mu (y_1).$$
Then
(2.7)
$$\underset{B}{}|uu_B|d\nu =_X\left|u(y_0)u_B\right|𝑑\mu (y_0)_{X\times X}|u(y_0)u(y_1)|𝑑\mu (y_0)𝑑\mu (y_1).$$
Next, we estimate $`|u(y_0)u(y_1)|`$ in terms of a geodesic path $`\gamma `$ joining $`y_0`$ to $`y_1`$, where $`y_0,y_1B`$. The length of such a geodesic path is clearly less than $`2r`$. Then, from the definition of an upper gradient,
(2.8)
$$|u(y_0)u(y_1)|\mathrm{\hspace{0.25em}2}r_0^1g(\gamma (t))𝑑t.$$
Now let $`\mathrm{\Pi }`$ be a dynamical democratic transference plan between $`\mu `$ and itself satisfying (1.7). Integrating (2.8) against $`\mathrm{\Pi }`$ gives, with $`\mu _t=(e_t)_{}\mathrm{\Pi }`$,
(2.9) $`{\displaystyle _{X\times X}}|u(y_0)u(y_1)|𝑑\mu (y_0)𝑑\mu (y_1)`$ $`{\displaystyle _\mathrm{\Gamma }}\left(2r{\displaystyle _0^1}g(\gamma (t))𝑑t\right)𝑑\mathrm{\Pi }(\gamma )`$
$`=2r{\displaystyle _0^1}\left({\displaystyle _\mathrm{\Gamma }}g(\gamma (t))𝑑\mathrm{\Pi }(\gamma )\right)𝑑t`$
$`=2r{\displaystyle _0^1}\left({\displaystyle _\mathrm{\Gamma }}(ge_t)𝑑\mathrm{\Pi }\right)𝑑t`$
$`=2r{\displaystyle _0^1}\left({\displaystyle _X}gd(e_t)_{}\mathrm{\Pi }\right)𝑑t`$
$`=2r{\displaystyle _0^1}{\displaystyle _X}g𝑑\mu _t𝑑t.`$
Combining this with (2.7), we conclude that
(2.10)
$$\underset{B}{}|uu_B|d\nu 2r_0^1_Xg𝑑\mu _t𝑑t.$$
However, a geodesic joining two points in $`B`$ cannot leave $`2B`$, so (2.10) and $`\mathrm{DM}(\mathrm{C})`$ together imply that
(2.11)
$$\underset{B}{}|uu_B|d\nu \frac{2Cr}{\nu [B]}_{2B}g𝑑\nu .$$
By the doubling property, $`\frac{1}{\nu [B]}\frac{D}{\nu [2B]}`$. The conclusion is that
(2.12)
$$\underset{B}{}|uu_B|d\nu \mathrm{\hspace{0.25em}2}CDr\underset{2B}{}gd\nu .$$
This proves the theorem. ∎
## 3. Nonnegative $`N`$-Ricci curvature and $`\mathrm{DM}`$
In this section we show that a measured length space with nonnegative $`N`$-Ricci curvature satisfies the condition $`\mathrm{DM}(2^\mathrm{N})`$ as soon as geodesics are almost-everywhere unique.
We use the notion of nonnegative $`N`$-Ricci curvature from \[10, Definition 5.12\]. This is the same as the case $`K=0`$ of Section 4. We will be concerned here with the case $`N<\mathrm{}`$.
###### Theorem 3.1.
Suppose that a compact measured length space $`(X,d,\nu )`$ has nonnegative $`N`$-Ricci curvature, and that minimizing geodesics in $`X`$ are almost-everywhere unique in the sense of (1.4). Then $`(X,d,\nu )`$ satisfies $`\mathrm{DM}(2^\mathrm{N})`$.
###### Remark 3.2.
If $`(X,d,\nu )`$ has nonnegative $`N`$-Ricci curvature and $`(X,d)`$ is nonbranching then minimizing geodesics in $`X`$ are almost-everywhere unique .
Before proving Theorem 3.1, we state a corollary:
###### Corollary 3.3.
If a compact measured length space $`(X,d,\nu )`$ has nonnegative $`N`$-Ricci curvature and almost-everywhere unique geodesics then it satisfies the local Poincaré inequality of Definition 2.3 with $`\lambda =\mathrm{\hspace{0.25em}2}`$ and $`P=\mathrm{\hspace{0.25em}2}^{2N+1}`$. More generally, if $`\{(X_i,d_i,\nu _i)\}_{i=1}^{\mathrm{}}`$ is a sequence of compact measured length spaces with nonnegative $`N`$-Ricci curvature and almost-everywhere unique geodesics, and it converges in the measured Gromov-Hausdorff topology to $`(X,d,\nu )`$, then $`(X,d,\nu )`$ satisfies the local Poincaré inequality of Definition 2.3 with $`\lambda =\mathrm{\hspace{0.25em}2}`$ and $`P=\mathrm{\hspace{0.25em}2}^{2N+1}`$.
###### Proof of Corollary 3.3.
First, \[10, Theorem 5.19\] implies that $`(X,d,\nu )`$ has nonnegative $`N`$-Ricci curvature. Then the generalized Bishop-Gromov inequality of \[10, Theorem 5.31\] implies that $`\nu [B]=\nu [\overline{B}]`$ for each ball $`B`$ whose center belongs to the support of $`\nu `$. It also implies that $`\nu `$ is doubling with constant $`D=2^N`$. Then the conclusion follows from Theorems 1.8, 2.5 and 3.1. ∎
As preparation for the proof of Theorem 3.1, we first prove a lemma concerning optimal transport to delta functions.
###### Lemma 3.4.
Under the hypotheses of Theorem 3.1, let $`B`$ be an open ball in $`X`$. Then for almost all $`x_0\mathrm{supp}(\nu )`$, the (unique) Wasserstein geodesic $`\{\mu _t\}_{t[0,1]}`$ joining $`\mu _0=\delta _{x_0}`$ to $`\mu _1=\frac{1_B}{\nu [B]}\nu `$ can be written as $`\mu _t=\rho _t\nu `$ with
(3.5)
$$\rho _t(x)\frac{1}{t^N\nu [B]}.$$
###### Proof of Lemma 3.4.
Let $`\mathrm{\Pi }`$ be the (unique) optimal dynamical transference plan giving rise to $`\{\mu _t\}_{t[0,1]}`$. Let $`Y_0`$ be the set of $`x_0\mathrm{supp}(\nu )`$ such that for $`\nu `$-almost every $`xX`$ there is a unique geodesic joining $`x_0`$ to $`x`$. By assumption, $`Y_0`$ has full $`\nu `$-measure. Consider $`x_0Y_0`$. Given $`t(0,1)`$, $`yX`$ and $`r>0`$, let $`Z`$ be the set of endpoints $`\gamma (1)`$ of geodesics $`\gamma `$ with $`\gamma (0)=x_0`$, $`\gamma (t)B_r(y)`$ and $`\gamma (1)B`$. Then
(3.6) $`\mu _t[B_r(y)]`$ $`=((e_t)_{}\mathrm{\Pi })[B_r(y)]=\mathrm{\Pi }[e_t^1(B_r(y))]=\mathrm{\Pi }[e_1^1(Z)]=((e_1)_{}\mathrm{\Pi })[Z]`$
$`=\mu _1[Z]={\displaystyle \frac{\nu [Z]}{\nu [B]}}.`$
If $`\nu [Z]=\mathrm{\hspace{0.25em}0}`$ then $`\mu _t[B_r(y)]=\mathrm{\hspace{0.25em}0}`$. Otherwise, put $`\mu _1^{}=\frac{1_Z}{\nu [Z]}\nu `$ and let $`\{\mu _t^{}\}_{t[0,1]}`$ be the (unique) Wasserstein geodesic joining $`\mu _0^{}=\delta _{x_0}`$ to $`\mu _1^{}`$. By the construction of $`Z`$, $`\mu _t^{}[B_r(y)]=\mathrm{\hspace{0.25em}1}`$.
Put
(3.7)
$$\varphi (s)=_X(\rho _s^{})^{11/N}𝑑\nu $$
where $`\rho _s^{}`$ is the density in the absolutely continuous part of the Lebesgue decomposition of $`\mu _s^{}`$ with respect to $`\nu `$. As $`(X,d,\nu )`$ has nonnegative $`N`$-Ricci curvature, $`\varphi `$ satisfies a convexity inequality on $`[0,1]`$. (We use here the uniqueness of the Wasserstein geodesic $`\{\mu _t^{}\}_{t[0,1]}`$. From the definition of nonnegative $`N`$-Ricci curvature, a priori one only has convexity along some Wasserstein geodesic from $`\mu _0^{}`$ to $`\mu _1^{}`$.)
As $`\varphi (0)=\mathrm{\hspace{0.25em}0}`$ and $`\varphi (1)=\nu [Z]^{\frac{1}{N}}`$, we obtain
(3.8)
$$\varphi (t)t\nu [Z]^{\frac{1}{N}}.$$
On the other hand, by Jensen’s inequality
(3.9) $`\varphi (t)`$ $`=\nu [B_r(y)]\left({\displaystyle \frac{1}{\nu [B_r(y)]}}{\displaystyle _{B_r(y)}}(\rho _t^{})^{1\frac{1}{N}}𝑑\nu \right)\nu [B_r(y)]\left({\displaystyle \frac{1}{\nu [B_r(y)]}}{\displaystyle _{B_r(y)}}\rho _t^{}𝑑\nu \right)^{1\frac{1}{N}}`$
$`\nu [B_r(y)]^{\frac{1}{N}}\mu _t^{}[B_r(y)]^{1\frac{1}{N}}=\nu [B_r(y)]^{\frac{1}{N}}.`$
This, combined with (3.8), gives
(3.10)
$$t\nu [Z]^{\frac{1}{N}}\nu [B_r(y)]^{\frac{1}{N}},$$
Then by (3.6),
(3.11)
$$\frac{\mu _t[B_r(y)]}{\nu [B_r(y)]}=\frac{\frac{\nu [Z]}{\nu [B_r(y)]}}{\nu [B]}\frac{1}{t^N\nu [B]}.$$
Since this is true for any ball centered at any $`y`$ and since balls generate the Borel $`\sigma `$-algebra, we deduce that $`\mu _t\frac{\nu }{t^N\nu [B]}`$; so $`\mu _t`$ is absolutely continuous with respect to $`\nu `$ and its density is bounded above by $`\frac{1}{t^N\nu [B]}`$. ∎
###### Proof of Theorem 3.1.
As in Definition 1.6, let $`B`$ be a ball in $`X`$ with $`\nu [B]>0`$ and put $`\mu =\frac{1_B}{\nu [B]}\nu `$. As in Example 1.3, there is a unique dynamical democratic transference plan from $`\mu `$ to itself. We want to show that the condition of Definition 1.6 is satisfied.
Define $`\mu _t^{x_0}`$ as in Lemma 3.4, with density $`\rho _t^{x_0}`$. From Lemma 3.4, $`\rho _t^{x_0}\frac{1}{t^N\nu [B]}`$. The key point is that this is independent of $`x_0`$.
We now want to integrate with respect to $`x_0`$. With $`\mu _t`$ as in Definition 1.6 and $`\phi C(X)`$, we have
(3.12) $`{\displaystyle _X}\phi 𝑑\mu _t`$ $`={\displaystyle _X}\phi d(e_t)_{}\mathrm{\Pi }={\displaystyle _X}(\phi e_t)𝑑\mathrm{\Pi }={\displaystyle _\mathrm{\Gamma }}\phi (\gamma (t))𝑑\mathrm{\Pi }(\gamma )`$
$`={\displaystyle _{X\times X}}\phi (\gamma _{x_0,x_1}(t))𝑑\mu (x_0)𝑑\mu (x_1)`$
and
(3.13)
$$_X\phi 𝑑\mu _t^{x_0}=_X\phi (\gamma _{x_0,x_1}(t))𝑑\mu (x_1).$$
These equations show that
(3.14)
$$\mu _t=_X\mu _t^{x_0}𝑑\mu (x_0).$$
In particular, $`\mu _t`$ admits a density $`\rho _t`$, which satisfies the equation
(3.15)
$$\rho _t(x)=_X\rho _t^{x_0}(x)𝑑\mu (x_0).$$
It follows immediately that
(3.16)
$$\rho _t(x)\frac{1}{t^N\nu [B]}.$$
As geodesics are almost-everywhere unique, we can apply the preceding arguments symmetrically with respect to the change $`t1t`$. This gives
(3.17)
$$\rho _t(x)\frac{1}{(1t)^N\nu [B]}.$$
Then
(3.18)
$$\rho _t(x)\mathrm{min}(\frac{1}{t^N},\frac{1}{(1t)^N})\frac{1}{\nu [B]}\frac{2^N}{\nu [B]}.$$
The theorem follows. ∎
###### Remark 3.19.
The above bounds (3.18) can be improved as follows. Let $`\mu =\rho \nu `$ be a measure that is absolutely continuous with respect to $`\nu `$, and otherwise arbitrary. Then there exists a probability measure $`\mathrm{\Pi }P(\mathrm{\Gamma })`$, with $`E_{}\mathrm{\Pi }=\mu \mu `$, such that $`\mu _t=(e_t)_{}\mathrm{\Pi }`$ admits a density $`\rho _t`$ with respect to $`\nu `$, and
(3.20)
$$\rho _t_{L^p}\mathrm{min}(\frac{1}{t^{N/p^{}}},\frac{1}{(1t)^{N/p^{}}})\rho _{L^p}$$
for all $`p(1,\mathrm{})`$, where $`p^{}=p/(p1)`$ is the conjugate exponent to $`p`$ and $`\rho _{L^p}=(_X\rho ^p𝑑\nu )^{1/p}`$. Condition (3.20) is also stable by measured Gromov-Hausdorff limits. Yet we prefer to focus on Condition $`\mathrm{DM}`$ because it is a priori weaker, and still implies the local Poincaré inequality.
## 4. Definition of $`N`$-Ricci curvature bounded below by $`K`$
We recall some more notation.
For $`N[1,\mathrm{})`$, the class $`𝒟𝒞_N`$ is the set of continuous convex functions $`U:[0,\mathrm{})`$, with $`U(0)=0`$, such that the function
(4.1)
$$\psi (\lambda )=\lambda ^NU\left(\lambda ^N\right)$$
is convex on $`(0,\mathrm{})`$. For $`N=\mathrm{}`$, the class $`𝒟𝒞_{\mathrm{}}`$ is the set of continuous convex functions $`U:[0,\mathrm{})`$, with $`U(0)=0`$, such that the function
(4.2)
$$\psi (\lambda )=e^\lambda U\left(e^\lambda \right)$$
is convex on $`(\mathrm{},\mathrm{})`$. In both cases, such a $`\psi `$ is automatically nonincreasing by the convexity of $`U`$. We write $`U^{}(\mathrm{})=lim_r\mathrm{}\frac{U(r)}{r}`$. If a reference probability measure $`\nu P(X)`$ is given, we define a function $`U_\nu :P(X)\{\mathrm{}\}`$ by
(4.3)
$$U_\nu (\mu )=_XU(\rho )𝑑\nu +U^{}(\mathrm{})\mu _s(X),$$
where $`\mu =\rho \nu +\mu _s`$ is the Lebesgue decomposition of $`\mu `$ with respect to $`\nu `$.
We now introduce some expressions that played a prominent role in and in . Given $`K`$ and $`N(1,\mathrm{}]`$, define
(4.4)
$$\beta _t(x_0,x_1)=\{\begin{array}{cc}e^{\frac{1}{6}K(1t^2)d(x_0,x_1)^2}\hfill & \text{if }N=\mathrm{},\hfill \\ \mathrm{}\hfill & \text{if }N<\mathrm{}\text{}K>0\text{ and }\alpha >\pi ,\hfill \\ \left(\frac{\mathrm{sin}(t\alpha )}{t\mathrm{sin}\alpha }\right)^{N1}\hfill & \text{if }N<\mathrm{}\text{}K>0\text{ and }\alpha [0,\pi ],\hfill \\ 1\hfill & \text{if }N<\mathrm{}\text{ and }K=0,\hfill \\ \left(\frac{\mathrm{sinh}(t\alpha )}{t\mathrm{sinh}\alpha }\right)^{N1}\hfill & \text{if }N<\mathrm{}\text{ and }K<0,\hfill \end{array}$$
where
(4.5)
$$\alpha =\sqrt{\frac{|K|}{N1}}d(x_0,x_1).$$
When $`N=1`$, define
(4.6)
$$\beta _t(x_0,x_1)=\{\begin{array}{cc}\mathrm{}\text{if }K>0,\hfill & \\ 1\text{if }K0.\hfill & \end{array}$$
Although we may not write it explicitly, $`\alpha `$ and $`\beta `$ depend on $`K`$ and $`N`$.
###### Definition 4.7.
We say that $`(X,d,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K`$ if the following condition is satisfied. Given $`\mu _0,\mu _1P(X)`$ with support in $`\mathrm{supp}(\nu )`$, write their Lebesgue decompositions with respect to $`\nu `$ as $`\mu _0=\rho _0\nu +\mu _{0,s}`$ and $`\mu _1=\rho _1\nu +\mu _{1,s}`$, respectively. Then there is some optimal dynamical transference plan $`\mathrm{\Pi }`$ from $`\mu _0`$ to $`\mu _1`$, with corresponding Wasserstein geodesic $`\mu _t=(e_t)_{}\mathrm{\Pi }`$, so that for all $`U𝒟𝒞_N`$ and all $`t[0,1]`$, we have
(4.8) $`U_\nu (\mu _t)`$ $`(1t){\displaystyle _{X\times X}}\beta _{1t}(x_0,x_1)U\left({\displaystyle \frac{\rho _0(x_0)}{\beta _{1t}(x_0,x_1)}}\right)𝑑\pi (x_1|x_0)𝑑\nu (x_0)+`$
$`t{\displaystyle _{X\times X}}\beta _t(x_0,x_1)U\left({\displaystyle \frac{\rho _1(x_1)}{\beta _t(x_0,x_1)}}\right)𝑑\pi (x_0|x_1)𝑑\nu (x_1)+`$
$`U^{}(\mathrm{})\left((1t)\mu _{0,s}[X]+t\mu _{1,s}[X]\right).`$
Here if $`\beta _t(x_0,x_1)=\mathrm{}`$ then we interpret $`\beta _t(x_0,x_1)U\left(\frac{\rho _1(x_1)}{\beta _t(x_0,x_1)}\right)`$ as $`U^{}(0)\rho _1(x_1)`$, and similarly $`\beta _{1t}(x_0,x_1)U\left(\frac{\rho _0(x_0)}{\beta _{1t}(x_0,x_1)}\right)`$ as $`U^{}(0)\rho _0(x_0)`$. It is not difficult to show that if $`N<\mathrm{}`$ and $`(X,d,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K>0`$, then the diameter of the support of $`\nu `$ is bounded above by $`\pi \sqrt{(N1)/K}`$. In that case, the quantity $`\alpha `$ defined in (4.5) will vary only in $`[0,\pi ]`$ as $`x_0`$, $`x_1`$ vary in the support of $`\nu `$.
###### Remark 4.9.
If $`\mu _0`$ and $`\mu _1`$ are absolutely continuous with respect to $`\nu `$ then the inequality can be rewritten in the more symmetric form
(4.10) $`U_\nu (\mu _t)`$ $`(1t){\displaystyle _{X\times X}}{\displaystyle \frac{\beta _{1t}(x_0,x_1)}{\rho _0(x_0)}}U\left({\displaystyle \frac{\rho _0(x_0)}{\beta _{1t}(x_0,x_1)}}\right)𝑑\pi (x_0,x_1)+`$
$`t{\displaystyle _{X\times X}}{\displaystyle \frac{\beta _t(x_0,x_1)}{\rho _1(x_1)}}U\left({\displaystyle \frac{\rho _1(x_1)}{\beta _t(x_0,x_1)}}\right)𝑑\pi (x_0,x_1).`$
###### Remark 4.11.
Note that (4.8) is unchanged by the addition of a linear function $`rcr`$ to $`U`$. Of course, the validity of (4.8) depends on the values of $`K`$ and $`N`$. The parameter $`\beta _t`$ is monotonically nondecreasing in $`K`$ and the function $`\beta \beta U(\rho /\beta )`$ is monotonically nonincreasing in $`\beta `$ (because of the convexity of $`U`$). It follows that if $`KK^{}`$ and $`(X,d,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K^{}`$ then it also has $`N`$-Ricci curvature bounded below by $`K`$, as one would expect. One can also show that if $`NN^{}`$ and $`(X,d,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K`$ then it has $`N^{}`$-Ricci curvature bounded below by $`K`$.
We now compare Definition 4.8 with earlier definitions in the literature, starting with the case $`N<\mathrm{}`$. If $`N<\mathrm{}`$ and $`K=0`$ then one recovers the $`N<\mathrm{}`$, $`K=0`$ definition of . If $`N<\mathrm{}`$ and one specializes to $`U(r)`$ being
(4.12)
$$U_N(r)=Nr\left(1r^{1/N}\right),$$
with corresponding entropy function
(4.13)
$$H_{N,\nu }(\mu )=NN_X\rho ^{1\frac{1}{N}}𝑑\nu ,$$
then one recovers the $`N<\mathrm{}`$ definition of Sturm . (In it was not required that $`\pi `$ and $`\{\mu _t\}_{t[0,1]}`$ be related in the sense that they both arise from an optimal dynamical transference plan $`\mathrm{\Pi }`$. We can make that requirement without loss of consistency.)
To deal with the $`N=\mathrm{}`$ case, we use the following lemma.
###### Lemma 4.14.
If $`N=\mathrm{}`$, with $`\psi `$ defined as in (4.2), put
(4.15)
$$\lambda (U)=\{\begin{array}{cc}K\psi ^{}(\mathrm{})\hfill & \text{ if }K>0,\hfill \\ 0\hfill & \text{ if }K=0,\hfill \\ K\psi ^{}(\mathrm{})\hfill & \text{ if }K>0.\hfill \end{array}$$
If $`\mu _0`$ and $`\mu _1`$ are absolutely continuous with respect to $`\nu `$ then
(4.16)
$$_{X\times X}\frac{\beta _t(x_0,x_1)}{\rho _1(x_1)}U\left(\frac{\rho _1(x_1)}{\beta _t(x_0,x_1)}\right)𝑑\pi (x_0,x_1)_XU(\rho _1)𝑑\nu \frac{1}{6}(1t^2)\lambda (U)W_2(\mu _0,\mu _1)^2.$$
###### Proof.
Suppose first that $`K>0`$. From the convexity of $`\psi `$, if $`\rho _1(x_1)>0`$ then
(4.17)
$$\frac{\psi (\mathrm{ln}\rho _1(x_1)+\frac{1}{6}K(1t^2)d(x_0,x_1)^2)\psi (\mathrm{ln}\rho _1(x_1))}{\frac{1}{6}K(1t^2)d(x_0,x_1)^2}\psi ^{}(\mathrm{}).$$
Then
(4.18)
$$\frac{\beta _t(x_0,x_1)}{\rho _1(x_1)}U\left(\frac{\rho _1(x_1)}{\beta _t(x_0,x_1)}\right)\frac{1}{\rho _1(x_1)}U\left(\rho _1(x_1)\right)+\frac{1}{6}K\psi ^{}(\mathrm{})(1t^2)d(x_0,x_1)^2.$$
The lemma follows upon integration with respect to $`d\pi (x_0,x_1)`$. The cases $`K=0`$ and $`K<0`$ are similar. ∎
Using Lemma 4.14, and the analogous inequality for $`_{X\times X}\frac{\beta _{1t}(x_0,x_1)}{\rho _0(x_0)}U\left(\frac{\rho _0(x_0)}{\beta _{1t}(x_0,x_1)}\right)𝑑\pi (x_0,x_1)`$, one finds that (4.10) implies
(4.19)
$$U_\nu (\mu _t)tU_\nu (\mu _1)+(1t)U_\nu (\mu _0)\frac{1}{2}\lambda (U)t(1t)W_2(\mu _0,\mu _1)^2,$$
which is exactly the inequality used in to define what it means for $`(X,d,\nu )`$ to have $`\mathrm{}`$-Ricci curvature bounded below by $`K`$, in the sense of . We have only shown that (4.19) holds when $`\mu _0`$ and $`\mu _1`$ are absolutely continuous with respect to $`\nu `$, but \[10, Proposition 3.21\] then implies that it holds for all $`\mu _0,\mu _1`$ with support in $`\mathrm{supp}(\nu )`$.
A consequence is that any result of concerning measured length spaces with $`\mathrm{}`$-Ricci curvature bounded below by $`K`$, in the sense of , also holds for measured length spaces with $`\mathrm{}`$-Ricci curvature bounded below by $`K`$ in the sense of Definition 4.7.
Finally, Sturm’s notion of having $`\mathrm{}`$-Ricci curvature bounded below by $`K`$ is the specialization of the definition of to the case $`U(r)=U_{\mathrm{}}(r)=r\mathrm{ln}(r)`$.
The notion of having $`N`$-Ricci curvature bounded below by $`K`$, in the sense of Definition 4.7, is preserved by measured Gromov-Hausdorff limits. We will present the proof, which is more complicated than that of the analogous statement in , elsewhere.
We now show that in the case of a Riemannian manifold equipped with a smooth measure, a lower Ricci curvature bound in the sense of Definition 4.7 is equivalent to a tensor inequality.
Let $`M`$ be a smooth compact connected $`n`$-dimensional manifold with Riemannian metric $`g`$. We let $`(M,g)`$ denote the corresponding metric space. Given $`\mathrm{\Psi }C^{\mathrm{}}(M)`$ with $`_Me^\mathrm{\Psi }\mathrm{dvol}_M=\mathrm{\hspace{0.25em}1}`$, put $`d\nu =e^\mathrm{\Psi }\mathrm{dvol}_M`$.
###### Definition 4.20.
For $`N[1,\mathrm{}]`$, let the $`N`$-Ricci tensor $`\mathrm{Ric}_N`$ of $`(M,g,\nu )`$ be defined by
(4.21)
$$\mathrm{Ric}_N=\{\begin{array}{cc}\mathrm{Ric}+Hess(\mathrm{\Psi })\hfill & \text{ if }N=\mathrm{},\hfill \\ \mathrm{Ric}+Hess(\mathrm{\Psi })\frac{1}{Nn}d\mathrm{\Psi }d\mathrm{\Psi }\hfill & \text{ if }n<N<\mathrm{},\hfill \\ \mathrm{Ric}+Hess(\mathrm{\Psi })\mathrm{}(d\mathrm{\Psi }d\mathrm{\Psi })\hfill & \text{ if }N=n,\hfill \\ \mathrm{}\hfill & \text{ if }N<n\text{,}\hfill \end{array}$$
where by convention $`\mathrm{}0=\mathrm{\hspace{0.25em}0}`$.
###### Theorem 4.22.
For $`N[1,\mathrm{}]`$, the measured length space $`(M,g,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K`$ if and only if $`\mathrm{Ric}_NKg`$.
###### Proof.
The proof is similar to that of \[10, Theorems 7.3 and 7.42\] and \[14, Theorem 1.9\]. We only sketch a few points of the proof, in order to clarify the role played by the function $`U`$.
Suppose that $`N(1,\mathrm{})`$ and $`\mathrm{Ric}_NKg`$. We want to show that the condition in Definition 4.7 holds. As in \[10, Theorem 7.3\], we can reduce to the case when $`\mu _0`$ and $`\mu _1`$ are absolutely continuous with respect to $`\nu `$. The unique Wasserstein geodesic between them is of the form $`\mu _t=(F_t)_{}\mu _0`$ for certain maps $`F_t:MM`$. Put
(4.23)
$$C(y,t)=e^{\frac{\mathrm{\Psi }(F_t(y))}{N}}\stackrel{\frac{1}{N}}{det}(dF_t)(y)$$
and
(4.24)
$$\eta _0=\frac{d\mu _0}{\mathrm{dvol}_M}.$$
Then in terms of the function $`\psi `$ of (4.1) there is an equation \[10, (7.19)\]
(4.25)
$$U_\nu (\mu _t)=_M\psi \left(C(y,t)\eta _0^{\frac{1}{N}}(y)\right)𝑑\mu _0(y).$$
With the notation (4.5) in use, define
(4.26)
$$\tau _{K,N}^{(t)}(d(x_0,x_1))=\{\begin{array}{cc}t^{\frac{1}{N}}\left(\frac{\mathrm{sin}(t\alpha )}{\mathrm{sin}\alpha }\right)^{1\frac{1}{N}}\hfill & \text{ if }K>0,\hfill \\ t\hfill & \text{ if }K=0,\hfill \\ t^{\frac{1}{N}}\left(\frac{\mathrm{sinh}(t\alpha )}{\mathrm{sinh}\alpha }\right)^{1\frac{1}{N}}\hfill & \text{ if }K<0,\hfill \end{array}$$
one can show by combining \[10, Section 7\] and \[14, Section 5\] that
(4.27)
$$C(y,t)\tau _{K,N}^{(1t)}(d(y,F_1(y)))C(y,0)+\tau _{K,N}^{(t)}(d(y,F_1(y)))C(y,1).$$
As $`\psi `$ is nonincreasing, we obtain
(4.28)
$$U_\nu (\mu _t)_M\psi \left(\frac{\tau _{K,N}^{(1t)}(d(y,F_1(y)))C(y,0)+\tau _{K,N}^{(t)}(d(y,F_1(y)))C(y,1)}{\eta _0^{\frac{1}{N}}(y)}\right)𝑑\mu _0(y).$$
As $`\psi `$ is convex by assumption, we now obtain
(4.29) $`U_\nu (\mu _t)`$ $`(1t){\displaystyle _M}\psi \left({\displaystyle \frac{\tau _{K,N}^{(1t)}(d(y,F_1(y)))}{1t}}{\displaystyle \frac{C(y,0)}{\eta _0^{\frac{1}{N}}(y)}}\right)𝑑\mu _0(y)+`$
$`t{\displaystyle _M}\psi \left({\displaystyle \frac{\tau _{K,N}^{(t)}(d(y,F_1(y)))}{t}}{\displaystyle \frac{C(y,1)}{\eta _0^{\frac{1}{N}}(y)}}\right)𝑑\mu _0(y).`$
After using the definition of $`\psi `$ again, along with (4.25) in the cases $`t=0`$ and $`t=1`$, one arrives at (4.10). ∎
The next result is an analog of \[10, Theorem 5.52\].
###### Theorem 4.30.
If $`(X,d,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K`$ then for any $`\mu _0,\mu _1P(X)`$ that are absolutely continuous with respect to $`\nu `$, the Wasserstein geodesic $`\{\mu _t\}_{t[0,1]}`$ of Definition 4.7 has the property that $`\mu _t`$ is absolutely continuous with respect to $`\nu `$, for all $`t[0,1]`$.
###### Proof.
The proof is along the lines of \[10, Theorem 5.52\]. ∎
Theorem 4.30 will be needed in equation (5.7) below. This is one reason why we require (4.8) to hold for all $`U𝒟𝒞_N`$, as opposed to just $`U_N`$. (We note that the distinction between these two definitions disappears in nonbranching spaces, as will be shown elsewhere.)
## 5. Sobolev inequality and global Poincaré inequality
###### Definition 5.1.
Given $`f\mathrm{Lip}(X)`$, put
(5.2)
$$|^{}f|(x)=\underset{yx}{lim\; sup}\frac{[f(y)f(x)]_{}}{d(x,y)}=\underset{yx}{lim\; sup}\frac{[f(x)f(y)]_+}{d(x,y)}.$$
Here $`a_+=\mathrm{max}(a,0)`$ and $`a_{}=\mathrm{max}(a,0)`$. Note that $`|^{}f|(x)|f|(x)`$, the latter being defined as in (2.2).
###### Theorem 5.3.
Given $`N(1,\mathrm{})`$ and $`K>0`$, suppose that $`(X,d,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K`$. Then for any positive Lipschitz function $`\rho _0\mathrm{Lip}(X)`$ with $`_X\rho _0𝑑\nu =\mathrm{\hspace{0.25em}1}`$, one has
(5.4)
$$NN_X\rho _0^{1\frac{1}{N}}𝑑\nu _X\theta ^{(N,K)}(\rho _0,|^{}\rho _0|)𝑑\nu ,$$
where
(5.5) $`\theta ^{(N,K)}(r,g)=r\underset{\alpha [0,\pi ]}{sup}[{\displaystyle \frac{N1}{N}}{\displaystyle \frac{g}{r^{1+\frac{1}{N}}}}`$ $`\sqrt{{\displaystyle \frac{N1}{K}}}\alpha +N\left(1\left({\displaystyle \frac{\alpha }{\mathrm{sin}(\alpha )}}\right)^{1\frac{1}{N}}\right)+`$
$`(N1)({\displaystyle \frac{\alpha }{\mathrm{tan}(\alpha )}}1)r^{\frac{1}{N}}].`$
###### Proof.
We recall the definitions of $`U_N`$ and $`H_{N,\nu }`$ from (4.12) and (4.13). Applying Definition 4.7 with $`U=U_N`$, any two probability measures $`\mu _0=\rho _0\nu `$ and $`\mu _1=\rho _1\nu `$ can be joined by a Wasserstein geodesic $`\{\mu _t\}_{t[0,1]}`$, arising from an optimal dynamical transference plan, along which the following inequality holds :
(5.6) $`H_{N,\nu }(\mu _t)`$
$`NN{\displaystyle _{X\times X}}\left[\tau _{K,N}^{(1t)}(d(x_0,x_1))\rho _0^{\frac{1}{N}}(x_0)+\tau _{K,N}^{(t)}(d(x_0,x_1))\rho _1^{\frac{1}{N}}(x_1)\right]𝑑\pi (x_0,x_1).`$
By Theorem 4.30, $`\mu _t`$ is absolutely continuous with respect to $`\nu `$.
Given a positive function $`\rho _0\mathrm{Lip}(X)`$, put $`\mu _0=\rho _0\nu `$ and $`\mu _1=\nu `$. Put $`\varphi (t)=H_{N,\nu }(\mu _t)`$. In the proof of \[10, Proposition 3.36\] it was shown that
(5.7) $`\underset{t0}{lim\; sup}{\displaystyle \frac{\varphi (t)\varphi (0)}{t}}`$ $`{\displaystyle _\mathrm{\Gamma }}U_N^{\prime \prime }(\rho (\gamma (0)))|^{}\rho _0|(\gamma (0))d(\gamma (0),\gamma (1))𝑑\mathrm{\Pi }(\gamma )`$
$`={\displaystyle \frac{N1}{N}}{\displaystyle _X}{\displaystyle \frac{|^{}\rho _0|(x_0)}{\rho _0(x_0)^{1+\frac{1}{N}}}}d(x_0,x_1)𝑑\pi (x_0,x_1).`$
On the other hand, from (5.6),
(5.8)
$$\varphi (t)NN_{X\times X}\left[\tau _{K,N}^{(1t)}(d(x_0,x_1))\rho _0^{\frac{1}{N}}(x_0)+\tau _{K,N}^{(t)}(d(x_0,x_1))\right]𝑑\pi (x_0,x_1)$$
and so
(5.9)
$$\frac{\varphi (t)\varphi (0)}{t}N_{X\times X}\left[\frac{\tau _{K,N}^{(1t)}(d(x_0,x_1))1}{t}\rho _0^{\frac{1}{N}}(x_0)+\frac{\tau _{K,N}^{(t)}(d(x_0,x_1))}{t}\right]𝑑\pi (x_0,x_1).$$
Then
(5.10) $`\underset{t0}{lim\; sup}{\displaystyle \frac{\varphi (t)\varphi (0)}{t}}N{\displaystyle _{X\times X}}\left({\displaystyle \frac{\sqrt{\frac{K}{N1}}d(x_0,x_1)}{\mathrm{sin}\left(\sqrt{\frac{K}{N1}}d(x_0,x_1)\right)}}\right)^{1\frac{1}{N}}𝑑\pi (x_0,x_1)+`$
$`{\displaystyle _{X\times X}}\left[1+(N1)\sqrt{{\displaystyle \frac{K}{N1}}}d(x_0,x_1)\mathrm{cot}\left(\sqrt{{\displaystyle \frac{K}{N1}}}d(x_0,x_1)\right)\right]\rho _0^{\frac{1}{N}}(x_0)𝑑\pi (x_0,x_1).`$
Combining (5.7) and (5.10), and slightly rewriting the result, gives
(5.11) $`{\displaystyle \frac{N1}{N}}{\displaystyle _X}{\displaystyle \frac{|^{}\rho _0|(x_0)}{\rho _0(x_0)^{1+\frac{1}{N}}}}d(x_0,x_1)𝑑\pi (x_0,x_1)`$
$`N{\displaystyle _{X\times X}}\left[1\left({\displaystyle \frac{\sqrt{\frac{K}{N1}}d(x_0,x_1)}{\mathrm{sin}\left(\sqrt{\frac{K}{N1}}d(x_0,x_1)\right)}}\right)^{1\frac{1}{N}}\right]𝑑\pi (x_0,x_1)+`$
$`(N1){\displaystyle _{X\times X}}\left[\sqrt{{\displaystyle \frac{K}{N1}}}d(x_0,x_1)\mathrm{cot}\left(\sqrt{{\displaystyle \frac{K}{N1}}}d(x_0,x_1)\right)\mathrm{\hspace{0.25em}1}\right]\rho _0^{\frac{1}{N}}(x_0)𝑑\pi (x_0,x_1)`$
$`H_{N,\nu }(\mu ),`$
or
(5.12) $`H_{N,\nu }(\mu ){\displaystyle \frac{N1}{N}}{\displaystyle _X}{\displaystyle \frac{|^{}\rho _0|(x_0)}{\rho _0(x_0)^{1+\frac{1}{N}}}}d(x_0,x_1)𝑑\pi (x_0,x_1)+`$
$`N{\displaystyle _{X\times X}}\left[1\left({\displaystyle \frac{\sqrt{\frac{K}{N1}}d(x_0,x_1)}{\mathrm{sin}\left(\sqrt{\frac{K}{N1}}d(x_0,x_1)\right)}}\right)^{1\frac{1}{N}}\right]𝑑\pi (x_0,x_1)+`$
$`(N1){\displaystyle _{X\times X}}\left[\sqrt{{\displaystyle \frac{K}{N1}}}d(x_0,x_1)\mathrm{cot}\left(\sqrt{{\displaystyle \frac{K}{N1}}}d(x_0,x_1)\right)\mathrm{\hspace{0.25em}1}\right]\rho _0^{\frac{1}{N}}(x_0)𝑑\pi (x_0,x_1).`$
Replacing $`\sqrt{\frac{K}{N1}}d(x_0,x_1)`$ by $`\alpha `$, we get only a weaker inequality by taking the sup over $`\alpha [0,\pi ]`$. The theorem follows. ∎
In order to clarify the nature of the inequality of Theorem 5.3, we derive a slightly weaker inequality. First, we prove an elementary estimate.
###### Lemma 5.13.
For $`x[0,\pi ]`$, one has
(5.14)
$$\frac{x}{\mathrm{tan}(x)}\mathrm{\hspace{0.25em}1}\frac{x^2}{3}$$
and
(5.15)
$$1\left(\frac{x}{\mathrm{sin}(x)}\right)^{1\frac{1}{N}}\left(1\frac{1}{N}\right)\frac{x^2}{6}.$$
###### Proof.
Put
(5.16)
$$F(x)=x\frac{\mathrm{sin}(x)\mathrm{cos}(x)}{1\frac{2}{3}\mathrm{sin}^2(x)}.$$
Then
(5.17)
$$F^{}(x)=\frac{4}{9}\frac{\mathrm{sin}^4(x)}{\left(1\frac{2}{3}\mathrm{sin}^2(x)\right)^2}\mathrm{\hspace{0.25em}0}.$$
As $`F(0)=0`$, it follows that $`F(x)0`$ for $`x[0,\pi ]`$, so
(5.18)
$$x\left(1\frac{2}{3}\mathrm{sin}^2(x)\right)\mathrm{sin}(x)\mathrm{cos}(x).$$
Putting
(5.19)
$$G(x)=\frac{x}{\mathrm{tan}(x)}+\frac{1}{3}x^2$$
and using (5.18), one obtains
(5.20)
$$G^{}(x)=\frac{x}{\mathrm{sin}^2(x)}+\frac{1}{\mathrm{tan}(x)}+\frac{2}{3}x\mathrm{\hspace{0.25em}0}.$$
As $`G(0)=1`$, we have
(5.21)
$$\frac{x}{\mathrm{tan}(x)}+\frac{1}{3}x^2\mathrm{\hspace{0.25em}1},$$
which proves (5.14).
Next, from (5.14) we have
(5.22)
$$\frac{1}{x}\frac{1}{\mathrm{tan}(x)}\frac{x}{3}.$$
Integrating gives
(5.23)
$$\mathrm{ln}\left(\frac{x}{\mathrm{sin}(x)}\right)\frac{x^2}{6},$$
so
(5.24)
$$\frac{x}{\mathrm{sin}(x)}e^{\frac{x^2}{6}}\left(1+\left(1\frac{1}{N}\right)\frac{x^2}{6}\right)^{\frac{1}{1\frac{1}{N}}}.$$
Thus
(5.25)
$$\left(\frac{x}{\mathrm{sin}(x)}\right)^{1\frac{1}{N}}\mathrm{\hspace{0.25em}1}+\left(1\frac{1}{N}\right)\frac{x^2}{6}.$$
This proves (5.15). ∎
We now prove a Sobolev-type inequality.
###### Theorem 5.26.
Given $`N(1,\mathrm{})`$ and $`K>0`$, suppose that $`(X,d,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K`$. Then for any nonnegative Lipschitz function $`\rho _0\mathrm{Lip}(X)`$ with $`_X\rho _0𝑑\nu =\mathrm{\hspace{0.25em}1}`$, one has
(5.27)
$$NN_X\rho _0^{1\frac{1}{N}}𝑑\nu \frac{1}{2K}\left(\frac{N1}{N}\right)^2_X\frac{\rho _0^{1\frac{2}{N}}}{\frac{1}{3}+\frac{2}{3}\rho _0^{\frac{1}{N}}}|^{}\rho _0|^2𝑑\nu .$$
###### Proof.
If $`\rho _0`$ is positive then using Lemma 5.13, we can estimate the function $`\theta ^{(N,K)}(r,g)`$ of (5.5) by
(5.28) $`\theta ^{(N,K)}(r,g)`$ $`r\underset{\alpha [0,\pi ]}{sup}\left[{\displaystyle \frac{N1}{N}}{\displaystyle \frac{g}{r^{1+\frac{1}{N}}}}\sqrt{{\displaystyle \frac{N1}{K}}}\alpha {\displaystyle \frac{N1}{6}}\alpha ^2\left(1+\mathrm{\hspace{0.25em}2}r^{\frac{1}{N}}\right)\right]`$
$`{\displaystyle \frac{1}{2K}}\left({\displaystyle \frac{N1}{N}}\right)^2{\displaystyle \frac{r^{1\frac{2}{N}}}{\frac{1}{3}+\frac{2}{3}r^{\frac{1}{N}}}}g^2.`$
The theorem in this case follows from Theorem 5.3. The case when $`\rho _0`$ is nonnegative can be handled by approximation with positive functions. ∎
To put Theorem 5.26 into a more conventional form, we prove a slightly weaker inequality.
###### Theorem 5.29.
Given $`N(2,\mathrm{})`$ and $`K>0`$, suppose that $`(X,d,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K`$. Then for any nonnegative Lipschitz function $`f\mathrm{Lip}(X)`$ with $`_Xf^{\frac{2N}{N2}}𝑑\nu =\mathrm{\hspace{0.25em}1}`$, one has
(5.30)
$$1\left(_Xf𝑑\nu \right)^{\frac{2}{N+2}}\frac{6}{KN}\left(\frac{N}{N2}\right)^2_X|^{}f|^2𝑑\nu .$$
###### Proof.
Put $`\rho _0=f^{\frac{2N}{N2}}`$. From (5.27) we have
(5.31)
$$NN_X\rho _0^{1\frac{1}{N}}𝑑\nu \frac{3}{2K}\left(\frac{N1}{N}\right)^2_X\rho _0^{1\frac{2}{N}}|^{}\rho _0|^2𝑑\nu ,$$
which gives
(5.32)
$$1_Xf^{2\left(\frac{N1}{N2}\right)}𝑑\nu \frac{6}{KN}\left(\frac{N1}{N2}\right)^2_X|^{}f|^2𝑑\nu .$$
By Hölder’s inequality,
(5.33)
$$_Xf^{2\left(\frac{N1}{N2}\right)}𝑑\nu \left(_Xf^{\frac{2N}{N2}}𝑑\nu \right)^{\frac{N}{N+2}}\left(_Xf𝑑\nu \right)^{\frac{2}{N+2}}.$$
The theorem follows. ∎
Putting (5.30) into a homogeneous form reveals the content of Theorem 5.29: there is a bound of the form $`f_{\frac{2N}{N2}}F(f_1,^{}f_2)`$ for some appropriate function $`F`$. This is an example of Sobolev embedding. Of course equation (5.30) is not sharp, due to the many approximations made.
Finally, we prove a sharp global Poincaré inequality.
###### Theorem 5.34.
Given $`N(1,\mathrm{})`$ and $`K>0`$, suppose that $`(X,d,\nu )`$ has $`N`$-Ricci curvature bounded below by $`K`$. Suppose that $`f\mathrm{Lip}(X)`$ has $`_Xf𝑑\nu =\mathrm{\hspace{0.25em}0}`$. Then
(5.35)
$$_Xf^2𝑑\nu \frac{N1}{KN}_X|^{}f|^2𝑑\nu .$$
###### Proof.
Without loss of generality we may assume that $`\mathrm{max}|f|\mathrm{\hspace{0.25em}1}`$. Given $`ϵ(1,1)`$, put $`\rho _0=\mathrm{\hspace{0.25em}1}+ϵf`$. Then $`\rho _0>0`$ and $`_X\rho _0𝑑\nu =\mathrm{\hspace{0.25em}1}`$. For small $`ϵ`$,
(5.36)
$$NN_X\rho _0^{1\frac{1}{N}}𝑑\nu =ϵ^2\frac{N1}{2N}_Xf^2𝑑\nu +O(ϵ^3)$$
and
(5.37)
$$\frac{1}{2K}\left(\frac{N1}{N}\right)^2_X\frac{\rho _0^{1\frac{2}{N}}}{\frac{1}{3}+\frac{2}{3}\rho _0^{\frac{1}{N}}}|^{}\rho _0|^2𝑑\nu =\frac{ϵ^2}{2K}\left(\frac{N1}{N}\right)^2_X|^{}f|^2𝑑\nu +O(ϵ^3).$$
Then the result follows from Theorem 5.26. ∎
###### Remark 5.38.
1. In the case of an $`N`$-dimensional Riemannian manifold with $`\mathrm{Ric}Kg`$, one recovers the Lichnerowicz inequality for the lowest positive eigenvalue of the Laplacian . It is sharp on round spheres.
2. The case $`N=\mathrm{}`$ was treated by similar means in \[10, Theorem 6.18\].
## Appendix: Stability of dynamical transference plans
In this appendix we prove a general compactness theorem for probability measures on geodesic paths. This theorem is used to show that the condition $`\mathrm{DM}(\mathrm{C})`$ is preserved under measured Gromov-Hausdorff limits.
###### Lemma A.39.
Let $`X`$ be a compact length space. Given $`ϵ>0`$, there is a $`\delta >0`$ with the following property. Suppose that $`Y`$ is a compact length space and $`f:YX`$ is a $`\delta `$-approximation. Let $`\gamma :[0,1]Y`$ be a geodesic. Then there is a geodesic $`T(\gamma ):[0,1]X`$ so that for all $`t[0,1]`$, $`d_X(T(\gamma )(t),f(\gamma (t)))ϵ`$.
###### Proof.
Suppose that the lemma is not true. Then there is some $`ϵ>0`$ along with
1. A sequence of compact metric spaces $`\{Y_i\}_{i=1}^{\mathrm{}}`$,
2. $`\frac{1}{i}`$-approximations $`f_i:Y_iX`$ and
3. Geodesics $`\gamma _i:[0,1]Y_i`$
so that for each geodesic $`\gamma ^{}:[0,1]X`$, there is some $`t_{i,\gamma ^{}}[0,1]`$ with
(A.40)
$$d_X(\gamma ^{}(t_{i,\gamma ^{}}),f_i(\gamma _i(t_{i,\gamma ^{}})))>ϵ.$$
After passing to a subsequence, we may assume that $`\{f_i\gamma _i\}_{i=1}^{\mathrm{}}`$ converges uniformly to a geodesic $`\gamma _{\mathrm{}}:[0,1]X`$. After passing to a further subsequence, we may assume that $`lim_i\mathrm{}t_{i,\gamma _{\mathrm{}}}=t_{\mathrm{}}`$ for some $`t_{\mathrm{}}[0,1]`$. Then
(A.41) $`d_X(\gamma _{\mathrm{}}(t_{i,\gamma _{\mathrm{}}}),f_i(\gamma _i(t_{i,\gamma _{\mathrm{}}})))`$ $`d_X(\gamma _{\mathrm{}}(t_{i,\gamma _{\mathrm{}}}),\gamma _{\mathrm{}}(t_{\mathrm{}}))+d_X(\gamma _{\mathrm{}}(t_{\mathrm{}}),f_i(\gamma _i(t_{\mathrm{}}))+`$
$`d_X(f_i(\gamma _i(t_{\mathrm{}}),f_i(\gamma _i(t_{i,\gamma _{\mathrm{}}})))`$
$`\mathrm{diam}(X)|t_{i,\gamma _{\mathrm{}}}t_{\mathrm{}}|+d_X(\gamma _{\mathrm{}}(t_{\mathrm{}}),f_i(\gamma _i(t_{\mathrm{}}))+`$
$`{\displaystyle \frac{1}{i}}+\mathrm{diam}(Y_i)|t_{i,\gamma _{\mathrm{}}}t_{\mathrm{}}|.`$
Then the right-hand side converges to 0 as $`i\mathrm{}`$, which contradicts (A.40) with $`\gamma ^{}=\gamma _{\mathrm{}}`$. ∎
###### Lemma A.42.
One can choose the map $`T`$ in Lemma A.39 to be a measurable map from $`\mathrm{\Gamma }(Y)`$ to $`\mathrm{\Gamma }(X)`$.
###### Proof.
This follows from \[16, Theorem A.5\], as
(A.43)
$$\{(\gamma _1,\gamma _2)\mathrm{\Gamma }(X)\times \mathrm{\Gamma }(Y):\text{ for all }t[0,1],d_X(\gamma _1(t),f(\gamma _2(t)))ϵ\}$$
is a Borel subset of $`\mathrm{\Gamma }(X)\times \mathrm{\Gamma }(Y)`$ and for each $`\gamma _2\mathrm{\Gamma }(Y)`$,
(A.44)
$$\{\gamma _1\mathrm{\Gamma }(X):\text{ for all }t[0,1],d_X(\gamma _1(t),f(\gamma _2(t)))ϵ\}$$
is compact. ∎
###### Theorem A.45.
Let $`\{(X_i,d_i)\}_{i=1}^{\mathrm{}}`$ be a sequence of compact length spaces that converges in the Gromov-Hausdorff topology to a compact length space $`(X,d)`$. Let $`f_i:X_iX`$ be $`\epsilon _i`$-approximations, with $`\epsilon _i0`$, that realize the Gromov-Hausdorff convergence. For each $`i`$, let $`\mathrm{\Pi }_i`$ be a Borel probability measure on $`\mathrm{\Gamma }(X_i)`$. Let $`\pi _i`$ and $`\{\mu _{i,t}\}_{t[0,1]}`$ be the associated transference plan and measure-valued path. Then after passing to a subsequence, there is a dynamical transference plan $`\mathrm{\Pi }`$ on $`X`$, with associated transference plan $`\pi `$, and measure-valued path $`\{\mu _t\}_{t[0.1]}`$, such that
(i) $`lim_i\mathrm{}(f_i,f_i)_{}\pi _i=\pi `$ in the weak-$``$ topology on $`P(X\times X)`$;
(ii) $`lim_i\mathrm{}(f_i)_{}\mu _{i,t}=\mu _t`$ for all $`t[0,1]`$.
###### Proof.
Let $`T_i:\mathrm{\Gamma }(X_i)\mathrm{\Gamma }(X)`$ be the map constructed by means of Lemma A.39 with $`\epsilon =\epsilon _i`$, $`Y=X_i`$, $`f=f_i`$. After passing to a convergent subsequence, we can assume that $`lim_i\mathrm{}(T_i)_{}\mathrm{\Pi }_i=\mathrm{\Pi }`$ in the weak-$``$ topology, for some $`\mathrm{\Pi }P(\mathrm{\Gamma }(X))`$. Given $`FC(X\times X)`$, we have
(A.46)
$$_{X\times X}F𝑑\pi =_{\mathrm{\Gamma }(X)}F(\gamma (0),\gamma (1))𝑑\mathrm{\Pi }(\gamma )=\underset{i\mathrm{}}{lim}_{\mathrm{\Gamma }(X_i)}F(T_i(\gamma _i)(0),T_i(\gamma _i)(1))𝑑\mathrm{\Pi }_i(\gamma _i).$$
By Lemma A.39 and the uniform continuity of $`F`$,
(A.47) $`\underset{i\mathrm{}}{lim}{\displaystyle _{\mathrm{\Gamma }(X_i)}}F(T_i(\gamma _i)(0),T_i(\gamma _i)(1))𝑑\mathrm{\Pi }_i(\gamma _i)`$ $`=\underset{i\mathrm{}}{lim}{\displaystyle _{\mathrm{\Gamma }(X_i)}}F(f_i(\gamma _i(0)),f_i(\gamma _i(1)))𝑑\mathrm{\Pi }_i(\gamma _i)`$
$`=\underset{i\mathrm{}}{lim}{\displaystyle _{X_i\times X_i}}F(f_i(x_i),f_i(x_i^{}))𝑑\pi _i(x_i,x_i^{})`$
$`=\underset{i\mathrm{}}{lim}{\displaystyle _{X\times X}}Fd(f_i,f_i)_{}\pi _i.`$
This proves (i). Similarly, for $`t[0,1]`$ and $`FC(X)`$,
(A.48)
$$_XF𝑑\mu _t=_{\mathrm{\Gamma }(X)}F(\gamma (t))𝑑\mathrm{\Pi }(\gamma )=\underset{i\mathrm{}}{lim}_{\mathrm{\Gamma }(X_i)}F(T_i(\gamma _i)(t))𝑑\mathrm{\Pi }_i(\gamma _i).$$
By Lemma A.39 and the uniform continuity of $`F`$,
(A.49) $`\underset{i\mathrm{}}{lim}{\displaystyle _{\mathrm{\Gamma }(X_i)}}F(T_i(\gamma _i)(t))𝑑\mathrm{\Pi }_i(\gamma _i)`$ $`=\underset{i\mathrm{}}{lim}{\displaystyle _{\mathrm{\Gamma }(X_i)}}F(f_i(\gamma _i(t)))𝑑\mathrm{\Pi }_i(\gamma _i)`$
$`=\underset{i\mathrm{}}{lim}{\displaystyle _{X_i}}F(f_i(x_i))𝑑\mu _{i,t}(x_i)`$
$`=\underset{i\mathrm{}}{lim}{\displaystyle _X}Fd(f_i)_{}\mu _{i,t}.`$
This proves (ii). ∎
For reference we give a slight variation of Lemma A.39, although it is not needed in the body of the paper.
###### Lemma A.50.
Let $`X`$ be a compact length space. Choose points $`x,x^{}X`$. Given $`ϵ>0`$, there is a $`\delta =\delta (x,x^{})>0`$ with the following property. Suppose that $`Y`$ is a compact length space and $`f:YX`$ is a $`\delta `$-approximation. Given $`yf^1(x)`$ and $`y^{}f^1(x^{})`$, let $`\gamma :[0,1]Y`$ be a geodesic joining them. Then there is a geodesic $`T(\gamma ):[0,1]X`$ from $`x`$ to $`x^{}`$ so that for all $`t[0,1]`$, $`d_X(T(\gamma )(t),f(\gamma (t)))ϵ`$.
###### Proof.
The proof is along the same lines as that of Lemma A.39. ∎
###### Remark A.51.
In general, one cannot take $`\delta `$ to be independent of $`x`$ and $`x^{}`$.
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# Flow equations in generalized braneworld scenarios
## I Introduction
The recent improvement in the determination of cosmological observables by the Wilkinson Microwave Anisotropy Probe wmap ; pei03 and other large-scale structure experiments has given a boost to the search for viable theoretical scenarios of the early Universe. An important problem cosmologists still have to address is cosmic confusion in inflationary scenarios, whereby different underlying physics leads to the same observables; deriving robust conclusions from data requires an understanding of such model degeneracies. Viable inflation models must predict the quasi-invariant density perturbation spectra and subdominant tensor perturbations that observations require, but having done so it is then difficult to discriminate between different details such as the precise shape of the potential. Indeed, even quite radical revisions of understanding of high-energy physics, such as the braneworld scenario, have yet to lead to appreciably characteristic predictions, at least in the simplest cases (see Ref. rub01 for some reviews).
Rather than choosing a particular class of potentials and calculating the preferred values in the parameter space, one can try to circumscribe the allowed regions within it through the general behaviour of the inflationary dynamics. This can be achieved, for instance, by the study of the consistency equations, which do not depend on the choice of the potential; see Ref. PhD and references therein. Another possibility is to consider the evolution of the so-called flow parameters in terms of the number of $`e`$-foldings HoT ; kin02 ; lid03 . The flow equations naturally select a subset in the observational plane defined by the scalar spectral index $`n_\mathrm{s}`$ and the tensor-to-scalar ratio $`r`$. Modifications to the underlying dynamics will lead to a shift in the location of that subset, helping to indicate whether a particular class of inflationary models is better able to produce observables in the region required by data.
In this paper we shall explore the flow approach in several braneworld scenarios, where the effective Friedmann equation governing the cosmological brane dynamics is modified with respect to the 4D evolution. The simplest such case, the Randall–Sundrum (RS) type II model for a normal scalar field, was already considered in Ref. RL2 and found not to introduce qualitatively different features. Here we shall extend the discussion to the Gauss–Bonnet (GB) braneworld by means of the “patch formalism” PhD , which provides a unified description of a wide range of possible gravity models. This has the advantage that we can simultaneously implement general inflationary scenarios with modified Friedmann equations and either an ordinary or Dirac–Born–Infeld (DBI) tachyonic scalar field. The goal is to seek new insights into the general structure of braneworld and tachyon scenarios via a technique that might give complementary information as compared to model-building approaches.
We will find that in the Gauss–Bonnet case with normal scalar field, the flow structure in the $`n_\mathrm{s}`$$`r`$ plane is quite different from the 4D and RS cases because of the details of the flow equations. The flow structure for the tachyon models, here presented for the first time, gives independent support to other results found in literature; in particular, the most characteristic prediction is the generation of spectra with a rather low tensor-to-scalar ratio. The stability of the fixed points in the plane is also analyzed at higher order in the flow parameters.
Finally, we shall consider how the general relativistic and braneworld scenarios are affected by the introduction of a noncommutative space-time algebra, which modifies the inflationary anisotropies while leaving the homogeneous background (and the associated flow equations) untouched BH . This skews the theoretical points in the $`n_\mathrm{s}`$$`r`$ plane towards the blue-tilted region, with models having a high tensor-to-scalar ratio suffering the largest shifts.
## II Formalism
### II.1 The patch universe
We shall assume that the universe is confined into a 3-brane embedded in a five-dimensional bulk. Matter lives on the brane only, while gravitons are free to propagate in the bulk. This is guaranteed as long as the effective energy density measured by the observer is smaller than the mass scale of the theory, $`\rho <m_5^4`$. By restricting our attention to the inflationary physics, we can disregard the contribution of the projected Weyl tensor at large scales.
If the energy density on the brane is comparable to the potential energy adopted to stabilize the extra dimension, the bulk backreacts and the effective cosmological evolution is modified by nonstandard terms. In particular, we can describe the primordial universe, at least in some finite time interval or energy patch, by the Friedmann equation
$$H^2=\beta _q^2\rho ^q,$$
(1)
where $`q`$ is constant and $`\beta _q>0`$ is a constant factor with energy dimension $`[\beta _q]=E^{12q}`$. This equation encodes a number of situations, including the pure 4D radion-stabilized regime ($`q=1`$, also known as the standard cosmology), the high-energy limit of the Randall–Sundrum braneworld ($`q=2`$), and the high-energy limit of the Gauss–Bonnet scenario ($`q=2/3`$). Defining the parameter
$$\theta 2\left(1q^1\right),$$
(2)
Eq. (1) can be recast as
$$H^{2\theta }=\beta _q^{2\theta }\rho .$$
(3)
Deviations from the 4D case $`\theta =0`$ will characterize exotic scenarios (braneworlds or modified gravities) according to their magnitude and sign.
Because of local conservation of the energy–momentum tensor, for a perfect fluid with pressure $`p`$ the continuity equation is
$$\dot{\rho }+3H(\rho +p)=0.$$
(4)
The dynamics of the patch inflationary universe has been extensively explored in the case of both an ordinary scalar inflaton $`\varphi `$ and a DBI tachyon $`T`$; for a complete review of results and a list of references see Ref. PhD . It is useful to generically label the inflaton field as $`\psi `$ and introduce a new parameter
$`\stackrel{~}{\theta }`$ $`=`$ $`\theta \text{for }\psi =\varphi \text{,}`$ (5a)
$`\stackrel{~}{\theta }`$ $`=`$ $`2\text{for }\psi =T\text{.}`$ (5b)
Let $`H_0`$ be the Hubble parameter evaluated at some reference time $`t_0`$ (later it will correspond to the integration starting time). The quantity
$$\alpha _q\frac{1}{\beta _q}\sqrt{\frac{2}{3q}\left(\frac{\beta _q}{H_0}\right)^{\stackrel{~}{\theta }}},$$
(6)
has dimension $`[\alpha _q]=E^{q(2\stackrel{~}{\theta })1}`$ ($`=E^1`$ for $`\psi =T`$ and $`=E`$ for $`\psi =\varphi `$), and can be absorbed in the normalization of the scalar field, so that the latter becomes dimensionless. In the following we use the redefinition $`\psi \psi /\alpha _q`$ everywhere; all dropped $`\alpha _q`$ factors can be recovered by counting the number of $`\psi `$ derivatives.
### II.2 Slow-roll parameters and flow equations
The amount of inflation is described by the number of $`e`$-foldings
$$N(t)\mathrm{ln}(a_\mathrm{f}/a)=_t^{t_\mathrm{f}}H(t^{})𝑑t^{},$$
(7)
this “backward” definition measuring the number of remaining $`e`$-foldings at the time $`t`$ before the end of inflation at $`t_\mathrm{f}`$. In the flow approach, the cosmological variables during inflation are written as functions of $`N`$ or $`\psi `$; one can shift from one picture to the other via
$$\frac{d}{dN}=\frac{d}{Hdt}=\left(\frac{H_0}{H}\right)^{\stackrel{~}{\theta }}\frac{H^{}}{H}\frac{d}{d\psi },$$
(8)
where a prime denotes differentiation with respect to $`\psi `$.
We define the flow parameters as
$`\lambda _0`$ $``$ $`ϵ{\displaystyle \frac{d\mathrm{ln}H}{dN}}`$ (9a)
$`=`$ $`\left({\displaystyle \frac{H}{H_0}}\right)^{\stackrel{~}{\theta }}{\displaystyle \frac{\dot{\psi }^2}{H^2}}=\left({\displaystyle \frac{H_0}{H}}\right)^{\stackrel{~}{\theta }}\left({\displaystyle \frac{H^{}}{H}}\right)^2,`$
$`\lambda _{\mathrm{}}`$ $``$ $`\left({\displaystyle \frac{H_0}{H}}\right)^{\stackrel{~}{\theta }\mathrm{}}{\displaystyle \frac{H^{(\mathrm{}+1)}(H^{})^\mathrm{}1}{H^{\mathrm{}}}},\mathrm{}1.`$ (9b)
where $`(n)`$ indicates the $`n`$-th $`\psi `$ derivative. Usually in literature they are dubbed $`{}_{}{}^{\mathrm{}}\lambda _{H}^{}`$. Note that this definition can be obtained from the potential slow-roll (SR) tower of Ref. PhD with the substitutions $`\varphi \psi `$, $`VH`$, and $`q1+\stackrel{~}{\theta }`$. In four dimensions and with an ordinary scalar field this definition coincides with that of Ref. kin02 .
It is also convenient to consider the following quantity:
$$\sigma 2\lambda _1(4+\theta +\stackrel{~}{\theta })ϵ.$$
(10)
The set of evolution equations reads, from Eqs. (8), (9), and (10),
$`{\displaystyle \frac{dϵ}{dN}}`$ $`=`$ $`ϵ[\sigma +(2+\theta )ϵ],`$ (11a)
$`{\displaystyle \frac{d\sigma }{dN}}`$ $`=`$ $`2\lambda _2(5+2\stackrel{~}{\theta }+\theta )ϵ\sigma `$ (11b)
$`(4+\theta +\stackrel{~}{\theta })(3+\theta +\stackrel{~}{\theta })ϵ^2,`$
$`{\displaystyle \frac{d\lambda _{\mathrm{}}}{dN}}`$ $`=`$ $`\lambda _{\mathrm{}+1}+{\displaystyle \frac{1}{2}}\lambda _{\mathrm{}}\{(\mathrm{}1)\sigma `$ (11c)
$`+[2\mathrm{}\overline{\theta }(4+\theta +\stackrel{~}{\theta })]ϵ\},\mathrm{}2,`$
where $`2\overline{\theta }2+\theta \stackrel{~}{\theta }`$.
As shown in Ref. lid03 , the flow parameters and equations do not encode the inflationary dynamics, since they are only a set of identities for the Hubble rate $`H`$ as a function of $`N`$. In fact, by definition the next-to-lowest-order SR terms appear with higher powers, and one can approximate the dynamics through a power truncation of the traditional SR tower. In the case of the flow parameters, Eq. (9), this is achieved by imposing a constraint such as $`H^{(\overline{n}+1)}=0`$ and $`H^{(\overline{n})}0`$, for some maximum $`\overline{n}`$. However, the method still provides an algorithm for generating inflationary models in the space of the cosmological observables, corresponding to a Taylor expansion of $`H(\psi )`$ lid03 .
For reference, the flow parameters are related to the Hubble-slow-roll parameters, defined by LPB ; PhD
$`ϵ_n`$ $``$ $`\left({\displaystyle \frac{H_0}{H}}\right)^{\stackrel{~}{\theta }}{\displaystyle \frac{H^{}}{H}}\left[{\displaystyle \frac{H^{\stackrel{~}{\theta }}}{H^{}}}\left({\displaystyle \frac{H^{}}{H^{\stackrel{~}{\theta }}}}\right)^{(n)}\right]^{1/n},n1,`$
via the relation
$$ϵ_n^n=\lambda _n\frac{H^{\stackrel{~}{\theta }}}{H^{(n+1)}}\left(\frac{H^{}}{H^{\stackrel{~}{\theta }}}\right)^{(n)}.$$
(13)
For a comparison between the two towers from an observational point of view, see Ref. mak05 .
### II.3 Attractors and observables
It is easy to find the attractors of the flow and project them on the $`\sigma `$$`ϵ`$ plane. One is the line corresponding to de Sitter (dS) fixed points
$`ϵ^{}`$ $`=`$ $`0,`$ (14a)
$`\sigma ^{}`$ $`=`$ $`\text{const},`$ (14b)
$`\lambda _{\mathrm{}}^{}`$ $`=`$ $`0,\mathrm{}2,`$ (14c)
and the other is the power-law inflationary attractor
$`ϵ^{}`$ $`=`$ $`\text{const},`$ (15a)
$`\sigma ^{}`$ $`=`$ $`(2+\theta )ϵ,`$ (15b)
$`\lambda _2^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+\stackrel{~}{\theta })(2+\stackrel{~}{\theta })ϵ^2,`$ (15c)
$`\lambda _{\mathrm{}+1}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[2+(1+\mathrm{})\stackrel{~}{\theta }\right]\lambda _{\mathrm{}}^{}ϵ^{},\mathrm{}2.`$ (15d)
Note that the slope of the line Eq. (15b) depends on the braneworld model one is considering, but not on the type of scalar field on the brane.
To first-order in slow roll the cosmological observables are
$$r\frac{ϵ}{\zeta _q},n_\mathrm{s}1\sigma ,$$
(16)
where $`rA_\mathrm{t}^2/A_\mathrm{s}^2`$ is the ratio between the tensor and scalar perturbation amplitudes and $`n_\mathrm{s}`$ is the scalar spectral index. The coefficient $`\zeta _q`$ is equal to 1 in the 4D and GB cases DLMS , while $`\zeta _2=2/3`$ in the RS scenario LMW . The expression for $`r`$ is valid in all cases considered, while the expression for $`n_\mathrm{s}`$ is for the commutative geometry case only and will be modified later to address the noncommutative case.
The stability analysis of the above fixed points against linear perturbations in the flow parameters involves an infinite-dimensional parameter space. However, one can truncate the flow tower at $`\mathrm{}`$-th order and study the system order-by-order, seeking convergence of dynamical properties.
Define the vector $`X(ϵ,\sigma ,\lambda _2,\mathrm{},\lambda _{\mathrm{}})^t`$. The perturbation equation reads
$$\frac{d\delta X}{dN}=M\delta X,$$
(17)
where $`\delta X(\delta ϵ,\delta \sigma ,\mathrm{})^t`$ encodes the linear perturbations and $`M`$ is a $`(\mathrm{}+1)\times (\mathrm{}+1)`$ matrix whose elements $`m_{ij}`$ ($`i,j=0,\mathrm{}\mathrm{}`$) are evaluated at a fixed point. For $`i<2`$ (the first two rows) one has
$`m_{00}`$ $`=`$ $`2(2+\theta )ϵ^{},m_{01}=ϵ^{},`$ (18a)
$`m_{10}`$ $`=`$ $`[(5+2\stackrel{~}{\theta }+\theta )\sigma ^{}+2(4+\theta +\stackrel{~}{\theta })(3+\theta +\stackrel{~}{\theta })ϵ^{}],`$
$`m_{11}`$ $`=`$ $`(5+2\stackrel{~}{\theta }+\theta )ϵ^{},m_{12}=2,`$ (18c)
$`m_{ij}`$ $`=`$ $`0,ji+2,`$ (18d)
while for $`i2`$
$`m_{i0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[2i\overline{\theta }(4+\theta +\stackrel{~}{\theta })]\lambda _i^{},`$ (18e)
$`m_{i1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(i1)\lambda _i^{},`$ (18f)
$`m_{ii}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{(i1)\sigma ^{}+[2i\overline{\theta }(4+\theta +\stackrel{~}{\theta })]ϵ^{}\},`$
$`m_{i,i+1}`$ $`=`$ $`1,`$ (18h)
$`m_{ij}`$ $`=`$ $`0,1<j<i,ji+2.`$ (18i)
The stability condition requires that the real part of all the eigenvalues $`\gamma `$ of the matrix $`M`$ is nonnegative.<sup>1</sup><sup>1</sup>1In the following the operation $`\gamma \mathrm{}(\gamma )`$ is understood. The sign for the stability depends on the convention set by Eq. (7); nonpositive eigenvalues correspond to fixed points that are stable in the past.
For the dS fixed points, the Jordan equation $`det(M\gamma )=0`$ is simply
$$\gamma ^2\left(\frac{\sigma ^{}}{2}\gamma \right)\mathrm{}\left[\frac{(\mathrm{}1)\sigma ^{}}{2}\gamma \right]=0,$$
(19)
and stability is guaranteed for $`\sigma ^{}>0`$ (blue-tilted spectra).
The power-law case is more complicated since one cannot easily diagonalize the matrix. The only notable exception is the Gauss–Bonnet braneworld with a normal scalar field. In that case, $`\lambda _{\mathrm{}}^{}=0`$ for $`\mathrm{}2`$ and the eigenvalue equation reads
$$\left(\gamma ^22ϵ^2\right)\left(\frac{ϵ^{}}{2}\gamma \right)\mathrm{}\left[\frac{(\mathrm{}1)ϵ^{}}{2}\gamma \right]=0.$$
(20)
The first two eigenvalues are equal in absolute value and opposite in sign. Therefore the power-law fixed points in GB $`\varphi `$-inflation are a repellor in both backward and forward integration.
In other cases of interest a full analysis is required; moreover, one should check the stability in both backward and forward time integration and along a “sufficient” number of directions in the parameter space. If all the eigenvalues have the same sign, then there is an attractor either in the past or in the future. Conversely, if some of them have relative opposite sign, there is no attractor. The truncation order is also important, since higher-order terms can change the sign of the eigenvalues.
To determine the minimum level at which the stability analysis is reliable, we have solved the Jordan equation at the power-law fixed points from first order ($`2\times 2`$ matrix) to $`10^{\mathrm{th}}`$ order in the flow parameters ($`\mathrm{}=10`$), and checked that the real parts of the eigenvalues have different signs at higher orders. Figure 1 shows the eigenvalues of the 4D and RS models with a normal scalar field and the GB tachyon case. The GB $`\varphi `$ case is as described above, while the other tachyon plots are quite similar to that presented here.
Qualitatively these results do not depend on the actual value of $`ϵ^{}`$; the only assumption is $`ϵ^{}>0`$. As regards the level of approximation, for the normal scalar field and the 4D tachyon the first order is already sufficient to give the result holding at all higher orders, while in the RS and GB tachyon cases all eigenvalues are negative up to $`2^{\mathrm{nd}}`$ and $`3^{\mathrm{rd}}`$ order, respectively. At higher orders one eigenvalue becomes and remains positive. Note that the degeneracy of one of the eigenvalues (the lowest one in the RS $`\varphi `$ and GB $`T`$ cases) is lifted every two orders (at even orders in the cited example).
The integration of the flow equations, performed at $`5^{\mathrm{th}}`$ order (see below), confirms these results. Therefore one can conclude that the power-law fixed points are unstable in both time directions, for all the present cosmological models, and at any (high) order in the flow truncation. The latter has been shown to be consistent in all cases from $`4^{\mathrm{th}}`$ order on.
## III Numerical integration of the flow equations
In order to study the evolution of the flow parameters, we have to establish a truncation level for the flow hierarchy and set suitable ranges of initial conditions for the parameters themselves within which they are randomly selected. Early flow papers solved the flow hierarchy numerically kin02 , but a more efficient approach RL3 exploits an analytic solution to the flow equations found in Ref. lid03 :
$$H(\psi )=H_0\left[1+\underset{\mathrm{}=1}{\overset{M+1}{}}A_{\mathrm{}}\psi ^{\mathrm{}}\right],$$
(21)
where $`H_0=H(t_0)`$ is evaluated at the initial integration time $`t_0`$; without loss of generality we choose $`\psi (t_0)=0`$. This solution continues to be valid in all cases considered here. With the above normalization for the scalar field the coefficients $`A_{\mathrm{}}`$ are dimensionless:
$$A_{\mathrm{}}=\frac{(1)^{\mathrm{}}\lambda _\mathrm{}1}{\mathrm{}!ϵ^{(\mathrm{}2)/2}}|_{\psi =0},\mathrm{}1,$$
(22)
where $`\lambda _{\mathrm{}1=0}=1`$. The choice of the $``$ sign, coming from the $`(H^{})^\mathrm{}1`$ term in Eq. (9b), is a convention determined by the rolling direction of the scalar field, which in this case is such that $`\dot{\psi }>0`$:
$$\frac{H^{}}{H}=\left(\frac{H}{H_0}\right)^{\stackrel{~}{\theta }/2}\sqrt{ϵ}.$$
(23)
Note the coefficients Eq. (22) have a different normalization with respect to Ref. RL3 , due to dimensionless factors such as $`\sqrt{4\pi }`$ absorbed into the definition of $`\psi `$.
To obtain the full evolution of the flow parameters, one needs only numerically integrate the relation between $`N`$ and $`\psi `$, which from Eq. (8) is
$$dN=\sqrt{\frac{1}{ϵ}\left(\frac{H}{H_0}\right)^{\stackrel{~}{\theta }}}d\psi .$$
(24)
For each model we used the same sequence of 40,000 initial conditions chosen within the ranges kin02
$`ϵ_0`$ $`=`$ $`[0,0.8],`$
$`\sigma _0`$ $`=`$ $`[0.5,0.5],`$
$`\lambda _{2,0}`$ $`=`$ $`[0.05,0.05],`$ (25)
$`\lambda _{3,0}`$ $`=`$ $`[0.005,0.005],`$
$`\mathrm{}`$
$`\lambda _{6,0}`$ $`=`$ $`0.`$
For an implementation at $`5^{\mathrm{th}}`$ order in slow-roll the value of the last parameter truncates the series and closes the hierarchy. Initial conditions which do not generate either a sufficient amount of inflation or a well-defined asymptotic behaviour are rejected. Conversely, when enough inflation is realized, then either
1. The evolution proceeds until the end of inflation, $`ϵ(t_\mathrm{f})=1`$. Then the flow equations are integrated backwards for a suitable number of $`e`$-foldings where the observables $`n_\mathrm{s}`$ and $`r`$ are evaluated and plotted, or
2. A late-time attractor is reached and the observables are read off at that point.
The reference number of $`e`$-folds for the backward/forward integration can be either fixed or randomly chosen within a range. Here we take $`N_0=50`$, although the results will not depend significantly upon its precise value kin02 . Also, we can distinguish points evolving to $`ϵ_\mathrm{f}=1`$ in two subcases:
1. More than 50 $`e`$-foldings of inflation are obtained from the initial point, meaning that the location where the observables are read off was reached by forwards integration from the initial condition.
2. Less than 50 $`e`$-foldings are obtained from the initial point, so that the point corresponding to the observables is obtained by integrating backwards in time from the initial condition.
The following results were all obtained via the approach described above. Additionally we checked that integrating the full flow hierarchy gives the same results point-by-point, which is a useful cross-check of the numerical implementation.
## IV Results
We divide our results into the (standard) commutative case and the noncommutative case. In the latter the flow equation evolution is unchanged, but the expression for the spectral index Eq. (16) is altered.
### IV.1 The commutative case
Figure 2 shows the flow points in the $`n_\mathrm{s}`$$`r`$ plane for an ordinary and a tachyonic scalar field in the 4D, RS, and GB scenarios. The classification of points is model dependent and summarized in Table 1. In all cases the attractor points are the dominant class, and hence the blue-tilted region ($`n_\mathrm{s}>1`$) is more populated than the red-tilted one.
The points that finish inflation with 50 $`e`$-foldings or more integrating forwards in time are a very small group of the three considered categories. Most of the points that finish inflation do so with less than 50 $`e`$-foldings. Note that what one does with the integration backwards is to find the point in parameter space where the initial condition should have been in order to give 50 $`e`$-foldings or more. This does not necessarily correspond to the process that gave rise to the physical initial conditions for inflation to start.
Figure 2 roughly shows that the swathe near the power-law attractor decreases for increasing $`\theta `$, so that its vertical extension decreases from GB to 4D to RS. This is clear from Eq. (11a), since the second positive term in square brackets becomes more and more important. This seems to indicate faster flow, which would naturally give less height to the swathe as the natural flow is towards the $`r=0`$ line. Note that the swathe corresponds to the flow being integrated backwards in time from the initial point in order to generate 50 $`e`$-foldings.
Comparing the percentages, it is clear that in the GB $`\varphi `$ case there is an increasing number of backward integrated points, eating off about 10% of the population of attractor points. Despite this fact, the swathe loses coherence and points spread in a large region in the observable plane, making the above descriptions rather approximate. Although there is no particular feature in the GB power-law attractor, the plot points do not lie close to it as in the other cases. This feature is a direct consequence of the flow equations. Equation (11) shows that the relevant parameter governing the cosmological evolution is $`\theta `$ rather than $`q`$; then the GB flow ($`\theta =1`$) would act in the opposite direction with respect to the RS ($`\theta =1`$) case, the 4D one being intermediate between the two. More precisely, while in RS the points are more concentrated at the region around $`r=0`$, $`n_\mathrm{s}=1`$, it is natural to expect more a widespread point distribution for GB. This does not necessarily imply that the GB braneworld is more severely constrained than the RS one.
In the GB $`\varphi `$ system, contrary to what happens in the others, all equations can be written in terms of $`\lambda _2`$ and integrals of it. Since
$$\lambda _2=\frac{d\lambda _1}{dN}+(1+\stackrel{~}{\theta })ϵ\lambda _1,$$
(26)
for $`\theta =\stackrel{~}{\theta }=1=\overline{\theta }`$ one has
$`\lambda _1`$ $`=`$ $`{\displaystyle \lambda _2𝑑N},`$ (27)
$`\lambda _{\mathrm{}}`$ $`=`$ $`{\displaystyle [\lambda _{\mathrm{}+1}+(\mathrm{}1)\lambda _1\lambda _{\mathrm{}}]𝑑N},\mathrm{}2.`$ (28)
This property makes the numerical integration of the GB equation peculiar relative to the other models. Also, as shown in Sec. II.3, the extra directions orthogonal to the $`n_\mathrm{s}r`$ plane completely decouple in the dS and power-law cases ($`\lambda _{\mathrm{}}^{}=0`$ for $`\mathrm{}2`$). Therefore the two-dimensional analysis presented above is expected to encode all the relevant dynamical features.
In the tachyon case, the percentages do not change appreciably in different braneworlds, nor with respect to the ordinary scalar situation. This is because Eq. (8) is insensitive to the type of braneworld in the tachyon case and the information on the extra dimension is carried only by the parameter $`\lambda _1`$. Therefore the tachyon plots are very similar to one another.
Also, in this case plot points are squeezed towards the dS attractor and prefer lower values of $`r`$. In order to understand why, let us define $`\varrho _\psi ^2\sigma ^2+r^2`$; the difference of the (squared) radii for a fixed braneworld $`\{\zeta _q,\theta \}`$ is, from Eq. (10),
$$\mathrm{\Delta }^2\varrho _\varphi ^2\varrho _T^2=(2\theta )[4\lambda _1(10+3\theta )ϵ]ϵ.$$
(29)
For a given input set $`\{ϵ,\lambda _1,\mathrm{}\}`$ of parameters, $`\mathrm{\Delta }^2>0`$ in any case of interest considered so far and for the main part of initial conditions, which explains why the tachyon points are concentrated near the scale-invariant origin of the $`\sigma `$$`r`$ plane, while $`\varphi `$ points are more widespread through the plane. We verified numerically that $`\mathrm{\Delta }^2<0`$ in less than 10% of the cases. A similar argument, with the SR parameters now depending on the scalar field, can be applied for the large-field model plots of Ref. CT .
Note that these results are completely determined by the choice of the patch parameter $`\stackrel{~}{\theta }`$. The parameter $`\zeta _q`$ would need to be extracted from the two-point correlation function of the graviton zero-mode within the original gravitational theory (3-brane in Einstein or Gauss–Bonnet gravity). However, one can redraw the plots in terms of the quantity $`\zeta _qr=O(r)`$, which differs from $`r`$ only in the RS scenario. In this case, the only effect is a further squeezing (by a factor $`2/3`$) of those points towards the dS attractor.
### IV.2 The noncommutative case
Beside extra dimensions, one can take into account a modification of the space-time geometry at the quantum level. In particular, a space-time uncertainty principle is believed to be a universal property of high-energy theories motivated by strings yon87 . This can be realized by a noncommutative algebra which preserves homogeneity and isotropy in the cosmological context BH :
$$[\tau ,x]=il_s^2,$$
(30)
where $`\tau a𝑑t`$ is a redefinition of time, $`x`$ is a spatial coordinate on the brane, and $`l_s`$ is the fundamental string length. In the simplest version of this model, moduli fields are fixed so that the extra dimension commutes with all the other directions.
The presence of a particular scale further breaks scale invariance and the perturbation spectra are affected accordingly. The detailed consequences of this ansatz were explored elsewhere PhD ; CT ; HL1 . Here we just recall that at large scales (corresponding to the strongly noncommutative limit) there are basically two different implementations for noncommutativity (“class 1” and “class 2”), depending on whether the Friedmann–Robertson–Walker 2-sphere is factored out from the measure of the action for the cosmological perturbations or not. In both cases the background, flow equations, percentage populations of Table 1, and stability analysis are not modified, and the new ingredient is encoded only in the expression for the scalar spectral index:
$$n_\mathrm{s}1\sigma +cϵ,$$
(31)
where $`\sigma `$ is defined in Eq. (10) and the constant $`c`$ is $`c=0`$ in the commutative case, $`c=2`$ in the class 2 noncommutative model, and $`c=6`$ in the class 1 model. Since the scalar and tensor amplitudes are multiplied by a same extra factor depending on $`l_s`$, their ratio $`r`$ is unchanged. Therefore the points in the upper region of the $`n_\mathrm{s}`$$`r`$ plane are shifted towards the right (blue-tilted spectra), while those lying on the $`r=0`$ line are unaffected. As one can see in Fig. 3, for $`c=6`$ the swathes point sharply rightwards instead of leftwards as in the commutative case. For $`c=2`$ the swathes are almost vertical. Thus the effect of noncommutativity is a skewing in the plane whose magnitude depends on the value of the parameter $`c`$ CT .
Note that since the noncommutative correction to the spectral index is positive, the stability condition for the dS fixed points is again satisfied for blue-tilted spectra.
For brevity we have not shown tachyon plots in the noncommutative case, but their shape can readily be pictured by applying the same skewing to the points shown in Fig. 2.
### IV.3 Relation to observations
We conclude by comparing the above flow plots with observations, with an analysis similar to that of Peiris et al. pei03 and Makarov mak05 for the standard cosmology. Present observations are highly restrictive in the $`n_\mathrm{s}`$$`r`$ plane, with only a small region near the origin still allowed in the plots we showed (very roughly, the allowed region corresponds to $`r<0.03`$ and $`0.95<n_\mathrm{s}<1.1`$, but one must allow for those parameters being correlated). The regions where our plots show major differences are thus already excluded. We take the 95% likelihood bound of Ref. sel04 , conservatively choosing the wider region obtained without including Ly$`\alpha `$-forest data, and compute in each case the fraction of points lying within that region (which we approximate by a suitably-oriented ellipse, noting that their $`r`$ is defined to be 16 times ours). We do not impose a constraint from the spectral index running.
The percentage of points inside the allowed region is shown in Table 2. In all cases the percentages are quite small, indicating that observations have already chopped off a substantial part of inflationary model space. Bearing in mind that the majority of the points are late-time attractor points with $`n_\mathrm{s}>1`$ and negligible $`r`$, the strong constraint $`n_\mathrm{s}<1.04`$ at $`r=0`$ sel04 plays a major role in ensuring the percentages are so small. If this limit were further strengthened in the future to exclude blue-tilted models, all those points would be lost.
There are a couple of trends evident in the numbers. Firstly, the tachyon cases tend to give higher percentages, due to their suppression of the tensor ratio $`r`$. Secondly, for the normal scalar field, increasing the noncommutativity parameter $`c`$ increases the percentages as it skews the swathes across to where the observations lie. Nevertheless, every cosmology proves capable of generating models within the allowed region, which prevents any useful conclusions from being drawn.
## V Conclusions
In this paper we have analyzed braneworld and tachyon scenarios in terms of the flow evolution equations. The equations of motion and the stability of the associated fixed points have been considered in detail in the general relativistic, Randall–Sundrum, and Gauss–Bonnet cases at high order in the flow parameters. Each model generates different predictions for the main cosmological observables (scalar spectral index $`n_\mathrm{s}`$ and tensor-to-scalar ratio $`r`$). In particular, Gauss–Bonnet gravity deeply modifies the flow structure in the parameter space for an ordinary inflaton field. In the cosmological tachyon case, the theoretical points in the $`n_\mathrm{s}`$$`r`$ plane are pressed towards scale invariance but the characteristic flow swathe does not lose coherence. We also discussed the imprint of a noncommutative geometry on the inflationary observables.
Some trends remain to be explored. In the braneworld inflationary scenarios the Weyl tensor is negligible at large scales and we have consistently neglected its contribution to the observables $`r`$ and $`n_\mathrm{s}`$. However, at small scales it can play an important role, as well as in the determination of other observables such as the bispectrum even at long wavelengths. Although we expect the flow approach to be unmodified by bulk physics, further investigation might clarify this point.
The flow equations approach is only one possible way in which one can generate an ensemble of inflationary models, corresponding to a Taylor expansion of $`H(\psi )`$. In Ref. RL3 , other methods were implemented in the general relativistic case to study the robustness of the flow equations structure, showing that there are significant variations in outcome between methods. It may be interesting to explore these differences in the braneworld and tachyon cases, too.
###### Acknowledgements.
The work of G.C. is supported by JSPS, A.R.L. by PPARC, and E.R. by Conacyt.
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# 1 Introduction
## 1 Introduction
The Toda lattice (TL) is one of the most important families of models in the theory of integrable systems. Its various generalizations and supersymmetric extensions, having deep implications in modern mathematical physics, have been the subject of intense investigations during the last decades.
The 2D TL hierarchy was first studied in , and at present two different nontrivial supersymmetric extensions of 2D TL are known. They are the $`N=(2|2)`$ - and $`N=(0|2)`$ supersymmetric TL hierarchies that possess a different number of supersymmetries and contain the $`N=(2|2)`$ and $`N=(0|2)`$ TL equations as subsystems. Quite recently, the 2D generalized fermionic TL equations have been introduced and their two reductions related to the $`N=(2|2)`$ and $`N=(0|2)`$ supersymmetric TL equations were considered. In the present paper, we describe a wide class of integrable two-dimensional fermionic Toda lattice hierarchies – 2D fermionic $`(K^+,K^{})`$-TL hierarchies, which includes the 2D $`N=(2|2)`$ and $`N=(0|2)`$ supersymmetric TL hierarchies as particular cases and contains the 2D generalized fermionic TL equations as a subsystem.
The Hamiltonian description of the 2D TL hierarchy has been constructed only quite recently in the framework of the R-matrix approach in , where the new R-matrix associated with splitting of algebra given by a pair of difference operators was introduced. In the present paper, we adapt this R-matrix to the case of $`Z_2`$-graded operators and derive the bi-Hamiltonian structure of the 2D fermionic $`(K^+,K^{})`$-TL hierarchy.
Remarkably, in solving this problem the generalized graded bracket (5) on the space of graded operators with an involution finds its new application. This bracket was introduced in , where it was observed that the $`N=(1|1)`$ supersymmetric 2D TL hierarchy had a natural Lax-pair representation in terms of this bracket which allowed one to derive the dispersionless $`N=(1|1)`$ 2D TL hierarchy and its Lax representation. In the present paper, the generalized graded bracket is used to describe the 2D fermionic $`(K^+,K^{})`$-TL hierarchy and define its two Hamiltonian structures. Moreover, we demonstrate that the classical graded Yang-Baxter equation has an equivalent operator representation in terms of the generalized graded bracket. All these facts attest that this bracket has a fundamental meaning and allows a broad spectrum of applications in modern mathematical physics.
This paper is the extended version of . The structure of the paper is as follows.
In Sec. 2, we define the space of the $`Z_2`$-graded difference operators with the involution and recall the generalized graded bracket and its properties.
In Sec. 3, we give a theoretical background of the $`R`$-matrix method generalized to the case of the $`Z_2`$-graded difference operators. We define the $`R`$-matrix on the associative algebra g of the $`Z_2`$-graded difference operators, derive the graded modified Yang-Baxter equation and using the generalized graded bracket obtain two Poisson brackets for the functionals on $`𝗀^{}=𝗀`$. The proper properties of the Poisson brackets thus obtained are provided by the properties of the generalized graded bracket. Thus, for the $`Z_2`$-graded difference operators of odd (even) parity this bracket defines odd (even) first Poisson bracket. The second Poisson bracket is found only for even difference operators which in this case are compatible with the first Poisson bracket. Using these Poisson brackets one can define the Hamiltonian equations that can equivalently be rewritten in terms of the Lax-pair representation. The basic results of Sec. 3 are formulated as Theorem.
In Sec. 4, using the generalized graded bracket we propose a new 2D fermionic $`(K^+,K^{})`$-TL hierarchy in terms of the Lax-pair representation and construct the algebra of its flows. Then we present the explicit expression for its flows and show that all known up to now 2D TL equations can be derived from this hierarchy as subsystems.
In Sec. 5, we consider the reduction of the 2D fermionic $`(K^+,K^{})`$-TL hierarchy to the 1D space and reproduce the 1D generalized fermionic TL equations as the first flow of the reduced hierarchy with additional constraint imposed.
In Secs. 6 and 7, we apply the results of Sec. 3 to derive the Hamiltonian structures of the 1D and 2D fermionic $`(K^+,K^{})`$-TL hierarchies. Following we use the $`R`$-matrix which acts nontrivially on the space of the direct sum of two difference operators and derive two different Hamiltonian structures of the 2D fermionic $`(K^+,K^{})`$-TL hierarchy. The first Hamiltonian structure is obtained for both even and odd values of $`(K^+,K^{})`$ while the second one is found for even values of $`(K^+,K^{})`$ only. We perform their Dirac reduction and demonstrate that in general the Dirac brackets for the second Hamiltonian structure are nonlocal but for the case of the fermionic $`(2,2)`$-TL hierarchy they become local. As an example, we give the explicit form of the first and second Hamiltonian structures for the fermionic 1D (2,2)-TL hierarchy.
In Sec. 8, we briefly summarize the main results obtained in this paper and point out open problems. In Appendices, we clarify some technical aspects.
## 2 Space of difference operators
In this section, we define the space of difference operators which will play an important role in our consideration. These operators can be represented in the following general form:
$`𝕆_m={\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}f_{k,j}^{(m)}e^{(km)},m,j,`$ (1)
parameterized by the functions $`f_{2k,j}^{(m)}`$ ($`f_{2k+1,j}^{(m)}`$) which are the $`Z_2`$-graded bosonic (fermionic) lattice fields with the lattice index $`j`$ $`(j)`$ and the Grassmann parity defined by index $`k`$
$`d_{f_{k,j}^{(m)}}=|k|\text{mod}2.`$
$`e^k`$ is the shift operator whose action on the lattice fields results into a discrete shift of a lattice index
$`e^lf_{k,j}^{(m)}=f_{k,j+l}^{(m)}e^l.`$ (2)
The shift operator has $`Z_2`$-parity defined as
$`d_{e^l}^{}=|l|\text{mod}2.`$
The operators $`𝕆_m`$ (1) admit the diagonal $`Z_2`$-parity
$`d_{𝕆_m}=d_{f_{k,j}^{(m)}}+d_{e^{(km)}}^{}=|m|\text{mod}2`$ (3)
and the involution
$`𝕆_m^{}={\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}(1)^kf_{k,j}^{(m)}e^{(km)}.`$
In what follows we also need the projections of the operators $`𝕆_m`$ defined as
$`(𝕆_m)_p={\displaystyle \underset{kp+m}{}}f_{k,j}e^{(km)},(𝕆_m)_p={\displaystyle \underset{kp+m}{}}f_{k,j}e^{(km)}`$
and we will use the usual notation for the projections $`(𝕆_m)_+:=(𝕆_m)_0`$ and $`(𝕆_m)_{}:=(𝕆_m)_{<0}`$. Note that $`e^l`$ is a conventional form for the shift operators defined in terms of infinite-dimensional matrices $`(e^l)_{i,j}\delta _{i,jl}`$, and there is an isomorphism between operators (1) and infinite-dimensional matrices (see e.g. )
$`𝕆_m={\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}f_{k,j}^{(m)}e^{(km)}(𝕆_m)_{j,i}{\displaystyle \underset{k=\mathrm{}}{\overset{\mathrm{}}{}}}f_{k,j}^{(m)}\delta _{j,ik+m}.`$
In the operator space (1) one can extract two subspaces which are of great importance in our further consideration
$`𝕆_{K^\pm }^\pm `$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}f_{k,j}^\pm {}_{}{}^{\pm (K^\pm k)},K^\pm .`$ (4)
The operators of the subspaces $`𝕆_{K^\pm }^\pm `$ form associative algebras with the multiplication (2). Using this fact we define on these subspaces the generalized graded algebra with the bracket
$`[𝕆,\stackrel{~}{𝕆}\}:=𝕆\stackrel{~}{𝕆}(1)^{d_𝕆d_{\stackrel{~}{𝕆}}}\stackrel{~}{𝕆}^{(d_𝕆)}𝕆^{(d_{\stackrel{~}{𝕆}})},`$ (5)
where the operators $`𝕆`$ and $`\stackrel{~}{𝕆}`$ belong to the subspaces $`𝕆_{K^+}^+(𝕆_K^{}^{})`$, and $`𝕆^{(m)}`$ denotes the $`m`$-fold action of the involution $``$ on the operator $`𝕆`$, ($`𝕆^{(2)}=𝕆`$). Bracket (5) generalizes the (anti)commutator in superalgebras and satisfies the following properties :
symmetry
$`[𝕆,\stackrel{~}{𝕆}\}=(1)^{d_𝕆d_{\stackrel{~}{𝕆}}}[\stackrel{~}{𝕆}^{(d_𝕆)},𝕆^{(d_{\stackrel{~}{𝕆}})}\},`$ (6)
derivation
$`[𝕆,\stackrel{~}{𝕆}\widehat{𝕆}\}=[𝕆,\stackrel{~}{𝕆}\}\widehat{𝕆}+(1)^{d_𝕆d_{\stackrel{~}{𝕆}}}\stackrel{~}{𝕆}^{(d_𝕆)}[𝕆^{(d_{\stackrel{~}{𝕆}})},\widehat{𝕆}\},`$ (7)
and Jacobi identity
$`(1)^{d_𝕆d_{\widehat{𝕆}}}[[𝕆,\stackrel{~}{𝕆}^{(d_𝕆)}\},\widehat{𝕆}^{(d_𝕆+d_{\stackrel{~}{𝕆}})}\}+(1)^{d_{\stackrel{~}{𝕆}}d_𝕆}[[\stackrel{~}{𝕆},\widehat{𝕆}^{(d_{\stackrel{~}{𝕆}})}\},𝕆^{(d_{\stackrel{~}{𝕆}}+d_{\widehat{𝕆}})}\}`$
$`+(1)^{d_{\widehat{𝕆}}d_{\stackrel{~}{𝕆}}}[[\widehat{𝕆},𝕆^{(d_{\widehat{𝕆}})}\},\stackrel{~}{𝕆}^{(d_{\widehat{𝕆}}+d_𝕆)}\}`$ $`=`$ $`0.`$ (8)
For the operators $`𝕆_m`$ (1) we define the supertrace
$`str𝕆={\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(1)^jf_{m,j}^{(m)}.`$ (9)
In what follows we assume suitable boundary conditions for the functions $`f_{k,j}^{(m)}`$ in order the main property of supertraces
$`str[𝕆,\stackrel{~}{𝕆}\}=0`$ (10)
be satisfied for the case of the generalized graded bracket (5).
## 3 R-matrix formalism
In this section, we develop a theoretical background of the R-matrix method adapted to the case of the operator space (1).
Let g be an associative algebra of the operators from the space (1) with the invariant non-degenerate inner product
$`<𝕆,\stackrel{~}{𝕆}>=str(𝕆\stackrel{~}{𝕆})`$
using which one can identify the algebra g with its dual g. We set the following Poisson bracket:
$`\{f,g\}(𝕆)=<𝕆,[g,(f)^{(d_g)}\}>,`$ (11)
where $`f,g`$ are functionals on g, and $`f`$ and $`g`$ are their gradients at the point $`𝕆`$ which are related with $`f,g`$ through the inner product
$`{\displaystyle \frac{f(𝕆+ϵ\delta 𝕆)}{ϵ}}|_{ϵ=0}=<\delta 𝕆,f(𝕆)>.`$
Note that the proper properties of the Poisson bracket (11) follow from the properties (62) of the generalized bracket (5) and are strictly determined by the $`Z_2`$-parity of the operator $`𝕆`$. Thus, one has symmetry
$`\{f,g\}`$ $`=`$ $`(1)^{(d_f+d_𝕆)(d_g+d_𝕆)}\{g,f\},`$ (12)
derivation
$`\{f,gh\}`$ $`=`$ $`\{f,g\}h+(1)^{d_g(d_f+d_𝕆)}g\{f,h\},`$ (13)
and Jacobi identity
$`(1)^{(d_f+d_𝕆)(d_h+d_𝕆)}\{\{f,g\},h\}+(1)^{(d_g+d_𝕆)(d_f+d_𝕆)}\{\{g,h\},f\}`$
$`+(1)^{(d_h+d_𝕆)(d_g+d_𝕆)}\{\{h,f\},g\}`$ $`=`$ $`0.`$ (14)
Therefore, for the even operator $`𝕆`$ one has a usual (even) $`Z_2`$-graded Poisson bracket, while for the operators with odd diagonal parity $`d_𝕆`$ (3) eq. (11) defines the odd $`Z_2`$-graded Poisson bracket (antibracket).
Having defined the Poisson bracket we proceed with the search for the hierarchy of flows generated by this bracket using Hamiltonians. Therefore, we need to determine an infinite set of functionals which should be in involution to play the role of Hamiltonians. For Poisson bracket (11) one can find an infinite set of Hamiltonians in a rather standard way
$`H_k={\displaystyle \frac{1}{k}}str𝕆_{}^k={\displaystyle \frac{1}{k}}{\displaystyle \underset{i=\mathrm{}}{\overset{\mathrm{}}{}}}(1)^if_{km,i}^{(km)},`$ (15)
where $`𝕆_{}^k`$ is defined as
$`(𝕆)_{}^{2k}:=(𝕆^{(d_𝕆)}𝕆)^k,(𝕆)_{}^{2k+1}:=𝕆(𝕆)_{}^{2k}.`$ (16)
For the odd operators $`𝕆`$ eq. (15) defines only fermionic nonzero functionals $`H_{2k+1}`$, since in this case even powers of the operators $`𝕆`$ have the following representation:
$`d_𝕆=1:(𝕆)_{}^{2k}=(1/2[(𝕆)^{},𝕆\})^k1/2[((𝕆)_{}^{2k1})^{},𝕆\}`$ (17)
and all the bosonic Hamiltonians are trivial ($`H_{2k}=0`$ ) like the supertrace of the generalized graded bracket.
The functionals (15) are obviously in involution but produce a trivial dynamics. Actually, the functionals $`H_k`$ (15) are the Casimir operators of the Poisson bracket (11), so the Poisson bracket of $`H_k`$ with any other functional is equal to zero as an output (due to the relation $`H_{k+1}=𝕆_{}^k`$). Nevertheless, it is possible to modify the Poisson bracket (11) in such a way that the new Poisson bracket would produce nontrivial equations of motion using the same Hamiltonians (15) and these Hamiltonians are in involution with respect to the modified Poisson bracket as well. Let us introduce the modified generalized graded bracket on the space (1)
$`[𝕆,\stackrel{~}{𝕆}\}_R:=[R(𝕆),\stackrel{~}{𝕆}\}+[𝕆,R(\stackrel{~}{𝕆})\},`$ (18)
where the $`R`$-matrix is a linear map $`R`$: g $``$ g such that the bracket (18) satisfies the properties (62). One can verify that the Jacobi identities (2) for the bracket (18) can equivalently be rewritten in terms of the generalized graded bracket (5)
$`(1)^{d_𝕆d_{\widehat{𝕆}}}[[𝕆,\stackrel{~}{𝕆}^{(d_𝕆)}\}_R,\widehat{𝕆}^{(d_𝕆+d_{\stackrel{~}{𝕆}})}\}_R+\text{cycle perm.}=`$
$`(1)^{d_𝕆d_{\widehat{𝕆}}}[R([𝕆,\stackrel{~}{𝕆}^{(d_𝕆)}\}_R)[R(𝕆),R(\stackrel{~}{𝕆}^{(d_𝕆)})\},\widehat{𝕆}^{(d_𝕆+d_{\stackrel{~}{𝕆}})}\}+\text{cycle perm.}=0.`$
Thus, one can conclude that a sufficient condition for $`R`$ to be the $`R`$-matrix is the validity of the following equation:
$`R([𝕆,\stackrel{~}{𝕆}\}_R)[R(𝕆),R(\stackrel{~}{𝕆})\}=\alpha [𝕆,\stackrel{~}{𝕆}\},`$ (19)
where $`\alpha `$ is an arbitrary constant. One can show (see Appendix A) that for the case of graded antisymmetric operators $`R`$ eq. (19) at $`\alpha =1`$ represents the operator form of the graded classical Yang-Baxter equation . Following the terminology of we call eq. (19) the graded modified Yang-Baxter equation. Equation (19) is the generalization of the graded modified classical Yang-Baxter equation discussed in paper for the case of space of graded operators (1).
With the new bracket (18) one can define the corresponding new Poisson bracket on dual $`𝗀^{}`$:
$`\{f,g\}_1(𝕆)`$ $`=`$ $`1/2<𝕆,[g,(f)^{(d_g)}\}_R>`$ (20)
$``$ $`{\displaystyle \frac{1}{2}}<(1)^{d_gd_𝕆}R(g)^{(d_𝕆)}[𝕆^{(d_g)},(f)^{(d_g)}\}`$
$`[𝕆,g\}R((f)^{(d_g)})>.`$
With respect to the dependence of the r.h.s of (20) on the point $`𝕆`$ this is a linear bracket. Without going into details we introduce also bi-linear bracket for bosonic graded operators $`𝕆_B`$ ( $`d_{𝕆_B}=0`$) as follows:
$`\{f,g\}_2(𝕆_B)`$ $`=`$ $`1/4<[𝕆_B,g\}R((f)^{(d_g)}𝕆_B^{(d_f+d_g)}`$ (21)
$`+𝕆_B^{(d_g)}(f)^{(d_g)})R(g𝕆_B^{(d_g)}`$
$`+𝕆_Bg)[𝕆_B^{(d_g)},(f)^{(d_g)}\}>.`$
We did not succeed in constructing the bi-linear bracket for the case of fermionic operators $`𝕆_F`$ ( $`d_{𝕆_F}=1`$). The bracket (20) is obviously the Poisson bracket if $`R`$ is an $`R`$-matrix on $`𝗀`$. The bi-linear bracket (21) becomes Poisson bracket under more rigorous constraints which can be found in the following
Theorem. a) Linear bracket (20) is the Poisson bracket if $`R`$ obeys the graded modified Yang-Baxter equation (19);
b) the bi-linear bracket (21) is the Poisson bracket if $`R`$ and its graded antisymmetric part $`1/2(RR^{})`$ obey the graded modified Yang-Baxter equation (19) with the same $`\alpha `$;
c) if $`𝕆=𝕆_B`$, then these two Poisson brackets are compatible;
d) the Casimir operators $`H_k`$ (15) of the bracket (11) are in involution with respect to both linear (20) and bi-linear (21) Poisson brackets;
e) the Hamiltonians $`H_k0`$ (15) generate evolution equations
$`_k𝕆`$ $`=`$ $`\{H_{k+1},𝕆\}_1=1/2[R((H_{k+1})^{(d_𝕆)}),𝕆\},`$
$`_k𝕆_B`$ $`=`$ $`\{H_k,𝕆_B\}_2=1/4[R(H_k𝕆_B+𝕆_BH_k),𝕆_B\}`$ (22)
via the brackets (20) and (21), respectively, which connect the Lax-pair and Hamiltonian representations.
Proof. a). This is a summary of the above discussion on the linear bracket.
b). Using the property of cyclic permutations inside the supertrace one can easily verify that the symmetry property $`\{f,g\}_2=(1)^{d_fd_g}\{g,f\}_2`$ holds. Verification of the Jacobi identities for the bi-linear bracket amounts to straightforward and tedious calculations which are presented in Appendix B.
c). These two Poisson brackets are obviously compatible. Indeed, a deformation of the point $`𝕆_B𝕆_B+b`$ on the dual $`𝗀^{}`$, where $`b`$ is an arbitrary constant operator, transforms (21) into the sum of two Poisson brackets
$`\{f,g\}_2(𝕆_B+b)=\{f,g\}_2(𝕆_B)+b\{f,g\}_1(𝕆_B).`$
d). Substituting the expressions for Hamiltonians (15) into eqs. (20) and (21) and taking into account that $`H_{k+1}=𝕆_{}^k`$ it is easily to check that the Casimirs of the bracket (11) are in involution with respect to both the Poisson structures (20) and (21).
e). Using cyclic permutations inside the supertrace operation let us rewrite both the Poisson brackets (20) and (21) in the following general form:
$`\{f,g\}_i(𝕆)=<P_i(𝕆)g,(f)^{(d_g)}>,i=1,2,`$
where $`P_i(𝕆)`$ is the Poisson tensor corresponding to the bracket $`\{..,..\}_i`$
$`P_1(𝕆)g`$ $`=`$ $`1/2([𝕆,R(g)\}+R^{}([𝕆,g]\})),`$
$`P_2(𝕆_B)g`$ $`=`$ $`1/4([R(g𝕆_B+𝕆_Bg),𝕆_B^{(d_g)}\}`$
$``$ $`𝕆_BR^{}([𝕆_B,g\})R^{}([𝕆_B,g\})𝕆_B^{(d_g)})`$
and the adjoint operator $`R^{}`$ acts on the dual $`𝗀^{}`$
$$<𝕆,R(\stackrel{~}{𝕆})>=<R^{}(𝕆),\stackrel{~}{𝕆}>.$$
The Hamiltonian vector field associated with Hamiltonian $`H_k`$ is given by $`_k𝕆=P_i(𝕆)H_k`$. Taking into account that $`[𝕆,H_k\}=0`$ we arrive at the Lax-pair representations (3). $`\mathrm{}`$
Note that a similar Theorem when the shift operators and functions parameterizing the difference operators $`𝕆`$ (1) have even $`Z_2`$-parity was discussed in .
For the graded modified Yang-Baxter equation (19) there is a particular class of solutions which are useful in application. Suppose that the algebra g can be represented as a vector space direct sum of two subalgebras
$$𝗀=𝗀{}_{+}{}^{}𝗀{}_{}{}^{}:[𝗀_+,𝗀_+\}𝗀_+,[𝗀_{},𝗀_{}\}𝗀_{}.$$
Let $`P_\pm `$ be the projection operators on these subalgebras, $`P_\pm 𝗀=𝗀_\pm `$, then one can easily verify that $`R=P_+P_{}`$ satisfies the graded modified Yang-Baxter equation (19) at $`\alpha =1`$ and, therefore, represents the $`R`$-matrix on g. Indeed, in this case the modified generalized graded bracket (18)
$`[𝕆,\stackrel{~}{𝕆}\}_R=2[(𝕆)_+,(\stackrel{~}{𝕆})_+\}2[(𝕆)_{},(\stackrel{~}{𝕆})_{}\}`$ (23)
obviously satisfies the Jacobi identities (2), since it determines the usual direct sum of two subalgebras
$$[𝗀_\pm ,𝗀_\pm \}_R𝗀_\pm ,[𝗀_+,𝗀_{}\}_R=0.$$
## 4 2D fermionic $`(K^+,K^{})`$-Toda lattice hierarchy
In this section, we introduce the two-dimensional fermionic $`(K^+,K^{})`$-Toda lattice hierarchy in terms of the Lax-pair representation.
Let us consider two difference operators $`L_{K^\pm }^\pm `$
$`L_{K^+}^+={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}u_{k,i}e^{(K^+k)},L_K^{}^{}={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}v_{k,i}e^{(kK^{})},`$ (24)
which obviously belong to the spaces (4). The lattice fields and the shift operator entering into these operators have the following length dimensions: $`[u_{k,i}]=1/2k`$, $`[v_{k,i}]=1/2(kK^+K^{})`$ and $`[e^k]=1/2k`$, respectively, so operators (24) are of equal length dimension, $`[L_{K^+}^+]=[L_K^{}^{}]=1/2K^+`$. The dynamics of the fields $`u_{k,i},v_{k,i}`$ are governed by the Lax equations expressed in terms of the generalized graded bracket (5)
$`D_s^\pm L_{K^\alpha }^\alpha =\alpha (1)^{sK^\alpha K^\pm }[(((L_{K^\pm }^\pm )_{}^s)_\alpha )^{(K^\alpha )},L_{K^\alpha }^\alpha \},\alpha =+,,s,`$ (25)
where $`D_s^\pm `$ are evolution derivatives with the $`Z_2`$-parity defined as
$$d_{D_s^\pm }=sK^\pm \text{ mod }2$$
and the length dimension $`[D_s^+]=[D_s^{}]=sK^+/2.`$ The Lax equations (25) generate non-Abelian (super)algebra of flows of the 2D fermionic $`(K^+,K^{})`$-TL hierarchy
$`[D_s^\pm ,D_p^\pm \}=(1(1)^{spK^\pm })D_{s+p}^\pm ,[D_s^+,D_p^{}\}=0.`$
The composite operators $`(L_{K^\pm }^\pm )_{}^s`$ entering into the Lax equations (25) are defined by eq. (16) and also belong to the spaces (4)
$`(L_{K^+}^+)_{}^r:={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}u_{k,i}^{(r)}e^{(rK^+k)},(L_K^{}^{})_{}^r:={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}v_{k,i}^{(r)}e^{(krK^{})}.`$
Here $`u_{k,i}^{(r)}`$ and $`v_{k,i}^{(r)}`$ are functionals of the original fields and there are the following recursion relations for them
$`u_{p,i}^{(r+1)}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{p}{}}}(1)^{kK^+}u_{k,i}^{(r)}u_{pk,ik+rK^+},u_{p,i}^{(1)}=u_{p,i},`$
$`v_{p,i}^{(r+1)}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{p}{}}}(1)^{kK^{}}v_{k,i}^{(r)}v_{pk,i+krK^{}},v_{p,i}^{(1)}=v_{p,i}.`$
Now using the Lax representation (25) and relations (7) and (16) one can derive the equations of motion for the composite Lax operators
$`D_s^\pm (L_{K^\alpha }^\alpha )_{}^r`$ $`=`$ $`\alpha (1)^{srK^\alpha K^\pm }[(((L_{K^\pm }^\pm )_{}^s)_\alpha )^{(rK^\alpha )},(L_{K^\alpha }^\alpha )_{}^r\}.`$ (26)
The Lax-pair representation (26) generates the following equations for the functionals $`u_{k,i}^{(r)},v_{k,i}^{(r)}`$:
$`D_s^+u_{k,i}^{(r)}`$ $`=`$ $`{\displaystyle \underset{p=1}{\overset{k}{}}}((1)^{rpK^++1}u_{p+sK^+,i}^{(s)}u_{kp,ip}^{(r)}`$ (27)
$`+`$ $`(1)^{(k+p)sK^+}u_{kp,i}^{(r)}u_{p+sK^+,i+pk+rK^+}^{(s)}),`$
$`D_s^{}u_{k,i}^{(r)}`$ $`=`$ $`{\displaystyle \underset{p=0}{\overset{sK^{}1}{}}}((1)^{(sK^{}+p)rK^+}v_{p,i}^{(s)}u_{p+ksK^{},i+psK^{}}^{(r)}`$ (28)
$``$ $`(1)^{(k+p+1)sK^{}}u_{p+ksK^{},i}^{(r)}v_{p,ipk+sK^{}+rK^+}^{(s)}),`$
$`D_s^+v_{k,i}^{(r)}`$ $`=`$ $`{\displaystyle \underset{p=0}{\overset{sK^+}{}}}((1)^{(sK^++p)rK^{}}u_{p,i}^{(s)}v_{p+ksK^+,ip+sK^+}^{(r)}`$ (29)
$``$ $`(1)^{(k+p+1)sK^+}v_{p+ksK^+,i}^{(r)}u_{p,i+p+ksK^+rK^{}}^{(s)}),`$
$`D_s^{}v_{k,i}^{(r)}`$ $`=`$ $`{\displaystyle \underset{p=0}{\overset{k}{}}}((1)^{rpK^{}+1}v_{p+sK^{},i}^{(s)}v_{kp,i+p}^{(r)}`$ (30)
$`+`$ $`(1)^{(k+p)sK^{}}v_{kp,i}^{(r)}v_{p+sK^{},i+kprK^{}}^{(s)}).`$
It is assumed that in the right-hand side of eqs. (2730) all the functionals $`u_{k,i}^{(r)}`$ $`v_{k,i}^{(r)}`$ with $`k<0`$ should be set equal to zero.
Let us demonstrate that all known up to now 2D supersymmetric Toda lattice equations can be derived from the system (2730).
First, the 2D generalized fermionic Toda lattice equation discussed in can be reproduced from the system of equations (2730) as a subsystem with additional reduction constraints imposed. In order to see this, let us introduce the notation $`v_{0,i}=d_i,v_{1,i}=\rho _i,u_{1,i}=\gamma _i,u_{2,i}=c_i`$ and consider eqs. (28) and (30) at $`K^+=K^{}=2`$, $`r=s=1`$. One obtains
$`D_1^+d_i=d_i(c_ic_{i2}),D_1^{}\gamma _i=\rho _iu_{0,i1}\rho _{i+2}u_{0,i},`$
$`D_1^{}c_i=d_iu_{0,i2}d_{i+2}u_{0,i}\gamma _i\rho _{i+1}\gamma _{i1}\rho _i,`$
$`D_1^+\rho _i=\rho _i(c_ic_{i1})+d_{i+1}\gamma _id_i\gamma _{i2},D_1^{}u_{0,i}=0.`$ (31)
It is easy to check that after reduction $`u_{0,i}=1`$ eqs. (4) coincide with the 2D generalized fermionic Toda lattice equations up to time redefinition $`D_1^{}D_1^{}`$.
Next, the N=$`(2|2)`$ supersymmetric Toda lattice equation also belongs to the system (2730). In order to see that, let us consider eqs. (28) at $`K^+=K^{}=s=k=r=1`$
$`D_1^{}u_{1,i}=u_{0,i1}v_{0,i}u_{0,i}v_{0,i+1}`$ (32)
and eqs. (29) at $`K^+=K^{}=s=r=1`$, $`k=0`$
$`D_1^+v_{0,i}=v_{0,i}(u_{1,i}u_{1,i1}).`$ (33)
Then imposing the constraint $`u_{0,i}=1`$ and eliminating the fields $`u_{1,i}`$ from eqs. (3233) one obtains the N=$`(1|1)`$ superfield form of the N=$`(2|2)`$ supersymmetric Toda lattice equation
$`D_1^+D_1^{}\text{ ln }v_{0,i}=v_{0,i+1}v_{0,i1}.`$ (34)
Analogously, one can show that the consideration of eqs. (29) and (30) at $`K^+=1`$, $`K^{}=2`$, $`s=r=1`$ and $`k=0,1`$ leads to the N=$`(0|1)`$ superfield form of the N=$`(0|2)`$ supersymmetric Toda lattice equation after imposing the reduction constrains $`u_{0,i}=1`$, $`v_{0,2i+1}=0`$.
We call equations (25) for arbitrary $`(K^+,K^{})`$ the 2D fermionic $`(K^+,K^{})`$-Toda lattice hierarchy.
## 5 1D fermionic $`(K^+,K^{})`$-Toda lattice hierarchy
In this section, we consider the reduction of the 2D fermionic $`(K^+,K^{})`$-Toda lattice hierarchy for even values of $`(K^+,K^{})`$ to the 1D space.
Let $`(K^+,K^{})`$ be even numbers. In this case the generalized graded bracket (5) between two $`Z_2`$-even operators turns into the usual commutator and eqs. (25) become
$`D_s^\pm L_{K^\alpha }^\alpha `$ $`=`$ $`[((L_{K^\pm }^\pm )^s)_\pm ,L_{K^\alpha }^\alpha ].`$ (35)
Following for even $`(K^+,K^{})`$ one can impose the reduction constraint on the Lax operators (24) as follows:
$`L_{K^+}^++(L_{K^+}^+)^1=L_K^{}^{}+(L_K^{}^{})^1L_{K^+,K^{}},`$ (36)
which leads to the following explicit form for the reduced Lax operator
$`L_{K^+,K^{}}={\displaystyle \underset{k=0}{\overset{K^+}{}}}u_{k,i}e^{(K^+k)}+{\displaystyle \underset{k=0}{\overset{K^{}1}{}}}v_{k,i}e^{(kK^{})}{\displaystyle \underset{k=0}{\overset{K^++K^{}}{}}}\stackrel{~}{u}_{k,i}e^{(K^+k)}.`$ (37)
Substituting the expressions for the reduced composite Lax operators $`L_{K^\pm }^\pm =L_{K^+,K^{}}(L_{K^\pm }^\pm )^1`$ into Lax equations (35) one can see that these equations become equivalent to the single Lax equation on the reduced Lax operator
$`D_sL_{K^+,K^{}}`$ $`=`$ $`[((L_{K^+,K^{}})^s)_+,L_{K^+,K^{}}]`$ (38)
with $`D_s^+=D_s^{}=D_s.`$ As a consequence of eq. (38) we have
$`D_s(L_{K^+,K^{}})^r`$ $`=`$ $`[((L_{K^+,K^{}})^s)_+,(L_{K^+,K^{}})^r].`$
At the reduction $`u_{0,i}=1`$ the 1D $`(2,2)`$-TL hierarchy becomes that studied in detail in . In this case, the representation (38) with the Lax operator
$`L_{2,2}=e^2+\gamma _ie^{}+c_i+\rho _ie^{}+d_ie^2`$
gives the following first flow:
$`D_1d_i=d_i(c_ic_{i2}),`$
$`D_1\rho _i=\rho _i(c_ic_{i1})+d_{i+1}\gamma _id_i\gamma _{i2},`$
$`D_1\gamma _i=\rho _{i+2}\rho _i,`$
$`D_1c_i=d_{i+2}d_i+\gamma _i\rho _{i+1}+\gamma _{i1}\rho _i.`$
These are the 1D generalized fermionic Toda lattice equations which possess the $`N=4`$ supersymmetry.
## 6 Bi-Hamiltonian structure of 1D fermionic $`(K^+,K^{})`$-TL hierarchy
In this section, we apply the R-matrix approach to build the bi-Hamiltonian structure of the 1D fermionic $`(K^+,K^{})`$-TL hierarchy and perform its Dirac reduction.
The space of operators $`𝕆_{K^+}^+`$ (4) can obviously be split into the vector space direct sum, $`𝕆_{K^+}=(𝕆_{K^+})_+(𝕆_{K^+})_{}`$. The $`R`$-matrix arising from this splitting
$`R=P_+P_{},R(𝕆_{K^+})=(𝕆_{K^+})_+(𝕆_{K^+})_{}`$
obviously solves the graded modified Yang-Baxter equation (19) at $`\alpha =1`$. This $`R`$-matrix is not graded antisymmetric, $`RR^{}`$; however, its graded antisymmetric part $`A=1/2(RR^{})`$ as well as the $`R`$-matrix itself satisfy the graded modified Yang-Baxter equation (19). According to the general Theorem of Section 3 this means that there exist two Poisson structures on $`𝗀^{}=𝗀.`$ Substituting the general form of operators $`L_{K^+}^+`$ (24) and
$`u_{n,\xi }=e^{(nK^+)}(1)^i\delta _{i,\xi }`$
into (20) and (21) one can find the explicit form of the first and second Poisson brackets, respectively,
$`\{u_{n,i},u_{m,j}\}_1`$ $`=`$ $`(1)^j(\delta _{n,K^+}^{}+\delta _{m,K^+}^{}1)(u_{n+mK^+,i}\delta _{i,j+nK^+}`$ (39)
$`(1)^{(m+K^+)(n+K^++1)}u_{n+mK^+,j},\delta _{i,jm+K^+}`$
and
$`\{u_{n,i},u_{m,j}\}_2`$ $`=`$ $`(1)^j{\displaystyle \frac{1}{2}}[u_{n,i}u_{m,j}(\delta _{i,j+nK^+}(1)^m\delta _{i,jm+K^+})`$ (40)
$`+{\displaystyle \underset{k=0}{\overset{n+m}{}}}(\delta _{m,k}^+\delta _{m,k}^{})((1)^{mk}u_{n+mk,i}u_{k,j}\delta _{i,j+nk}`$
$`(1)^{m(n+k+1)}u_{k,i}u_{n+mk,j}\delta _{i,jm+k}\left)\right],`$
where
$$\delta _{n,m}^+=\{\begin{array}{ccc}\hfill 1,& \text{if}& n>m\hfill \\ \hfill 0,& \text{if}& nm\hfill \end{array},\delta _{n,m}^{}=\{\begin{array}{ccc}\hfill 1,& \text{if}& n<m\hfill \\ \hfill 0,& \text{if}& nm.\hfill \end{array}$$
Let us remind that the second Poisson brackets are defined for even values of $`K^+`$ only.
Our next goal is to perform the reduction of Poisson brackets (3940) for the functions parameterizing the operators $`L_{K^+}^+`$ (24) to the Poisson brackets corresponding to the reduced operators (37)
$`L_{K^+,K^{}}^{red}=e^{K^+}+{\displaystyle \underset{k=1}{\overset{K^++K^{}}{}}}u_{k,i}e^{(K^+k)},`$
where $`K^+,K^{}`$ are even numbers. Therefore, one needs to modify the Poisson brackets (3940) according to the reduction constraints
$`\begin{array}{ccc}\hfill u_{k,i}& =& 0,k>K^++K^{},\hfill \\ \hfill u_{0,i}& =& 1\hfill \end{array}\}\text{ for any }i.`$ (43)
We apply these reduction constraints in two steps. First, we note that for the first constraint in (43) the reduction simply amounts to imposing constraint $`u_{k,i}=0`$ $`(k>K^++K^{})`$ due to the observation that
$`\{u_{n,i},u_{m,j}\}_p|_{\stackrel{u_{k,i}=0,}{k>K^++K^{}}}=0,0nK^++K^{},m>K^++K^{},p=1,2.`$ (44)
For the first Poisson brackets (39) relation (44) is obvious. One can derive eq. (44) for the second Poisson brackets (40) if one divides the sum in (40) into three pieces
$`{\displaystyle \underset{k=0}{\overset{n+m}{}}}={\displaystyle \underset{k=0}{\overset{\text{max}(0,n+mK^+K^{}1)}{}}}+{\displaystyle \underset{k=\text{max}(1,n+mK^+K^{})}{\overset{\text{min}(n+m1,K^++K^{})}{}}}+{\displaystyle \underset{k=\text{min}(n+m,K^++K^{}+1)}{\overset{n+m}{}}}.`$ (45)
Now it is easy to verify that the second sum in the r.h.s. of eq. (45) is the only sum which could give a nonzero contribution to eq. (44), but it is equal to zero if $`0nK^++K^{},m>K^++K^{}.`$
Now let us consider the second reduction constraint in eq. (43), $`u_{0,i}=1`$. Following the standard Dirac reduction prescription we obtain for the Dirac brackets
$`\{u_{n,i},u_{m,j}\}_p^{red}=(\{u_{n,i},u_{m,j}\}_p)|_{u_{0,i}=1}\mathrm{}_p(u_{n,i},u_{m,j}),p=1,2`$
with the correction term
$`\mathrm{}_p(u_{n,i},u_{m,j})=({\displaystyle \underset{i^{^{}},j^{^{}}}{}}\{u_{n,i},u_{0,i^{^{}}}\}_p\{u_{0,i^{^{}}},u_{0,j^{^{}}}\}_p^1\{u_{0,j^{^{}}},u_{m,j}\}_p)|_{u_{0,i}=1}.`$
In the case of the first Poisson brackets $`\{u_{0,i},u_{m,j}\}_1=0`$ for any $`m`$. Thus, one can conclude that the first Poisson brackets (39) are not modified, and they can simply be restricted by imposing the constraints (43).
Before investigating the second Poisson brackets (40) we supply the fields $`u_{n,i}`$ and $`v_{n,i}`$ with the boundary conditions
$`\underset{i\pm \mathrm{}}{lim}u_{n,i}=\underset{i\pm \mathrm{}}{lim}v_{n,i}=0,n0`$ (46)
and introduce a new notation $`\delta _{i,jn}:=(\mathrm{\Lambda }^n)_{i,j}`$ which is useful in what follows. One can verify that in the new notation the multiplication of matrices results in adding powers of the operators $`\mathrm{\Lambda }`$:
$$\delta _{i,j^{^{}}+n}\delta _{j^{^{}},j+m}:=\mathrm{\Lambda }_{i,j^{^{}}}^n\mathrm{\Lambda }_{j^{^{}},j}^m=\mathrm{\Lambda }_{i,j}^{mn}.$$
Then, one can represent the correction term for the reduced second Poisson brackets as follows:
$`\mathrm{}_2(u_{n,i},u_{m,j})`$ $`=`$ $`1/2(1)^ju_{n,i}[(1\mathrm{\Lambda }^{K^+n})(1+\mathrm{\Lambda }^{K^+})`$
$`(\mathrm{\Lambda }^{K^+}\mathrm{\Lambda }^{K^+})^1(1+\mathrm{\Lambda }^{K^+})(1(1)^m\mathrm{\Lambda }^{mK^+})]_{i,j}u_{m,j}.`$
In general, the reduced second Poisson brackets are nonlocal since the inverse matrix
$$(\mathrm{\Lambda }^{2\nu }\mathrm{\Lambda }^{2\nu })^1=(1+\mathrm{\Lambda }^{2\nu })^1(1\mathrm{\Lambda }^\nu )^1(1+\mathrm{\Lambda }^\nu )^1,$$
being considered as an operator acting in the space of functionals with boundary conditions (46) can be expressed via infinite sums
$`(1\mathrm{\Lambda }^\nu )^1`$ $`=`$ $`\lambda _1{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\mathrm{\Lambda }^{k\nu }(1\lambda _1){\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\mathrm{\Lambda }^{k\nu },`$
$`(1+\mathrm{\Lambda }^\nu )^1`$ $`=`$ $`\lambda _2{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1)^k\mathrm{\Lambda }^{k\nu }+(1\lambda _2){\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(1)^k\mathrm{\Lambda }^{k\nu },`$
where $`\lambda _1`$ and $`\lambda _2`$ are arbitrary parameters. However, in the particular case $`K^+=2`$ the Dirac bracket becomes local, since the nonlocality is eliminated due to the contraction of the matrix with its inverse matrix. Indeed, one has
$`1\mathrm{\Lambda }^\nu `$ $`=`$ $`(1\mathrm{\Lambda }^1)(\delta _{\nu ,0}^{}{\displaystyle \underset{k=\nu +1}{\overset{0}{}}}\mathrm{\Lambda }^k\delta _{\nu ,0}^+{\displaystyle \underset{k=1}{\overset{\nu }{}}}\mathrm{\Lambda }^k),`$
$`1(1)^\nu \mathrm{\Lambda }^\nu `$ $`=`$ $`(1+\mathrm{\Lambda }^1)(\delta _{\nu ,0}^{}{\displaystyle \underset{k=\nu +1}{\overset{0}{}}}(1)^k\mathrm{\Lambda }^k\delta _{\nu ,0}^+{\displaystyle \underset{k=1}{\overset{\nu }{}}}(1)^k\mathrm{\Lambda }^k)`$
and for $`K^+=2`$ the Dirac brackets are local
$`\mathrm{}_2(u_{n,i},u_{m,j})=1/2(1)^ju_{n,i}[(1+\mathrm{\Lambda }^2)`$
$`(\delta _{n,2}^+{\displaystyle \underset{k=3n}{\overset{0}{}}}\mathrm{\Lambda }^k\delta _{n,2}^{}{\displaystyle \underset{k=1}{\overset{2n}{}}}\mathrm{\Lambda }^k)(\delta _{m,2}^{}{\displaystyle \underset{s=m1}{\overset{0}{}}}(1)^s\mathrm{\Lambda }^s\delta _{m,2}^+{\displaystyle \underset{s=1}{\overset{m2}{}}}(1)^s\mathrm{\Lambda }^s)]_{i,j}u_{m,j}.`$ (47)
As an example, we finish this section discussing the explicit form of the second Hamiltonian structure of the 1D $`(2,2)`$-Toda lattice hierarchy. The Lax operator defining this hierarchy is parameterized as follows:
$$L_{2,2}=u_ie^2+\gamma _ie^{}+c_i+\rho _ie^{}+d_ie^2.$$
Using eqs. (39), (40) and (6) one can derive the corresponding first Hamiltonian structure
$`\{d_i,c_j\}_1`$ $`=`$ $`(1)^jd_i(\delta _{i,j}\delta _{i,j+2}),`$
$`\{c_i,\rho _j\}_1`$ $`=`$ $`(1)^j\rho _j(\delta _{i,j1}+\delta _{i,j}),`$
$`\{\rho _i,\rho _j\}_1`$ $`=`$ $`(1)^j(d_i\delta _{i,j1}d_j\delta _{i,j+1}),`$
$`\{\gamma _i,\gamma _j\}_1`$ $`=`$ $`(1)^j(\delta _{i,j1}\delta _{i,j+1})`$ (48)
and the second Hamiltonian structure
$`\{d_i,d_j\}_2`$ $`=`$ $`1/2(1)^jd_id_j(1+\mathrm{\Delta })(\delta _{i,j2}\delta _{i,j+2}),`$
$`\{d_i,\rho _j\}_2`$ $`=`$ $`1/2(1)^jd_i\rho _j((1\mathrm{\Delta })(\delta _{i,j}+\delta _{i,j+1})(1+\mathrm{\Delta })(\delta _{i,j1}+\delta _{i,j+2})),`$
$`\{d_i,c_j\}_2`$ $`=`$ $`(1)^jd_ic_j(\delta _{i,j}\delta _{i,j+2}),`$
$`\{d_i,\gamma _j\}_2`$ $`=`$ $`1/2(1)^jd_i\gamma _j((1\mathrm{\Delta })(\delta _{i,j}+\delta _{i,j+3})(1+\mathrm{\Delta })(\delta _{i,j+1}+\delta _{i,j+2})),`$
$`\{d_i,u_j\}_2`$ $`=`$ $`1/2(1)^jd_iu_j(1\mathrm{\Delta })(\delta _{i,j}\delta _{i,j4}),`$
$`\{c_i,c_j\}_2`$ $`=`$ $`(1)^j(u_id_j\delta _{i,j2}u_jd_i\delta _{i,j+2}+\gamma _i\rho _j\delta _{i,j1}+\gamma _j\rho _i\delta _{i,j+1}),`$
$`\{c_i,\rho _j\}_2`$ $`=`$ $`(1)^j(d_i\gamma _j\delta _{i,j+1}+d_j\gamma _i\delta _{i,j2}c_i\rho _j(\delta _{i,j}+\delta _{i,j1})),`$
$`\{c_i,\gamma _j\}_2`$ $`=`$ $`(1)^j(u_i\rho _j\delta _{i,j1}+u_j\rho _i\delta _{i,j+2}),`$
$`\{\rho _i,\rho _j\}_2`$ $`=`$ $`(1)^j(c_id_j\delta _{i,j1}c_jd_i\delta _{i,j+1}1/2\rho _i\rho _j(1+\mathrm{\Delta })(\delta _{i,j+1}+\delta _{i,j1})),`$
$`\{\rho _i,\gamma _j\}_2`$ $`=`$ $`(1)^j(u_id_j\delta _{i,j1}u_jd_i\delta _{i,j+3}\rho _i\gamma _j(\delta _{i,j+1}1/2(1\mathrm{\Delta })(\delta _{i,j}+\delta _{i,j2}))),`$
$`\{\rho _i,u_j\}_2`$ $`=`$ $`1/2(1)^j\rho _iu_j(1\mathrm{\Delta })(\delta _{i,j}\delta _{i,j+1}+\delta _{i,j+2}\delta _{i,j+3}),`$
$`\{\gamma _i,\gamma _j\}_2`$ $`=`$ $`(1)^j(u_ic_j\delta _{i,j1}u_jc_i\delta _{i,j+1}1/2\gamma _i\gamma _j(1\mathrm{\Delta })(\delta _{i,j+1}+\delta _{i,j1})),`$
$`\{\gamma _i,u_j\}_2`$ $`=`$ $`1/2(1)^j\gamma _iu_j(1\mathrm{\Delta })(\delta _{i,j}\delta _{i,j+1}+\delta _{i,j+2}\delta _{i,j1})),`$
$`\{u_i,u_j\}_2`$ $`=`$ $`1/2(1)^ju_iu_j(1\mathrm{\Delta })(\delta _{i,j+2}\delta _{i,j2}),`$ (49)
where only nonzero brackets are written down. Here we have introduced the parameter $`\mathrm{\Delta }`$ which for the unreduced brackets is equal to zero, $`\mathrm{\Delta }=0`$, and for the Dirac reduced brackets with the reduction constraint $`u_i=1`$ is equal to one, $`\mathrm{\Delta }=1`$. In the latter case, algebras (6) and (6) reproduce, respectively, the first and second Hamiltonian structures of the 1D generalized fermionic Toda lattice hierarchy found in by a heuristic approach.
## 7 Bi-Hamiltonian structure of 2D fermionic $`(K^+,K^{})`$-TL hierarchy
In this section, we construct the bi-Hamiltonian structure of the 2D fermionic $`(K^+,K^{})`$-TL hierarchy. This hierarchy is associated with two Lax operators (24) belonging to the operator space (4). Following we consider the associative algebra on the space of the direct sum of two difference operators
$`𝗀:=𝕆_{K^+}^+𝕆_K^{}^{}.`$ (50)
However, in contrast to the case of pure bosonic 2D TL hierarchy, the difference operators in the direct sum (50) can be of both opposite and equal diagonal $`Z_2`$-parity. It turns out that the Poisson brackets can correctly be defined only for the latter case. In what follows we restrict ourselves to the case when both operators in g (50) have the same diagonal parity.
We denote $`(x^+,x^{})`$ elements of such algebra $`𝗀=𝗀^{}`$ with the product
$`(x_1^+,x_1^{})(x_2^+,x_2^{})=(x_1^+x_2^+,x_1^{}x_2^{}),`$ (51)
and define the inner product as follows:
$`<(x^+,x^{})>:=str(x^++x^{}),`$ (52)
where $`x^+𝕆_{K_1}^+`$, $`x^{}𝕆_{K_2}^{}`$. Using this definition we set the Poisson bracket to be
$`\{f_1,f_2\}=<(𝕆_{K^+}^+,𝕆_K^{}^{}),[f_1,f_2\}^{}>,`$ (53)
where
$$[f_1,f_2\}^{}:=([f_1^+,(f_2^+)^{(d_{f_1^+})}\},[f_1^{},(f_2^{})^{(d_{f_1^{}})}\}),$$
$`f_k`$ are functionals on g (50), and $`f_k[(𝕆_{K^+}^+,𝕆_K^{}^{})]=(f_k^+,f_k^{})`$ are their gradients which can be found from the definition
$`{\displaystyle \frac{f_k[(𝕆_{K^+}^+,𝕆_K^{}^{})+ϵ(\delta 𝕆_{K^+}^+,\delta 𝕆_K^{}^{})]}{ϵ}}|_{ϵ=0}`$ $`=`$ $`<(\delta 𝕆_{K^+}^+,\delta 𝕆_K^{}^{}),(f_k^+,f_k^{})>`$
$`=`$ $`<\delta 𝕆_{K^+}^+,f_k^+>+<\delta 𝕆_K^{}^{},f_k^{}>.`$
In order to obtain nontrivial Hamiltonian dynamics, one needs to modify the bracket (53) applying the $`R`$-matrix
$`[f_1,f_2\}^{}[f_1,f_2\}_R^{}=[R(f_1),f_2\}^{}+[f_1,R(f_2)\}^{}.`$
The $`R`$-matrix acts on the space (50) in the nontrivial way and mixes up the elements from two subalgebras in the direct sum with each other
$`R(x^+,x^{})=(x_+^+x_{}^++2x_{}^{},x_{}^{}x_+^{}+2x_+^+)`$ (54)
which is a crucial point of the $`R`$-matrix approach in the two-dimensional case . This $`R`$-matrix allows one to find two compatible Poisson structures and rewrite the Lax-pair representation (25) in the Hamiltonian form.
By construction the $`R`$-matrix (54) satisfies the graded modified Yang-Baxter equation
$`R([(x^+,x^{}),(y^+,y^{})\}_R)[R(x^+,x^{}),R(y^+,y^{})\}=\alpha [(x^+,x^{}),(y^+,y^{})\}`$ (55)
with $`\alpha =1`$. In order to show this, we simply repeat the arguments of representing the $`R`$-matrix as the difference $`R=\mathrm{\Pi }\overline{\mathrm{\Pi }}`$ of two projection operators
$$\mathrm{\Pi }(x^+,x^{})=(x_+^++x_{}^{},x_+^++x_{}^{}),\overline{\mathrm{\Pi }}(x^+,x^{})=(x_{}^+x_{}^{},x_+^{}x_+^+)$$
$$\mathrm{\Pi }(x^+,x^{})+\overline{\mathrm{\Pi }}(x^+,x^{})=(x^+,x^{}),\mathrm{\Pi }^2=\mathrm{\Pi },\overline{\mathrm{\Pi }}^2=\overline{\mathrm{\Pi }},\overline{\mathrm{\Pi }}\mathrm{\Pi }=\mathrm{\Pi }\overline{\mathrm{\Pi }}=0.$$
Therefore, the $`R`$-matrix (54) provides the splitting of the algebra g and solves the modified graded Yang-Baxter equation (55).
The two-dimensional $`R`$-matrix is not graded antisymmetric, its adjoint counterpart $`R^{}`$ looks like
$`R^{}(x^+,x^{})=(x_0^+x_{>0}^++2x_0^{},x_{>0}^{}x_0^{}+2x_{>0}^+)=\mathrm{\Pi }^{}\overline{\mathrm{\Pi }}^{},`$
where the dual projections are
$`\mathrm{\Pi }^{}(x^+,x^{})`$ $`=`$ $`(x_0^++x_0^{},x_{>0}^++x_{>0}^{}),`$
$`\overline{\mathrm{\Pi }}^{}(x^+,x^{})`$ $`=`$ $`(x_{>0}^+x_0^{},x_0^{}x_{>0}^+).`$
The direct verification by substitution in (55) shows that the graded antisymmetric part
$`1/2(R(x^+,x^{})R^{}(x^+,x^{}))=(x_{>0}^+x_{<0}^+x_0^{},x_{<0}^{}x_{>0}^{}+x_0^+)`$
also satisfies the graded modified Yang-Baxter equation (55). Therefore, by Theorem of Section 3 there exist two Poisson structures on g (50).
Using eqs. (2021), (5152), (54) and cyclic permutations inside the supertrace (9) we obtain the following general form of the first and second Poisson brackets:
$`\{f,g\}_i`$ $`=`$ $`<P_i^+(g^+,g^{}),(f^+)^{(d_g)}>`$ (56)
$`+`$ $`<P_i^{}(g^+,g^{}),(f^{})^{(d_g)}>,i=1,2,`$
where $`d_g:=d_{g^+}=d_g^{}`$. The Poisson tensors in eq. (56) are found for any values of $`(K^+,K^{})`$ for the first Hamiltonian structure
$`P_1^+(g^+,g^{})`$ $`=`$ $`[(g_{}^{}g_{}^+)^{(K^+)},(L_{K^+}^+)^{(d_g)}\}`$
$``$ $`([L_{K^+}^+,g^+\}+[L_K^{}^{},g^{}\})_0,`$
$`P_1^{}(g^+,g^{})`$ $`=`$ $`[(g_+^+g_+^{})^{(K^{})},(L_K^{}^{})^{(d_g)}\}`$
$``$ $`([L_{K^+}^+,g^+\}+[L_K^{}^{},g^{}\})_{>0},`$
while for the second Hamiltonian structure we constructed the explicit expression of the Poisson tensors for even values of $`(K^+,K^{})`$ only
$`P_2^+(g^+,g^{})`$ $`=`$ $`1/2([(g^{}(L_K^{}^{})^{(d_g)}+L_K^{}^{}g^{}`$
$``$ $`g^+(L_{K^+}^+)^{(d_g)}L_{K^+}^+g^+)_{},(L_{K^+}^+)^{(d_g)}\}`$
$``$ $`L_{K^+}^+([L_{K^+}^+,g^+\}+[L_K^{}^{},g^{}\})_0`$
$``$ $`([L_{K^+}^+,g^+\}+[L_K^{}^{},g^{}\})_0(L_{K^+}^+)^{(d_g)}),`$
$`P_2^{}(g^+,g^{})`$ $`=`$ $`1/2([(g^+(L_{K^+}^+)^{(d_g)}+L_{K^+}^+g^+`$
$``$ $`g^{}(L_K^{}^{})^{(d_g)}L_K^{}^{}g^{})_+,(L_K^{}^{})^{(d_g)}\}`$
$``$ $`L_K^{}^{}([L_{K^+}^+,g^+\}+[L_K^{}^{},g^{}\})_{>0}`$
$``$ $`([L_{K^+}^+,g^+\}+[L_K^{}^{},g^{}\})_{>0}(L_K^{}^{})^{(d_g)}).`$
The Poisson brackets for the functions $`u_{n,i}`$ and $`v_{n,i}`$ parameterizing the Lax operators (24) can explicitly be derived from (56) if one takes into account that
$`u_{n,\xi }`$ $``$ $`(u_{n,\xi }^+,u_{n,\xi }^{})=(e^{(nK^+)}(1)^i\delta _{i,\xi },0),`$
$`v_{n,\xi }`$ $``$ $`(v_{n,\xi }^+,v_{n,\xi }^{})=(0,e^{(K^{}n)}(1)^i\delta _{i,\xi }).`$
In such a way one can obtain the following expressions:
$`\{u_{n,i},u_{m,j}\}_1`$ $`=`$ $`(1)^j(\delta _{n,K^+}^{}+\delta _{m,K^+}^{}1)(u_{n+mK^+,i}\delta _{i,j+nK^+}`$
$`(1)^{(m+K^+)(n+K^++1)}u_{n+mK^+,j}\delta _{i,jm+K^+}),`$
$`\{u_{n,i},v_{m,j}\}_1`$ $`=`$ $`(1)^j[\delta _{m,K^{}}^+((1)^{(m+K^{})(n+K^++1)}u_{nm+K^{},j}\delta _{i,j+mK^{}}`$
$`u_{nm+K^{},i}\delta _{i,j+nK^+})+(\delta _{n,K^+}^{}1)(v_{mn+K^+,i}\delta _{i,j+nK^+}`$
$`(1)^{(m+K^{})(n+K^++1)}v_{mn+K^+,j}\delta _{i,j+mK^{}})],`$
$`\{v_{n,i},v_{m,j}\}_1`$ $`=`$ $`(1)^j(1\delta _{n,K^{}}^+\delta _{m,K^{}}^+)(v_{n+mK^{},i}\delta _{i,jn+K^{}}`$ (57)
$`(1)^{(m+K^{})(n+K^{}+1)}u_{n+mK^{},j}\delta _{i,j+mK^{}})`$
for the first Hamiltonian structure and
$`\{u_{n,i},u_{m,j}\}_2`$ $`=`$ $`(1)^j{\displaystyle \frac{1}{2}}[u_{n,i}u_{m,j}(\delta _{i,j+nK^+}(1)^m\delta _{i,jm+K^+})`$
$`+{\displaystyle \underset{k=0}{\overset{n+m}{}}}(\delta _{m,k}^+\delta _{m,k}^{})((1)^{mk}u_{n+mk,i}u_{k,j}\delta _{i,j+nk}`$
$`(1)^{m(n+k+1)}u_{k,i}u_{n+mk,j}\delta _{i,jm+k}\left)\right],`$
$`\{u_{n,i},v_{m,j}\}_2`$ $`=`$ $`(1)^j{\displaystyle \frac{1}{2}}\left[u_{n,i}v_{m,j}\right(\delta _{i,j}+\delta _{i,j+nK^+}`$
$`(1)^m(\delta _{i,j+mK^{}}+\delta _{i,j+n+mK^+K^{}}))`$
$`+2{\displaystyle \underset{k=\text{max}(0,mn)}{\overset{m1}{}}}(u_{nm+k,i}v_{k,j+mk}\delta _{i,j+nK^+}`$
$`(1)^{(n+1)m}v_{k,j}u_{nm+k,i+km}\delta _{i,j+m+K^{}}\left)\right],`$
$`\{v_{n,i},v_{m,j}\}_2`$ $`=`$ $`(1)^j{\displaystyle \frac{1}{2}}[v_{n,i}v_{m,j}(\delta _{i,jn+K^{}}(1)^m\delta _{i,j+mK^{}})`$ (58)
$`{\displaystyle \underset{k=0}{\overset{n+m}{}}}(\delta _{m,k}^+\delta _{m,k}^{})((1)^{mk}v_{n+mk,i}v_{k,j}\delta _{i,jn+k}`$
$`(1)^{m(n+k+1)}v_{k,i}v_{n+mk,j}\delta _{i,j+mk}\left)\right]`$
for the second Hamiltonian structure; the latter is valid for even values of $`(K^+,K^{})`$ only.
The reduction, according the reduction constraint $`u_{0,i}=1`$, does not require any correction terms for the first Hamiltonian structure (7) because $`\{u_{0,i},u_{n,j}\}=\{u_{0,i},v_{n,j}\}=0`$. For the second Hamiltonian structure the correction terms are
$`\mathrm{}_2(u_{n,i},u_{m,j})`$ $`=`$ $`1/2(1)^ju_{n,i}(1\mathrm{\Lambda }^{K^+n})(1+\mathrm{\Lambda }^{K^+})`$
$`(\mathrm{\Lambda }^{K^+}\mathrm{\Lambda }^{K^+})^1(1+\mathrm{\Lambda }^{K^+})(1(1)^m\mathrm{\Lambda }^{mK^+})u_{m,j},`$
$`\mathrm{}_2(u_{n,i},v_{m,j})`$ $`=`$ $`1/2(1)^ju_{n,i}(1\mathrm{\Lambda }^{K^+n})(1+\mathrm{\Lambda }^{K^+})`$
$`(\mathrm{\Lambda }^{K^+}\mathrm{\Lambda }^{K^+})^1(1+\mathrm{\Lambda }^{K^+})(1(1)^m\mathrm{\Lambda }^{K^{}m})v_{m,j},`$
$`\mathrm{}_2(v_{n,i},v_{m,j})`$ $`=`$ $`1/2(1)^jv_{n,i}(1\mathrm{\Lambda }^{nK^{}})(1+\mathrm{\Lambda }^{K^+})`$ (59)
$`(\mathrm{\Lambda }^{K^+}\mathrm{\Lambda }^{K^+})^1(1+\mathrm{\Lambda }^{K^+})(1(1)^m\mathrm{\Lambda }^{K^{}m})v_{m,j}`$
and they are nonlocal due to the presence of $`(\mathrm{\Lambda }^{K^+}\mathrm{\Lambda }^{K^+})^1`$ in the r.h.s. of eq. (7). However, there are unique values of $`(K^+,K^{})`$ when nonlocal terms are eliminated. Indeed, for $`K^+=K^{}=2`$ eqs. (7) become local
$`\mathrm{}_2(u_{n,i},u_{m,j})=1/2(1)^ju_{n,i}(1+\mathrm{\Lambda }^2)`$
$`(\delta _{n,2}^+{\displaystyle \underset{k=3n}{\overset{0}{}}}\mathrm{\Lambda }^k\delta _{n,2}^{}{\displaystyle \underset{k=1}{\overset{2n}{}}}\mathrm{\Lambda }^k)(\delta _{m,2}^{}{\displaystyle \underset{s=m1}{\overset{0}{}}}(1)^s\mathrm{\Lambda }^s\delta _{m,2}^+{\displaystyle \underset{s=1}{\overset{m2}{}}}(1)^s\mathrm{\Lambda }^s)u_{m,j},`$
$`\mathrm{}_2(u_{n,i},v_{m,j})=1/2(1)^ju_{n,i}(1+\mathrm{\Lambda }^2)`$
$`(\delta _{n,2}^+{\displaystyle \underset{k=3n}{\overset{0}{}}}\mathrm{\Lambda }^k\delta _{n,2}^{}{\displaystyle \underset{k=1}{\overset{2n}{}}}\mathrm{\Lambda }^k)(\delta _{m,2}^+{\displaystyle \underset{s=3m}{\overset{0}{}}}(1)^s\mathrm{\Lambda }^s\delta _{m,2}^{}{\displaystyle \underset{s=1}{\overset{2m}{}}}(1)^s\mathrm{\Lambda }^s)v_{m,j},`$
$`\mathrm{}_2(v_{n,i},v_{m,j})=1/2(1)^jv_{n,i}(1+\mathrm{\Lambda }^2)`$
$`(\delta _{n,2}^{}{\displaystyle \underset{k=n1}{\overset{0}{}}}\mathrm{\Lambda }^k\delta _{n,2}^+{\displaystyle \underset{k=1}{\overset{n2}{}}}\mathrm{\Lambda }^k)(\delta _{m,2}^+{\displaystyle \underset{s=3m}{\overset{0}{}}}(1)^s\mathrm{\Lambda }^s\delta _{m,2}^{}{\displaystyle \underset{s=1}{\overset{2m}{}}}(1)^s\mathrm{\Lambda }^s)v_{m,j}.`$
The Hamiltonian structures thus obtained possess the properties (123) with $`d_𝕆=d_{L_{K^+}^+}=d_{L_K^{}^{}}`$. Using them one can rewrite flows (2730) for even values of $`(K^+,K^{})`$ in the bi-Hamiltonian form
$`D_s^\pm (\begin{array}{c}u_{n,i}^{(r)}\\ v_{n,i}^{(r)}\end{array})=\{(\begin{array}{c}u_{n,i}^{(r)}\\ v_{n,i}^{(r)}\end{array}),H_{s+1}^\pm \}_1=\{(\begin{array}{c}u_{n,i}^{(r)}\\ v_{n,i}^{(r)}\end{array}),H_s^\pm \}_2,`$ (66)
with Hamiltonians
$`H_s^+`$ $`=`$ $`{\displaystyle \frac{1}{s}}str(L_{K^+}^+)_{}^s={\displaystyle \frac{1}{s}}{\displaystyle \underset{i=\mathrm{}}{\overset{\mathrm{}}{}}}(1)^iu_{sK^+,i}^{(s)},`$
$`H_s^{}`$ $`=`$ $`{\displaystyle \frac{1}{s}}str(L_K^{}^{})_{}^s={\displaystyle \frac{1}{s}}{\displaystyle \underset{i=\mathrm{}}{\overset{\mathrm{}}{}}}(1)^iv_{sK^{},i}^{(s)}.`$ (67)
For odd values of $`(K^+,K^{})`$ one can reproduce the bosonic flows of (2730) only. In this case eqs. (7), due to relation (17), give only fermionic nonzero Hamiltonians using which the bosonic flows can be generated via odd first Hamiltonian structure (7)
$`D_{2s}^\pm (\begin{array}{c}u_{n,i}^{(r)}\\ v_{n,i}^{(r)}\end{array})=\{(\begin{array}{c}u_{n,i}^{(r)}\\ v_{n,i}^{(r)}\end{array}),H_{2s+1}^\pm \}_1.`$ (72)
One remark is in order. In Sec. 6, we derived the bi-Hamiltonian structure for the 1D fermionic $`(K^+,K^{})`$-TL hierarchy in the $`R`$-matrix approach applying the $`R`$-matrix formalism developed in Sec. 3. However, the 1D fermionic $`(K^+,K^{})`$-TL hierarchy is obtained as a reduction of the 2D fermionic $`(K^+,K^{})`$-TL hierarchy with the reduction constraint (36). Therefore, this reduction constraint can be carried over into Hamiltonian structures and the bi-Hamiltonian structure for the 1D fermionic $`(K^+,K^{})`$-TL hierarchy can equivalently be derived from that of the 2D hierarchy just by reduction with the corresponding constraint. Actually, this reduction amounts to the extraction of subalgebras in the Hamiltonian structures for the 2D fermionic $`(K^+,K^{})`$-TL hierarchy. Indeed, one can verify that the fields $`u_{n,i}(0nK^+)`$ and $`v_{n,i}(0nK^{}1)`$ form subalgebras for even values of $`(K^+,K^{})`$ in both the first (7) and the second (7) Hamiltonian structures and these subalgebras are the first (39) and the second (40) Hamiltonian structures, respectively, if one redefines the fields as follows:
$`u_{n,i}`$ $`=`$ $`u_{n,i},0nK^+`$
$`u_{K^++K^{}n,i}`$ $`=`$ $`v_{n,i},0nK^{}1.`$
## 8 Conclusion
In this paper, we have generalized the $`R`$-matrix method to the case of $`Z_2`$-graded operators with an involution and found that there exist two Poisson bracket structures. The first Poisson bracket is defined for both odd and even operators with $`Z_2`$-grading while the second one is found for even operators only. It was shown that properties of the Poisson brackets were provided by the properties of the generalized graded bracket. We have deduced the operator form of the graded modified Yang-Baxter equation and demonstrated that for the class of graded antisymmetric $`R`$-matrices it was equivalent to the tensor form of the graded classical Yang-Baxter equation. Then we have proposed the Lax-pair representation in terms of the generalized graded bracket of the new 2D fermionic $`(K^+,K^{})`$-Toda lattice hierarchy and demonstrated that this hierarchy included all known up to now 2D supersymmetric TL equations as subsystems. Next we have considered the reduction of this hierarchy to the 1D space and reproduced the 1D generalized fermionic TL equations . Finally, we have applied the developed R-matrix formalism to derive the bi-Hamiltonian srtucture of the 1D and 2D fermionic $`(K^+,K^{})`$-TL hierarchies. For even values of $`(K^+,K^{})`$ both even first and second Hamiltonian structures were obtained and for this case all the flows of the 2D fermionic $`(K^+,K^{})`$-TL hierarchy can be rewritten in a bi-Hamiltonian form. For odd values of $`(K^+,K^{})`$ odd first Hamiltonian structure was found and for this case only bosonic flows of the 2D fermionic $`(K^+,K^{})`$-TL hierarchy can be represented in a Hamiltonian form using fermionic Hamiltonians.
Thus, the problem of Hamiltonian description of the fermionic flows of the 2D fermionic $`(K^+,K^{})`$-TL hierarchy is still open. Other problems yet to be answered are the construction of the second Hamiltonian structure (if any) for odd Lax operators and of the Hamiltonian structures (if any) for Lax operators $`L_K^+`$ and $`L_M^{}`$ of opposite $`Z_2`$-parities. All these questions are a subject for future investigations.
Acknowledgments. We would like to thank A.P. Isaev, P.P. Kulish, and A.A. Vladimirov for useful discussions. This work was partially supported by RFBR-DFG Grant No. 04-02-04002, DFG Grant 436 RUS 113/669-2, the NATO Grant PST.GLG.980302, and by the Heisenberg-Landau program.
Appendix A
Here we show that the graded modified Yang-Baxter equation (19) for the case of graded antisymmetric operators $`R`$ is equivalent to the tensor form of the graded classical Yang-Baxter equation introduced in the pioneer paper
$`[r_{12},r_{13}+r_{23}]+[r_{13},r_{23}]=0.`$ (A.1)
Let $`𝒢`$ be a superalgebra with the generators $`e_\mu `$ $`(\mu =1,\mathrm{},n+m)`$, the structure constants $`C_{\mu \nu }^\rho `$ and the graded Lie bracket
$`\{e_\mu ,e_\nu \}=e_\mu e_\nu (1)^{d_\mu d_\nu }e_\nu e_\mu =C_{\mu \nu }^\rho e_\rho ,`$
where $`d_\mu `$ is the Grassmann parity of the generator $`e_\mu `$ and $`d_\mu =0`$, if $`\mu =1,\mathrm{},n`$ and $`d_\mu =1`$, if $`\mu =n+1,\mathrm{},n+m`$.
For the generators of the algebra $`𝒢`$ the graded modified Yang-Baxter equation (19) at $`\alpha =1`$ takes the form
$`R(\{Re_\xi ,e_\gamma \})+R(\{e_\xi ,Re_\gamma \})\{Re_\xi ,Re_\gamma \}=\{e_\xi ,e_\gamma \}`$
which can equivalently be rewritten as follows:
$`(R_\xi ^\beta C_{\beta \gamma }^\alpha R_\alpha ^\mu +R_\gamma ^\beta C_{\xi \beta }^\alpha R_\alpha ^\mu R_\gamma ^\beta R_\xi ^\alpha C_{\alpha \beta }^\mu )e_\mu =C_{\xi \gamma }^\mu e_\mu .`$ (A.2)
Let us define an invariant supersymmetric non-degenerate bi-linear form associated with some representation of $`𝒢`$
$`<x,y>:=str(xy)\text{ for any }x,y𝒢`$
using which one can introduce the dual basis $`e^\mu `$ in superalgebra $`𝒢`$
$`<e^\mu ,e_\nu >=\delta _\nu ^\mu `$
and the supermetric
$`\eta _{\mu \nu }=<e_\mu ,e_\nu >,\eta ^{\mu \nu }=<e^\mu ,e^\nu >,\eta ^{\mu \nu }=(1)^{d_\mu }\eta ^{\nu \mu }=(\eta _{\nu \mu })^1`$
by which one can raise and lower indices as follows:
$`e_\alpha =\eta _{\alpha \beta }e^\beta ,e^\alpha =(1)^{d_\alpha }\eta ^{\alpha \beta }e_\beta =e_\beta \eta ^{\beta \alpha },\eta _{\alpha \nu }C_{\gamma \mu }^\nu \eta ^{\beta \mu }=C_{\alpha \gamma }^\beta .`$ (A.3)
Note that supermetric $`\eta _{\mu \nu }`$ and $`R_\nu ^\mu `$ are even matrices, i.e., for any their nonzero entry one has $`d_\mu +d_\nu =0`$.
Now we take the graded tensor product of both sides of eq. (A.2) with $`e^\xi e^\gamma `$ and using relations (A.3) rewrite it in the following form:
$`((1)^{d_\lambda }R^{\nu \beta }R^{\alpha \mu }C_{\alpha \beta }^\lambda +(1)^{d_\mu }R^{\lambda \beta }R^{\alpha \mu }C_{\beta \alpha }^\nu `$
$`(1)^{d_\mu }R^{\lambda \beta }R^{\nu \alpha }C_{\alpha \beta }^\mu )e_\mu e_\nu e_\lambda `$ $`=`$ $`\eta ^{\nu \xi }\eta ^{\lambda \gamma }C_{\xi \gamma }^\mu e_\mu e_\nu e_\lambda ,`$ (A.4)
where
$`R^{\alpha \beta }=(1)^{d_\alpha }\eta ^{\alpha \gamma }R_\gamma ^\beta ,R_\alpha ^\beta =\eta _{\alpha \gamma }R^{\gamma \beta }.`$
Here we use the graded tensor product
$`(e_\mu e_\nu )(e_\lambda e_\xi )=(1)^{d_\nu d_\lambda }(e_\mu e_\lambda e_\nu e_\xi ),`$ (A.5)
$`(e_\mu e_\nu )_{j_1j_2}^{i_1i_2}=(e_\mu )_{i_1j_1}(e_\nu )_{i_2j_2}(1)^{d_{i_2}(d_{i_1}+d_{j_1})},`$ (A.6)
where $`d_i(d_j)`$ means the Grassmann parity of the row (column) of the supermatrix element $`(e_\mu )_{ij}`$ and one has $`d_i+d_j=d_\mu `$ for any nonzero $`(e_\mu )_{ij}`$.
Assume further that the operator $`R`$ is graded antisymmetric, i.e.,
$`R^{\mu \nu }=(1)^{d_\mu }R^{\nu \mu }<Rx,y>=<x,Ry>`$
and define two $`(n+m)\times (n+m)`$ matrices
$`\stackrel{~}{r}=R^{\alpha \beta }e_\alpha e_\beta ,t=(1)^{d_\alpha }\eta ^{\alpha \beta }e_\alpha e_\beta ,`$
where $`\stackrel{~}{r}`$ is antisymmetric, $`\stackrel{~}{r}_{j_1j_2}^{i_1i_2}=\stackrel{~}{r}_{j_2j_1}^{i_2i_1}`$ and $`t`$ is the tensor Casimir element invariant with respect to the adjoint action
$`[t,x1+1x]=0\text{ for any }x𝒢.`$
Note that $`\stackrel{~}{r}`$ and $`t`$ are even matrices, i.e., $`d_{i_1}+d_{i_2}+d_{j_1}+d_{j_2}=0`$ for any $`\stackrel{~}{r}_{j_1j_2}^{i_1i_2}0`$ or $`t_{j_1j_2}^{i_1i_2}0`$.
Defining triple graded tensor products
$`\stackrel{~}{r}_{12}=R^{\alpha \beta }(e_\alpha e_\beta 1),\stackrel{~}{r}_{13}=R^{\alpha \beta }(e_\alpha 1e_\beta ),\stackrel{~}{r}_{23}=R^{\alpha \beta }(1e_\alpha e_\beta ),`$
$`t_{12}=\eta ^{\beta \alpha }(e_\alpha e_\beta 1),t_{13}=\eta ^{\beta \alpha }(e_\alpha 1e_\beta ),t_{23}=\eta ^{\beta \alpha }(1e_\alpha e_\beta )`$
and using eqs. (8) and (A.5) one can find that eq. (8) reads
$`[\stackrel{~}{r}_{12},\stackrel{~}{r}_{13}+\stackrel{~}{r}_{23}]+[\stackrel{~}{r}_{13},\stackrel{~}{r}_{23}]=[t_{12},t_{13}]`$ (A.7)
that is the tensor form of the graded modified classical Yang-Baxter equation (see e.g. ). In order to reproduce the graded classical Yang-Baxter equation with the zero in the r.h.s one needs to introduce a new matrix
$`r=\stackrel{~}{r}+t,\stackrel{~}{r}_{12}=1/2(r_{12}r_{21}),t_{12}=1/2(r_{12}+r_{21})`$
for which eq. (A.7) takes the form (A.1). Thus, for the case of graded antisymmetric operators $`R`$ we have established the equivalence of equations (19) at $`\alpha =1`$ and (A.1) which are, respectively, the operator form of the graded modified classical Yang-Baxter equation and the tensor form of the graded classical Yang-Baxter equation. Note that the former equation is a more general one, since it admits solutions which are not graded antisymmetric.
For completeness we give here the component form of eq. (A.1)
$`r_{kj_2}^{i_1i_2}r_{j_1j_3}^{ki_3}(1)^{d_{j_2}(d_{i_3}+d_{j_3})}r_{kj_3}^{i_1i_3}r_{j_1j_2}^{ki_2}(1)^{d_{i_2}(d_{i_3}+d_{j_3})}`$
$`+`$ $`r_{j_1k}^{i_1i_3}r_{j_2j_3}^{i_2k}(1)^{d_{i_2}(d_{i_1}+d_{j_1})}r_{j_2k}^{i_2i_3}r_{j_1j_3}^{i_1k}(1)^{d_{j_2}(d_{i_1}+d_{j_1})}`$
$`+`$ $`r_{j_1k}^{i_1i_2}r_{j_2j_3}^{ki_3}r_{kj_3}^{i_2i_3}r_{j_1j_2}^{i_1k}=0,`$
where eq. (A.6) is used when calculating the triple tensor products.
Appendix B
In this Appendix we investigate bi-linear bracket (21) and establish conditions on the $`R`$-matrix and its graded antisymmetric part which are necessary in order the bracket (21) be the Poisson bracket. Therefore, we need to verify the Jacobi identities for any $`f,g`$ and $`h`$ in g
$`(1)^{d_hd_g}\{h,\{f,g\}_2\}_2`$ $`+`$ $`\text{c. p.}=1/4(1)^{d_hd_g}<[𝕆,(\{f,g\}_2)\}`$ (B.1)
$`R\left((h)^{(d_g+d_f)}𝕆^{(d_f+d_g+d_h)}+𝕆^{(d_g+d_f)}(h)^{(d_g+d_f)}\right)`$
$`R((\{f,g\}_2)𝕆^{(d_g+d_f)}+𝕆(\{f,g\}_2))`$
$`[𝕆^{(d_g+d_f)},(h)^{(d_g+d_f)}\}>+\text{c. p.}=0,`$
where
$`(\{f,g\}_2)`$ $`=`$ $`1/4[gR((f)^{(d_g)}𝕆^{(d_f+d_g)}+𝕆^{(d_g)}(f)^{(d_g)})`$ (B.2)
$`+R(𝕆g+g𝕆^{(d_g)})(f)^{(d_g)})`$
$`+gR^{}((f)^{(d_g)}𝕆^{(d_f+d_g)}𝕆^{(d_g)}(f)^{(d_g)})`$
$`+R^{}(𝕆gg𝕆^{(d_g)})(f)^{(d_g)})`$
$`(1)^{d_fd_g}(R(f𝕆^{(d_f)}+𝕆f)(g)^{(d_f)}`$
$`+fR((g)^{(d_f)}𝕆^{(d_f+d_g)}+𝕆^{(d_f)}(g)^{(d_f)})`$
$`+fR^{}((g)^{(d_f)}𝕆^{(d_f+d_g)}𝕆^{(d_f)}(g)^{(d_f)})`$
$`+R^{}(f𝕆^{(d_f)}𝕆f)(g)^{(d_f)})].`$
Inserting (B.2) into (B.1) after tedious but straightforward calculations we get for the Jacobi identities
$`(1)^{d_f(d_h+d_g)}<[𝕆,h\}([R^{}F_{},R^{}G_{}\}+[RF_+,RG_+\}R([F_+,G_+\}_R))>+\text{c.p.}=0,`$
where we introduce the notation
$`F_\pm `$ $`=`$ $`𝕆^{(d_h)}(f)^{(d_h)}\pm (f)^{(d_h)}𝕆^{(d_h+d_f)}`$
$`G_\pm `$ $`=`$ $`𝕆^{(d_h+d_f)}(g)^{(d_h+d_f)}\pm (g)^{(d_h+d_f)}𝕆^{(d_h+d_f+d_g)}`$
for brevity. For any linear map $`R`$ and its graded antisymmetric part $`A=1/2(RR^{})`$ one has the identity
$`(1)^{d_f(d_h+d_g)}<[𝕆,h\}[R^{}F_{},R^{}G_{}\}>+\text{c.p.}`$
$`=(1)^{d_f(d_h+d_g)}<[𝕆,h\}(4/3([AF_{},AG_{}\}A([F_+,G_+\}_A)`$
$`([RF_{},RG_{}\}R([F_{},G_{}\}_R))>+\text{c.p.}`$ (B.3)
Now using (8) and the following identity
$`(1)^{d_f(d_h+d_g)}<[𝕆,h\}([F_{},G_{}\}+3[F_+,G_+\}))>+\text{c.p.}=0,`$
which can directly be verified, we finally rewrite the Jacobi identities with an arbitrary parameter $`\alpha `$ as follows:
$`(1)^{d_f(d_h+d_g)}<[𝕆,h\}([RF_+,RG_+\}R([F_+,G_+\}_R)+\alpha [F_+,G_+\}`$
$`([RF_{},RG_{}\}R([F_{},G_{}\}_R)+\alpha [F_{},G_{}\}`$
$`+4/3([AF_{},AG_{}\}A([F_{},G_{}\}_A)+\alpha [F_{},G_{}\}))>+\text{c.p.}=0.`$
Now it is obvious that the Jacobi identities are satisfied if $`R`$ and its graded antisymmetric part $`A`$ obey the graded modified Yang-Baxter equation (19) with the same $`\alpha `$.
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# The VIMOS VLT Deep Survey based on data obtained with the European Southern Observatory Very Large Telescope, Paranal, Chile, program 070.A-9007(A), and on data obtained at the Canada-France-Hawaii Telescope, operated by the CNRS of France, CNRC in Canada and the University of Hawaii
## 1 Introduction
An unbiased and detailed characterization of the luminosity function (LF) of field galaxies is a basic requirement in many extragalactic issues.
At present the local luminosity function is well constrained by the results obtained by the 2dF Galaxy Redshift Survey (2dFGRS, Norberg et al. norberg02 (2002)) and by the Sloan Digital Sky Survey (SDSS, Blanton et al. blanton03 (2003)). These surveys measure redshifts for $`10^510^6`$ galaxies over a large area, and therefore explore well the properties of the local ($`z<0.3`$) Universe. Such large numbers of objects also allow to study the luminosity functions (as well as the correlation functions and other properties) for galaxies of different types, defined on the basis of colours and/or spectral properties. A critical analysis of the luminosity functions depending on galaxy type as measured from the various redshift surveys, as well as a comparison of the different results, can be found in de Lapparent (delapparent03b (2003)).
Madgwick et al. (madgwick02 (2002)), analyzing 2dFGRS data, find a systematic steepening of the faint end slope and a faintening of $`M^{}`$ of the luminosity function as one moves from passive to active star forming galaxies. Similar results are found by Blanton et al. (blanton01 (2001)) for the SDSS sample, moving from the redder to the bluer galaxies.
For what concerns the high redshift Universe, several studies in the past ten years have aimed to map the evolution of the luminosity function. However, because of the long exposure times required to obtain spectra of high redshift galaxies, spectroscopic surveys were limited to a few $`10^2`$ objects. The Canadian Network for Observational Cosmology field galaxy redshift survey (CNOC-2, Lin et al. lin99 (1999)) and the ESO Sculptor Survey (ESS, de Lapparent et al. delapparent03a (2003)) derived the luminosity function up to $`z0.5`$ using $`2000`$ and $`600`$ redshifts, respectively. de Lapparent et al. (delapparent03a (2003)) find a behaviour of the LF by type similar to the local one derived from 2dFGRS and a strong evolution of a factor 2 in the volume density of the late type galaxies with respect to the early type galaxies. Lin et al. (lin99 (1999)) find for early type galaxies a positive luminosity evolution with increasing redshift, which is nearly compensated by a negative density evolution. On the contrary, for late type galaxies they find a strong positive density evolution, with nearly no luminosity evolution. At higher redshift, the Canada France Redshift Survey (CFRS, Lilly et al. lilly95 (1995)) allowed to study the luminosity function up to $`z1.1`$ with a sample of $`600`$ redshifts. From this survey, the LF of the red population shows small changes with redshift, while the LF of the blue population brightens by about one magnitude from $`z0.5`$ to $`z0.75`$. Other results suggest a strong number density evolution of early type galaxies (Bell et al. bell04 (2004), Faber et al. faber05 (2006)); conversely, in the K20 survey (Cimatti et al. k20 (2002)), Pozzetti et al. (pozzetti03 (2003)) found that red and early type galaxies dominate the bright end of the LF and that their number density shows at most a small decrease ($`<30\%`$) up to $`z1`$ (see also Saracco et al. saracco05 (2006) and Caputi et al. caputi06 (2006)).
Luminosity function estimates at higher redshift and/or with larger samples are up to now based only on photometric redshifts, like the COMBO-17 survey (Wolf et al. wolf03 (2003)) and the analysis of the FORS Deep Field (FDF, Gabasch et al. gabasch04 (2004)) and the Hubble Deep Fields (HDF-N and HDF-S, see e.g. Sawicki et al. sawicki97 (1997); Poli et al. poli01 (2001), poli03 (2003)); most of these projects derived also the luminosity function for different galaxy types. Wolf et al. (wolf03 (2003)) find that early type galaxies show a decrease of a factor $`10`$ in $`\varphi ^{}`$ up to $`z=1.2`$. Latest type galaxies show a brightening of about one magnitude in $`M^{}`$ and a $`\varphi ^{}`$ increase of a factor $`1.6`$ in their highest redshift bin ($`z1.1`$) in the blue band. Giallongo et al. (giallongo05 (2005)), using HDFs data, find that the B band number densities of red and blue galaxies have a different evolution, with a strong decrease of the red population at $`z=23`$ and a corresponding increase of the blue population. Dahlen et al. (dahlen05 (2005)), using GOODS data, claim that the starburst population fraction increases with redshift by a factor of 3 at $`z=2`$ in the U band.
Although photometric redshifts represent a powerful tool for deep surveys, their precision strongly relies on the number of used photometric bands, on the templates and on the adopted training procedure; moreover, they are affected by the problem of “catastrophic errors”, i.e. objects with a large difference between the spectroscopic and the photometric redshift.
A major improvement in this field is obtained with the VIMOS VLT Deep Survey (VVDS, Le Fèvre et al. 2003b ) and the DEEP-2 Galaxy Redshift Survey (Davis et al. davis03 (2003)). The VVDS is an ongoing program to map the evolution of galaxies, large scale structures and AGNs from the redshift measurements of $`10^5`$ objects down to a magnitude I$`{}_{AB}{}^{}=24`$, in combination with a multiwavelength dataset from radio to X-rays.
From the analysis of the evolution of the global luminosity function from the first epoch VVDS data (Ilbert et al. vvdsLF (2005)), we found a significant brightening of the $`M^{}`$ parameter in the U, B, V, R and I rest frame bands, going from $`z=0.05`$ to $`z=2`$. Moreover, we measured an increase of the comoving density of bright galaxies: this increase depends on the rest frame band, being higher in the bluest bands.
Among the other results of this survey, we recall the study of the radio selected objects (Bondi et al. radio1 (2003)) and of their optical counterparts (Ciliegi et al. radio2 (2005)), the evolution of the clustering properties (Le Fèvre et al. 2005b , Pollo et al. clus2 (2005)) and of the bias parameter (Marinoni et al. marinoni05 (2005)). Moreover, from the joined GALEX-VVDS sample, we derived the evolution of the far UV luminosity function (Arnouts et al. galexLF (2005)) and luminosity density (Schiminovich et al. galexLD (2005)).
In this paper we study the evolution of the luminosity functions of galaxies of different spectral types based on the VVDS data. This sample allows to perform this analysis for the first time with excellent statistical accuracy over a large redshift range ($`0.05<z<1.5`$).
The plan of the paper is the following: in sect. 2 we briefly present the first epoch VVDS sample, in sect. 3 we describe the galaxy classification and in sect. 4 we illustrate the method we used to estimate the luminosity functions. In sect. 5 we compare the luminosity functions of the different galaxy types and in sect. 6 we show the evolution with redshift of the luminosity functions by type. Finally in sect. 7 we compare our results with previous literature estimates and in sect. 8 we summarize our results.
Throughout the paper we adopt the cosmology $`\mathrm{\Omega }_m=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, with $`h=H_0/100`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, Magnitudes are given in the AB system and are expressed in the five standard bands U (Bessel), B and V (Johnson), R and I (Cousins).
## 2 The first epoch VVDS sample
The VVDS is described in detail in Le Fèvre et al. (2005a ): here we report only the main characteristics of the sample used for the analysis presented in this paper.
The entire VVDS is formed by a wide part on 4 fields (which is not used in this paper), and by a deep part, with spectroscopy in the range $`17.5`$ I$`{}_{AB}{}^{}24`$ on the field 0226-04. Multicolour photometry is available for each field (Le Fèvre et al. 2004a ): in particular, the B, V, R, I photometry for the 0226-04 deep field is described in detail in McCracken et al. (photom2 (2003)). Moreover, U band (Radovich et al. photomU (2004)) and J and K band (Iovino et al. photomK (2005)) data are available for smaller areas of these fields.
Starting from these photometric catalogues, spectroscopic observations were performed with the VIsible Multi–Object Spectrograph (VIMOS, Le Fèvre et al. 2003a ) mounted on the ESO Very Large Telescope (UT3). The selection of objects for spectroscopic observations was based only on magnitude, without any other colour or shape criteria.
Deep spectroscopic observations ($`17.5`$ I$`{}_{AB}{}^{}24`$) were performed also on the Chandra Deep Field South (VVDS-CDFS, Le Fèvre et al. 2004b ), starting from the EIS I band photometry and astrometry (Arnouts et al. eis (2001)). Multicolour U, B, V, R and I photometry for this sample is available from the COMBO-17 survey (Wolf et al. wolf03 (2003)).
Spectroscopic data were reduced with the VIMOS Interactive Pipeline Graphical Interface (VIPGI, Scodeggio et al. vipgi (2005), Zanichelli et al. ifu (2005)) and redshift measurements were performed with the KBRED package (Scaramella et al. kbred (2006)) and then visually checked. Each redshift measurement was assigned a quality flag, ranging from 0 (failed measurement) to 4 (100% confidence level); flag 9 indicates spectra with a single emission line, for which multiple solutions are possible. Further details on the quality flags are given in Le Fèvre et al. (2005a ).
The analysis presented in this paper is based on the first epoch VVDS deep sample, which has been obtained from the first observations (fall 2002) on the fields VVDS-02h and VVDS-CDFS, which cover 1750 and 450 arcmin<sup>2</sup>, respectively. We eliminated from the sample spectroscopically confirmed stars and broad line AGNs, remaining with 6477 + 1236 galaxy spectra with secure spectroscopic identification (flag 2, 3, 4, 9), corresponding to a confidence level higher than 75%. Redshifts with flags 0 and 1 are taken into account statistically (see sect. 4). This spectroscopic sample, which is purely magnitude selected, has a median redshift of $`0.76`$.
## 3 Galaxy classification
Galaxies have been classified using all the multicolour information available; in the VVDS-02h field B, V, R and I band magnitudes are available for all galaxies, while U band data are available for 83% of the galaxies. For the VVDS-CDFS sample U, B, V, R and I photometry from the COMBO-17 survey is used.
Absolute magnitudes are computed following the method described in the Appendix of Ilbert et al. (vvdsLF (2005)). The K-correction is computed using a set of templates and all the photometric information (UBVRI) available. However, in order to reduce the template dependency, the rest frame absolute magnitude in each band is derived using the apparent magnitude from the closest observed band, redshifted at the redshift of the galaxy. With this method, the applied K-correction is as small as possible as possible.
For each galaxy the rest frame magnitudes were matched with the empirical set of SEDs described in Arnouts et al. (arnouts99 (1999)), composed of four observed spectra (CWW, Coleman et al. cww (1980)) and two starburst SEDs computed with GISSEL (Bruzual & Charlot bc93 (1993)). The match is performed minimizing a $`\chi ^2`$ variable on these templates at the spectroscopic redshift of each galaxy. The same procedure has been applied by Lin et al. (lin99 (1999)) to the CNOC-2 survey up to $`z0.55`$. This approach is also similar to that adopted by Wolf et al. (wolf03 (2003)) for the COMBO-17 survey, but we have the advantage of using spectroscopic redshifts, while they had to rely on photometric redshifts.
Galaxies have been divided in four types, corresponding to the E/S0 template (type 1), early spiral template (type 2), late spiral template (type 3) and irregular template (type 4). These types are based on the four CWW templates: type 4 includes also the two starburst templates. The numbers of galaxies for each type are listed in Table 1.
In order to have an idea of the correspondence of these types with colours, we report here the rest frame colours for each template: type 1, 2, 3 and 4 have $`B_{AB}I_{AB}=`$1.58, 1.11, 0.79 and 0.57, respectively. Given these colours, a rough colour subdivision for each class is $`1.3<B_{AB}I_{AB}`$ for type 1, $`0.95<B_{AB}I_{AB}<1.3`$ for type 2, $`0.68<B_{AB}I_{AB}<0.95`$ for type 3 and $`B_{AB}I_{AB}<0.68`$ for type 4. However, we remind that these colour ranges are only indicative and our classification scheme is based on the whole multicolour coverage. In Fig.1 we show the $`UB`$ and $`BI`$ colour distributions for the galaxies of our sample divided according to type. From this figure it is clear that, although the different types have different colour distributions, they present significant overlaps. This fact is a consequence of classification schemes using template fitting on multicolour data.
Note that, in order to avoid to be model dependent, we did not apply to the templates any correction aimed at taking into account colour evolution with redshift. It is well known that the colour of a simple stellar population subject to passive evolution was bluer in the past. In principle, this could imply that galaxies classified as type 1 at low redshift might be classified differently at higher redshift. Indeed, this effect has been invoked by the authors who found negative evolution in the luminosity function of “red” galaxies (see f.i. Wolf et al. wolf03 (2003)). In order to verify this hypothesis, we applied our classification scheme to synthetic spectra (Bruzual & Charlot bc93 (1993)) of ellipticals (i.e. simple stellar populations and exponentially declining star formation with time scales of 0.1 Gyr and 0.3 Gyr) with formation redshift between $`z_{form}=`$2 and 20. We find that all ellipticals with $`z_{form}>2`$ would be classified as type 1 objects, even at $`z1`$.
In order to check, at least on a statistical basis, the consistency between this photometric classification and average spectral properties, we summed the normalized spectra of all galaxies in each of the four types. The resulting average spectra are shown in Fig.2 for each type in various redshift bins. This figure confirms the robustness of our classification scheme: moving from type 1 to type 4 objects, the composite spectra show an increasingly blue continuum, with emission lines of increasing strength. This confirms that the four types show different spectral features and therefore represent different classes of objects.
For the VVDS-CDFS sample, HST-ACS images are available. Using these data, Lauger et al. (lauger06 (2006)) classified the galaxies in this sample using an asymmetry-concentration diagram. Plotting our type 1 galaxies in this diagram, we find that $`91\%`$ of them lie in the region of bulge dominated objects, showing an excellent consistency also between our photometric classification and a morphological one.
In Fig.3 we plot the observed fraction of bright galaxies of each type as a function of redshift. We selected objects with $`M_{B_{AB}}5log(h)<20`$ because these galaxies are visible in the whole redshift range. From this figure it is clear the growing importance of bright late type objects with increasing redshift and the corresponding strong decrease of the fraction of bright early type galaxies.
## 4 Luminosity function estimate
Luminosity functions were derived using the Algorithm for Luminosity Function (ALF), a dedicated tool which uses various estimators: the non-parametric $`1/V_{max}`$ (Schmidt schmidt68 (1968)), $`C^+`$ (Lynden Bell lyndenbell71 (1971)), $`SWML`$ (Efstathiou et al. swml (1988)) and the parametric $`STY`$ (Sandage, Tammann & Yahil sty (1979)), for which we assumed a Schechter function (Schechter schechter76 (1976)). The tool and these estimators, as well as their specific use in the context of the VVDS, are described in detail in Ilbert et al. (vvdsLF (2005)).
Ilbert et al. (ilbert04 (2004)) have shown that the estimate of the global luminosity function can be biased, mainly in its faint end, when the band in which it is measured is far from the rest frame band in which galaxies are selected. This is due to the fact that, because of the K-corrections, different galaxy types are visible in different absolute magnitude ranges at a given redshift. When computing the global luminosity functions (Ilbert et al. vvdsLF (2005)), we avoided this bias by using for the $`STY`$ estimate, in each redshift range, only galaxies within the absolute magnitude range where all the SEDs are potentially observable.
Even if this bias is much less important when estimating the luminosity function of galaxies divided by types, we have, however, taken it into account. The absolute magnitude limits for the $`STY`$ estimate are indicated with vertical dashed lines in the figures, and in the tables where the best fit parameters are reported (Table 2 and 3) we give both the total number of objects and the number of galaxies within this magnitude limit.
In order to take into account the unknown redshifts (not observed objects and failed spectra), a weight was applied to each galaxy, following the procedure described in detail in Ilbert et al. (vvdsLF (2005)).
This weight is a combination of two different contributions: the target sampling rate and the spectroscopic success rate. The target sampling rate, i.e. the fraction of observed galaxies, corrects for the selection effects due to the procedure used for the mask preparation (Bottini et al. vmmps (2005)): to maximize the number of slits, the procedure tends to select objects with smaller angular size on the x-axis of the image, corresponding to the direction in which the slits are placed. As a consequence, the final spectroscopic sample has a bias against large objects, which produces a mild dependence of the target sampling rate on the apparent magnitude. The target sampling rate is $`25\%`$ for most of the sample and is computed as a function of the object size (see Ilbert et al. vvdsLF (2005) for further details). The spectroscopic success rate takes into account the fraction of objects without a good redshift determination (i.e. flags 0 and 1). As shown in Ilbert et al. (vvdsLF (2005)), these objects are expected to have a different redshift distribution with respect to that of the sample with measured redshift, as confirmed by the use of their photometric redshifts. Given the fact that the spectroscopic success rate decreases for faint apparent magnitudes, we derived it in four magnitude bins as a function of redshift (using photometric redshifts, see Fig.3 in Ilbert et al. vvdsLF (2005)). The shape of the spectroscopic success rate is similar in all magnitude bins, showing a maximum at $`z0.7`$ and two minima for $`z<0.5`$ and $`z>1.5`$. Since the number of galaxies for each type is not large enough to reliably estimate the spectroscopic success rate as a function also of the galaxy type, we have used the global spectroscopic success rate for all galaxy types.
In order to check the effect of the “cosmic variance”, i.e. variations in the luminosity function due to fluctuations in the large scale structure, we applied the following test on the VVDS-02h deep area. For this field we derived photometric redshifts (Ilbert et al. zphot (2006)) based on both VVDS photometry (BVRIJK) and on new CFHT Legacy Survey photometry (ugriz), which has now become available in a field covering 1 sq. deg. (http://www.cfht.hawaii.edu/Science/CFHLS/) which includes the 1700 arcmin<sup>2</sup> area covered by the VVDS spectroscopic survey. Then we divided the field in two non overlapping regions (the sub-area where spectroscopic data are available and the remaining area) and we compared the luminosity distributions of galaxies in these two samples ($``$0.5 sq.deg. each), in the same redshift bins in which the luminosity functions were derived. In each redshift bin the two distributions show average differences of the order of 10%, with some larger fluctuations due to Poisson statistics, without any systematic trend. Therefore the influence of the “cosmic variance” is expected to be limited.
## 5 Comparison of the luminosity functions of different types
As a first step, we compare the luminosity functions for galaxies of different types, in order to see which is the relative behaviour of the various populations. To perform this comparison, we selected galaxies in the redshift range $`[0.40.9]`$. About $`50\%`$ of the objects of our sample are included in this redshift interval, covering a wide range of luminosities (i.e. absolute magnitudes in the B band are in the range $`[23.7;16.8]5`$ log$`(h)`$). Moreover, the spectroscopic success rate of our survey reaches a maximum at $`z0.7`$ and therefore the possible dependency of the estimated LF on the weighting scheme described above is minimized in this redshift interval.
In Fig.4 we report the luminosity functions estimated with the $`STY`$ method in the rest frame bands U, B, R and I, with the corresponding confidence ellipses for the $`\alpha `$ and $`M^{}`$ parameters. The values of the parameters with their $`1\sigma `$ errors, as well as the number of galaxies for each type, are reported in Table 2. The $`\varphi ^{}`$ parameters listed in this table are derived adopting the density estimator of Efstathiou et al. (swml (1988)), following the procedure described in the Appendix of Ilbert et al. (vvdsLF (2005)).
Table 2 shows that our estimates are based on several hundreds of galaxies for each type, and are therefore well statistically constrained, as can be seen also from the sizes of the confidence ellipses in the figures. The first, very clear result which appears from Fig.4 is the significant strong steepening of the luminosity functions going from early to late types. In all bands the power law slope steepens by $`\mathrm{\Delta }\alpha 1.31.5`$ going from type 1 to type 4 galaxies and galaxies of late types are the dominant population at faint magnitudes.
Systematic trends are also seen in the $`M^{}`$ parameter. In the reddest bands (lower panels in Fig.4), $`M^{}`$ is significantly fainter for late type galaxies and this faintening is particularly apparent for type 4 objects. The brighter $`M^{}`$ for early type galaxies reflects the fact that most of the more massive objects belong to this population.
The difference of the $`M^{}`$ values for different types decreases in the B band and disappears or even changes sign in the U band. This behaviour is explained by the fact that the luminosity in the bluer bands is dominated by the light of young stars, produced during the star formation activity. Galaxies of later types, which are still actively forming stars, are therefore more luminous in the bluer bands.
These results are qualitatively in agreement with previous results from the literature, most of which at lower redshift (see de Lapparent delapparent03b (2003) for a review of the results from a number of surveys in the redshift range $`0z0.6`$). In particular, in almost all surveys the luminosity function of late type galaxies is steeper and with a fainter $`M^{}`$ with respect to that of early type galaxies. However, a quantitative comparison with previous results is difficult, because of the different classification schemes adopted in the various surveys, the different redshift ranges and selection criteria.
## 6 Evolution with redshift of the luminosity functions by type
We derived luminosity functions for each type in redshift bins in the U, B, V, R and I rest frame bands. Given our multicolour coverage and the explored redshift range, the estimate of the absolute magnitudes in the U and B rest frame bands are those which require less extrapolations (see Appendix A and Figure A.1 in Ilbert et al. vvdsLF (2005)). Therefore, to limit the number of figures in the paper, we show the results in the B rest frame band.
Figures 5, 6, 7 and 8 show the luminosity function for type 1, 2, 3 and 4 galaxies in redshift bins, obtained with $`C^+`$ and $`STY`$ methods. The luminosity functions derived with the other two methods ($`1/V_{max}`$ and $`SWML`$) are consistent with those shown in the figures, but are not drawn for clarity. The dotted line in each panel represents the fit derived in the redshift range $`[0.40.9]`$ (see previous section), while the dashed line is the estimate derived by fixing the slope $`\alpha `$ to the value obtained in the range $`[0.40.9]`$.
In Table 3 we report the Schechter parameters, with their $`1\sigma `$ errors, estimated for the various redshift bins, from $`z=0.2`$ to $`z=1.5`$ for each type; as a reference, in the last line we give the parameters derived in the redshift bin $`[0.40.9]`$.
We do not show the results for bins where the number of objects is too small (less than $`30`$) to constrain the parameters of the luminosity function (these are the bin $`[0.050.2]`$ for type 1 and 2 and the bin $`[1.21.5]`$ for type 1). Given the bright magnitude limit of the survey (I$`{}_{AB}{}^{}17.5`$) and the small sampled volume, bright galaxies are not sampled in the redshift bin $`[0.050.2]`$ and therefore we can not constrain the $`M^{}`$ parameter even for type 3 and 4 galaxies, where the number of objects is relatively high ($`80`$). For this reason we show in the figures the luminosity function estimates in this bin, but we do not report the $`STY`$ parameters for this redshift range in Table 3.
In Fig.9 we show the confidence ellipses of the parameters $`\alpha `$ and $`M^{}`$ in different redshift bins for the different types. From this figure it is possible to see that, within each type, the estimated slopes $`\alpha `$ in the various redshift bins are always consistent (within 90% confidence level) with each other and with the value derived in the redshift range $`[0.40.9]`$. Therefore, there is no evidence of a significant change with redshift of the luminosity function slope within each galaxy type. Note also that the uncertainties on the slope estimates become quite large for $`z>1`$: this is due to the fact that, even with the faint limit (I$`{}_{AB}{}^{}24`$) of this survey, the number of galaxies fainter than $`M^{}`$ is too low to well constrain the slope.
In each panel of Figures 5, 6, 7 and 8 we draw, as a reference, the luminosity function derived in the redshift bin $`[0.40.9]`$ (dotted line). Comparing this curve with the estimates of the luminosity function in the different redshift bins an evolution can be seen, strongly depending on the galaxy type. This evolution is particularly evident for type 4 galaxies: going from low to high redshift there is an almost continuous brightening of $`M^{}`$ and at fixed luminosity the density of these galaxies was much higher in the past.
The observed evolution of the luminosity function could be due to an evolution in luminosity and/or in density or both. One of the advantages of the $`STY`$ method is that of allowing to derive the $`\alpha `$ and $`M^{}`$ parameters independently from $`\varphi ^{}`$, which is not possible when one directly fits a Schechter function on the $`1/V_{max}`$ points.
Given the fact that we found that $`\alpha `$ is consistent with being constant for each type, we can fix it at the reference value derived in the redshift range $`[0.40.9]`$ and then study the variations of the parameters $`M^{}`$ and $`\varphi ^{}`$ as a function of the redshift (see upper and middle panels of Fig.10). These estimates are reported in Table 3.
From Fig.10 we can see a mild evolution of $`M^{}`$ from the lowest to the highest redshift bin for each type. In particular, this brightening ranges from $`\stackrel{<}{}0.5`$ mag for early type galaxies to $`1`$ mag for the latest type galaxies. The only exception with respect to the general trend is that of type 3 objects in the bin $`[0.20.4]`$, for which the best fit value of the $`M^{}`$ parameter is significantly fainter than expexted. The reason for this discontinuity in $`M^{}`$ for type 3 galaxies at low redshift is not clear.
On the contrary, the $`\varphi ^{}`$ parameter shows a very different behaviour for type 1, 2 and 3 galaxies with respect to type 4 galaxies. The first three types show a rapid decrease of $`\varphi ^{}`$ at low redshifts (between $`z0.3`$ and $`z0.5`$), then $`\varphi ^{}`$ remains roughly constant up to $`z0.9`$ and finally slowly decreases up to $`z=1.5`$.
Type 4 objects, on the contrary, show an increase in $`\varphi ^{}`$ at low redshift, then $`\varphi ^{}`$ is nearly constant up to $`z0.8`$ and shows a rapid increase of a factor $`2`$ at $`z=1.1`$. Then there seems to be a decrease from $`z=1.1`$ to $`z=1.3`$. However, this decrease can likely be a spurious effect, due to the fact that in this bin the estimated $`M^{}`$ is very close to the bias limit (see Sect. 4). If, for example, we fix $`M^{}`$ to the value obtained in the previous redshift bin $`[1.01.2]`$, we derive a significantly higher value for $`\varphi ^{}`$, i.e. $`\varphi ^{}=5.83\times 10^3h^3Mpc^3`$ (corresponding to $`\varphi ^{}/\varphi _{ref}^{}=1.31`$). Therefore the last $`\varphi ^{}`$ value for type 4 galaxies is likely to be a lower limit of the true density value.
This analysis of the different trends with redshift of the $`\varphi ^{}`$ parameter indicates that the importance of type 4 galaxies is increasing with redshift. However, since $`M^{}`$ is changing with redshift (see above), the trends of $`\varphi ^{}`$ can not be immediately interpreted in terms of density at a given absolute magnitude. For this reason, we have computed the density of bright galaxies as a function of redshift. We integrated the best fit luminosity function down to $`M_{B_{AB}}5log(h)<20`$. This limit approximately corresponds to the faintest galaxies which are visible in the whole redshift range. In the lowest panel of Fig.10 we plot the density of bright galaxies of each type as a function of redshift.
The main results shown in this plot can be summarized in the following way:
a. the density of bright early type galaxies (type 1) decreases with increasing redshift, however this decrease is rather modest, being of the order of $`40\%`$ from $`z0.3`$ to $`z1.1`$;
b. the density of bright late type galaxies (type 4) is instead significantly increasing, by a factor $`6.6`$ from $`z0.3`$ to $`z1.3`$.
The behaviour of type 4 galaxies is also responsible of the evolution of the global luminosity function measured by Ilbert et al. (vvdsLF (2005)). In fact, the increasing number of both faint and bright type 4 galaxies leads to the steepening of the global LF (due to the very steep slope of type 4 LF) and to the brightening of $`M^{}`$ (due to the increasing fraction of bright blue objects). This fact has been directly checked summing the LF of all types and comparing the result with the global LF estimate.
## 7 Comparison with previous literature results
Although various estimates of the luminosity function by galaxy type are available in the literature, a quantitative comparison of our results with previous analyses is not straightforward, because of the different classification schemes adopted in the various surveys, the different numbers of galaxy types, the different redshift ranges and selection criteria.
The strong density evolution of late type galaxies we find in the redshift range $`[0.21.5]`$ extends to significantly higher redshift the results found by de Lapparent et al. (delapparent04 (2004)) for their latest type at $`z0.5`$.
Among the other surveys, CNOC-2 (Lin et al. lin99 (1999)) and COMBO-17 (Wolf et al. wolf03 (2003)) adopted a classification scheme somewhat similar to ours. Lin et al. (lin99 (1999)) divided their sample of galaxies with $`z<0.55`$ in three classes (early, intermediate and late type) using CWW templates. For early type galaxies they found a positive luminosity evolution, which is nearly compensated by a negative density evolution. On the contrary, for late type galaxies they found a strong positive density evolution, with nearly no luminosity evolution. We see at higher redshift the same trend in density, but our evolution in $`M^{}`$ is contained within one magnitude for each type. Wolf et al. (wolf03 (2003)) used a sample of $`25000`$ galaxies with photometric redshifts, applying a classification scheme in four classes similar to ours but using the Kinney et al. (kinney96 (1996)) templates instead of the CWW templates.
In Figure 11 we compare our results (solid lines) with those from COMBO-17 (dashed lines). Note that, since our data extend always to fainter absolute magnitudes, in particular for type 1 galaxies, our estimates of the faint end slope are likely to be better determined: these estimates are derived in each redshift bin, while the COMBO-17 slopes are fixed to the value determined in the redshift range \[0.2 - 0.4\]. This figure shows that there are significant differences in both shapes and evolution of the LF estimates. The slope of the COMBO-17 LF is flatter than ours for type 1 galaxies, while it is steeper for types 2 and 3 (at least up to $`z=1.0`$). The most significant difference between the two surveys regards the evolution of type 1 galaxies, for which we do not have any evidence of the very strong density decrease with increasing redshift present in COMBO-17 data. The reason for this difference is unclear: it could be due to the use of different templates in the definition of the galaxy types or to a degeneracy between photometric redshift and classification, which might affect the COMBO-17 data. Bell et al. (bell04 (2004)) explained the strong negative evolution of type 1 galaxies as a consequence of the blueing with increasing redshift of elliptical galaxies with respect to the template used for classification. In this way, at increasing redshift an increasing number of “ellipticals” would be assigned a later type, therefore producing the detected density decrease observed in COMBO-17 data. However, as already mentioned in sect.3, we verified that this effect does not affect our classification scheme, at least up to $`z1`$ and for simple stellar populations with $`z_{form}>2`$.
As a test of this possible effect, we compared the LF obtained by adding together type 1 and 2 objects. In this way we can check whether the differences between VVDS and COMBO-17 LF of type 1 galaxies is due to the fact that a significant fraction of high redshift type 1 galaxies in COMBO-17 are classified as type 2 galaxies. This comparison is shown in Fig.12. The discrepancy between VVDS and COMBO-17 LF is now reduced, but there are still significant differences both in slope (being the COMBO-17 LF steeper than that of VVDS) and in normalization, especially in the highest redshift bin, where the VVDS LF is more than a factor 2 higher than the COMBO-17 LF.
## 8 Conclusions
In this paper we studied the evolution of the luminosity function of different galaxy types up to $`z=1.5`$, using 7713 spectra with $`17.5`$ I$`{}_{AB}{}^{}24`$ from the first epoch VVDS deep sample.
The VVDS data allow for the first time to study with excellent statistical accuracy the evolution of the luminosity functions by galaxy type from relatively low redshift up to $`z=1.5`$ from a single purely magnitude selected spectroscopic sample. The faint limiting magnitude of the VVDS sample allows to measure the slope of the faint end of the luminosity function with unprecedented accuracy up to $`z1.2`$. The use of spectroscopic redshifts implies low “catastrophic” failure rate compared to the photometric redshifts and therefore rare populations are sampled with a better accuracy, like in the bright end on the luminosity function. Moreover, the use of spectroscopic redshifts allow to classify galaxies avoiding a possible degeneracy between photometric redshift and classification.
VVDS galaxies were classified in four spectral classes using their colours and redshift, from early type to irregular galaxies, and luminosity functions were derived for each type in redshift bins, from $`z=0.05`$ to $`z=1.5`$, in the U, B, V, R and I rest frame bands.
We find a significant strong steepening of the luminosity function going from early to late types: in all bands the power law slope steepens by $`\mathrm{\Delta }\alpha 1.31.5`$ going from type 1 to type 4. Moreover, the $`M^{}`$ parameter of the Schechter function is significantly fainter for late type galaxies. As expected, this difference increases in the redder bands, reaching $`1.4`$ mag in the I band.
Studying the variations with redshift of the luminosity function for each type, we find that there is no evidence of a significant change of the slope, while we find a brightening of $`M^{}`$ with increasing redshift, ranging from $`\stackrel{<}{}0.5`$ mag for early type galaxies to $`1`$ mag for the latest type galaxies. We also find a strong evolution in the normalization of the luminosity function of latest type galaxies, with an increase of more than a factor $`2`$ in the $`\varphi ^{}`$ parameter going from $`z0.3`$ to $`z1.3`$. The density of bright ($`M_{B_{AB}}5log(h)<20`$) galaxies shows a modest decrease ($`40\%`$) for early type objects from $`z0.3`$ to $`z1.1`$; on the contrary, the number of bright late type galaxies increases of a factor $`6.6`$ from $`z0.3`$ to $`z1.3`$.
Our results indicate that the importance of type 4 galaxies is increasing with redshift, with an important contribution of both bright and faint blue objects. This fact is also largely responsible of the evolution of the global luminosity function measured by Ilbert et al. (vvdsLF (2005)), which shows a brightening of $`M^{}`$ and a steepening of $`\alpha `$ with increasing redshift. Moreover, the increasing contribution of blue galaxies has been seen in the evolution of the GALEX-VVDS luminosity function at 1500Å(Arnouts et al. galexLF (2005)). We are therefore pinpointing that the galaxies responsible for most of the evolution quoted in the literature belong to the population of the latest spectral type. The epoch at which a transition between a Universe dominated by late type galaxies and a Universe dominated by old massive objects occurs is at a redshift of $`z0.70.8`$.
The fact that type 1 galaxies show only a mild evolution both in luminosity (positive) and in density (negative) is consistent with the fact that most of the objects in this class are old ($`z_{form}>2`$, see sect. 3) galaxies, experiencing only a passive evolution in the explored redshift range. More intriguing is the density evolution of type 4 galaxies, which corresponds to an increasing number of bright star forming galaxies towards high redshift which could be connected to various populations of high redshift objects seen in multiwavelength surveys.
###### Acknowledgements.
This research has been developed within the framework of the VVDS consortium.
This work has been partially supported by the CNRS-INSU and its Programme National de Cosmologie (France), and by Italian Ministry (MIUR) grants COFIN2000 (MM02037133) and COFIN2003 (num.2003020150).
The VIMOS VLT observations have been carried out on guaranteed time (GTO) allocated by the European Southern Observatory (ESO) to the VIRMOS consortium, under a contractual agreement between the Centre National de la Recherche Scientifique of France, heading a consortium of French and Italian institutes, and ESO, to design, manufacture and test the VIMOS instrument.
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# Hubbard-like Hamiltonian for ultracold atoms in a 1D optical lattice
## I Introduction
Since the Bose-Einstein condensation of alkali atoms in magnetic traps Anderson et al. (1995); Davis et al. (1995), a massive experimental and theoretical effort has been dedicated to the investigation of confined atoms in the extremely low-temperature regime (for a review see Leggett (2001); Dalfovo et al. (1999); Jaksch and Zoller (2005); Pitaevskii and Stringari (2003)).
The flexibility of optical trapping techniques has suggested the devise of different configurations (lattices Anderson and Kasevich (1998); Orzel et al. (2001); Morsch et al. (2001); Jaksch and Zoller (2005); Jaksch et al. (1998), superlattices, etc. Buonsante et al. (2004); Buonsante and Vezzani (2003); Peil et al. (2003); Guidoni and Verkerk (1998); Roth and Burnett (2003); Santos et al. (2004)), opening a vast scenario of research. The ability to tune atomic interactions via a magnetic field (Feshbach resonance Kokkelmans et al. (2002)), along with the proposal of single atom trap loading techniques Rabl et al. (2003), has proven to be of capital importance for ultracold fermions physics, yielding the possibility to study fundamental aspects of superfluidity (BCS-BEC crossover, see e.g. Milstein et al. (2002); Giorgetti et al. (2004); Greiner et al. (2003); Astrakharchik et al. (2004)) and envisaging new perspectives in quantum information processing Calarco et al. (2004); Zanardi (2002).
The present work focuses on the theoretical investigation of the properties of (few) fermionic ultracold atoms loaded into a 1D optical lattice, where global confinement is ensured by a magnetic trap. The description of such a physical system can be naturally performed in terms of a generalised Hubbard Hamiltonian (gHH) which is deduced from a general field-theoretic Hamiltonian with two body interaction Greiner W. (1998). At this stage, particular care must be taken in the choice of the function basis for the field operator expansion. Although the symmetries of the system can provide selection rules that reduce the involvement of the gHH, the resulting coefficient structure is very rich and, as a direct consequence, the Hamiltonian hardly tractable. Nevertheless, the generality of the model gives rise to a wealth of sub-models, depending upon different approximations and regimes. The guideline to find simplified Hamiltonians is given by the thorough analysis of the gHH coefficient structure.
From this perspective the analytical knowledge of the coefficients is a powerful tool to establish the physical relevance of different sub-models in the various situations that may be conceived in the framework of the trapped ultracold atoms physics. Moreover, the nontrivial dependence of the coefficient from controllable external parameters provides the possibility to use these parameters to control the dynamics of the atoms trapped in the optical lattice. Thus the key aspect of this paper is the analytical determination of the hopping and interaction coefficients as a function of experimental parameters such as magnetic trap frequency, laser intensity, wavelength, angle between laser sources, $`s`$-wave scattering length etc.
We would like to stress that the procedure followed here for the determination of the coefficients is statistics-independent: the bosonic or fermionic nature of the atoms loaded into the trap is completely taken into account by the commutators of raising and lowering operators that will be described in the paper. For example with the calculations performed here it seems feasible to go beyond the approximations that lead to the Bose-Hubbard model in the description of the BEC dynamics in optical lattices, taking into account the specific nature of the interaction between alkali atoms in a low density regime. Even for a ground-state calculation it can be shown that it is necessary to include levels beyond the single particle ground state (see Abrikosov A. A. (1977)).
The confinement model considered here has a direct experimental relevance (see e.g. Modugno et al. (2002); Pezzé et al. (2004)). However, while in Modugno et al. (2002); Pezzé et al. (2004) a number of atoms of the order of $`10^4`$ is considered, allowing thus the adoption of a semiclassical model, we focus on a low occupation-number regime similarly to what is done in Hofstetter et al. (2002) and Albus et al. (2003), yet extending to a multi-band model whose correctness is limited by the validity limit of the space-mode approximation.
Challenging tasks for the future will include the determination of tractable yet interesting models for different aspects for theoretical condensed matter physics and quantum mechanics. On the other hand the experimental realisation of systems that exhibit a behaviour which can be described in the framework of the various models here proposed, would represent an important achievement for both condensed matter experimentalist and theoreticians: the main difficulties seem to arise form the nearly-single atom trap loading and, quite naturally, from the coupling with the external environment.
Throughout the paper we have tried to emphasise the generality of the procedure followed. However we have decided to write down and plot few numerical values of the coefficients to stress the fact that this calculation is a direct and relatively simple tool to shape out simplified and approximate Hamiltonians for different physical situations.
In section II we depict the potential configuration of the system, moving then to the description of our field-theoretical approach. The field operators are written in terms of mode raising and lowering operators. Each mode corresponds to a set of quantum numbers, one of them identifies the lattice site (hence space-mode approximation) while the others describe on-site quantum numbers (local-mode)Giampaolo et al. (2004). As previously stated this choice is not unique, but symmetry constraints suggest expansions that emphasise conservation laws and selection rules.
In section III we evaluate the expression of the Hamiltonian hopping and interaction coefficients and we try to describe the interaction coefficient symmetry properties into some detail.
The purpose of section IV is twofold. One the one hand we show how, with suitable approximations, the Hamiltonian of the system reduces to known cases, such as the Hubbard Hamiltonian or a trivial non-interacting Hamiltonian. On the other we introduce a novel Rotational Hubbard Hamiltonian, as a first instance of the involvement of higher order approximations. For this case, by means of established mean-field approaches Bach et al. (1994); Zhang et al. (1990), we suggest a possible path of research involving general group-theoretical procedures Zhang et al. (1990). It will be shown that these procedures, even if the explicit solution for the ground state is not given, allow to grasp interesting aspects of the physics of the model here discussed.
We have included two Appendices where the relatively simple but lengthy calculations of the tunnelling and interaction coefficients are provided explicitly. In Appendix A there are various plots of multilevel hopping parameters that supply a good example of the scenario that we are moving in and may constitute a good starting point for further investigation.
## II Fermions trapped in 1D optical lattices
### II.1 General features
The general field-theoretic Hamiltonian (see e.g. Greiner W. (1998)) with 2-body interaction can be written as
$$\widehat{𝖧}=𝑑𝐫\widehat{\mathrm{\Psi }}^{}(𝐫)𝖧_{1b}(𝐫)\widehat{\mathrm{\Psi }}(𝐫)+𝑑𝐫𝑑𝐫^{}\widehat{\mathrm{\Psi }}^{}(𝐫)\widehat{\mathrm{\Psi }}^{}(𝐫^{})𝖧_{2b}(𝐫,𝐫^{})\widehat{\mathrm{\Psi }}(𝐫^{})\widehat{\mathrm{\Psi }}(𝐫)$$
(1)
where $`𝖧_{1b}(𝐫)`$ represents the 1-body term of the Hamiltonian (kinetic + external potential term) while $`𝖧_{2b}(𝐫,𝐫^{})`$ the 2-body interaction potential term, $`\widehat{\mathrm{\Psi }}(𝐫)`$ is the field operator and $`\widehat{\mathrm{\Psi }}^{}(𝐫)`$ its adjoint.
As previously mentioned we will stick to neutral fermionic atoms loaded into a 1D optical lattice. The lattice is generated by two lasers counter-propagating along the $`x`$-axis, with wavenumber $`K`$. The depth – or height, depending if red or blue detuning of the laser is considered – in each point $`x`$ the potential is proportional to the intensity of the laser and thus, according to the considered setup, to $`\mathrm{sin}^2(2Kx)`$, for the evaluation of the multiplicative constant see e.g. Leggett (2001). Here we set the multiplicative constant equal to $`m\omega ^2/(2K^2)`$ where $`\omega `$ represents the harmonic oscillator frequency in the second order expansion of the term $`V_{ext}`$.
Global confinement is ensured by a cigar-shaped magnetic trap with principal axis along the $`x`$-direction (see e.g. Modugno et al. (2003)). This trap can be modelled by a 3D harmonic anisotropic trap of axial and radial frequencies equal to $`\mathrm{\Omega }_x`$ and $`\mathrm{\Omega }_{}`$ respectively ($`\mathrm{\Omega }_x\mathrm{\Omega }_{}`$).
The magneto-optical trap can be thought as if the constituents of the system were trapped in the cigar shaped potential with a “slicing” effect of the laser, giving rise to a linear array of 3D prolate harmonic oscillators. Besides, the radial trapping frequency has a deep influence on the interaction among the constituents of the system, allowing to control the volume of each “disk”. With the previous assumptions $`𝖧_{1b}`$ becomes
$$𝖧_{1b}=E_{kin}+V_{ext}$$
where
$`E_{kin}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2^2}{2m}}`$
$`V_{ext}`$ $`=`$ $`{\displaystyle \frac{m}{2}}\left[\mathrm{\Omega }_x^2x^2+\mathrm{\Omega }_{}^2\rho ^2\right]+{\displaystyle \frac{m\omega ^2}{2K^2}}\mathrm{sin}^2(Kx),`$ (2)
and the second term of $`V_{ext}`$ represents the harmonic confinement of the magnetic trap, while the third one corresponds to the optical potential and $`\rho ^2=y^2+z^2`$. For future convenience, we write equation (II.1) as
$$𝖧_{1b}=E_{kin}+\underset{j}{}V_j+\left(V_{ext}\underset{j}{}V_j\right)$$
(3)
with
$$V_j=\mathrm{\Pi }_j(x)\frac{m}{2}\left[\mathrm{\Omega }_x^2j^2\frac{\pi ^2}{k^2}+\omega ^2x_j^2+\mathrm{\Omega }_{}^2\rho ^2\right]$$
(4)
$`\mathrm{\Pi }_j(x)=\mathrm{\Pi }\left(\frac{Kx}{\pi }j\right)`$ where $`\mathrm{\Pi }(x)`$ is the rectangle function ($`\mathrm{\Pi }(x)=1`$ for $`1x<1`$, $`\mathrm{\Pi }(x)=0`$ elsewhere), $`k=l_{}K`$ (with $`l_{}=\sqrt{\mathrm{}/(m\omega _{})})`$ and $`x_j=\left(xj\frac{\pi }{k}\right)`$. Here the harmonic axial confinement of the magnetic field has been considered as a site-dependent – with $`j`$ site index – constant addictive term, merely shifting the local minima of the optical potential. From Eq. (3) with the properties of the rectangle function we obtain
$$𝖧_{1b}=\underset{j}{}\mathrm{\Pi }_j(x)\left[\left(E_{kin}+V_j\right)+\left(V_{ext}V_j\right)\right]$$
(5)
Hereafter the axial confinement of the magnetic trap will be neglected (small $`\mathrm{\Omega }_x`$).
We are now led to consider two different terms in Eq. (5). The first represents a local harmonic-oscillator Hamiltonian
$`𝖧_j^{ho}`$ $`=`$ $`\mathrm{\Pi }_j(x)\left(E_{kin}+V_j\right)`$ (6)
$`=`$ $`\mathrm{\Pi }_j(x)\left[{\displaystyle \frac{\mathrm{}^2^2}{2m}}+{\displaystyle \frac{m\omega ^2}{2}}x_j^2+{\displaystyle \frac{m\mathrm{\Omega }_{}^2}{2}}\rho ^2\right]`$
and a hopping one
$`𝖧_j^{tunn}`$ $`=`$ $`\mathrm{\Pi }_j(x)\left(V_{ext}V_j\right)`$ (7)
$`=`$ $`\mathrm{\Pi }_j(x)\left[{\displaystyle \frac{m\omega ^2}{2K^2}}\mathrm{sin}^2(Kx){\displaystyle \frac{m\omega ^2}{2}}x_j^2\right].`$
The term $`V_j`$ is the local second-order expansion of the optical potential, thus equation (7) represents the discrepancy between an harmonic potential and the true optical potential, describing hopping of atoms between neighbouring sites.
Neutrality of the atoms, ensuring a finite-range interaction allows us to introduce a pseudo-potential approximation (see e.g. Kerson Huang (1987))
$$U(𝐫)=\underset{j}{}\mathrm{\Pi }_j(x)\stackrel{~}{a}_s\delta (𝐫)\frac{}{𝐫}𝐫,\stackrel{~}{a}_s:=\frac{4\pi \mathrm{}^2a_s}{m},$$
(8)
where $`𝐫`$ is the interatomic distance and $`a_s`$ the $`s`$wave scattering length ($`a_s`$ in our approximation is considered constant). The validity of this model is ensured by the low energies involved in these interactions, direct consequence of both low temperature limit (virtually zero) and diluteness (low Fermi energy). Besides, the form of Eq. (8) shows that on-site terms only will contribute to the interaction Hamiltonian. Thus Eq. (1) can be rewritten in the form
$$\begin{array}{c}\widehat{𝖧}=\underset{j}{}[d𝐫\widehat{\mathrm{\Psi }}^{}(𝐫)(𝖧_j^{ho}(𝐫)+𝖧_j^{tunn}(𝐫))\widehat{\mathrm{\Psi }}(𝐫)+\hfill \\ \hfill +\stackrel{~}{a}_sd𝐫d𝐫^{}\widehat{\mathrm{\Psi }}^{}(𝐫)\widehat{\mathrm{\Psi }}^{}(𝐫^{})\delta (𝐫𝐫^{})\widehat{\mathrm{\Psi }}(𝐫^{})\widehat{\mathrm{\Psi }}(𝐫)].\end{array}$$
(9)
### II.2 The (space+local)-modes expansion
The choice of the basis for the expansion of the field operators is crucial. As already suggested by the grouping of terms in Eq. (6), we will choose a basis constituted by local harmonic oscillator eigenfunctions. In addition, because of the symmetry of the system we have chosen central-symmetric 2D h.o. eigenfunctions for the 2D isotropic radial h.o.Wünsche (1998), instead of decomposing it in 1D h.o. eigenfunctions, this will give us deeper insight into conservation laws and selection rules imposed by the symmetries of the system. We then have
$$\widehat{\mathrm{\Psi }}(𝐱)=\underset{i,n_x,J,m,\sigma }{}u_{n_x}(xx_i)_{J,m}(\rho ,\varphi )\xi (\sigma )\widehat{c}_{n_x,J,m,i,\sigma }$$
(10)
with
$`u_n(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2^nn!\sqrt{\pi }l_x}}}H_n(x/l_x)e^{\frac{x^2}{2l_x^2}},`$ (11)
$`_{J,m}(\rho ,\varphi )`$ $`=`$ $`{\displaystyle \frac{e^{2im\varphi }}{\sqrt{\pi }l_{}}}C_{Jm}\left({\displaystyle \frac{\rho }{l_{}}}\right)^{2m}L_{Jm}^{2J}\left(\rho /l_{}\right),`$ (12)
$`C_{Jm}=\sqrt{(J+m)!/(Jm)!}`$, $`\xi (\sigma )`$ is a spin function and $`l_x=\sqrt{\mathrm{}/(m\omega _x)}`$. In this decomposition $`u_n(x)`$ is a 1D harmonic-oscillator eigenfunction ($`H_n`$ represent the $`n`$th Hermite polynomial) and $`_{J,m}`$ a 2D harmonic-oscillator eigenfunction Wünsche (1998) with $`L_{Jm}^{2J}(x)`$ a generalised Laguerre polynomial.
Fermionic operators will thus have 5 indexes: 3 of them ($`n_x,J,m`$) identify (2+1)D local harmonic oscillator states, while $`i`$ identifies the site and $`\sigma `$ the spin. While $`n_x`$ has its usual interpretation of 1D harmonic oscillator number operator eigenvalue, $`J`$ and $`m`$ can be construed as angular momentum and $`x`$axis component of the angular momentum, respectively.
This decomposition can be thought as a generalised space-mode approximation, with additional local modes that, in the present case, correspond to the local (2+1)D harmonic-oscillators quantum numbers. If not explicitly required, we will use $`\varphi _\alpha =u_{n_\alpha }(xx_{i_\alpha })_{J_\alpha ,m_\alpha }(\rho ,\varphi )`$ $`\xi (\sigma _\alpha )`$, with $`\alpha =\{n_\alpha ,J_\alpha ,m_\alpha ,i_\alpha ,\sigma _\alpha \}`$ to simplify the index notation. We wish to stress that decomposition (10) is an approximation of field $`\widehat{\mathrm{\Psi }}(𝐱)`$: there is a non-nil overlapping between wavefunctions belonging to different sites, thus orthogonality is not fulfilled. Nevertheless these overlapping integrals are supposed to be small, ensuring the consistency of this choice Jaksch et al. (1998).
In the forthcoming calculation of the interaction term, Eq. (11) allows us to easily recognise that $`m`$ is a conserved quantity. If we come back to (1), with the decomposition (10) we obtain
$`\widehat{𝖧}={\displaystyle \underset{j}{}}{\displaystyle \underset{\alpha ,\beta }{}}[{\displaystyle }d𝐫\varphi _\alpha ^{}(𝐫)𝖧_j^{ho}\varphi _\beta (𝐫)\widehat{c}_\alpha ^{}\widehat{c}_\beta +{\displaystyle }d𝐫\varphi _\alpha ^{}(𝐫)𝖧_j^{tunn}\varphi _\beta (𝐫)\widehat{c}_\alpha ^{}\widehat{c}_\beta +`$
$`\stackrel{~}{a}_s{\displaystyle \underset{\gamma ,\delta }{}}{\displaystyle }d𝐫d𝐫^{}\varphi _\alpha ^{}(𝐫)\varphi _\gamma ^{}(𝐫^{})\times \delta (𝐫𝐫^{})\varphi _\beta (𝐫)\varphi _\delta (𝐫^{})\widehat{c}_\alpha ^{}\widehat{c}_\beta ^{}\widehat{c}_\delta \widehat{c}_\gamma ]`$ (13)
## III Hamiltonian Coefficients
We are now in the position to calculate all the coefficients in Hamiltonian (II.2). The first term becomes
$$\widehat{𝖧}^{ho}=\underset{j,\alpha ,\beta }{}\lambda _\beta _{\mathrm{}}^{\mathrm{}}𝑑𝐫\mathrm{\Pi }_j(x)\varphi _\alpha ^{}(𝐫)\varphi _\beta (𝐫)\widehat{c}_\alpha ^{}\widehat{c}_\beta ,$$
(14)
where $`\lambda _\beta `$ is the (2+1)D harmonic-oscillator eigenvalue
$$\lambda _{n_\beta ,J_\beta ,m_\beta ,i_\beta ,\sigma _\beta }=\left[\mathrm{}\omega _x\left(n_\beta +\frac{1}{2}\right)+\mathrm{}\left(2J_\beta +1\right)\right].$$
(15)
Eq. (14) can be written as
$$\underset{j,\alpha ,\beta }{}\lambda _\beta \delta _{\alpha ,\beta }\delta _{i_\alpha ,j}\widehat{c}_\alpha ^{}\widehat{c}_\beta =\underset{\alpha }{}\lambda _\alpha \widehat{n}_\alpha $$
(16)
where the second Kronecker delta is a consequence of the space-mode approximation, i.e. we consider only superposition of wavefunctions among which at least one is a local harmonic-oscillator eigenfunction, while the first one stems from the orthogonality of the $`\varphi _\gamma (x)`$ functions.
We move now to the evaluation of the integral in the second term of equation (II.2). Namely
$$\widehat{𝖧}^{tunn}=\underset{j,\alpha ,\beta }{}𝑑𝐫\varphi _\alpha ^{}(𝐫)\mathrm{\Pi }_j(x)𝖧_j^{tunn}(x)\varphi _\beta (𝐫)$$
(17)
being $`𝖧_j^{tunn}(x)`$ independent of radial and spin degrees of freedom, we can rewrite equation (17) as
$$\widehat{𝖧}^{tunn}=\underset{j,\alpha ,\beta }{}\delta _{J_\alpha ,J_\beta }\delta _{m_\alpha ,m_\beta }\delta _{\sigma _\alpha ,\sigma _\beta }\delta _{i_\alpha ,j}𝑑xu_{n_\alpha ,i_\alpha }^{}(x)𝖧_j^{tunn}(x)u_{n_\beta ,i_\beta }(x)\widehat{c}_\alpha ^{}\widehat{c}_\beta .$$
(18)
With the same assumptions of the local harmonic-oscillator case the integral in equation (18) becomes
$$K_{n_\alpha ,n_\beta }\mathrm{}\omega _x𝑑ye^{\frac{(y\tau )^2}{2}}H_{n_\alpha }\left(y\tau \right)𝖧_j^{tunn}(y)e^{\frac{y^2}{2}}H_{n_\beta }\left(y\right),$$
(19)
where $`K_{n_\alpha ,n_\beta }=[2^{n_\alpha +n_\beta }n_\alpha !n_\beta !\pi ]^{1/2}`$ and we have put $`y=(xi_\beta d)/l_x`$ (where $`d=\pi /k`$), $`\tau =(i_\beta i_\alpha )`$ and $`\mathrm{\Omega }=Kl_x`$. By substituting the expression of $`𝖧_j^{tunn}(y)`$ from equation Eq. (7) we obtain
$$K_{n_\alpha ,n_\beta }\mathrm{}\omega _x𝑑ye^{\frac{(y\tau )^2+y^2}{2}}H_{n_\alpha }\left(y\tau \right)H_{n_\beta }\left(y\right)\left[\frac{1\mathrm{cos}(2\mathrm{\Omega }y)}{4\mathrm{\Omega }^2}\frac{y^2}{2}\right]\widehat{c}_\alpha ^{}\widehat{c}_\beta .$$
(20)
If we define
$$\begin{array}{c}T_{\alpha ,\beta }=\frac{K_{n_\alpha ,n_\beta }}{2}\delta _{J_\alpha ,J_\beta }\delta _{m_\alpha ,m_\beta }\delta _{\sigma _\alpha ,\sigma _\beta }\mathrm{}\omega _x𝑑ye^{\frac{(y\tau )^2+y^2}{2}}\hfill \\ \hfill \times H_{n_\beta }\left(y\right)H_{n_\alpha }\left(y\tau \right)\left[\frac{y^2}{2}\frac{1\mathrm{cos}(2\mathrm{\Omega }y)}{4\mathrm{\Omega }^2}\right]\end{array}$$
(21)
equation (14) becomes
$$\widehat{𝖧}^{tunn}=\underset{\alpha ,\beta }{}T_{\alpha ,\beta }\widehat{c}_\alpha ^{}\widehat{c}_\beta .$$
(22)
Then the term $`T_{\alpha ,\alpha }\widehat{c}_\alpha ^{}\widehat{c}_\alpha `$ can be incorporated into the $`\widehat{𝖧}^{h.o.}`$ term, giving
$$\mu _\alpha =\lambda _\alpha T_{\alpha ,\alpha }.$$
(23)
We will here skip the explicit solution of the integral in Eq. (21), along with the analytic expression of T, which can be found in Appendix A. These calculations allow us to write
$$T_{\alpha ,\beta }=\delta _{J_\alpha ,J_\beta }\delta _{m_\alpha ,m_\beta }\delta _{\sigma _\alpha ,\sigma _\beta }T_{n_\alpha ,n_\beta ,i_\alpha ,i_\beta }.$$
(24)
In fig. 1 we have the plot the coefficient $`T_{n_\alpha ,n_\beta ,i_\alpha ,i_\beta }`$ as a function of the ratio between distance and the period of the optical lattice, for $`n_\alpha ,n_\beta =0,1`$. In boldface we have marked the points corresponding to discrete values of the ratio $`x/d`$, i.e. the points with a relevant physical meaning, The values of $`T`$ plotted here are in arbitrary units. Even if the correctness of the above procedure seems undoubted, it must be remembered that it is entirely based on the space-mode approximation, whose validity depends on the overlapping of wavefunctions belonging to different sites and thus might be violated.
These plots show how the tunnelling amplitude varies with the distance. In particular it is clear how, for long-distance tunnelling, there is a negative exponential dependence. Nevertheless, if the experimental conditions are properly chosen (i.e. angle between counterpropagating laser beams and their power), it is possible to obtain conditions where, for instance nearest-neighbour and next-to-nearest neighbour tunnelling coefficients have opposite signs (see e.g. Fig.2 $`T_{0,0,i_\alpha ,i_\beta }`$), and thus the model, in that case, might exhibit frustration.
We will now move to the determination of the interaction term, namely the last term of equation (II.2). As a first step, we can write the integral in cylindrical coordinates
$`\stackrel{~}{a}_s{\displaystyle 𝑑𝐫𝑑𝐫^{}\varphi _\alpha ^{}(𝐫)\varphi _\gamma ^{}(𝐫^{})\delta (𝐫𝐫^{})\varphi _\beta (𝐫)\varphi _\delta (𝐫^{})}=`$ (25)
$`\stackrel{~}{a}_s{\displaystyle 𝑑x𝑑x^{}u_{n_\alpha }^{}(xx_{i_\alpha })u_{n_\gamma }^{}(x^{}x_{i_\gamma })\delta (xx^{})u_{n_\beta }(xx_{i_\beta })u_{n_\delta }(x^{}x_{i_\delta })}`$
$`\times {\displaystyle }d\stackrel{~}{\rho }d\stackrel{~}{\rho }^{}{\displaystyle \frac{\stackrel{~}{\rho }}{\pi }}{\displaystyle }d\varphi d\varphi ^{}_{J_\alpha ,m_\alpha }^{}(\stackrel{~}{\rho },\varphi )_{J_\gamma ,m_\gamma }^{}(\stackrel{~}{\rho }^{},\varphi ^{})\delta (\rho \rho ^{})`$
$`\delta (\varphi \varphi ^{})_{J_\beta ,m_\beta }(\stackrel{~}{\rho },\varphi )_{J_\delta ,m_\delta }(\stackrel{~}{\rho }^{},\varphi ^{})`$
with $`\stackrel{~}{\rho }=\rho /l_\rho `$ and the identity
$$\delta (𝐫)=\frac{\delta (\rho )\delta (\varphi )}{\pi \rho }.$$
As we are dealing with a short range interaction modelled by a $`\delta (𝐫𝐫^{})`$ function, we will consider on-site interaction only ($`\stackrel{~}{x}_{i_\alpha }=\stackrel{~}{x}_{i_\beta }=\stackrel{~}{x}_{i_\delta }=\stackrel{~}{x}_{i_\gamma }`$).
This choice is completely justified because the interaction term is modelled by a pseudopotential term for which nearest-neighbours interactions become negligible. In this case the first integral on the left-hand side of Eq. (25) becomes
$`U_x={\displaystyle \frac{1}{\pi l_x}}\sqrt{{\displaystyle \frac{2^{(n_\alpha +n_\beta +n_\gamma +n_\delta )}}{n_\alpha !n_\beta !n_\gamma !n_\delta !}}}{\displaystyle 𝑑\stackrel{~}{x}H_{n_\alpha }(\stackrel{~}{x})H_{n_\beta }(\stackrel{~}{x})H_{n_\gamma }(\stackrel{~}{x})H_{n_\delta }(\stackrel{~}{x})e^{2\stackrel{~}{x}^2}}`$ (26)
with $`\stackrel{~}{x}=x/l_x`$, whose explicit calculation is given in Appendix B. Here we just give the final result
$$U_x=\frac{\delta _{\overline{n},2}}{\pi l_x}\underset{\overline{s}}{\overset{\overline{n}}{}}\frac{\mathrm{\Xi }(\overline{s})}{\sqrt{2}^{\overline{s}+3}}\mathrm{\Gamma }\left[\frac{\left(\overline{n}\overline{s}\right)}{2}+1\right]$$
(27)
with
$$\mathrm{\Xi }\left(\overline{s}\right)=\{\begin{array}{cc}0\hfill & \text{if }s_\theta \text{ odd}\hfill \\ _\theta \frac{1}{n_\theta !}\left(\genfrac{}{}{0pt}{}{n_\theta }{s_\theta }\right)H_{s_\theta }(0)\hfill & \text{if }s_\theta \text{ even}\hfill \end{array}$$
(28)
The summation is to be intended as 4 separate summations over the components of a vector $`\overline{s}=\{s_\alpha ,s_\beta ,s_\gamma ,s_\delta \}`$ from $`\{0,0,0,0\}`$ to $`\overline{n}=\{n_\alpha ,n_\beta ,n_\gamma ,n_\delta \}`$. The norm $`\overline{x}`$ is a 1-norm ($`\overline{x}=_\theta |x_\theta |`$, $`\theta =\alpha ,\beta ,\gamma ,\delta `$) and the $`\delta `$ in Eq. (26) represents the parity selection rule, obtained from the explicit calculation of the integral.
For the radial part of the integral we have
$$U_\rho =𝑑\rho 𝑑\varphi \frac{\rho }{\pi }_{J_\alpha ,m_\alpha }^{}(\rho ,\varphi )_{J_\gamma ,m_\gamma }^{}(\rho ,\varphi )_{J_\beta ,m_\beta }(\rho ,\varphi )_{J_\delta ,m_\delta }(\rho ,\varphi )$$
(29)
with the definition given by Eq. (11), we can easily perform the angular integration and we obtain
$$U_\rho =\frac{2\delta _{m_\alpha +m_\gamma ,m_\beta +m_\delta }}{\pi ^2}_0^{\mathrm{}}𝑑\rho \rho R_{J_\alpha ,m_\alpha }^{}(\rho )R_{J_\gamma ,m_\gamma }^{}(\rho )R_{J_\beta ,m_\beta }(\rho )R_{J_\delta ,m_\delta }(\rho ).$$
(30)
The reader is again addressed to Appendix B for the explicit evaluation of the integral in Eq. (30). The result is given by
$$U_\rho =\frac{\delta _{m_\alpha +m_\gamma ,m_\beta +m_\delta }}{\pi ^2l_{}^2}\underset{\overline{q}=\overline{|m|}}{\overset{\overline{J}}{}}\mathrm{\Lambda }(\overline{J},\overline{m},\overline{q})\frac{\mathrm{\Gamma }\left(\overline{q}+3/2\right)}{2^{\overline{q}+3/2}}$$
(31)
with
$$\mathrm{\Lambda }(\overline{J},\overline{m},\overline{q})=\underset{\theta =\alpha ,\beta ,\gamma ,\delta }{}\frac{(1)^{J_\theta q_\theta }\sqrt{(J_\theta +m_\theta )!(J_\theta m_\theta )!}}{(J_\theta q_\theta )!(q_\theta +m_\theta )!\left(q_\theta m_\theta \right)!}$$
(32)
and, following previous notation, we obtain $`\overline{q}=\{q_\alpha ,q_\beta ,q_\gamma ,q_\delta \}`$, $`\overline{J}=\{J_\alpha ,J_\beta ,J_\gamma ,J_\delta \}`$ and $`\overline{|m|}=\{|m_\alpha |,|m_\beta |,|m_\gamma |,|m_\delta |\}`$. The overall interaction coefficient can then be written as the product of Eqs. (31) and (27)
$$\begin{array}{c}U_{\alpha ,\beta ,\gamma ,\delta }=\delta _{\overline{n},2}\delta _{m_\alpha +m_\beta ,m_\gamma +m_\delta }\frac{\stackrel{~}{a}_s}{4l_x\pi ^3l_{}^2}\underset{\overline{q}=\overline{|m|}}{\overset{\overline{J}}{}}\underset{\overline{s}}{\overset{\overline{n}}{}}\frac{\mathrm{\Lambda }(\overline{J},\overline{m},\overline{q})\mathrm{\Xi }\left(\overline{s}\right)}{\sqrt{2}^{\overline{s}+2\overline{q}}}\hfill \\ \hfill \times \mathrm{\Gamma }\left(\overline{q}+\frac{3}{2}\right)\mathrm{\Gamma }\left[\frac{(\overline{n}\overline{s}+1)|}{2}\right].\end{array}$$
(33)
We are thus enabled to rewrite Hamiltonian (II.2) in terms of the calculated coefficients obtaining
$$\widehat{𝖧}=\underset{j}{}\left[\underset{\alpha }{}\lambda _{\alpha ,\beta }\widehat{n}_\alpha ^{}+\underset{\alpha ,\beta }{}T_{\alpha ,\beta }\widehat{c}_\alpha ^{}\widehat{c}_\beta +\underset{\alpha ,\beta ,\gamma ,\delta }{}U_{\alpha ,\beta ,\gamma ,\delta }\widehat{c}_\alpha ^{}\widehat{c}_\beta ^{}\widehat{c}_\delta \widehat{c}_\gamma \right].$$
(34)
We will refer to (34) as the generalised Hubbard Hamiltonian. For sake of simplicity, in Eq. (34) we have not written down explicitly the selection rules imposed by symmetry constraints (see below).
### III.1 Symmetry properties of the interaction term
In addition to global symmetry properties, such as 1) rotational symmetry along the $`x`$-axis and 2) left-right symmetry, reflected by momentum $`x`$-component conservation and parity conservation for the 1D harmonic oscillators along the $`x`$-axis, it is clear from equation (33) that the coefficient $`U_{\alpha ,\beta ,\gamma ,\delta }`$ has some symmetry properties: a) $`U`$ does not depend on the sign of $`m_\chi `$ with $`\chi =\alpha ,\beta ,\gamma ,\delta `$, provided the conservation $`m`$ during the interaction ($`m_\alpha +m_\beta =m_\gamma +m_\delta `$); b) $`U`$ possesses a permutational symmetry, namely
$$U_{\alpha ,\beta ,\gamma ,\delta }=U_{\beta ,\alpha ,\gamma ,\delta }=U_{\alpha ,\beta ,\delta ,\gamma }=U_{\beta ,\alpha ,\delta ,\gamma }.$$
(35)
We would like to draw reader’s attention to the two $`\delta `$ functions in equation (33) which make explicit the conservation laws that might have been expected by simply considering the symmetry of the problem. The first one represents parity conservation, while the second conservation of the $`x`$ component of the angular momentum. In table 1, we give the analytical value of $`U`$ for interaction between particles belonging to the first three shells of the 2D radial harmonic oscillator and to the first level for the axial harmonic oscillator. These symmetry constraints allow to class the possible quantum numbers of the interacting particles according to the value of the corresponding $`U`$. For example for
$$\alpha =\{0,1,0,i,\sigma \},\beta =\{0,0,0,i,\sigma ^{}\},$$
$$\delta =\{0,0,0,i,\sigma \},\gamma =\{0,0,0,i,\sigma ^{}\},$$
and
$$\alpha =\{0,0,0,i,\sigma \},\beta =\{0,0,0,i,\sigma ^{}\},$$
$$\delta =\{0,1,0,i,\sigma \},\gamma =\{0,0,0,i,\sigma ^{}\},$$
we have the same value of $`U`$, henceforth the class definition of table 1. We would like to point out two aspects of this example. First of all it may be noticed that the angular momentum $`J`$ is not conserved throughout the interaction: this is a general feature of the system considered, there is not global rotational symmetry but only in the plane orthogonal to the 1D optical lattice. Moreover in this particular interaction the value for the coefficient $`U`$ is negative, this appears to be a rare (but not unique) situation. The implications of this condition will be pointed out in section IV
## IV Special Cases
In this section we derive three model Hamiltonians for fermions in optical lattices. We consider, for both cases, only the lowest-state axial quantum number (i.e. $`n_\alpha =n_\beta =0`$). Hence $`T_{\alpha ,\beta }`$ can be written as
$$T_{\alpha ,\beta }=\delta _{J_\alpha ,J_\beta }\delta _{m_\alpha ,m_\beta }\delta _{\sigma _\alpha ,\sigma _\beta }T_{0,0,i_\alpha ,i_\beta }$$
(36)
As far as an ultracold gas is considered, it seems feasible to restrict our analysis to the first few levels above the ground state (i.e. $`J_\alpha =0,1/2,\mathrm{}`$). As a first example, we consider the case having $`J_\alpha =0`$ as the only radial level allowed and the fermionic gas is spin unpolarised. From Eq. (33), along with the previous assumptions, we obtain Ruuska and Törmä (2004)
$$\widehat{H}=\underset{i,\sigma }{}\mu _i\widehat{n}_{i,\sigma }T\underset{i,\sigma }{}\left(\widehat{c}_{i,\sigma }^{}\widehat{c}_{i+1,\sigma }+\widehat{c}_{i+1,\sigma }^{}\widehat{c}_{i,\sigma }\right)+U\underset{i,\sigma ,\sigma ^{}}{}\widehat{n}_{i,\sigma }\widehat{n}_{i,\sigma ^{}}$$
(37)
with $`T=T_{0,0,i_\alpha ,i_\alpha +1}`$ which is easily recognised as the Hubbard Hamiltonian, whose role in the ultracold atoms physics has been pointed out elsewhere Jaksch and Zoller (2005); Hofstetter et al. (2002). Note that in this example we have made the assumption that the tunnelling coefficient is significantly different from zero only for nearest-neighbouring sites. Nevertheless more involved situations may arise, suggesting interesting physical features, as it is shown in Appendix A.
If we now consider a spin-polarised gas in a (radial) multi-level system we obtain Ruuska and Törmä (2004)
$$\widehat{H}=\underset{\overline{n},i}{}\mu _{\overline{n},i}\widehat{n}_{\overline{n},i}T\underset{\overline{n},i}{}\left(\widehat{c}_{\overline{n},i}^{}\widehat{c}_{\overline{n},i+1}+\widehat{c}_{\overline{n},i+1}^{}\widehat{c}_{\overline{n},i}\right)$$
(38)
where the absence of the interaction term is related to the symmetry properties of the coefficient $`U_{\alpha \beta \gamma \delta }`$. The Hamiltonian (38) is readily diagonalised to yield
$$\widehat{H}=\underset{\overline{n}}{}\widehat{H}_{\overline{n}}=\underset{\overline{n},k,\sigma }{}\left[\mu _{\overline{n}}2T\mathrm{cos}(k)\right]\widehat{n}_{\overline{n},k,\sigma }$$
(39)
with the same procedure followed in the strong coupling limit in the Hubbard Hamiltonian.
### Rotational Hubbard Hamiltonian
As the last example, we derive a third Hamiltonian that may give the reader a first insight on the increasing complexity if higher single-particle levels are taken into account. Here we consider a situation where we allow $`J_\alpha =0,1/2`$ (but always $`n_\alpha =0`$, leading to $`T_{n_\alpha ,i_\alpha ,n_\beta ,i_{\alpha +1}}=T`$ as already pointed out)
$$\begin{array}{c}H_{2level}=\underset{i,\sigma }{}\underset{a=1}{\overset{1}{}}\left[\lambda _an_{i,a,\sigma }T\left(c_{i,a,\sigma }^{}c_{i+1,a,\sigma }+c_{i+1,a,\sigma }^{}c_{i,a,\sigma }\right)\right]+\hfill \\ \hfill \underset{i}{}\underset{a,b,c,d}{}U_{a,b,c,d}\underset{\sigma ,\sigma ^{}}{}c_{i,a,\sigma }^{}c_{i,b,\sigma ^{}}^{}c_{i,c,\sigma ^{}}c_{i,d,\sigma }\end{array}$$
(40)
where $`U_{a,b,c,d}=U_{a,i_\alpha ;b,i_\beta ;c,i_\gamma ;d,i_\delta }`$. The label $`a`$ (as well as $`b`$, $`c`$, and $`d`$) has been introduced to represent the triplet of harmonic-oscillator numbers $`(n_\alpha ,J_\alpha ,m_\alpha )`$, hence $`\alpha =\{a,i_\alpha ,\sigma _\alpha \}`$. One should recall that originally $`\alpha =(n_\alpha ,J_\alpha ,m_\alpha ,i_\alpha ,\sigma _\alpha )`$. Here, however, it is convenient to write in an explicit way both spin indices $`\sigma _\alpha `$’s and site indices $`i_\alpha `$’s.
The triplet $`a=(n_\alpha ,J_\alpha ,m_\alpha )`$ is such that the value $`a=0`$ corresponds to $`(0,0,0)`$, $`a=1(0,1/2,+1/2)`$ and $`a=1(0,1/2,1/2)`$. The axial quantum number $`n_x`$ has been “frozen” to 0 due to the disk-shaped potential form (i.e. $`\omega _x\omega _{}`$) while the radial quantum number $`J`$ has been limited to the values $`\{0,1/2\}`$ as a first approximation beyond the $`J=0`$ (Hubbard Hamiltonian see (37)). The present model thus enriches the dynamical scenario by introducing modes that takes into account the simplest possible rotational processes for fermions confined in a well.
The wealth of the scenario depicted in Eq. (40) arises from the level-dependence of the interaction coefficient $`U_{a,b,c,d}`$. In fact $`U_{a,b,c,d}`$, as a function of the energy levels may provide a useful tool to simplify Eq. (40) hinting the best strategy for both numerical and analytical analysis of this model. Two main aspects concerning these coefficients are worth repeating here: a) the $`m_\alpha `$-conserving nature of the interaction, related to the symmetry properties of the confining potential and of the interaction coefficient (see section III.1), reduces the number of possible processes; b) the symmetry properties of $`U_{\alpha ,\beta ,\gamma ,\delta }`$ (see Eq. (35)) allow the grouping of interaction terms, accordingly to what has been done in table 1.
From a general point of view, in Hamiltonian (40) the hopping factor may be construed as a multichannel tunnelling coefficient, where the radial quantum numbers identify the channel label, in the same spirit of Hamiltonian (39). Incidentally, this is true if axial degrees of freedom are “frozen” to $`n_\alpha =n_\beta =0`$, otherwise there is tunnelling among levels with $`n_\alpha n_\beta `$, for some $`\alpha `$ and $`\beta `$.
Hamiltonian (40) can represent a situation where single traps are loaded with a small number of atoms, as to fill the first two radial levels of the local harmonic oscillator. To experimentally obtain one of the different simplified Hamiltonians – like (40)– it is necessary to have control of four experimental parameters: laser intensity, angle between counterpropagating lasers, axial magnetic trapping frequency, scattering length. With these parameters it is possible to gain full knowledge of “lattice constant”, interaction parameter, shape and depth of the 3D harmonic traps. The most critical point seems the few-atoms loading of the trap but a technique involving a 3D anisotropic array –a sort of 2D array of 1D arrays– might overcome the problem.
In this picture, the interaction coefficient $`U`$ can then be used as a source of entanglement between different channels. Moreover the possibility of experimental control of the scattering length, and thus of the interaction term, via an applied magnetic field may provide an useful tool of external manipulation of the state of the system in the rich scenario here depicted.
To outline future paths of research, we will here sketch a way to set up a mean-field procedure for the Rotational Hubbard Hamiltonian. The main interest of this approach resides in the possibility of a general discussion of some features of the model which have a direct experimental relevance. For example it is possible to state that, according to what is usually affirmed in the literature Bach et al. (1994) , no BCS-like ground-state is possible for repulsive interaction, while for an attractive two-body potential a paired ground state is possible. The flexibility of experimental techniques involved in the study of ultracold atom physics allows to envisage experimental conditions where these two different regimes are attained. For example exploiting a Feshbach resonance it is possible to drive the scattering length $`a_s`$ from positive to negative values leading thus the system through a quantum phase transition.
The analytic procedure adopted hereafter deeply relies on the concept of quasi-free state Bach et al. (1994). In our situation the following definition of quasi-free state can be adopted:
* all correlation functions can be computed from Wick’s theorem;
* four fermionic expectation values over a quasi-free state have the form
$`<\varphi |e_1e_2e_3e_4|\varphi >=<\varphi |e_1e_2|\varphi ><\varphi |e_3e_4|\varphi >`$
$`<\varphi |e_1e_3|\varphi ><\varphi |e_2e_4|\varphi >+<\varphi |e_1e_4|\varphi ><\varphi |e_2e_3|\varphi >`$
with $`e_i=c_i,c_i^{}`$. In particular we would like point out how the three terms on the right hand-side will lead to the direct, the exchange and pairing energy term of a Hatree-Fock-Bogoliubov mean field Hamiltonian, which, for our RHH becomes
$`\widehat{H}_{2level}^{HFB}`$ $`=`$ $`\widehat{H}_0+{\displaystyle \underset{\begin{array}{c}i,a,b,c,d,\\ \sigma \sigma ^{}\end{array}}{}}U_{a,b,c,d}[\chi _{i,a,\sigma ,c,\sigma }\widehat{c}_{i,b,\sigma ^{}}^{}\widehat{c}_{i,d,\sigma ^{}}`$
$`+`$ $`\chi _{i,d,\sigma ^{},b,\sigma ^{}}\widehat{c}_{i,a,\sigma }^{}\widehat{c}_{i,c,\sigma }\chi _{i,c,\sigma ,b,\sigma ^{}}\widehat{c}_{i,a,\sigma }^{}\widehat{c}_{i,d,\sigma ^{}}`$
$``$ $`\chi _{i,d,\sigma ^{},a,\sigma }\widehat{c}_{i,b,\sigma ^{}}^{}\widehat{c}_{i,c,\sigma }+\xi _{i,b,\sigma ^{},a,\sigma }^{}\widehat{c}_{i,d,\sigma ^{}}\widehat{c}_{i,c,\sigma }`$
$`+`$ $`\xi _{i,c,\sigma ,d,\sigma ^{}}\widehat{c}_{i,a,\sigma }^{}\widehat{c}_{i,b,\sigma ^{}}^{}]`$
with
$`\widehat{H}_0={\displaystyle \underset{i,a,\sigma }{}}[\lambda _a\widehat{n}_{i,a,\sigma }+T(\widehat{c}_{i+1,a,\sigma }^{}\widehat{c}_{i,a,\sigma }+h.c.)]`$
$`\chi _{i,a,\sigma ,b,\sigma ^{}}=<\varphi _{HFB}|\widehat{c}_{i,a,\sigma }^{}\widehat{c}_{i,b,\sigma ^{}}|\varphi _{HF}>`$
$`\xi _{i,a,\sigma ,b,\sigma ^{}}=<\varphi _{HFB}|\widehat{c}_{i,a,\sigma }\widehat{c}_{i,b,\sigma ^{}}|\varphi _{HFB}>`$
The set of generators $`\{\widehat{c}_\alpha ^{}\widehat{c}_\beta \frac{1}{2}\delta _{\alpha \beta }(1\alpha \beta r),`$ $`\widehat{c}_\alpha \widehat{c}_\beta ,\widehat{c}_\alpha ^{}\widehat{c}_\beta ^{}(1\alpha \beta r)\}`$ obeys the following commutation relations
$`[\widehat{c}_i^{}\widehat{c}_j{\displaystyle \frac{1}{2}}\delta _{ij},\widehat{c}_k^{}\widehat{c}_l{\displaystyle \frac{1}{2}}\delta _{kl}]=\delta _{jk}(\widehat{c}_i^{}\widehat{c}_l{\displaystyle \frac{1}{2}}\delta _{il})`$
$`[\widehat{c}_i^{}\widehat{c}_j{\displaystyle \frac{1}{2}}\delta _{ij},\widehat{c}_k^{}\widehat{c}_l^{}]=\delta _{jk}\widehat{c}_i^{}\widehat{c}_l^{}\delta _{jl}\widehat{c}_i^{}\widehat{c}_k^{}`$
$`[\widehat{c}_i\widehat{c}_j,\widehat{c}_k^{}\widehat{c}_l^{}]=\delta _{ik}(\widehat{c}_i^{}\widehat{c}_j{\displaystyle \frac{1}{2}}\delta _{ij})+`$
$`\delta _{ij}(\widehat{c}_k^{}\widehat{c}_i1/2\delta _{ki})\delta _{li}(\widehat{c}_k^{}\widehat{c}_j1/2\delta _{kj})`$
$`\delta _{ki}(\widehat{c}_l^{}\widehat{c}_i1/2\delta _{li})`$ (41)
allowing to state that the dynamical algebra of this new Hamiltonian, which is now quadratic in terms of $`\widehat{c}_i,\widehat{c}_i^{}`$ can be easily recognised to be $`so(2r)`$ Zhang et al. (1990).
Having determined the dynamical algebra of the model Hamiltonian, enables us – at least in principle- to find the ground state of the system with a straightforward procedure. As it will be clear form the subsequent discussion, the main difficulties arise as the number the generators of the $`so(2r)`$ algebra grows with $`r(2r1)`$. For instance for the two-site, $`J=0,1/2`$ model, the Hamiltonian dynamical algebra will have 276 generators.
In spite of the technical difficulties (both analytical and numerical), it is appropriate to apply algebraic techniques to diagonalise $`\widehat{H}_{2level}^{HFB}`$. As stated before, this general approach will give some insight to the ground state properties of the system. If we consider a unitary transformation $`gSO(2r)`$ we can write
$$\widehat{H}_d=g\widehat{H}_{HFB}g^1$$
(42)
where $`\widehat{H}_d`$ is diagonal. As a direct consequence the ground state $`|\varphi _{HF}>`$ of $`\widehat{H}_{2level}^{HFB}`$ can be written as
$$|\varphi _{HF}>=g|0,0,0,\mathrm{},0,0>=g|0>$$
(43)
where $`|0>`$ can be defined as the Bogoliubov particle vacuum (ground state of $`\widehat{H}_d`$). Following Zhang et al. (1990), $`|0>`$ represents a possible choice for the extremal state for the $`SO(2r)`$ group with $`U(r)`$ as the corresponding maximum stability subgroup. Leading to
$$g|0>=\mathrm{\Omega }h|0>\mathrm{\Omega }|0>e^{i\varphi (h)}$$
(44)
where:
$$\mathrm{\Omega }=exp\underset{1\alpha \beta r}{}(\eta _{\alpha ,\beta }\widehat{c}_\alpha ^{}\widehat{c}_\beta ^{}H.c.)\frac{SO(2r)}{U(r)}.$$
(45)
The phase appearing in Eq.(44) has no relevance for our purposes, as we are interested in the evaluation of observable expectation values.
The problem mentioned above about the size of the dynamical algebra, appears here with all its implications. It is necessary to exponentiate the operator $`_{1\alpha \beta r}(\eta _{\alpha ,\beta }\widehat{c}_\alpha ^{}\widehat{c}_\beta ^{}H.c.)`$ which is a $`2r\times 2r`$ matrix in the faithful matrix representation.
Nevertheless, for a repulsive two-body potential, the pairing term can be neglected Bach et al. (1994), thus the dynamical algebra of the system becomes $`U(r)`$. Following Zhang et al. (1990), we can express the Hamiltonian ground state as
$$|\varphi _{HF}>=\mathrm{exp}\underset{\begin{array}{c}k+1\alpha r\\ 1jk\end{array}}{}(\eta _{\alpha ,\beta }\widehat{c}_\alpha ^{}\widehat{c}_\beta H.c.)|0>$$
(46)
where
$$|0>=|\underset{k}{\underset{}{1,1,\mathrm{},1}},0,\mathrm{},0>$$
(47)
which is, in fact, the ground state of the non interacting Hamiltonian. It is worth noticing that this general procedure can be greatly simplified if further constraints, related to symmetries of the problem, are imposed onto the coefficients $`\eta _{\alpha ,\beta }`$. For example, if we consider the two-site($`A`$,$`B`$), $`J=0,1/2`$ case, due to equation (36) the matrix $`\overline{\eta }`$ with elements $`\eta _{\alpha ,\beta }`$ will have the form
$$\overline{\eta }=\left(\begin{array}{cccccccccccc}0& \eta _{1,2}& \eta _{1,3}& \eta _{1,4}& \eta _{1,5}& \eta _{1,6}& \eta _{1,7}& 0& 0& 0& 0& 0\\ \eta _{1,2}& 0& \eta _{2,3}& \eta _{2,4}& \eta _{2,5}& \eta _{2,6}& 0& \eta _{2,8}& 0& 0& 0& 0\\ \eta _{1,3}& \eta _{2,3}& 0& \eta _{3,4}& \eta _{3,5}& \eta _{3,6}& 0& 0& \eta _{3,9}& 0& 0& 0\\ \eta _{1,4}& \eta _{2,4}& \eta _{3,4}& 0& \eta _{4,5}& \eta _{4,6}& 0& 0& 0& \eta _{4,10}& 0& 0\\ \eta _{1,5}& \eta _{2,5}& \eta _{3,5}& \eta _{4,5}& 0& \eta _{5,6}& 0& 0& 0& 0& \eta _{5,11}& 0\\ \eta _{1,6}& \eta _{2,6}& \eta _{3,6}& \eta _{4,6}& \eta _{5,6}& 0& 0& 0& 0& 0& 0& \eta _{6,12}\\ \eta _{1,7}& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& \eta _{2,8}& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& \eta _{3,9}& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& \eta _{4,10}& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& \eta _{5,11}& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& \eta _{6,12}& 0& 0& 0& 0& 0& 0\end{array}\right)$$
(48)
thus impressively reducing the computational effort needed to evaluate the exponential in equation (45). In Eq. (48) we have have assumed the following convention
$`\{n_\alpha =0,J_\alpha =0,m_\alpha =0,i_\alpha =A,\sigma _\alpha =\}1`$
$`\{n_\alpha =0,J_\alpha =0,m_\alpha =0,i_\alpha =A,\sigma _\alpha =\}2`$
$`\mathrm{}`$
$`\{n_\alpha =0,J_\alpha =0,m_\alpha =0,i_\alpha =B,\sigma _\alpha =\}7`$
$`\{n_\alpha =0,J_\alpha =0,m_\alpha =0,i_\alpha =B,\sigma _\alpha =\}8`$
$`\mathrm{}`$
## V Conclusions
In this paper we have investigated the complex structure of fermion interactions for a fermion gas distributed in a linear periodic array of potential wells. Based on the standard many-fermion quantum field theory endowed with a potential distribution mimicking a realistic experimental setup, we have calculated analytically the hopping and interaction coefficients that describe the interactions of fermions within a generalised multimode Hubbard Hamiltonian. Their dependence on the external controllable parameters (such as laser intensity, magnetic trap frequency, wavelength, and scattering length) has been determined.
Our analysis shows that, except for two particularly simple cases (the gas of spin unpolarised fermions and the gas of noninteracting spin polarised fermions), models with different degree of complexity can be derived depending on the interaction processes one decides to account for or to neglect \[consider, e. g., that, in principle, one might introduce an unlimited number of (local) rotational levels\]. In this respect, our simplest nontrivial model (40), which is able to account for the (local) rotational activity of fermions, appears to be far more complex than the Hubbard model or the spin-polarised noninteracting model derived in section IV.
Therefore, the first objective of our future work is to perform a systematic study of model (40). Based on the present analysis and exploiting the interaction-parameter scenario here depicted, the second objective is to recognise the significant regimes characterising the confined fermion gas and to derive the relevant models from Eq. (34).
We would like to stress once again how the analytical knowledge of the coefficients in principle allows us to tailor Hamiltonians performing specific tasks.
An aspect that certainly deserves attention is the study of the zero-temperature phase diagram of model (40) (and, more in general, of sufficiently simple –and thus tractable– models derived from the gHH) and of the relevant phenomenology aimed at suggesting new possible experiments. To achieve a reliable description of these systems, several established analytical and numerical approaches (see e.g. Montorsi A. et al. (1987); Montorsi and Penna (1999), and White S. R. (1993); Batrouni G. G. et al. (1995); Freerikcs J.K. and Monien H. (1996), respectively) can be implemented in analogy to what has been done for bosons Batrouni G. G. et al. (1995); Freerikcs J.K. and Monien H. (1996). Moreover in the recent past several authors (see e.g. Zanardi (2004); Somma et al. (2004)) have proposed to use entanglement measures as a quantum phase transition identifier. We think that our model can represent a good test-field for this new approach to quantum-phase transitions.
## Appendix A Tunnelling coefficient calculation
In the following calculation we will fix $`n_\beta n_\alpha `$, without loss of generality, as it can be easily verified.
The integral in Eq. (21) can be decomposed in the sum of three terms
$$\mathrm{\Theta }_1^{n_\alpha ,n_\beta }+\mathrm{\Theta }_2^{n_\alpha ,n_\beta }+\mathrm{\Theta }_3^{n_\alpha ,n_\beta }=𝑑ye^{\frac{(y\tau )^2}{2}}H_{n_\alpha }\left(y\tau \right)f(y)e^{\frac{y^2}{2}}H_{n_\beta }\left(y\right)$$
(49)
with
$$f(y)=\left[\frac{1\mathrm{cos}(2\mathrm{\Omega }y)}{4\mathrm{\Omega }^2}\frac{y^2}{2}\right],$$
integral (49) becomes
$$\mathrm{\Theta }_1^{n_\alpha ,n_\beta }=\frac{dy}{4\mathrm{\Omega }^2}e^{\frac{y^2+(y\tau )^2}{2}}H_{n_\beta }(y\tau )H_{n_\alpha }(y),$$
(50)
$$\mathrm{\Theta }_2^{n_\alpha ,n_\beta }=\frac{dy}{4\mathrm{\Omega }^2}e^{\frac{y^2+(y\tau )^2}{2}}\mathrm{cos}(2\mathrm{\Omega }y)\times H_{n_\alpha }\left(y\right)H_{n_\beta }\left(y\tau \right),$$
(51)
$$\mathrm{\Theta }_3^{n_\alpha ,n_\beta }=\frac{1}{2}𝑑yy^2e^{\frac{y^2+(y\tau )^2}{2}}\times H_{n_\alpha }\left(y\right)H_{n_\beta }\left(y\tau \right),$$
(52)
The substitution $`\zeta =y\tau /2`$ yields
$$\mathrm{\Theta }_1^{n_\alpha ,n_\beta }=C_\mathrm{\Omega }^\tau 𝑑\zeta e^{\zeta ^2}H_{n_\alpha }\left(\zeta +\frac{\tau }{2}\right)H_{n_\beta }\left(\zeta \frac{\tau }{2}\right)$$
(53)
where $`C_\mathrm{\Omega }^\tau =e^{\tau ^2/4}/(4\mathrm{\Omega }^2)`$. We then use the Hermite polynomial identity
$$H_n(x+y)=\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)H_k(x)(2y)^{nk}$$
(54)
to obtain
$$\mathrm{\Theta }_1^{n_\alpha ,n_\beta }=\frac{e^{\frac{\tau ^2}{4}}}{4\mathrm{\Omega }^2}𝑑\zeta e^{\zeta ^2}\underset{l,k=0}{\overset{n_\alpha ,n_\beta }{}}\left(\genfrac{}{}{0pt}{}{n_\alpha }{l}\right)\left(\genfrac{}{}{0pt}{}{n_\beta }{k}\right)\tau ^{n_\alpha +n_\beta (l+k)}(1)^{n_\beta k}H_k(\zeta )H_l(\zeta )$$
(55)
With the orthogonality of Hermite polynomials
$$_{\mathrm{}}^{\mathrm{}}𝑑xH_n(x)H_m(x)e^{x^2}=\delta _{n,m}2^nn!\sqrt{\pi }$$
(56)
we are able to perform the $`\zeta `$ integration
$$\mathrm{\Theta }_1^{n_\alpha ,n_\beta }=\frac{\sqrt{\pi }}{4\mathrm{\Omega }^2}e^{\frac{\tau ^2}{4}}(1)^{n_\beta }\underset{l=0}{\overset{n_\alpha }{}}\left(\genfrac{}{}{0pt}{}{n_\alpha }{l}\right)\left(\genfrac{}{}{0pt}{}{n_\beta }{l}\right)\tau ^{n_\beta +n_\alpha 2l}(2)^ll!.$$
(57)
It is worth noting that the summation extends to $`n_\alpha `$, that is $`\left(\genfrac{}{}{0pt}{}{a}{b}\right)=0`$ if $`a<b`$. From Wünsche (1998) it can be verified that the last summation is related to generalised Laguerre polynomials, giving
$$\mathrm{\Theta }_1^{n_\alpha ,n_\beta }=\frac{\sqrt{\pi }n_\alpha !2^{n_\alpha }}{4\mathrm{\Omega }^2(\tau )^{n_\alpha n_\beta }}e^{\frac{\tau ^2}{4}}L_{n_\alpha }^{n_\beta n_\alpha }\left(\frac{\tau ^2}{2}\right).$$
(58)
We move now to the calculation of $`\mathrm{\Theta }_2^{n_\alpha ,n_\beta }`$which is given by
$$\mathrm{\Theta }_2^{n_\alpha ,n_\beta }=\frac{1}{4\mathrm{\Omega }^2}𝑑y\mathrm{exp}\left[\frac{y^2}{2}\right]H_{n_\alpha }\left(y\right)\mathrm{exp}\left[\frac{(y\tau )^2}{2}\right]H_{n_\beta }\left(y\tau \right)\mathrm{cos}(\mathrm{\Omega }y).$$
(59)
Eq. (59), with the substitution $`\zeta =y\tau /2`$, can be written as
$$\mathrm{\Theta }_2^{n_\alpha ,n_\beta }=\frac{1}{4\mathrm{\Omega }^2}\mathrm{exp}\left(\tau ^2/4\right)𝑑\zeta e^{\zeta ^2}H_{n_\alpha }\left(\zeta +\frac{\tau }{2}\right)H_{n_\beta }\left(\zeta +\frac{\tau }{2}\right)\mathrm{cos}\left(\mathrm{\Omega }\zeta \right).$$
(60)
Again, using Eq.(54) gives
$$\begin{array}{c}\mathrm{\Theta }_2^{n_\alpha ,n_\beta }=\frac{e^{\tau ^2/4}}{4\mathrm{\Omega }^2}\underset{l,k=0}{\overset{n_\alpha ,n_\beta }{}}\left(\genfrac{}{}{0pt}{}{n_\alpha }{l}\right)\left(\genfrac{}{}{0pt}{}{n_\beta }{k}\right)\tau ^{n_\beta +n_\alpha (l+k)}\hfill \\ \hfill \times (1)^{n_\beta k}𝑑\zeta e^{\zeta ^2}H_k(\zeta )H_l(\zeta )\mathrm{cos}\left(2\mathrm{\Omega }\zeta \right)\end{array}$$
(61)
The integral in equation (61) can be interpreted as the real Fourier transform of the function
$$e^{\zeta ^2}H_k(\zeta )H_l(\zeta ).$$
Recalling that
$$\left[f(x)g(x)\right]=\left[f(x)\right]\left[g(x)\right],$$
where $`[]`$ indicates Fourier transform and $``$ convolution product, we obtain
$$\begin{array}{c}\mathrm{\Theta }_2^{n_\alpha ,n_\beta }=\frac{e^{\frac{\tau ^2}{4}}}{4\mathrm{\Omega }^2}\underset{l,k=0}{\overset{n_\alpha ,n_\beta }{}}\left(\genfrac{}{}{0pt}{}{n_\alpha }{l}\right)\left(\genfrac{}{}{0pt}{}{n_\beta }{k}\right)\tau ^{n_\beta +n_\alpha lk}\hfill \\ \hfill \times (1)^{n_\beta k}e\left[\left[e^{\frac{\zeta ^2}{2}}H_k(\zeta )\right]\left[e^{\frac{\zeta ^2}{2}}H_l(\zeta )\right]\right]\end{array}$$
(62)
giving
$$\begin{array}{c}\mathrm{\Theta }_2^{n_\alpha ,n_\beta }=\frac{e^{\tau ^2/4}}{4\mathrm{\Omega }^2}\underset{l,k=0}{\overset{n_\alpha ,n_\beta }{}}\left(\genfrac{}{}{0pt}{}{n_\alpha }{l}\right)\left(\genfrac{}{}{0pt}{}{n_\beta }{k}\right)\tau ^{n_\beta +n_\alpha (l+k)}(1)^{n_\beta k+l}\hfill \\ \hfill \times e\left[i^{k+l}𝑑ϵe^{\frac{ϵ^2+(ϵ2\mathrm{\Omega })^2}{2}}H_k(ϵ)H_l(ϵ2\mathrm{\Omega })\right].\end{array}$$
(63)
The integral on the left-hand side of equation (62) can be solved following the same procedure used for $`\mathrm{\Theta }_1^{n_\alpha ,n_\beta }`$
$$𝑑ϵe^{ϵ^2/2}H_k(ϵ)e^{\left(ϵ2\mathrm{\Omega }\right)^2/2}H_l\left(ϵ2\mathrm{\Omega }\right)=(1)^{kl}\sqrt{\pi }e^\mathrm{\Omega }\left(2\mathrm{\Omega }\right)^{kl}l!2^lL_l^{kl}\left(4\mathrm{\Omega }^2\right)$$
(64)
giving
$$\begin{array}{c}\mathrm{\Theta }_2^{n_\alpha ,n_\beta }=\frac{e^{\frac{\tau ^2}{4}}}{4\mathrm{\Omega }^2}\underset{l,k=0}{\overset{n_\alpha ,n_\beta }{}}\left(\genfrac{}{}{0pt}{}{n_\alpha }{l}\right)\left(\genfrac{}{}{0pt}{}{n_\beta }{k}\right)\tau ^{n_\beta +n_\alpha (l+k)}(1)^{n_\beta }\hfill \\ \hfill \times 2^ll!e\left[i^{k+l}\right]\frac{\sqrt{\pi }}{(2\mathrm{\Omega })^{lk}}L_l^{kl}\left(2\mathrm{\Omega }^2\right)e^{\mathrm{\Omega }^2}.\end{array}$$
(65)
The calculation of $`\mathrm{\Theta }_3^{n_\alpha ,n_\beta }`$ is quite straightforward. Applying twice the identity
$$xH_n(x)=\frac{1}{2}H_{n+1}(x)+nH_{n1}(x)$$
(66)
we can write $`\mathrm{\Theta }_3^{n_\alpha ,n_\beta }`$ as
$$\begin{array}{c}\mathrm{\Theta }_3^{n_\alpha ,n_\beta }=\frac{e^{\frac{\tau ^2}{4}}}{2}d\zeta e^{\zeta ^2}[\frac{1}{4}H_{n_\alpha +2}(\zeta +\frac{\tau }{2})+\frac{2n_\alpha +1}{2}H_{n_\alpha }(\zeta +\frac{\tau }{2})+n_\alpha (n_\alpha 1)\hfill \\ \hfill \times H_{n_\alpha 2}(\zeta +\frac{\tau }{2})\left]H_{n_\beta }\right(\zeta \frac{\tau }{2}).\end{array}$$
(67)
With the same procedure used for $`\mathrm{\Theta }_1^{n_\alpha ,n_\beta }`$, $`\mathrm{\Theta }_3^{n_\alpha ,n_\beta }`$ is given by
$$\begin{array}{c}\mathrm{\Theta }_3^{n_\alpha ,n_\beta }=\frac{(1)^{n_\beta n_\alpha +1}}{2}\sqrt{\pi }2_\alpha ^nn_\alpha !\tau ^{n_\beta n_\alpha }e^{\frac{\tau ^2}{4}}[\frac{(n_\alpha +1)(n_\alpha +2)}{\tau ^2}L_{n_\alpha +2}^{n_\beta n_\alpha 2}(\tau ^2/2)+\hfill \\ \hfill +\frac{\tau ^2}{4}L_{n_\alpha 2}^{n_\beta n_\alpha +2}\left(\frac{\tau ^2}{2}\right)+\frac{2n_\alpha +1}{2}L_{n_\alpha }^{n_\beta n_\alpha }\left(\frac{\tau ^2}{2}\right)].\end{array}$$
(68)
Hence $`T_{\alpha ,\beta }`$ becomes
$$T_{\alpha ,\beta }=\frac{\mathrm{}\omega _x\delta _{J_\alpha ,J_\beta }\delta _{m_\alpha ,m_\beta }\delta _{\sigma _\alpha ,\sigma _\beta }}{\sqrt{2^{n_\alpha +n_\beta +2}n_\alpha !n_\beta !\pi }}\left[\mathrm{\Theta }_1^{n_\alpha ,n_\beta }+\mathrm{\Theta }_2^{n_\alpha ,n_\beta }+\mathrm{\Theta }_3^{n_\alpha ,n_\beta }\right]$$
(69)
with $`\mathrm{\Theta }_1^{n_\alpha ,n_\beta }`$, $`\mathrm{\Theta }_2^{n_\alpha ,n_\beta }`$, and $`\mathrm{\Theta }_3^{n_\alpha ,n_\beta }`$ given by formulas (58), (65), and (68), respectively.
We have here the plot of $`T_{n_\alpha ,n_\beta ,i_\alpha ,i_\alpha +1}`$ as a function of the difference $`i_\alpha i_\beta `$ for values of $`n_\alpha `$ and $`n_\beta `$ ranging from 0 to 2.
The long distance exponential decay is common to all tunnelling coefficients, regardless of the energy level. On the other hand its detailed shape has deep relevance for nearest-neighbours and next-to-nearest-neighbours (i.e. there may be sign changes passing from $`T_{n,m,i,i+1}`$ and $`T_{n,m,i,i+2}`$) as shown in figures (2-4). Another interesting feature is that an extra-term, due to “on site” tunnelling coefficients, bust be added to the harmonic-oscillator energy term.
## Appendix B Interaction coefficients
We provide here the detailed calculation for the interaction term matrix elements. To solve integral (26)
$$U_x=\frac{1}{\pi l_x}\sqrt{\frac{2^{(n_\alpha +n_\beta +n_\gamma +n_\delta )}}{n_\alpha !n_\beta !n_\gamma !n_\delta !}}\times 𝑑xH_{n_\alpha }(x)H_{n_\beta }(x)H_{n_\gamma }(x)H_{n_\delta }(x)e^{2x^2}$$
we exploit again Eq. (54), obtaining
$$\begin{array}{c}U_x=\frac{1}{\pi l_x}\sqrt{\frac{2^{(n_\alpha +n_\beta +n_\gamma +n_\delta )}}{n_\alpha !n_\beta !n_\gamma !n_\delta !}}\left(\genfrac{}{}{0pt}{}{n_\alpha }{i_\alpha }\right)\left(\genfrac{}{}{0pt}{}{n_\beta }{i_\beta }\right)\left(\genfrac{}{}{0pt}{}{n_\gamma }{i_\gamma }\right)\left(\genfrac{}{}{0pt}{}{n_\delta }{i_\delta }\right)H_{i_\alpha }(0)H_{i_\beta }(0)H_{i_\gamma }(0)H_{i_\delta }(0)\hfill \\ \hfill \times d\zeta (2\zeta )^{n_\alpha +n_\beta +n_\gamma +n_\delta (i_\alpha +i_\beta +i_\gamma +i_\delta )}e^{2\zeta ^2}.\end{array}$$
(70)
In previous equation the summation must be intended over four independent of $`s_\theta =0\mathrm{}n_\theta `$ with $`\theta =\alpha ,\beta ,\gamma ,\delta `$. With the substitution
$$𝑑\zeta (2\zeta )^\alpha e^{2\zeta ^2}=\delta _{\alpha ,2}\sqrt{2}^{\alpha 3}\mathrm{\Gamma }\left[(\alpha +1)/2\right]$$
(71)
($`\delta _{\alpha ,2}`$ indicates that $`\alpha `$ must be an even number) Eq. (70) becomes
$$U_x=\frac{1}{\pi l_x}\underset{\overline{s}}{\overset{\overline{n}}{}}\frac{\mathrm{\Xi }(\overline{s})}{\sqrt{2}^{\overline{s}+3}}\mathrm{\Gamma }\left[\frac{\overline{n}\overline{s}+1}{2}\right]\delta _{n,2}$$
(72)
with
$$\overline{a}=\{a_\alpha ,a_\beta ,a_\gamma ,a_\delta \}\text{ 1-norm: }\overline{a}=\underset{\theta }{}a_\theta $$
(73)
and
$$\mathrm{\Xi }\left(\overline{s}\right)=\underset{\theta }{}\frac{1}{\sqrt{n_\theta !}}\left(\genfrac{}{}{0pt}{}{n_\theta }{s_\theta }\right)H_{s_\theta }(0)$$
(74)
where $`\mathrm{\Xi }\left(\overline{s}\right)=0`$ for odd $`i_\theta `$. The $`\delta `$ function in Eq. (72) should be written as $`\delta _{(\overline{n}\overline{s},2}`$. However the condition $`i_\theta =`$even already implies $`\overline{s}=`$even. We are then allowed to write in Eq. (72):
$$\delta _{\overline{n}\overline{s},2}=\delta _{\overline{n},2}.$$
We solve now the radial part of the interaction term integral written in Eq. (29)
$$U_\rho =\frac{d^2\eta }{\pi }_{J_\alpha ,m_\alpha }^{}(\eta )_{J_\gamma ,m_\gamma }^{}(\eta )_{J_\beta ,m_\beta }(\eta )_{J_\delta ,m_\delta }(\eta ),$$
where $`\eta =(\rho ,\varphi )`$ and $`d^2\eta =\rho d\rho d\varphi `$. Following Wünsche (1998), we express $`_{J_\alpha ,m_\alpha }(\rho ,\varphi )`$ in terms of a finite sum:
$$\begin{array}{c}_{J_\alpha ,m_\alpha }(\rho ,\varphi )=e^{2im_\alpha \varphi }e^{\frac{\rho ^2}{l_{}^2}}\sqrt{(J_\alpha +m_\alpha )!(J_\alpha m_\alpha )!}\hfill \\ \hfill \times \frac{1}{\pi }\underset{q_\alpha =|m_\alpha |}{\overset{J_\alpha }{}}\frac{(1)^{J_\alpha q_\alpha }\left(\frac{\rho }{l_{}}\right)^{2q_\alpha }}{(J_\alpha q_\alpha )!(q_\alpha +m_\alpha )!\left(q_\alpha m_\alpha \right)!}\end{array}$$
(75)
Hence the radial part of $`_{J_\alpha ,m_\alpha }(\stackrel{~}{\rho },\varphi )`$ can be written as
$$R_{J_\alpha ,m_\alpha }(\stackrel{~}{\rho })=\frac{e^{\left(\stackrel{~}{\rho }/l_{}\right)^2/2}}{l_{}}\underset{q_\alpha =|m_\alpha |}{\overset{J_\alpha }{}}\mathrm{\Lambda }_\alpha \left(\stackrel{~}{\rho }/l_{}\right)^{2q_\alpha }$$
(76)
where
$$\mathrm{\Lambda }_\alpha =\frac{(1)^{J_\alpha q_\alpha }\sqrt{(J_\alpha +m_\alpha )!(J_\alpha m_\alpha )!}}{(J_\alpha q_\alpha )!(q_\alpha +m_\alpha )!\left(q_\alpha m_\alpha \right)!}.$$
(77)
Substituting Eq. (76) into Eq. (29) we have
$$\begin{array}{c}U_\rho =\frac{2\delta _{m_\alpha +m_\gamma ,m_\beta +m_\delta }}{\pi l_{}^4}\underset{q_\alpha =|m_\alpha |}{\overset{J_\alpha }{}}\underset{q_\beta =|m_\beta |}{\overset{J_\beta }{}}\underset{q_\gamma =|m_\gamma |}{\overset{J_\gamma }{}}\underset{q_\delta =|m_\delta |}{\overset{J_\delta }{}}\hfill \\ \hfill \times \mathrm{\Lambda }_\alpha \mathrm{\Lambda }_\beta \mathrm{\Lambda }_\gamma \mathrm{\Lambda }_\delta _0^{\mathrm{}}𝑑\rho \frac{\rho }{\pi l_{}^2}\left(\frac{\rho }{l_{}}\right)^{2\overline{q}}e^{2\stackrel{~}{\rho }^2}\end{array}$$
(78)
that, with the same notation of Eq. (72), becomes
$$U_\rho =\frac{\delta _{m_\alpha +m_\gamma ,m_\beta +m_\delta }}{\pi l_{}^2}\underset{\overline{q}=\overline{|m|}}{\overset{\overline{J}}{}}\mathrm{\Lambda }(\overline{J},\overline{m},\overline{q})\frac{\mathrm{\Gamma }\left(\overline{q}+\frac{3}{2}\right)}{2^{\overline{q}+3/2}}$$
(79)
with
$$\mathrm{\Lambda }(\overline{J},\overline{m},\overline{q})=\underset{\theta =\alpha ,\beta ,\gamma ,\delta }{}\mathrm{\Lambda }_\theta .$$
(80)
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# 1 Introduction
## 1 Introduction
The advantages of supersymmetry are somewhat obscured in the Ramond-Neveu-Schwarz formalism, as is the case for supersymmetric particle theories when not formulated in superspace. For example, cancellations of divergences are not obvious, and amplitudes with many fermions are difficult to calculate.
Some of these problems were resolved with the Green-Schwarz formalism, but it proved difficult to quantize except in the lightcone gauge, where some manifest supersymmetry is retained in trade for the loss of some manifest Lorentz invariance. (Similar remarks apply to the Casalbuoni-Brink-Schwarz superparticle.) For example, higher-point diagrams of any type are difficult to calculate because longitudinal polarizations and momenta introduce nonlinearities, and in particular cancellation of anomalies (or any $`ϵ`$-tensor contribution) is difficult to check.
Covariant quantization of the Green-Schwarz action was attempted . A class of derivative gauges was introduced that led to a pyramid of ghosts. Counting arguments showed that the conformal anomaly canceled, and summation of ghost determinants agreed with the lightcone result due to the “identity” $`12+3\mathrm{}=1/4`$. Unfortunately, due to a noninvertible transformation the gauge-fixed action found by this method proved not to be invariant under the Becchi-Rouet-Stora-Tyutin transformations derived by the same method . This problem already appeared for the Casalbuoni-Brink-Schwarz superparticle.
In the meantime, an alternative approach to the quantum superparticle was developed , based on adding extra dimensions to the lightcone, a method that had successfully given free gauge-invariant actions for arbitrary representations of the Poincaré group in arbitrary dimensions . This approach directly gave a BRST operator with the right cohomology. Using the relation between this BRST operator and Zinn-Justin-Batalin-Vilkovisky first-quantization , a manifestly supersymmetric classical mechanics action for this superparticle followed, including a BRST-invariant gauge-fixed action . A crucial difference from the previous method was that “nonminimal” fields were required: There was necessarily a “pyramid” of ghosts, not just a linear tower. However, because of a required Fierz identity, the method of adding extra dimensions could not be directly applied to the lightcone Green-Schwarz superstring.
Various alternatives for a manifestly supersymmetric superstring have since been tried; the most successful is the pure spinor formalism . It has proven somewhat more useful than RNS or lightcone GS approaches in calculating tree amplitudes ; its application to loop amplitudes is in progress . If the formalism for all loops is developed, it should provide a simpler proof of finiteness, which previously required a combination of RNS and lightcone GS results (and equivalence of the two approaches). The pure spinor approach has two main shortcomings:
The first problem is the lack of a manifestly supersymmetric (and Lorentz covariant) path-integration measure. This is a problem in all known superspace approaches to first-quantizing superparticles and superstrings. One consequence is that Green functions (or the effective action in the superparticle case) are not manifestly supersymmetric off shell. Another is that gauge fixing the string field theory (with ghost fields) is not simple. We will not address this problem here.
The other problem is that the pure-spinor BRST operator lacks the $`c`$ and $`b`$ ghosts associated with the usual 2D coordinate invariances (and their associated Virasoro constraints). This is directly related to the lack of a corresponding action with worldsheet metric; the action is known only in the conformal gauge. Furthermore, the moduli that are the remnants of the metric in the conformal gauge must be inserted by hand. Another consequence of the lack of these ghosts as fundamental variables is that they must be reconstructed as complicated composite operators for use as insertions in loop diagrams. The (gauge-fixed) action, BRST operator, moduli, and operator insertions are thus separate postulates of the formalism, rather than all following from a gauge-invariant action as in other formalisms.
In this paper we will formulate the superstring with the ghost structure indicated by the original attempt of and the successful treatment of the superparticle in : the usual $`c`$ and $`b`$ ghosts, and a pyramid of spinors labeled by ghost number and generation. The main result is the BRST operator (from which the gauge-invariant action follows), which takes the form
$$Q_{sstring}=U\left(cT+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi |_>\right)U^1$$
(1.1)
with
$$U=e^{{\scriptscriptstyle \stackrel{~}{\theta }D}}e^{i{\scriptscriptstyle R^a}|_>P_a^{(\pm )}}e^{{\scriptscriptstyle (R^{}+\theta \stackrel{~}{\gamma }^{}\theta /2)}|_>b}$$
(1.2)
where $`T`$ is essentially the energy-momentum tensor, $`D`$ and $`P`$ are the usual “covariant derivatives” in the affine Lie algebra of the classical superstring, $`\stackrel{~}{\theta }`$ is a certain linear combination of ghost $`\theta `$’s, $`R^i`$ are certain expressions quadratic in $`\theta `$’s, $`\pi `$ is conjugate to $`\theta `$, $`\stackrel{~}{\gamma }^{}`$ is a ghost partner to the gamma matrices $`\gamma ^a`$ (which act only on $`\theta `$ and $`\pi `$), and $`|_>`$ picks out the ghost contributions. The gauge-fixed Hamiltonian is just $`\{Q,b\}=T`$. The unitary transformations are necessary because they change the Hilbert space, and so cannot be dropped: A simple analog is the BRST operator for the spinning (Dirac) particle in an external gauge field:
$$Q_{Dirac}=e^{c\gamma ^a_a/\gamma ^{}}(\gamma ^2b)e^{c\gamma ^a_a/\gamma ^{}}=\gamma ^2b+\gamma ^{}\gamma ^a_a\frac{1}{2}c(\gamma ^a_a)^2$$
where $`1/\gamma ^{}`$ doesn’t exist on the correct Hilbert space, but cancels when the “unitary” transformation is evaluated.
We begin in section 2 by reviewing the free superparticle, which has already been quantized (and its BRST cohomology checked) in this approach. Because of the similarity of the algebra of super Yang-Mills to that of the superstring , in section 3 we couple this superparticle to external super Yang-Mills superfields. We use an almost identical method to derive the BRST operator for the superstring in section 4. We finish with our conclusions in section 5. (Mathematical details are relegated to the Appendices.)
## 2 Review of free superparticle
We will start from the free super BRST operator derived in . The generic BRST operator for arbitrary fields (massless, or massive by dimensional reduction) is constructed by starting with a representation of the lightcone SO(D$``$2) (which defines the theory) and adding 4 bosonic and 4 fermionic dimensions to obtain a covariant representation, including all auxiliary fields and ghosts. (This is somewhat redundant for bosons, but necessary for fermions.) The resulting generators $`S^{AB}`$ of OSp(D,2$`|`$4) spin carry vector indices $`A,B`$ that are separated into the usual SO(D$``$1,1) indices $`a,b`$ and the rest as
$$A=(+,,a;\mu ,\stackrel{~}{\mu })=(+,,i),\mu =(,)$$
(2.1)
where $`+,`$ belong to an SO(1,1) subgroup and $`\mu ,\stackrel{~}{\mu }`$ to two Sp(2)’s, of which only the diagonal subgroup will be useful. The BRST operator then takes the generic form
$$Q_{free}^{}=\frac{1}{2}c\mathrm{}+S^a_a+S^{}b+S^\stackrel{~}{}(\mathrm{}=^a_a)$$
(2.2)
In the case of the superparticle, the spin operators are
$$S^{AB}=\frac{1}{4}\overline{\eta }\mathrm{\Gamma }^{[A}\mathrm{\Gamma }^{B\}}\eta $$
(2.3)
in terms of self-conjugate variables $`\eta `$, which arose from the usual self-conjugate SO(D$``$2) fermionic spinor of lightcone superspace. We decompose the OSp(D,2$`|`$4) gamma-matrices $`\mathrm{\Gamma }^A`$ in terms of those of the subgroup SO(1,1) and those ($`\gamma `$) of the subgroup OSp(D$``$1,1$`|`$4) as
$$\mathrm{\Gamma }^i=\left(\begin{array}{cc}\gamma ^i& 0\\ 0& \gamma ^i\end{array}\right),\mathrm{\Gamma }^+=\left(\begin{array}{cc}0& I\\ 0& 0\end{array}\right),\mathrm{\Gamma }^{}=\left(\begin{array}{cc}0& 0\\ I& 0\end{array}\right)$$
(2.4)
with (anti)commutation relations
$`\{\gamma ^a,\gamma ^b\}=2\eta ^{ab},\eta ^{ab}=(+++\mathrm{})`$
$`\{\gamma ^a,\gamma ^\mu \}=\{\gamma ^a,\stackrel{~}{\gamma }^\mu \}=0`$ (2.5)
$`[\gamma ^\mu ,\gamma ^\nu ]=[\stackrel{~}{\gamma }^\mu ,\stackrel{~}{\gamma }^\nu ]=2C^{\mu \nu },[\gamma ^\mu ,\stackrel{~}{\gamma }^\nu ]=0`$ (2.6)
where $`C^{\mu \nu }`$ is the Sp(2) metric with convention
$$C^{}=C_{}=i=C^{}=C_{}$$
(2.7)
and we have denoted $`\gamma ^{\stackrel{~}{\mu }}\stackrel{~}{\gamma }^\mu `$ for legibility. The generalization of the fermionic superspace coordinate $`\theta `$ and its conjugate momentum appear through the analogous decomposition
$$\eta =\left(\begin{array}{c}\pi \\ \theta \end{array}\right),\pi =\frac{}{\theta }$$
(2.8)
We begin with a chiral (Weyl) spinor $`\eta `$, and multiplication by any $`\mathrm{\Gamma }`$ changes the chirality: not just $`\mathrm{\Gamma }^a`$ ($`\gamma ^a`$) as usual, but also $`\mathrm{\Gamma }^\pm `$, which shows that $`\pi `$ and $`\theta `$ have opposite chirality (as expected, since they are conjugate), and $`\mathrm{\Gamma }^\mu `$ ($`\gamma ^\mu `$).
This BRST operator is supersymmetric and also has an infinite pyramid of ghosts. To see these ghosts we need to define creation and annihilation operators from $`\gamma ^\mu `$ and $`\stackrel{~}{\gamma }^\mu `$ as follows:
$`\gamma ^\mu =a^\mu +a^\mu `$ , $`\stackrel{~}{\gamma }^\mu =i(a^\mu a^\mu )`$ (2.9)
$`[a^\mu ,a^\nu ]`$ $`=`$ $`C^{\mu \nu }`$ (2.10)
Then $`\theta `$ can be expanded giving the usual physical supersymmetry fermionic coordinate $`\theta _0`$ at the top of the infinite pyramid of ghosts:
$`|{}_{}{}^{p,q}`$ $``$ $`i^{{\scriptscriptstyle \frac{(p+q)(p+q+1)}{2}}}{\displaystyle \frac{1}{\sqrt{p!}\sqrt{q!}}}(a^{})^p(a^{})^q|0`$
$`{}_{p,q}{}^{}|`$ $``$ $`(i)^{{\scriptscriptstyle \frac{(p+q)(p+q+1)}{2}}}{\displaystyle \frac{1}{\sqrt{p!}\sqrt{q!}}}0|(a_{})^p(a_{})^q`$
$`{}_{p,q}{}^{}|{}_{}{}^{r,s}`$ $`=`$ $`\delta _p^r\delta _q^s`$
$`\theta ^{p,q}`$ $``$ $`\theta |{}_{}{}^{p,q}=\theta ^{p,q}`$
$`\pi _{p,q}`$ $``$ $`{}_{p,q}{}^{}|\pi `$ (2.11)
where
$$\theta _0\theta ^{0,0}.$$
(2.12)
A power of $`i`$ has been inserted to make $`\theta ^{p,q}`$ real: The product of $`n`$ real fermions gets a sign $`(1)^{n(n1)/2}`$ under Hermitian conjugation, because of the reverse ordering. The ghost $`a`$’s and $`a^{}`$’s are fermions, because they take fermions to bosons, and vice versa (in contrast to ordinary $`\gamma `$ matrices, which take fermions to fermions). Thus $`\theta ^{p,q}`$ is the product of $`p+q+1`$ fermions, including $`\theta |`$ itself. Then $`\pi _{p,q}`$ is not necessarily Hermitian, but has been defined to give 0 or 1 in graded commutators. (But $`\pi ^{p,q}`$ is always Hermitian, like $`|\pi `$ and $`\pi _0`$.) We will sometimes also use a notation where $`\theta ^{p,q}`$ carries instead $`p`$ $``$’s and $`q`$ $``$’s: For example, $`\theta ^{1,0}\theta ^{}`$. Note that the ghosts alternate in both statistics and chirality with each ghost level.
So in superspace notation the free super BRST operator is
$$Q_{free}^{}=\frac{1}{2}c\mathrm{}\frac{1}{2}\overline{\pi }\gamma ^2\theta b+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi \frac{i}{2}\overline{\pi }\gamma ^{}\text{/}p\theta ,\text{/}pi_a\gamma ^a$$
(2.13)
We can make a unitary transformation on $`Q_{free}^{}`$ to give a convenient form with which to work. Specifically, the unitary transformation
$$Q_{free}=U_0Q_{free}^{}U_0^{}$$
(2.14)
with
$$U_0=e^{\overline{\theta }\stackrel{~}{\gamma }^{}\theta b/2}$$
(2.15)
gives $`Q_{free}`$ in terms of the supersymmetry generator $`q_0`$, spinor covariant derivative $`d_0`$ and all their nonminimal versions:
$`Q_{free}`$ $`=`$ $`\frac{1}{2}c\mathrm{}2\overline{\pi }a^{}a^{}\theta b\frac{i}{2}\overline{q}a^{}d`$ (2.16)
$`q`$ $`=`$ $`\pi \text{/}p\theta ,d=\pi +\text{/}p\theta `$ (2.17)
Actually, $`q_0`$ is the only part of $`q`$ that does not appear in this form of the BRST operator: Because of the creation and annihilation operators, $`\theta _0`$ and $`\pi _0`$ appear only as their supersymmetry invariant combination $`d_0`$. Thus the supersymmetry generator that anticommutes with this form of $`Q`$ is just the usual one $`q_0`$. (This can also be derived in a straightforward way by starting with the lightcone $`q`$.) Then the supersymmetry generator for $`Q_{free}^{}`$ can be obtained by inverting the unitary transformation on $`q_0`$:
$`q_0^{}`$ $`=`$ $`U_0^{}q_0U_0`$ (2.18)
$`=`$ $`\pi _0\text{/}p\theta _0\theta ^{}b`$
## 3 Interacting superparticle
The (D=3,4,6,10) superparticle BRST operator in a super Yang-Mills background (with constant superfield strength) is closely related to the superstring BRST operator. The introduction of the SYM background can be established by gauge covariantizing the super covariant derivatives $`p_a`$ and $`d_{0\alpha }`$:
$`p_a`$ $``$ $`_a`$ (3.1)
$`d_0`$ $``$ $`_0`$ (3.2)
Then the graded algebra among the covariant derivatives is
$`[_a,_b]=F_{ab}`$ (3.3)
$`\{_{0\alpha },_{0\beta }\}=2\gamma _{a\alpha \beta }^a`$ (3.4)
$`[_{0\alpha },_a]=\gamma _{a\alpha \beta }W^\beta `$ (3.5)
The Bianchi identity from the above algebra gives
$$\gamma _{a\alpha \beta }[^a,W^\beta ]=0$$
(3.6)
and the D=3,4,6,10 dimensional gamma matrix (which is symmetric in those cases) identity
$$\gamma _{a(\alpha \beta }\gamma _{}^{a}{}_{\gamma )}{}^{}{}_{}{}^{\delta }=0.$$
(3.7)
We begin at linear order in the fields, where the background satisfies the equations of motion
$`\{_\alpha ,W^\alpha \}=0`$ (3.8)
$`[^a,F_{ab}]=0`$ (3.9)
### 3.1 Constant YM background
One way to build this interacting super BRST operator is by considering an ordinary constant YM background first, and next supersymmetrizing it by including a constant fermionic field strength (not yet superfield) $`\stackrel{}{w}^\alpha `$. Then we extend the result to a nonconstant SYM background in the next subsection.
Making the gauge choice
$$\stackrel{}{A}{}_{a}{}^{}=\frac{i}{2}x^b\stackrel{}{F}_{ba}$$
for constant field strength, the super BRST operator can be written in the form
$$\stackrel{}{Q}{}_{YMB}{}^{}=Q_{free}^{}+\frac{1}{2}\stackrel{}{F}{}_{ab}{}^{}V_{}^{ab}$$
(3.1.1)
We then find
$$V^{ab}=\frac{i}{2}cx^{[a}p^{b]}+\frac{i}{2}R^{}R^{[a}p^{b]}(c+R^{})\overline{\pi }\gamma ^{ab}\theta +\frac{1}{4}(x+R)^{[a}\overline{\pi }\gamma ^{}\gamma ^{b]}\theta $$
(3.1.2)
in terms of an expression $`R^i`$ defined below, where we use the notation
$`C^{[a}D^{b]}`$ $``$ $`C^aD^bC^bD^a`$ (3.1.3)
$`\gamma ^{ab}`$ $``$ $`\frac{1}{4}\gamma ^{[a}\gamma ^{b]}`$ (3.1.4)
We can also write
$`V^{ab}=V^{ab}`$
$`V^{ijk}=i(x^ix^j+R^iR^j)p^k+\frac{1}{2}(x+R)^{(i}\overline{\pi }\gamma ^{j]}\gamma ^k\theta `$
$`x^{}=c,p^{}=0`$ (3.1.5)
(There is further antisymmetrization in the last two indices upon contraction with $`F`$, following the graded symmetrization in the first two indices shown above: The tensor $`V^{ijk}`$ has mixed symmetry.)
The expression $`R^i`$ is given by
$$R^i(\theta )\frac{1}{2}\overline{\theta }𝒪\gamma ^i\theta $$
(3.1.6)
where the operator $`𝒪`$ is defined to satisfy
$`[\gamma ^{},𝒪]`$ $`=`$ $`0`$
$`\{\stackrel{~}{\gamma }^{},𝒪\}`$ $`=`$ $`2\gamma ^{}`$
$`[a^{}a_{}a^{}a_{},𝒪]`$ $`=`$ $`0`$
$`0|𝒪\stackrel{~}{\gamma }^{}`$ $`=`$ $`i0|\stackrel{~}{\gamma }^{}=0|\gamma ^{}`$ (3.1.7)
As an explicit form of $`𝒪`$ we find
$$𝒪=\frac{1}{2}\{\frac{1}{\stackrel{~}{\gamma }^{}},\gamma ^{}\}$$
(3.1.8)
where
$`{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}}`$ $`=`$ $`{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}[\mathrm{\Theta }_{N_{}N_{}}{\displaystyle \frac{N_{}!}{(N_{}+p+1)!}}ia_{}(ia_{}a_{})^p`$ (3.1.9)
$`a^{}(ia^{}a^{})^p\mathrm{\Theta }_{N_{}N_{}}{\displaystyle \frac{N_{}!}{(N_{}+p+1)!}}]`$
with
$$\mathrm{\Theta }_x=\{\begin{array}{cc}1& x0\\ 0& x<0\end{array},N_\mu =a^\mu a_\mu $$
(not summed over $`\mu `$). This representation satisfies (3.1) if we regularize indefinite norm states. (See the Appendices for details.)
### 3.2 Constant SYM background
From this $`\stackrel{}{Q}_{YMB}^{}`$ we can construct a BRST operator for a supersymmetric constant SYM background $`\stackrel{}{Q}_{SYMB}^{}`$ in the form
$$\stackrel{}{Q}{}_{SYMB}{}^{}=Q_{free}^{}+\frac{1}{2}\stackrel{}{F}{}_{ab}{}^{}V_{}^{ab}+\stackrel{}{w}{}_{}{}^{\alpha }V_{\alpha }^{}$$
(3.2.1)
In addition to first-quantized transformations we take $`q_0^{}`$ (2.18) to also generate the second-quantized transformations of $`\stackrel{}{w}^\alpha `$ and $`\stackrel{}{F}_{ab}`$
$`\{q_{0\beta }^{},\stackrel{}{w}{}_{}{}^{\alpha }\}`$ $`=`$ $`\gamma _{}^{ab}{}_{\beta }{}^{}{}_{}{}^{\alpha }\stackrel{}{F}_{ab}`$ (3.2.2)
$`\{q_0^{},\stackrel{}{F}{}_{ab}{}^{}\}`$ $`=`$ $`0`$ (3.2.3)
so that they cancel up to a gauge transformation (generated by $`Q_{free}^{}`$):
$`\{q_0^{},\stackrel{}{Q}{}_{SYMB}{}^{}\}`$ $`=`$ $`\{Q_{free}^{},\mathrm{\Psi }\}`$ (3.2.4)
We then have
$$\frac{1}{2}\stackrel{}{F}{}_{ab}{}^{}\{q_{0\beta }^{},V^{ab}\}\stackrel{}{w}{}_{}{}^{\alpha }[q_{0\beta }^{},V_\alpha ]=\stackrel{}{F}{}_{ab}{}^{}\gamma _{}^{ab}{}_{\beta }{}^{}{}_{}{}^{\alpha }V_\alpha +\{Q_{free}^{},\mathrm{\Psi }_\beta \}$$
(3.2.5)
This is true if we define $`V_\alpha `$ by
$$\{q_{0\beta }^{},V^{ab}\}|_{\gamma ^{ab}}2\gamma ^{ab}{}_{\beta }{}^{}{}_{}{}^{\alpha }V_{\alpha }^{}$$
(3.2.6)
which means we define $`V_\alpha `$ from the left-hand side by selecting only terms with an explicit $`\gamma ^{ab}`$.
With this definition we find $`V_\alpha `$ and $`\mathrm{\Psi }_\alpha `$
$`V_\alpha `$ $``$ $`(c+R^{})q_{0\alpha }^{}`$ (3.2.7)
$`\mathrm{\Psi }_\alpha `$ $``$ $`\frac{i}{2}(x^b+R^b)(\gamma _{\alpha \beta }^a\theta _0^\beta \stackrel{}{F}{}_{ab}{}^{}2\gamma _{b\alpha \beta }\stackrel{}{w}{}_{}{}^{\beta })+\frac{1}{2}(\gamma _{\alpha \beta }^b\theta _0^\beta +\gamma _{\alpha \beta }^b\stackrel{~}{\theta }_0^\beta )iR^a\stackrel{}{F}{}_{ab}{}^{}`$ (3.2.8)
where
$$\stackrel{~}{\theta }=i0|𝒪|\theta =0|e^{ia_{}a_{}}1|\theta $$
(3.2.9)
which contains all nonminimal ghost-number-zero ghosts.
To obtain a gauge independent and explicitly supersymmetric expression we perform a unitary transformation
$$U_1=e^\mathrm{\Lambda }$$
where
$`\mathrm{\Lambda }`$ $`=`$ $`iR^b(\theta _0\gamma _b\stackrel{}{w}+\stackrel{~}{\theta }\gamma _b\stackrel{}{w}+\frac{1}{2}\theta _0\gamma _b\gamma ^{ac}\theta _0\stackrel{}{F}{}_{ac}{}^{}+\frac{1}{2}\stackrel{~}{\theta }\gamma _b\gamma ^{ac}\theta _0\stackrel{}{F}_{ac}`$ (3.2.10)
$`+\frac{1}{2}\stackrel{~}{\theta }\gamma _b\gamma ^{ac}\stackrel{~}{\theta }\stackrel{}{F}{}_{ac}{}^{}\frac{1}{2}\theta _0\gamma ^c\stackrel{~}{\theta }\stackrel{}{F}{}_{bc}{}^{})`$
$`\theta _0\gamma ^b\stackrel{~}{\theta }(\frac{1}{3}\theta _0\gamma _b\stackrel{}{w}+\frac{2}{3}\stackrel{~}{\theta }\gamma _b\stackrel{}{w}+\frac{1}{4}\theta _0\gamma _b\gamma ^{ac}\theta _0\stackrel{}{F}_{ac}`$
$`+\frac{3}{4}\stackrel{~}{\theta }\gamma _b\gamma ^{ac}\theta _0\stackrel{}{F}_{ac}+\frac{5}{12}\stackrel{~}{\theta }\gamma _b\gamma ^{ac}\stackrel{~}{\theta }\stackrel{}{F}_{ac})`$
After another unitary transformation $`U_0`$ (2.15), $`\stackrel{}{Q}_{SYMB}^{}`$ becomes (at the linearized level)
$`\stackrel{}{Q}_{SYMB}`$ $`=`$ $`\frac{1}{2}(c+R^{}+\frac{1}{2}\overline{\theta }\stackrel{~}{\gamma }^{}\theta )|_>(\mathrm{}+W_0\overline{\pi }\gamma ^{ab}\theta |_>F_{ab})`$ (3.2.11)
$``$ $`2\overline{\pi }a^{}a^{}\theta |_>b+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi |_>`$
$``$ $`\frac{i}{2}(_a+\stackrel{~}{\theta }\gamma _aW+\frac{1}{2}\stackrel{~}{\theta }\gamma _a\{W,_0\}\stackrel{~}{\theta })\overline{\pi }\gamma ^{}\gamma ^a\theta |_>`$
$``$ $`\frac{1}{2}\overline{\pi }^{}_0\frac{1}{2}(_a+\stackrel{~}{\theta }\gamma _aW+\frac{1}{2}\stackrel{~}{\theta }\gamma _a\{W,_0\}\stackrel{~}{\theta }\frac{i}{2}R^b|_>F_{ab})\overline{\theta }^{}\gamma ^a_0`$
$``$ $`\frac{1}{3}\overline{q}^{}\gamma ^a\stackrel{~}{\theta }\stackrel{~}{\theta }\gamma _aW+\frac{5}{24}\overline{q}^{}\gamma ^a\stackrel{~}{\theta }\stackrel{~}{\theta }\gamma _a\{W,_0\}\stackrel{~}{\theta }`$
$`+`$ $`\frac{i}{4}\overline{q}^{}\gamma ^b\stackrel{~}{\theta }R^a|_>F_{ab}`$
$`+`$ $`\frac{1}{4}[i_bR^b|_>,i_a\overline{\pi }\gamma ^{}\gamma ^a\theta |_>]`$
$`+`$ $`R^{}|_>P_a(\stackrel{~}{\theta }\gamma ^aW+\frac{1}{2}\stackrel{~}{\theta }\gamma ^a\{W,_0\}\stackrel{~}{\theta })`$
$`+`$ $`\frac{1}{23!}[i_cR^c|_>,[i_bR^b|_>,i_a\overline{\pi }\gamma ^{}\gamma ^a\theta |_>]]`$
where $`|_>`$ means that we drop $`\theta _0`$ contributions, and
$`\mathrm{}`$ $`=`$ $`^a_a`$ (3.2.12)
$`_a`$ $`=`$ $`p_a+A_a`$ (3.2.13)
$`_{0\alpha }`$ $`=`$ $`\pi _{0\alpha }+(\text{/}p\theta _0)_\alpha +A_\alpha `$ (3.2.14)
$`q^\alpha `$ $`=`$ $`\pi ^\alpha +(\text{/}p\theta ^{})^\alpha `$ (3.2.15)
The superfields have the $`\theta _0`$ expansions
$`F_{ab}`$ $`=`$ $`\stackrel{}{F}_{ab}`$ (3.2.16)
$`W^\alpha `$ $`=`$ $`\stackrel{}{w}{}_{}{}^{\alpha }+(\gamma ^{ab}\theta _0)^\alpha \stackrel{}{F}_{ab}`$ (3.2.17)
$`A_a`$ $`=`$ $`\stackrel{}{A}{}_{a}{}^{}+\overline{\theta }_0\gamma _a\stackrel{}{w}+\frac{1}{2}\overline{\theta }_0\gamma _a\gamma ^{bc}\theta _0\stackrel{}{F}_{bc}`$ (3.2.18)
$`A_\alpha `$ $`=`$ $`(\gamma ^a\theta _0)_\alpha \stackrel{}{A}{}_{a}{}^{}+\frac{2}{3}(\gamma ^a\theta _0)_\alpha \overline{\theta }_0\gamma _a\stackrel{}{w}+\frac{1}{4}(\gamma ^a\theta _0)_\alpha \overline{\theta }_0\gamma _a\gamma ^{bc}\theta _0\stackrel{}{F}{}_{bc}{}^{}`$ (3.2.19)
in the gauge
$`\stackrel{}{A}_a`$ $`=`$ $`\frac{i}{2}x^b\stackrel{}{F}_{ba}`$ (3.2.20)
$`\stackrel{}{A}_\alpha `$ $`=`$ $`0`$ (3.2.21)
used above, but (3.2.11) is manifestly gauge independent and supersymmetric.
### 3.3 Arbitrary SYM background
After making a final unitary transformation
$$U_3=e^{(R^{}+\theta \stackrel{~}{\gamma }^{}\theta /2)|_>b}$$
the above BRST operator can be written in the simple form
$$Q_{SYMB}^{\prime \prime }=U\left[\frac{1}{2}c\left(\mathrm{}+W_0\overline{\pi }\gamma ^{ab}\theta |_>F_{ab}\right)+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi |_>\right]U^1$$
$$U=e^{\stackrel{~}{\theta }_0}e^{iR^a|_>_a}e^{(R^{}+\theta \stackrel{~}{\gamma }^{}\theta /2)|_>b}$$
(3.3.1)
which can be applied directly to the case of an arbitrary, nonlinear SYM background.
In fact, the nilpotence of this BRST operator does not seem to require that the background be on shell. This is contradictory to the usual result that any description of linearized “quantum” Yang-Mills in a Yang-Mills background must have the background on shell, since nonabelian gauge invariance relates kinetic and interaction terms . (Similar remarks apply to any nonabelian gauge theory, such as gravity or strings.) This paradox is probably due to the fact that we have not required an “integrability” condition on the background: For a generic self-interacting field theory, an action (or ZJBV action) of the form
$$S=\frac{1}{2}\varphi ^j\varphi ^iK_{ij}+\frac{1}{6}\varphi ^k\varphi ^j\varphi ^iV_{ijk}+\mathrm{}$$
(3.3.2)
results in the kinetic operator (or BRST operator) in a background
$$Q_{ij}=K_{ij}+\varphi ^kV_{kij}+\mathrm{}$$
(3.3.3)
From $`S`$ we can see that $`K,V,\mathrm{}`$ must be totally (graded) symmetric. In $`Q`$, this condition on $`K`$ is seen to follow simply from hermiticity, but the condition on $`V`$ is not so obvious. Since we are ultimately concerned with the BRST operator for the superstring without background, and are using the SYM case in a background only as an analogy, we will not consider this obscurity further here.
As explained in the Introduction, in the above expression for the BRST operator (3.3.1) we are not allowed to remove the exponential factors, since that would lead to a trivial result. This fact can be understood already in the free case: The BRST operator that would result from dropping the background and exponentials has the wrong cohomology, since the remaining two terms have no dependence on $`\theta _0`$, so one would obtain an ordinary superfield satisfying only the Klein-Gordon equation. In this case the exponentials are required for $`Q`$ to be regularizable: Certain poorly defined quantities cancel upon their expansion. (See Appendices B-C.)
## 4 Superstring
The superstring is described by a 2D field theory whose algebra of covariant derivatives (currents) resembles that of interacting particle covariant derivatives for a constant SYM background:
$`\{D_\alpha ^{(\pm )}(1),D_\beta ^{(\pm )}(2)\}`$ $`=`$ $`2\delta (21)\gamma _{\alpha \beta }^aP_a^{(\pm )}(1)`$
$`[D_\alpha ^{(\pm )}(1),P_a^{(\pm )}(2)]`$ $`=`$ $`2\delta (21)\gamma _{a\alpha \beta }\mathrm{\Omega }^{(\pm )\beta }(1)`$
$`\{D_\alpha ^{(\pm )}(1),\mathrm{\Omega }^{(\pm )\beta }(2)\}`$ $`=`$ $`\pm i\delta ^{}(21)\delta _\alpha ^\beta `$
$`[P_a^{(\pm )}(1),P_b^{(\pm )}(2)]`$ $`=`$ $`\pm i\delta ^{}(21)\eta _{ab}`$
$`[P^{(\pm )},\mathrm{\Omega }^{(\pm )}]`$ $`=`$ $`\{\mathrm{\Omega }^{(\pm )},\mathrm{\Omega }^{(\pm )}\}=0`$ (4.1)
where
$`D_\alpha ^{(\pm )}`$ $`=`$ $`\pi _{0\alpha }+(\gamma ^a\theta _0)_\alpha \widehat{P}_a^{(\pm )}\pm i\frac{1}{2}(\gamma ^a\theta _0)_\alpha \theta _0\gamma _a\theta _0^{}`$
$`P_a^{(\pm )}`$ $`=`$ $`\widehat{P}_a^{(\pm )}\pm i\theta _0\gamma _a\theta _0^{}`$
$`\mathrm{\Omega }^{(\pm )\alpha }`$ $`=`$ $`\pm i\theta _0^{}`$ (4.2)
and
$$\widehat{P}^{(\pm )}=\frac{1}{\sqrt{2}}\left(i\frac{\delta }{\delta X}\pm X^{}\right)$$
(4.3)
in the Hamltonian formalism correspond to the left(right)-moving combinations of $`P_0`$ and $`P_1`$ of the first-order formalism after using the equation of motion for $`P_1`$ (see below). (In the definitions above, $`(\pm )`$’s on $`\pi `$ and $`\theta `$ are understood.) Also, means a $`\sigma `$ derivative as usual. $`D,P,\mathrm{\Omega }`$ (anti)commute with the supersymmetry generator
$$q_\alpha =q_\alpha ^{(+)}+q_\alpha ^{()},q_\alpha ^{(\pm )}=\pi _{0\alpha }(\gamma ^a\theta _0)_\alpha \widehat{P}_a^{(\pm )}i\frac{1}{6}(\gamma ^a\theta _0)_\alpha \theta _0\gamma _a\theta _0^{}$$
(4.4)
So we can see the analogy between the covariant derivatives of the free superstring and the superparticle with SYM background.
$$(D_\alpha ,P_a,\mathrm{\Omega }^\alpha )(_\alpha ,_a,W^\alpha )$$
(4.5)
as well as the less precise analogy
$${}_{}{}^{}F_{ab}$$
(4.6)
### 4.1 BRST
Now we can guess the result for the superstring BRST operator from the result of the superparticle in a constant SYM background:
$$Q_{sstring}=U\left(cT+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi |_>\right)U^1$$
(4.1.1)
where
$$U=e^{{\scriptscriptstyle \stackrel{~}{\theta }D}}e^{{\scriptscriptstyle iR^a}|_>P_a}e^{{\scriptscriptstyle (R^{}+\theta \stackrel{~}{\gamma }^{}\theta /2)}|_>b}$$
(4.1.2)
and $`Q_{sstring}=Q_{sstring}^{(+)}+Q_{sstring}^{()}`$. From now on we will suppress the $`\sigma `$-integral symbol for convenience. Since $`\theta _0`$ and $`\pi _0`$ appear only in $`D`$, $`P`$, and $`T`$, this $`Q`$ is automatically supersymmetric under the above supersymmetry generator.
There are two major differences in $`T`$ as compared to the superparticle: Firstly the string has a $`c^{}b`$ ghost contribution. Secondly the superstring has $`\mathrm{\Omega }D`$ as an analog of $`W_0\overline{\pi }\gamma ^{ab}\theta F_{ab}`$, from the correspondence above. So our trial form of $`T`$ is
$`T^{(\pm )}`$ $`=`$ $`\frac{1}{2}\mathrm{}_\pm i\{c^{}b+\overline{\theta }^{}\pi +w^\pm (\overline{\theta }\pi )^{}`$ (4.1.3)
$`+A_1^\pm \left[\overline{\theta }(a^{}a_{}a^{}a_{})\pi \right]^{}+A_2^\pm \left[\overline{\theta }(a^{}a_{}+a^{}a_{})\pi \right]^{}\}`$
where $`\mathrm{}_\pm =\widehat{P}^{(\pm )a}\widehat{P}_a^{(\pm )}`$. (The true energy-momentum tensor is actually $`Ti(cb)^{}`$.) The constants $`w^\pm `$$`A_1^\pm `$ and $`A_2^\pm `$ will be determined by 3 conditions: (1) The conformal weight of $`\overline{\pi }\stackrel{~}{\gamma }^{}\pi `$ should be $`1`$. (2) The conformal anomaly should cancel in $`D=10`$. (3) $`\theta _0`$ should have conformal weight $`0`$ due to supersymmetry. (The $`A_1`$ term is ghost number, while the $`A_2`$ term is ghost level.)
Satisfying these constraints we find
$`A_1^\pm `$ $`=`$ $`1,A_2^\pm =w^\pm =0`$
$`T^{(\pm )}`$ $`=`$ $`\frac{1}{2}\mathrm{}_\pm i\left\{c^{}b+\overline{\theta }^{}\pi +\left[\overline{\theta }(a^{}a_{}a^{}a_{})\pi \right]^{}\right\}`$ (4.1.4)
and the gauge-fixed Hamiltonian is $`\{Q,b^{(+)}+b^{()}\}=T^{(+)}+T^{()}`$.
This $`Q`$ has four interesting quantum numbers:(1) ghost number; (2) conformal weight, which is “momentum number” (1 for $`P,b,\pi ,^{}`$) minus ghost number; (3) (10D) engineering dimension ($`1`$ for $`x,c`$, $`\frac{1}{2}`$ for $`\theta `$, 2 for ); and (4) a mysterious “field weight”, which is 1 for all fields, but for which we attribute a 1 for $`\stackrel{~}{\gamma }^{}`$. (Thus, $`Q`$ is quadratic in momenta and primes, and cubic in fields and $`\stackrel{~}{\gamma }^{}`$’s.)
From now on let’s concentrate on one chirality. After expanding the exponential factor and regularizing (as explained in Appendix C) we find
$`Q_{sstring}^{(+)}`$ $`=`$ $`\left(c+R^{}+\frac{1}{2}\overline{\theta }\stackrel{~}{\gamma }^{}\theta \right)|_>`$ (4.1.5)
$`\times \left(\frac{1}{2}\mathrm{}_+ic^{}bi\overline{\theta }^{}\pi i[\overline{\theta }(a^{}a_{}a^{}a_{})\pi ]^{}\right)`$
$``$ $`2\overline{\pi }a^{}a^{}\theta |_>b+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi |_>`$
$``$ $`\frac{i}{2}\left(\widehat{P}_a+i\theta _0\gamma _a\theta _0^{}+2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{1}{2}R_a^{}|_>\right)\overline{\pi }\gamma ^{}\gamma ^a\theta |_>`$
$``$ $`\frac{1}{2}\overline{\pi }^{}\left(\pi _0+(\gamma ^a\theta _0)\widehat{P}_a+i\frac{1}{2}(\gamma ^a\theta _0)\theta _0\gamma _a\theta _0^{}\right)`$
$``$ $`\frac{1}{2}\left(\widehat{P}_a+i\theta _0\gamma _a\theta _0^{}+2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{1}{2}R_a^{}|_>\right)`$
$`\times \overline{\theta }^{}\gamma ^a\left(\pi _0+(\gamma ^b\theta _0)\widehat{P}_b+i\frac{1}{2}(\gamma ^b\theta _0)\theta _0\gamma _b\theta _0^{}\right)`$
$``$ $`\frac{1}{3}\overline{q}^{}\gamma ^a\stackrel{~}{\theta }\left(2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\frac{5}{4}\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{3}{4}R_a^{}|_>\right)`$
$`+`$ $`\frac{1}{2}\left(\widehat{P}_a+i\theta _0\gamma _a\theta _0^{}+2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{1}{2}R_a^{}|_>\right)^2R^{}|_>`$
$`+`$ $`ic^{}b\left(R^{}+\frac{1}{2}\overline{\theta }\stackrel{~}{\gamma }^{}\theta \right)|_>`$
This $`Q`$ satisfies $`Q^2=0`$, as can be checked directly.
### 4.2 Constraints
The constraints of the gauge-invariant action (see following subsection) can be obtained directly from the BRST operator by taking its (graded) commutator and keeping just ghost-number-zero terms: The Virasoro constraints $`𝒜`$ follow as usual from $`b`$ (with the gauge-fixed action from $`b`$), while generalizations $``$ of the $`\gamma pd`$ constraint ($`\kappa `$ symmetry generator) follow from $`\theta ^{p,p+1}`$, and first-class generalizations $``$ of the second-class constraint $`d`$ follow from $`\pi ^{p,p+1}`$ :
$`𝒜`$ $`=`$ $`\frac{1}{2}\mathrm{}_+i{\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}\overline{\theta }^{q,q}\pi _{q,q}`$ (4.2.1)
$`_0`$ $`=`$ $`\gamma ^a\mathrm{\Pi }𝒫_a(\text{d}_0+\pi ^{1,1})2\theta ^{1,1}𝒜+2\vartheta ^0(\frac{1}{2}𝒫^2+𝒜)`$ (4.2.2)
$`_p`$ $`=`$ $`\gamma ^a\mathrm{\Pi }𝒫_a(\pi ^{p,p}+\pi ^{p+1,p+1})+2(\theta ^{p,p}\theta ^{p+1,p+1})𝒜+2\vartheta ^p(\frac{1}{2}𝒫^2+𝒜)`$ (4.2.3)
$`_0`$ $`=`$ $`\text{d}_0\pi ^{1,1}+\gamma ^a\mathrm{\Pi }𝒫_a\theta ^{1,1}`$ (4.2.4)
$`_p`$ $`=`$ $`\mathrm{\Pi }(\pi ^{p,p}\pi ^{p+1,p+1})+\gamma ^a\mathrm{\Pi }𝒫_a(\theta ^{p,p}+\theta ^{p+1,p+1})`$ (4.2.5)
where
$`𝒫_a`$ $``$ $`\widehat{P}_a+i\theta _0\gamma _a\theta _0^{}+2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{1}{2}\stackrel{~}{R}_a^{}|_>`$
$`\text{d}_0`$ $``$ $`\pi _0+\gamma ^a\widehat{P}_a\theta _0+i\frac{1}{2}(\gamma ^a\theta _0)\theta _0\gamma _a\theta _0^{}+\frac{2}{3}\gamma ^a\stackrel{~}{\theta }\left(2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\frac{5}{4}\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{3}{4}\stackrel{~}{R}_a^{}|_>\right)`$
$`\vartheta ^p`$ $``$ $`2\theta ^{1,1}+2\theta ^{2,2}+\mathrm{}+2\theta ^{p1,p1}+\theta ^{p,p}\theta ^{p+1,p+1}2\theta ^{p+2,p+2}2\theta ^{p+3,p+3}\mathrm{}`$
and $`\stackrel{~}{R}^a`$ indicates that only ghost-number-zero $`\theta `$’s are selected.
Since this procedure requires the component expression, we explicitly use the projection operator
$$\mathrm{\Pi }=\frac{1}{\stackrel{~}{\gamma }^{}}(\stackrel{~}{\gamma }^{})_{reg}$$
(4.2.7)
as, e.g., $`\mathrm{\Pi }|\theta `$, in some terms, as explained in Appendix C. Also, because $`R^{}`$ only interacts with $`\pi \mathrm{\Pi }`$ we express $`[\pi ^{p,p+1},R^{}]`$ as a projected expression $`\vartheta ^p`$. (For the full expression, see Appendix D.)
These constraints are closed classically (after regularization: see Appendix D)
$`[𝒜(1),𝒜(2)]`$ $`=`$ $`\delta ^{}(21)[𝒜(1)+𝒜(2)]`$
$`[𝒜(1),_p(2)]`$ $`=`$ $`\delta ^{}(21)_p(1)`$
$`[𝒜(1),_p(2)]`$ $`=`$ $`\delta ^{}(21)[_p(1)+_p(2)]`$
$`\{_p(1),_q(2)\}`$ $`=`$ $`0`$
$`\{_p^\alpha (1),_q^\beta (2)\}`$ $`=`$ $`8\delta ^{}(21)\vartheta ^{p\alpha }𝒫_a(1)\gamma ^{a\beta \delta }E_{q\delta }(2)`$
$`+4\delta ^{}(21)\gamma _a^{\alpha \lambda }(\pi ^{p,p}+\pi ^{p+1,p+1})_\lambda (1)\gamma ^{a\beta \delta }E_{q\delta }(2)`$
$`2\delta ^{}(21)(\theta ^{p,p}\theta ^{p+1,p+1}+\vartheta ^p)^\alpha _q^\beta [(1)+(2)]`$
$`+((p,\alpha ,1)(q,\beta ,2))`$
$`\{_{p\alpha }(1),_q^\delta (2)\}`$ $`=`$ $`4\delta _\alpha ^\beta \delta (21)𝒜(1)`$ (4.2.8)
$`2\delta ^{}(21)(\theta ^{q,q}\theta ^{q+1,q+1}+\vartheta ^p)^\beta _{p\alpha }(2)`$
where
$`E_q`$ $``$ $`{\displaystyle \underset{r}{}}i^{rq}{\displaystyle \frac{1}{\sqrt{r}}}(\mathrm{\Theta }_{rq1}\mathrm{\Theta }_{rq2})_{r1}`$
### 4.3 Action
From this Hamiltonian form of the BRST operator for the superstring we can find the ZJBV form :
$$Q^{ZJBV}=\frac{i}{2}\dot{\varphi }^A\varphi _AH\stackrel{ˇ}{\varphi }_A\{\varphi ^A,Q^H]$$
(4.3.1)
where
$`[\varphi ^A,\varphi ^B\}`$ $`=`$ $`\mathrm{\Omega }^{AB}`$
$`\varphi ^A=\mathrm{\Omega }^{AB}\varphi _B,`$ $`\varphi _B=\varphi ^A\mathrm{\Omega }_{AB}`$ (4.3.2)
and $`\stackrel{ˇ}{\varphi }_A`$ is the antifield which is canonically conjugate to the field $`\varphi ^A`$ by the antibracket. ZJBV is useful for Lagrangian quantization, but since $`Q`$ is sufficient for Hamiltonian quantization, we leave the details for Appendix E.
Constraints appearing in the gauge-invariant Hamiltonian have ghost number 0; their ghosts have ghost number 1; their antighosts have ghost number $`1`$; the antifields of their antighosts have ghost number $`0`$, and we can identify them in the ZJBV BRST operator with the Lagrange multipliers of the gauge-invariant Hamiltonian. (More generally, we interpret all the negative-ghost-number fields as antifields.) Let
$`\mathrm{\Phi }_{p,p+1}`$ $``$ $`(1)^{p+1}i\frac{1}{2}\sqrt{p+1}\stackrel{ˇ}{\theta }_{p,p+1}`$
$`\mathrm{\Psi }_{p,p+1}`$ $``$ $`(1)^{p+1}i\frac{1}{2}\sqrt{p+1}\stackrel{ˇ}{\pi }_{p,p+1}`$
$`\stackrel{ˇ}{b}`$ $``$ $`g`$ (4.3.3)
and similarly for their antifields. Then we find the gauge invariant action in Hamiltonian form $`S_H`$ either from the usual Hamiltonian procedure (using the constraints of the previous subsection), or as the antifield-free part of $`Q_{sstring}^{ZJBV}`$:
$$S_H=\dot{X}\widehat{P}^0+i\underset{\pm ,p}{}\dot{\theta }^{p,p}\pi _{p,p}\underset{\pm }{}g_\pm 𝒜_\pm +\underset{\pm ,p0}{}\left(\stackrel{~}{\mathrm{\Phi }}_{\pm ,p,p+1}_p+\stackrel{~}{\mathrm{\Psi }}_{\pm ,p,p+1}_p\right)$$
(4.3.4)
(Again, for each sign $`\pm `$ we use fermions of the corresponding chirality.)
If we consider only the quadratic terms and the $`𝒜`$ term, keeping only the physical fields, and introducing $`\widehat{P}^1`$ as an independent variable, we find the first-order, 2$`D`$ world-sheet covariant form (with world-sheet metric $`\eta _{mn}=(+)`$)
$$S_0^{phys}=\widehat{P}^m_mX\frac{1}{2}g_{mn}\widehat{P}^m\widehat{P}^n+i\sqrt{2}\underset{\pm }{}_\pm \theta _0^\pm \pi _0^\pm $$
(4.3.5)
where $`\theta _0^\pm \theta _{0L,R}`$, $`\pi _0^\pm \pi _{0L,R}`$, and $`_\pm \frac{1}{\sqrt{2}}(e_0{}_{}{}^{m}\pm e_1{}_{}{}^{m})_m`$. By introducing supersymmetric variables
$`P^m`$ $`=`$ $`\widehat{P}^m+ϵ^{mn}(\eta _{(0)n}^+\eta _{(0)n}^{})`$
$`D_0^\pm `$ $`=`$ $`\pi _0^\pm +[\widehat{P}^\pm \pm \frac{1}{2}(\eta _{(0)}^+\eta _{(0)}^{})]\gamma \theta _0^\pm `$
$`\eta _{(0)m}^\pm `$ $``$ $`\frac{i}{\sqrt{2}}(_m\theta _0^\pm )\gamma \theta _0^\pm `$ (4.3.6)
(where we suppress spacetime indices for simplicity) and plugging this into (4.3.5) we find
$`S_0^{phys}`$ $`=`$ $`\frac{1}{2}g_{mn}P^mP^n+P^m[_mX(\eta _{(0)m}^++\eta _{(0)m}^{})]`$ (4.3.7)
$`ϵ^{mn}[(_mX)(\eta _{(0)n}^++\eta _{(0)n}^{})\eta _{(0)m}^+\eta _{(0)n}^{}]+i\sqrt{2}{\displaystyle \underset{\pm }{}}_\pm \theta _0^\pm D_0^\pm `$
Except for the last term this is the Green-Schwarz supersting action. To extend this redefinition to the whole action one can further define (we use $`p,q`$ for ghost level and $`(p),(q)`$ when there is confusion with world sheet indices $`l,m,n`$)
$`𝒫^m`$ $`=`$ $`\widehat{P}^m+ϵ^{mn}(\eta _{(0)n}^+\eta _{(0)n}^{})+ϵ^{mn}(\chi _n^+\chi _n^{})`$
$`𝒟_0^\pm `$ $`=`$ $`\pi _0^\pm +\left\{[\widehat{P}^\pm \pm \frac{1}{2}(\eta _{(0)}^+\eta _{(0)}^{})]\gamma \theta _0^\pm \frac{2}{3}(\xi _{}^+\xi _{}^{})\gamma \stackrel{~}{\theta }^\pm \right\}`$
$`𝒟_p^\pm `$ $`=`$ $`\pi _p^\pm +𝒫^\pm \gamma \theta _p^\pm (p1)`$
$`\pi _p^\pm `$ $``$ $`\pi _{L,R}^{p,p}`$
$`\theta _p^\pm `$ $``$ $`\theta _{L,R}^{p,p}`$
$`\vartheta _p^\pm `$ $``$ $`\vartheta _{L,R}^p`$
$`\chi _m^\pm `$ $``$ $`\frac{1}{\sqrt{2}}\left(2i(_m\theta _0^\pm )\gamma \stackrel{~}{\theta }^\pm i(_m\stackrel{~}{\theta }^\pm )\gamma \stackrel{~}{\theta }^\pm +\frac{1}{2}_m\stackrel{~}{R}^\pm \right)`$
$`\xi _m^\pm `$ $``$ $`\frac{1}{\sqrt{2}}\left(2i(_m\theta _0^\pm )\gamma \stackrel{~}{\theta }^\pm i\frac{5}{4}(_m\stackrel{~}{\theta }^\pm )\gamma \stackrel{~}{\theta }^\pm +\frac{3}{4}_m\stackrel{~}{R}^\pm \right)`$
$`\eta _{(p)m}^\pm `$ $``$ $`\frac{i}{\sqrt{2}}(_m\theta _p^\pm )\gamma \theta _p^\pm `$
$`\mathrm{\Phi }_{p\pm }`$ $``$ $`\mathrm{\Phi }_{p,p+1}^{L,R}`$
$`\mathrm{\Psi }_{p\pm }`$ $``$ $`\mathrm{\Psi }_{p,p+1}^{L,R}`$ (4.3.8)
Then our manifestly worldsheet-covariant action reads
$$S_0=\stackrel{~}{S}_{GS}+i\sqrt{2}\underset{\pm ,p0}{}_\pm \theta _p^\pm 𝒟_p^\pm +S_A$$
(4.3.9)
where $`S_A`$ consists of Lagrange multipliers times all the (first-class) constraints other than Virasoro
$$S_A=\underset{\pm ,p0}{}\left(\mathrm{\Psi }_{p\pm }_p^\pm +\mathrm{\Phi }_{p\pm }_p^\pm \right)$$
(4.3.10)
and $`\stackrel{~}{S}_{GS}`$ is an extension of the usual $`GS`$ action to the fields $`\theta _p^\pm `$ at nonzero ghost levels
$`\stackrel{~}{S}_{GS}`$ $`=`$ $`\frac{1}{2}g_{mn}𝒫^m𝒫^n`$ (4.3.11)
$`+{\displaystyle \underset{\pm }{}}𝒫^\pm \left[_\pm X(\eta _{(0)\pm }^++\eta _{(0)\pm }^{})+{\displaystyle \underset{p1}{}}\eta _{(p)\pm }^\pm \pm (\chi _\pm ^+\chi _\pm ^{})\right]`$
$`ϵ^{mn}\left[(_mX)\left\{(\eta _{(0)n}^+\eta _{(0)n}^{})+(\chi _n^+\chi _n^+)\right\}+\eta _{(0)m}^+\eta _{(0)n}^{}\right]`$
$`+\frac{2}{3}{\displaystyle \underset{\pm }{}}\pm \frac{i}{\sqrt{2}}(_\pm \theta _0^\pm \gamma \stackrel{~}{\theta }^\pm )(\xi _{}^+\xi _{}^{})`$
This is a first-order action in terms of the coordinates $`X,\theta _p^\pm `$, momenta $`𝒫^m,𝒟_p^\pm `$, worldsheet metric $`g_{mn}`$, and Lagrange multipliers $`\mathrm{\Phi }_{p\pm },\mathrm{\Psi }_{p,\pm }`$. Now $`_p`$ and $`_p`$ are expressed in terms of these new variables as
$`_p^\pm `$ $`=`$ $`𝒟_p^\pm 𝒟_{p+1}^\pm +2𝒫^\pm \gamma \theta _{p+1}^\pm `$
$`_p^\pm `$ $`=`$ $`𝒫^\pm \gamma (𝒟_p^\pm +𝒟_{p+1}^\pm 2𝒫^\pm \gamma \theta _{p+1}^\pm )`$
$`+`$ $`(\mathrm{\Theta }_{p1}\theta _p^\pm \theta _{p+1}^\pm +\vartheta _p^\pm )\times \{{\displaystyle \frac{}{}}𝒫^{\pm 2}[𝒫^\pm (\eta _{(0)}^+\eta _{(0)}^{}+\chi _{}^+\chi _{}^{})]^2`$
$`\pm \frac{1}{\sqrt{2}}[{\displaystyle \underset{q1}{}}i_\pm \theta _q^\pm (𝒟_q^\pm 𝒫^\pm \gamma \theta _q^\pm )+i_\pm \theta _0^\pm `$
$`\times (𝒟_0^\pm \{[\widehat{P}^\pm \pm \frac{1}{2}(\eta _{(0)}^+\eta _{(0)}^{})]\gamma \theta _0^\pm \frac{2}{3}(\xi _{}^+\xi _{}^{})\gamma \stackrel{~}{\theta }^\pm \})]\}`$
Elimination of $`𝒫^1`$ by its equation of motion reproduces the previous Hamiltonian form of the action except for terms quadratic in $``$, which can be eliminated by a redefinition of $`\mathrm{\Phi }`$. The gauge-fixed action with ghosts is most easily obtained from the Hamiltonian formalism as $`H=\{Q,b\}=T`$: Then (with the full $`\theta ^{p,q}`$)
$$S_{GF}=\widehat{P}^m_mX\frac{1}{2}\eta _{mn}\widehat{P}^m\widehat{P}^n+i\sqrt{2}\underset{\pm }{}_\pm c^\pm b_{\pm \pm }+i\sqrt{2}\underset{\pm }{}_\pm \theta ^\pm \pi ^\pm $$
(4.3.13)
## 5 Future
We have given a gauge-invariant action for the superstring and its corresponding BRST operator. The BRST-invariant gauge-fixed action is the obvious quadratic expression following from $`\{Q,b\}`$ (and is thus BRST invariant since $`Q^2=0`$). This is sufficient to perform S-matrix calculations (with vertex operators of the type given for the superparticle above), but a naive application would require a measure that breaks manifest supersymmetry. (For example, solving for the cohomology of the superparticle with this BRST operator in required using the equivalent of the lightcone gauge.) In principle, a covariant measure that avoids picture changing altogether (in particular, for the bosonic ghosts) can be found by methods similar to those used in ; we hope to return to this problem. The cohomology of this BRST operator should also be checked: The massless level follows from the previous analysis for the superparticle; the massive levels should follow from a similar lightcone analysis.
## Acknowledgment
This work was supported in part by the National Science Foundation Grant No. PHY-0354776.
## Appendix A Sp(2) components
The matrix elements of the Sp(2) operators are
$`p,q|\gamma ^{}|r,s`$ $`=`$ $`C_{p,q}^{r,s}(\sqrt{p}\delta _{p,r+1}\delta _{q,s}+i\sqrt{q+1}\delta _{p,r}\delta _{q+1,s})`$ (A.1)
$`p,q|\stackrel{~}{\gamma }^{}|r,s`$ $`=`$ $`C_{p,q}^{r,s}(i\sqrt{p}\delta _{p,r+1}\delta _{q,s}\sqrt{q+1}\delta _{p,r}\delta _{q+1,s})`$ (A.2)
where $`C_{p,q}^{r,s}=i^{\frac{(r+s)(r+s+1)(p+q)(p+q+1)}{2}}`$. From these we can find “inverse” operators, especially
$$p,q|\frac{1}{\stackrel{~}{\gamma }^{}}|r,s=i^{rp}C_{p,q}^{r,s}\sqrt{\frac{p!s!}{q!r!}}\delta _{qp+rs,1}[\mathrm{\Theta }_{pq}\mathrm{\Theta }_{sq}\mathrm{\Theta }_{sr}\mathrm{\Theta }_{pr}]$$
(A.3)
It satisfies
$$\stackrel{~}{\gamma }^{}\frac{1}{\stackrel{~}{\gamma }^{}}\stackrel{~}{\gamma }^{}=\stackrel{~}{\gamma }^{},\stackrel{~}{\gamma }^{}\frac{1}{\stackrel{~}{\gamma }^{}}I,\frac{1}{\stackrel{~}{\gamma }^{}}\stackrel{~}{\gamma }^{}I$$
(A.4)
The arrows means there is cancellation among the multiplied matrix elements to the infinite ghost level. The subtle point of this cancellation will be studied in Appendix B.
Then we find that $`\gamma ^{}`$ and $`\frac{1}{\stackrel{~}{\gamma }^{}}`$ don’t (anti)commute but give
$`p,q|\{{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}},\gamma ^{}\}|r,s`$ $`=`$ $`i^{rp+1}C_{p,q}^{r,s}\sqrt{{\displaystyle \frac{p!s!}{q!r!}}}\delta _{qp+rs,0}[(\mathrm{\Theta }_{pq}+\mathrm{\Theta }_{pq1})(\mathrm{\Theta }_{sq}+\mathrm{\Theta }_{sq1})`$ (A.5)
$`(\mathrm{\Theta }_{sr1}+\mathrm{\Theta }_{sr})(\mathrm{\Theta }_{pr1}+\mathrm{\Theta }_{pr})]`$
$`p,q|[{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}},\gamma ^{}]|r,s`$ $`=`$ $`i^{rp+1}C_{p,q}^{r,s}\sqrt{{\displaystyle \frac{p!s!}{q!r!}}}\delta _{qp+rs,0}[(\mathrm{\Theta }_{pq}\mathrm{\Theta }_{pq1})(\mathrm{\Theta }_{sq}+\mathrm{\Theta }_{sq1})`$ (A.6)
$`(\mathrm{\Theta }_{sr1}\mathrm{\Theta }_{sr})(\mathrm{\Theta }_{pr1}+\mathrm{\Theta }_{pr})]`$
$`=`$ $`2i\delta _{p,q}\delta _{r,s}`$
Some interesting and useful commutators are
$`p,q|[\gamma ^{},\{{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}},\gamma ^{}\}]|r,s`$ $`=`$ $`4i^{rp}C_{p,q}^{r,s}(\sqrt{r+1}\delta _{p,q}\delta _{r+1,s}+\sqrt{p}\delta _{p,q+1}\delta _{r,s})`$ (A.7)
$`p,q|[\gamma ^{}\gamma ^{},\{{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}},\gamma ^{}\}]|r,s`$ $`=`$ $`8i^{rp+1}C_{p,q}^{r,s}(\sqrt{(r+2)(r+1)}\delta _{p,q}\delta _{r+2,s}`$
$`+`$ $`\sqrt{p(p1)}\delta _{p,q+2}\delta _{r,s}2\sqrt{p(r+1)}\delta _{p,q+1}\delta _{r+1,s})`$
$`p,q|[\gamma ^{}\gamma ^{},\{{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}},\gamma ^{}\}]\gamma ^{}|r,s`$ $`=`$ $`16i^{rp}C_{p,q}^{r,s}(\sqrt{(r+3)(r+2)(r+1)}\delta _{p,q}\delta _{r+3,s}`$ (A.9)
$`+`$ $`\sqrt{p(p1)(r+1)}\delta _{p,q+2}\delta _{r+1,s}`$
$``$ $`2\sqrt{p(r+2)(r+1)}\delta _{p,q+1}\delta _{r+1,s})`$
Using (A.1),(A.2) and (A.5) we find
$`\overline{\pi }\gamma ^{}\gamma ^{}\theta `$ $`=`$ $`{\displaystyle \underset{pq}{}}[\sqrt{p(p1)}(1)^{p+1}i^{p+q+1}\overline{\pi }^{p,q}\theta ^{q,p2}`$ (A.10)
$`+2\sqrt{p(q+1)}(1)^{q+1}i^{p+q+1}\overline{\pi }^{p,q}\theta ^{q+1,p1}`$
$`+\sqrt{(q+1)(q+2)}(1)^{p+1}i^{p+q+1}\overline{\pi }^{p,q}\theta ^{q+2,p}]`$
$$\frac{1}{2}\overline{\pi }\stackrel{~}{\gamma }^{}\pi =\underset{pq}{}\sqrt{p}(1)^{q+1}\overline{\pi }^{p,q}\pi ^{q,p1}$$
(A.11)
$`\overline{\pi }\gamma ^{}\gamma ^a\theta `$ $`=`$ $`i{\displaystyle \underset{pq}{}}[\sqrt{p}(1)^p\overline{\pi }^{p,q}\gamma ^a\theta ^{q,p1}`$ (A.12)
$`+\sqrt{q+1}(1)^{q+1}\overline{\pi }^{p,q}\gamma ^a\theta ^{q+1,p}]`$
$`R^a`$ $`=`$ $`\frac{1}{4}\overline{\theta }\{{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}},\gamma ^{}\}\gamma ^a\theta `$ (A.13)
$`=`$ $`\frac{1}{2}{\displaystyle \underset{pqr}{}}(1)^{(qp+r+1)rqpp}i^{qp+1}\sqrt{{\displaystyle \frac{p!(qp+r)!}{q!r!}}}\overline{\theta }^{p,q}\gamma ^a\theta ^{qp+r,r}`$
$`[(\mathrm{\Theta }_{pq}+\mathrm{\Theta }_{pq1})(\mathrm{\Theta }_{rp}+\mathrm{\Theta }_{rp1})`$
$`(\mathrm{\Theta }_{qp1}+\mathrm{\Theta }_{qp})(\mathrm{\Theta }_{pr1}+\mathrm{\Theta }_{pr})]`$
$`R^{}`$ $`=`$ $`\frac{1}{4}\overline{\theta }\{{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}},\gamma ^{}\}\gamma ^{}\theta `$ (A.14)
$`=`$ $`\frac{1}{2}{\displaystyle \underset{pqr}{}}(1)^{(qp+r+2)rqpp}\sqrt{{\displaystyle \frac{p!(qp+r+1)!}{q!r!}}}\overline{\theta }^{p,q}\theta ^{qp+r+1,r}`$
$`[(\mathrm{\Theta }_{pq}+\mathrm{\Theta }_{pq1})(\mathrm{\Theta }_{rp+1}+2\mathrm{\Theta }_{rp}+\mathrm{\Theta }_{rp1})`$
$`(\mathrm{\Theta }_{qp1}+\mathrm{\Theta }_{qp})(\mathrm{\Theta }_{pr2}+2\mathrm{\Theta }_{pr1}+\mathrm{\Theta }_{pr})]`$
The component fields defined above satisfy
$$\{\pi _{p,q},\theta ^{r,s}]=\delta _p^r\delta _q^s$$
(A.15)
## Appendix B Subtle points in Sp(2) operators
In this appendix we explain some subtle points about $`\frac{1}{\stackrel{~}{\gamma }^{}}`$, due to the infinite dimensional structure of Sp(2) operators.
Consider the commutator
$$\{\overline{\pi }\stackrel{~}{\gamma ^{}}\theta ,[\stackrel{(1)}{\theta }\stackrel{(0)}{\theta },R^a]\}$$
(B.1)
where
$`\stackrel{(0)}{\theta }`$ $`=`$ $`0|e^{ia_{}a_{}}|\theta `$
$`=`$ $`\theta _0+\stackrel{~}{\theta }`$
$`=`$ $`\theta _0+\theta ^{}+\theta ^{}+\theta ^{}+\mathrm{}`$
$`\stackrel{(1)}{\theta }`$ $`=`$ $`0|e^{ia_{}a_{}}(ia_{})|\theta `$ (B.2)
$`=`$ $`2i(\theta ^{}\sqrt{2}\theta ^{}+\sqrt{3}\theta ^{}\sqrt{4}\theta ^{}+\mathrm{})`$
We will also need:
for $`n0`$
$`\stackrel{(n)}{\theta }`$ $`=`$ $`0|e^{ia_{}a_{}}(ia_{})^n|\theta `$ (B.3)
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1)^{k(n+k+1)}i^{{\scriptscriptstyle \frac{n(n+1)}{2}}}(1+\mathrm{\Theta }_{n{\scriptscriptstyle \frac{1}{2}}})\sqrt{{\displaystyle \frac{(n+k)!}{k!}}}\theta ^{n+k,k}`$
for $`n<0`$
$`\stackrel{(n)}{\theta }`$ $`=`$ $`0|e^{ia_{}a_{}}(ia^{})^{|n|}|\theta `$ (B.4)
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1)^{k(|n|+k+1)}i^{{\scriptscriptstyle \frac{|n|(|n|+1)}{2}}}(1+\mathrm{\Theta }_{|n|{\scriptscriptstyle \frac{1}{2}}})\sqrt{{\displaystyle \frac{k!}{(|n|+k)!}}}\theta ^{k,|n|+k}`$
The above double commutator should vanish because the inner one involves only $`\theta `$. However, if we apply the Jacobi identity we see
$`\{\overline{\pi }\stackrel{~}{\gamma ^{}}\pi ,[\stackrel{(1)}{\theta }\stackrel{(0)}{\theta },R^a]\}`$ $`=`$ $`0=[\{\overline{\pi }\stackrel{~}{\gamma ^{}}\pi ,\stackrel{(1)}{\theta }\stackrel{(0)}{\theta }\},R^a]\{\stackrel{(1)}{\theta }\stackrel{(0)}{\theta },[\overline{\pi }\stackrel{~}{\gamma ^{}}\pi ,R^a]\}`$ (B.5)
$`=[\stackrel{~}{0},R^a]2\{\stackrel{(1)}{\theta }\stackrel{(0)}{\theta },\overline{\pi }\gamma ^{}\gamma ^a\theta \}`$
$`=[\stackrel{~}{0},R^a]+2\stackrel{(1)}{\theta }\gamma ^a\stackrel{(1)}{\theta }+4\stackrel{(0)}{\theta }\gamma ^a\stackrel{(2)}{\theta }`$
where we have introduced the scalar $`\stackrel{~}{0}`$ defined by
$$[\frac{1}{2}\overline{\pi }\gamma ^{}\pi ,\stackrel{(n)}{\theta }](n1)i^{{\scriptscriptstyle \frac{n(n+1)}{2}}+n}(1+\mathrm{\Theta }_{|n|{\scriptscriptstyle \frac{1}{2}}})\sqrt{\frac{(n+1)!}{0!}}(\overline{\pi }^{n+1,0}\overline{\pi }^{n+1,0})$$
$`+(1)^ni^{{\scriptscriptstyle \frac{n(n+1)}{2}}+n}(1+\mathrm{\Theta }_{|n|{\scriptscriptstyle \frac{1}{2}}})\sqrt{{\displaystyle \frac{(n+2)!}{1!}}}(\overline{\pi }^{n+2,1}\overline{\pi }^{n+2,1})`$
$`+\mathrm{}`$
$`+(1)^{(k+2)(n+k+2)1}i^{{\scriptscriptstyle \frac{n(n+1)}{2}}+n}(1+\mathrm{\Theta }_{|n|{\scriptscriptstyle \frac{1}{2}}})\sqrt{{\displaystyle \frac{(n+k+1)!}{k!}}}(\overline{\pi }^{n+k+1,k}\overline{\pi }^{n+k+1,k})`$
$`+\mathrm{}`$
$`\stackrel{~}{0}0`$ (B.6)
and
$$[\stackrel{(n)}{\theta },\overline{\pi }\gamma ^{}\gamma ^a\theta ]=i^n(1+\mathrm{\Theta }_{n{\scriptscriptstyle \frac{1}{2}}})\gamma ^a\stackrel{(n+1)}{\theta }$$
(B.7)
Unfortunately, $`\stackrel{~}{0}`$ is not zero in the presence of $`R^a`$. Let’s assume the collective $`\stackrel{\left(n\right)}{\theta }_k`$ has $`k+1`$ terms instead of an infinite number of terms. Then $`[\frac{1}{2}\overline{\pi }\gamma ^{}\pi ,\stackrel{\left(n\right)}{\theta }{}_{k}{}^{}]`$ gives
$$(1)^{(k+2)(n+k+2)1}i^{{\scriptscriptstyle \frac{n(n+1)}{2}}+n}(1+\mathrm{\Theta }_{|n|{\scriptscriptstyle \frac{1}{2}}})\sqrt{\frac{(n+k+1)!}{k!}}\overline{\pi }^{n+k+1,k}$$
Finally, one can see that
$$(1)^{(k+2)(n+k+2)1}i^{{\scriptscriptstyle \frac{n(n+1)}{2}}+n}(1+\mathrm{\Theta }_{|n|{\scriptscriptstyle \frac{1}{2}}})\sqrt{\frac{(n+k+1)!}{k!}}[\overline{\pi }^{n+k+1,k},R^a]$$
gives $`i^n(1+\mathrm{\Theta }_{|n|{\scriptscriptstyle \frac{1}{2}}})\stackrel{\left(n+1\right)}{\theta }_k`$. So if we make the collective $`\stackrel{\left(n\right)}{\theta }_k`$ have an infinite number of terms by $`k\mathrm{}`$ then $`\stackrel{~}{0}`$ produces a nonzero result in the commutator with $`R^a`$. This exactly cancels the remaining terms in (B.5) to make the double commutator consistent.
Keeping this subtle point in mind let’s consider the following transformation
$`Q_{free}^{}`$ $`=`$ $`e^{iR^ap_a}\stackrel{~}{Q}_{free}e^{iR^ap_a}`$ (B.8)
$`=`$ $`e^{iR^ap_a}(\frac{1}{2}c\mathrm{}\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\theta b+\stackrel{~}{0}b+[iR^ap_a,\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\theta +\stackrel{~}{0}]b`$
$`+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi +[iR^ap_a,\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi ]+\frac{1}{2}[iR^ap_a,[iR^ap_a,\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi ]]`$
$`\frac{i}{2}\overline{\pi }\gamma ^{}\text{/}p\theta +[iR^ap_a,\frac{i}{2}\overline{\pi }\gamma ^{}\text{/}p\theta ])e^{iR^ap_a}`$
where $`\stackrel{~}{0}b`$ vanishes in $`Q_{free}^{}`$ but will give a nontrivial contribution in the presence of $`R^a`$. We will determine this term from nilpotency of $`Q_{free}^{}`$.
In terms of $`R^{}`$ and $`\stackrel{(n)}{\theta }`$, $`\stackrel{~}{Q}_{free}`$ is
$`\stackrel{~}{Q}_{free}`$ $`=`$ $`\frac{1}{2}c\mathrm{}\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\theta b\stackrel{(2)}{\theta }\gamma ^a\stackrel{(0)}{\theta }p_ab+\frac{1}{2}\stackrel{(1)}{\theta }\gamma ^a\stackrel{(1)}{\theta }p_ab`$ (B.9)
$`+\stackrel{~}{0}b+[iR^a,\stackrel{~}{0}]p_ab+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi \stackrel{~}{0}^ap_a`$
$`+\frac{1}{2}[iR^ap_a,\stackrel{~}{0}^bp_b]+\frac{1}{2}R^{}\mathrm{}`$
where
$$[iR^ap_a,\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\theta ]b=\stackrel{(2)}{\theta }\gamma ^a\stackrel{(0)}{\theta }p_ab+\frac{1}{2}\stackrel{(1)}{\theta }\gamma ^a\stackrel{(1)}{\theta }p_ab$$
(B.10)
and we have defined the vector $`\stackrel{~}{0}^a`$
$$\stackrel{~}{0}^a[iR^a,\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi ]+\frac{i}{2}\overline{\pi }\gamma ^{}\gamma ^a\theta $$
(B.11)
The nilpotency of $`Q_{free}^{}`$ implies
$$\{\stackrel{~}{0}^a,\frac{1}{2}R^{}\}p_a\mathrm{}=\frac{1}{2}(\stackrel{(2)}{\theta }\gamma ^a\stackrel{(0)}{\theta }+\frac{1}{2}\stackrel{(1)}{\theta }\gamma ^a\stackrel{(1)}{\theta }+[iR^a,\stackrel{~}{0}])p_a\mathrm{}$$
(B.12)
$`[\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\theta ,\frac{1}{2}[iR^ap_a,\stackrel{~}{0}^bp_b]]b`$ $`=`$ $`[\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\theta +\stackrel{~}{0},\frac{1}{2}R^{}]\mathrm{}b`$
$`=`$ $`\frac{1}{4}(3i\stackrel{(1)}{\theta }\stackrel{(2)}{\theta }2i\stackrel{(0)}{\theta }\stackrel{(3)}{\theta }2[\stackrel{~}{0},R^{}])\mathrm{}b`$
From (B.11) we can find $`\stackrel{~}{0}^a`$ as
$$\stackrel{~}{0}^a=\frac{i}{4}\stackrel{\left(1\right)}{\pi }{}_{\mathrm{}}{}^{}\gamma _{}^{a}\stackrel{(0)}{\theta }\frac{i}{4}\stackrel{(0)}{\pi _{\mathrm{}}}\gamma ^a\stackrel{(1)}{\theta }$$
(B.14)
where
$$\stackrel{(n)}{\chi _{\mathrm{}}}\underset{k\mathrm{}}{lim}(1)^{k(n+k+1)}(i)^{{\scriptscriptstyle \frac{n(n+1)}{2}}}\sqrt{\frac{(n+k+1)!}{k!}}\chi ^{n+k,k}$$
(B.15)
for $`\chi =\theta ,\pi `$. Inserting this into (B.12) we find
$`[iR^a,\stackrel{~}{0}]`$ $`=`$ $`\stackrel{(1)}{\theta }\gamma ^a\stackrel{(1)}{\theta }`$ (B.16)
$`[\stackrel{~}{0},R^{}]`$ $`=`$ $`2i\stackrel{(1)}{\theta }\stackrel{(2)}{\theta }`$ (B.17)
One solution for $`\stackrel{~}{0}`$ is
$$\stackrel{~}{0}=\stackrel{\left(1\right)}{\pi }{}_{\mathrm{}}{}^{}\stackrel{(1)}{\theta }$$
(B.18)
Now $`\stackrel{~}{Q}_{free}`$ is nilpotent, as it should be. However, the origin of collective nonminimal fields is that $`R^a`$ produces $`\overline{\pi }\gamma ^{}\gamma ^a\theta `$ using $`[\overline{\pi }\stackrel{~}{\gamma }^{}\pi ,\stackrel{(n)}{\theta }]0`$. (The key feature of $`R^a`$ is $`\stackrel{(n)}{\theta }`$, see (A.13).) And this implies that $`R^a`$ and $`R^{}`$ commute with $`\overline{\pi }\gamma ^{}\gamma ^{}\theta `$ only up to these collective nonminimal fields (see (A.7)). So these collective fields are just mathematical objects to compensate terms like $`[R^i,\stackrel{~}{0}]`$. Therefore, a physical (but not mathematical) equivalent is to drop $`\stackrel{~}{0}^{(a)}`$ and $`\stackrel{(n)}{\theta }`$ and regard $`R^i`$ as terms commuting with $`\overline{\pi }\gamma ^{}\gamma ^{}\theta `$ in (B.9). The resulting $`\stackrel{~}{Q}_{free}`$ is simply
$$\stackrel{~}{Q}_{free}=\frac{1}{2}c\mathrm{}\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\theta b+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi +\frac{1}{2}R^{}\mathrm{}$$
(B.19)
with $`[\overline{\pi }\gamma ^{}\gamma ^{}\theta ,R^{}]0`$.
Our results for $`\stackrel{}{Q}_{YMB}^{}`$(3.1.1), $`\stackrel{}{Q}_{SYMB}`$(3.2.11) and $`Q_{SYMB}^{\prime \prime }`$(3.3.1) all reflect this prescription. $`\stackrel{~}{\theta }`$ in these BRST operators will produce only $`\pi ^{}`$$`\theta ^{}`$ and $`\theta ^{}b`$ ($`\theta ^{}b`$ will only appear in (3.1.1)) dropping $`\stackrel{(1)}{\theta }`$ and $`\stackrel{(2)}{\theta }b`$. If we want to be rigorous $`\stackrel{(n)}{\theta }`$ and $`\stackrel{~}{0}^{(a)}`$ should be kept and the commutator $`[R^i,\overline{\pi }\gamma ^{}\gamma ^{}\theta ]`$ should be calculated for both to compensate terms from $`\stackrel{~}{0}^{(a)}`$.
If we consider this mathematical rigor for $`Q_{YMB}^{}`$(3.1.1) we find
$`Q_{YMB}^{collective}`$ $`=`$ $`e^{iR^a_a}[\frac{1}{2}c(\mathrm{}\overline{\pi }\gamma ^{ab}\theta F_{ab})+\frac{1}{2}R^{}(\mathrm{}\overline{\pi }\gamma ^{ab}\theta F_{ab})`$ (B.20)
$`+`$ $`\frac{1}{2D}[iR^a,\stackrel{~}{0}_a](\mathrm{}\overline{\pi }\gamma ^{ab}\theta F_{ab})\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\theta b+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi `$
$`+`$ $`\stackrel{~}{0}b\stackrel{(2)}{\theta }\gamma ^a\stackrel{(0)}{\theta }_ab+\frac{1}{2}\stackrel{(1)}{\theta }\gamma ^a\stackrel{(1)}{\theta }_ab[iR^a,\stackrel{~}{0}]_ab`$
$``$ $`\stackrel{~}{0}^a_a`$
$``$ $`\frac{1}{2}\{(cb\frac{1}{2})+R^{}b+\frac{1}{2D}[iR^a,\stackrel{~}{0}_a]b\}`$
$`\times \overline{\theta }{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}}\gamma ^{ab}\theta [F_{ab},_c](\stackrel{(2)}{\theta }\gamma ^c\stackrel{(0)}{\theta }+\frac{1}{2}\stackrel{(1)}{\theta }\gamma ^c\stackrel{(1)}{\theta }[iR^c,\stackrel{~}{0}])`$
$`+`$ $`\frac{1}{2}\{c+R^{}+\frac{1}{2D}[iR^a,\stackrel{~}{0}_a]\}`$
$`\times \overline{\theta }{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}}\gamma ^{ab}\theta [F_{ab},_c]\stackrel{~}{0}^c]e^{iR^a_a}|_{\text{linear in }F,[,F]=0}`$
where
$`D`$ $`=`$ $`\text{dimension of space-time}(\text{10 here})`$
$`\text{and}\mathrm{}`$ $`=`$ $`^a_a`$ (B.21)
If we drop $`\stackrel{(n)}{\theta }`$ and $`\stackrel{~}{0}^{(a)}`$ and regard $`[\gamma ^{},\frac{1}{\stackrel{~}{\gamma }^{}}]=0`$ (which is equivalent to $`[R^i,\overline{\pi }\gamma ^{}\gamma ^{}\theta ]]=0`$) we come back to $`Q_{YMB}^{}`$(3.1.1). The last four lines do not contribute to $`Q_{YMB}^{}`$ but are there for nonconstant Yang-Mills background.
## Appendix C Regularization
In this appendix we will consider a regularization procedure which will give the prescription of the previous appendix. The motivation is the fact that $`[\overline{\pi }\stackrel{~}{\gamma }^{}\pi ,\stackrel{(n)}{\theta }]0`$. However, this does not exactly vanish, but leaves a piece of $`\stackrel{\left(n+1\right)}{\pi }_{\mathrm{}}`$. This remnant gives a nontrivial contribution in the presence of $`\frac{1}{\stackrel{~}{\gamma }^{}}`$, i.e.,
$`0|e^{ia_{}a_{}}(ia_{})^n\stackrel{~}{\gamma }^{}|\theta `$ $``$ $`0`$
$`0|e^{ia_{}a_{}}(ia_{})^n\stackrel{~}{\gamma }^{}{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}}|\theta `$ $``$ $`\stackrel{(n)}{\theta }`$ (C.1)
Now if we introduce some regularization parameter $`z`$ as
$$p,q|\stackrel{~}{\gamma }^{}|r,s\stackrel{regularized}{}z^{p+r+s+r}p,q|\stackrel{~}{\gamma }^{}|r,s,z1$$
(C.2)
then (C) becomes
$`0|e^{ia_{}a_{}}(ia_{})^n\stackrel{~}{\gamma }^{}|\theta _k`$ $``$ $`z^{2n+2k+1}\theta ^{n+k+1,k}`$
$`0|e^{ia_{}a_{}}(ia_{})^n\stackrel{~}{\gamma }^{}{\displaystyle \frac{1}{\stackrel{~}{\gamma }^{}}}|\theta _k`$ $``$ $`z^{2n+2k+1}\stackrel{\left(n\right)}{\theta }_k`$ (C.3)
where $`k`$ means the collective field has $`k+1`$ terms. The regularization is to send $`k`$ to infinity with fixed $`z`$ $`(<1)`$. Then the operator $`\stackrel{~}{\gamma }^{}\frac{1}{\stackrel{~}{\gamma }^{}}`$ is a projection operator which will project out the collective fields. The free terms in any $`Q_{backgroud}`$ all have this projection operator which goes to the identity in the absence of $`\frac{1}{\stackrel{~}{\gamma }^{}}`$.
$`\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\theta b+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\pi `$ $``$ $`i\frac{1}{2}_a\overline{\pi }\gamma ^{}\gamma ^a\theta `$ (C.4)
$``$ $`\frac{1}{2}\overline{\pi }\gamma ^{}\gamma ^{}\mathrm{\Pi }\theta b+\frac{1}{4}\overline{\pi }\stackrel{~}{\gamma }^{}\mathrm{\Pi }\pi i\frac{1}{2}_a\overline{\pi }\gamma ^{}\gamma ^a\mathrm{\Pi }\theta `$
The arrow means inserting the projection operator and dropping collective fields after expansion of exponential factors.
If a collective field is truncated, i.e., with incomplete beginning or ending components, then we cannot remove it by this regularization procedure, but we can still avoid its contribution. This situation occurs when we consider the commutator
$$\{\eta \pi ^{p,p+1},[_0\stackrel{~}{\theta },\overline{\pi }\gamma ^{}\gamma ^a\theta ]\}=\{\eta \pi ^{p,p+1},i_0\gamma ^a\theta ^{}\}=0(p>0)$$
(C.5)
where $`\eta `$ is a constant fermionic field. This implies
$$[_0\stackrel{~}{\theta },(\pi ^{p,p}+\pi ^{p+1,p+1})\eta ]=0$$
(C.6)
The second argument in the commutator is an example of a truncated collective field. Actually, this commutator indeed vanishes if we consider $`\stackrel{\left(1\right)}{\theta }`$, which we drop in regularization. So for consistent regularization we take this as vanishing.
This fact is applied for closure of the algebra
$$[\eta \theta ^{q,q+1},\{[\eta ^{}\pi ^{p,p+1},Q_R],Q_R\}]$$
(C.7)
where $`Q_R`$ is any version of the BRST operator including $`R^i`$. After canceling ghost-number-nonzero components this commutator reduces to
$$[\eta ^{}C_p^{},\eta C_q^+]$$
(C.8)
where $`C_p^\pm `$ will be superstring constraints if we use $`Q_{sstring}`$. But the original commutator is just
$$\frac{1}{2}[\eta \theta ^{q,q+1},[\eta ^{}\pi ^{p,p+1},\{Q_R,Q_R\}]]$$
(C.9)
and it is just zero due to nilpotency of $`Q_R`$. However, $`\eta ^{}C_p^{}`$ has a term like $`\eta \gamma ^a(\pi ^{p,p}+\pi ^{p+1,p+1})`$, which comes from $`[\eta ^{}\pi ^{p,p+1},\overline{\pi }\gamma ^{}\gamma ^a\theta ]`$. But this combination of $`\pi ^{p,p}`$ is just an example of truncated collective fields. This truncated collective field will interact with $`\stackrel{~}{\theta }`$ in $`\eta C_q^+`$ giving a nonzero contribution in this apparently vanishing commutator. This should be canceled by a $`[\eta ^{}\pi ^{p,p+1},\stackrel{(1)}{\theta }]`$ contribution, which we projected out by regularization. What this means is that for consistent regularization we should take the commutator with truncated collective fields, $`\stackrel{~}{\theta }`$$`R^i`$ as vanishing. For example, we should take $`[(\pi ^{p,p}+\pi ^{p+1,p+1}),\stackrel{~}{\theta }]`$ as zero but we should calculate $`\{(\theta ^{p,p}+\theta ^{p+1,p+1}),(\pi ^{p,p}+\pi ^{p+1,p+1})\}`$ in $`[\eta ^{}C_p^{},\eta C_p^+]`$, both of which come from $`\overline{\pi }\gamma ^{}\gamma ^a\theta `$.
## Appendix D Closure of constraints
First of all, if one directly calculates $`[\pi ^{p,p+1},R^{}]`$ one gets
$`[\pi ^{p,p+1},R^{}]`$ $`=`$ $`\frac{3}{4}{\displaystyle \underset{r}{}}(1)^p\sqrt{p+1}\theta ^{r,r}(\mathrm{\Theta }_{rp}+2\mathrm{\Theta }_{rp1}+\mathrm{\Theta }_{rp2})`$ (D.1)
$`+\frac{1}{4}{\displaystyle \underset{r}{}}(1)^p\sqrt{p+1}\theta ^{r,r}[(\mathrm{\Theta }_{pr+1}+2\mathrm{\Theta }_{pr}+\mathrm{\Theta }_{pr1})`$
$``$ $`\frac{3}{4}A+\frac{1}{4}B`$
$`=`$ $`\frac{1}{2}(AB)\frac{1}{4}(A+B)`$
But $`A+B`$ is just $`(1)^p\sqrt{p+1}\stackrel{(0)}{\theta }`$, and it will vanish when it acts on the projection operator $`\mathrm{\Pi }`$. Also,
$$AB=2(1)^p\sqrt{p+1}\vartheta _p$$
in terms of $`\vartheta `$ (4.2). When we calculate closure of the constraints, there are two types of terms related to the above expression, i.e.,
$`\{[\eta \pi ^{p,p+1},R^{}],[\eta ^{}\theta ^{q,q+1},\overline{\pi }\stackrel{~}{\gamma }^{}\pi ]\}`$
and
$`\{[\eta \pi ^{p,p+1},R^{}],[\eta ^{}\pi ^{q,q+1},\overline{\pi }\gamma ^{}\gamma ^a\theta ]\}`$
$`\{[\eta \pi ^{p,p+1},\overline{\pi }\gamma ^{}\gamma ^a\theta ],[\eta ^{}\pi ^{q,q+1},R^{}]\}`$
where $`\eta `$ and $`\eta ^{}`$ are constant spinors.
The first type is always zero due to the symmetry of $`\gamma ^{}`$ and $`\stackrel{~}{\gamma }^{}`$. If $`pq`$, $`pq+1`$ and $`p+1q`$ the second type cancels because of the $`AB`$ sign difference. If $`p=q`$ we can express $`[\eta \pi ^{p,p+1},R^{}]`$ as just $`\sqrt{p+1}\eta (\theta ^{p,p}\theta ^{p+1,p+1})`$. If $`p=q+1`$ or $`p+1=q`$ we can use $`\sqrt{p+1}\eta (\theta ^{p,p}\theta ^{p+1,p+1})`$. This becomes clearer in a simpler situation,
$$\{R^{},\overline{\pi }\gamma ^{}\gamma ^a\theta \}=0$$
Then we have
$$\{\eta ^{}\pi ^{q,q+1},[\eta \pi ^{p,p+1},\{R^{},\overline{\pi }\gamma ^{}\gamma ^a\theta \}]\}=0$$
But this implies
$$[\{\eta \pi ^{p,p+1},R^{}\},\{\eta ^{}\pi ^{q,q+1},\overline{\pi }\gamma ^{}\gamma ^a\theta \}]+[\{\eta ^{}\pi ^{q,q+1},R^{}\},\{\eta \pi ^{p,p+1},\overline{\pi }\gamma ^{}\gamma ^a\theta \}]=0$$
This is one part of the closure of constraints. A similar analysis shows
$`[\eta \widehat{P}_a\gamma ^a(\theta ^{p,p}+\theta ^{p+1,p+1}),\eta ^{}\widehat{P}_a\gamma ^a(\theta ^{p,p}+\theta ^{p+1,p+1})]`$ (D.2)
$``$ $`[\eta (\pi ^{p,p}\pi ^{p+1,p+1}),\frac{i}{4}\stackrel{~}{R}_a]\eta ^{}\gamma ^a(\theta ^{p,p}+\theta ^{p+1,p+1})`$
$`+\eta \gamma ^a(\theta ^{p,p}+\theta ^{p+1,p+1})[\frac{i}{4}\stackrel{~}{R}_a,\eta ^{}(\pi ^{p,p}\pi ^{p+1,p+1})]`$
where $``$ means we should drop truncated collective fields.
Secondly, from the fact that $`[\stackrel{~}{\theta },\overline{\pi }\gamma ^{}\gamma ^a\theta ]`$=$`\stackrel{(1)}{\theta }i\theta ^{}`$ $``$ $`i\theta ^{}`$ due to regularization when $`\stackrel{~}{\theta }`$ acts on terms from $`\overline{\pi }\gamma ^{}\gamma ^a\theta `$, we should change the sign of $`\stackrel{~}{\theta }`$. ( More precisely, they are all zero except for $`\pi ^{1,1}`$ from $`\overline{\pi }\gamma ^{}\gamma ^a\theta `$. For only this term we can see this sign-change effect as explained in the previous appendix.) This fact was implicitly expressed with the projection operator $`\mathrm{\Pi }`$ in the constraints.
Finally, one should be cautious about $`\{\overline{\pi }\gamma ^{}\gamma ^a\theta `$ , $`\overline{\pi }\gamma ^{}\gamma ^b\theta \}`$ $`=`$ $`(1)(2)g^{ab}`$ $`\overline{\pi }\gamma ^{}\gamma ^{}\theta `$. The “$``$” sign comes from the fact that OSp(2) gamma matrices (and therefore $`a^{}`$ and $`a`$) anticommute with ordinary gamma matrices. This gives an additional sign when one calculates terms like $`[\eta \gamma ^a\theta ^{p,p}`$ , $`\eta ^{}\gamma ^b\pi ^{p,p}]`$.
## Appendix E ZJBV form of BRST
The ZJBV form of the BRST operator follows from the Hamiltonian form of the BRST operator
$$Q^{ZJBV}=\frac{i}{2}\dot{\varphi }^A\varphi _A\stackrel{ˇ}{\varphi }_A\{\varphi ^A,Q^H]$$
(E.1)
Then in our case,
$$Q_{sstring}^{ZJBV}=\dot{X}P^0+\underset{\pm }{}Q_{sstring}^{(\pm )ZJBV}$$
(E.2)
where, e.g., for the $`(+)`$ term (for $`()`$, just add a $``$ for each )
$`Q_{sstring}^{ZJBV}`$ $`=`$ $`\dot{c}b+\dot{\theta }\pi `$
$`+`$ $`(\stackrel{ˇ}{X}_a\stackrel{ˇ}{P}_a^{})[(c+\frac{1}{2}\overline{\theta }\stackrel{~}{\gamma }^{}\theta )\widehat{P}^a\frac{i}{2}\overline{\pi }\gamma ^{}\gamma ^a\theta `$
$`\left(i\theta _0\gamma _a\theta _0^{}+2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{1}{2}R_a^{}|_>\right)R^{}+i\frac{1}{2}\overline{\theta }^{}\theta _0\theta _0\gamma ^a\theta _0^{}`$
$`\frac{1}{2}\overline{\theta }^{}\gamma ^b\gamma ^a\theta _0\left(i\frac{1}{2}\theta _0\gamma _b\theta _0^{}+2i\stackrel{~}{\theta }\gamma _b\theta _0^{}+i\stackrel{~}{\theta }\gamma _b\stackrel{~}{\theta }^{}\frac{1}{2}R_a^{}|_>\right)`$
$`\frac{1}{3}\overline{\theta }^{}\gamma ^a\gamma ^b\stackrel{~}{\theta }(2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\frac{5}{4}\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{3}{4}R_a^{}|_>)]`$
$``$ $`\stackrel{ˇ}{c}\left[icc^{}+2\overline{\pi }a^{}a^{}\theta |_>\right]`$
$``$ $`\stackrel{ˇ}{b}\left(\frac{1}{2}\widehat{P}^22ic^{}bicb^{}i\overline{\theta }^{}\pi i[\overline{\theta }(a^{}a_{}a^{}a_{})\pi ]^{}\right)`$
$`+`$ $`{\displaystyle \underset{p>1}{}}\stackrel{ˇ}{\theta }_{p,p+1}[\left\{i(c+R^{}+\frac{1}{2}\overline{\theta }\stackrel{~}{\gamma }^{}\theta )|_>\theta ^{p,p+1}\right\}`$
$`2i[(p+1)\sqrt{p+1}{\displaystyle \frac{\sqrt{p}+\sqrt{p+2}}{2}}]\mathrm{\Pi }\theta ^{p+1,p}b`$
$`+(1)^{p+1}i\frac{1}{2}\sqrt{p+1}\mathrm{\Pi }(\pi ^{p,p}\pi ^{p+1,p+1})`$
$`+(1)^{p+1}i\frac{1}{2}\sqrt{p+1}\gamma ^a\mathrm{\Pi }(\theta ^{p,p}+\theta ^{p+1,p+1})𝒫_a]`$
$``$ $`\frac{i}{2}\stackrel{ˇ}{\theta }_{0,1}\{\pi _0+(\gamma ^a\theta _0)\widehat{P}_a+i\frac{1}{2}(\gamma ^a\theta _0)\theta _0\gamma _a\theta _0^{}`$
$`+\frac{2}{3}\gamma ^a\stackrel{~}{\theta }(2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\frac{5}{4}\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{3}{4}R_a^{}|_>)\pi ^{1,1}+\gamma ^a\mathrm{\Pi }\theta ^{1,1}𝒫_a\}`$
$``$ $`{\displaystyle \underset{p>1}{}}\stackrel{ˇ}{\pi }_{p,p+1}[\left\{i(c+R^{}+\frac{1}{2}\overline{\theta }\stackrel{~}{\gamma }^{}\theta )|_>\pi ^{p,p+1}\right\}^{}`$
$`+(1)^{p+1}i\sqrt{p+1}\vartheta ^p\left\{\frac{1}{2}𝒫^2+\frac{1}{2}\widehat{P}^2+i\overline{\theta }^{}\pi +i[\overline{\theta }(a^{}a_{}a^{}a_{})\pi ]^{}\right\}`$
$`+\frac{1}{2}R^{}|_>{\displaystyle \underset{r}{}}(1)^{r^2+p+1}\sqrt{{\displaystyle \frac{p+1}{r}}}\gamma ^a\theta ^{r1,r}(\mathrm{\Theta }_{rp1}+\mathrm{\Theta }_{rp2})𝒫_a`$
$`(1)^{p+1}i\sqrt{p+1}(\theta ^{p,p}\theta ^{p+1,p+1})`$
$`\times \left(\frac{1}{2}\widehat{P}^2+i\overline{\theta }^{}\pi +i[\overline{\theta }(a^{}a_{}a^{}a_{})\pi ]^{}\right)`$
$`+2i[(p+1)\sqrt{p+1}{\displaystyle \frac{\sqrt{p}+\sqrt{p+2}}{2}}]\pi ^{p+1,p}b`$
$`(1)^{p+1}i\frac{1}{2}\sqrt{p+1}\gamma ^a\mathrm{\Pi }(\pi ^{p,p}+\pi ^{p+1,p+1})𝒫_a`$
$`+\frac{1}{4}{\displaystyle \underset{r}{}}(1)^{r^2+p+1}\sqrt{{\displaystyle \frac{p+1}{r}}}\gamma ^a\theta ^{r1,r}(\mathrm{\Theta }_{rp1}+\mathrm{\Theta }_{rp2})`$
$`\times (\overline{\pi }\gamma ^{}\gamma ^a\theta |_>+i\overline{\theta }^{}\gamma ^a(\pi _0+(\gamma ^b\theta _0)\widehat{P}_b+i\frac{1}{2}(\gamma ^b\theta _0)\theta _0\gamma _b\theta _0^{})+i\overline{q}^{}\gamma _a\stackrel{~}{\theta })]`$
$``$ $`\frac{i}{2}\stackrel{ˇ}{\pi }_{0,1}\{\gamma ^a(\pi _0+(\gamma ^b\theta _0)\widehat{P}_b+i\frac{1}{2}(\gamma ^b\theta _0)\theta _0\gamma _b\theta _0^{})`$ (E.3)
$`\times \left(\widehat{P}_a+i\theta _0\gamma _a\theta _0^{}+2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{1}{2}R_a^{}|_>\right)`$
$`+\frac{2}{3}\gamma ^b\gamma ^a\stackrel{~}{\theta }\left(2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\frac{5}{4}\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{3}{4}R_a^{}|_>\right)`$
$`\times \left(\widehat{P}_b+i\theta _0\gamma _b\theta _0^{}+2i\stackrel{~}{\theta }\gamma _b\theta _0^{}+i\stackrel{~}{\theta }\gamma _b\stackrel{~}{\theta }^{}\frac{1}{2}R_b^{}|_>\right)`$
$`\gamma ^a\mathrm{\Pi }\pi ^{1,1}\left(\widehat{P}_a+i\theta _0\gamma _a\theta _0^{}+2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{1}{2}R_a^{}|_>\right)`$
$`+2\theta ^{1,1}\left(\frac{1}{2}\widehat{P}^2i\overline{\theta }^{}\pi i[\overline{\theta }(a^{}a_{}a^{}a_{})\pi ]^{}\right)`$
$`+2\vartheta ^0\times \{\frac{1}{2}(\widehat{P}_a+i\theta _0\gamma _a\theta _0^{}+2i\stackrel{~}{\theta }\gamma _a\theta _0^{}+i\stackrel{~}{\theta }\gamma _a\stackrel{~}{\theta }^{}\frac{1}{2}R_a^{}|_>)^2`$
$`+\frac{1}{2}\widehat{P}^2+i\overline{\theta }^{}\pi +i[\overline{\theta }(a^{}a_{}a^{}a_{})\pi ]^{}\}`$
$`+\frac{1}{2}R^{}|_>{\displaystyle \underset{r}{}}(1)^{r^2+1}{\displaystyle \frac{1}{\sqrt{r}}}\gamma ^a\theta ^{r1,r}(\mathrm{\Theta }_{r1}+\mathrm{\Theta }_{r2})𝒫_a`$
$`\frac{i}{4}{\displaystyle \underset{r}{}}(1)^{r^2+1}\frac{1}{\sqrt{r}}\gamma ^a\theta ^{r1,r}(\mathrm{\Theta }_{r1}+\mathrm{\Theta }_{r2})`$
$`\times (\overline{\pi }\gamma ^{}\gamma ^a\theta |_>+i\overline{\theta }^{}\gamma ^a(\pi _0+(\gamma ^b\theta _0)\widehat{P}_b+i\frac{1}{2}(\gamma ^b\theta _0)\theta _0\gamma _b\theta _0^{})+i\overline{q}^{}\gamma _a\stackrel{~}{\theta })\}`$
$``$ $`{\displaystyle \underset{qp+1}{}}\stackrel{ˇ}{\theta }_{p,q}[\theta ^{p,q},Q_{sstring}\}`$
$``$ $`{\displaystyle \underset{qp+1}{}}\stackrel{ˇ}{\pi }_{p,q}[\pi ^{p,q},Q_{sstring}\}`$
where
$$\mathrm{\Theta }_x=\{\begin{array}{cc}1& x0\\ 0& x<0\end{array}$$
(E.4)
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# Gaussian measures of entanglement versus negativities: the ordering of two–mode Gaussian states
## I Introduction
Quantum information with continuous variables (CV) book ; brareview is a flourishing field dedicated to the manipulation of the information using quantum states governed by the laws of quantum mechanics. This approach contrasts with the usual methods involving discrete-spectrum observables (such as, e. g., polarization, spin, energy level) of single photons, atoms or ions. The ability of quantum states with continuous spectra to implement quantum cryptography cvcrypto , quantum teleportation cvtelep ; furuscience ; dernier ; bra00 ; telepoppate , entanglement swapping cvswap ; dernier , dense coding cvdense , quantum state storage memory , and, to some extent, quantum computation cvqc processes, brings up new and exciting perspectives.
The crucial resource enabling a better-than-classical manipulation and processing of information is CV entanglement, introduced for the first time in the landmark paper by Einstein, Podolski and Rosen epr in 1935. There, it was shown that the simultaneous eigenstate of relative position and total momentum of two particles (or of a two modes of the radiation field) contains perfect quantum correlations, *i.e.* infinite CV entanglement. While this state is clearly an unphysical, unnormalizable state, it can be approximated arbitrarily well by two-mode squeezed Gaussian states with large enough squeezing parameter. The special class of Gaussian states (which includes thermal, coherent, and squeezed states), thus emerges quite naturally in the CV scenario. These entangled states can be easily produced and manipulated experimentally, and moreover their mathematical description is greatly simplified due to the fact that, while still living in a infinite-dimensional Hilbert space, their relevant properties (such as entanglement and mixedness) are completely determined by the finite-dimensional covariance matrix of two-point correlations between the canonically conjugated quadrature operators. Therefore, clarifying the characterization and quantification of CV entanglement in two-mode and, eventually, multimode Gaussian states stands as a major issue in the field of CV quantum information, as the amount of entanglement contained in a certain state directly quantifies its usefulness for information and communication tasks like teleportation telepoppate .
For the prototypical entangled states of a CV system, the two–mode Gaussian states, much is known about entanglement qualification, as the separability is completely characterized by the necessary and sufficient PPT criterion (positivity of the partially transposed state) simon00 ), and also with regard to its quantifcation. Concerning the latter aspect, the negativity (quantifying the violation of the necessary and sufficient PPT condition for separability) is computable for all two–mode Gaussian states logneg . Moreover, for symmetric two-mode Gaussian states also the entanglement of formation is computable giedke03 , and it turns out to be completely equivalent to the negativity for these states.
Another measure of CV entanglement, adapted for the class of Gaussian states, has been introduced in Ref. GEOF , where the Gaussian entanglement of formation (an upper bound to the true entanglement of formation) was defined as the cost of producing an entangled mixed state out of an ensemble of pure, Gaussian states. While the Gaussian entanglement of formation coincides with the true one for symmetric states, at present it is not known whether this equality holds for nonsymmetric states as well openprob .
In this work, aimed at sheding new light on the quantification of entanglement in two–mode Gaussian states, we compute the Gaussian entanglement of formation and, in general, the family of Gaussian entanglement measures, for two different classes of two–mode Gaussian states, namely the states of extremal, maximal and minimal, negativities at fixed global and local purities prl ; extremal . We find that the two families of entanglement measures (negativities and Gaussian measures) are not equivalent for nonsymmetric states. Remarkably, they may induce a completely different ordering on the set of entangled two–mode Gaussian state: a nonsymmetric state $`\varrho _A`$ can be more entangled than another state $`\varrho _B`$, with respect to negativities, and less entangled than the same state $`\varrho _B`$, with respect to Gaussian measures of entanglement. However, the inequivalence between the two families of measures is somehow bounded: we show that, at fixed negativities, the Gaussian entanglement measures are rigorously bounded from below. Moreover, we provide strong evidence hinting that they should be bounded from above as well.
The paper is organized as follows. In Sec. II we set up the notation and review the basic properties of Gaussian states of CV systems. In Sec. III we review the main results on the characterization of separability in Gaussian states, introducing also two families of measures of entanglement, respectively the negativities and the Gaussian entanglement measures. In Sec. IV we compute the latter for two–mode Gaussian states, solving the problem explicitely for the states of extremal negativities at fixed purities, described in Sec. IV.1. In Sec. V we compare the orderings induced by negativities and Gaussian measures on the set of extremal two–mode Gaussian states. In Sec. VI we compare the two families of measures for generic two–mode Gaussian states, finding lower and upper bounds on one of them, when keeping the other fixed. Finally, in Sec. VII we summarize our results and discuss future perspectives.
## II Gaussian states: Definitions and notation
A continuous variable (CV) system is described by a Hilbert space $`=_{i=1}^n_i`$ resulting from the tensor product structure of infinite dimensional Fock spaces $`_i`$’s. Let $`a_i`$ be the annihilation operator acting on $`_i`$, and $`\widehat{q}_i=(a_i+a_i^{})`$ and $`\widehat{p}_i=(a_ia_i^{})/i`$ be the related quadrature phase operators. The corresponding phase space variables will be denoted by $`q_i`$ and $`p_i`$. Let $`\widehat{X}=(\widehat{q}_1,\widehat{p}_1,\mathrm{},\widehat{q}_n,\widehat{p}_n)`$ denote the vector of the operators $`\widehat{q}_i`$ and $`\widehat{p}_i`$. The canonical commutation relations for the $`\widehat{X}_i`$ can be expressed in terms of the symplectic form $`\mathrm{\Omega }`$
$$[\widehat{X}_i,\widehat{X}_j]=2i\mathrm{\Omega }_{ij},$$
$$\mathrm{with}\mathrm{\Omega }\underset{i=1}{\overset{n}{}}\omega ,\omega \left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).$$
The states of a CV system can be equivalently described by a positive trace-class operator (the density matrix $`\varrho `$) or by quasi–probability distributions such as the Wigner function barnett . States with Gaussian characteristic functions and quasi–probability distributions are referred to as Gaussian states. Such states are at the heart of information processing in CV systems brareview and are the subject of our analysis. By definition, a Gaussian state $`\varrho `$ is completely characterized by the first and second statistical moments of the quadrature field operators, which will be denoted, respectively, by the vector of first moments $`\overline{X}(\widehat{X}_1,\widehat{X}_1,\mathrm{},\widehat{X}_n,\widehat{X}_n)`$ and the covariance matrix (CM) $`𝝈`$ of elements
$$𝝈_{ij}\frac{1}{2}\widehat{X}_i\widehat{X}_j+\widehat{X}_j\widehat{X}_i\widehat{X}_i\widehat{X}_j,$$
(1)
where, for any observable $`\widehat{o}`$, the expectation value $`\widehat{o}\mathrm{Tr}(\varrho \widehat{o})`$. Notice that the entries of the CM can be expressed as energies by multiplying them by the quantity $`\mathrm{}\omega `$, where $`\omega `$ is the frequency of the considered mode. In fact, for any $`n`$-mode state (even non Gaussian) the quantity $`\mathrm{}\omega \mathrm{Tr}𝝈/4`$ is simply the average of the non interacting Hamiltonian $`_{i=1}^n(a_i^{}a_i+1/2)`$. First moments can be arbitrarily adjusted by local unitary operations (displacements), which cannot affect any property related to entropy or entanglement. Therefore, they will be unimportant to the present scope and we will set them to $`0`$ in the following, without any loss of generality.
The canonical commutation relations and the positivity of the density matrix $`\varrho `$ imply
$$𝝈+i\mathrm{\Omega }0,$$
(2)
Inequality (2) is the necessary and sufficient constraint the matrix $`𝝈`$ has to fulfill to be a CM corresponding to a physical Gaussian state simon87 ; simon . More in general, the previous condition is necessary for the CM of any, generally non Gaussian, state. We note that such a constraint implies $`𝝈0`$.
A major role in the theoretical and experimental manipulation of Gaussian states is played by unitary operations which preserve the Gaussian character of the states on which they act. Such operations are all those generated by Hamiltonian terms at most quadratic in the field operators. As a consequence of the Stone-Von Neumann theorem, any such unitary operation at the Hilbert space level corresponds, in phase space, to a symplectic transformation, *i.e.* to a linear transformation $`S`$ which preserves the symplectic form $`\mathrm{\Omega }`$, so that $`\mathrm{\Omega }=S^T\mathrm{\Omega }S`$. Symplectic transformations on a $`2n`$-dimensional phase space form the (real) symplectic group $`Sp_{(2n,)}`$. Such transformations act linearly on first moments and by congruences on covariance matrices: $`𝝈S^𝖳𝝈S`$. One has $`\mathrm{Det}S=1`$, $`SSp_{(2n,)}`$. Ideal beam splitters, phase shifters and squeezers are all described by some kind of symplectic transformation. A particularly important symplectic transformation is the one realizing the decomposition of a Gaussian state in normal modes. Through this decomposition, thanks to Williamson theorem williamson36 , the CM of a $`n`$–mode Gaussian state can always be written in the so-called Williamson normal, or diagonal form
$$𝝈=S^T𝝂S,$$
(3)
where $`SSp_{(2n,)}`$ and $`𝝂`$ is the CM
$$𝝂=\mathrm{diag}(\nu _1,\nu _1,\mathrm{},\nu _n,\nu _n),$$
(4)
corresponding to a tensor product of thermal states with a diagonal density matrix $`\varrho ^{_{}}`$ given by
$$\varrho ^{_{}}=\underset{i}{}\frac{2}{\nu _i+1}\underset{k=0}{\overset{\mathrm{}}{}}\left(\frac{\nu _i1}{\nu _i+1}\right)|k_i_ik|,$$
(5)
where $`|k_i`$ denotes the number state of order $`k`$ in the Fock space $`_i`$.
The quantities $`\nu _i`$’s form the symplectic spectrum of the CM $`𝝈`$, and they can be computed as the eigenvalues of the matrix $`|i\mathrm{\Omega }𝝈|`$. Such eigenvalues are in fact invariant under the action of symplectic transformations on the matrix $`𝝈`$. The symplectic eigenvalues $`\nu _i`$ encode essential informations on the Gaussian state $`𝝈`$ and provide powerful, simple ways to express its fundamental properties. For instance, in terms of the symplectic eigenvalues $`\nu _i`$, the uncertainty relation (2) reads
$$\nu _i1.$$
(6)
Moreover, the entropic quantities of Gaussian states can be as well expressed in terms of their symplectic eigenvalues and invariants extremal . Notably, the purity $`\mathrm{Tr}\varrho ^2`$ of a Gaussian state $`\varrho `$ is simply given by the symplectic invariant $`\mathrm{Det}𝝈=_{i=1}^n\nu _i^2`$, being mga
$$\mu \mathrm{Tr}\varrho ^2=\frac{1}{\sqrt{\mathrm{Det}𝝈}}.$$
(7)
### II.1 Two–mode states
This work is focused on two–mode Gaussian states: we thus briefly review here some of their basic properties. The expression of the two–mode CM $`𝝈`$ in terms of the three $`2\times 2`$ matrices $`𝜶`$, $`𝜷`$, $`𝜸`$, that will be useful in the following, takes the form
$$𝝈\left(\begin{array}{cc}𝜶& 𝜸\\ 𝜸^𝖳& 𝜷\end{array}\right).$$
(8)
For any two–mode CM $`𝝈`$ there is a local symplectic operation $`S_l=S_1S_2`$ which brings $`𝝈`$ in the so called standard form $`𝝈_{sf}`$ simon00 ; duan00
$$S_l^T𝝈S_l=𝝈_{sf}\left(\begin{array}{cccc}a& 0& c_+& 0\\ 0& a& 0& c_{}\\ c_+& 0& b& 0\\ 0& c_{}& 0& b\end{array}\right).$$
(9)
States whose standard form fulfills $`a=b`$ are said to be symmetric. Let us recall that any pure state ($`\mu =1`$) is symmetric and fulfills $`c_+=c_{}=\sqrt{a^21}`$. The correlations $`a`$, $`b`$, $`c_+`$, and $`c_{}`$ are determined by the four local symplectic invariants $`\mathrm{Det}𝝈=(abc_+^2)(abc_{}^2)`$, $`\mathrm{Det}𝜶=a^2`$, $`\mathrm{Det}𝜷=b^2`$, $`\mathrm{Det}𝜸=c_+c_{}`$. Therefore, the standard form corresponding to any CM is unique (up to a common sign flip in $`c_{}`$ and $`c_+`$).
For two–mode states, the uncertainty principle Ineq. (2) can be recast as a constraint on the $`Sp_{(4,)}`$ invariants $`\mathrm{Det}𝝈`$ and $`\mathrm{\Delta }(𝝈)=\mathrm{Det}𝜶+\mathrm{Det}𝜷+2\mathrm{Det}𝜸`$ serafozzi :
$$\mathrm{\Delta }(𝝈)1+\mathrm{Det}𝝈.$$
(10)
The symplectic eigenvalues of a two–mode Gaussian state will be denoted as $`\nu _{}`$ and $`\nu _+`$, with $`\nu _{}\nu _+`$, with the uncertainty relation (6) reducing to
$$\nu _{}1.$$
(11)
A simple expression for the $`\nu _{}`$ can be found in terms of the two $`Sp_{(4,)}`$ invariants (invariants under global, two–mode symplectic operations) logneg ; serafozzi
$$2\nu _{}^2=\mathrm{\Delta }(𝝈)\sqrt{\mathrm{\Delta }^2(𝝈)4\mathrm{Det}𝝈}.$$
(12)
## III Entanglement of Gaussian states
In this section we recall the main results on the qualification and quantification of entanglement for Gaussian states of CV systems.
### III.1 Qualification: PPT criterion
The positivity of the partially transposed state (Peres-Horodecki PPT criterion PPT ) is necessary and sufficient for the separability of two–mode Gaussian states simon00 and, more generally, of all $`(1+n)`$–mode Gaussian states under $`1\times n`$-mode bipartitions werwolf and of symmetric and bisymmetric $`(m+n)`$–mode Gaussian states under $`m\times n`$-mode bipartitions unitarily . In general, the partial transposition $`\stackrel{~}{\varrho }`$ of a bipartite quantum state $`\varrho `$ is defined as the result of the transposition performed on only one of the two subsystems in some given basis. In phase space, the action of partial transposition amounts to a mirror reflection of one of the four canonical variables simon00 . The CM $`𝝈`$ is then transformed into a new matrix $`\stackrel{~}{𝝈}`$ which differs from $`𝝈`$ by a sign flip in $`\mathrm{Det}𝜸`$. Therefore the invariant $`\mathrm{\Delta }(𝝈)`$ is changed into $`\stackrel{~}{\mathrm{\Delta }}(𝝈)\mathrm{\Delta }(\stackrel{~}{𝝈})=\mathrm{Det}𝜶+\mathrm{Det}𝜷2\mathrm{Det}𝜸`$. Now, the symplectic eigenvalues $`\stackrel{~}{\nu }_{}`$ of $`\stackrel{~}{𝝈}`$ read
$$\stackrel{~}{\nu }_{}=\sqrt{\frac{\stackrel{~}{\mathrm{\Delta }}(𝝈)\sqrt{\stackrel{~}{\mathrm{\Delta }}^2(𝝈)4\mathrm{Det}𝝈}}{2}}.$$
(13)
The PPT criterion for separability thus reduces to a simple inequality that must be satisfied by the smallest symplectic eigenvalue $`\stackrel{~}{\nu }_{}`$ of the partially transposed state
$$\stackrel{~}{\nu }_{}1,$$
(14)
which is equivalent to
$$\stackrel{~}{\mathrm{\Delta }}(𝝈)\mathrm{Det}𝝈+1.$$
(15)
Moreover, the above inequalities imply $`\mathrm{Det}𝜸=c_+c_{}<0`$ as a necessary condition for a two–mode Gaussian state to be entangled. Therefore, the quantity $`\stackrel{~}{\nu }_{}`$ encodes all the qualitative characterization of the entanglement for arbitrary (pure or mixed) two–mode Gaussian states.
### III.2 Negativities
From a quantitative point of view, a measure of entanglement which can be computed for general Gaussian states is provided by the negativity $`𝒩`$, first introduced in Ref. zircone , later thoroughly discussed and extended in Refs. logneg ; jenstesi to CV systems. The negativity of a quantum state $`\varrho `$ is defined as
$$𝒩(\varrho )=\frac{\stackrel{~}{\varrho }_11}{2},$$
(16)
where $`\stackrel{~}{\varrho }`$ is the partially transposed density matrix and $`\widehat{o}_1=\mathrm{Tr}|\widehat{o}|`$ stands for the trace norm of the hermitian operator $`\widehat{o}`$. The quantity $`𝒩(\varrho )`$ is equal to $`|_i\lambda _i|`$, the modulus of the sum of the negative eigenvalues of $`\stackrel{~}{\varrho }`$, quantifying the extent to which $`\stackrel{~}{\varrho }`$ fails to be positive. Strictly related to $`𝒩`$ is the logarithmic negativity $`E_𝒩`$, defined as $`E_𝒩\mathrm{log}\stackrel{~}{\varrho }_1`$, which constitutes an upper bound to the distillable entanglement of the quantum state $`\varrho `$ and is related to the entanglement cost under PPT preserving operations auden03 . Both the negativity and the logarithmic negativity have been proven to be monotone under LOCC (local operations and classical communications) logneg ; jenstesi ; plenio05 , a crucial property for a bona fide measure of entanglement. Moreover, the logarithmic negativity possesses the nice property of being additive.
For any two–mode Gaussian state $`\varrho `$ it is easy to show that both the negativity and the logarithmic negativity are simple decreasing functions of $`\stackrel{~}{\nu }_{}`$ logneg ; extremal
$$\stackrel{~}{\varrho }_1=\frac{1}{\stackrel{~}{\nu }_{}}𝒩(\varrho )=\mathrm{max}[0,\frac{1\stackrel{~}{\nu }_{}}{2\stackrel{~}{\nu }_{}}],$$
(17)
$$E_𝒩(\varrho )=\mathrm{max}[0,\mathrm{log}\stackrel{~}{\nu }_{}].$$
(18)
These expressions directly quantify the amount by which the necessary and sufficient PPT condition (14) for separability is violated. The symplectic eigenvalue $`\stackrel{~}{\nu }_{}`$ thus completely qualifies and quantifies (in terms of negativities) the entanglement of a two–mode Gaussian state $`𝝈`$: for $`\stackrel{~}{\nu }_{}1`$ the state is separable, otherwise it is entangled. Finally, in the limit of vanishing $`\stackrel{~}{\nu }_{}`$, the negativities grow unboundedly.
### III.3 Entanglement of Formation
In the special instance of symmetric two–mode Gaussian states, the entanglement of formation (EoF) bennet96 , can be computed as well giedke03 . We recall that the EoF $`E_F`$ of a quantum state $`\varrho `$ is defined as
$$E_F(\varrho )=\underset{\{p_i,|\psi _i\}}{\mathrm{min}}\underset{i}{}p_iE(|\psi _i),$$
(19)
where the minimum is taken over all the pure states realizations of $`\varrho `$:
$$\varrho =\underset{i}{}p_i|\psi _i\psi _i|.$$
The asymptotic regularization of the entanglement of formation coincides with the entanglement cost $`E_C(\varrho )`$, defined as the minimum number of singlets (maximally entangled antisymmetric two-qubit states) which is needed to prepare the state $`\varrho `$ through LOCC ecost .
The optimal convex decomposition of Eq. (19) has been found for symmetric two–mode Gaussian states, and turns out to be Gaussian, that is, the absolute minimum is realized within the set of pure two–mode Gaussian states giedke03 , yielding
$$E_F=\mathrm{max}[0,h(\stackrel{~}{\nu }_{})],$$
(20)
with
$$h(x)=\frac{(1+x)^2}{4x}\mathrm{log}\left[\frac{(1+x)^2}{4x}\right]\frac{(1x)^2}{4x}\mathrm{log}\left[\frac{(1x)^2}{4x}\right].$$
(21)
Such a quantity is, again, a monotonically decreasing function of $`\stackrel{~}{\nu }_{}`$, thus providing a quantification of the entanglement of symmetric states equivalent to the one provided by the negativities.
As a consequence of this equivalence, it is tempting to conjecture that there exists a unique quantification of entanglement for two–mode Gaussian states, embodied by the smallest symplectic eigenvalue $`\stackrel{~}{\nu }_{}`$ of the partially transposed CM, and that the different measures simply provide trivial rescalings of the same unique quantification. In particular, the ordering induced on the set of entangled Gaussian state is uniquely defined for the subset of symmetric two–mode states, and it is independent of the chosen measure of entanglement. However, regrettably, in Sec. V we will indeed show that different measures of entanglement induce, in general, different orderings on the set of nonsymmetric two–mode Gaussian states.
### III.4 Gaussian convex-roof extended measures
In this subsection we consider a family of entanglement measures exclusively defined for Gaussian states of CV systems. The formalism of Gaussian entanglement measures (Gaussian EMs) has been introduced in Ref. GEOF where the Gaussian EoF has been defined and analyzed. Furthermore, the framework developed in Ref. GEOF is general and enables to define generic Gaussian EMs of bipartite entanglement by applying the Gaussian convex roof, that is, the convex roof over pure Gaussian decompositions only, to any bona fide measure of bipartite entanglement defined for pure Gaussian states. The original motivation for the introduction of Gaussian EMs stems from the unfortunate fact that the optimization problem Eq. (19) for the computation of the EoF of nonsymmetric two–mode Gaussian states has not yet been solved, and it stands as an open problem in the theory of entanglement openprob . However, the task can be somehow simplified by restricting to decompositions into pure Gaussian states only. The resulting measure, named as Gaussian EoF in Ref. GEOF , is an upper bound to the true EoF and coincides with it for symmetric two–mode Gaussian states.
In general, we can define a Gaussian EM $`G_E`$ as follows. For any pure Gaussian state $`\psi `$ with CM $`𝝈^P`$, one has
$$G_E(𝝈^P)E(\psi ),$$
(22)
where $`E`$ can be any proper measure of entanglement of pure states, defined as a monotonically increasing function of the entropy of entanglement (*i.e.* the von Neumann entropy of the reduced density matrix of one party).
For any mixed Gaussian state $`\varrho `$ with CM $`𝝈`$, one has GEOF
$$G_E(𝝈)\underset{𝝈^P𝝈}{inf}G_E(𝝈^P).$$
(23)
If the function $`E`$ is taken to be exactly the entropy of entanglement, then the corresponding Gaussian EM is known as Gaussian EoF GEOF . In Ref. jensata the properties of the Gaussian EoF have been further investigated, and interesting connections with the capacity of bosonic Gaussian channels have been established.
In general, the definition Eq. (23) involves an optimization over all pure Gaussian states with CM $`𝝈^P`$ smaller than the CM $`𝝈`$ of the mixed state whose entanglement one wishes to compute. Despite being a simpler optimization problem than that appearing in the definition Eq. (19) of the true EoF (which, in CV systems, would imply considering decompositions over all, Gaussian and non-Gaussian pure states), the Gaussian EMs cannot be expressed in a simple closed form, not even in the simplest instance of (nonsymmetric) two–mode Gaussian states. It is the aim of the present paper to compute Gaussian EMs for two relevant classes of, generally nonsymmetric, two–mode Gaussian states, namely the states of extremal (maximal and minimal) negativity at fixed global and local purities prl ; extremal , which will be reviewed in Sec. IV.1. This will provide an insight into the problem of the ordering ordering of two–mode Gaussian states with respect to different measures of entanglement, leading to results somehow similar to those obtained for systems of two qubits frank , where in general the EoF and the negativity are found to be inequivalent.
Before moving on to the explicit computations, let us recall, as an important side remark, that any Gaussian EM is an entanglement monotone under Gaussian LOCC. The proof given in Sec. IV of Ref. GEOF for the Gaussian EoF, in fact, automatically extends to every Gaussian EM constructed via the Gaussian convex roof of any proper measure $`E`$ of pure-state entanglement.
## IV Gaussian entanglement measures <br>for two–mode Gaussian states
The problem of evaluating Gaussian EMs for a generic two–mode Gaussian state has been solved in Ref. GEOF . However, the explicit result contains so “cumbersome” expressions (involving the solutions of a fourth-order algebraic equation), that the authors of Ref. GEOF considered them not particularly useful to be reported explicitely in their paper.
We recall here the computation procedure GEOF that we will need in the following. For any two-mode Gaussian state with CM $`𝝈𝝈_{sf}`$ in standard form Eq. (9), a generic Gaussian EM $`G_E`$ is given by the entanglement $`E`$ of the least entangled pure state with CM $`𝝈^P𝝈`$. Denoting by $`\gamma _q`$ (respectively $`\gamma _p`$) the $`2\times 2`$ submatrix obtained from $`𝝈`$ by canceling the even (resp. odd) rows and columns, we have, explicitely
$$\gamma _q=\left(\begin{array}{cc}a\hfill & c_+\hfill \\ c_+\hfill & b\hfill \end{array}\right),\gamma _p=\left(\begin{array}{cc}a\hfill & c_{}\hfill \\ c_{}\hfill & b\hfill \end{array}\right).$$
(24)
All the covariances relative to the “position” operators of the two modes are grouped in $`\gamma _q`$, and analogously for the “momentum” operators in $`\gamma _p`$. The total CM can then be written as a direct sum $`𝝈=\gamma _q\gamma _p`$. Similarly, the CM of a generic pure two–mode Gaussian state in standard form (it has been proven that the CM of the optimal pure state has to be in standard form as well GEOF ) can be written as $`𝝈^P=\gamma _q^P\gamma _p^P`$, where the global purity of the state imposes $`(\gamma _p^P)^1=\gamma _q^P\mathrm{\Gamma }`$. The pure states involved in the definition of the Gaussian EM must thus fulfill the condition
$$\gamma _p^1\mathrm{\Gamma }\gamma _q.$$
(25)
This problem is endowed with a nice geometric description GEOF . Writing the matrix $`\mathrm{\Gamma }`$ in the basis constituted by the identity matrix and the three Pauli matrices,
$$\mathrm{\Gamma }=\left(\begin{array}{cc}x_0+x_3& x_1\\ x_1& x_0x_3\end{array}\right),$$
(26)
the expansion coefficients $`(x_0,x_1,x_3)`$ play the role of space-time coordinates in a three-dimensional Minkowski space. In this picture, for example, the rightmost inequality in Eq. (25) is satisfied by matrices $`\mathrm{\Gamma }`$ lying on a cone, which is equivalent to the (backwards) light cone of $`C_q`$ in the Minkowski space; and similarly for the leftmost inequality. Indeed, one can show that, for the optimal pure state $`𝝈_{opt}^P`$ realizing the minimum in Eq. (23), the two inequalities in Eq. (25) have to be simultaneously saturated GEOF . From a geometrical point of view, the optimal $`\mathrm{\Gamma }`$ has then to be found on the rim of the intersection of the forward and the backward cones of $`\gamma _p^1`$ and $`\gamma _q`$, respectively. This is an ellipse, and one is left with the task of minimizing the entanglement $`E`$ of $`𝝈^P=\mathrm{\Gamma }\mathrm{\Gamma }^1`$ (see Eq. (22)) for $`\mathrm{\Gamma }`$ lying on this ellipse noteole .
At this point, let us pause to briefly recall that any pure two–mode Gaussian state $`𝝈^P`$ is locally equivalent to a two–mode squeezed state with squeezing parameter $`r`$, described by a CM
$$𝝈_{sq}^P=\left(\begin{array}{cccc}\mathrm{cosh}(2r)& 0& \mathrm{sinh}(2r)& 0\\ 0& \mathrm{cosh}(2r)& 0& \mathrm{sinh}(2r)\\ \mathrm{sinh}(2r)& 0& \mathrm{cosh}(2r)& 0\\ 0& \mathrm{sinh}(2r)& 0& \mathrm{cosh}(2r)\end{array}\right).$$
(27)
The following statements are then equivalent: (i) $`E`$ is a monotonically increasing function of the entropy of entanglement; (ii) $`E`$ is a monotonically increasing function of the single–mode determinant $`m\mathrm{Det}𝜶\mathrm{Det}𝜷`$ (see Eq. (8)); (iii) $`E`$ is a monotonically decreasing function of the local purity $`\mu _i\mu _1\mu _2`$ (see Eq. (7)); (iv) $`E`$ is a monotonically decreasing function of the smallest symplectic eigenvalue $`\stackrel{~}{\nu }_{}^P`$ of the partially transposed CM $`\stackrel{~}{𝝈}^P`$; (v) $`E`$ is a monotonically increasing function of the squeezing parameter $`r`$. This chain of equivalences is immediately proven by simply recalling that a pure state is completely specified by its single–mode marginals, and that for a single–mode Gaussian state there is a unique symplectic invariant (the determinant), so that all conceivable entropic quantities are monotonically increasing functions of this invariant extremal . In particular, statement (ii) allows us to minimize directly the single–mode determinant over the ellipse:
$$m=1+\frac{x_1}{\mathrm{Det}\mathrm{\Gamma }},$$
(28)
with $`\mathrm{\Gamma }`$ given by Eq. (26).
To simplify the calculations, one can move to the plane of the ellipse with a Lorentz boost which preserves the relations between all the cones; one can then choose the transformation so that the ellipse degenerates into a circle (with fixed radius), and introduce polar coordinates on this circle. The calculation of the Gaussian EM for any two–mode Gaussian state is thus finally reduced to the minimization of $`m`$ from Eq. (28), at given standard-form covariances of $`𝝈`$, as a function of the polar angle $`\theta `$ on the circle noteole . So far, this technique has been applied to the computation of the Gaussian EoF by minimizing Eq. (28) numerically GEOF (see also oleposter ). In addition to that, as already mentioned, the Gaussian EoF has been exactly computed for symmetric states, and it has been proven that in this case the Gaussian EoF is the true EoF giedke03 .
In this work we present new analytical calculations of the Gaussian EMs for two relevant classes of nonsymmetric two–mode Gaussian states: the states of extremal negativities at fixed global and local purities extremal , which will be introduced in the next subsection. We begin by writing the general expression of the single–mode determinant Eq. (28) in terms of the covariances of a generic two–mode state (see Eq. (9)) and of the polar angle $`\theta `$. After some tedious but straightforward algebra, one finds
$`m_\theta (a,b,c_+,c_{})=1`$ $`+`$ $`\left\{\left[c_+(abc_{}^2)c_{}+\mathrm{cos}\theta \sqrt{\left[ab(abc_{}^2)\right]\left[ba(abc_{}^2)\right]}\right]^2\right\}`$ (29)
$`\times `$ $`\{2(abc_{}^2)(a^2+b^2+2c_+c_+)`$
$`{\displaystyle \frac{\mathrm{cos}\theta \left[2abc_{}^3+\left(a^2+b^2\right)c_+c_{}^2+\left(\left(12b^2\right)a^2+b^2\right)c_{}ab\left(a^2+b^22\right)c_+\right]}{\sqrt{\left[ab(abc_{}^2)\right]\left[ba(abc_{}^2)\right]}}}`$
$`+\mathrm{sin}\theta (a^2b^2)\sqrt{1{\displaystyle \frac{\left[c_+(abc_{}^2)+c_{}\right]^2}{\left[ab(abc_{}^2)\right]\left[ba(abc_{}^2)\right]}}}\}^1,`$
where we have assumed $`c_+|c_{}|`$ without any loss of generality. This implies that, for any entangled state, $`c_+>0`$ and $`c_{}<0`$. The Gaussian EM (defined in terms of the function $`E`$ on pure states, see Eq. (22)) of a generic two–mode Gaussian state coincides then with the entanglement $`E`$ computed on the pure state with $`m=m_{opt}`$, with $`m_{opt}\mathrm{min}_\theta (m_\theta )`$. Accordingly, the symplectic eigenvalue $`\stackrel{~}{\nu }_{}`$ of the partial transpose of the corresponding optimal pure-state CM $`𝝈_{opt}^P`$, realizing the infimum in Eq. (23), would read (see Eq. Eq. (13))
$$\stackrel{~}{\nu }_{opt}^P\stackrel{~}{\nu }_{}(𝝈_{opt}^P)=\sqrt{m_{opt}}\sqrt{m_{opt}1}.$$
(30)
As an example, for the Gaussian EoF one has
$$G_{E_F}(𝝈)=h\left(\stackrel{~}{\nu }_{opt}^P(m_{opt})\right),$$
(31)
with $`h(x)`$ defined by Eq. (21).
Finding the minimum of Eq. (29) analytically for a generic state is a difficult task. Numerical investigations show that the equation $`_\theta m_\theta =0`$ can have from one to four physical solutions (in a period) corresponding to extremal points, and the global minimum can be attained in any of them depending on the parameters of the CM $`𝝈`$ under inspection. However, a closed solution can be found for two important classes of nonsymmetric two–mode Gaussian states, as we will now show.
### IV.1 Parametrization of two–mode covariance matrices and definition of extremal states
We have shown in Refs. prl ; extremal that, at fixed global purity $`\mu \mathrm{Tr}\varrho ^2`$ of the global state $`\varrho `$, and at fixed local purities $`\mu _{1,2}\mathrm{Tr}\varrho _{1,2}^2`$ of each of the two reduced single–mode states $`\varrho _i=\mathrm{Tr}_{ji}\varrho `$, the smallest symplectic eigenvalue $`\stackrel{~}{\nu }_{}`$ of the partial transpose of the CM $`𝝈`$ of a generic two–mode Gaussian state (which qualifies its separability by the PPT criterion, and quantifies its entanglement in terms of the negativities) is strictly bounded from above and from below. This entails the existence of two disjoint classes of extremal states, namely the states of maximum negativity for fixed global and local purities (GMEMS), and the states of minimum negativity for fixed global and local purities (GLEMS) extremal . The negativities of the two extremal classes of Gaussian states, moreover, turn out to remain very close to each other for all the possible assignments of the three purities, allowing for a reliable experimental estimate of the negativity of a generic two-mode Gaussian state in terms of the average negativity prl . The latter is determined by knowledge of the three purities alone, which, in turn, may be experimentally measured in direct, possibly efficient, ways fiuracerf .
Recalling these results, one can provide a very useful and insightful parametrization of the entangled two–mode Gaussian states in standard form (see also polacchi ). In fact, the coefficients appearing in Eq. (9) can be rewritten, in general, according to the following, useful parametrization:
$`a`$ $`=`$ $`s+d,b=sd,`$ (32)
$`c_\pm `$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{s^2d^2}}}\{\sqrt{\left[4d^2+{\displaystyle \frac{1}{2}}\left(g^2+1\right)(\lambda 1)\left(2d^2+g\right)(\lambda +1)\right]^24g^2}`$ (33)
$`\pm `$ $`\sqrt{\left[4s^2+{\displaystyle \frac{1}{2}}\left(g^2+1\right)(\lambda 1)\left(2d^2+g\right)(\lambda +1)\right]^24g^2}\},`$
where the two local purities are regulated by the parameters $`s`$ and $`d`$, being $`\mu _1=(s+d)^1,\mu _2=(sd)^1`$, and the global purity is $`\mu =g^1`$. The coefficient $`\lambda `$ embodies the only remaining degree of freedom needed for the complete determination of the negativities, once the three purities have been fixed. It ranges from the minimum $`\lambda =1`$ (corresponding to the GLEMS) to the maximum $`\lambda =+1`$ (corresponding to the GMEMS). Therefore, as it varies, $`\lambda `$ encompasses all possible entangled two–mode Gaussian states compatible with a given set of assigned values of the purities. The constraints that the parameters $`s,d,g`$ must obey for Eq. (9) to denote a proper CM of a physical state are: $`s1`$, $`|d|s1`$, and
$$g2|d|+1,$$
(34)
If the global purity is large enough so that Ineq. (34) is saturated, GMEMS and GLEMS coincide, the CM becomes independent of $`\lambda `$, and the two classes of extremal states coalesce into a unique class, completely determined by the marginals $`s`$ and $`d`$. We denote these states as GMEMMS extremal , that is, Gaussian two–mode states of maximal negativity at fixed local purities. Their CM is simply characterized by $`c_\pm =\pm \sqrt{s^2(d+1)^2}`$, where we have assumed, without any loss of generality, that $`d0`$ (corresponding to choose, for instance, mode $`1`$ as the more mixed one: $`\mu _1\mu _2`$).
In general prl , a GMEMS ($`\lambda =+1`$) is entangled for
$$g<2s1,$$
(35)
while a GLEMS ($`\lambda =1`$) is entangled for smaller $`g`$, namely
$$g<\sqrt{2(s^2+d^2)1}.$$
(36)
To have a physical insight on these peculiar two–mode states, let us recall extremal that GMEMS are simply nonsymmetric thermal squeezed states, usually referred to as maximally entangled mixed states in CV systems. On the other hand, GLEMS are mixed states of partial minimum uncertainty, in the sense that the smallest symplectic eigenvalue of their CM is equal to $`1`$, saturating the uncertainty inequality (11).
We are now equipped with the necessary tools, and in the next subsection we move on to compute Gaussian EMs for the two extremal classes of nonsymmetric two–mode Gaussian states, the GLEMS and the GMEMS.
### IV.2 Gaussian entanglement of minimum-negativity states (GLEMS)
We want to find the optimal pure state $`𝝈_{opt}^P`$ entering in the definition Eq. (23) of the Gaussian EM. To do this, we have to minimize the single–mode determinant of $`𝝈_{opt}^P`$, given by Eq. (29), over the angle $`\theta `$. It turns out that, for a generic GLEMS, the coefficient of $`\mathrm{sin}\theta `$ in the last line of Eq. (29) vanishes, and the expression of the single–mode determinant reduces to the simplified form
$$m_\theta ^{_{\mathrm{GLEMS}}}=1+\frac{[A\mathrm{cos}\theta +B]^2}{2(abc_{}^2)[(g^21)\mathrm{cos}\theta +g^2+1]},$$
(37)
with $`A=c_+(abc_{}^2)+c_{},B=c_+(abc_{}^2)c_{},`$ and $`a,b,c_\pm `$ are the covariances of GLEMS, obtained from Eqs. (32,33) setting $`\lambda =1`$.
The only relevant solutions (excluding the unphysical and the trivial ones) of the equation $`_\theta m_\theta =0`$ are $`\theta =\pi `$ and
$$\theta =\pm \theta ^{}\mathrm{arccos}\left[\frac{3+g^2}{1g^2}\frac{2c_{}}{c_+(abc_{}^2)+c_{}}\right].$$
Studying the second derivative $`_\theta ^2m_\theta `$ for $`\theta =\pi `$ one finds immediately that, for
$$g\sqrt{\frac{2c_+(abc_{}^2)+c_{}}{c_{}}}$$
(38)
(remember that $`c_{}0`$), the solution $`\theta =\pi `$ is a minimum. In this range of parameters, the other solution $`\theta =\theta ^{}`$ is unphysical (in fact $`|\mathrm{cos}\theta ^{}|1`$), so $`m_{\theta =\pi }`$ is the global minimum. When, instead, Ineq. (38) is violated, $`m_\theta `$ has a local maximum for $`\theta =\pi `$ and two minima appear at $`\theta =\pm \theta ^{}`$. The global minimum is attained in any of the two, given that, for GLEMS, $`m_\theta `$ is invariant under reflection with respect to the axis $`\theta =\pi `$. Collecting, substituting, and simplifying the obtained expressions, we arrive at the final result for the optimal $`m`$:
$$m_{opt}^{_{\mathrm{GLEMS}}}=\{\begin{array}{cc}1,\hfill & g\sqrt{2(s^2+d^2)1}\text{[separable state]};\hfill \\ & \\ \frac{16s^2d^2}{(g^21)^2},\hfill & \sqrt{\frac{\left(4s^2+1\right)d^2+s^2+4s\sqrt{\left(s^2+1\right)d^2+s^2}|d|}{d^2+s^2}}g<\sqrt{2(s^2+d^2)1};\hfill \\ & \\ \frac{g^4+2\left(2d^2+2s^2+1\right)g^2\left(4d^21\right)\left(4s^21\right)\sqrt{\delta }}{8g^2},\hfill & 2|d|+1g<\sqrt{\frac{\left(4s^2+1\right)d^2+s^2+4s|d|\sqrt{\left(s^2+1\right)d^2+s^2}}{d^2+s^2}}.\hfill \end{array}$$
(39)
Here $`\delta (2dg1)(2dg+1)(2d+g1)(2d+g+1)(g2s1)(g2s+1)(g+2s1)(g+2s+1)`$.
Immediate inspection crucially reveals that $`m_{opt}^{_{\mathrm{GLEMS}}}`$ is not in general a function of the symplectic eigenvalue $`\stackrel{~}{\nu }_{}`$ alone. Therefore, unfortunately, the Gaussian EMs, and in particular, the Gaussian EoF, are not equivalent to the negativities for GLEMS. Further remarks will be given in the following, when the Gaussian EMs of GLEMS and GMEMS will be compared and their relationship with the negativities will be elucidated.
### IV.3 Gaussian entanglement of maximum-negativity states (GMEMS)
The minimization of $`m_\theta `$ from Eq. (29) can be carried out in a simpler way in the case of GMEMS, whose covariances can be retrieved from Eq. (33) setting $`\lambda =1`$. First of all, one can notice that, when expressed as a function of the Minkowski coordinates $`(x_0,x_1,x_3)`$, corresponding to the submatrix $`\mathrm{\Gamma }`$ Eq. (26) of the pure state $`𝝈^P=\mathrm{\Gamma }\mathrm{\Gamma }^1`$ entering in the optimization problem Eq. (23), the single–mode determimant $`m`$ of $`𝝈^P`$ is globally minimized for $`x_3=0`$. In fact, from Eq. (28), $`m`$ is minimal, with respect to $`x_3`$, when $`\mathrm{Det}\mathrm{\Gamma }=x_0^2x_1^2x_3^2`$ is maximal. Next, one can show that for GMEMS there always exists a matrix $`\mathrm{\Gamma }`$, with $`x_3=0`$, which is a simultaneous solution of the two matrix equations obtained by imposing the saturation of the two sides of inequality (25). As a consequence of the above discussion, this matrix would denote the optimal pure state $`𝝈_{opt}^P`$. Solving the system of equations $`\mathrm{Det}(\gamma _q\mathrm{\Gamma })=\mathrm{Det}(\mathrm{\Gamma }\gamma _p^1)=0`$, where the matrices involved are explicitely defined combining Eq. (24) and Eq. (33) with $`\lambda =1`$, one finds the following two solutions for the coordinates $`x_0`$ and $`x_1`$:
$$\begin{array}{cc}\hfill x_0^\pm & =\frac{(g+1)s\pm \sqrt{\left[(g1)^24d^2\right]\left(d^2+s^2g\right)}}{2\left(d^2+g\right)},\hfill \\ \hfill x_1^\pm & =\frac{(g+1)\sqrt{d^2+s^2g}\pm s\sqrt{(g1)^24d^2}}{2\left(d^2+g\right)}.\hfill \end{array}$$
(40)
The corresponding pure state $`𝝈^{P\pm }=\mathrm{\Gamma }^\pm \mathrm{\Gamma }_{}^{\pm }{}_{}{}^{1}`$ turns out to be, in both cases, a two–mode squeezed state described by a CM of the form Eq. (27), with $`\mathrm{cosh}(2r)=x_0^\pm `$. Because the single–mode determinant $`m=\mathrm{cosh}^2(2r)`$ for these states, the optimal $`m`$ for GMEMS is simply equal to $`(x_0^{})^2`$. Summarizing,
$$m_{opt}^{_{\mathrm{GMEMS}}}=\{\begin{array}{cc}1,g2s1\text{[separable state]};\hfill & \\ & \\ \frac{\left\{(g+1)s\sqrt{\left[(g1)^24d^2\right]\left(d^2+s^2g\right)}\right\}^2}{4\left(d^2+g\right)^2},\hfill & \\ \mathrm{\hspace{0.33em}2}|d|+1g<2s1.\hfill & \end{array}$$
(41)
Once again, also for the class of GMEMS the Gaussian EMs are not simple functions of the symplectic eigenvalue $`\stackrel{~}{\nu }_{}`$ alone. Consequently, they provide a quantification of CV entanglement of GMEMS inequivalent to the one determined by the negativities. Furthermore, we will now show how these result raise the problem of the ordering of two–mode Gaussian states according to their degree of entanglement, as quantified by different families of entanglement measures.
## V Extremal ordering of <br>two–mode Gaussian states
Entanglement is a physical quantity. It has a definite mathematical origin within the framework of quantum mechanics, and its conceptual meaning in the end stems from and is rooted in the existence of the superposition principle. Further, entanglement has a fundamental operative interpretation as that resource that in principle enables information processing and communication in better-than-classical realizations telepoppate . One would then expect that, picking two states $`\varrho _A`$ and $`\varrho _B`$ out of a certain (subset of) Hilbert space, the question “Is $`\varrho _A`$ more entangled than $`\varrho _B`$?” should have a unique, well-defined answer, independent of the measure that one chooses to quantify entanglement. But, contrary to the common expectations, this is generally not the case for mixed states. Different measures of entanglement will in general induce different, inequivalent orderings on the set of entangled states belonging to a given Hilbert space ordering , as they usually measure different aspects of quantum correlations existing in generic mixed states.
In the context of CV systems, when one restricts to symmetric, two–mode Gaussian states, which include all pure states, the known computable measures of entanglement all correctly induce the same ordering on the set of entangled states. We will now show that, indeed, this nice feature is not preserved moving to mixed, nonsymmetric two-mode Gaussian states. We aim at comparing Gaussian EMs and negativities on the two extremal classes of two–mode Gaussian states extremal , introducing thus the concept of extremal ordering. At fixed global and local purities, the negativity of GMEMS (which is the maximal one) is obviously always greater than the negativity of GLEMS (which is the minimal one). If for the same values of purities the Gaussian EMs of GMEMS are larger than those of GLEMS, we will say that the extremal ordering is preserved. Otherwise, the extremal ordering is inverted. In this latter case, which is clearly the most intriguing, the states of minimal negativities are more entangled, with respect to Gaussian EMs, than the states of maximal negativities, and the inequivalence of the orderings, induced by the two different families of entanglement measures, becomes manifest.
The problem can be easily stated. By comparing $`m_{opt}^{_{\mathrm{GLEMS}}}`$ from Eq. (39) and $`m_{opt}^{_{\mathrm{GMEMS}}}`$ from Eq. (41), one has that in the range of global and local purities, or, equivalently, of parameters $`\{s,d,g\}`$, such that
$$m_{opt}^{_{\mathrm{GMEMS}}}m_{opt}^{_{\mathrm{GLEMS}}},$$
(42)
the extremal ordering is preserved. When Ineq. (42) is violated, the extremal ordering is inverted. The boundary between the two regions, which can be found imposing the equality $`m_{opt}^{_{\mathrm{GMEMS}}}=m_{opt}^{_{\mathrm{GLEMS}}}`$, yields the range of global and local purities such that the corresponding GMEMS and GLEMS, despite having different negativities, have equal Gaussian EMs. This boundary surface can be found numerically, and the result is shown in the 3D plot of Fig. 1.
One can see, as a crucial result, that a region where the extremal ordering is inverted does indeed exist. The Gaussian EMs and the negativities are thus definitely not equivalent for the quantification of entanglement in nonsymmetric two–mode Gaussian states. The interpretation of this result is quite puzzling. On the one hand, one could think that the ordering induced by the negativities is a natural one, due to the fact that such measures of entanglement are directly inspired by the necessary and sufficient PPT criterion for separability. Thus, one would expect that the ordering induced by the negativities should be preserved by any bona fide measure of entanglement, especially if one considers that the extremal states, GLEMS and GMEMS, have a clear physical interpretation extremal . Therefore, as the Gaussian EoF is an upper bound to the true EoF, one could be tempted to take this result as an evidence that the Gaussian EoF overestimates the true EoF, at least for GLEMS, and that, moreover, the true EoF of GLEMS should be lower than the true EoF of GMEMS, at fixed values of the purities. If this were the case, the true EoF would not coincide with the Gaussian EoF, whose evaluation would consequently necessarily involve a decomposition over non-Gaussian states. However, this is only a qualitative/speculative argument: proving or disproving that the Gaussian EoF is the true EoF for any two–mode Gaussian state is still an open question under lively debate openprob .
On the other hand, one could take the simplest discrete-variable instance, constituted by a two–qubit system, as a test-case for comparison. There, although for pure states the negativity coincides with the concurrence, an entanglement monotone equivalent to the EoF for all states of two qubits Wootters , the two measures cease to be equivalent for mixed states, and the orderings they induce on the set of entangled states can be different frank . This analogy seems to support again the stand that, in the arena of mixed states, a unique measure of entanglement is a chimera and cannot really be expected, due to the different operative meanings and physical processes (in the cases when it has been possible to identify them) that are associated to each definition: one could think, for instance, of the operative difference existing between the definitions of distillable entanglement and entanglement cost. In other words, from this point of view, each inequivalent measure of entanglement introduced for mixed states should capture physically distinct aspects of quantum correlations existing in these states. Then, joining this kind of outlook, one could hope that the Gaussian EMs might still be considered as proper measures of CV entanglement, especially if one were able to prove the conjecture that the Gaussian EoF is the true EoF for a broader class of Gaussian states beyond the symmetric ones. One could then live on with the existence of inverted orderings of entangled states, and see it as a not so annoying problem.
Whatever be the case, we have shown that two different families of measures of CV entanglement can induce different orderings on the set of two–mode entangled states. This is more clearly illustrated in Fig. 2, where we keep fixed one of the local mixednesses and we classify, in the space of the other local mixedness and of the global mixedness, the different regions related to entanglement and extremal ordering of two–mode Gaussian states, improving and completing a similar diagram previously introduced in Ref. prl to describe separability in the space of purities.
## VI Gaussian measures of entanglement <br>versus negativities
In this section we wish to give a more direct comparison of the two families of entanglement measures for two–mode Gaussian states. In particular, we are interested in finding the maximum and minimum values of one of the two measures, if the other is kept fixed. A very similar analysis has been performed by Verstraete et al. frank , in their comparative analysis of the negativity and the concurrence for states of two-qubit systems.
Here it is useful to perform the comparison directly between the symplectic eigenvalue $`\stackrel{~}{\nu }_{}(𝝈)`$ of the partially transposed CM $`\stackrel{~}{𝝈}`$ of a generic two–mode Gaussian state with CM $`𝝈`$, and the symplectic eigenvalue $`\stackrel{~}{\nu }_{}(𝝈_{opt}^P)`$ of the partially transposed CM $`\stackrel{~}{𝝈}_{opt}^P`$ of the optimal pure state with CM $`𝝈_{opt}^P`$, which minimizes Eq. (23). In fact, the negativities are all monotonically decreasing functions of $`\stackrel{~}{\nu }_{}(𝝈)`$, while the Gaussian EMs are all monotonically decreasing functions of $`\stackrel{~}{\nu }_{}(𝝈_{opt}^P)`$.
To start with, let us recall once more that for pure states and for mixed symmetric states (in the set of two–mode Gaussian states), the two quantities coincide. For nonsymmetric states, one can immediately prove the following bound
$$\stackrel{~}{\nu }_{}(𝝈_{opt}^P)\stackrel{~}{\nu }_{}(𝝈).$$
(43)
In fact, from Eq. (23), $`𝝈_{opt}^P𝝈`$ GEOF . For positive matrices, $`AB`$ implies $`a_kb_k`$, where the $`a_k`$s (resp. $`b_k`$s) denote the ordered symplectic eigenvalues of $`A`$ (resp. $`B`$) giedkeqic . Because the ordering $`AB`$ is preserved under partial transposition, Ineq. (43) holds true. This fact induces a characterization of symmetric states, which saturate Ineq. (43), as the two–mode Gaussian states with minimal Gaussian EMs at fixed negativities.
It is then natural to raise the question whether an upper bound on the Gaussian EMs at fixed negativities exists as well. It seems hard to address this question directly, as one lacks a closed expression for the Gaussian EMs of generic states. But we can promptly give partial answers if we restrict to the classes of GLEMS and of GMEMS, for which the Gaussian EMs have been explicitely computed in the previous section.
Let us begin with the GLEMS. We can compute the squared symplectic eigenvalue $`\stackrel{~}{\nu }_{}^2(𝝈^{_{\mathrm{GLEMS}}})=\left[4(s^2+d^2)g^21\sqrt{\left(4(s^2+d^2)g^21\right)^24g^2}\right]/2`$. Next, we can reparametrize the CM (obtained by Eq. (33) with $`\lambda =1`$) to make $`\stackrel{~}{\nu }_{}`$ appear explicitely, namely $`g=\sqrt{\stackrel{~}{\nu }_{}^2[4(s^2+d^2)1\stackrel{~}{\nu }_{}^2]/(1+\stackrel{~}{\nu }_{}^2)}`$. At this point, one can study the piecewise function $`m_{opt}^{_{\mathrm{GLEMS}}}`$ from Eq. (39), and find out that it is a convex function of $`d`$ in the whole space of parameters corresponding to entangled states. Hence, $`m_{opt}^{_{\mathrm{GLEMS}}}`$, and thus the Gaussian EM, is maximized at the boundary $`|d|=(2\stackrel{~}{\nu }_{}s\stackrel{~}{\nu }_{}^21)/2`$, resulting from the saturation of Ineq. (34). The states maximizing Gaussian EMs at fixed negativities, if we restrict to the class of GLEMS, have then to be found in the subclass of GMEMMS (states of maximal negativity for fixed marginals extremal , defined after Ineq. (34)), depending on the parameter $`s`$ and on the eigenvalue $`\stackrel{~}{\nu }_{}`$ itself, which completely determines the negativity). For these states,
$$m_{opt}^{_{\mathrm{GMEMMS}}}(s,\stackrel{~}{\nu }_{})=\left(\frac{2s}{1\stackrel{~}{\nu }_{}^2+2\stackrel{~}{\nu }_{}s}\right)^2.$$
(44)
The further optimization over $`s`$ is straightforward because $`m_{opt}^{_{\mathrm{GMEMMS}}}`$ is an increasing function of $`s`$, so its global maximum is attained for $`s\mathrm{}`$. In this limit, one has simply
$$m_{\mathrm{max}}^{_{\mathrm{GMEMMS}}}(\stackrel{~}{\nu }_{})=\frac{1}{\stackrel{~}{\nu }_{}^2}.$$
(45)
From Eq. (30), one thus finds that for all GLEMS the following bound holds
$$\stackrel{~}{\nu }_{}(𝝈_{opt}^P)\frac{1}{\stackrel{~}{\nu }_{}(𝝈)}\left(1\sqrt{1\stackrel{~}{\nu }_{}^2(𝝈)}\right).$$
(46)
One can of course perform a similar analysis for GMEMS. But, after analogous reasonings and computations, what one finds is exactly the same result. This is not so surprising, keeping in mind that GMEMS, GLEMS and all two–mode Gaussian states with generic $`s`$ and $`d`$ but with global mixedness $`g`$ saturating Ineq. (34), collapse into the same family of two–mode Gaussian states, the GMEMMS, completely determined by the local single–mode properties (they can be viewed as a generalization of the pure two–mode states: the symmetric GMEMMS are in fact pure). Hence, the bound of Ineq. (46), limiting the Gaussian EMs from above at fixed negativities, must hold for all GMEMS as well.
At this point, it is tempting to conjecture that Ineq. (46) holds for all two–mode Gaussian states. Unfortunately, the lack of a closed, simple expression for the Gaussian EM of a generic state makes the proof of this conjecture impossible, at the present time. However, one can show, by analytical power-series expansions of Eq. (29), truncated to the leading order in the infinitesimal increments, that, for any infinitesimal variation of the parameters of a generic CM around the limiting values characterizing GMEMMS, the Gaussian EMs of the resulting states lie always below the boundary imposed by the corresponding GMEMMS with the same $`\stackrel{~}{\nu }_{}`$. In this sense, the GMEMMS are, at least, a local maximum for the Gaussian EM versus negativity problem. Furthermore, extensive numerical investigations of up to a million CMs of randomly generated two–mode Gaussian states, provide confirmatory evidence that GMEMMS attain indeed the global maximum (see Fig. 3). We can thus quite confidently conjecture, however, at the moment, without a complete formal proof of the statement, that GMEMMS, in the limit of infinite average local mixedness ($`s\mathrm{}`$), are the states of maximal Gaussian EMs at fixed negativities, among all two–mode Gaussian states.
A direct comparison between the two prototypical representatives of the two families of entanglement measures, respectively the Gaussian EoF $`G_{E_F}`$ and the logarithmic negativity $`E_𝒩`$, is plotted in Fig. 4. For any fixed value of $`E_𝒩`$, Ineq. (43) provides in fact a rigorous lower bound on $`G_{E_F}`$, namely
$$G_{E_F}h[\mathrm{exp}(E_𝒩)],$$
(47)
while Ineq. (46) provides the conjectured lower bound
$$G_{E_F}h\left[\mathrm{exp}(E_𝒩)\left(1\sqrt{1\mathrm{exp}(2E_𝒩)}\right)\right],$$
(48)
where we exploited Eqs. (18,31) and $`h[x]`$ is given by Eq. (21).
The existence of lower and upper bounds on the Gaussian EMs at fixed negativities (the latter strictly proven only for extremal states), limits to some extent the inequivalence arising between the two families of entanglement measures, for nonsymmetric two–mode Gaussian states.
## VII Summary and Outlook
In this work we focused on the simplest conceivable states of a bipartite CV system: two–mode Gaussian states. We have shown that, even in this simple instance, the theory of quantum entanglement hides several subtleties and reveals some surprising aspects. In particular, we have studied the relations existing between different computable measures of entanglement, showing how the negativities (including the standard logarithmic negativity) and the Gaussian convex-roof extended measures (Gaussian EMs, including the Gaussian entanglement of formation GEOF ) are inequivalent entanglement quantificators for nonsymmetric two–mode Gaussian states. We have computed Gaussian EMs explicitely for the two classes of two-mode Gaussian states having extremal (maximal and minimal) negativities at fixed purities extremal . We have highlighted how, in a certain range of values of the global and local purities, the ordering on the set of entangled states, as induced by the Gaussian EMs, is inverted with respect to that induced by the negativities. The question whether a certain Gaussian state is more entangled than another, thus, has no definite answer, not even when only extremal states are considered, as the answer comes to depend on the measure of entanglement one chooses. Extended comments on the possible meanings and consequences of the existence of inequivalente orderings of entangled states have been given in Section V and in Section VI. Furthermore, we have proven the existence of a lower bound holding for the Gaussian EMs at fixed negativities, and that this bound is saturated by two–mode symmetric Gaussian states. Finally, we have provided some strong numerical evidence, and partial analytical proofs restricted to extremal states, that an upper bound on the Gaussian EMs at fixed negativities exists as well, and is saturated by states of maximal negativity for given marginals, in the limit of infinite average local mixedness.
We believe that our results will raise renewed interest in the problem of the quantification of entanglement in CV systems, which seemed fairly well understood in the special instance of two–mode Gaussian states. Moreover, we hope that the present work may constitute a first step toward the solution of more general problems concerning the entanglement of Gaussian states, such as the computation of the entanglement of formation for generic two–mode Gaussian states openprob , and the proof of its identity with the Gaussian EoF in a larger class of Gaussian states beyond the symmetric instance. On the other hand, the explicit expressions, computed in the present work, now available for the Gaussian EoF of GMEMS and GLEMS, might serve as well as a basis to find an explicit counterexample to the conjecture that the decomposition over all pure Gaussian states, in the definition of the EoF, is the optimal one for all two–mode Gaussian states.
Finally, the results collected in the present work might prove useful as well in the task of quantifying multipartite entanglement of Gaussian states. For instance, we should mention here that any two–mode reduction of a pure three–mode Gaussian state is a GLEMS, as a consequence of the Schmidt decomposition operated at the CM level holewer . Therefore, thanks to the results that we have derived here, its Gaussian EoF can be explicitely computed, and can be compared with the entropy of entanglement between one reference mode and the remaining two in the global state. One has then available the tools and can apply them to investigate the sharing structure of multipartite CV entanglement of three-mode, and, more generally, multimode Gaussian states contangle .
###### Acknowledgements.
We are grateful to O. Krüger and M. M. Wolf for fruitful discussions and for providing us with supplementary material on the computation of the Gaussian entanglement of formation. Financial support from CNR-Coherentia, INFN, and MIUR is acknowledged.
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# The Einstein–scalar field constraints on asymptotically Euclidean manifolds
## 1 Introduction.
To explain recent observations of far away stars and galaxies, as well as the possible origin of matter elements, it has become more and more relevant in Einsteinian cosmology to admit the existence of a scalar field with a potential which remains to be estimated. On the other hand various considerations, in particular the search for the unification of all the fundamental fields, including gravitation, leads to the belief that the universe has extra dimensions, beyond the usual three space and one time. These extra dimensions would be spacelike, and their extent so small that we don’t perceive them at the usual scales of our experiments.
The relevant equations for cosmology would then be the Einstein equations on an $`n+1`$ dimensional manifold $`V`$, with source a scalar field $`\psi `$ of potential $`V(\psi )`$. These equations are, for a metric $`g`$ on $`V`$ of Lorentzian signature<sup>1</sup><sup>1</sup>1We choose the signature to be $`++\mathrm{}.+.`$,
$$\text{Ei}nstein(g)Ricci(g)\frac{1}{2}R(g)=T;$$
(1.1)
that is, in a local frame
$$S_{\alpha \beta }R_{\alpha \beta }\frac{1}{2}g_{\alpha \beta }R=T_{\alpha \beta }$$
(1.2)
where $`T`$ is the stress energy tensor of a scalar field $`\psi `$ with potential $`V(\psi ),`$ i.e.,
$$T_{\alpha \beta }_\alpha \psi _\beta \psi \frac{1}{2}g_{\alpha \beta }_\lambda \psi ^\lambda \psi g_{\alpha \beta }V(\psi ).$$
(1.3)
The Einstein tensor satisfies the contracted Bianchi identities
$$_\alpha S^{\alpha \beta }0.$$
(1.4)
The field $`\psi `$ is supposed to satisfy the semi linear wave equation
$$^\alpha _\alpha \psi V^{}(\psi )=0,\text{ }V^{}(\psi ):=\frac{dV(\psi )}{d\psi }.$$
(1.5)
The tensor $`T`$ is then divergence free
$$_\alpha T^{\alpha \beta }=0.$$
(1.6)
As a consequence of condition 1.6, equations 1.2 are compatible.
The Cauchy problem for the Einstein equations, determination of an Einsteinian spacetime from initial data on a spacelike $`n`$ dimensional manifold, is a geometric analysis problem. Its solution does not exist for arbitrary initial data, and is not unique from the point of view of analysis due to the invariance of the equations under diffeomorphisms. The geometric initial data are a triple $`(M,\overline{g},K)`$ with $`M`$ an $`n`$ dimensional maifold, which we suppose to be smooth, $`\overline{g}`$ a Riemannian metric on $`M`$, and $`K`$ a symmetric 2 - tensor on $`M`$. The Cauchy data for the scalar field are two functions $`\overline{\psi }`$ and $`\overline{\pi }`$. An $`n+1`$ dimensional spacetime $`(V,g)`$ together with a scalar function $`\psi `$ on $`V`$ is called an Einstein scalar development of these initial data if $`M`$ can be embedded in $`V,`$ so that $`g`$ induces on $`M`$ the metric $`\overline{g}`$ and $`K`$ can be identified with the extrinsic curvature of $`M`$ as submanifold of $`(V,g),`$ while $`\overline{\psi }`$ is the value of $`\psi `$ on $`M,`$ and $`\overline{\pi }`$ is the value on $`M`$ of the derivative of $`\psi `$ in the direction of the unit normal to $`M`$ in $`(V,g).`$
In sections 1 to 5 of this article we use the conformal method to obtain an elliptic system for the constraints satisfied by the initial data of an Einstein - scalar field system. In the following sections we prove some existence and uniqueness theorems for their solution in the case where $`(M,\overline{g})`$ is asymptotically euclidean, under low regularity hypothesis.
The cases of compact $`M`$ and of $`(M,\overline{g})`$ asymptotically hyperbolic will be treated elsewhere.
## 2 Constraints for the Einstein - scalar field equations.
The constraint equations are a consequence of the Gauss Codazzi identities satisfied by the Ricci tensor of any pseudo riemannian manifold.
It is convenient to suppose that $`V=M\times R`$ and to choose on $`V`$ a moving frame $`\theta ^\alpha ,`$ $`\alpha =0,1,\mathrm{}n,`$called a Cauchy adapted frame as long as $`\theta ^0`$ annihilates vectors tangent to submanifolds $`M\times \{t\}.`$ The space time metric is then decomposed as follows
$$gN^2(\theta ^0)^2+g_{ij}\theta ^i\theta ^j\text{ with }\theta ^0dt,\text{ }\theta ^idx^i+\beta ^idt,\text{ }i=1,\mathrm{},n.$$
(2.1)
The function $`N`$ is called the lapse and the time dependent spatial vector $`\beta `$ the shift of the chosen representation of the spacetime metric. In this frame the unit normal $`n`$ to a submanifold $`M\times \{t\}`$ has components
$$n^0=N^1,\text{ }n_0=N,\text{ }n^i=n_i=0.$$
(2.2)
The derivative of the function $`\psi `$ in the direction of $`n`$is
$$\pi N^1_0\psi ,$$
(2.3)
with $`_0`$ the Pfaff derivative with respect to the 1-form $`\theta ^0`$ in the frame $`\theta ^\alpha ,`$ i.e.
$$_0\frac{}{t}\beta ^i_i,\text{ }_i\frac{}{x^i}.$$
(2.4)
In a Cauchy adapted frame the constraints read as the following equations, where we overbar values induced on $`M`$ by spacetime quantities, and we set
$$\tau tr_{\overline{g}}K:=\overline{g}^{ij}K_{ij},\text{ }|K|_{\overline{g}}^2:=\overline{g}^{ih}\overline{g}^{jk}K_{ij}K_{hk}.$$
(2.5)
* Hamiltonian constraint.
$$R(\overline{g})|K|_{\overline{g}}^2+\tau ^2=2\rho 2\overline{N}^2\overline{T}_{00}.$$
(2.6)
* Momentum constraint:
$$\overline{}_iK^{ij}\overline{g}^{ij}_i\tau =J^j\overline{N}^1\overline{T}_0^j.$$
(2.7)
In the case under study, where a source is a scalar field $`\psi `$ we find that
$$2\overline{N}^2\overline{T}_{00}=\overline{\pi }^2+|D\overline{\psi }|_{\overline{g}}^2+2V(\overline{\psi })$$
and
$$\overline{N}^1\overline{T}_0^i=\overline{\pi }\overline{g}^{ij}_j\overline{\psi }.$$
## 3 Conformal formulation.
### 3.1 Hamiltonian constraint.
In order to turn the Hamiltonian constraint into a semilinear elliptic equation to be solved for a scalar function, one considers the metric $`\overline{g}`$ as determined only up to a conformal factor. One sets for $`n>2`$
$$\overline{g}=\phi ^{\frac{4}{n2}}\gamma ,\text{ i.e. }\overline{g}_{ij}=\phi ^{\frac{4}{n2}}\gamma _{ij}$$
(3.1)
with $`\gamma `$ a given Riemannian metric on $`M`$. This particular conformal weight turns into a linear operator the differential operator on $`\phi `$ appearing in the parenthesis of 3.2 below.
The scalar curvatures $`R(\overline{g})`$ and $`R(\gamma )`$ of the conformal metrics $`\overline{g}`$ and $`\gamma `$ are linked by the formula, where $`\mathrm{\Delta }_\gamma `$ is the Laplace operator in the metric $`\gamma ,`$
$$R(\overline{g})\phi ^{\frac{n+2}{n2}}(\phi R(\gamma )\frac{4(n1)}{n2}\mathrm{\Delta }_\gamma \phi ).$$
(3.2)
The Hamiltonian constraint becomes, when $`\gamma `$ and $`K`$ are known, a semi linear elliptic equation for $`\phi `$ with a non linearity of a fairly simple type:
$$\mathrm{\Delta }_\gamma \phi k_nR(\gamma )\phi +k_n(|K|_{\overline{g}}^2\tau ^2+2\rho )\phi ^{\frac{n+2}{n2}}=0$$
(3.3)
with
$$k_n=\frac{n2}{4(n1)}.$$
(3.4)
### 3.2 Momentum constraint.
We can express the momentum constraint in terms of $`\gamma ,`$ $`K,`$ $`\rho ,J`$ and $`\phi `$ by using the relations between the connections of two conformally related metrics.
###### Lemma 3.1
On an $`n`$ dimensional manifold, if $`\overline{g}=\phi ^{\frac{4}{n2}}\gamma ,`$ and if the covariant derivatives in $`\overline{g}`$ and $`\gamma `$ are written respectively as $`\overline{}`$ and $`D,`$ then the divergences in the metric $`\overline{g}`$ and $`\gamma `$ of an arbitrary contravariant 2- tensor $`P^{ij}`$ are linked by the identity
$$\overline{}_iP^{ij}\phi ^{\frac{2(n+2)}{n2}}D_i\{\phi ^{\frac{2(n+2)}{n2}}P^{ij}\}\frac{2}{n2}\phi ^1\gamma ^{ij}_i\phi tr_\gamma P.$$
(3.5)
Proof. The proof follows from a simple computation using the identity which links the coefficients of the connections $`\overline{\mathrm{\Gamma }}`$ of $`\overline{g}`$ and $`C`$ of $`\gamma `$:
$$\overline{\mathrm{\Gamma }}_{jh}^i=C_{jh}^i+\frac{2}{n2}\phi ^1\{\delta _j^i_h\phi +\delta _h^i_j\phi \gamma ^{ik}\gamma _{jh}_k\phi \}.$$
(3.6)
One sees from the identity 3.5 that it is convenient to split the unknown $`K`$ into a weighted traceless part and its trace, namely we set
$$K^{ij}=\phi ^{\frac{2(n+2)}{n2}}\stackrel{~}{K}^{ij}+\frac{1}{n}\overline{g}^{ij}\tau .$$
(3.7)
Here $`\stackrel{~}{K}^{ij}`$ is a symmetric traceless two tensor, in the sense that
$$tr\stackrel{~}{K}\overline{g}_{ij}\stackrel{~}{K}^{ij}=\gamma _{ij}\stackrel{~}{K}^{ij}=0,$$
(3.8)
while $`\tau `$ is the trace.
The momentum constraint 2.7 then becomes
$$D_i\stackrel{~}{K}^{ij}=\frac{n1}{n}\phi ^{\frac{2n}{n2}}\gamma ^{ij}_i\tau +\phi ^{\frac{2(n+2)}{n2}}J^j.$$
(3.9)
It follows from an elementary computation that
$$|K|_{\overline{g}}^2\overline{g}_{ih}\overline{g}_{jk}K^{ij}K^{hk}=\phi ^{\frac{3n+2}{n2}}|\stackrel{~}{K}|_\gamma ^2+\frac{1}{n}\tau ,\text{ with }|\stackrel{~}{K}|_\gamma ^2\gamma _{ih}\gamma _{jk}\stackrel{~}{K}^{ij}\stackrel{~}{K}^{hk}.$$
(3.10)
The Hamiltonian constraint therefore reads
$`\mathrm{\Delta }_\gamma \phi k_nR(\gamma )\phi +k_n\phi ^{\frac{3n+2}{n2}}|\stackrel{~}{K}|_\gamma ^2{\displaystyle \frac{n2}{4n}}\phi ^{\frac{n+2}{n2}}\tau ^2`$ (3.11)
$`={\displaystyle \frac{n2}{2(n1)}}\rho \phi ^{\frac{n+2}{n2}}`$
If $`\gamma ,\stackrel{~}{K},\tau `$ and $`\rho `$ are specified, this is a semilinear elliptic equation for $`\phi `$ when $`\stackrel{~}{K}`$ is known, called a Lichnerowicz equation.<sup>2</sup><sup>2</sup>2This equation was derived by Lichnerowicz 1944 for n=3, \[Li\]. In 1972 York \[Yo72\] introduced the scaling of the sources, and in 1987 Choquet-Bruhat extended the analysis to general n. In view of this history we refer to 3.11 as the Lichnerowicz equation..
## 4 Scaling of $`\overline{\pi }`$.
We denote by an overbar the values induced on $`M`$ by spacetime quantities. The initial data of the scalar field $`\psi `$ is the value $`\overline{\psi }`$ induced by $`\psi `$ on $`M.`$ It is independent on the choice of the conformal metric $`\gamma ,`$ but there is an ambiguity for the data of the initial data for $`\pi ,`$ because $`\pi `$ depends on the lapse $`N:`$ it holds that $`\overline{\pi }=\overline{N}^1\overline{_0\psi }.`$ We associate to the unphysical metric $`\gamma `$ an unphysical lapse $`\stackrel{~}{N},`$ such that $`\overline{N}`$ and $`\stackrel{~}{N}`$ have the same associated densities respectively for $`\overline{g}`$ and $`\gamma ,`$ that is:
$$\overline{N}(Det\overline{g})^{\frac{1}{2}}=\stackrel{~}{N}(Det\gamma )^{\frac{1}{2}}$$
(4.1)
i.e.
$$\overline{N}=\phi ^{2n/(n2)}\stackrel{~}{N}.$$
(4.2)
and we suppose that the given initial data is
$$\stackrel{~}{\pi }=\stackrel{~}{N}^1\overline{_0\psi }=\phi ^{2n/(n2)}\overline{N}^1\overline{_0\psi }=\phi ^{2n/(n2)}\overline{\pi }.$$
### 4.1 Hamiltonian constraint.
The energy density on $`M`$ of a scalar field $`\psi `$ with potential $`V(\psi ),`$ for an observer at rest in the physical metric $`\overline{g}`$ reads as follows in terms of the given data:
$$\rho =\frac{1}{2}(\phi ^{\frac{4n}{n2}}|\stackrel{~}{\pi }|^2+\phi ^{\frac{4}{n2}}\gamma ^{ij}_i\overline{\psi }_j\overline{\psi })+V(\overline{\psi }).$$
(4.3)
We see that the term $`|\stackrel{~}{\pi }|^2`$ adds in the Hamiltonian constraint to $`|`$$`\stackrel{~}{K}|_\gamma ^2`$, while the term $`V(\overline{\psi })`$ remains unscaled by a power of $`\phi `$. The $`\psi `$ term adds a positive contribution to the $`\phi `$ term, adding to $`R(\gamma )`$. The Hamiltonian constraint now reads
$$\mathrm{\Delta }_\gamma \phi f(\phi )=0,\text{ }$$
(4.4)
with
$$f(\phi )r\phi a\phi ^{\frac{3n2}{n2}}+b\phi ^{\frac{n+2}{n2}},$$
where we have again set $`k_n=\frac{n2}{4(n1)}`$ and where
$$rk_n[R(\gamma )|D\overline{\psi }|_\gamma ^2],\text{ }ak_n(|\stackrel{~}{K}|_\gamma ^2+|\stackrel{~}{\pi }|^2),\text{ }b\frac{n2}{4n}\tau ^2\frac{n2}{(n1)}V(\overline{\psi }).$$
(4.5)
We observe that $`a0,`$ while $`b0`$ if $`V(\overline{\psi })0`$ and $`\tau =0`$ (maximal slicing).
We call the equation 4.4 the conformally formulated Hamiltonian constraint, or the Lichnerowicz equation for the Einstein - scalar field theory..
### 4.2 Momentum constraint
The expression of the scalar field momentum density in terms of the new data is:
$$J^i=\overline{g}^{ij}(_j\overline{\psi })\overline{\pi }=\phi ^{\frac{2(n+2)}{n2}}\gamma ^{ij}(_j\overline{\psi })\stackrel{~}{\pi }.$$
(4.6)
The momentum constraint now reads
$$^jD_i\stackrel{~}{K}^{ij}F^j=0$$
(4.7)
with
$$F^j\frac{n1}{n}\phi ^{2n/(n2)}\gamma ^{ij}_i\tau \gamma ^{ij}_j\overline{\psi }\stackrel{~}{\pi }.$$
We call this equation the conformally formulated momentum constraint.
We have proved the following theorem.
###### Theorem 4.1
The conformally formulated momentum constraint of the Einstein - scalar field system 4.7 is a linear system for $`\stackrel{~}{K}`$ when $`\gamma ,`$ $`\tau ,\overline{\psi }`$ and $`\stackrel{~}{\pi }`$ are given and the function $`\phi `$ is known. It does not contain $`\phi `$ if $`\tau `$ is a constant.
### 4.3 Conformal covariance of the constaint equations.
It follows from the analysis above that that if ($`\phi ,\stackrel{~}{K}`$) satisfies the conformally formulated constraints (3.9, 4.6, 3.11), for a specified choice of the free data ($`\gamma ,\tau ,\stackrel{~}{\psi },\stackrel{~}{\pi }),`$ then
$$\overline{g}_{ij}=\phi ^{\frac{4}{n2}}\gamma _{ij},\text{ }K^{ij}=\phi ^{2(n+2)/(n2)}\stackrel{~}{K}^{ij}+\frac{1}{n}\overline{g}^{ij}\tau ,\text{ }\overline{\psi },\text{ }\overline{\pi }=\phi ^{\frac{2n}{n2}}\stackrel{~}{\pi }.$$
(4.8)
is a solution of the original Einstein - scalar field constraints.
The following conformal covariance result is an immediate corollary:
###### Theorem 4.2
Let ($`\phi ,\stackrel{~}{K})`$ be a solution of the conformally formulated constraints in the metric $`\gamma `$ with data $`\tau ,`$ $`\overline{\psi }`$ and $`\stackrel{~}{\pi }.`$ Then ($`\phi ^{}=\theta ^1\phi ,`$ $`\stackrel{~}{K}^{}\theta ^{2(n+2)/(n2)}\stackrel{~}{K})`$ is a solution of the conformally formulated constraints in the metric $`\gamma ^{}=\theta ^{\frac{4}{n2}}\gamma `$ with data $`\tau ^{}=\tau ,`$ $`\overline{\psi }=\overline{\psi }^{},`$ $`\stackrel{~}{\pi }^{}=\theta ^{\frac{2n}{n2}}\stackrel{~}{\pi }.`$
## 5 Solution of the conformal momentum constraint.
The general solution of a non homogeneous linear system is obtained by adding a particular solution to the general solution of the associated linear homogeneous system, which, in the case of 3.9, is the following:
$$D_j\stackrel{~}{K}^{ij}=0,\text{ }\gamma _{ij}\stackrel{~}{K}^{ij}=0.$$
(5.1)
Symmetric 2- tensors satisfying 5.1 are called TT tensors (transverse, traceless). As a consequence of lemma 3.1 the space of TT tensors is the same for two conformal metrics.
We may obtain both the particular solution to 3.9 and the general solution to 5.1 by essentially the same ansatz. One can look for the particular solution of 4.7 as the conformal Lie derivative of a vector field $`Z,`$ an element of the formal $`L^2`$ dual of the space of TT tensors defined by
$$(_{\gamma ,conf}Z)_{ij}:=D_iZ_j+D_jZ_i\frac{2}{n}\gamma _{ij}D_hZ^h.$$
(5.2)
We look for $`\stackrel{~}{K}_{TT}`$ as the sum of the conformal Lie derivative of a vector $`Y`$ and an arbitary traceless symmetric 2 tensor $`U.`$ Then, setting $`X:=Z+Y`$ it holds that:
$$\stackrel{~}{K}^{ij}=(_{\gamma ,conf}X)^{ij}U^{ij},$$
(5.3)
with $`X`$ a vector field solution of the linear system
$$(\mathrm{\Delta }_{\gamma ,conf}X)^j:=D_i(_{\gamma ,conf}Z)^{ij}=D_iU^{ij}+\frac{n1}{n}\phi ^{2n/(n2)}\gamma ^{ij}_i\tau \gamma ^{ij}_j\overline{\psi }\stackrel{~}{\pi }.$$
(5.4)
The arbitrary data in the traceless tensor $`\stackrel{~}{K}`$ is the symmetric traceless tensor $`U.`$
It has been noted by York \[Yo99\] that, though the formulation 4.4, 4.7 is invariant in the sense of Theorem 4.2, the splitting of the solution $`\stackrel{~}{K}`$ into a given traceless tensor $`U`$ and the conformal Lie derivative of an unknown vector $`X`$ cannot be made conformally invariant. To try to obtain $`\stackrel{~}{K}^{}\theta ^{2(n+2)/(n2)}\stackrel{~}{K}`$ given $`\gamma ^{}`$=$`\theta ^{\frac{4}{n2}}\gamma ,`$ we can impose the relation between the given traceless tensors $`U`$ and $`U^{}:`$
$$U^{ij}\theta ^{2(n+2)/(n2)}U^{ij};$$
(5.5)
however for an arbitrary vector $`X`$ one has
$$(_{\gamma ^{},conf}X)^{ij}\theta ^{4/(n2)}(_{\gamma ,conf}X)^{ij}.\text{ }$$
There is no scaling of $`X`$ by a power of $`\phi `$ that leads to a vector $`X^{}`$ and results in the desired scaling of its conformal Lie derivative. York has proposed to remedy this defect by what he called ”the conformal thin sandwich formulation” . Inspired by his work, and by the expression $`K^{ij}N^1\overline{}_0g^{ij},`$ we replace the search for a particular solution as a conformal Lie derivative by the following. For $`\stackrel{~}{N}`$ is a given scalar we define:
$$\stackrel{~}{}_{\gamma ,conf}X:=\stackrel{~}{N}^1_{\gamma ,conf}X,$$
(5.6)
$$(\stackrel{~}{\mathrm{\Delta }}_{\gamma ,conf}X)^j:=D_i(\stackrel{~}{}_{\gamma ,conf}X)^{ij}.$$
(5.7)
The mathematical properties of $`\mathrm{\Delta }_{\gamma ,conf}`$ and $`\stackrel{~}{\mathrm{\Delta }}_{\gamma ,conf}`$ are essentially the same.
We choose $`X`$ to be a solution of the equation
$$(\stackrel{~}{\mathrm{\Delta }}_{\gamma ,conf}X)^j=D_iU^{ij}+\frac{n1}{n}\phi ^{2n/(n2)}\gamma ^{ij}_i\tau \gamma ^{ij}_j\overline{\psi }\stackrel{~}{\pi }.$$
(5.8)
(instead of 5.4). The tensor $`\stackrel{~}{K}`$ solution of 3.9 is now, instead of 5.3,
$$\stackrel{~}{K}^{ij}(\stackrel{~}{}_{\gamma ,conf}X)^{ij}U^{ij}.$$
(5.9)
Noting that if we conformally change the metric via $`\gamma ^{}=\theta ^{\frac{4}{n2}}\gamma `$ and the lapse via $`\stackrel{~}{N}=\theta ^{\frac{2n}{n2}}\stackrel{~}{N}^{}`$ then
$$(\stackrel{~}{}_{\gamma ^{},conf}X)^{ij}=\theta ^{2(n+2)/(n2)}(\stackrel{~}{}_{\gamma ,conf}X)^{ij},$$
(5.10)
we find that $`K`$ has the required scaling.
We are thus led to the following corollary to the theorem 4.2, under otherwise the same hypothesis.
###### Corollary 5.1
If the tensor $`\stackrel{~}{K},`$ a solution of the momentum constraint conformally formulated in a metric $`\gamma ,`$ is obtained as the sum of a given traceless tensor $`U`$ and the product by a given function $`\stackrel{~}{N}`$ of a conformal Lie derivative of a vector $`X:`$
$$\stackrel{~}{K}^{ij}(\stackrel{~}{}_{\gamma ,conf}X)^{ij}U^{ij},$$
(5.11)
then the tensor
$$\stackrel{~}{K}^{ij}(\stackrel{~}{}_{\gamma ^{},conf}X)^{ij}U^{ij},\text{ }U^{ij}=\theta ^{2(n+2)/(n2)}U^{ij},\text{ }\stackrel{~}{N}^{}=\theta ^{\frac{2n}{n2}}\stackrel{~}{N}$$
(5.12)
is a solution of the momentum constraint conformally formulated in the metric $`\gamma ^{}.`$
## 6 Asymptotically Euclidean Manifolds.
### 6.1 Definitions.
In the following sections we will study the solution of the conformally formulated constraints 4.4 and 4.7 on asymptotically euclidean manifolds of dimension $`n3.`$
The Euclidean space $`E^n`$ is the manifold $`R^n`$ endowed with the Euclidean metric, which is $`(dx^i)^2`$ in canonical coordinates.
A $`C^{\mathrm{}}`$, $`n`$-dimensional, Riemannian manifolds $`(M,e)`$ is called ”Euclidean at infinity” if there exists a compact subset $`S`$ of $`M`$ such that $`MS`$ is the disjoint union of a finite number of open sets $`U_i`$, with each $`(U_i,e)`$ being isometric to the exterior of a ball in $`R^n`$. Each open set $`U_iM`$ is sometimes called an ”end” of $`M`$. If $`M`$ is diffeomorphic to $`R^n`$, it has only one end; and we can then take for $`e`$ the Euclidean metric. Unless otherwise specified our manifolds are without boundary; hence the manifold $`(M,e)`$ is complete<sup>3</sup><sup>3</sup>3For studies on asymptotically Euclidean manifolds with boundary see Chrusciel and Delay \[Chru-De\], \[Ma03\], \[Ma04a\] The articles \[Da\], \[Da-Fr\] also consider such manifolds, using the Friedrich’s conformal compactification..
A Riemannian manifold $`(M,\gamma )`$ is called asymptotically Euclidean if there exists a Riemannian manifold $`(M,e),`$ Euclidean at infinity, and if $`\gamma `$ tends to $`e`$ at infinity in each end. Consider one end $`U`$ and the canonical coordinates $`x^i`$ in the space $`R^n`$ which contains the exterior of the ball to which $`U`$ is diffeomorphic. Set $`r\{(x^i)^2\}^{1/2}`$. In the coordinates $`x^i`$ the metric $`e`$ has components $`e_{ij}=\delta _{ij}`$. The metric $`\gamma `$ tends to $`e`$ at infinity if in these coordinates $`\gamma _{ij}\delta _{ij}`$ tends to zero. A possible way of making this statement mathematically precise is to use the Nirenberg - Walker weighted Sobolev spaces. One can also use in these elliptic constraint problems weighted Hölder spaces<sup>4</sup><sup>4</sup>4See \[CB-CS\]., but they are not well adapted to the related evolution questions.
A weighted Sobolev space $`W_{s,\delta }^p,`$with $`1p\mathrm{},`$ with $`s`$ a positive or zero integer, $`\delta `$ a real number, for tensors of some given type on the manifold $`(M,e)`$ euclidean at infinity is the space of tensors of that type which admit generalized $`e`$ \- covariant derivatives of order up to $`s`$ and for which the following norm is finite:
$$u_{W_{s,\delta }^p}=\left\{\underset{0ms}{}_V^mu^p(1+d^2)^{\frac{1}{2}p(\delta +m)}𝑑\mu \right\}^{1/p}.$$
(6.1)
Here $``$, $`||`$ and $`d\mu `$ denote the covariant derivative, norm and volume element corresponding to the metric $`e`$, and $`d`$ is the distance in the metric $`e`$ from a point of $`M`$ to a fixed point. If $`(M,e)`$ is a euclidean space one can choose $`d=r`$, the euclidean distance to the origin. The space $`𝒟`$ of $`C^{\mathrm{}}`$ tensors with compact support is dense in $`W_{s,\delta }^p,`$ regardless of what $`s`$ and $`\delta `$ are, so long as $`p<\mathrm{}.`$
If $`s`$ and $`\delta `$ are large enough, a function (or tensor field) in $`W_{s,\delta }^p`$ is continuous and tends to zero at infinity. Specifically if we define $`C_\beta ^m`$ to be the Banach space of weighted $`C^m`$ functions (or tensor fields) on $`(M,e)`$ with norm given by
$$u_{C_\beta ^m}\underset{0\mathrm{}m}{}\underset{M}{sup}(|^{\mathrm{}}u|(1+d^2)^{\frac{1}{2}(\beta +\mathrm{})}),$$
then the following inequality holds<sup>5</sup><sup>5</sup>5For proofs of this embedding and the multiplication rule 6.3, see \[CB-Ch\] 1981, or \[CB- DM\] II p. 396., with $`C`$ a number depending only on $`(M,e),`$
$$u_{C_\beta ^m}Cu_{W_{s,\delta }^p},\text{ if }s>m+\frac{n}{p},\text{ }\delta >\beta \frac{n}{p}.$$
(6.2)
We see that $`uW_{s,\delta }^p`$ implies that $`u`$ is continuous and tends to zero at infinity if $`s>\frac{n}{p}`$ and $`\delta >\frac{n}{p}.`$
Let $`(M,e)`$ be a manifold which is Euclidean at infinity. The Riemannian manifold $`(M,\gamma )`$ is said to be $`(p,\sigma ,\rho )`$ asymptotically euclidean if $`\gamma eW_{\sigma ,\rho }^p`$. If $`\gamma eW_{\sigma ,\rho }^p`$ with $`\sigma >\frac{n}{p}`$, and $`\rho >\frac{n}{p},`$ then $`\gamma `$ is $`C^0`$ and $`\gamma e`$ tends to zero at infinity. The set of Riemannian metrics (i.e. positive definite symmetric 2-tensors) such that $`\gamma eW_{\sigma ,\rho }^p`$ is denoted $`M_{\sigma ,\rho }^p.`$
We recall the multiplication lemma
$$W_{s_1,\delta _1}^p\times W_{s_{2,\delta _2}}^pW_{s,\delta }^p\text{ if }s<s_1+s_2\frac{n}{p},\text{ }\delta <\delta _1+\delta _2+\frac{1}{p},$$
(6.3)
and the interpolation<sup>6</sup><sup>6</sup>6CB, to appear. inequality: for any $`\epsilon >0,`$ there is a $`C(\epsilon )`$ such that, for all $`uW_{m,\delta }^p,`$ $`1pq\mathrm{},`$ and $`j<m,`$ one has
$$^ju_{_{W_{0,\delta +\frac{n}{p}\frac{n}{q}+j}^q}}C\{\epsilon ||u||_{W_{m,\delta }^p}+C_\epsilon u_{W_{0,\delta }^p}\}.$$
(6.4)
### 6.2 Linear elliptic systems.
We state<sup>7</sup><sup>7</sup>7The theorem relies on previous results of Nirenberg and Walker \[N-W\], Cantor \[Ca\], Choquet-Bruhat and Christodoulou \[CB-Ch\], and the interpolation lemma to lower the regularity required of coefficiemts \[CB\]. the following existence theorem for solutions of linear ellipticPDE’s.
###### Theorem 6.1
Hypotheses:
Let $`(M,e)`$ be a smooth Riemannian manifold Euclidean at infinity. Let
$$Lua_2^2u+a_1u+a_0u$$
(6.5)
be a second order linear elliptic operator acting on tensor fields on $`(M,e),`$ which in terms of components $`u^A,`$ $`A=1,\mathrm{}p,`$ of the tensor $`u`$ takes the form
$$(Lu)^Aa_{2,B}^{ij,A}_{ij}^2u^B+a_{1,B}^{i,A}_iu^B+a_{0,B}^Au^B.$$
Let the principal symbol $`a_2^{ij}\xi _i\xi _j`$ of $`L`$ be an isomorphism from $`R^p`$ onto $`R^p`$ for $`\xi 0.`$ Suppose the coefficients of $`L`$ satisfy the following hypotheses
$$a_2AW_{2,\delta }^p,a_1W_{1,\delta +1}^p,\text{ }a_0W_{0,\delta +2}^p.$$
(6.6)
where $`A^2`$ is an elliptic operator with $`C^{\mathrm{}}`$ coefficients, constant in each end of $`(M,e)`$ and
$$p>\frac{n}{2},\frac{n}{p}<\delta <\frac{n}{p}+n2.$$
(6.7)
Conclusions:
1. The operator $`L`$ is a continuous mapping from $`W_{2,\delta }^p`$ into $`W_{0,\delta +2}^p`$
2. a. There exists a number $`C_L>0,`$ depending only on $`A`$ and on the norms of $`a_2A,`$ $`a_1,a_0`$, and a number $`\delta ^{}>\delta `$ such that the following inequality holds for all $`uW_{2,\delta }^p:`$
$$u_{W_{2,\delta }^p}C_L\{Lu_{W_{0,\delta +2}^p}+u_{W_{1,\delta ^{}}^p}\}.$$
(6.8)
b. If in addition $`L`$ is injective there exists a number $`C`$ such that the following inequality holds for all $`uW_{2,\delta }^p:`$
$$u_{W_{2,\delta }^p}CLu_{W_{0,\delta +2}^p}.$$
(6.9)
c. The operator $`L`$ has finite dimensional kernel and closed range.
3. a. If the adjoint operator $`L^{}`$ is injective on $`W_{2,\delta }^p,`$ then $`L`$ is surjective from $`W_{2,\delta }^p`$ onto $`W_{0,\delta +2}^p.`$
b. If $`L`$ and $`L^{}`$ are both injective then they are isomorphisms from $`W_{2,\delta }^p`$ onto $`W_{0,\delta +2}^p.`$
###### Corollary 6.2
If, in addition to the previous hypothesis (including injectivity) it holds that
$$a_2AW_{s+2,\delta }^p,a_1W_{s+1,\delta +1}^p,\text{ }a_0W_{s,\delta +2}^p.$$
(6.10)
then $`L`$ is an isomorphism from $`W_{s+2,\delta }^p`$ onto $`W_{s,\delta +2}^p.`$
We recall also the following lemma (lemma 5.2 of CB-Ch).
###### Lemma 6.3
Suppose that $`uW_{2,\delta }^p`$ is a solution of the equation
$$Lua_2^2u+a_1u+a_0u=f$$
(6.11)
where $`L`$ satisfies the hypotheses of the above theorem and where $`fW_{0,\stackrel{~}{\delta }+2}^p,`$ $`\delta \stackrel{~}{\delta }.`$ Then $`u`$ is in fact in $`W_{2,\stackrel{~}{\delta }}^p,`$ so long as $`\stackrel{~}{\delta }<n2\frac{n}{p}.`$
### 6.3 The Poisson operator, $`\mathrm{\Delta }_\gamma a`$
###### Theorem 6.4
Let $`(M,\gamma )`$ be a $`M_{2,\delta }^p`$ manifold<sup>8</sup><sup>8</sup>8The hypothesis p$`>`$n/2 is stronger than necessary but simplifies the proof, and is needed later in our treatment of non linear equations., $`p>\frac{n}{2},`$ $`\delta >\frac{n}{p}.`$ Let $`aW_{0,\delta +2}^p`$be given. The Poisson operator $`\mathrm{}_\gamma a`$ is an isomorphism from $`W_{2,\delta }^p`$ onto $`W_{0,\delta +2}^p`$ if
$$_M\{|u|^2+au^2\}\mu _\gamma >0,\text{ }$$
(6.12)
for any $`uW_{2,\stackrel{~}{\delta }}^p,`$with $`\stackrel{~}{\delta }`$ some number such that $`n2\frac{n}{p}>\stackrel{~}{\delta }>1+\frac{n}{2}\frac{n}{p}(\stackrel{~}{\delta }=1`$ if $`p=2),`$ with $`u0.`$
Proof. The operator $`\mathrm{}_\gamma a`$ is self adjoint. It is an isomorphism $`W_{2,\delta }^pW_{0,\delta +2}^p`$ if injective. By lemma 6.3 it is sufficient to prove the injectivity on $`W_{2,\stackrel{~}{\delta }}^p`$ for some $`\stackrel{~}{\delta }`$ such that $`\stackrel{~}{\delta }<`$ $`n2\frac{n}{p}.`$ The theorem is obtained by integration on $`M`$ of $`u(\mathrm{}_\gamma uau),`$ trivially in the case $`p=2,`$ $`\stackrel{~}{\delta }=1`$ (compatible with $`p>\frac{n}{2}`$ if and only if $`n=3)`$; and by using either Sobolev embeddings or the Holder inequality in the case $`p2`$ and $`\stackrel{~}{\delta }>1+\frac{n}{2}\frac{n}{p}.`$
###### Theorem 6.5
Let $`u`$ satisfy the equation
$$\mathrm{\Delta }_\gamma uau=f$$
(6.13)
with $`\gamma M_{2,\delta }^p,\delta >\frac{n}{p},`$ and $`p>\frac{n}{2}.`$ Suppose $`aW_{0,\delta +2}^p,`$ $`ucW_{2,\stackrel{~}{\delta }}^p`$, where $`c`$ is a given number, $`1+\frac{n}{2}\frac{n}{p}<\stackrel{~}{\delta }`$ , ($`\stackrel{~}{\delta }1`$ if $`p=2).`$ Suppose$`a0.`$ Then $`u0`$ on $`M`$ if $`f0`$ and $`c0.`$ If $`f0`$ and $`c0,`$ then $`u0`$ on $`M.`$ The lower bound of $`\stackrel{~}{\delta }`$ can be weakened to $`\delta >\frac{n}{p}`$ if $`c=0`$ and $`fW_{2,\stackrel{~}{\delta }}^p,`$ $`\stackrel{~}{\delta }>\frac{n}{p}+\frac{n}{2}1,`$ ($`\stackrel{~}{\delta }=1`$ if $`p=2).`$
Proof. The integration on $`M`$ of $`v(\mathrm{}_\gamma uau),`$ and the choice $`v=u^+=Sup(u,0),`$ gives $`u^+=constant,`$ therefore $`u^+0`$ since $`u^+`$ tends to zero at infinity.
## 7 Solution of the momentum constraint.
Given the riemannian metric $`\gamma `$ and the scalar field $`\stackrel{~}{N}`$ the conformally formulated momentum constraint reads
$$D_j(\stackrel{~}{}X)^{ij}(\stackrel{~}{\mathrm{\Delta }}_{\gamma ,conf}X)^i=F^i(\phi )$$
(7.1)
with
$$F^i(\phi )D_jU^{ij}+\frac{n1}{n}\phi ^{2n/(n2)}\gamma ^{ij}_j\tau +\gamma ^{ij}(_j\overline{\psi })\stackrel{~}{\pi }.$$
(7.2)
where $`\tau `$ is a given function on $`M`$ and $`U`$ is a given symmetric traceless 2 - tensor field. The sources $`\overline{\psi }`$ and $`\stackrel{~}{\pi }`$ are given. We suppose momentarily that $`\phi `$ is also a known function. In fact it disappears from the equation if $`\tau 0`$.
###### Lemma 7.1
Let $`(M,\gamma )`$ be a $`W_{2,\delta }^p`$ asymptotically euclidean manifold, and let $`\stackrel{~}{N}=1+\nu ,`$ $`\nu W_{2,\delta }^p,`$ $`\stackrel{~}{N}>0,`$ be given. Suppose that $`p>\frac{n}{2},`$ $`\delta >\frac{n}{p};`$ then
1. The operator $`\stackrel{~}{\mathrm{\Delta }}_{\gamma ,conf}`$ is elliptic.
2. Its kernel in W$`{}_{}{}^{p}{}_{2,\delta }{}^{}`$ is the space of W$`{}_{}{}^{p}{}_{2,\delta }{}^{}`$ conformal Killing vector fields of the metric $`\gamma .`$
Proof. It holds that
$$(\stackrel{~}{\mathrm{\Delta }}_{\gamma ,conf}Y)^j\stackrel{~}{N}^1D_i[D^iY^j+D^jY^i\frac{2}{n}\gamma ^{ij}D_kY^k]\stackrel{~}{N}^2(_{\gamma ,conf}Y)^{ij}D_i\stackrel{~}{N}.$$
(7.3)
Using the Ricci identity we find that the principal part is
$$\stackrel{~}{N}^1[(\mathrm{\Delta }_\gamma Y)^j+(1\frac{2}{n})D^jD_iY^i].$$
(7.4)
The principal symbol is easily checked to be an isomorphism of $`R^n`$ , for any $`n2.`$
2. We prove the second part of this lemma using integration by parts, using lemma 6.3.
We can now prove the following theorem.
###### Theorem 7.2
Let $`(M,\gamma )`$ be a $`M_{2,\delta }^p`$ asymptotically euclidean manifold, with $`p>\frac{n}{2},`$ and $`\delta >\frac{n}{p}`$ Let $`\overline{\psi }W_{2,\delta }^p`$ and $`U`$, $`\tau ,\stackrel{~}{\pi }W_{1,\delta +1}^p`$ be given. Suppose also that $`\phi `$ is known, with $`\phi >0,`$ and $`(1\phi )W_{2,\delta }^p`$. Then the momentum constraint 4.7 has one and only one solution $`XW_{2,\delta }^p`$ if, in addition, $`\delta <n2\frac{n}{p}`$.
If $`\tau 0,`$ the condition on $`\phi `$ is irrelevant.
###### Corollary 7.3
If in addition $`\gamma M_{2+s,\delta }^p,`$ $`\overline{\psi }W_{2+s,\delta }^p`$ and $`U,\tau ,\stackrel{~}{\pi }W_{1+s,\delta +1}^p`$ and $`(1\phi )W_{s+2,\delta }^p`$ then the solution $`X`$ belongs to $`W_{s+2,\delta }^p.`$
Proof. The given hypothesis and the Sobolev embedding and multiplication properties imply that the coefficients of the operator $`\stackrel{~}{\mathrm{\Delta }}_{\gamma ,conf}`$ satisfy the hypotheses of the theorem 8.2 and that $`F(\phi )W_{0,\delta +2}^p.`$ The operator $`\stackrel{~}{\mathrm{\Delta }}_{\gamma ,conf\text{ }}`$ is self adjoint, and its kernel in $`W_{2,\delta }^p`$ is empty, because there are no such conformal Killing fields on $`(M,\gamma ).`$
## 8 Solution of the Lichnerowicz equation.
We consider 4.4 (the Lichnerowicz equation)
$$(x,X,\phi )\mathrm{\Delta }_\gamma \phi f(x,\phi )=0,\text{ }$$
(8.1)
where
$$f(x,\phi )r\phi a\phi ^{\frac{3n2}{n2}}+b\phi ^{\frac{n+2}{n2}},$$
and
$$rk_n(R(\gamma )|D\overline{\psi }|_\gamma ^2),\text{ }ak_n(|\stackrel{~}{K}|_\gamma ^2+|\stackrel{~}{\pi }|_\gamma ^2)0,\text{ }b\frac{n2}{4n}\tau ^2\frac{n2}{n1}V(\overline{\psi }).$$
(8.2)
We first prove the following lemma.
###### Lemma 8.1
If $`(M,\gamma )`$ is a $`M_{2,\delta }^p`$ manifold with $`p>\frac{n}{2},`$ $`\delta >\frac{n}{p},`$ and $`\overline{\psi }W_{2,\delta }^p,`$ then $`r(\gamma ,\overline{\psi })W_{0,\delta +2}^p.`$
Proof. $`R(\gamma )`$ is a sum of terms of the form $`\gamma ^2\gamma ,`$ and $`\gamma \gamma \gamma .`$ with $`\gamma eW_{2,\delta }^p,`$ $`\gamma W_{1,\delta +1}^p,`$ and $`^2\gamma W_{\sigma 2,\delta +2}^p.`$ Under the hypotheses made on $`p`$ and $`\delta ,`$ the Sobolev embedding theorem shows that $`\gamma e`$ is continuous and bounded on $`M;`$ the multiplication theorem completes the proof, also for $`|D\overline{\psi }|_\gamma ^2`$.
### 8.1 General existence theorem.
The following theorem extends to asymptotically Euclidean manifolds a theorem which has been proved for data on compact<sup>9</sup><sup>9</sup>9Results of this sort on compact manifolds were proven by Choquet-Bruhat and Leray \[CB-Le\] 1972, using Leray - Schauder degree techniques. In later work \[Is\], sub and super solution techniques have been used.manifolds. It can be proved by similar methods.
###### Theorem 8.2
Let $`(M,\gamma )`$ be in $`M_{2,\delta }^p,`$ $`\delta >\frac{n}{p},`$ $`p>\frac{n}{2}`$. Suppose that $`a,`$ $`b,`$ $`rW_{0,\delta +2}^p,`$ and $`1+\frac{n}{2}`$ $`\frac{n}{p}<\delta `$ (if $`p=2`$ then $`\delta =1`$ is admissible). Suppose the Lichnerowicz equation 4.4 admits a subsolution $`\phi _{}`$ and a supersolution $`\phi _+,`$ which are continuous and bounded functions with $`\phi _+,`$ $`\phi _{}W_{1,\delta +1}^p,`$ such that
$$\mathrm{}_\gamma \phi _{}f(.,\phi _{}),\mathrm{}_\gamma \phi _+f(.,\phi _+),$$
(8.3)
$$\underset{\mathrm{}}{lim}\phi _{}1,\underset{\mathrm{}}{lim}\phi _+1$$
(8.4)
and for which there exist numbers $`\mathrm{}`$ and $`m,`$ with $`\mathrm{}>0`$ if $`a0,`$ such that on $`M`$
$$\mathrm{}\phi _{}\phi _+m.\text{ }$$
(8.5)
Then the equation admits a solution $`\phi `$ such that:
$$\phi _{}\phi \phi _+,1\phi W_{2,\delta }^p$$
(8.6)
for
$$\delta <n2\frac{n}{p}.$$
(8.7)
If moreover $`\gamma M_{2+s,\delta }^p,`$ and $`a,`$ $`bW_{s,\delta +2}^p,`$ then the solution is such that $`1\phi W_{s+2,\delta }^p`$.
Note that constant sub and super solutions are not natural in the asymptotically Euclidean case. In our application of this theorem to the Lichnerowicz equation, we introduce some intermediate steps to obtain non constant sub and supersolutions.
### 8.2 Uniqueness theorem.
The uniqueness of a solution $`\phi `$ of the Lichnerowicz equation follows from monotonicity if we assume that $`r0,a0,`$ and $`b0.`$ A proof of uniqueness can be given under the same hypothesis on $`a`$ and $`b,`$ but with no restriction on the sign of $`r.`$
###### Theorem 8.3
The Lichnerowicz equation 4.4 on $`(M,\gamma ),`$ with $`\gamma M_{2,\delta }^p,`$ $`p>\frac{n}{2},`$ $`\delta >\frac{n}{p}`$ has at most one positive solution $`\phi ,`$ $`\phi 1W_{2,\delta }^p,`$ if $`a,b,rW_{0,\delta +2}^p,`$ and if $`a0`$, and $`b0`$.
Proof. Suppose it admits two solutions $`\phi _1>0`$ and $`\phi _2>0`$. Using the identity 3.2 we find that, with $`\gamma _i:=\phi _i^{4/(n2)}\gamma ,`$ and $`r_i:=k_n(R(\gamma _i)|D\overline{\psi }|_{\gamma _i}^2),`$ $`i=1,2`$
$$\mathrm{\Delta }_{\gamma _2}(\phi _1\phi _2^1)(\phi _1\phi _2^1)r_2(\phi _1\phi _2^1)^{(n+2)/(n2)}r_1$$
(8.8)
Since $`\phi _1`$ is a solution of 8.1 we have
$`r_1`$ $`\phi _1^{(n+2)/(n2)}\{\mathrm{\Delta }_\gamma \phi _1\phi _1r(\gamma ,\overline{\psi })\}`$
$`=\phi _1^{(n+2)/(n2)}\{a\phi _1^{(3n+2)/(n2)}\phi ^{(n+2)/(n2)}b\}`$
and an analogous equation for $`r_2.`$ Inserting these results in the previous equation gives an equation of the form
$$\mathrm{\Delta }_{\gamma _2}(\phi _1\phi _2^11)\lambda \{(\phi _1\phi _2^11\}=0,$$
(8.9)
with
$`\lambda `$ $`a\phi _1^{(3n+2)/(n2)}\phi _2^{(n+2)/(n2)}{\displaystyle \frac{(\phi _1\phi _2^1)^{4(n1)/(n2)}1}{\phi _1\phi _2^11}}+`$
$`b\phi _1\phi _2^1{\displaystyle \frac{(\phi _1\phi _2^1)^{4/(n2)}1}{\phi _1\phi _2^11}}.`$ (8.10)
So long as $`\phi _1`$ and $`\phi _2`$ are continuous and positive functions on $`M,`$ the fractions with denominator $`\phi _1\phi _2^11`$ are continuous and positive functions as well, since the powers of $`\phi _1\phi _2^1`$ appearing in their numerators are greater than 1, Therefore $`\lambda W_{0,\delta +2}^p`$. Noting that by definition $`a0,`$ it follows that if $`\tau `$ and $`V(\overline{\psi })`$ are such that $`b0,`$ then $`\lambda 0.`$ Hence using $`\phi _1\phi _2^11W_{2,\delta }^p`$, and the injectivity of $`\mathrm{\Delta }_\gamma \lambda `$ on $`W_{2,\delta }^p,`$ we have $`\phi _1\phi _2^110.`$
### 8.3 Generalized Brill-Cantor Theorem.
For compact smooth Riemannian manifolds the solutions of the Lichnerowicz equations have been classified by Isenberg \[Is95\] through the use of the Yamabe theorem. The Yamabe conformal invariant is defined by
$$Inf_{f𝒟,f0}(_M\{|Df|^2+k_nR(\gamma )f^2\}\mu _\gamma /||f||_{L^{2n/(n2)}}^2\text{}$$
The Yamabe theorem, proved for smooth metrics in an increasing number of cases by Trudinger, Aubin and Schoen, says that any compact Riemannian manifold is conformal to a manifold with constant scalar curvature, $`+1,1,`$ or $`0`$ according to the sign of the Yamabe invariant. It is easy to see that this theorem extends to $`W_2^p`$ metrics, $`p>\frac{n}{2},`$ in the negative or zero case. In the positive case only a weaker form (proved by Yamabe himself in the smooth case) is proved to hold, namely that the $`W_2^p`$ manifold is conformal to a manifold with strictly positive scalar curvature. This property is used in \[CB02\] and \[Ma04b\]. Maxwell in particular establishes the classification of solutions of the Lichnerowicz equation using only the sign of the Yamabe invariant and not the full Yamabe theorem. The definition of the Yamabe conformal invariant extends to non compact manifolds \[Ma03\]<sup>10</sup><sup>10</sup>10The definition used by Brill - Cantor in their theorem, carried over in CB-I-Y,
$$_M\left\{|Df|^2+k_nR(\gamma )f^2\right\}\mu _\gamma /f_{L^{2n/(n2)}}^2>0\text{}f𝒟,f0\text{ }$$
was incorrect, because it did not imply this inequality for all $`fW_{2,\delta }^p,`$ since the limit of positive functions is not necessarily positive. but there is no theorem for asymptotically euclidean manifolds analogous to the Yamabe theorem, and the denomination of ”positive Yamabe class” for asymptotically Euclidean manifolds with a positive Yamabe invariant is somewhat misleading, as shown by the following theorem, proved<sup>11</sup><sup>11</sup>11Under more restrictive hypothesis on regularity, and in the case n=3. by Brill and Cantor 1981 \[Br-Ca\] , and generalized in the presence of a scalar field as follows.
###### Theorem 8.4
Let $`(M,\gamma )`$ be a $`(p,2,\delta )`$ asymptotically Euclidean manifold with $`p>\frac{n}{2},`$ $`\delta >\frac{n}{p},`$ and let $`\overline{\psi }W_{2,\delta }^p`$ be a scalar field on $`M.`$ There exists on $`M`$ a $`(p,2,\delta )`$ asymptotically euclidean metric $`\gamma ^{}`$ conformal to $`\gamma `$ such that $`r(\gamma ^{},\overline{\psi })=0`$ if and only if $`(M,\gamma ,\overline{\psi })`$ satisfy the following inequality<sup>12</sup><sup>12</sup>12This condition, already used to prove injectivity, is implied by the positivity of the Yamabe invariant, because $`𝒟`$ is dense in $`W_{2,\delta }^p.`$
$$_M\left\{|Df|^2+r(\gamma ,\overline{\psi })f^2\right\}\mu _\gamma >0$$
(8.11)
for every function $`f`$ on $`M`$ with $`fW_{2,\stackrel{~}{\delta }}^p`$, $`\stackrel{~}{\delta }>\frac{n}{p}+\frac{n}{2}1`$ ($`\stackrel{~}{\delta }1`$ if $`p=2),f0.`$
Proof. ($`M,\gamma )`$ is conformal to $`(M,\gamma ^{})M_{2,\delta }^p`$ with $`r(\gamma ^{},\overline{\psi })=0`$ if and only if there exists a function $`\phi >0,`$ such that $`\gamma ^{}=\phi ^{4/(n2)}\gamma M_{2,\delta }^p`$ and
$$\mathrm{}_\gamma \phi r(\gamma ,\overline{\psi })\phi =0,$$
(8.12)
equivalently, setting $`\phi 1+u,`$ 8.12 reads
$$\mathrm{\Delta }_\gamma ur(\gamma ,\overline{\psi })u=r(\gamma ,\overline{\psi }).$$
(8.13)
1. Suppose that the condition 8.11 is satisfied.The equation 8.13. is linear and elliptic, with $`\mathrm{\Delta }_\gamma r(\gamma ,\overline{\psi })`$ an injective operator on $`W_{2,\delta }^p`$, and satisfies the hypotheses of the theorem 8.2; therefore it admits a solution $`uW_{2,\delta }^pC_\alpha ^0.`$ It remains to prove that $`\phi 1+u`$ is positive, then $`\phi ^{4/(n2)}\gamma M_{2,\delta }^p.`$ One cannot use directly the maximum principle because $`r(\gamma ,\overline{\psi })`$ is not necessarily positive. Inspired by Brill and Cantor (see also \[Ma03a\]) we consider the family of equations
$$\mathrm{}_\gamma \phi kr(\gamma ,\overline{\psi })\phi =0\text{ i.e. }\mathrm{}_\gamma ukr(\gamma ,\overline{\psi })u=kr(\gamma ,\overline{\psi })$$
(8.14)
with $`k[0,1]`$ a number. Each of these equations satisfy the condition 8.11 hence admits a solution $`u_kW_{2,\delta }^pC_\alpha ^0,`$ and the $`C_\alpha ^0`$ norm of $`u_k`$ depends continuously on $`k.`$ The set $`S:=\{u_kC_\alpha ^0,`$ $`u_k>1\}`$ is open in $`C_\alpha ^0`$ and non empty because for $`k=0`$ it holds that $`u_0=0`$ (i.e. $`\phi _0=1).`$ To show that it is closed, suppose that $`u_{k^{}\text{ }}`$ belongs to its boundary $`S;`$ then $`u_k^{}1,`$ $`\phi _k^{}0.`$ Suppose that $`\phi _k^{},`$ solution of the elliptic equation 8.12, vanishes at a point of $`M.`$ Then by the weak Harnack inequality (Trudinger 1973) there is a ball $`B_R`$ of center $`x`$ and a number $`C`$ such that
$$\phi _k^{}_{L^q(B_{2R})}CInf_{B_R}\phi _k^{}=0,$$
(8.15)
hence $`\phi _k^{}=0`$ in $`B_R`$ and also, by continuity, on $`M.`$ This is impossible because $`\phi _k^{}`$ tends to 1 at infinity. Hence $`\phi _k^{}>0,`$ The subset $`S`$ of $`C_\alpha ^0`$ being both open and closed is all of $`C_\alpha ^0`$.
2. Conversely suppose that $`\phi >0`$ exists and solves the equation satisfying the hypothesis of the theorem. Then we will show that for any $`f0`$, $`fW_{2,\delta }^p`$ the inequality 8.11 holds. We set $`\theta =f\phi ^1,`$ then $`\theta W_{2,\delta }^pC_\alpha ^0.`$ We have by elementary calculus:
$$|Df|^2=|D\theta |^2\phi ^2+\phi D\phi .D(\theta ^2)+\theta ^2|D\phi |^2.$$
(8.16)
The following integration by parts holds for the functions under consideration:
$$_M\phi D\phi .D(\theta ^2)\mu _\gamma =_M\theta ^2D(\phi D\phi )\mu _\gamma .$$
(8.17)
Therefore,
$$_M\phi D\phi .D(\theta ^2)\mu _\gamma =_M\theta ^2(\phi \mathrm{}_\gamma \phi +|D\phi |^2)\mu _\gamma $$
(8.18)
and
$$_M|Df|^2\mu _\gamma =_M\{|D\theta |^2\phi ^2\theta ^2\phi \mathrm{}_\gamma \phi \}\mu _\gamma .$$
(8.19)
We have $`D\theta 0`$ since $`\theta C_\alpha ^0`$ tends to zero at infinity and cannot be a constant without being identically zero, which is ruled out by the hypothesis $`f0`$. Hence when $`\phi >0`$ satisfies the equation 8.12 the function $`fW_{2,\delta }^p,`$ $`f0`$ satisfies the inequality
$$_M\left\{|Df|^2+r(\gamma ,\overline{\psi })f^2\right\}\mu _\gamma >0.$$
(8.20)
Remark. The same sort of proof shows that, under the same hypothesis, there exists on $`M`$ a metric $`\gamma ^{}`$ conformal to $`\gamma `$ such that $`r(\gamma ^{},\overline{\psi })0.`$
### 8.4 Existence theorems.
###### Theorem 8.5
Let $`(M,\gamma )`$ be a $`M_{2,\delta }^p`$ manifold with $`p>\frac{n}{2}.`$ Let $`\overline{\psi }`$ be a scalar field on $`M`$ with potential $`V(\overline{\psi }),`$ such that $`\overline{\psi }W_{2,\delta }^p`$ and $`V(\overline{\psi })W_{0,\delta +2}^p.`$ Suppose that 8.11 is satisfied and $`b0.`$ The Lichnerowicz equation
$$\mathrm{\Delta }_\gamma \phi r\phi +a\phi ^{\frac{3n2}{n2}}b\phi ^{\frac{n+2}{n2}}=0,$$
(8.21)
$$a,bW_{0,\delta +2}^p,\text{ }\delta >1+\frac{n}{2}\frac{n}{p},\text{ }\delta 1\text{ if }p=2,$$
(8.22)
has one and only one solution, $`\phi =1+u,`$ $`uW_{2,\delta }^p,`$ if $`n2\frac{n}{p}>\delta >1+\frac{n}{2}\frac{n}{p}`$ ( extended to $`\delta 1`$ if $`p=2).`$ The solution can be obtained by iteration.
###### Corollary 8.6
If moreover $`\gamma M_{2+s,\delta }^p`$ and $`a,bW_{s,\delta +2}^p,`$ then $`uW_{s+2,\delta }^p.`$
Proof. Uniqueness: This follows from the general theorem 8.3. It can also be proved directly using the monotonicity of the non linear term.
Existence. Since it follows from theorem 4.2 that the Lichnerowicz equation is conformally invariant, we may, without loss of generality, conformally transform equation to a metric such that $`r(\gamma ,\overline{\psi })=0:`$
$$\mathrm{\Delta }_\gamma \phi +a\phi ^{\frac{3n2}{n2}}b\phi ^{\frac{n+2}{n2}}=0,$$
(8.23)
1. We first consider equation 8.23 with $`b=0:`$
$$\mathrm{\Delta }_\gamma \phi +a\phi ^{\frac{3n2}{n2}}=0$$
(8.24)
This equation admits a constant subsolution $`\phi _{}=1`$ but no finite constant supersolution. However, it admits a non constant supersolution, namely the function $`\phi _+=1+u_+`$ with $`u_+W_{2,\delta }^p`$ a solution of the linear equation
$$\mathrm{\Delta }_\gamma u_+=a;$$
(8.25)
indeed the maximum principle shows that $`u_+0,`$ hence $`\phi _+1`$ and
$$\mathrm{\Delta }_\gamma \phi _+=aa\phi _+^{\frac{3n2}{n2}}.$$
(8.26)
We can apply the general existence theorem 8.2 to prove the existence of a solution $`\phi _1.`$
2. We next consider the equation with $`a=0:`$
$$\mathrm{\Delta }_\gamma \phi b\phi ^{\frac{n+2}{n2}}=0.$$
(8.27)
This equation admits the subsolution $`\phi _{}=0`$ and the supersolution $`\phi _+=1.`$ It admits therefore a solution $`\phi _2,`$ with $`1\phi _2W_{2,\delta }^p,`$ and $`0\phi _21.`$ We prove that $`\phi _2>0`$ by an argument similar to the one used in the proof of the Brill -Cantor theorem: We consider the family of equations
$$\mathrm{}_\gamma \phi kb\phi ^{\frac{n+2}{n2}}=0\text{ }$$
(8.28)
with $`k[0,1]`$ a number. Each of these equations admits one solution $`\phi _k=1+u_k0,`$ with $`u_kW_{2,\delta }^pC_\alpha ^0,`$ and the $`C_\alpha ^0`$ norm of $`u_k`$ depends continuously on $`k.`$ The proof continues as in the proof of theorem 8.4.
3. Consider the general equation 8.23. By the above results this admits $`\phi _1`$ as a supersolution and $`\phi _2`$ as a subsolution. Therefore the existence of a solution follows again from the general existence theorem 8.2. The proof of the corollary also follows from this result.
The proof of the corollary follows from that of theorem 8.2.
We now state two theorems which suppose $`b0.`$ They can be applied in particular when the scalar field has a non negative potential $`V(\overline{\psi })`$ and the initial manifold is maximal or has an appropriately small mean extrinsic curvature.
These theorems can also be applied if there exists a density of matter $`q`$ which is unscaled and non negative. Such a term $`q`$ adds to $`V(\overline{\psi }).`$
We first prove a calculus lemma.
###### Lemma 8.7
Consider the following algebraic function of $`y`$ with $`a>0,r>0`$ and $`d0:`$
$$f(y)dy^{\frac{n}{n2}}ry^{\frac{n1}{n2}}+a$$
(8.29)
There are two real numbers $`y_1`$ and $`y_2,`$ such that $`0<y_1y_2`$ and
$$f(y_1)0,\text{ }f(y_2)0$$
(8.30)
if
$$\text{ }ad^{n1}<[\frac{(n1)^{n1}}{n^n}]r^n$$
(8.31)
Proof. Suppose $`d>0.`$ The function $`f`$ starts from $`a>0`$ for $`y=0,`$ decreases when $`y`$ increases from $`0`$ to $`y_m=[\frac{(n1)r}{nd}]^{n2},`$ then increases up to infinity with $`y.`$ The numbers $`y_1`$ and $`y_2`$ exist with the indicated properties if $`f(y_m)<0;`$ that is if the inequality 8.31 is satisfied. This inequality always holds if $`d=0:`$ $`f(y)`$ starts then from $`a>0`$ and decreases to $`\mathrm{},`$ so we can then verify that the numbers $`y_1`$ and $`y_2`$ exist.
We use this lemma to prove the following result.
###### Theorem 8.8
Let $`(M,\gamma )`$ be a $`M_{2,\delta }^p`$ manifold with $`p>\frac{n}{2}.`$ Let $`a,b,rW_{0,\delta +2}^p`$ be given on $`(M,\gamma )`$, with $`a0,`$ $`r0,`$ $`b0`$ and $`\delta >\frac{n}{p}+\frac{n}{2}1`$ ($`\delta 1`$ if $`p=2`$). The equation
$$\mathrm{}_\gamma \phi r\phi +a\phi ^{\frac{3n2}{n2}}b\phi ^{\frac{n+2}{n2}}=0$$
(8.32)
has a solution $`\phi >0`$, with $`1\phi W_{2,\delta }^p,`$ if $`\delta <n2\frac{n}{p}`$ and if the inequality 8.31 is satisfied on $`M,`$ with $`d=b,`$ and so long as
$$\underset{xM}{inf}y_1(x)>0,\underset{xM}{inf}y_2(x)\mathrm{max}\{1,\underset{xM}{sup}z_1(x)\},$$
(8.33)
where $`y_1(x)`$ and $`y_2(x)`$ are the two positive numbers which annul the algebraic function<sup>13</sup><sup>13</sup>13Polynomial in the case n=3.
$$f_x(z)b(x)y^{\frac{n}{n2}}r(x)y^{\frac{n1}{n2}}+a(x).$$
(8.34)
Proof. The equation admits a constant subsolution $`\phi _{}=\mathrm{}>0`$ and a constant supersolution $`\phi _+=m1,`$ $`\mathrm{},`$ and therefore a solution $`\phi `$ with the given properties, so long as almost every $`xM`$ it holds that
$$f_x(\mathrm{}^4)0,f_x(m^4)0.$$
(8.35)
The lemma, and the inequalities 8.33 insure the existence of such numbers $`\mathrm{}`$ and $`m,`$ given by
$$\mathrm{}=\mathrm{min}\{1,\underset{xM}{inf}z_1(x)\},m=\underset{xM}{inf}z_2(x).$$
(8.36)
The next theorem does not rely on the sub - super solution method. It supposes that $`r0,`$ hence applies in particular to data satisfying the generalized positive Yamabe condition, after their conformal transformation to the case $`r=0.`$ It has a simpler formulation, but it restricts the size of the coefficients $`a,r`$ and $`b.`$
###### Theorem 8.9
Let $`(M,\gamma )`$ be a $`M_{2,\delta }^p`$ manifold with $`p>\frac{n}{2}.`$ Let $`a,bW_{0,\delta +2}^p`$ be given on $`(M,\gamma )`$, $`a0,`$ while $`b0,`$ $`r0,`$ $`\delta >\frac{n}{p}+\frac{n}{2}1`$ ($`\delta 1`$ if $`p=2`$). The equation
$$\mathrm{}_\gamma \phi r\phi +a\phi ^{\frac{3n2}{n2}}b\phi ^{\frac{n+2}{n2}}=0$$
(8.37)
has a solution $`\phi >0`$, with $`1\phi W_{2,\delta }^p,`$ if $`\delta <n2\frac{n}{p}`$ and if $`a,b`$ and $`r`$ are small enough in the $`W_{0,\delta +2}^p`$ norm.
Proof. The equation admits the subsolution $`\phi =1.`$ We solve it by iteration, starting from $`u_0=1\phi _0=0.`$ We set
$$\mathrm{}_\gamma u_1=a+b+r0,$$
(8.38)
and we see that $`u_1`$ exists, $`u_10,u_1W_{2,\delta }^p`$ with
$$u_1_{W_{2,\delta }^p}C_EM$$
(8.39)
where $`C_E`$ is a number depending only on $`\gamma ,`$ through the constant $`C_E`$ of the elliptic estimate, and where we have set
$$M:=A+B+R\text{}Aa_{W_{0,\delta +2}^p},\text{ }Bb_{W_{0,\delta +2}^p},\text{ }Rr_{W_{0,\delta +2}^p}.$$
(8.40)
The Sobolev embedding theorem $`W_{2,\delta }^pC_\alpha ^0`$ implies then the following inequality where $`C_S`$ is a Sobolev constant
$$u_1_{C_\alpha ^0}C_Su_1_{W_{2,\delta }^p}CM,\text{ with }C:=C_SC_E$$
(8.41)
This inequality implies that
$$\phi _1_{C^0}1+M.$$
(8.42)
Recursively, we suppose $`u_{n1}0`$ and $`u_{n1}_{W_{2,\delta }^p}C_EM;`$ hence $`u_{n1}_{C_\alpha ^0}CM`$. The equation defining $`u_n,`$
$$\mathrm{}_\gamma u_n=r\phi _{n1}+a\phi _{n1}^{\frac{3n2}{n2}}b\phi _{n1}^{\frac{n+2}{n2}},$$
(8.43)
implies $`u_n0`$ and also that
$$u_n_{W_{2,\delta }^p}C_E\{A+R(1+CM)+B(1+CM)^{\frac{n+2}{n2}}\}.$$
(8.44)
Hence $`u_n_{W_{2,\delta }^p}C_EMC_E(A+B+R)`$ if
$$A+R(1+M)+B(1+M)^{\frac{n+2}{n2}}A+B+R,$$
(8.45)
that is:
$$RM+B[(1+M)^{\frac{n+2}{n2}}1]0.$$
(8.46)
This inequality is satisfied if $`A,B,R`$ are small enough. The sequence $`u_n`$ is then uniformly bounded in $`W_{2,\delta }^p.`$ The proof can be completed by the usual methods of functional analysis.
## 9 Uncoupled system of constraints.
The conformally formulated momentum and hamiltonian constraints for the Einstein - scalar field system decouple, in the asymptotically Euclidean case if the initial manifold $`M`$ is maximal. When the constraints decouple the theorems of the previous sections are sufficient to give existence, non-existence or uniqueness theorems of the systems of constraints. The previously obtained results give, for example, the following theorems under a common hypothesis on the a priori given conformal data.
###### Theorem 9.1
Let $`(M,\gamma )`$ be a $`M_{2,\delta }^p`$ manifold; $`\overline{\psi }W_{2,\delta }^p`$ a scalar field with potential $`V(\overline{\psi })W_{0,\delta +2}^p;`$ $`\stackrel{~}{\pi }W_{1,\delta +1}^p`$ a second scalar field, and $`UW_{\delta +1\text{ }}^p`$ a symmetric 2 - tensor. We assume that $`p>\frac{n}{2},`$ $`\delta >1+\frac{n}{2}\frac{n}{p}`$ and $`\delta <2+n\frac{n}{p}`$ ($`\delta =1`$ is admissible if $`p=2).`$ Then the conformally formulated constraints 7.1 and 8-1 on a maximal submanifold ($`\tau =0)`$ admit a solution $`X,`$ $`\phi =1+u>0,`$ with $`X,uW_{2,\delta }^p`$ if either the hypothesis of the theorem 8.5, or 8.8, or 8.9 are satisfied.
Proof. We have already proven that under the given hypotheseses the constraint 7.1 has a unique solution, $`XW_{2,\delta }^p,`$ therefore $`\stackrel{~}{K}W_{1,\delta +1}^p`$ and $`aW_{0,\delta +2}^p`$ (Sobolev embedding and multiplication 6.2, 6.3). We know also (lemma 8.1) that $`rW_{0,\delta +2}^p.`$ Therefore the coefficients of the Lichnerowicz equation (given by equation 4.5 with $`\tau =0)`$ satisfy the hypothesis required in the quoted theorems. It has a solution $`\phi >0,`$ $`\phi 1W_{2,\delta }^p`$, and the pair $`X,\phi `$ satisfies the conformally formulated constraints.
This solution is unique in the cases for which the solution of the Lichnerowicz equation is unique.
###### Remark 9.2
The theorem still holds if in addition to the scalar field $`\overline{\psi }`$ there exists unscaled sources with zero momentum and energy density $`q,`$ and we set $`b\frac{n2}{n1}\{V(\overline{\psi })+q\}.`$ This is so because the constraint equations still decouple (assuming $`\tau 0)`$ when unscaled matter sources are present if these sources have a zero momentum<sup>14</sup><sup>14</sup>14Dain and Nagy \[D-N\] 2002 consider unscaled sources with scaled momentum on a maximal submanifold, using H. Friedrich conformal compactification...
## 10 Coupled system of constraints.
In this section we prove a theorem for the case in which the constraints do not decouple. This result is in the spirit of a stability theorem.The use of the implicit function theorem is the simplest way of proving existence of solutions of equations in the neighbourhood of a given one.
We consider as given the $`M_{2,\delta }^p`$ manifold $`(M,\gamma )`$ together with the scalar functions $`\overline{\psi },`$ $`V(\overline{\psi }),`$ $`\stackrel{~}{\pi }`$ and the traceless symmetric 2-tensor $`U`$, with $`\overline{\psi }W_{2,\delta }^p,`$ $`V(\overline{\psi })W_{0,\delta +2}^p,`$ $`\stackrel{~}{\pi },`$ $`UW_{1,\delta +1}^p`$. We consider the existence of a solution $`\phi `$ and $`X`$ of the constraints 4.4, 4.7 as we perturb $`\tau W_{1,\delta +1}^p`$ away from zero.
We define as follows a mapping $``$ from open sets of a pair of Banach spaces into another Banach space:
$$\text{: (}W_{1,\delta +1}^p;W_{2,\delta }^p\times W_{2,\delta }^p\phi >0)W_{0,\delta +2}^p\times W_{0,\delta +2}^p,\text{ }p>\frac{n}{2},\text{ }\delta >\frac{n}{p}$$
(10.1)
by
$$(\tau ;X,u\phi 1)((\tau ;\phi ,X),(\tau ;\phi ,X))$$
(10.2)
where $``$ and $``$ are the left hand sides of the conformal formulation 4.4, 4.7 of the constraints.
The multiplication properties of weighted Sobolev spaces show that $``$ is a $`C^1`$ mapping. The partial derivative $`_{X,u}^{}`$ at a point $`(0;X,u)`$ is the linear mapping from $`W_{2,\delta }^p\times W_{2,\delta }^p`$ into $`W_{0,\delta +2}^p`$ given by
$$(\delta X,\delta u)(\delta ,\delta )$$
(10.3)
where ($`\delta b=0`$ because $`\tau =0`$ at the considered point and $`V(\overline{\psi })`$ is fixed)
$$\delta \mathrm{\Delta }_\gamma \delta u\alpha \delta u+\phi ^{\frac{3n2}{n2}}\delta a,\text{ }$$
(10.4)
with
$$\alpha =r+\frac{3n2}{n2}a\phi ^{4\frac{n1}{n2}}+\frac{n+2}{n2}b\phi ^{\frac{4}{n2}},$$
(10.5)
and, using the expressions for $`a`$ and $`\stackrel{~}{K},`$
$$\delta a=\frac{n2}{2(n1}\stackrel{~}{K}(\stackrel{~}{}_{\gamma ,conf})\delta X.$$
(10.6)
On the other hand
$$\delta (\stackrel{~}{\mathrm{\Delta }}_{\gamma ,conf}\delta X)^i.$$
(10.7)
###### Theorem 10.1
Specify on the $`M_{2,\delta }^p`$ manifold $`(M,\gamma )`$ the scalar functions $`\overline{\psi },`$ $`V(\overline{\psi }),`$ $`\stackrel{~}{\pi },`$ $`\stackrel{~}{N}`$ and the traceless symmetric 2-tensor $`U`$, with $`\overline{\psi }W_{2,\delta }^p,`$ $`V(\overline{\psi })W_{0,\delta +2}^p,`$ $`\stackrel{~}{\pi },`$ $`UW_{1,\delta +1}^p,`$ $`\stackrel{~}{N}1W_{s+2,\delta }^p`$, $`\stackrel{~}{N}>0,p>\frac{n}{2},`$ $`\frac{n}{p}<\delta <n2\frac{n}{p}`$. Let $`(X_0,\phi _0)`$ be a solution of the corresponding constraints with $`\tau _0=0`$. Suppose that for some $`\stackrel{~}{\delta }>1+\frac{n}{2}\frac{n}{p}`$ ($`\stackrel{~}{\delta }=1`$ if $`p=2)`$ it holds that:
$$_M\{|Df|_\gamma ^2+\alpha _0f^2\}𝑑\mu _\gamma >0\text{ for all }fW_{2,\stackrel{~}{\delta }}^p,\text{ }f0$$
(10.8)
with
$$\alpha _0:=r+\frac{3n2}{n2}a_0\phi _0^{4\frac{n1}{n2}}\frac{n+2}{n2}V(\overline{\psi })\phi _0^{\frac{4}{n2}}0.$$
(10.9)
Then there exists a neighbourhood $`\mathrm{\Omega }`$ of zero in $`W_{1,\delta +1}^p`$ such that, if $`\tau \mathrm{\Omega }`$, the coupled constraints have one and only one solution $`(X,\phi ),`$with $`\phi >0,`$ and$`X,u\phi 1W_{2,\delta }^p.`$
Proof. Under the hypotheses that we have made, due to properties of elliptic equations discussed above, the partial derivative of $``$ with respect to the pair $`(u,X)`$ determines an isomorphism from $`W_{2,\delta }^p\times W_{2,\delta }^p`$ onto $`W_{0,\delta +2}^p,`$ given by
$$(\delta u,\delta X)(\delta ,\delta ).$$
(10.10)
A straightforward application of the implicit function theorem then completes the proof.
###### Remark 10.2
The conclusion of the theorem holds in particular if $`(M,\gamma ,\overline{\psi })`$ satisfy the inequality 8.11 and $`V(\overline{\psi })0,`$ since then one has always $`a_00.`$
###### Remark 10.3
An analogous method can be used to prove the existence of solutions of the coupled system in the additional presence of matter sources with momentum small enough in $`W_{0,\delta +2}^p`$ norm.
Acknowledgments. This work was partially done in the Department of Mathematics of the University of Washington, whose hospitality is gratefully aknowledged by YCB and JI.
This research was partially supported by NSF Grant PHY-0354659 at the University of Oregon, and by NSF Grant DMS-0305048 at the University of Washington.
YCB: [email protected] Académie des Sciences, 23 Quai Conti, 75270 Paris 06, France
JI: [email protected] Department of Mathematics, University of Oregon, Eugene, OR. 97403 USA
DP: [email protected] Department of Mathematics, University of Washington, Seattle, WA, 98195-4350, USA
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# 1 Introduction and Notations.
## 1 Introduction and Notations.
The transposition mirror construction has been proposed by P.Berglund and T. Hübsch as a trial to generalize so called Greene-Plesser mirror construction that comprises mirror pairs of Fermat type hypersurfaces. Later, in the article , in relying on the transposition method, the authors have proposed a natural hypothesis on the period integrals associated to the complete intersection (CI) Calabi-Yau variety $`X`$ that is supposed to be a mirror symmetry to the generic multi-quasihomogeneous Calabi-Yau variety $`Y`$ of codimension $`\mathrm{}`$ defined in the product of the quasihomogeneous projective spaces $`𝐏_{(g_1^{(1)},\mathrm{},g_{\tau _1+1}^{(1)})}^{(\tau _1)}\times \mathrm{}\times 𝐏_{(g_1^{(k)},\mathrm{},g_{\tau _k+1}^{(k)})}^{(\tau _k)}`$. They mean under the notion of the mirror symmetry between $`X`$ and $`Y`$ an interchange between geometric symmetry $`(𝒢_X,`$ $`𝒢_Y)`$ and quantum symmetry $`(𝒬_Y,`$ $`𝒬_X)`$ groups of each varieties (See Definition 1, Theorem 4.2 below). That is to say, for the mirror symmetry pair $`X`$ and $`Y`$, the following isomorphisms holds,
$$𝒬_X𝒢_Y,$$
$$𝒬_Y𝒢_X.$$
In this article we propose certain sufficient conditions on $`X`$ and $`Y`$ so that their hypothesis ( §3.3) on the period integrals holds (Theorem 3.1). Namely the period integrals defined on $`X`$ can be expressed by means of quasihomogeneous weights of its mirror symmetry $`Y`$ and vice versa. It will be shown that our sufficient conditions entail the mirror symmetry between $`X`$ and $`Y`$ in the above sense of (see Theorem 4.2).
AMS Subject Classification: 14J32 primary), 32S25, 32C65 (secondary).
Key words and phrases: mirror symmetry, hypergeometric functions, monodromy.
partially supported by Hokkaido University, ICTP (Trieste).
During the entire article we shall restrict ourselves to the case $`\mathrm{}=k.`$
In accordance with the suggestion on the mirror symmetry to the generic multi-quasihomogeneous Calabi-Yau made in , we shall consider the following system of equations on $`𝐓^n:=(𝐂^\times )^n`$ as defining equations of $`X=X_\mathrm{𝟏}`$,
$$X_s=\{x𝐓^n;f_1(x)+s_1(f_2(x)1)=\mathrm{}=f_{2k1}(x)+s_k(f_{2k}(x)1)=0\},$$
with system $`(1.1)`$ below. Here we use the notation $`s=(s_1,\mathrm{},s_k)𝐓^k,`$ $`\mathrm{𝟏}=(1,1,\mathrm{},1)𝐑^k.`$
$$f_1(x)=x^{\stackrel{}{v}_1^{(1)}}+\mathrm{}+x^{\stackrel{}{v}_{\tau _1}^{(1)}},$$
$`(1.1)`$
$$f_2(x)=\underset{jI^{(1)}}{}x_j+1,$$
$$\mathrm{},$$
$$f_{2i1}(x)=x^{\stackrel{}{v}_1^{(i)}}+\mathrm{}+x^{\stackrel{}{v}_{\tau _i}^{(i)}},$$
$$f_{2i}(x)=\underset{jI^{(i)}}{}x_j+1,$$
$$\mathrm{},$$
$$f_{2k1}(x)=x^{\stackrel{}{v}_1^{(k)}}+\mathrm{}+x^{\stackrel{}{v}_{\tau _k}^{(k)}},$$
$$f_{2k}(x)=\underset{jI^{(k)}}{}x_j+1.$$
Here $`I^{(j)},1jk`$ are sets of indices which are complementary one another in such a way that $`_{q[1,k]}I^{(q)}=\{1,\mathrm{},n\}`$ and $`I^{(q)}I^{(q^{})}=\mathrm{}`$ if $`qq^{}.`$ From now on we shall make use of the notations $`\stackrel{~}{\tau }_\nu :=|I^{(\nu )}|`$ and $`b^q:=_{\nu =1}^q\tau _\nu .`$ Additionally we suppose that
$$\underset{\nu =1}{\overset{k}{}}\tau _\nu =b^k=n.$$
The equation $`f_{2j1}(x)`$ (resp.$`f_{2j}(x)`$) is defined by the monomials with powers $`\stackrel{}{v}_1^{(j)},\mathrm{},\stackrel{}{v}_{\tau _j}^{(j)}`$ $`𝐙^n`$ (resp. $`\stackrel{}{v}_{\tau _j+1}^{(j)}`$ $`𝐙^n`$) such that for the weight vector $`\stackrel{}{g}^{(q)}=(\stackrel{b^{q1}}{\stackrel{}{0,\mathrm{},0}},g_1^{(q)},\mathrm{},g_{\tau _q}^{(q)},\stackrel{nb^q}{\stackrel{}{0,\mathrm{},0}})`$ $`𝐙_0^n,`$ $`1qk`$ the following quasihomogeneiety condition holds,
$$Q_j^{(q)}:=<\stackrel{}{v}_1^{(j)},\stackrel{}{g}^{(q)}>=\mathrm{}=<\stackrel{}{v}_{\tau _j}^{(j)},\stackrel{}{g}^{(q)}>=<\stackrel{}{v}_{\tau _{j+1}}^{(j)},\stackrel{}{g}^{(q)}>,1jk.$$
$`(1.2)`$
This means that the point $`\stackrel{}{v}_{\tau _j+1}^{(j)}`$ belongs to the $`(\tau _j1)`$-dimensional hyperplane generated by $`\stackrel{}{v}_1^{(j)},\stackrel{}{v}_1^{(j)},\mathrm{},\stackrel{}{v}_{\tau _j}^{(j)}.`$ The following condition shall be imposed if we suppose that $`X_s`$ is a Calabi-Yau variety:
$$\underset{j=1}{\overset{k}{}}Q_j^{(q)}=\underset{i=1}{\overset{\tau _q}{}}g_i^{(q)}=<\stackrel{}{g}^{(q)},(\stackrel{b^{q1}}{\stackrel{}{0,\mathrm{},0}},\stackrel{\tau _q}{\stackrel{}{1,\mathrm{},1}},\stackrel{nb^q}{\stackrel{}{0,\mathrm{},0}})>,1qk.$$
$`(1.3)`$
We will impose this condition on $`(1.1)`$ in the further arguments.
In addition to that we assume that for each element $`\lambda Aut(X_j),`$ of the group automorphism of the hypersurface $`X_j=\{x𝐓^n;f_{2j1}(x)=0\}`$ the relation $`(\lambda _{}f_{2j1})(x)=\lambda _{}(x^{\stackrel{}{v}_{\tau _j+1}^{(j)}})`$ holds. More precisely, every $`\lambda Aut(X_s)`$ admits the following decomposition
$$\lambda =𝐪𝐠𝐡,$$
where $`𝐪𝒬_X`$, $`𝐠𝒢_X`$ and $`𝐡`$ each of which is a non cyclic element of $`Aut(X_s)`$. The subgroup $`𝒬_X_{q=1}^k𝐙_{\overline{𝐐}^{(𝐪)}},`$ with $`\overline{𝐐}^{(𝐪)}=L.C.M.(Q_1^{(q)},\mathrm{},Q_k^{(q)})`$ is a cyclic group generated by $`k`$ different cyclic actions $`𝐪^{(\nu )},1\nu k,`$ corresponding to the quasihomogeneiety,
$$𝐪_{}^{(\nu )}:(x_1,\mathrm{},x_n)(x_1,\mathrm{},x_{b^{\nu 1}},e^{\frac{2\pi ig_1^{(\nu )}}{\overline{𝐐}^{(\nu )}}}x_{b^{\nu 1}+1},\mathrm{},e^{\frac{2\pi ig_{\tau _\nu }^{(\nu )}}{\overline{𝐐}^{(\nu )}}}x_{b^\nu },x_{b^\nu +1},\mathrm{},x_n).$$
The group (called geometric symmetry) $`𝒢_X`$ consists of elements $`𝐠𝒬_X`$ of the following form
$$𝐠_{}:(x_1,\mathrm{},x_n)(e^{\frac{2\alpha _1\pi i}{d}}x_1,e^{\frac{2\alpha _2\pi i}{d}}x_2,\mathrm{},e^{\frac{2\alpha _n\pi i}{d}}x_n),$$
for some $`d>0`$ and $`(\alpha _1,\mathrm{},\alpha _n)𝐙^n.`$ Here we remark that some of $`\alpha _i`$ can be zero. The group $``$ is the non cyclic part of the group $`Aut(X_s)`$.
###### Definition 1
In the following decomposition,
$$Aut(X_s)𝒬_X\times 𝒢_X\times ,$$
we call $`𝒬_X`$ (resp. $`𝒢_X`$) the quantum symmetry (resp. geometric symmetry) of $`X_s`$.
Further we apply so called Cayley trick to $`(1.1)`$ to get a polynomial
$$F(x,s,y)=\underset{j=1}{\overset{k}{}}y_{2j1}(f_{2j1}(x)+s_j)+\underset{j=1}{\overset{k}{}}y_{2j}f_{2j}(x),$$
$`(1.4)`$
with $`L=n+3k`$ terms. The procedures $`(2.3)`$, $`(2.4)`$ below explain why we consider this polynomial $`F(x,s,y)`$ to calculate the period integrals associated to $`X_s`$. One may consult and for more details on the utility of the Cayley trick in the calculus of period integrals.
To manipulate the polynomial $`F(x,s,y)`$ we introduce the notation $`a^\nu :=\tau _1+\mathrm{}+\tau _\nu +3\nu =b^\nu +3\nu .`$ In particular, the $`a^{i1}+1`$th term of $`F(x,s,y)`$ corresponds to $`y_{2i1}x^{\stackrel{}{v}_1^{(i)}}`$ and the $`a^i3=(a^{i1}+\tau _i)`$th term - $`y_{2i1}x^{\stackrel{}{v}_{\tau _i}^{(i)}}.`$ The $`(a^i1)`$th term - $`y_{2i}_{jI^{(i)}}x_j.`$ The$`(a^i2)`$th term - $`y_{2i}s_i.`$ The $`a^i`$th term - $`y_{2i}`$.
From the polynomial $`F(x,s,y)`$ we construct a matrix $`𝖫`$ that consists of the row vectors $`r`$th term of which corresponds to the power of the $`r`$th monomial term present in $`F(x,s,y).`$ For example, the row vector $`\overline{v}_q^{(\nu )}`$ corresponds to the row number $`a^{q1}+\nu `$ of the matrix $`𝖫`$,
$$\overline{v}_q^{(\nu )}=(\stackrel{}{v}_q^{(\nu )},\stackrel{2\nu 2}{\stackrel{}{0,\mathrm{},0}},1,\stackrel{3k2\nu +1}{\stackrel{}{0,\mathrm{},0}}),$$
$`1\nu k,1q\tau _\nu .`$
Next we look at the system of linear equations $`\mathrm{\Xi }=(\xi _1,\mathrm{},\xi _L),`$
$${}_{}{}^{t}𝖫\mathrm{\Xi }=^t(1,\mathrm{},1,z_1,\mathrm{},z_k),$$
$`(1.5)`$
that is equivalent to the relation,
$$\mathrm{\Xi }=^t𝖫^1^t(1,\mathrm{},1,z_1,\mathrm{},z_k).$$
$`(1.6)`$
Thus we get linear functions $`(\xi _1(z),\mathrm{},\xi _L(z))`$ that will be later denoted by $`(_1(0,z,0),\mathrm{},_L(0,z,0))`$ in the notation $`(2.5)`$ below. These linear functions will play essential role in the calculus of the period integrals.
We define a $`(n\times k)`$ matrix $`V^\mathrm{\Lambda }`$ as follows:
$$V^\mathrm{\Lambda }:=(^t\stackrel{}{v}_{\tau _1+1}^{(1)},\mathrm{},^t\stackrel{}{v}_{\tau _k+1}^{(k)}),$$
$`(1.7)`$
where $`\stackrel{}{v}_{\tau _q+1}^{(q)}`$ is a $`n`$row vector that corresponds to $`supp(f_{2q})\{0\}.`$ By virtue of the quasihomogeneiety $`(1.2)`$ we can define a $`k\times k`$ matrix as follows:
$$\widehat{Q}:=\left(\begin{array}{ccc}Q_1^{(1)}& \mathrm{}& Q_1^{(k)}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ Q_k^{(1)}& \mathrm{}& Q_k^{(k)}\end{array}\right)=^tV^\mathrm{\Lambda }({}_{}{}^{t}\stackrel{}{g}_{}^{(1)},\mathrm{},^t\stackrel{}{g}^{(k)}).$$
$`(1.8)`$
For the simplicity of the formulation, we will use a diagonal matrix
$$G=diag(g_1^{(1)},\mathrm{},g_{\tau _1}^{(1)},g_1^{(2)},\mathrm{},g_{\tau _2}^{(2)},\mathrm{},g_1^{(k)},\mathrm{},g_{\tau _k}^{(k)}).$$
$`(1.9)`$
We introduce a $`n\times n`$ matrix:
$$𝖫_\mathrm{\Lambda }:=\left(\begin{array}{c}\stackrel{}{v}_1^{(1)}\\ \stackrel{}{v}_2^{(1)}\\ \mathrm{}\\ \stackrel{}{v}_{\tau _1}^{(1)}\\ \mathrm{}\\ \stackrel{}{v}_{\tau _k}^{(k)}\end{array}\right).$$
$`(1.10)`$
We construct a matrix $`{}_{}{}^{T}𝖫`$ constructed from the transposed matrix $`{}_{}{}^{t}𝖫`$ after some proper permutations of the rows and columns such that each row of $`{}_{}{}^{T}𝖫`$ corresponds to a vertex of a polynomial
$${}_{}{}^{T}F(x,s,y)=\underset{j=1}{\overset{k}{}}y_{2j1}(^Tf_{2j1}(x)+s_j)+\underset{j=1}{\overset{k}{}}y_{2j}^Tf_{2j}(x),$$
$`(1.4)^T`$
for the polynomials,
$${}_{}{}^{T}f_{2q1}^{}(x)=x^{{}_{}{}^{T}\stackrel{}{v}_{1}^{(q)}}+\mathrm{}+x^{{}_{}{}^{T}\stackrel{}{v}_{\stackrel{~}{\tau }_q}^{(q)}},$$
$`(1.1)^T`$
$${}_{}{}^{T}f_{2q}^{}(x)=\underset{\mathrm{}^TI^{(q)}}{}x_{\mathrm{}}+1,,\mathrm{\hspace{0.33em}1}qk.$$
Here we impose the condition
$$\{\tau _1,\mathrm{},\tau _k\}=\{|^TI^{(1)}|,\mathrm{},|^TI^{(k)}|\}=\{\stackrel{~}{\tau }_1,\mathrm{},\stackrel{~}{\tau }_k\}=\{|I^{(1)}|,\mathrm{},|I^{(k)}|\}.$$
$`(1.11)`$
Thus we can define an one to one mapping
$$\nu :[1,k][1,k],$$
such that $`|^TI^{(\nu (j))}|=\tau _{\nu (j)}=\stackrel{~}{\tau }_j=|I^{(j)}|`$. Further we impose a condition $`\nu ^2=id`$ so that $`{}_{}{}^{T}(^TF(x,s,y))=F(x,s,y)`$ holds. In a way parallel to $`(1.7)(1.10)`$, we can define the weight system
$$(^T\stackrel{}{g}^{(1)},\mathrm{},^T\stackrel{}{g}^{(k)}),$$
$${}_{}{}^{T}Q_{j}^{(q)}=^T\stackrel{}{v}_r^{(j)},^T\stackrel{}{g}^{(q)},\mathrm{\hspace{0.33em}1}r\stackrel{~}{\tau }_q,1j,qk.$$
$`(1.2)^T`$
We impose a condition necessary for Calabi-Yau property of $`Y`$
$$\underset{j=1}{\overset{k}{}}{}_{}{}^{T}Q_{j}^{(q)}=\underset{i=1}{\overset{\stackrel{~}{\tau }_q}{}}{}_{}{}^{T}g_{i}^{(q)},1qk.$$
$`(1.3)^T`$
It is easy to see that the equations of $`{}_{}{}^{T}(1.1)`$ define a CI in $`𝐏_{(^Tg_1^{(1)},\mathrm{},^Tg_{\stackrel{~}{\tau }_1+1}^{(1)})}^{(\stackrel{~}{\tau }_1)}\times \mathrm{}\times 𝐏_{(^Tg_1^{(k)},\mathrm{},^Tg_{\stackrel{~}{\tau }_k+1}^{(k)})}^{(\stackrel{~}{\tau }_k)}`$ with $`{}_{}{}^{T}g_{\stackrel{~}{\tau }_q+1}^{(q)}=^TQ_q^{(q)}.`$
In analogy with the Definition 1 we define the following decomposition
$$Aut(Y_s)𝒬_Y\times 𝒢_Y\times ^T,$$
for
$$Y_s=\{x𝐓^n;^Tf_1(x)+s_1(^Tf_2(x)1)=\mathrm{}=^Tf_{2k1}(x)+s_k(^Tf_{2k}(x)1)=0\}.$$
The quantum symmetry $`𝒬_Y_{q=1}^k𝐙_{{}_{}{}^{𝐓}\overline{𝐐}_{}^{(𝐪)}},`$ with $`{}_{}{}^{𝐓}\overline{𝐐}_{}^{(𝐪)}=L.C.M.(^TQ_1^{(q)},\mathrm{},^TQ_k^{(q)})`$ is a cyclic group corresponding to the quasihomogeneiety and $`𝒢_Y`$ the remaining cyclic part called geometric symmetry. The group $`{}_{}{}^{T}`$ is the remaining non cyclic part of $`Aut(Y_s).`$
We introduce matrices analogous to the case of $`(1.1),`$
$${}_{}{}^{T}V_{}^{\mathrm{\Lambda }}:=(^t(^T\stackrel{}{v}_{\tau _1+1}^{(\nu (1))}),^t(^T\stackrel{}{v}_{\tau _2+1}^{(\nu (2))}),\mathrm{},^t(^T\stackrel{}{v}_{\tau _k+1}^{(\nu (k))})),$$
$`(1.7)^T`$
where $`{}_{}{}^{T}\stackrel{}{v}_{\tau _q+1}^{(\nu (q))}`$ is a $`n`$row vector which corresponds to $`supp(^Tf_{2\nu (q)})\{0\}.`$ More precisely, the $`j`$th column of the matrix $`{}_{}{}^{T}V_{}^{\mathrm{\Lambda }}`$ equals to
$${}_{}{}^{t}(^T\stackrel{}{v}_{\tau _j+1}^{(\nu (j))})=^t(\stackrel{b^{\nu (j)1}}{\stackrel{}{0,\mathrm{},0}},\stackrel{\tau _{\nu (j)}}{\stackrel{}{1,\mathrm{},1}},\stackrel{nb^{\nu (j)}}{\stackrel{}{0,\mathrm{},0}}),\tau _{\nu (j)}=\stackrel{~}{\tau _j},$$
here $`b^{\nu (j)}:=_{r=1}^{\nu (j)}\tau _r.`$
$${}_{}{}^{T}\widehat{Q}:=\left(\begin{array}{ccc}{}_{}{}^{T}Q_{\nu (1)}^{(1)}& \mathrm{}& {}_{}{}^{T}Q_{\nu (1)}^{(k)}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ {}_{}{}^{T}Q_{\nu (k)}^{(1)}& \mathrm{}& {}_{}{}^{T}Q_{\nu (k)}^{(k)}\end{array}\right)=^t(^TV^\mathrm{\Lambda })({}_{}{}^{t}(^T\stackrel{}{g}^{(1)}),\mathrm{},^t(^T\stackrel{}{g}^{(k)})).$$
$`(1.8)^T`$
$${}_{}{}^{T}G=diag(^Tg_1^{(1)},\mathrm{},^Tg_{\stackrel{~}{\tau }_1}^{(1)},^Tg_1^{(2)},\mathrm{},^Tg_{\stackrel{~}{\tau }_2}^{(2)},\mathrm{},^Tg_1^{(k)},\mathrm{},^Tg_{\stackrel{~}{\tau }_k}^{(k)}).$$
$`(1.9)^T`$
$${}_{}{}^{T}𝖫_{\mathrm{\Lambda }}^{}:=\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{v}_{1}^{(\nu (1))}\\ {}_{}{}^{T}\stackrel{}{v}_{2}^{(\nu (1))}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{v}_{\stackrel{~}{\tau }_{\nu (1)}}^{(\nu (1))}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{v}_{\stackrel{~}{\tau }_{\nu (k)}}^{(\nu (k))}\end{array}\right)=\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{v}_{1}^{(\nu (1))}\\ {}_{}{}^{T}\stackrel{}{v}_{2}^{(\nu (1))}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{v}_{\tau _1}^{(\nu (1))}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{v}_{\tau _k}^{(\nu (k))}\end{array}\right).$$
$`(1.10)^T`$
In an analogous way to $`(1.6)`$, we introduce the linear functions $`{}_{}{}^{T}\mathrm{\Xi }:=(^T\xi _1(z),\mathrm{},^T\xi _L(z)),`$ defined by the relation,
$${}_{}{}^{T}\mathrm{\Xi }=^t({}_{}{}^{T}𝖫)^1^t(1,\mathrm{},1,z_1,\mathrm{},z_k).$$
$`(1.6)^T`$
Finally we remark that due to the property $`\nu ^2=id`$ and the definition $`(1.7)^T`$, there exists a permutation matrix $`\lambda SL(n,𝐙)`$ such that
$$\lambda V^\mathrm{\Lambda }=^TV^\mathrm{\Lambda },\lambda {}_{}{}^{t}𝖫_{\mathsf{\Lambda }}^{}={}_{}{}^{T}𝖫_{\mathsf{\Lambda }}^{}.$$
$`(1.12)`$
## 2 Mellin transform of period integrals
In this section we review the results on the period integrals to be used for the verification of the hypothesis in in the subsequent section. See for the detail of proofs , .
Let us consider the Leray’s coboundary (see ) to define the period integral that is equivalent to the period integral of the variety $`X_s`$, $`\gamma H_n(𝐓^n_{\mathrm{}=1}^k\{x𝐓^n:f_{2\mathrm{}1}(X)+s_{\mathrm{}}=0\}_{\mathrm{}=1}^k\{x𝐓^n:f_2\mathrm{}(X)=0\})`$ such that $`\mathrm{}(f_{2\mathrm{}1}(X)+s_{\mathrm{}})|_\gamma <0`$, $`\mathrm{}f_2\mathrm{}(X)|_\gamma <0`$. Further on central object of our study is the following fibre integral,
$$I_{x^𝐢,\gamma }^\zeta (s)=_\gamma (f_1(x)+s_1)^{\zeta _11}(f_2(x))^{\zeta _21}\mathrm{}(f_{2k1}(x)+s_k)^{\zeta _{2k1}1}f_{2k}(x)^{\zeta _{2k}1}x^{𝐢+\mathrm{𝟏}}\frac{dx}{x^\mathrm{𝟏}},$$
$`(2.1)`$
and its Mellin transform,
$$M_{𝐢,\gamma }^\zeta (𝐳):=_\mathrm{\Pi }s^𝐳I_{x^𝐢,\gamma }^\zeta (s)\frac{ds}{s^\mathrm{𝟏}},$$
$`(2.2)`$
for certain cycle $`\mathrm{\Pi }`$ homologous to $`𝐑^k`$ which avoids the singular loci of $`I_{x^𝐢,\gamma }^\zeta (s)`$ (cf. ). Thus the fibre integral $`I_{x^𝐢,\gamma }^\zeta (s)`$ is a ramified function on the torus $`𝐓^k.`$ We introduce the notation $`\gamma ^\mathrm{\Pi }:=_{(s)\mathrm{\Pi }}((s),\gamma ).`$ One shall not confuse it with the thimble of Lefschetz, because $`\gamma ^\mathrm{\Pi }`$ is rather a tube without thimble. It is useful to understand the calculus of the Mellin transform in connection with the notion of the generalized HGF in the sense of Mellin-Barnes-Pincherle . After this formulation, the classical HGF of Gauss can be expressed by means of the integral,
$${}_{2}{}^{}F_{1}^{}(\alpha ,\beta ,\gamma |s)=\frac{1}{2\pi i}_{z_0i\mathrm{}}^{z_0+i\mathrm{}}(s)^z\frac{\mathrm{\Gamma }(z+\alpha )\mathrm{\Gamma }(z+\beta )\mathrm{\Gamma }(z)}{\mathrm{\Gamma }(z+\gamma )}𝑑z,\mathrm{}\alpha ,\mathrm{}\beta <z_0.$$
We can introduce $`(w_1^{\prime \prime },\mathrm{},w_{2k}^{\prime \prime })`$ natural quasihomogeneous weight of $`(y_1,\mathrm{}y_k)`$ so that $`F(x,0,y)`$ of $`(1.4)`$ gets the quasihomogeneous zero weight with respet to the variables $`(x,y).`$ Next we modify the Mellin transform
$$M_{𝐢,\gamma }^\zeta (𝐳)=c(\zeta )_{S_+^{k1}(w^{\prime \prime })\times \gamma ^\mathrm{\Pi }}\frac{x^𝐢\omega ^\zeta s^{𝐳\mathrm{𝟏}}dx\mathrm{\Omega }_0(\omega )ds}{(\omega _1(f_1(x)+s_1)+\mathrm{}+\omega _{2k}(f_{2k}(x)))^{\zeta _1+\mathrm{}+\zeta _{2k}+2k}}$$
$$=c(\zeta )_{𝐑_+}\sigma ^{\zeta _1+\mathrm{}+\zeta _{2k}+2k}\frac{d\sigma }{\sigma }_{S_+^{2k1}(w^{\prime \prime })}\omega ^\zeta \mathrm{\Omega }_0(\omega )_\gamma x^𝐢𝑑x_\mathrm{\Pi }s^𝐳e^{\sigma (\omega _1(f_1(X)+s_1)+\mathrm{}+\omega _{2k}(f_{2k}(x))}\frac{ds}{s^\mathrm{𝟏}},$$
with $`c(\zeta )=\frac{\mathrm{\Gamma }(\zeta _1+\mathrm{}+\zeta _{2k}+2k)}{\mathrm{\Gamma }(\zeta _1+1)\mathrm{}\mathrm{\Gamma }(\zeta _{2k}+1)}.`$ Here we made use of notations $`S_+^{2k1}(w^{\prime \prime })=\{(\omega _1,\mathrm{},\omega _{2k}):\omega _1^{|\frac{𝐰^{\prime \prime }}{w_1^{\prime \prime }}|}+\mathrm{}+\omega _{2k}^{|\frac{𝐰^{\prime \prime }}{w_{2k}^{\prime \prime }}|}=1,\omega _{\mathrm{}}>0`$ for all $`\mathrm{},`$ $`𝐰^{\prime \prime }`$ $`=`$ $`_{1i2k}w_i^{\prime \prime }\}`$ and $`\mathrm{\Omega }_0(\omega )`$ the $`(2k1)`$ volume form on $`S_+^{2k1}(w^{\prime \prime }),`$
$$\mathrm{\Omega }_0(\omega )=\underset{\mathrm{}=1}{\overset{2k}{}}(1)^{\mathrm{}}w_{\mathrm{}}^{\prime \prime }\omega _{\mathrm{}}d\omega _1\text{ }^{\genfrac{}{}{0pt}{}{\mathrm{}}{^{}}}\mathrm{}d\omega _{2k}.$$
In the above transformation we used a classical interpretation of Dirac’s delta function as a residue:
$$_\gamma _{𝐑_+}e^{y_j(f_j(x)+s_j)}y_j^{\zeta _j}𝑑y_jdx=\mathrm{\Gamma }(\zeta _j+1)_\gamma (f_j(x)+s_j)^{\zeta _j1}𝑑x.$$
We will rewrite, up to constant multiplication, the expression obtained as a modification of $`M_{𝐢,\gamma }^\zeta (𝐳)`$ into the following form,
$$_{(𝐑_+)^{2k}\times \gamma ^\mathrm{\Pi }}e^{\mathrm{\Psi }(T)}x^{𝐢+\mathrm{𝟏}}y^{\zeta +\mathrm{𝟏}}s^𝐳\frac{dx}{x^\mathrm{𝟏}}\frac{dy}{y^\mathrm{𝟏}}\frac{ds}{s^\mathrm{𝟏}}$$
where
$$\mathrm{\Psi }(T)=T_1(X,s,y)+\mathrm{}+T_L(x,s,y)=F(x,s,y),$$
$`(2.3)`$
in which each term $`T_i(x,s,y)`$ stands for a monomial in variables $`(x,s,y)`$ of the phase function $`(1.4).`$ We transform the above integral into the following form,
$$_{(𝐑_+)^{2k}\times \gamma ^\mathrm{\Pi }}e^{\mathrm{\Psi }(T(x,s,y))}x^{𝐢+\mathrm{𝟏}}s^𝐳y^{\zeta +\mathrm{𝟏}}\frac{dx}{x^\mathrm{𝟏}}\frac{dy}{y^\mathrm{𝟏}}\frac{ds}{s^\mathrm{𝟏}}$$
$`(2.4)`$
$$=(det𝖫)^1_{L_{}(𝐑_{+}^{}{}_{}{}^{2k}\times \gamma ^\mathrm{\Pi })}e^{_{aI}T_a}\underset{aI}{}T_a^{_a(𝐢,𝐳,\zeta )}\underset{aI}{}\frac{dT_a}{T_a}$$
$$=(1)^{\zeta _1+\mathrm{}+\zeta _{2k}+2k}(det𝖫)^1_{L_{}(𝐑_{+}^{}{}_{}{}^{2k}\times \gamma ^\mathrm{\Pi })}e^{_{aI}T_a}\underset{aI}{}T_a^{_a(𝐢,𝐳,\zeta )}\underset{aI}{}\frac{dT_a}{T_a}.$$
Here $`L_{}(𝐑_{+}^{}{}_{}{}^{2k}\times \gamma ^\mathrm{\Pi })`$ means the image of the chain in $`𝐂_X^M\times 𝐂_s^k\times 𝐂_y^{2k}`$ into that in $`𝐂_T^L`$ induced by the transformation $`(2.3).`$ We define $`L_{}(𝐑_{+}^{}{}_{}{}^{2k}\times \gamma ^\mathrm{\Pi })`$ $`=\{(T_1,\mathrm{},T_L)`$ $`𝐂^L;`$ $`(T_1,\mathrm{},T_L)`$ $`L_{}(𝐑_{+}^{}{}_{}{}^{2k}\times \gamma ^\mathrm{\Pi }),\mathrm{}T_a>0,a[1,L]\}.`$ The second equality of $`(2.4)`$ follows from Proposition 2.1, 3) below that can be proven in a way independent of the argument to derive $`(2.4).`$ We will denote the set of columns and rows of the matrix $`𝖫`$ by $`I,`$
$$I:=\{1,\mathrm{},L\}.$$
Here we remember the condition $`L=n+3k`$ imposed on $`(1.4)`$.
The following notion helps us to formulate the result in a compact manner.
###### Definition 2
A meromorphic function $`g(𝐳)`$ is called $`\mathrm{\Delta }`$periodic for $`\mathrm{\Delta }𝐙_{>0},`$ if
$$g(𝐳)=h(e^{2\pi \sqrt{}1\frac{z_1}{\mathrm{\Delta }}},\mathrm{},e^{2\pi \sqrt{}1\frac{z_k}{\mathrm{\Delta }}}),$$
for some rational function $`h(\zeta _1,\mathrm{},\zeta _k).`$
For the CI $`(1.1)`$ (i.e. we can construct $`F(x,s,y)`$ for which the matrix $`𝖫`$ is non-degenerate), we have the following statement.
###### Proposition 2.1
1)For any cycle $`\mathrm{\Pi }H_k(𝐓^kS.S.I_{x^𝐢,\gamma }^\zeta (s))`$ the Mellin transform $`(2.1)`$ can be represented as a product of $`\mathrm{\Gamma }`$ function factors up to a $`\mathrm{\Delta }`$periodic function factor $`g(𝐳)`$,
$$M_{𝐢,\gamma }^\zeta (𝐳)=g(𝐳)\underset{aI}{}\mathrm{\Gamma }\left(_a(𝐢,𝐳,\zeta )\right),$$
with
$$_a(𝐢,𝐳,\zeta )=\frac{\left(_{j=1}^nA_j^a(i_j+1)+_{\mathrm{}=1}^kB_{\mathrm{}}^az_{\mathrm{}}+_{\mathrm{}=1}^{2k}D_{\mathrm{}}^a(\zeta _{\mathrm{}}+1)\right)}{\mathrm{\Delta }},aI.$$
$`(2.5)`$
Here the following matrix $`\mathrm{\Delta }^1𝖳=(𝖫)^\mathrm{𝟣}`$ has integer elements,
$${}_{}{}^{t}𝖳=(A_1^a,\mathrm{},B_1^a,\mathrm{},B_k^a,D_1^a,\mathrm{},D_{2k}^a)_{1aL},$$
$`(2.6)`$
with $`G.C.D.(A_1^a,\mathrm{},A_n^a,B_1^a,\mathrm{},B_k^a,D_1^a,\mathrm{},D_{2k}^a)=1,`$ for all $`aI.`$ In this way $`\mathrm{\Delta }>0`$ is uniquely determined.
The coefficients of (2.5) satisfy the following properties for each index $`aI`$ :
$`𝐚`$ Either $`_a(𝐢,𝐳,\zeta )=\frac{\mathrm{\Delta }}{\mathrm{\Delta }}z_{\mathrm{}},`$ i.e. $`A_1^a=\mathrm{}=A_n^a=0,`$ $`B_1^a=\mathrm{}\text{ }^{\genfrac{}{}{0pt}{}{\mathrm{}}{^{}}}\mathrm{}=B_k^a=0,`$ $`B_{\mathrm{}}^a=1.`$
$`𝐛`$ Or $`_a(𝐢,𝐳,\zeta )=\frac{\mathrm{\Delta }}{\mathrm{\Delta }}(\zeta _{2\mathrm{}1}+\zeta _2\mathrm{}z_{\mathrm{}}),`$
$`𝐜`$ Or
$$_a(𝐢,𝐳,\zeta )=\frac{_{j=1}^nA_j^a(i_j+1)+_{\mathrm{}=1}^kB_{\mathrm{}}^a(z_{\mathrm{}}\zeta _{2\mathrm{}1}1)}{\mathrm{\Delta }}$$
2) For each fixed index $`1\mathrm{}n,1qk,`$ the following equalities take place:
$$\underset{aI}{}A_{\mathrm{}}^a=0,\underset{aI}{}B_q^a=0.$$
$`(2.7)`$
3) The following relation holds among the linear functions $`_a,`$ $`aI`$:
$$\underset{aI}{}_a(𝐢,𝐳,\zeta )=\zeta _1+\mathrm{}+\zeta _{2k}+2k.$$
In the sequel, especially from $`(3.2)`$ of the next section, we will make use of the notation as follows,
$$(w_1^{(a)},w_2^{(a)},\mathrm{},w_n^{(a)},p_1^{(a)},\mathrm{},p_k^{(a)},q_1^{(a)},\mathrm{},q_{2k}^{(a)})=(\frac{A_1^a}{\mathrm{\Delta }},\mathrm{},\frac{A_n^a}{\mathrm{\Delta }},\frac{B_1^a}{\mathrm{\Delta }},\mathrm{},\frac{B_k^a}{\mathrm{\Delta }},\frac{D_1^a}{\mathrm{\Delta }},\mathrm{},\frac{D_{2k}^a}{\mathrm{\Delta }})_{,1aL}.$$
$`(2.8)`$
In view of the Proposition 2.1, we introduce the subsets of indices $`aI=\{1,2,\mathrm{},L\}`$ as follows.
###### Definition 3
The subset $`I_q^+I`$ (resp. $`I_q^{},I_q^0`$) consists of the indices $`a`$ such that the coefficient $`B_q^a`$ of $`_a(𝐢,𝐳,\zeta )`$ (2.5) is positive (resp. negative, zero).
From this proposition we get the following.
###### Corollary 2.2
The integral $`I_{x^𝐢,\gamma }^\zeta (s)`$ satisfies the hypergeometric system of Horn type as follows:
$$L_q(\vartheta _s,s)I_{x^𝐢,\gamma }^\zeta (s):=\left[P_{q,𝐢}(\vartheta _s,\zeta )s_q^\mathrm{\Delta }Q_{q,𝐢}(\vartheta _s,\zeta )\right]I_{x^𝐢,\gamma }^\zeta (s)=0,1qk$$
$`(2.9)_1`$
with
$$P_{q,𝐢}(\vartheta _s,\zeta )=\underset{aI_q^+}{}\underset{j=0}{\overset{B_q^a1}{}}\left(_a(𝐢,\vartheta _s,\zeta )+j\right),$$
$`(2.9)_2`$
$$Q_{q,𝐢}(\vartheta _s,\zeta )=\underset{\overline{a}I_q^{}}{}\underset{j=0}{\overset{B_q^{\overline{a}}1}{}}\left(_{\overline{a}}(𝐢,\vartheta _s,\zeta )+j\right),$$
$`(2.9)_3`$
where $`I_q^+,I_q^{},1qk`$ are the sets of indices defined in Definition 3.
The degree of two operators $`P_{q,𝐢}(\vartheta _s,\zeta ),`$ $`Q_{q,𝐢}(\vartheta _s,\zeta )`$ are equal. Namely,
$$degP_{q,𝐢}(\vartheta _s,\zeta )=\underset{aI_q^+}{}B_q^a=\underset{\overline{a}I_q^{}}{}B_q^{\overline{a}}=degQ_{q,𝐢}(\vartheta _s,\zeta ).$$
$`(2.10)`$
## 3 Hypothesis by Berglund, Candelas et alii
In this section, we apply the results from §2 on the period integrals to a class of Calabi-Yau varieties $`(1.1)`$ studied in the framework of Landau-Ginzburg vacua theory.
Our aim is to find out sufficient conditions so that the (Mellin transform of) the period integrals on $`X_s`$ can be expressed by means of quasihomogeneous weight data of $`Y.`$ The main theorem of this section is Theorem 3.1.
Before proceeding to the proof of the Theorem 3.1, we write down concretely the matrix $`𝖫`$ in taking $`(1.1),(1.4)`$ into account, and we have the following row vectors of $`𝖫`$:
$$\overline{v}_q^{(\nu )}=(v_{q,1}^{(\nu )},v_{q,2}^{(\nu )},\mathrm{},v_{q,\tau _1}^{(\nu )},\mathrm{},v_{q,n}^{(\nu )},\stackrel{2\nu 2}{\stackrel{}{0,\mathrm{},0}},1,\stackrel{3k2\nu +1}{\stackrel{}{0,\mathrm{},0}}),$$
$`(3.1)`$
$`1\nu k,1q\tau _\nu .`$ In using the notations $`(2.8),`$ $`\mathrm{\S }2`$ for the column vectors of the matrix $`𝖫^1,`$ one can deduce the following system for each $`\nu `$ and $`q`$.
$$\begin{array}{ccc}v_{q,1}^{(\nu )}w_1^{(a^\nu )}+v_{q,2}^{(\nu )}w_2^{(a^\nu )}+\mathrm{}+v_{q,n}^{(\nu )}w_n^{(a^\nu )}& =& 1fora^{\nu 1}+1qa^{\nu 1}+\tau _\nu ,\\ & =& 0\mathrm{otherwise}.\end{array}$$
$`(3.2)`$
$$v_{1,j}^{(1)}p_q^{(1)}+v_{2,j}^{(1)}p_q^{(2)}+\mathrm{}+v_{\tau _1,j}^{(1)}p_q^{(\tau _1)}+v_{1,j}^{(2)}p_q^{(\tau _1+4)}+\mathrm{}+v_{\tau _2,j}^{(2)}p_q^{(\tau _1+\tau _2+3)}+\mathrm{}+v_{\tau _k,j}^{(k)}p_q^{(a^k3)}$$
$`(3.3)`$
$$\begin{array}{ccc}& =& 1forjI^{(q)},\\ & =& 0forjI^{(r)},rq.\end{array}$$
Here we shall remark that the system $`(3.2)`$ for a fixed $`\nu `$ (resp. $`(3.3)`$ for a fixed $`q`$) consists of $`n`$ linear independent equations with respect to unknowns $`\{w_1^{(a^\nu )},\mathrm{},w_n^{(a^\nu )}\},`$ (resp. $`\{p_q^{(j)};jI_\mathrm{\Lambda }\}`$ ). Here we made use of the notation $`I_\mathrm{\Lambda }`$ the set of indices $`\{1,\mathrm{},L\}_{\nu =1}^k\{a^\nu 2,a^\nu 1,a^\nu \}.`$ One sees other necessary conditions on $`w_j^{(a^{\nu 1})}`$,
$$\underset{jI^{(\nu )}}{}w_j^{(a^\nu 1)}=1,$$
$`(3.4)`$
$$\underset{jI^{(\nu )}}{}w_j^{(a^\nu )}=1.$$
$`(3.5)`$
This can be seen from the fact that the product of the $`(a^\nu 1)`$th column of the matrix $`𝖫^1`$ with the $`(a^\nu 1)`$th row of $`𝖫`$ is equal to $`1.`$ On the other hand $`n+2\nu `$th column of $`𝖫`$ with the $`j`$th row ($`1jn+2`$) of $`𝖫^1`$ is equal to $`0`$ which entails $`w_j^{(a^\nu 1)}+w_j^{(a^\nu )}=0.`$ In addition to that, we see
$$\underset{jI^{(\nu )}}{}w_j^{(a^\nu ^{}1)}=\underset{jI^{(\nu )}}{}w_j^{(a^\nu ^{})}=0,$$
for $`\nu ^{}\nu .`$
In this situation we deduce the following system from $`(3.3)`$,
$$\left(\begin{array}{c}\stackrel{}{P}_1\\ \stackrel{}{P}_2\\ \mathrm{}\\ \stackrel{}{P}_k\end{array}\right)𝖫_\mathrm{\Lambda }=^tV^\mathrm{\Lambda },$$
$`(3.6)`$
with $`\stackrel{}{P}_q=(p_q^{(1)},p_q^{(2)},\mathrm{},p_q^{(\tau _1)},\mathrm{},p_q^{(a^k3)}).`$ If we denote by $`\stackrel{ˇ}{𝖫}_j`$ the $`j`$th column vector of the matrix $`𝖫_\mathrm{\Lambda },`$ we have the following equation derived directly from $`(3.6)`$:
$$z_1\stackrel{}{P}_1+z_2\stackrel{}{P}_1+\mathrm{}+z_k\stackrel{}{P}_k,\stackrel{ˇ}{𝖫}_j+z_q=0\mathrm{if}jI^{(q)}.$$
$`(3.7)`$
In view of $`(1.12)`$, we get
$$({}_{}{}^{T}𝖫_{\mathrm{\Lambda }}^{})^1^TV^\mathrm{\Lambda }=({}_{}{}^{t}𝖫_{\mathrm{\Lambda }}^{})^1V^\mathrm{\Lambda },$$
$`(3.8)`$
that yields
$$\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{g}_{}^{(1)}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{g}_{}^{(k)}\end{array}\right)^TG^1({}_{}{}^{T}𝖫_{\mathrm{\Lambda }}^{})^1^TV^\mathrm{\Lambda }$$
$`(3.9)`$
$$=\left(\begin{array}{cccc}p_1^{(1)}+\mathrm{}+p_1^{(\stackrel{~}{\tau }_1)}& p_2^{(1)}+\mathrm{}+p_2^{(\stackrel{~}{\tau }_1)}& \mathrm{}& p_k^{(1)}+\mathrm{}+p_k^{(\stackrel{~}{\tau }_1)}\\ p_1^{(\stackrel{~}{a}_1+1)}+\mathrm{}+p_1^{(\stackrel{~}{a}_1+\stackrel{~}{\tau }_2)}& p_2^{(\stackrel{~}{a}_1+1)}+\mathrm{}+p_2^{(\stackrel{~}{a}_1+\stackrel{~}{\tau }_2)}& \mathrm{}& p_k^{(\stackrel{~}{a}_1+1)}+\mathrm{}+p_k^{(\stackrel{~}{a}_1+\stackrel{~}{\tau }_2)}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ p_1^{(\stackrel{~}{a}_{k1}+1)}+\mathrm{}+p_1^{(n+3k3)}& p_2^{(\stackrel{~}{a}_{k1}+1)}+\mathrm{}+p_2^{(n+3k3)}& \mathrm{}& p_k^{(\stackrel{~}{a}_{k1}+1)}+\mathrm{}+p_k^{(n+3k3)}\end{array}\right),$$
where we used the notation $`\stackrel{~}{a}_q:=_{\nu =1}^q(\stackrel{~}{\tau }_\nu +3).`$ Let us introduce a permutation matrix $`\nu ^{}SL(k,𝐙)`$ whose $`j`$th column equals to
$${}_{}{}^{t}(\stackrel{\nu (j)1}{\stackrel{}{0,\mathrm{},0}}\text{ }^{\genfrac{}{}{0pt}{}{\nu (j)}{^{}}},1,\stackrel{n\nu (j)}{\stackrel{}{0,\mathrm{},0}}).$$
The last expression $`(3.9)`$ turns out to be $`\nu ^{}`$. This fact can be seen by the relation $`𝖫𝖫^\mathrm{𝟣}=\mathrm{𝗂𝖽}_𝖫`$ and the Proposition 2.1, 1), $`𝐚,𝐛,𝐜`$.
We can define a set of indices $`{}_{}{}^{T}I_{\mathrm{\Lambda }}^{}`$ for $`(1.4)^T`$ analogous to $`I_\mathrm{\Lambda }.`$ We have the indices of $`I_\mathrm{\Lambda }`$, $`1=j_1<\mathrm{}<j_n=L4`$ and those of $`{}_{}{}^{T}I_{\mathrm{\Lambda }}^{},`$ $`1=i_1<\mathrm{}<i_n=L4.`$ In this situation we consider two conditions on the CI $`(1.1),`$
$$\left(\begin{array}{c}{}_{}{}^{T}\xi _{i_1}^{}(z)\\ \mathrm{}\\ {}_{}{}^{T}\xi _{i_n}^{}(z)\end{array}\right)=GV^\mathrm{\Lambda }\left(\begin{array}{c}{}_{}{}^{T}\xi _{}^{(1)}(z)\\ \mathrm{}\\ {}_{}{}^{T}\xi _{}^{(k)}(z)\end{array}\right),$$
$`(3.10)`$
for linear functions $`(^T\xi ^{(1)}(z),\mathrm{},^T\xi ^{(k)}(z))`$ and
$$\left(\begin{array}{c}\xi _{j_1}(z)\\ \mathrm{}\\ \xi _{j_n}(z)\end{array}\right)=^TG^TV^\mathrm{\Lambda }\left(\begin{array}{c}\xi ^{(1)}(z)\\ \mathrm{}\\ \xi ^{(k)}(z)\end{array}\right),$$
$`(3.10)^T`$
for possibly another $`k`$tuple of linear functions $`(\xi ^{(1)}(z),\mathrm{},\xi ^{(k)}(z)).`$
Due to the condition $`(1.11)`$ on the set of indices $`{}_{}{}^{T}I_{\mathrm{\Lambda }}^{}`$, we have for some permutation matrices $`\rho ,^T\rho SL(n,𝐙)`$,
$$\rho V^\mathrm{\Lambda }=G^1({}_{}{}^{t}(\stackrel{}{g}^{(1)}),\mathrm{},^t(\stackrel{}{g}^{(k)})),$$
$${}_{}{}^{T}\rho ^TV^\mathrm{\Lambda }=^TG^1({}_{}{}^{t}(^T\stackrel{}{g}^{(1)}),\mathrm{},^t(^T\stackrel{}{g}^{(k)})).$$
On these $`\rho `$ and $`{}_{}{}^{T}\rho ,`$ we impose the following conditions,
$$G\rho =^t(G\rho ),$$
$`(3.11)`$
$${}_{}{}^{T}G^T\rho =^t(^TG^T\rho ).$$
$`(3.11)^T`$
Under these conditions we formulate the following Theorem that verifies an hypothesis proposed by , §3.3 under certain conditions.
###### Theorem 3.1
The Mellin transform $`M_{0,\gamma }^0(z)`$ of the period integral $`I_{x^0,\gamma }^0(s)`$ for the Calabi-Yau CI $`(1.1)`$ has the following form up to a $`\mathrm{\Delta }`$ periodic function in the sense of Definition 2, if it satisfies the conditions $`(3.10),`$$`(3.10)^T,`$$`(3.11),`$ $`(3.11)^T,`$
$$M_{0,\gamma }^0(z(\xi ))=\frac{_{\nu =1}^k_{j=1}^{\stackrel{~}{\tau }_\nu }\mathrm{\Gamma }(^Tg_j^{(\nu )}\xi ^{(\nu )})}{_{q=1}^k\mathrm{\Gamma }(_{\nu =1}^k{}_{}{}^{T}Q_{q}^{(\nu )}\xi ^{(\nu )})}.$$
$`(3.12)`$
Here $`z(\xi )=(z_1(\xi ),\mathrm{},z_k(\xi ))`$ is a $`k`$tuple of linear functions in variables $`\xi =(\xi ^{(1)},\mathrm{},\xi ^{(k)})`$ defined by the relation $`(3.10)^T.`$
In a symmetric way the Mellin transform $`M_{0,^T\gamma }^0(z(^T\xi ))`$ for the Calabi-Yau CI, $`(1.1)^T`$ admits an expression as follows up to a $`\mathrm{\Delta }`$ periodic function,
$$M_{0,^T\gamma }^0(z(\xi ))=\frac{_{\nu =1}^k_{j=1}^{\tau _\nu }\mathrm{\Gamma }(g_j^{(\nu )}{}_{}{}^{T}\xi _{}^{(\nu )})}{_{q=1}^k\mathrm{\Gamma }(_{\nu =1}^kQ_q^{(\nu )}{}_{}{}^{T}\xi _{}^{(\nu )})}.$$
$`(3.12)^T`$
The functions $`{}_{}{}^{T}\xi _{}^{(\nu )}`$ are defined by the relation $`(3.10).`$
To prove the Theorem we prepare the following lemma.
###### Lemma 3.2
Under the conditions imposed on $`(1.1)`$ in the theorem 3.1, the Mellin transform $`M_{0,\gamma }^0(z)`$ of the period integral $`I_{x^0,\gamma }^0(s)`$ for the CI $`(1.1)`$ admits up to a $`\mathrm{\Delta }`$periodic function (in the sense of Definition 2) an expression as follows,
$$M_{0,\gamma }^0(z)=\underset{i=1}{\overset{k}{}}\mathrm{\Gamma }(z_i)\underset{jI_\mathrm{\Lambda }}{}\mathrm{\Gamma }(\xi _j(z)).$$
Proof It is necessary to show that special solutions of the system $`(1.5)`$ satisfy,
$$\xi _{a^\nu 2}(z)=\xi _{a^\nu 1}(z)=z_\nu ,\xi _{a^\nu }(z)=1z_\nu .$$
$`(3.13)`$
To see this, we remark that the system below is a direct consequence of the relation $`\mathrm{𝖫𝖫}^1=id_L,`$
$$\underset{i=1}{\overset{n}{}}v_{q,i}^{(\nu )}w_i^{(a^\nu )}+q_{2\mu 1}^{(a^\nu )}=0,a^{\mu 1}+1qa^{\mu 1}+\tau _\mu ,\mu [1,k].$$
$`(3.14)`$
$$\underset{i=1}{\overset{n}{}}w_i^{(a^\nu )}=\underset{i=1}{\overset{n}{}}w_i^{(a^\nu 1)}=1.$$
The last equality can be deduced from $`(3.4)`$ and $`(3.5).`$ Let $`\mu (j)[1,k]`$ be an index such that $`j`$ belongs to $`I^{(\mu (j))}.`$ Then we have
$$\underset{iI_\mathrm{\Lambda }}{}v_{i,j}^{(\mu (j))}p_q^{(i)}+p_q^{(a^{\mu (j)1})}=0,$$
$`(3.15)`$
for $`q[1,k].`$
We get the following relation,
$$q_{2\nu 1}^{(a^\nu )}=q_{2\nu }^{(a^\nu )}=1,$$
$`(3.16)`$
$$q_{2\nu 1}^{(a^\nu ^{})}=q_{2\nu }^{(a^\nu ^{})}=0,\mathrm{for}\nu ^{}\nu .$$
$`(3.17)`$
The equality $`q_{2\nu }^{(a^\nu ^{})}=0`$ can be derived from the fact that the product of the $`a^{\nu 1}`$th column of $`𝖫^1`$ with the $`a^\nu ^{}`$th row of $`𝖫`$ is equal to $`0.`$ On the other hand, the product of the $`(n+2\nu )`$th column of $`𝖫`$ with the $`a^\nu ^{}`$th row of $`𝖫^1`$ is equal to $`q_{2\nu 1}^{(a^\nu ^{})}+q_{2\nu }^{(a^\nu ^{})}=0.`$ This proves $`(3.17).`$ The product of the $`a^\nu `$th column of $`𝖫^1`$ with the $`a^\nu ^{}`$th row of $`𝖫`$ is equal to $`q_{2\nu 1}^{(a^\nu ^{})}+p_\nu ^{(a^\nu ^{})}=0.`$ One deduces the equality $`p_\nu ^{(a^\nu ^{})}=0`$ for $`\nu \nu ^{}.`$ The product of $`(a^\nu 1)`$th column of $`𝖫^1`$ with the $`(a^\nu ^{}2)`$th row of $`𝖫`$ is equal to $`q_{2\nu ^{}1}^{(a^\nu 1)}+p_\nu ^{}^{(a^\nu 1)}=0.`$ As a consequence, we have
$$p_\nu ^{}^{(a^\nu 1)}=0\mathrm{for}\nu \nu ^{},$$
$`(3.18)`$
$$p_\nu ^{(a^\nu 1)}=q_{2\nu 1}^{(a^\nu 1)}=1,p_\nu ^{(a^\nu )}=1,$$
because $`p_\nu ^{(a^\nu )}+p_\nu ^{(a^\nu 1)}=0.`$ Additionally we see that
$$q_r^{(a^\nu 1)}=0,r2\nu 1.$$
$`(3.19)`$
In summary, by virtue of $`(3.14),(3.16),(3.18),`$
$$\xi _{a^\nu }(z)=\underset{i=1}{\overset{n}{}}w_i^{(a^\nu )}+\underset{r=1}{\overset{2k}{}}q_r^{(a^\nu )}+\underset{q=1}{\overset{k}{}}p_q^{(a^\nu )}z_q=1+1+1z_\nu =1z_\nu .$$
On the other hand $`(3.14),(3.19),(3.18),`$ yield
$$\xi _{a^\nu 1}(z)=1+01+z_\nu =z_\nu .$$
As for the function $`\xi _{a^\nu 2}(z)`$ it is easy to see that all elements of the $`(a^\nu 2)`$th column of $`𝖫^1`$ consist of zeros except the $`(n+2k+\nu )`$th element which is equal to $`1.`$
We thus have the equality,
$$M_{0,\gamma }^0(z)=\underset{i=1}{\overset{k}{}}\mathrm{\Gamma }(z_i)^2\mathrm{\Gamma }(1z_i)\underset{jI_\mathrm{\Lambda }}{}\mathrm{\Gamma }(\xi _j(z))=\underset{i=1}{\overset{k}{}}\frac{\pi }{sin\pi z_i}\mathrm{\Gamma }(z_i)\underset{jI_\mathrm{\Lambda }}{}\mathrm{\Gamma }(\xi _j(z)).$$
Q.E.D.
Proof of the Theorem 3.1 Our main task is to show the following relation,
$$\nu ^{}\left(\begin{array}{c}1z_1\\ \mathrm{}\\ 1z_k\end{array}\right)=^T\widehat{Q}\left(\begin{array}{c}\xi ^{(1)}(z)\\ \mathrm{}\\ \xi ^{(k)}(z)\end{array}\right),$$
$`(3.20)`$
for the permutation matrix $`\nu ^{}SL(k,𝐙)`$ introduced just after the formula $`(3.9)`$. To do this, first of all we modify the relation,
$${}_{}{}^{T}\widehat{Q}=^t(^TV^\mathrm{\Lambda })({}_{}{}^{t}(^T\stackrel{}{g}^{(1)}),\mathrm{},^t(^T\stackrel{}{g}^{(k)}))=\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{g}_{}^{(1)}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{g}_{}^{(k)}\end{array}\right)^TG^1^t(^T\rho )^1({}_{}{}^{t}(^T\stackrel{}{g}^{(1)}),\mathrm{},^t(^T\stackrel{}{g}^{(k)})),$$
$`(3.21)`$
that can be derived from $`(1.8)^T.`$ Here we remark that after the definition of $`{}_{}{}^{T}\rho SL(n,𝐙)`$ intrduced just before $`(3.11)`$ the following relation holds,
$${}_{}{}^{T}V_{}^{\mathrm{\Lambda }}=^T\rho ^1^TG^1({}_{}{}^{t}(^T\stackrel{}{g}^{(1)}),\mathrm{},^t(^T\stackrel{}{g}^{(k)})).$$
$`(3.22)`$
From this relation and $`(3.11)^T`$ we see that
$$\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{g}_{}^{(1)}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{g}_{}^{(k)}\end{array}\right)^T\rho ^1^TG^1({}_{}{}^{t}(^T\stackrel{}{g}^{(1)}),\mathrm{},^t(^T\stackrel{}{g}^{(k)}))=\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{g}_{}^{(1)}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{g}_{}^{(k)}\end{array}\right)^TV^\mathrm{\Lambda }.$$
$`(3.23)`$
By virtue of the condition $`(3.10)^T`$, the following equality holds
$$\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{g}_{}^{(1)}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{g}_{}^{(k)}\end{array}\right)^TV^\mathrm{\Lambda }\left(\begin{array}{c}\xi ^{(1)}(z)\\ \mathrm{}\\ \xi ^{(k)}(z)\end{array}\right)=\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{g}_{}^{(1)}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{g}_{}^{(k)}\end{array}\right)^TG^1\left(\begin{array}{c}\xi _{j_1}(z)\\ \mathrm{}\\ \xi _{j_n}(z)\end{array}\right).$$
$`(3.24)`$
The combination of $`(3.6)`$ and $`(3.8)`$ entails that the expression $`(3.24)`$ is equal to
$$\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{g}_{}^{(1)}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{g}_{}^{(k)}\end{array}\right)^TG^1^t(𝖫_\mathrm{\Lambda })^1V^\mathrm{\Lambda }\left(\begin{array}{c}1z_1\\ \mathrm{}\\ 1z_k\end{array}\right)=\left(\begin{array}{c}{}_{}{}^{T}\stackrel{}{g}_{}^{(1)}\\ \mathrm{}\\ {}_{}{}^{T}\stackrel{}{g}_{}^{(k)}\end{array}\right)^TG^1(^T𝖫_\mathrm{\Lambda })^1{}_{}{}^{T}V_{}^{\mathrm{\Lambda }}\left(\begin{array}{c}1z_1\\ \mathrm{}\\ 1z_k\end{array}\right).$$
From $`(3.9)`$ it follows that the last expression equals to
$$\nu ^{}\left(\begin{array}{c}1z_1\\ \mathrm{}\\ 1z_k\end{array}\right).$$
This means $`(3.20)`$ and consequently,
$$\underset{\nu =1}{\overset{k}{}}{}_{}{}^{T}Q_{q}^{(\nu )}\xi ^{(\nu )}(z)=1z_q,\mathrm{\hspace{0.33em}\hspace{0.33em}1}qk.$$
From the last equality we can directly derive the relation to be proved
$$\underset{q=1}{\overset{k}{}}\mathrm{\Gamma }\left(z_q\right)=\underset{q=1}{\overset{k}{}}\mathrm{\Gamma }\left(1\underset{\nu =1}{\overset{k}{}}{}_{}{}^{T}Q_{q}^{\left(\nu \right)}\xi ^{\left(\nu \right)}\left(z\right)\right)=\frac{\pi }{_{q=1}^k\mathrm{\Gamma }\left(_{\nu =1}^k{}_{}{}^{T}Q_{q}^{\left(\nu \right)}\xi ^{\left(\nu \right)}\left(z\right)\right)sin\left(\pi _{\nu =1}^k{}_{}{}^{T}Q_{q}^{\left(\nu \right)}\xi ^{\left(\nu \right)}\left(z\right)\right)}.$$
On the other hand the lemma 3.2 and the condition $`(3.10)^T`$ gives us
$$\xi _{a^{\nu 1}+j}=^Tg_j^{(\nu )}\xi ^{(\nu )},j[1,\tau _\nu ]$$
that means
$$\underset{jI_\mathrm{\Lambda }}{}\mathrm{\Gamma }(\xi _j(z))=\underset{\nu =1}{\overset{k}{}}\underset{j=1}{\overset{\tau _\nu }{}}\mathrm{\Gamma }(^Tg_j^{(\nu )}\xi ^{(\nu )}),$$
which proves $`(3.12)`$. The formula $`(3.12)^T`$ can be proven in a parallel way. Q.E.D.
## 4 Duality between monodromy data and Poincaré polynomials
In connection with the mirror symmetry, we consider the structural algebra of the CI $`(1.1)`$ of dimension $`nk`$ denoted by $`X`$ $`=X_\mathrm{𝟏}`$,
$$A_X:=\frac{𝐂[x]}{(f_1+f_21,\mathrm{}f_{2k1}+f_{2k}1)𝐂[x]},$$
and a natural filtration on it,
$$A_X^j:=𝐂\{x^\alpha A_X;\alpha ,\stackrel{}{g}^{(1)}=j_1,\mathrm{},\alpha ,\stackrel{}{g}^{(k)}=j_k\},$$
with the Poincaré polynomial,
$$P_{A_X}(\lambda )=\underset{j𝐙_0}{}dim(A_X^j)\lambda _1^{j_1}\mathrm{}\lambda _k^{j_k}.$$
In an analogous way, we define corresponding notions of the CI $`Y`$ defined by $`(1.1)^T,`$
$$A_Y:=\frac{𝐂[x]}{(^Tf_1+^Tf_21,\mathrm{}^Tf_{2k1}+^Tf_{2k}1)𝐂[x]},$$
$$A_Y^j:=\{x^\alpha A_Y;\alpha ,^T\stackrel{}{g}^{(1)}=j_1,\mathrm{},\alpha ,^T\stackrel{}{g}^{(k)}=j_k\},$$
$$P_{A_Y}(\lambda )=\underset{j𝐙_0}{}dim(A_Y^j)\lambda ^j.$$
In this situation the classical result due to gives us,
$$P_{A_X}(\lambda )=\frac{_{q=1}^k_{\nu =1}^k(1\lambda _\nu ^{Q_q^{(\nu )}})}{_{\nu =1}^k_{j=1}^{\tau _\nu }(1\lambda _\nu ^{g_j^{(\nu )}})},P_{A_Y}(\lambda )=\frac{_{q=1}^k_{\nu =1}^k(1\lambda _\nu ^{{}_{}{}^{T}Q_{q}^{(\nu )}})}{_{\nu =1}^k_{j=1}^{\stackrel{~}{\tau }_\nu }(1\lambda _\nu ^{{}_{}{}^{T}g_{j}^{(\nu )}})}.$$
$`(4.1)`$
Further we introduce the variables $`(t_1,\mathrm{},t_k)𝐓^k`$ such that
$$\underset{\nu =1}{\overset{k}{}}s_\nu ^{{}_{}{}^{T}Q_{\nu }^{(q)}}=t_q,\mathrm{\hspace{0.33em}\hspace{0.33em}1}qk.$$
If we assume that $`rank^T\widehat{Q}=k`$, this equation is always solvable with respect to the variables $`s=s(t).`$ We consider the Mellin inverse transform of $`M_{0,\gamma }^0(z(\xi ))`$ associated to the CI (1.1),
$$U_\alpha (s)=_{\stackrel{ˇ}{\mathrm{\Pi }}_\alpha }\frac{_{\nu =1}^k_{j=1}^{\tau _\nu }\mathrm{\Gamma }(^Tg_j^{(\nu )}\xi ^{(\nu )}(z))}{_{\nu =1}^k\mathrm{\Gamma }(_{q=1}^k{}_{}{}^{T}Q_{q}^{(\nu )}\xi ^{(\nu )}(z))}s^{𝐳+\mathrm{𝟏}}𝑑𝐳,$$
where $`\stackrel{ˇ}{\mathrm{\Pi }}_\alpha 𝐓^k`$ is a cycle avoiding the singular loci of the integrand. It is easy to check
$$U_\alpha (s(t))=det(^T\widehat{Q})^1_{{}_{}{}^{T}\widehat{Q}_{}^{1}(\stackrel{ˇ}{\mathrm{\Pi }}_\alpha )}\frac{_{\nu =1}^k_{j=1}^{\stackrel{~}{\tau }_\nu }\mathrm{\Gamma }(^Tg_j^{(\nu )}\xi ^{(\nu )})}{_{\nu =1}^k\mathrm{\Gamma }(_{q=1}^k{}_{}{}^{T}Q_{q}^{(\nu )}\xi ^{(\nu )})}t_1^{\xi ^{(1)}}\mathrm{}t_k^{\xi ^{(k)}}d\xi ^{(1)}\mathrm{}d\xi ^{(k)}.$$
The following system of differential equations annihilates the inverse Mellin transform $`U_\alpha (s(t))`$,
$$L_\nu (t_\nu ,\vartheta _t)U_\alpha (s(t))=0,\mathrm{\hspace{0.33em}\hspace{0.33em}1}\nu k,$$
where
$$L_\nu (t_\nu ,\vartheta _t)=\left(\underset{j=1}{\overset{\stackrel{~}{\tau }_\nu }{}}\underset{r=0}{\overset{{}_{}{}^{T}g_{j}^{(\nu )}1}{}}(^Tg_j^{(\nu )}\vartheta _{t_\nu }+r)t_\nu \underset{q=1}{\overset{k}{}}\underset{r=0}{\overset{{}_{}{}^{T}Q_{q}^{(\nu )}1}{}}(\underset{\mu =1}{\overset{k}{}}{}_{}{}^{T}Q_{q}^{(\mu )}\vartheta _{t_\mu }r)\right),\mathrm{\hspace{0.33em}\hspace{0.33em}1}\nu k.$$
$`(4.2)`$
We denote by $`\chi _\nu `$ the degree of the operator $`L_\nu (t,\vartheta _t):`$ $`\chi _\nu =_{q=1}^k{}_{}{}^{T}Q_{q}^{(\nu )}=_{j=1}^{\stackrel{~}{\tau }_\nu }{}_{}{}^{T}g_{j}^{(\nu )}=^Tg_{\stackrel{~}{\tau }_\nu +1}^{(\nu )}`$ that has already been introduced in $`(1.3)^T.`$
Here we define the restriction of the operator $`L_\nu (t,\vartheta _t)`$ onto the torus $`𝐓=\{t𝐂^k;t_i=0,i\nu \}\{t_\nu =0\}`$ as follows,
$$\stackrel{~}{L}_\nu (t_\nu ,\vartheta _{t_\nu }):=(\underset{j=1}{\overset{\stackrel{~}{\tau }_\nu }{}}\underset{r=0}{\overset{{}_{}{}^{T}g_{j}^{(\nu )}1}{}}(^Tg_j^{(\nu )}\vartheta _{t_\nu }+r)t_\nu \underset{q=1}{\overset{k}{}}\underset{r=0}{\overset{{}_{}{}^{T}Q_{q}^{(\nu )}1}{}}(^TQ_q^{(\nu )}\vartheta _{t_\nu }r)).$$
$`(4.2)^{}`$
On the $`\chi _q`$ dimensional solution space of the operator $`\stackrel{~}{L}_q(t_q,\vartheta _{t_q})`$, we consider the monodromy $`M_q^{(0)}GL(\chi _q,𝐂)`$ (resp. $`M_q^{(\mathrm{})}GL(\chi _q,𝐂)`$) around the point $`t_q=0`$ (resp. $`t_q=\mathrm{}`$). Then we have the following characteristic polynomials of the monodromy that can be easily calculated from the expression $`\stackrel{~}{L}_\nu (t_\nu ,\vartheta _{t_\nu })`$,
$$det(M_q^{(\mathrm{})}\lambda _qid_{\chi _q})=\underset{\nu =1}{\overset{k}{}}(1\lambda _q^{{}_{}{}^{T}Q_{\nu }^{(q)}}),det(M_q^{(0)}\lambda _qid_{\chi _q})=\underset{j=1}{\overset{\tau _q}{}}(1\lambda _q^{{}_{}{}^{T}g_{j}^{(q)}}).$$
As a consequence the rational function defined by
$$M_X(\lambda _1,\mathrm{}\lambda _k):=\underset{q=1}{\overset{k}{}}\frac{det(M_q^{(\mathrm{})}\lambda _qid_{\chi _q})}{det(M_q^{(0)}\lambda _qid_{\chi _q})}$$
has a form
$$M_X(\lambda _1,\mathrm{}\lambda _k)=\frac{_{q=1}^k_{\nu =1}^k(1\lambda _q^{{}_{}{}^{T}Q_{\nu }^{(q)}})}{_{q=1}^k_{j=1}^{\stackrel{~}{\tau }_q}(1\lambda _q^{{}_{}{}^{T}g_{j}^{(q)}})}.$$
$`(4.3)`$
For the rational function $`M_Y(\lambda _1,\mathrm{}\lambda _k)`$ defined in a parallel way to the function $`M_X(\lambda _1,\mathrm{}\lambda _k),`$ we have
$$M_Y(\lambda _1,\mathrm{}\lambda _k)=\frac{_{q=1}^k_{\nu =1}^k(1\lambda _q^{Q_\nu ^{(q)}})}{_{q=1}^k_{j=1}^{\tau _q}(1\lambda _q^{g_j^{(q)}})}.$$
Let $`\overline{Y}`$ be the compactification of CI $`Y𝐓^n`$ in the product of quasihomogeneous projective spaces $`𝐏:=𝐏_{(^Tg_1^{(1)},\mathrm{},^Tg_{\tau _1+1}^{(1)})}^{(\tau _1)}\times \mathrm{}\times 𝐏_{(^Tg_1^{(k)},\mathrm{},^Tg_{\tau _k+1}^{(k)})}^{(\tau _k)}.`$ The coherent sheaf on $`𝐏`$ $`𝒪_𝐏(\zeta ),`$ $`\zeta =(\zeta _1,\mathrm{},\zeta _k)`$ is defined by sections on the open set $`U_I=\{x𝐂^n;x_i0,iI\}`$ that are given by
$$\mathrm{\Gamma }(U_I,𝒪_𝐏(\zeta )):=𝐂\{x^\alpha ;\alpha =(\alpha _1,\mathrm{},\alpha _n)𝐙^n,\alpha _i0foriI,^Tg^{(q)},\alpha =\zeta _q,1qk\}.$$
We define the coherent sheaf $`𝒪_{\overline{Y}}(\zeta )`$ by sections on the open set $`U_IY,`$
$$\mathrm{\Gamma }(U_I,𝒪_{\overline{Y}}\left(\zeta \right)):=𝐂\{x^\alpha A_Y;\alpha =(\alpha _1,\mathrm{},\alpha _n)𝐙^n,\alpha _i0foriI,^Tg^{\left(q\right)},\alpha =\zeta _q,1qk\}.$$
We introduce the Euler characteristic for this sheaf,
$$\chi (𝒪_{\overline{Y}}(\zeta )):=\underset{i=0}{\overset{nk}{}}(1)^idimH^i(𝒪_{\overline{Y}}(\zeta )).$$
After , the Poincaré polynomial of the Euler characteristics,
$$P𝒪_{\overline{Y}}(t_1,\mathrm{},t_k):=\underset{\zeta (𝐙_0)^k}{}\chi (𝒪_{\overline{Y}}(\zeta ))t_1^{\zeta _1}\mathrm{}t_k^{\zeta _k}$$
admits an expression as follows,
$$P𝒪_{\overline{Y}}(t_1,\mathrm{},t_k)=\frac{_{q=1}^k_{\nu =1}^k(1t_q^{{}_{}{}^{T}Q_{\nu }^{(q)}})}{_{q=1}^k_{j=1}^{\stackrel{~}{\tau }_\nu }(1t_q^{{}_{}{}^{T}g_{j}^{(q)}})}.$$
$`(4.4)`$
For the sheaf $`𝒪_{\overline{X}}(\zeta )`$ defined analogously to $`𝒪_{\overline{Y}}(\zeta )`$, we consider the Poincaré polynomial of the Euler characteristics
$$P𝒪_{\overline{X}}(t_1,\mathrm{},t_k):=\underset{\zeta (𝐙_0)^k}{}\chi (𝒪_{\overline{X}}(\zeta ))t_1^{\zeta _1}\mathrm{}t_k^{\zeta _k}.$$
We have the following expression after and ,
$$P𝒪_{\overline{X}}(t_1,\mathrm{},t_k)=\frac{_{q=1}^k_{\nu =1}^k(1t_q^{Q_\nu ^{(q)}})}{_{q=1}^k_{j=1}^{\tau _\nu }(1t_q^{g_j^{(q)}})}.$$
If we compare $`(4.1),`$ $`(4.3)`$ and $`(4.4)`$ we get the following statement.
###### Theorem 4.1
For the Calabi-Yau CI’s $`X`$ and $`Y`$ defined by $`(1.1)`$ and $`(1.1)^T`$ we have the following relations,
$$M_X(\lambda _1,\mathrm{},\lambda _k)=P𝒪_{\overline{Y}}(\lambda _1,\mathrm{},\lambda _k)=P_{A_Y}(\lambda _1,\mathrm{},\lambda _k),$$
$$M_Y(\lambda _1,\mathrm{},\lambda _k)=P𝒪_{\overline{X}}(\lambda _1,\mathrm{},\lambda _k)=P_{A_X}(\lambda _1,\mathrm{},\lambda _k),$$
if they satisfy sufficient conditions of Theorem 3.1.
In this situation one can easily derive the existence of the following symmetry from the Theorem 3.1.
###### Theorem 4.2
Assume that $`X`$, $`(1.1)`$ and $`Y`$, $`(1.1)^T`$ satisfy all conditions imposed on them in Theorem 3.1, then there is a symmetry between geometric symmetry and quantum symmetry as follows.
$$𝒬_X\underset{q=1}{\overset{k}{}}𝐙_{\overline{𝐐}^{(𝐪)}}𝒢_Y,$$
$$𝒬_Y\underset{q=1}{\overset{k}{}}𝐙_{{}_{}{}^{𝐓}\overline{𝐐}_{}^{(𝐪)}}𝒢_X.$$
Proof The isomorphism $`𝒬_X_{q=1}^k𝐙_{\overline{𝐐}^{(𝐪)}}`$ and its analogy on $`Y_s`$ is clear from the quasihomoeneiety of the systems $`(1.1)`$ and $`(1.1)^T`$. The existence of the cyclic action $`_{q=1}^k𝐙_{{}_{}{}^{𝐓}\overline{𝐐}_{}^{(𝐪)}}`$ on $`X_s`$ can be read off from the monodromy actions of $`M_q^{(\mathrm{})},1qk`$ on the solutions to the system $`(4.2)`$ that are defined along vanishing cycles on $`X_s.`$ Here we remark that the monodromy group is a subgroup of $`Aut(X_s)`$ in view of the integer power transform from $`(s_1,\mathrm{},s_k)`$ to $`(t_1,\mathrm{},t_k)`$ defined just after $`(4.1)`$. We remark here that the monodromy actions of $`M_q^{(0)},1qk`$ consist part of the cyclic action $`_{q=1}^k𝐙_{{}_{}{}^{𝐓}\overline{𝐐}_{}^{(𝐪)}}`$ because every $`{}_{}{}^{T}g_{j}^{(q)},`$ $`j=1,\mathrm{},\stackrel{~}{\tau }_q`$ divides $`{}_{}{}^{𝐓}\overline{𝐐}_{}^{(𝐪)}`$. As for the monodromy of the solution to the system $`(4.2)^{}`$ an analogous argument holds. Q.E.D.
###### Remark 1
As we have not calculated the global monodromy of the integrals $`U_\alpha (s(t))`$, $`(4.3)`$, it is not proper to talk about the existence of a mirror symmetry between $`X`$ and $`Y.`$ We gathered, however, the monodromy data which correspond to certain limit of the values $`t_{\mathrm{}}0,\mathrm{}\nu .`$ Roughly speaking, this procedure can be interpreted as the selection of a special face of the Newton polyhedron $`\mathrm{\Delta }(F(x,s,y))`$ (1.4) for the calculus of the monodromy data.
In the article , authors studied the period integral of $`X`$ at certain limit of the parameter values and they deduced informations on the analogies Gromov-Witten invariants of $`Y`$. They call this duality “local mirror symmetry”. It is probable that our Theorem 4.1 is one of numerous aspects of the local mirror symmetry.
## 5 Dual nef-partition interpretation
In this section we show that under certain condition the transposition mirror construction corresponds to the notion of the dual nef-partion due to L.Borisov .
First of all we consider the following set of vectors defined after $`(1.1)`$, $`(1.2).`$
$$\begin{array}{ccc}\stackrel{}{v}_1^{(q)}\stackrel{}{v}_{\tau _q+1}^{(q)}& & \\ \mathrm{}& & 1qk\\ \stackrel{}{v}_{\tau _q}^{(q)}\stackrel{}{v}_{\tau _q+1}^{(q)}& & \end{array}$$
$`(5.1)`$
This gives rise to a partition of $`I_\mathrm{\Lambda }`$ with $`\mathrm{}I_\mathrm{\Lambda }=n.`$ We set
$$\mathrm{\Delta }_q:=\mathrm{convex}\mathrm{hull}\mathrm{of}(\{0\}\underset{j=1}{\overset{\tau _q}{}}\{\stackrel{}{v}_j^{(q)}\stackrel{}{v}_{\tau _q+1}^{(q)}\}).$$
$`(5.2)`$
Further on we impose the following condition on the Minkowski sum of $`\mathrm{\Delta }_q`$,
$$dim(\mathrm{\Delta }_1+\mathrm{}+\mathrm{\Delta }_k)=nk.$$
$`(5.3)`$
We introduce the $`(nk)`$ dimensional integral lattice
$$V_𝐙=\{x𝐙^n;x,\stackrel{}{g}^{(q)}=0,1qk\}.$$
After this notation each $`\mathrm{\Delta }_q`$ is located on $`V_𝐑=V_𝐙\times 𝐑𝐑^{nk}.`$ Consequently $`\mathrm{\Delta }_1+\mathrm{}+\mathrm{\Delta }_kV_𝐑.`$ There exists a piecewise linear function $`\varphi _q`$ such that $`\varphi _q(y_1+y_2)\varphi _q(y_1)+\varphi _q(y_2).`$ Namely we shall define it as follows,
$$\varphi _q(y)=min_{x\mathrm{\Delta }_q}x,y.$$
$`(5.4)`$
We construct a set of polyhedra dual to the set $`\{\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_k\},`$
$$\mathrm{\Delta }_q^{}=\{\stackrel{}{m}𝐑^n;(5.5)_{\mathrm{}},(5.5)_{\mathrm{}}^{}\},\mathrm{\hspace{0.33em}\hspace{0.33em}1}\mathrm{}k.$$
$`(5.5)_0`$
The conditions $`(5.5)_{\mathrm{}},(5.5)_{\mathrm{}}^{}`$ look like the following,
$$\stackrel{}{m},\stackrel{}{v}_j^{(\mathrm{})}\stackrel{}{v}_{\tau _{\mathrm{}}+1}^{(\mathrm{})}1,1j\tau _{\mathrm{}}.$$
$`(5.5)_{\mathrm{}}`$
$$\stackrel{}{m},\stackrel{}{v}_i^{(q)}\stackrel{}{v}_{\tau _q+1}^{(q)}0,1i\tau _q,q\mathrm{}.$$
$`(5.5)_{\mathrm{}}`$
We can construct the set $`\{\mathrm{\Delta }_1^{},\mathrm{},\mathrm{\Delta }_k^{}\}`$ by means of the vertices vectors $`\stackrel{}{m}_1^{(\mathrm{})},\mathrm{},\stackrel{}{m}_\tau _{\mathrm{}}^{(\mathrm{})},1\mathrm{}k`$ which are defined by the following set of equalities and inequalities.
$$\stackrel{}{m}_r^{(\mathrm{})},\stackrel{}{v}_j^{(\mathrm{})}\stackrel{}{v}_{\tau _{\mathrm{}}+1}^{(\mathrm{})}=1,jr,1j\stackrel{~}{\tau }_{\mathrm{}}.$$
$`(5.6)_1`$
$$\stackrel{}{m}_r^{(\mathrm{})},\stackrel{}{v}_r^{(\mathrm{})}\stackrel{}{v}_{\tau _{\mathrm{}}+1}^{(\mathrm{})}=1.$$
$`(5.6)_2`$
$$\stackrel{}{m}_r^{(\mathrm{})},\stackrel{}{v}_j^{(q)}\stackrel{}{v}_{\tau _q+1}^{(q)}=0.$$
$`(5.6)_3`$
for $`q[1,k]`$ and $`j[1,\tau _q]j_q`$ for some index $`j_q`$ associated to $`\stackrel{}{m}_r^{(\mathrm{})}.`$
$$\stackrel{}{m}_r^{(\mathrm{})},\stackrel{}{v}_{j_q}^{(q)}\stackrel{}{v}_{\tau _q+1}^{(q)}0.$$
$`(5.6)_4`$
By virtue of the condition $`(5.3)`$ we have a set of $`(nk)`$ independent equations corresponding to the equalities $`(5.6)_1`$,$`(5.6)_3`$ above.
We see that we may choose as $`\stackrel{}{m}_r^{(\mathrm{})}M_𝐑=M_𝐙𝐑`$ induced from the mapping
$$pr:𝐙^n\frac{𝐙^n}{_{q=1}^k𝐙\stackrel{}{g}^{(q)}}:=M_𝐙𝐙^{nk},$$
in view of the relation $`(1.2).`$
Therefore we may uniquely detemine the set of vectors $`\{\stackrel{}{m}_r^{(\mathrm{})}\}M_𝐑`$ as solutions to a system of $`(nk)`$ equations $`(5.6)_1`$,$`(5.6)_3`$ under certain compatibility condition.
To formulate this compatibility condition, let us denote by
$$𝖯=^t(\stackrel{}{m}_1^{(1)},\mathrm{},\stackrel{}{m}_{\tau _1}^{(1)},\stackrel{}{m}_1^{(2)},\mathrm{},\stackrel{}{m}_{\tau _k}^{(k)}),$$
a $`n\times n`$ matrix whose $`b^\mathrm{}1+r`$ th column corresponds to $`m_r^{(\mathrm{})}`$. The system $`(5.6)_{}`$ can be realized by the following matrix equation if it is solvable,
$$(𝖫_\mathrm{\Lambda }\rho V^\mathrm{\Lambda }^tV^\mathrm{\Lambda })𝖯=^T𝖫_\mathrm{\Lambda }^T\rho ^TV^\mathrm{\Lambda }^t(^TV^\mathrm{\Lambda }).$$
$`(5.7)`$
The solvability of this equation can be understood as the compatibility mentioned above. Let us formulate a sufficient condition for the solvability of $`(5.7)`$.
###### Lemma 5.1
Let us assume that all conditions imposed on $`(1.1)`$ and $`{}_{}{}^{T}(1.1)`$ in Theorem 3.1 are satisfied. Furthermore we assume that it is possible to make $`\lambda =id_n`$ in $`(1.12)`$ by means of the rearrangements of rows and columns in $`𝖫_\mathrm{\Lambda }`$. Assume that $`G=id_n`$ (resp. $`{}_{}{}^{T}G=id_n`$) in $`(1.9)`$ (resp. $`{}_{}{}^{T}(1.9)`$). Then the equation $`(5.7)`$ is solvable with respect to $`𝖯`$.
proof The existence of $`k`$ linearly independent eigenvectors $`\stackrel{}{v}_1^{(0)},\mathrm{},\stackrel{}{v}_k^{(0)}`$ $`𝐑^n`$ such that
$${}_{}{}^{t}\stackrel{}{v}_{\mathrm{}}^{(0)}(𝖫_\mathrm{\Lambda }\rho V^\mathrm{\Lambda }^tV^\mathrm{\Lambda })=^t\stackrel{}{v}_{\mathrm{}}^{(0)}(^T𝖫_\mathrm{\Lambda }^T\rho ^TV^\mathrm{\Lambda }^t(^TV^\mathrm{\Lambda }))=\stackrel{}{0}𝐑^k,$$
is a necessary condition for the solvability of the equation $`(5.7)`$. It is also a sufficient condition for the solvability as the following relations show,
$${}_{}{}^{t}(𝖫_\mathrm{\Lambda }\rho V^\mathrm{\Lambda }^tV^\mathrm{\Lambda })GL(W),$$
$${}_{}{}^{t}𝖯^t(𝖫_\mathrm{\Lambda }\rho V^\mathrm{\Lambda }^tV^\mathrm{\Lambda })=^t(^T𝖫_\mathrm{\Lambda }^T\rho ^TV^\mathrm{\Lambda }^t(^TV^\mathrm{\Lambda }))GL(W),$$
for the vector space $`W:=\frac{𝐑^n}{_{\mathrm{}=1}^k𝐑\stackrel{}{v}_{\mathrm{}}^{(0)}}𝐑^{nk}.`$
We will see further that
$${}_{}{}^{t}(𝖫_\mathrm{\Lambda }\rho V^\mathrm{\Lambda }^tV^\mathrm{\Lambda })V^\mathrm{\Lambda }=^t(^T𝖫_\mathrm{\Lambda }^T\rho ^TV^\mathrm{\Lambda }^t(^TV^\mathrm{\Lambda }))V^\mathrm{\Lambda }=0End(𝐑^n,𝐑^k),$$
under the imposed conditions.
Firt we remark that $`{}_{}{}^{t}(\rho V^\mathrm{\Lambda })V^\mathrm{\Lambda }=^t\widehat{Q}`$ in view of $`(1.8)`$ due to the condition $`G=id_n.`$ Thus it is enough to show the equality
$$(^t𝖫_\mathrm{\Lambda })^1V^\mathrm{\Lambda }{}_{}{}^{t}\widehat{Q}=V^\mathrm{\Lambda }.$$
$`(5.8)`$
The left hand side of $`(5.8)`$ is, in its turn, equal to,
$$(^t\stackrel{}{P}_1,\mathrm{},^t\stackrel{}{P}_k)^t\widehat{Q},$$
by virtue of $`(3.6).`$ Let us introduce a matrix,
$$\stackrel{~}{𝖯}=\left(\begin{array}{ccc}\stackrel{~}{p}_1^{(1)}& \mathrm{}& \stackrel{~}{p}_1^{(k)}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ \stackrel{~}{p}_k^{(1)}& \mathrm{}& \stackrel{~}{p}_k^{(k)}\end{array}\right),$$
$`(5.9)`$
for $`{}_{}{}^{T}\xi _{}^{(q)}(z)=_{r=1}^k\stackrel{~}{p}_r^{(q)}(1z_r).`$ With this notation the matrix $`(5.9)`$ equals to
$$V^\mathrm{\Lambda }\stackrel{~}{𝖯}^t\widehat{Q},$$
$`(5.10)`$
if we assume $`(3.10)`$ for $`G=id_n.`$ Under the conditions imposed on $`(1.1)`$ and $`{}_{}{}^{T}(1.1)`$ in Theorem 3.1 we have
$$\left(\begin{array}{c}1z_1\\ \mathrm{}\\ 1z_k\end{array}\right)=\widehat{Q}\left(\begin{array}{c}{}_{}{}^{T}\xi _{}^{(1)}(z)\\ \mathrm{}\\ {}_{}{}^{T}\xi _{}^{(k)}(z)\end{array}\right),$$
which is a mere analogy to $`(3.10).`$ In making use of this equality we see that $`(5.10)`$ is equal to $`V^\mathrm{\Lambda }.`$
To see the equality
$${}_{}{}^{t}(^T𝖫_\mathrm{\Lambda }^T\rho ^TV^\mathrm{\Lambda }^t(^TV^\mathrm{\Lambda }))V^\mathrm{\Lambda }=0End(𝐑^n,𝐑^k),$$
first we prove the equality
$${}_{}{}^{t}(^T𝖫_\mathrm{\Lambda }^T\rho ^TV^\mathrm{\Lambda }^t(^TV^\mathrm{\Lambda }))^TV^\mathrm{\Lambda }=0,$$
in a way parallel to the proof of $`(5.8).`$ Further we see that $`{}_{}{}^{T}V_{}^{\mathrm{\Lambda }}=V^\mathrm{\Lambda }`$ under the condition $`\lambda =id_n.`$ Q.E.D.
As the matrices $`𝖫_\mathrm{\Lambda }\rho V^\mathrm{\Lambda }^tV^\mathrm{\Lambda }`$ and $`{}_{}{}^{T}𝖫_{\mathrm{\Lambda }}^{}^T\rho ^TV^\mathrm{\Lambda }^t(^TV^\mathrm{\Lambda })`$ have rank $`(nk)`$ because of the condition $`(5.3)`$, the columns of the matrix $`𝖯`$ are determined as elements in $`M_𝐑.`$ We recall that the matrices $`\rho V^\mathrm{\Lambda }`$ and $`{}_{}{}^{T}\rho ^TV^\mathrm{\Lambda }`$ have been introduced to formulate the conditions $`(3.11)`$ and $`(3.11)^T.`$ It is clear that the columns of the solution $`𝖯`$ to this equation satisfy the equations $`(5.6)_{}`$ and, in particular, determine $`j_q`$ of $`(5.6)_4`$.
The vectors $`(5.1)`$ admit another interpretation. We introduce unit vectors
$$ϵ_q=(0,\mathrm{},\text{ }^\genfrac{}{}{0pt}{}{q}{^{}}1,0,\mathrm{},0)𝐙^k,1qk.$$
Let us consider a $`(n+k)`$ dimensional cone $`\sigma `$ with $`(n+k)`$ generators
$$\stackrel{~}{v}_0^{(q)}=(0,\mathrm{},0,ϵ_q)𝐙^{n+k},$$
$$\stackrel{~}{v}_j^{(q)}=(\stackrel{}{v}_j^{(q)}\stackrel{}{v}_{\tau _q+1}^{(q)},ϵ_q)𝐙^{n+k},1j\tau _q,1qk.$$
The dual cone $`\stackrel{ˇ}{\sigma }`$ to the cone $`\sigma =𝐑_0\stackrel{~}{v}_0^{(1)},\stackrel{~}{v}_1^{(1)},\mathrm{},\stackrel{~}{v}_{\tau _1}^{(1)},\mathrm{},\stackrel{~}{v}_{\tau _k}^{(k)}`$ is defined as
$$\stackrel{ˇ}{\sigma }=\{y𝐑^{n+k};x,y0\mathrm{for}\mathrm{all}x\sigma \}.$$
The generators of the dual cone $`\stackrel{ˇ}{\sigma }`$ are given by the vectors,
$$\stackrel{~}{m}_0^{(\mathrm{})}:=(0,\mathrm{},0,ϵ_{\mathrm{}}),$$
$$\stackrel{~}{m}_r^{(\mathrm{})}:=(\stackrel{}{m}_r^{(\mathrm{})},ϵ_{\mathrm{}}),1r\tau _{\mathrm{}},1\mathrm{}k$$
with $`\stackrel{}{m}_r^{(\mathrm{})}`$ satisfying the equations $`(5.6)_{}.`$ It is evident that the following inequalities hold for the above vectors,
$$\stackrel{~}{v}_j^{(q)},\stackrel{~}{m}_r^{(\mathrm{})}0.$$
Conversely all vectors $`\stackrel{~}{m}`$ satisfying the conditions
$$\stackrel{~}{v}_j^{(q)},\stackrel{~}{m}0$$
for every $`\stackrel{~}{v}_j^{(q)}`$ must be a linear combination of $`\stackrel{~}{m}_r^{(\mathrm{})}`$’s with positive coefficients as they form a basis of $`M_𝐑\times 𝐑^k`$ the dual space to $`V_𝐑\times 𝐑^k.`$ Thus we get the generators of the dual cone $`\stackrel{ˇ}{\sigma }`$ by means of the vertices of $`\{\mathrm{\Delta }_1^{},\mathrm{},\mathrm{\Delta }_k^{}\}.`$ This is an example of realization of , Theorem 4.6.
Let us formulate a statement on the dual partition $`\{\mathrm{\Delta }_1^{},\mathrm{},\mathrm{\Delta }_k^{}\}`$ that can be characterized as a dual nef-partition to $`\{\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_k\}`$. We spare the space to formulate the definition of the notion of dual nef-partition (see Definition 4.2) by making it clear in the proof of the folowing statement.
###### Proposition 5.2
We assume the following two conditions on $`(1.1)`$ and $`{}_{}{}^{T}(1.1)`$. $`𝐚`$. $`\mathrm{\Delta }_q\mathrm{\Delta }_{\mathrm{}}=\{0\}`$ for $`q\mathrm{}`$. $`𝐛`$. It is possible to choose an integer entry matrix $`𝖯`$ in the matrix equation $`(5.7)`$. Then the following statements hold.
1. $`\{\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_k\}`$ is a nef-partition of the Minkowski sum $`\mathrm{\Delta }_1+\mathrm{}+\mathrm{\Delta }_k`$.
2. $`\{\mathrm{\Delta }_1^{},\mathrm{},\mathrm{\Delta }_k^{}\}`$ is a nef-partition of the Minkowski sum $`\mathrm{\Delta }_1^{}+\mathrm{}+\mathrm{\Delta }_k^{}`$ dual to $`\{\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_k\}`$.
3.The transposed polynomials $`{}_{}{}^{T}f_{2q1}^{}/(^Tf_{2q}1)`$ are obtained from $`\mathrm{\Delta }_q^{}`$ by means of the torus closed embedding,
$$\begin{array}{ccc}𝐏_{(^Tg_1^{(q)},\mathrm{},^Tg_{\stackrel{~}{\tau }_q}^{(q)})}^{\stackrel{~}{\tau }_q1}& & 𝐏^{n1}\\ (x_1,\mathrm{},x_n)& & (y_1,\mathrm{},y_n)\\ & & =(_{\mathrm{}=1}^k_{1j\tau _{\mathrm{}}}x_{b^\mathrm{}1+j}^{\stackrel{}{v}_j^{(\mathrm{})}\stackrel{}{v}_{\tau _{\mathrm{}}+1}^{(\mathrm{})},𝐞_1},\mathrm{},_{\mathrm{}=1}^k_{1j\tau _{\mathrm{}}}x_{b^\mathrm{}1+j}^{\stackrel{}{v}_j^{(\mathrm{})}\stackrel{}{v}_{\tau _{\mathrm{}}+1}^{(\mathrm{})},𝐞_n}),\end{array}$$
where $`𝐞_i=(0,\mathrm{},\text{ }^\genfrac{}{}{0pt}{}{i}{^{}}1,0,\mathrm{},0)`$ $`𝐙^n.`$
Proof
1. The condition $`𝖯End(𝐙)`$ entails the integral linear property of the function $`\varphi _{\mathrm{}}(y)`$ defined in $`(5.4).`$ The convexity of $`\varphi _{\mathrm{}}:M_𝐑𝐑`$ is guaranteed by the following facts. First the number of vertices in $`\mathrm{\Delta }_{\mathrm{}}`$ is less than $`(nk+1)=dimM_𝐑+1`$. Secondly the condition $`𝐚`$.
It is easy to see from the equation $`(5.7)`$ that
$$\varphi _q(\stackrel{}{m}_r^{(\mathrm{})})=\delta _{q,\mathrm{}},1r\tau _{\mathrm{}}.$$
$`(5.11)`$
The existence of an integral upper convex piecewise linear function $`\varphi _q,`$ $`1qk`$ satisfying $`(5.11)`$ is equivalent to the definition of nef- partition $`\{\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_k\}`$ of the Minkowski sum $`\mathrm{\Delta }_1+\mathrm{}+\mathrm{\Delta }_k`$ whose dimension is equal to $`dimV_𝐑`$ after $`(5.3)`$.
2. After the condition $`𝐛`$,
$$\psi _{\mathrm{}}(x)=min_{y\mathrm{\Delta }^{}}x,y,$$
is an integral piecewise linear function on $`V_𝐑.`$ The relation $`\mathrm{\Delta }_q^{}\mathrm{\Delta }_{\mathrm{}}^{}=\{0\}`$ for $`q\mathrm{}`$ is clear from the existence of the funciton $`(5.11).`$ The number of vertices in $`\mathrm{\Delta }_{\mathrm{}}^{}`$ is less than $`(nk+1)=dimV_𝐑+1`$. Thus we see that $`\psi _{\mathrm{}}(x)`$ is an upper convex function. This shows that $`\{\mathrm{\Delta }_1^{},\mathrm{},\mathrm{\Delta }_k^{}\}`$ is a nef-partition of the Minkowski sum $`\mathrm{\Delta }_1^{}+\mathrm{}+\mathrm{\Delta }_k^{}`$. The existence of an integral upper convex piecewise linear function $`(5.11)`$ shown before means that $`\{\mathrm{\Delta }_1^{},\mathrm{},\mathrm{\Delta }_k^{}\}`$ is a dual nef-partition to $`\{\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_k\}.`$
3.First of all we remark that the transposition construction entails,
$$\underset{q=1}{\overset{k}{}}{}_{}{}^{T}\stackrel{~}{g}_{j}^{(q)}(v_j^{(\nu (q))}\stackrel{}{v}_{\stackrel{~}{\tau }_q+1}^{(\nu (q))})=0.$$
By virtue of the convexity of the function $`_{\mathrm{}=1}^k\varphi _{\mathrm{}}`$ (see §2.3) or equivalently
$$\mathrm{\Delta }_1+\mathrm{}+\mathrm{\Delta }_k=\{xV_𝐑;x,y\underset{\mathrm{}=1}{\overset{k}{}}\varphi _{\mathrm{}}(y),yM_𝐑\},$$
the mapping $`(5.7)`$ is a closed torus embedding. Thus a polynomial whose Newton diagram equals to $`\mathrm{\Delta }_{\mathrm{}}^{}`$ in $`(y_1,\mathrm{},y_n)`$ variables coincides with a polynomial obtained as a sum of monomials with exponents from the rows of the RHS of $`(5.7)`$ i.e. $`_{\mathrm{}=1}^k_{1j\stackrel{~}{\tau }_{\mathrm{}}}x_{b^{\nu (\mathrm{}1)}+j}^{{}_{}{}^{T}\stackrel{}{v}_{j}^{(\mathrm{})}^T\stackrel{}{v}_{\stackrel{~}{\tau }_{\mathrm{}}+1}^{(\mathrm{})}}`$. Further argument is parallel to that in , §3. Q.E.D.
## 6 Examples
Example 6.1, Schimmrigk variety As a simple, but non-trivial example we recall the following example whose period integral has been studied in ,
$$f_1(x)=\underset{i=0}{\overset{3}{}}x_i^3,f_2(x)=x_1x_2x_3+1,$$
$$f_3(x)=\underset{i=1}{\overset{3}{}}x_ix_{i+3}^3,f_4(x)=x_0x_4x_5x_6+1.$$
We then have the matrices below after the notation $`(1.5)`$,
$$𝖫=\left(\begin{array}{ccccccccccccc}3& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 3& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 0& 3& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 0& 0& 3& 0& 0& 0& 1& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 1& 0\\ \multicolumn{13}{c}{}\\ 0& 1& 1& 1& 0& 0& 0& 0& 1& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 1& 0& 0& 3& 0& 0& 0& 0& 1& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 0& 1& 0& 0& 3& 0& 0& 0& 1& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 0& 0& 1& 0& 0& 3& 0& 0& 1& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 1\\ \multicolumn{13}{c}{}\\ 1& 0& 0& 0& 1& 1& 1& 0& 0& 0& 1& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0\end{array}\right).$$
$$𝖫^1=\left(\begin{array}{ccccccccccccc}1/3& 1/9& 1/9& 1/9& 0& 1/3& 1/3& 0& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 2/9& 1/9& 1/9& 0& 1/3& 1/3& 0& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 1/9& 2/9& 1/9& 0& 1/3& 1/3& 0& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 1/9& 1/9& 2/9& 0& 1/3& 1/3& 0& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 1/9& 1/27& \frac{2}{27}& \frac{2}{27}& 0& 1/9& 1/9& 2/9& 1/9& 1/9& 0& 1/3& 1/3\\ \multicolumn{13}{c}{}\\ 1/9& \frac{2}{27}& 1/27& \frac{2}{27}& 0& 1/9& 1/9& 1/9& 2/9& 1/9& 0& 1/3& 1/3\\ \multicolumn{13}{c}{}\\ 1/9& \frac{2}{27}& \frac{2}{27}& 1/27& 0& 1/9& 1/9& 1/9& 1/9& 2/9& 0& 1/3& 1/3\\ \multicolumn{13}{c}{}\\ 0& 1/3& 1/3& 1/3& 0& 1& 1& 0& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 1/3& 1/9& 1/9& 1/9& 0& 0& 0& 1/3& 1/3& 1/3& 0& 1& 1\\ \multicolumn{13}{c}{}\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1\\ \multicolumn{13}{c}{}\\ 0& 1/3& 1/3& 1/3& 1& 1& 1& 0& 0& 0& 0& 0& 0\\ \multicolumn{13}{c}{}\\ 1/3& 1/9& 1/9& 1/9& 0& 0& 0& 1/3& 1/3& 1/3& 1& 1& 1\end{array}\right).$$
After the construction $`(1.4)^T`$ based on the transposed matrix $`{}_{}{}^{t}(𝖫)`$ it is easy to see that $`{}_{}{}^{T}f_{\mathrm{}}^{}(x)=f_{\mathrm{}}(x),1\mathrm{}4.`$ The matrix $`{}_{}{}^{T}𝖫_{}^{1}=𝖫^1`$ gives us linear functions,
$$\xi ^{(1)}=\frac{1}{3}(z_11)+\frac{1}{9}(z_21),\xi ^{(2)}=\frac{1}{3}(z_21).$$
For $`f(x)`$ as well as for $`{}_{}{}^{T}f(x)`$ we can calculate the Mellin transform of the period integral,
$$M_{0,\gamma }^0(z)=\mathrm{\Gamma }(\frac{1}{3}(z_11)+\frac{1}{9}z_2)^3\mathrm{\Gamma }(\frac{1}{3}(z_21))^4\mathrm{\Gamma }(z_1)\mathrm{\Gamma }(z_2)=\frac{\mathrm{\Gamma }(\xi ^{(1)})^3\mathrm{\Gamma }(\xi ^{(2)})^4}{\mathrm{\Gamma }(3\xi ^{(1)}+\xi ^{(2)})\mathrm{\Gamma }(3\xi ^{(2)})},$$
up to $`27`$periodic functions.
Analogously we can look at the CI defined on $`𝐓^{2n+1}`$
$$f_1(x)=\underset{i=0}{\overset{n}{}}x_i^n,f_2(x)=x_1x_2\mathrm{}x_n+1,$$
$$f_3(x)=\underset{i=1}{\overset{n}{}}x_ix_{i+n}^n,f_4(x)=x_0x_{n+1}x_{n+2}\mathrm{}x_{2n}+1,$$
whose period integral can be expressed through its Mellin transform,
$$M_{0,\gamma }^0(z)=\frac{\mathrm{\Gamma }(\xi ^{(1)})^n\mathrm{\Gamma }(\xi ^{(2)})^{n+1}}{\mathrm{\Gamma }(n\xi ^{(1)}+\xi ^{(2)})\mathrm{\Gamma }(n\xi ^{(2)})},$$
up to $`n^n`$periodic functions.
It is worthy to notice that our matrix $`𝖫`$ satisfies an interesting condition below.
The magic square condition For each fixed $`q[1,k]`$, we can find a single valued mapping $`\sigma :bI_\mathrm{\Lambda }\{1,\mathrm{},n\}`$ such that
$$p_q^b=w_{\sigma (b)}^{a^q},\mathrm{for}\mathrm{all}bI_\mathrm{\Lambda }.$$
This condition plays central rôle in the interpretation of the strange duality found by Arnol’d on the interchange between Gabrielov number and Dolgachev number from the point of view of the mirror symmetry .
Example 6.2 We consider an example of an hypersurface studied in . We have the following data after the notations above,
$$f_1(x)=x_1^7+x_2^7x_4+x_3^7x_5+x_4^3+x_5^3,f_2=x_1x_2x_3x_4x_5.$$
$$𝖫=\left[\begin{array}{cccccccc}7& 0& 0& 0& 0& 1& 0& 0\\ \multicolumn{8}{c}{}\\ 0& 7& 0& 1& 0& 1& 0& 0\\ \multicolumn{8}{c}{}\\ 0& 0& 7& 0& 1& 1& 0& 0\\ \multicolumn{8}{c}{}\\ 0& 0& 0& 3& 0& 1& 0& 0\\ \multicolumn{8}{c}{}\\ 0& 0& 0& 0& 3& 1& 0& 0\\ \multicolumn{8}{c}{}\\ 0& 0& 0& 0& 0& 1& 0& 1\\ \multicolumn{8}{c}{}\\ 1& 1& 1& 1& 1& 0& 1& 0\\ \multicolumn{8}{c}{}\\ 0& 0& 0& 0& 0& 0& 1& 0\end{array}\right],$$
$$𝖫^1=\left[\begin{array}{cccccccc}\frac{6}{49}& 1/49& 1/49& \frac{2}{49}& \frac{2}{49}& 0& 1/7& 1/7\\ \multicolumn{8}{c}{}\\ \frac{2}{147}& \frac{19}{147}& \frac{2}{147}& \frac{11}{147}& \frac{4}{147}& 0& 2/21& 2/21\\ \multicolumn{8}{c}{}\\ \frac{2}{147}& \frac{2}{147}& \frac{19}{147}& \frac{4}{147}& \frac{11}{147}& 0& 2/21& 2/21\\ \multicolumn{8}{c}{}\\ 1/21& 1/21& 1/21& \frac{5}{21}& 2/21& 0& 1/3& 1/3\\ \multicolumn{8}{c}{}\\ 1/21& 1/21& 1/21& 2/21& \frac{5}{21}& 0& 1/3& 1/3\\ \multicolumn{8}{c}{}\\ 1/7& 1/7& 1/7& 2/7& 2/7& 0& 1& 1\\ \multicolumn{8}{c}{}\\ 0& 0& 0& 0& 0& 0& 0& 1\\ \multicolumn{8}{c}{}\\ 1/7& 1/7& 1/7& 2/7& 2/7& 1& 1& 1\end{array}\right]$$
We have therefore the Mellin transform of the period integral associated to the CI $`\{x`$ $``$ $`(𝐂^\times )^5`$ ; $`f_1(x)`$ $`+`$ $`s_1=0,f_2(x)+1=0\}`$ up to $`7`$periodic functions,
$$M_{0,\gamma }^0(z)=\mathrm{\Gamma }(\frac{1}{7}(z_11))^3\mathrm{\Gamma }(\frac{2}{7}(z_11))^2\mathrm{\Gamma }(z_1)=\frac{\mathrm{\Gamma }(\xi ^{(1)})^3\mathrm{\Gamma }(2\xi ^{(1)})^2}{\mathrm{\Gamma }(7\xi ^{(1)})}.$$
After the construction $`(1.4)^T`$ we see that
$${}_{}{}^{T}f_{1}^{}(x)=x_1^7+x_2^7+x_3^7+x_2x_4^3+x_3x_5^3,^Tf_2=x_1x_2x_3x_4x_5.$$
We then have the Mellin transform of the period integral associated to the CI $`\{x(𝐂^\times )^5;^Tf_1(x)+s_1=0,^Tf_2(x)+1=0\}`$ up to $`147`$periodic functions,
$$M_{0,\gamma }^0(z)=\mathrm{\Gamma }(\frac{1}{7}(z_11))\mathrm{\Gamma }(\frac{2}{21}(z_11))^2\mathrm{\Gamma }(\frac{1}{3}(z_11))^2\mathrm{\Gamma }(z_1)=\frac{\mathrm{\Gamma }(3\xi ^{(1)})\mathrm{\Gamma }(2\xi ^{(1)})^2\mathrm{\Gamma }(7\xi ^{(1)})^2}{\mathrm{\Gamma }(21\xi ^{(1)})}.$$
Indepent University of Moscow
Bol’shoj Vlasievskij pereulok 11,
Moscow, 121002,
Russia
E-mails: [email protected], [email protected]
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# Citation Statistics From 110 Years of Physical Review0footnote 00footnote 0©2005 American Institute of Physics. This article may be downloaded for personal use only. Any other use requires prior permission of the author and the American Institute of Physics. This article appeared in Physics Today, volume 58, page 49, June 2005 and the online version may be found at http://www.physicstoday.org/vol-58/iss-6/p49.shtml.
## Introduction
The first particle published in the Physical Review was received in 1893; the journal’s first volume included 6 issues and 24 articles. In the 20th century, the Physical Review branched into topical sections and spawned new journals. Today, all articles in the Physical Review family of journals (PR) are available online and, as a useful byproduct, all citations in PR articles are electronically available.
The citation data provide a treasure trove of quantitative information. As individuals who write scientific papers, most of us are keenly interested in how often our own work is cited. As dispassionate observers, we can use the citation data to identify influential research, new trends in research, unanticipated connections across fields, and in subfields that are exhausted. A certain pleasure can also can be gleaned from the data when they reveal the idiosyncratic features in the citation histories of individual publications.
The investigation of citation statistics has a long history ER90 in which a particularly noteworthy contribution was a 1965 study by Derek John de Solla Price P65 . In his study, Price built upon original models by George Yule and Herbert Simon Y25 to argue that the distribution in the number of citations to individual publications had a power-law form. Price also noted that well-cited papers continue to be referenced more frequently than less-cited papers, and coined the term cumulative advantage to describe the mechanism that causes a persistently higher rateM73 . Cumulative advantage means that the probability that a publication is cited is an increasing function of its current number of citations.
In the framework of current fashionable evolving network models, the mechanism is called preferential attachment BA99 . Linear preferential attachment models provide appealing explanations for the power-law distributions of connections that are observed in many social systems, natural networks, and manmade networks such as the World Wide Web rev . One fundamental motivation for studying citation statistics is to determine whether they exhibit some of the universal features that have been ascribed to prototypical models of evolving networks BA99 ; KRL00 ; DMS00 .
Before examining the citation data, I offer several caveats: First, the data include only internal citations — that is, citations from PR articles to other PR articles — and are perforce incomplete. For highly cited papers, a previous study Re04 found that total citations typically outnumber internal ones by a factor of 3 to 5, a result that gives a sense of the incompleteness of the PR data. Second, some 5–10% of citations appear to be erroneous Re04 ; error , although the recent practice by PR of crosschecking references when manuscripts are submitted has significantly reduced the error rate. Third, papers can be highly cited for many reasons — some substantive and some dubious. Thus the number of citations is merely an approximate proxy for scientific quality.
## Citation distribution and attachment rate
The PR citation cover 353,268 papers and 3,110,839 citations from July 1893 through June 2003. The 329,847 publications with at least 1 citation may be broken down as follows:
| 11 | publications with | $`>`$ | 1000 citations |
| --- | --- | --- | --- |
| 79 | publications with | $`>`$ | 500 citations |
| 237 | publications with | $`>`$ | 300 citations |
| 2,340 | publications with | $`>`$ | 100 citations |
| 8,073 | publications with | $`>`$ | 50 citations |
| 245,459 | publications with | $`<`$ | 10 citations |
| 178,019 | publications with | $`<`$ | 5 citations |
| 84,144 | publications with | | 1 citation |
A somewhat depressing observation is that nearly 70% of all PR articles have been cited less than 10 times. (The average number of citations is 8.8.) Also evident is the small number of highly cited publications; table LABEL:tab-top-10 lists the 11 publications with more than 1000 citations.
Citations have grown rapidly with time, a feature that mirrors the growth of PR family of journals. From 1893 until World War II, the number of annual citations, from PR publications doubled approximately every 5.5 years. The number of PR articles published in a given year also doubled every 5.5 years. Following the publication crash of the war years, the number of articles published annually doubled approximately every 15 years.
The citation data naturally raise the question, What is the distribution of citations? That is, what is the probability $`P(k)`$ that a paper gets cited $`k`$ times? This question was investigated by Price, who posited the power law $`P(k)k^\nu `$, with $`\nu `$ positive. A power-law form is exciting for statistical physicists because it implies the absence of a characteristic scale for citations — the influence of a publication may range from useless to earth-shattering. The absence of such a characteristic scale in turn implies that citations statistics should exhibit many of the intriguing features associated with phase transitions, which display critical phenomena on all length scales.
Somewhat surprisingly, the probability distribution derived from more than 3 million PR citations still has significant statistical fluctuations. It proves useful to study the cumulative distribution, $`C(k)=_k^{\mathrm{}}P(k^{})𝑑k^{}`$, the probability that a paper is cited at least $`k`$ times, to reduce these fluctuations.
On a double logarithmic scale, $`C(k)`$ has a modest negative curvature everywhere. That behavior, illustrated in Fig. 1, suggests that the distribution decays faster than a power law and is at variance with previous, smaller-scale studies that suggested either a power-law P65 ; R98 or a stretched exponential form LS98 , $`P(k)\mathrm{exp}(k^\beta )`$, with $`\beta <1`$. It is intriguing that a good fit over much of the range of the range of the distribution is the log-normal form $`C(k)=Ae^{b\mathrm{ln}kc(\mathrm{ln}k)^2}`$. Log-normal forms typically underlie random multiplicative processes. The describe, for example, the distribution of fragment sizes that remain after a rock has been hammered many times.
The development of citations may be characterized by the attachment rate $`A_k`$, which gives the likelihood that a paper with $`k`$ citations will be cited by a new article. To measure the attachment rate, first count the number of times each paper is cited during a specified time range; this gives $`k`$. Then, to get $`A_k`$, count the number of times each paper with a given $`k`$ in this time window was cited in a subsequent window. As shown in Fig. 2, the data suggest that $`A_k`$ is a linear function of $`k`$, especially for $`k<150`$, a condition that applies to nearly all PR papers JNB01 . Thus linear preferential attachment appears to account for the propagation of citations.
Linear attachment, though, leads to two paradoxes. First, a linear rate implies a power-law citation distribution, but Fig. 1 indicates that the data are better described by a log-normal form. While a log-normal distribution does arise from the nearly linear attachment rate $`A_k=k/(1+a\mathrm{ln}k)`$, with $`a`$ positive, Fig. 2 hints that $`A_k`$ may be growing slightly faster than linearly with $`k`$. Second, to implement linear preferential attachment consciously, a citer must be aware of the references to every existing paper. That’s clearly not realistic. A more realistic process that can lead to linear preferential attachment is the redirection mechanism K99 ; KRL00 . In redirection, an author who is writing the reference list for a paper figuratively picks a random paper. Then the author cites either the random selected paper (with probability $`1r`$) or one of the references within that paper (with probability $`r`$). This purely local mechanism generates the linear form $`A_k=k+(\frac{1}{r}2)`$ KRL00 . Still a mystery is why the myriad of attributes that influences whether a paper gets cited manifests itself as a citation rate that is a nearly linear function of the number of citations.
## Age structure
A common adage says that nobody cites classic papers anymore. Is this really true? How long does a paper continue to get cited?
The age of a citation is the difference between the year when a citation occurs and the publication year of the cited paper. Typically, unpopular papers are cited soon after publication, if at all, and then disappear. The converse is also true. For example, the average citation age $`a`$ over the entire PR data set is 6.2 years, but articles published before 2000 for which $`a`$ is less than 2 years receive, on average, only 3.6 citations. On the other hand, highly cited papers usually continue to be cited for a long time, and vice versa. Papers with more than 100 citations have $`a=11.7`$ years, and the 11 publications with more than 1000 citations have $`a=18.9`$ years. For all PR publications with500 or fewer citations, the average citation age grows with the number of citations roughly as $`a=k^\alpha `$, with $`\alpha 0.3`$.
The citation age distributions — that is, the number of citations a a function of age — reveal a fact that is surprising at first sight: The exponential growth of PR articles strongly skews the form of the age distributions! There are, in fact two distinct age distributions. One is the distribution of citing ages, defined as the number of citations of a given age from a paper. Remarkably, citing memory is independent of when a paper is published. An author publishing now seems to have the same range of memory as an author who published an article 50 years ago. The citing age distribution roughly decays exponentially in age, except for a sharp decrease in citations during the period of World War II. However, citing an old paper is difficult a priori simply because relatively few old publications exist. As noted by Hideshiro Nakamoto N88 , a more meaningful citing age distribution is obtained by rescaling the distribution by the total number of publications in each citing year. So, if one is interested in journal citations from 2005, the number of cited papers that are, say, four years old should be scaled by the total number of papers published in 2001. The rescaling has a dramatic effect: The nearly exponential citing age distribution is transformed into a power-law! An analogous skewing due to the rapid growth of PR articles also occurs in the distribution of cited ages, that is, the number of citation of a given age to an article.
## Individual citation histories
The citation histories of well-cited publications are diverse from the collective citation history of all PR articles. The varied histories roughly fall into classes that include revived classic works or “sleeping beauties” R04 , major discoveries, and hot publications. It’s fun to examine examples of each class.
Sometimes a publication will become in vogue after a long dormancy — a revival of an old classic. I arbitrarily define a revived classic as a nonreview PR articles, published before 1961, that has received more than $`250`$ citations and has a ratio of the average citation age to the age of the paper greater than 0.7. Thus, revived classics are well-cited old papers with the bulk of their citations occurring long after publication. Only the 12 papers listed in table LABEL:tab-classic fit these criteria.
The clustered citation histories of the five articles plotted in in Fig. 3 are particularly striking. These articles, published between 1951 and 1960 (with three in the same issue of Physical Review), investigated the double exchange mechanism in Perovskite manganites, the mechanism responsible for the phenomenon of colossal magnetoresistance. This topic became in vogue in the 1990’s because of the confluence of new synthesis and measurement techniques in thin-film transition-metal oxides, the sheer magnitude of the effect, and the clever use of the term “colossal”. The citation burst more than 40 years after the publication of these five articles is unique in the history of the PR journals.
The other seven papers have different claims to fame. Eugene Wigner’s 1932 paper had 115 citations before 1980 and 447 through June 2003. Similarly, the Albert Einstein, Boris Podolsky and Nathan Rosen (EPR) paper had 36 citations before 1980 and 456 more through June 2003. With average citation ages of 55.8 and 59.6 respectively, these are the longest-lived articles with more than 35 citations in the PR family. Those papers, as well as the one by Yakir Aharonov and David Bohm, owe their renewed popularity to the upsurge of interest in quantum information phenomena. Wigner’s 1934 papers deals with correlations in an electron gas, a problem of enduring interest in condensed matter physics. Julian Schwinger’s work is a classic contribution to quantum electrodynamics. The 1958 publication by Philip Anderson helped launched the field of localization. And Richard Feynman’s paper presented a widely applicable method for calculating molecular forces. Feynman’s article is noteworthy because it is cited by all PR journals (except the accelerators and beams special topics journal).
Publications that announce discoveries often receive a citation spike when the contribution becomes recognized. I arbitrarily define a discovery paper has having more than 500 citations and a ratio of average citation age to publication age less than 0.4; I exclude articles published in Reviews of Modern Physics and compilations by the Particle Data Group. Table LABEL:tab-discovery lists the 11 such discovery papers; all were published between 1962 and 1991. A trend in this group of papers is the shift from elementary-particle physics (the six articles published before 1976) to condensed-matter physics (the five articles published after 1983). The earlier discovery papers reflected major developments in elementary-particle physics, including $`SU(3)`$ symmetry, the prediction of charm, and grand unified theories. The condensed matter articles are on quasicrystals, multifractrals, and high-temperature superconductivity. If the citation threshold is relaxed to 300, an additional seven papers fit the discovery criteria. All of these are concerned with high-temperature superconductivity, and all but one appear during the golden age of the field, 1987–89.
It is not clear whether the shift in the field of discovery publications stems from a sea change in research direction or because of from prosaic concerns. The past 15 years have seen a major upsurge in quantum condensed-matter physics that perhaps stems from the discovery of high-temperature superconductivity. But recent elementary-particle physics discoveries may be underrepresented in PR because many CERN-based discoveries have been published in journals outside the PR family.
A number of classic, highly cited publications have noteworthy citations histories. Fig. 4 illustrates some of these histories. Citations to “Theory of Superconductivity”, Phys. Rev. 108, 1175 (1957) by John Bardeen, Leon Cooper and J. Robert Schrieffer (BCS) closely track the activity in superconductivity; the paper received its fewest citations in 1985, the year before the discovery of high-temperature superconductivity. The BCS paper is the earliest with more than 1000 citations in the PR family. Steven Weinberg’s paper (W) “A Model of Leptons”, on the electroweak theory, (Phys. Rev. Lett. 19, 1264 (1967)), has a broad citation peak followed by a relatively slow decay as befits this seminal paper’s long-term influence. On the other hand, the average citation age for the 1974 publications that announced the discovery of the $`J/\psi `$ particle – Phys. Rev. Lett. 33, 1404 and 1406 (1974) – is less than 3 years!
An unusual example is “Scaling Theory of Localization: Absence of Quantum Diffusion in Two Dimensions”, Phys. Rev. Lett. 42, 673 (1979) by Elihu Abrahams, Anderson, Don Licciardello and T. V. Ramakrishnan (the so-called gang of four; G4). Since publication, the G4 paper has been cited 30–60 times annually, a striking testament to its long-term impact. The paper with the most citations in all PR journals is “Self-Consistent Equations Including Exchange and Correlation Effects”, Phys. Rev. 140, A1133 (1965) by Walter Kohn and Lu Sham (KS). Amazingly, citations to this publication have been steadily increasing for nearly 40 years.
The KS paper is also an example of what may be called a hot paper, defined as a nonreview paper with 350 or more citations, a ratio of average citation age to publication age greater than two-thirds, and a citation rate increasing with time. Ten papers, listed in table LABEL:tab-hot, fit these criteria. The 1932 Wigner and 1935 EPR articles, both more than 70 years old, and the two most-cited PR papers of all time, KS and the 1964 article by Pierre Hohenberg Kohn, are all still hot. Astounding!
Of the remaining six hot papers, five are in quantum condensed-matter physics. Three of them build on the formalism introduced in the seminal articles by Hohenberg and Kohn and by Kohn and Sham. Another, Anderson’s 1958 localization paper, can be viewed both as hot and as the revival of a classic. The newest hot article, by Charles Bennett and coauthors, is concerned with quantum information theory, a research area that has recently become fashionable and also led to the sharp increase in citations to Wigner’s 1932 paper and the 1935 EPR paper.
## A unique window
A small number of physicists have played a remarkably large role in top-cited PR publications. Two individuals have coauthored five papers from among the top 100 cited PR articles Re04 : Kohn, who occupies positions 1, 2, 24, 96, and 100, and Anderson, with positions 9, 19, 20, 35, and 41. Wigner appears four times (4, 8, 53, and 55), and Lars Onsager (16, 64, and 68) and John Slater (12, 27, and 40) each appear three times.
The PR citation data provide a unique window with which to study the development of citations, and the work I have described can be extended and applied in many ways. For example, constructing a graphical representation of the entire dynamically changing citation network should be revealing. Such a graph could show how fields develop and could expose unexpected connections between disparate areas. A practical, if more controversial, use of citation data would be to construct retrospective journals that include only highly cited papers. Such journals would provide a welcome reduction in the total literature volume, because only 30% of all articles have more than 10 citations and a mere 2.3% have more than 50 citations. A repository for all publications would still be necessary, as sleeping beauties do emerge long after publication.
###### Acknowledgements.
I thank Mark Doyle of the American Physical Society editorial office for providing the citation data, Jon Kleinberg for initial collaboration, Andy Cohen and Andy Millis for literature advice, an anonymous referee for pointing out Nakamoto’s work, Paul Krapivsky and Mark Newman for helpful manuscript suggestions, and Claudia Bondila and Guoan Hu for writing Perl scripts for some data processing. Finally, I am grateful to NSF grant DMR0227670 (BU) and DOE grant W-7405-ENG-36 (LANL) for financial support.
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# Lyubeznik’s invariants for cohomologically isolated singularities.
## 1. Introduction
Let $`A=R/I`$ for $`I`$ an ideal in a regular (local) ring $`(R,m)`$ of dimension $`n`$ and containing a field $`k`$. The main results of \[Lyu93, HS93\] state that the local cohomology module $`H_m^a(H_I^{ni}(R))`$ is injective and supported at $`m`$. Therefore it is a finite direct sum of $`e=e(H_m^a(H_I^{ni}(R)))`$ many copies of the injective hull $`E_{R/m}`$ of the residue field of $`R`$. Lyubeznik shows in \[Lyu93\] that this number
$$\lambda _{a,i}(A)\stackrel{\mathrm{def}}{=}e(H_m^a(H_I^{ni}(R)))$$
does not depend on the auxiliary choice of $`R`$ and $`I`$. If $`A`$ is a complete intersection, these invariants are essentially trivial (all are zero except $`\lambda _{d,d}=1`$ where $`d=dimA`$). The goal of this paper is to describe these invariants for a large class of rings, including those which are complete intesections away from the closed point. Alternatively this class of rings can be viewed as consisting of the rings which behave cohomologically like an isolated singularity.
###### Theorem 1.1.
Let $`A=𝒪_{𝒴,𝓍}`$ for $`Y`$ a closed $`k`$–subvariety of a smooth variety $`X`$. If for $`id`$ the modules $`H_{[Y]}^{ni}(𝒪_𝒳)`$ are supported in the point $`x`$ then
1. For $`2ad`$ one has
$$\lambda _{a,d}(A)\delta _{a,d}=\lambda _{0,da+1}(A)$$
and all other $`\lambda _{a,i}(A)`$ vanish.
2. $$\lambda _{a,d}(A)\delta _{a,d}=\{\begin{array}{cc}dim_{𝔽_p}H_{\{x\}}^{da+1}(Y_{\text{ét}},𝔽_p)\hfill & \text{if }\mathrm{char}k=p\hfill \\ dim_{}H_{\{x\}}^{da+1}(Y_{\mathrm{an}},)\hfill & \text{if }k=\hfill \end{array}$$
where $`\delta _{a,d}`$ is the Kroneker delta function.
In the case that $`A`$ has only an isolated singularity and $`k=`$, this was shown by Garcia Lopez and Sabbah in \[GLS98\]. In \[BB04\] Bondu and myself proved part (1) in all characteristics but part (2) only under an additional assumption. The point of this note is to show that this additional assumption was unnecessary, an observation which also makes the proof much clearer.
## 2. Proof of the result
Let us first fix our notation. Let $`YX`$ be a closed subscheme of $`X`$. Let $`X`$ be smooth of dimension $`n`$ and $`Y`$ is of dimension $`d`$. Let $`xY`$ be a point. We take the freedom to shrink $`X`$ (and $`Y`$) since everything is local at $`x`$. We denote the inclusions as follows:
To ease notation we carry out the argument in the case when $`\mathrm{char}k`$ is positive. The proof of the characteristic zero result is exactly the same, replacing the Emerton–Kisin correspondence by the Riemann–Hilbert correspondence.
Building upon \[BB04\] it remains to show that
$$\lambda _{a,d}(A)\delta _{a,d}=dimH_{\{x\}}^{da+1}(Y_{\text{ét}},𝔽_p)$$
for $`2id`$. By definition of the invariants
$$\lambda _{a,i}(A)=e(H_{[x]}^a(H_{[Y]}^{ni}(𝒪_𝒳)))$$
this comes down to computing $`dim\mathrm{Sol}(H_{[x]}^a(H_{[Y]}^{ni}(𝒪_𝒳))`$. The only trick in the proof is to replace $`H_{[Y]}^{ni}(𝒪_𝒳)`$ by something more accessible – that is by something whose solutions $`\mathrm{Sol}`$ can readily be computed. The rest is mere computation.
Quite generally one has a short exact sequence
$$0\stackrel{}{}K\stackrel{}{}j_!^{}H_{[Y]}^{ni}(𝒪_𝒳)|_{𝒳𝓍}\stackrel{}{}_{[𝒴]}^{𝓃𝒾}(𝒪_𝒳)\stackrel{}{}𝒞\stackrel{}{}\mathcal{0}$$
where, by construction, the kernel $`K`$ and quotient $`C`$ are both supported on $`x`$.<sup>1</sup><sup>1</sup>1The functors we call $`j_{}`$ and $`j_!`$ are denoted by $`j_+`$ and $`j_!+`$ in \[EK04\]. In order to adjust to the notation used in the Riemann–Hilbert correpondence I changed this notation here.<sup>2</sup><sup>2</sup>2The intermediate extension $`j_!M`$ of a module ($`𝒟_{(𝒳𝓍)}`$– or $`𝒪_{,(𝒳𝓍)}`$–module if in characteristic $`0`$ or $`p`$ respectively) is defined as the unique smallest submodule $`M^{}`$ of $`H^0j_{}(M)`$ such that $`j^!M^{}M`$, where $`j:Xx\stackrel{}{}X`$ is the open inclusion of the complement of $`x`$ into $`X`$. The intermediate extension appears here as a substitute for $`j_!`$ which does not exist in the characteristic $`p`$ context. It was this realization (replace $`j_!`$ with $`j_!`$) which made it possible to make the argument work in all characteristics. Splitting this four term sequence into two short exact sequences and using the long exact sequence for $`\mathrm{\Gamma }_{[x]}(\underset{¯}{})`$ we get that for $`a2`$
(2.1)
$$H_{[x]}^a(H_{[Y]}^{ni}(𝒪_𝒳))_{[𝓍]}^𝒶(𝒿_!^{}(_{[𝒴]}^{𝓃𝒾}(𝒪_𝒳)|_{𝒳𝓍})).$$
This first substitution, combined with \[BB04, Lemma 2.3\], we record as a Lemma:
###### Lemma 2.1.
With notation as above, and without any assumptions on the singularities of $`A=\mathrm{Spec}𝒪_{𝒴,𝓍}`$
$$\lambda _{a,d}(A)=dim(H^aj_!^{}\mathrm{Sol}H_{[Yx]}^{nd}(𝒪_{𝒳𝓍}))_𝓍$$
for $`2ad`$.
###### Proof.
By definition we have
$$\begin{array}{cc}\hfill \lambda _{a,d}(A)& =e(H_{[x]}^aH_{[Y]}^{nd}(𝒪_𝒳))\hfill \\ & =dim(\mathrm{Sol}H_{[x]}^aH_{[Y]}^{nd}(𝒪_𝒳)))_𝓍\hfill \\ & =dim(k_!k^1H^a\mathrm{Sol}H_{[Y]}^{nd}(𝒪_𝒳)))_𝓍\text{[BB04, Lemma 2.3]}\hfill \\ & =dim(H^a\mathrm{Sol}j_!^{}(H_{[Y]}^{nd}(𝒪_𝒳)|_{𝒳𝓍}))_𝓍\hfill \end{array}$$
where $`k:x\stackrel{}{}X`$ is the inclusion of the point. The commutation of $`\mathrm{Sol}`$ with $`j_!^{}`$ now finishes the argument. ∎
The assumption that $`H_{[Y]}^{ni}(𝒪_𝒳)`$ is supported at $`x`$ for $`id`$ we rephrase by saying that one has a quasi-isomorphism of complexes
$$H_{[Yx]}^{nd}(𝒪_{𝒳𝓍})𝐑\mathcal{\Gamma }_{[𝒴𝓍]}(𝒪_{𝒳𝓍})[𝓃𝒹].$$
and the solutions of the latter can be computed easily<sup>3</sup><sup>3</sup>3In the case $`k=`$ this computation yields $`i_!_{Yx}[d]`$ instead. (and is preverse!) as done in \[BB04\]:
$$\mathrm{Sol}(𝐑\mathrm{\Gamma }_{[Yx]}(𝒪_{𝒳𝓍})[𝓃𝒹])=𝒾_!^{}(𝔽_𝓅)_{𝒴𝓍}[𝒹]$$
Since $`j`$ is just the inclusion of complement of a point we have that
$$j_!(\underset{¯}{})\tau _{d1}𝐑j_{}^{}(\underset{¯}{})$$
by \[Bor84, V.2.2 (2)\]. Continuing the computation of Lemma 2.1 using these observations we get for $`a1`$ that
$$\begin{array}{cc}\hfill \lambda _{a,i}(A)& =dim(H^aj_!^{}\mathrm{Sol}H_{[Yx]}^{nd}(𝒪_{𝒳𝓍}))_𝓍\hfill \\ & =dim(H^aj_!^{}\mathrm{Sol}(𝐑\mathrm{\Gamma }_{[Yx]}(𝒪_{𝒳𝓍})[𝓃𝒹]))_𝓍\hfill \\ & =dim(H^aj_!^{}i_!^{}(𝔽_p)_{Yx}[d])_x\hfill \\ & =dim(H^ai_!j_!(𝔽_p)_{Yx}[d])_x\hfill \\ & =dim(H^{da}\tau _{d1}𝐑j_{}(𝔽_p)_{Yx})_x\hfill \\ & =dim(H^{da}𝐑j_{}(𝔽_p)_{Yx})_x\hfill \end{array}$$
The latter was computed in \[BB04, Lemma 2.7\] to be equal to $`H_{\{x\}}^{da+1}(Y_{\text{ét}},𝔽_p)+\delta _{a,d}`$ as required.
###### Remark 2.2.
A slight refinement of the same techniques yield also to a more general statement. Namely, if one requires vanishing of $`H_m^aH_I^{ni}(R)|_{\mathrm{Spec}R\mathrm{pt}}`$ below the diagonal $`a=i+m`$ then the result remains true in the range $`dm+2ad`$.
## 3. Examples
This section is to provide some examples of the bad behavior of the invariants $`\lambda _{a,i}`$ under reduction to positive characteristic. The uniformity of the Theorem for all characteristics seems, on the first sight, to suggest that that one can expect a good behavior of the invariants under reduction mod $`p`$. This impression is however quickly shattered, essentially for reasons that local cohomology is well known to behave poorly under reduction. Alternatively one also can observe that the cohomology theory corresponding to $`H_x^i(Y_{an},)`$ under reduction is not $`H_x^i(Y_{\text{ét}},𝔽_p)`$ but rather $`p`$-adic rigid cohomology or crystalline cohomology, of which $`H_x^i(Y_{\text{ét}},𝔽_p)`$ is only a small part, namely the slope zero part.
The examples that follow are standard examples for the bad behaviour of local cohomology under reduction mod $`p`$. I learned them from a talk by Anurag Singh at Oberwofach in March 2005. Our general setup is as follows: Let $`A=R/I`$ where $`R`$ is a polynomial ring over $``$ and $`I`$ is a homogeneous ideal. We denote by $`A_0=A_{}`$ the generic characteristic zero model and by $`A_p=A𝔽_p`$ for all $`p`$ prime the positive characteristic models.
###### Example 3.1.
Let $`R=\left[\begin{array}{ccc}u& v& w\\ x& y& z\end{array}\right]`$ and $`I=(\delta _1,\delta _2,\delta _3)`$ be the ideal of $`2\times 2`$ minors of the displayed matrix of variables. Then $`A=R/I`$ has a free resolution (as an $`R`$–module) given by
$$0\stackrel{}{}R^2\stackrel{\left(\begin{array}{cc}u& x\\ v& y\\ w& z\end{array}\right)}{}R^3\stackrel{\left(\delta _1\delta _2\delta _3\right)}{}R\stackrel{}{}0$$
This shows that $`\mathrm{Ext}^3(R/I,R)=0`$ and therefore, reducing mod $`p`$, that $`\mathrm{Ext}^3(R_p/I^{[p^e]},R_p)=0`$ for all $`e`$ by the flatness of the Frobenius. This implies that $`H_I^3(R_p)=0`$ as well since in positive characteristic $`H_I^i(R_p)=lim\mathrm{Ext}^i(R_p/I^{[p^e]},R_p)`$. On the other hand, it is well known that in zero characteristic, $`H_I^3(R_{})`$ is not zero. Hence this provides an example where for the characteristic zero model we have $`\lambda _{0,3}=\lambda _{2,4}0`$ whereas in all positive characteristics $`\lambda _{0,3}=\lambda _{2,4}=0`$.
The next example even shows that the vanishing of $`\lambda _{a,i}`$ can vary in an arithmetic progression:
###### Example 3.2.
Let $`A`$ be the homogeneous coordinate ring of $`^1\times E`$ where $`E`$ is the elliptic curve $`\mathrm{Proj}\frac{[x,y,z]}{x^3+y^3+z^3}`$. Then $`A`$ is given as the quotient of $`R`$ (as above) by the ideal
$$I=(\delta _1,\delta _2,\delta _3,x^3+y^3+z^3,ux^2+vy^2+zw^2,u^2x+v^2y+w^2z,u^3+v^3+w^3).$$
The resolution of $`A=R/I`$ can be computed to be equal to
$$0\stackrel{}{}R\stackrel{}{}R^6\stackrel{}{}R^{11}\stackrel{}{}R^7\stackrel{}{}R\stackrel{}{}0$$
and one verifies that $`H_I^4(R_p)=0`$ if and only if $`\mathrm{char}k2mod3`$ (this essentially follows from the fact that depending on the modulus of $`p`$ mod $`3`$ the elliptic curve is supersingular or not). Hence we have that $`\lambda _{0,2}=\lambda _{2,3}=0`$ if and only if $`\mathrm{char}k2mod3`$.
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# Spectral state transitions in low-mass X-ray binaries - the effect of hard and soft irradiation
## 1 Introduction
Of great interest in black hole physics are the properties of the matter flow towards a compact object. Viscous differentially rotating gas flows determine the accretion processes around black holes, where viscous transport of angular momentum outward allows the gas to spiral inward towards the compact object. For a detailed discussion of the different possible modes of accretion see the work by Chen et al. (1995) and the review by Narayan et al. (1998).
The difference in observed spectra arising from the innermost accretion flow, can be used to study these accretion modes which are the same in stellar black holes or supermassive black holes in the centers of galaxies. Applications provide information for truncated thin disks in stellar mass systems as e.g. the black hole binary nova Muscae 1991 (Esin et al. 1997) as well as in low luminosity active galactic nuclei, e.g. M81 (Quataert et al. 1999). A rich documentation is present for accretion in low-mass X-ray binaries, containing a black hole or a neutron star primary (Chen et al. 1997, McClintock & Remillard 2005). The observations display the spectra of the X-ray radiation for different modes of accretion and correspondingly different spectral states: (1) a very hard spectrum (up to 100 keV) originating from the very energetic particles of a hot spherical advection-dominated flow (ADAF-type) (2) a soft spectrum (a few keV) radiated from the much cooler geometrically thin accretion disk. This is a simplified picture. In addition to the mentioned X-ray states of low and high luminosity also an “intermediate” state and a “very high” state can be found. McClintock & Remillard (2005) present an overview of the emission states of black-hole binaries where the combined X-ray and multi-frequency spectral characteristics of the sources are taken into account and a clearer definition of state names is suggested.
An interesting documentation of spectral changes are the transitions between hard and soft state observed for low-mass X-ray binaries, both neutron star and black hole systems (Tanaka & Shibazaki 1996), and also for high-mass X-ray binaries. The first observed change for a neutron star was found by Mitsuda et al. (1989) for 1608-522, and for a black hole system by Ebisawa et al. (1994), for GRS 1124-684, Nova Muscae 1991. A first modeling of the spectra was presented by Esin et al. (1997) using the concept of an inner advection-dominated accretion flow. Up to now the number of observations and the quality of spectra has very much increased, so that these data allow us to learn from the comparison of theoretical models with observations.
There is one special feature in the appearance of the transition of spectral states: considering an X-ray nova outburst light curve we see that the hard-soft transition does not occur at the same luminosity as the soft-hard transition. The latter happens at a luminosity lower by a factor of about 5. In recent work (Meyer-Hofmeister, Liu and Meyer 2005, hereafter MLM05) this surprising hysteresis could be explained as arising from the different amount of Compton cooling or heating acting on the accretion disk corona at the time of the state transition. During phases of low mass flow in the disk, as in quiescence and the early rise to outburst, the inner disk region is filled with an ADAF, the radiation is very hard and the hot photons lead to mainly Compton heating of the corona. If otherwise during the time around outburst maximum the mass flow rate in the disk is high, the thin disk reaches inward to the last stable orbit, the radiation is soft and the Compton effect acts as cooling of the corona. This difference in irradiation leads to a different vertical structure of the corona and a different amount of evaporation of gas from the disk.
In Sect.2 we summarize the observations for hysteresis. In Sect.3 we briefly describe the physics of our model for the interaction of disk and corona and the evaporation of gas from the cool disk. The computation of the vertical structure of the corona allows to study the Compton effect on the corona in detail and to evaluate the dependence on the hardness of the hard radiation. In Sect. 4 the results of the model calculations are shown. Since the irradiation from the inner region affects the evaporation rate, the accretion rate and the efficiency of radiation enter in the results. We compare with observations in Sect.5. The intermediate state between hard and soft spectral state is discussed in Sect.6, further in the discussion a possible effect of the inclination under which we observe a binary system in connection with flaring or warping of the disk. In Sect.8 the conclusions follow .
## 2 Observed hysteresis in lightcurves of X-ray binaries
Data for spectral state transitions in individual sources were already discussed in the first paper on hysteresis (MLM05). In Table 1 we summarize data on hysteresis now available in the literature. The values for the hysteresis given there are approximate values, either derived in the work referenced or deduced for the present paper from light curves and hardness ratios HR2 given in the literature (state transitions taken as occurring when HR2=(5-12keV)/(3-5 keV)= 1.5, McClintock & Remillard 2005). The existence of a hysteresis is well documented in the hardness-intensity diagrams but it is difficult to deduce the ratio between the luminosities at the transitions. The hardness used in these diagrams is different for different observations and it is difficult to deduce from this information the value of hysteresis. To deduce a reliable value for the hysteresis is difficult also due to the fact that the time of transition is not clearly documented, the hard/soft and even more the soft/hard transition take some time (as discussed in following sections). This means the values in Table 1 only show a range. We also refer to the description of spectral states in the review of McClintock & Remillard (2005, Sect.4.3): the low/hard (LH) and high/soft (HS) state correspond to our nomenclature hard and soft.
Historically the hysteresis effect was first pointed out by Mijamoto et al. (1995). A first detailed analysis of hysteresis in transient low-mass X-ray binaries was performed by Maccarone and Coppi (2003a), who discussed in detail the observations of Aql X-1 during the 1999 outburst, one of the best example, with two spectral transitions clearly documented at different count rates. For data and references of further interesting sources see Table 1. The source 4U 1543-47 (Park et al. 2004) was included because it has a classical X-ray nova outburst light curve. McClintock & Remillard (2005) show in their Fig. 4.2 a hard phase in the very beginning. Since the luminosity at outburst maximum is in the range of the Eddington luminosity the luminosity at the hard/soft transition, (expected at a few percent of this value) cannot be extracted from the data. The transition soft/hard is well documented.
### 2.1 The hardness of observed spectra
Since the hardness of the radiation from the inner region has an essential influence on the evaporation process the hardness at the time of the change from hard to soft state is of interest. But only a limited number of hard spectra are available for the sources listed in Table 1.
Maccarone & Coppi (2003b) analyzed data from the May/June outburst 1999 of Aql X-1 and present a spectrum taken shortly before transition to the soft state. This spectrum ideally documents the irradiation at the state transition. We evaluated the mean photon energy as about $`h\overline{\nu }`$ 90 keV. The true value might be higher depending on the poorly known spectrum at $`100`$ keV. In the work by Muno et al. (2002, Fig.1) hard and soft colors for different sources are shown. The hard color of 4U 1705-440 at the hard/soft state transition seems comparable or slightly less than that of Aql X-1.
Wardziński et al. (2002) analyzed hard spectra of GX 339-4 taken by the soft $`\gamma `$-ray OSSE detector on board CGRO simultaneous with Ginga and RXTE observations. The energy range reaches up to several hundred keV. The source flux at the time the spectra were taken was about 1 percent of the Eddington luminosity, therefore close to a value where we expect the state transition. (Further discussion of these spectra by Zdziarski et al., 1998 and Zdziarski & Gierliński, 2004). The evaluation of the spectrum from 1997 leads to $`h\overline{\nu }`$ 100 keV. The large errors at high energies cause quite an uncertainty in the result.
For XTE J1550-664 a hard spectrum (PCA + HEXTE data, 18-200 keV) taken 16 days before the hard/soft transition of the outburst in 2000 is shown by Rodriguez et al. (2003). We evaluate about 90 keV for the mean photon energy (assuming no significant contributions at higher energy). An averaged hard spectrum taken during the 2000 outburst (Arefiev et al. 2004) is less hard.
Rossi et al. (2004) present a study of the 2001/2002 outburst of the transient source XTE J1650-500. The beginning of the outburst was well covered, but the hard/soft transition discussed could already be the transition to the very high state. A spectrum taken right before this transition (Rossi, private communication 2005) is not as hard as the spectra from other sources discussed above.
## 3 The spectral state transitions
### 3.1 The general picture
For the analysis of the spectral state transitions in low-mass X-ray binaries we take the commonly accepted picture for the two accretion modes: (1) accretion via an optically thick geometrically thin cool disk (cool compared to a coronal flow) or (2) an optically thin spherically extended hot flow, ADAF-type two-temperature solution. Usually at larger distance from the compact object the accretion flow always is of the first type mentioned, closer to the central object both types are possible. The radiation is dominated by the flow in the innermost region, that is a soft spectrum arises from disk accretion, a hard spectrum from the hot flow of energetic particles. X-ray novae are good candidates to study changes between the two modes, a hard spectrum in quiescence during mass accumulation in the disk and a soft one after a dwarf-nova type instability had triggered an outburst (Meyer-Hofmeister & Meyer 1999). The importance of irradiation had already been pointed out by de Kool & Wrickramasinghe (1999).
Caused by the interaction of disk and corona gas evaporates from the disk into the corona and flows in form of a hot advection-dominated flow. Thereby the mass flow rate in the thin disk is diminished. The evaporation rate increases with decreasing distance from the compact object, but reaches a maximum at certain distance. This maximum rate determines the switch between the accretion modes in the inner region. Only if the mass flow in the outer disk is larger than this value will the disk “survive” this reduction in mass flow and continues inward. Otherwise the disk becomes truncated at a certain distance depending on the mass flow rate, and from then on all mass flows inward via the coronal/ADAF flow. This truncation of the inner disk during phases of low mass flow rate has long been recognized to be an essential feature of disk evolution (Mineshige et al. 1998) which also appeared in numerical simulations (Cannizzo 1998, 2000 and Dubus et al. 2001). A recent systematic analysis of Done & Gierliński (2004) using all data available from galactic binary systems on changes of spectra and truncation radii as a function of the accretion rate confirms this picture.
### 3.2 The effect of irradiation on state transitions
In the work by Meyer et al. (2000) the interaction of disk and corona was approximated by a one-zone model incorporating the standard equations of viscous hydrodynamics (see also Liu et al. 2002). As became clear in the recent work on hysteresis (MLM05) Compton cooling and heating by radiation from the innermost region acting on the vertically extended corona is an important process. The hysteresis is mainly caused by the Compton cooling in the soft state, but irradiation in the hard state also plays an important role
Evaporation has three important features: (a) the rates increase towards smaller distance $`R`$, (b) the rates have a maximum value at about several hundred Schwarzschild radii, and (c) from that distance on inward the coupling between electrons and ions becomes poor. Hard irradiation leads to higher mass evaporation rates, a higher maximal value, and therefore the spectral state transition at a higher luminosity than in the case of soft irradiation. (compare Fig.1 in MLM05). The harder the hard radiation is, the higher is the evaporation rate at maximum. Our recent computations were based on irradiation as hard as 100keV. Since the spectra observed for X-ray transients show quite a difference in hardness we take this into account for the evaluation of evaporation rates. As described in the foregoing paper (Sect. 4.3) the Compton cooling/heating rate per unit volume taken is the sum of Compton cooling and heating (inverse Compton and Compton effect).
$$q_{\mathrm{Comp}}=\frac{4kT_eh\overline{\nu }}{m_ec^2}n_e\sigma _Tcu,$$
(1)
with $`k`$ the Boltzmann constant, $`T_e`$ electron temperature, $`m_e`$ electron mass, $`c`$ velocity of light, $`n_e`$ electron particle density, $`\sigma _T`$ Thomson cross section and $`u`$ the energy density of the photon field. At transition the dominant contribution comes from photons from the central source, those from the disk underneath can be neglected here. Also Compton cooling by photons of the secondary stars is generally negligible.
For the flux from the central region in the soft state we now take a slightly different form compared to that in MLM05. We replace $`H/R`$ ($`H`$ scaleheight) by the term $`2\mathrm{cos}\delta =2z/(r^2+z^2)^{1/2}`$, where $`\delta `$ is the inclination angle under which the inner disk appears at height $`z`$. The formulae then is
$$F=\frac{L}{4\pi R^2}\frac{2z}{(r^2+z^2)^{1/2}},$$
(2)
where $`L`$ is the luminosity of the central source, which is related to the central mass accretion rate $`\dot{M}`$ as $`L=\eta \dot{M}c^2`$. The more detailed formula results in a less strong Compton effect than the earlier used expression and therefore a higher maximal evaporation rate.
## 4 Model calculations for hysteresis
In our computations we took a black hole mass $`M=6M_{}`$ (the results can be scaled for other masses and we show the results as measured in Eddington accretion rate $`\dot{M}_{\mathrm{Edd}}=L_{\mathrm{Edd}}/0.1c^2`$ with $`L_{\mathrm{Edd}}=4\pi GMc/\kappa `$, $`\kappa `$ electron scattering opacity and in Schwarzschild radius $`R_\mathrm{S}=2GM/c^2`$). The viscosity has a strong influence on the results as shown in earlier work (Meyer-Hofmeister & Meyer 2001 and Liu et al. 2002) and also pointed out by Różańska & Cerny (2000). We use the value $`\alpha `$=0.3 supported by modeling of X-ray binary spectra (Esin et al. 1997), and also application to accretion disk evolution (Meyer-Hofmeister & Meyer 1999).
### 4.1 Evaporation rates – dependence on hardness of irradiation and radiation efficiency
We determined the maximal evaporation rates (the rate determining the state transition) in the hard state for 30 to 120 keV mean photon energy of the radiation from the central source. The peak value has to be consistent with the accretion rate taken for the irradiation. For this consistency the efficiency $`\eta `$ with which radiation is produced from accretion enters. In the recent paper we had taken the value 0.05 for the efficiency in the hard state. In Fig. 1 we show the earlier results for 100 keV mean photon energy as dashed lines to illustrate the procedure to find the peak evaporation rate for the given efficiency $`\eta =0.05`$ or 0.025, marked by a filled or open dots. For the determination of each dot a series of evaporation rates at different distances $`R`$ from the center has to be evaluated to find the maximum (and the maximum has to be consistent with respect to the radiation produced in the central region). With this procedure we determined the maximal evaporation rates for different hardness and efficiency. The rates for $`\eta =0.05`$ and 0.025 are marked by filled and open dots. Along each solid curve the efficiency decreases towards zero, without irradiation (square).
The cross marks the peak evaporation rate for soft irradiation when the disk reaches inward to the last stable orbit. This is the case at the soft/hard transition. The difference between this value and the peak evaporation rate for hard radiation yields the hysteresis value.
### 4.2 Results for spectral state transitions
The peak evaporation rates for given hardness and efficiency of radiation provide the full picture of accretion rates which correspond to the spectral transitions from hard to soft state. Taking e.g. the radiation efficiency as 0.05, the filled dots in Fig. 1 give the truncation radii expected at the time of change from ADAF-type accretion to thin disk accretion. These radii are larger than those without Compton effect.
## 5 Comparison with observations
### 5.1 Hysteresis
In our section on observations (Sect.2) we pointed already out that hysteresis is clearly observed for several sources, but it is difficult to deduce the ratio of the luminosities corresponding to the two transition. There are several reasons. The hard/soft transition often is not well sampled, the spectral change takes some time (days to weeks), the exact time of transition is an interpretation of the data. A further uncertainty enters since different energy bands are used for the observations. The hardness-intensity diagrams presented by Corbel et al.(2004) and Homan & Belloni (2005) show the effect of the different hardness bands used for the hardness ratios (4.5-7.9)/(2.5-4.6)keV and (2-200)/(3-20)keV respectively. The observationally deduced ratios of luminosities at transition therefore only can show a range of values.
In Table 1 for GX 339-4 two hysteresis values are listed, 3.5 for the 1998 outburst and 10 for the 2002 outburst. The soft/hard transitions both took place at $`0.02L_{\mathrm{Edd}}`$, a typical value, the hard/soft transitions at $`0.07`$ (1998) and $`0.2L_{\mathrm{Edd}}`$ (in 2002), respectively (Zdziarski & Gierliński 2004). How could such a difference arise? In our picture of state transition this always should occur at the same luminosity and should not depend on the history as argued by Zdziarski & Gierliński (2004). A possible explanation could be that the hard/soft transition at $`0.2L_{\mathrm{Edd}}`$ was already a change to the very high state and the then earlier transition from the low/hard to the soft state had occurred at a lower luminosity before the observations began (comparable to the situation in XTE J1650, see Sect. 2.1).
Our model calculations yield a clear dependence on hardness of the irradiation. In Fig. 2 we show this dependence together with luminosity ratios for some sources. As discussed in Sect. 2.1 only a few spectra taken at the time of the hard/soft transition are available. The results in Fig. 2 show the wide scatter. The hardness might be underestimated if radiation at higher energy bands, not included, would be present. Recent work of Ling (2005) reports observations of the gamma-ray emitting source GRO J1719-24 at energy bands up to several hundred keV. Similarities with other black hole X-ray transients, GRO J0422+32, and with Cyg X-1 are mentioned and it is not clear whether sources discussed here could also have X-ray emission in the gamma-ray energy region. Then a larger hysteresis would be expected.
### 5.2 Truncation radii
The fact that the state transitions occur at different luminosities leads to a dichotomy in the truncation radii for this luminosity interval where either states can be realized: the soft state in which the disk reaches to the last stable orbit, or, the hard state in which the disk is truncated according to the mass-flow rate \- radius relation. The theoretically expected truncation radii during outburst rise and decline are shown in Fig. 2 of MLM05. Probably related to the truncation radii are the variations in power density spectra as discussed by Olive et al. (2003) for 4U 1705-44.
If we compare our results for the truncation radii obtained from the Compton effect on the coronal structure taken into account (MLM05 and this work) we find radii larger than expected from observations. Originally, the one-zone model for evaporation without including the Compton effect yielded a disk truncation at a few hundred Schwarzschild radii when the spectral transition occurs (Meyer et al.2000). In the paper by Yuan & Narayan (2004) the transition radii for several galaxies and galactic black hole sources are compared. Note that also sources in the hard state were included, for which the disk truncation at spectral state transition should be less or equal to the obtained value. The source XTE J1118 is a good example, the disk edge moves inward from values of $`10^4`$ to $`10^3R_\mathrm{S}`$ during quiescence to about 100 $`R_\mathrm{S}`$ during outburst (where the source remained in the hard state, as can occur in case of low accretion rates (Meyer-Hofmeister 2004)). Gilfanov & al. (2000) estimated using frequency-resolved spectroscopy of Cyg X-1 that the inner radius of the reflector, presumably an optically thick accretion disk, lies at about 100 $`R_\mathrm{S}`$ in the hard spectral state. Based on observations of Nova Mus 1991 Zycki & Done (1998) suggested that optically thick material within about 20-100 $`R_\mathrm{S}`$ is generally present in the hard (low) state. The question arises whether a left over inner disk is possible together with the ADAF, a problem also connected with the intermediate state (see next section).
There is one parameter in our analysis which influences the truncation radii, the heat conduction. In forthcoming work we will evaluate the effect of a reduced heat conduction on the interaction of disk and corona. A farther effect on the coronal structure and the location of the disk truncation might come from magnetic fields reaching into the corona.
## 6 The intermediate state
In the foregoing sections we discussed only two spectral states, hard and soft. The observations document an intermediate state which can persist for quite some time and sometimes does not even lead to a complete transition. Consistent with the explanation of the other two states this state would correspond to an inner region with a comparable amount of mass flow accreting in the two forms, via a hot extended corona and via a standard thin accretion disk. To understand how this can happen a closer look at the coronal evaporation model is necessary: On top of a standard thin accretion disk always a corona forms in interaction with the disk, and this corona itself carries a part of the accretion flow. Analysis of the formation of the coronal flow/ADAF (Meyer-Hofmeister & Meyer 2003) yielded how this coronal accretion flow increases with decreasing distance from the black hole until it reaches a maximum at a fixed distance between 100 and 1000 Schwarzschild radii. Inside of this distance either a pure coronal flow exists (if the disk became already fully evaporated further outward) or the thin disk continues. What happens to the coronal flow above such a disk continuing inward after the hot flow has passed the distance where it reached its maximum? According to the calculations a disk region farther in has a weaker corona and carries less coronal flow. Therefore the part that can no longer be carried inward must condense into the cool disk underneath (Liu et al. 2004).
This condensation process allows interesting aspects for the intermediate state, which we will discuss here and in forthcoming work. Can all the former coronal flow condense so that the accretion in the innermost region occurs completely through the cool disk? This would say that in the interior we have either a pure coronal flow in the hard state or a pure disk flow in the soft state with no room for the various forms of the observed intermediate state. We suggest here a solution for this apparent contradiction between the theoretical model and the observations. This solution relies on the strong temperature dependence of the collisional coupling between the electrons and ions in the corona. For a cool corona these two species are well coupled and efficiently exchange energy during the thermal timescale for heating and cooling of the corona. In this case frictional heat released in the ion component is readily transfered to the electrons and conducted to lower denser levels and radiated away. This keeps the corona in a cool dense state as one goes to smaller radii until it would formally disappear altogether yielding a perfect self consistent soft state in the interior. If however the coronal temperature in the beginning is high, coupling between electrons and ions is poor and the electrons can not remove the frictional heat of the ions which will then stay near virial temperatures. The farther one gets inward the hotter the ions become and the poorer is the coupling between ions and electrons. This then will allow an intermediate state in which a hot coronal flow passes above and below, on both sides of a regular accretion disk with not much interaction between the two accretion streams.
Now it is very important that the regular coronal structure calculations just yield a beginning weakening of the coupling between electrons and ions at the distance from the black hole where the maximum coronal flow rate occurs (compare the shape of the hatched evaporation rates curve in Fig.1). If the accretion rate is higher than this maximum the disk does not fully evaporate and the coronal flow condenses effectively into the cool disk further in. If however the accretion rate drops to a value lower than the maximal possible coronal flow rate the disk evaporates over some distance outside of the radius of the evaporation maximum and a gap is formed. From the outer border of the gap on through the whole gap region the accretion continues inward in form of an ADAF. ADAF temperatures are near virial and are higher than those of coronae above a cool accretion disk at the same distance from the central black hole. Thus where the ADAF reaches the inner boundary of the gap i.e. the inner disk, it has a rather high temperature. This temperature is higher the farther in the gap extends because the virial temperatures increases with $`1/R`$ as the radius $`R`$ decreases. Correspondingly poor becomes the coupling between the electrons and ions and the flow is therefore only partially capable to cool off and to condense back into the disk. Though this partial condensation will feed the disk and keep it stretching all the way inward another part of the accretion flow will stay in the hot near virial flow above and below the disk and on arriving in the interior will give the hard spectral component to this intermediate state, the soft component being provided by the remaining disk flow.
If the accretion rate is quite close to the critical value (that is the maximal evaporation rate) or above it the hard power law tail will only constitute a small fraction of the spectrum and with the fluctuating accretion rate vary significantly. This would then account for the observed hard power law tail in the soft state.
## 7 Discussion
A number of factors can complicate a direct comparison between theoretical model and observations.
In the soft state the luminosity is radiated from a thin innermost accretion disk. When this disk lies in the equatorial plane the observer sees it foreshortened. The ”observed” luminosity is then inclination dependent and must be corrected for this projection in order to obtain the true luminosity at the soft/hard transition.
Further, should the inner disk be tilted out of the equatorial plane and precess, an additional and time varying aspect of the accretion disk enters and would simulate an apparent variation of the transition luminosity even when the true transition luminosity remains the same. Tilted and precessing accretion disks might be indicated by periodic shifts in color-intensity plots of low mass X-ray binaries.
Recently Narayan and McClintock (2005) investigated the effects of the inclination angles of black hole X-ray binary systems on their observational properties. Within their sample of 20 sources they found none with an inclination angle larger than $`75^{}`$ and suggested that this absence of eclipsing sources is due to a flaring of the accretion disk by about $`15^{}`$, so that the disk permanently occults the X-ray source for the observer for all inclination angles that would allow us to see eclipses. Noisy light curves of sources with inclination between 70 and $`75^{}`$ support this picture and indicate the presence of partial and time varying coverage of the central source. Such a partial occultation could affect the determination of the luminosities at both hard/soft and soft/hard transitions.
A perhaps more remote possibility which would however have a direct impact on the coronal evaporation mechanism itself and thus on the hysteresis is a warping of the inner accretion disk. This would change the aspect under which the corona sees the irradiating innermost disk surface and thereby change the strength of the irradiation.
## 8 Conclusions
As one follows the light curve of an X-ray nova from early rise to late decline, observations show a profound change in the source spectrum from a hard to a soft one on the rise and back from a soft to hard one on the decline. These changes occur at characteristic luminosities but remarkably, they are not the same for the two transitions. This hysteresis in the light curve could be understood as arising from the different type of irradiation coming from the innermost region, which is hard at the hard/soft transition but soft at the soft/hard transition: This difference leads to a different Compton cooling or heating of the coronal layers and results in a different coronal density and mass flow. The latter determines whether a disk can be truncated (hard state) or continues all the way to the black hole (soft state)(MLM05).
Now we have carried out further analysis and show how this hysteresis depends on the hardness of the irradiation. Unfortunately, only a few spectra taken at the moment of state transition are available. Also the contributions at the important high energies are uncertain. Thus the observational basis for a comparison is still fairly sparse. In addition theoretical uncertainties arising from the rough ”one-zone” modeling and the choice of its parameters together with the assumed planar inner disk geometry will enter into comparisons between observations and theoretical models.
Thus to verify in detail the theoretically well understood dependence on the hardness of the hard irradiation in the comparison with the data is difficult. It is difficult to determine when and at which luminosity the transition actually happened. Different observations are taken in different energy bands. Since the change from one state to the other takes some time this clearly also demands a more detailed understanding of the intermediate state between the hard and the soft state. In our first discussion here we find a remarkable coincidence of the distance where the coronal flow rate reaches its maximum with that of the least thermal coupling between electrons and ions which could perhaps explain the features of an intermediate state.
In conclusion one might take the explanation of the hysteresis as a further strong support for the coronal evaporation model. It is surprising how this very simplified model appears to catch essential features of the accretion flows around compact objects. More observations of spectral state transitions in different systems, broad spectral energy coverage and refinement of the theoretical modeling in the future might prove very fruitful.
###### Acknowledgements.
One of the authors, BFL, thanks the Alexander-von-Humboldt Foundation for the award of a research fellowship during which this investigation started. We would like to thank Sabrina Rossi and collaborators for providing a spectrum of XTE J1650 at the time of transition during outburst rise.
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# Numerical simulations of stiff fluid gravitational singularities
## I Introduction
A longstanding problem in general relativity has been to find the general behavior of singularities. Several results, both analyticalalanreview and numericalbevreview have been obtained. Though most of the results are for the case where the spacetimes have one or more symmetries, recent work has been done on the general case where there are no symmetries.larsandalan ; harmonic ; dgprl There is a longstanding conjecturebkl due to Belinski, Lifschitz and Khalatnikov (BKL) that states that the generic singularity is spacelike and local. This conjecture has been reformulated and put more precisely by Uggla et al.jw The type of local dynamics conjectured by BKL depends on the type of matter. For vacuum and for many types of matter, the BKL conjecture is that the local dynamics is oscillatory, corresponding to the dynamics of a Bianchi type IX spacetime. However, for stiff fluid (i.e. fluid with pressure equal to energy density) the BKL conjecture is that the local dynamics is asymptotically velocity term dominated corresponding to the dynamics of a Bianchi type I spacetime. The vacuum version of the BKL conjecture has been supported by the numerical simulations of dgprl which show local and oscillatory dynamics in vacuum spacetimes with no symmetry. The stiff fluid version of the BKL conjecture has been supported by the theorem of Andersson and Rendalllarsandalan which shows the local existence in a neighborhood of the singularity of solutions of the Einstein equations with stiff fluid matter with the expected asymptotic behavior and with enough degrees of freedom to be the generic solutions. What is not known is whether generic stiff fluid initial data evolves to a solution of the Andersson and Rendall class.
To address this issue, we perform numerical simulations of the approach to the singularity for stiff fluid matter with no symmetries. Our methods are those of dgprl using the system of jw . Section II presents the equations and numerical methods used. The results are given in section III and conclusions in section IV.
## II Equations
The system evolved here is essentially that of referencejw but specialized to the stiff fluid case and with a slightly different choice of gauge. Here the spacetime is described in terms of a coordinate system ($`t,x^i`$) and a tetrad ($`𝐞_0,𝐞_\alpha `$) where both the spatial coordinate index $`i`$ and the spatial tetrad index $`\alpha `$ go from 1 to 3. It is assumed that $`𝐞_0`$ is hypersurface orthogonal and that the relation between tetrad and coordinates is of the form $`𝐞_0=N^1_t`$ and $`𝐞_\alpha =e_{\alpha }^{}{}_{}{}^{i}_i`$ Here $`N`$ is the lapse and we have chosen the shift to be zero. We choose the spatial frame $`\{𝐞_\alpha \}`$ to be Fermi propagated along the integral curves of $`𝐞_0`$. The commutators of the tetrad components are decomposed as follows:
$`[𝐞_0,𝐞_\alpha ]`$ $`=`$ $`\dot{u}_\alpha 𝐞_0(H\delta _{\alpha }^{}{}_{}{}^{\beta }+\sigma _{\alpha }^{}{}_{}{}^{\beta })𝐞_\beta `$ (1)
$`[𝐞_\alpha ,𝐞_\beta ]`$ $`=`$ $`(2a_{[\alpha }\delta _{\beta ]}^{}{}_{}{}^{\gamma }+ϵ_{\alpha \beta \delta }n^{\delta \gamma })𝐞_\gamma `$ (2)
where $`n^{\alpha \beta }`$ is symmetric, and $`\sigma ^{\alpha \beta }`$ is symmetric and trace free. Square brackets denote the antisymmetric part of a tensor.
Define $`u^ae_{0}^{}{}_{}{}^{a}`$ and $`h_{ab}g_{ab}+u_au_b`$, that is $`u^a`$ is the timelike vector of the tetrad and $`h_{ab}`$ is the spatial metric corresponding to the choice of $`u^a`$ as the time direction. Then the stress-energy tensor can be decomposed as
$$T_{ab}=\mu u_au_b+2q_{(a}u_{b)}+ph_{ab}+\pi _{ab}$$
(3)
where $`q_a`$ and $`\pi _{ab}`$ are orthogonal to $`u^a`$ and where $`\pi _{ab}`$ is symmetric and trace-free. Round brackets denote the symmetric part of a tensor.
Scale invariant variables are defined as follows:
$`\{_0,_\alpha \}\{𝐞_0,𝐞_\alpha \}/H`$ (4)
$`\{E_{\alpha }^{}{}_{}{}^{i},\mathrm{\Sigma }_{\alpha \beta },A^\alpha ,N_{\alpha \beta }\}\{e_{\alpha }^{}{}_{}{}^{i},\sigma _{\alpha \beta },a^\alpha ,n_{\alpha \beta }\}/H`$ (5)
$`q+1_0\mathrm{ln}H`$ (6)
$`r_\alpha _\alpha \mathrm{ln}H`$ (7)
$`\{\mathrm{\Omega },P,Q^\alpha ,\mathrm{\Pi }_{\alpha \beta }\}\{\mu ,p,q^\alpha ,\pi _{\alpha \beta }\}/(3H^2)`$ (8)
The matter variables are not all independent, because we assume that the stress-energy is that of a stiff fluid
$$T_{ab}=\stackrel{~}{\mu }\left(2\stackrel{~}{u}_a\stackrel{~}{u}_b+g_{ab}\right)$$
(9)
Here $`\stackrel{~}{\mu }`$ is the rest frame energy density of the fluid and $`\stackrel{~}{u}^a`$ is the fluid four-velocity, which can be decomposed as $`\stackrel{~}{u}^a=\mathrm{\Gamma }(u^a+v^a)`$ where $`v^a`$ is orthogonal to $`u^a`$. Comparison of equations (3) and (9) yields
$`Q_\alpha ={\displaystyle \frac{2\mathrm{\Omega }}{G_+}}v_\alpha `$ (10)
$`\mathrm{\Pi }_{\alpha \beta }={\displaystyle \frac{2\mathrm{\Omega }}{G_+}}v_{<\alpha }v_{\beta >}`$ (11)
$`P={\displaystyle \frac{\mathrm{\Omega }}{G_+}}\left(1\frac{1}{3}v^2\right)`$ (12)
Here $`v^2=v^\alpha v_\alpha `$ and $`G_+=1+v^2`$ and angle brackets denote the symmetric trace-free part of a tensor. Thus, all scale invariant matter variables can be expressed in terms of $`\mathrm{\Omega }`$ and $`v^\alpha `$.
Finally choose the lapse to be $`N=H^1`$. The relation between scale invariant frame derivatives and coordinate derivatives is $`_0=_t`$ and $`_\alpha =E_{\alpha }^{}{}_{}{}^{i}_i`$. From the Einstein field equations and the conservation of stress-energy one obtains the following evolution equations:
$`_tE_{\alpha }^{}{}_{}{}^{i}`$ $`=`$ $`F_{\alpha }^{}{}_{}{}^{\beta }E_{\beta }^{}{}_{}{}^{i}`$ (13)
$`_tr_\alpha `$ $`=`$ $`F_{\alpha }^{}{}_{}{}^{\beta }r_\beta +_\alpha q`$ (14)
$`_tA^\alpha `$ $`=`$ $`F_{}^{\alpha }{}_{\beta }{}^{}A^\beta +\frac{1}{2}_\beta \mathrm{\Sigma }^{\alpha \beta }`$ (15)
$`_t\mathrm{\Sigma }^{\alpha \beta }`$ $`=`$ $`(q2)\mathrm{\Sigma }^{\alpha \beta }2N_{}^{<\alpha }{}_{\gamma }{}^{}N^{\beta >\gamma }+N_{\gamma }^{}{}_{}{}^{\gamma }N^{<\alpha \beta >}`$ (16)
$`+`$ $`^{<\alpha }r^{\beta >}^{<\alpha }A^{\beta >}+2r^{<\alpha }A^{\beta >}`$
$`+`$ $`ϵ^{\gamma \delta <\alpha }(_\gamma 2A_\gamma )N_{}^{\beta >}{}_{\delta }{}^{}+3\mathrm{\Pi }^{\alpha \beta }`$
$`_tN^{\alpha \beta }`$ $`=`$ $`qN^{\alpha \beta }+2\mathrm{\Sigma }_{}^{(\alpha }{}_{\delta }{}^{}N^{\beta )\delta }ϵ^{\gamma \delta (\alpha }_\gamma \mathrm{\Sigma }_{}^{\beta )}{}_{\delta }{}^{}`$ (17)
$`_t\mathrm{\Omega }`$ $`=`$ $`(2q1)\mathrm{\Omega }3P_\alpha Q^\alpha +2Q^\alpha A_\alpha `$ (18)
$``$ $`\mathrm{\Pi }^{\alpha \beta }\mathrm{\Sigma }_{\alpha \beta }`$
$`_tv^\alpha `$ $`=`$ $`{\displaystyle \frac{G_+}{2G_{}\mathrm{\Omega }}}[(G_{}\delta _{}^{\alpha }{}_{\beta }{}^{}+2v^\alpha v_\beta )(_tQ^\beta 2[q+1]Q^\beta )`$ (19)
$``$ $`2v^\alpha (_t\mathrm{\Omega }2[q+1]\mathrm{\Omega })]`$
$`_tq`$ $`=`$ $`\left[2(q2)+\frac{1}{3}\left(2A^\alpha r^\alpha \right)_\alpha \frac{1}{3}^\alpha _\alpha \right]q`$ (20)
$``$ $`\frac{4}{3}_\alpha r^\alpha +\frac{8}{3}A^\alpha r_\alpha +\frac{2}{3}r_\beta _\alpha \mathrm{\Sigma }^{\alpha \beta }2\mathrm{\Sigma }^{\alpha \beta }W_{\alpha \beta }`$
$`+`$ $`{\displaystyle \frac{2}{G_+}}[2\mathrm{\Omega }\mathrm{\Sigma }^{\alpha \beta }v_\alpha v_\beta 2(q2)\mathrm{\Omega }+_t\mathrm{\Omega }`$
$``$ $`{\displaystyle \frac{2\mathrm{\Omega }}{G_+}}v_\alpha _tv^\alpha ]`$
Here we are using units where $`c=8\pi G=1`$. Furthermore the quantities $`G_{}`$, $`F_{\alpha \beta }`$, $`W_{\alpha \beta }`$ and $`_tQ_\alpha `$ are given by
$`G_{}`$ $``$ $`1v^2`$ (21)
$`F_{\alpha \beta }`$ $``$ $`q\delta _{\alpha \beta }\mathrm{\Sigma }_{\alpha \beta }`$ (22)
$`W_{\alpha \beta }`$ $``$ $`\frac{2}{3}N_{\alpha \gamma }N_{\beta }^{}{}_{}{}^{\gamma }\frac{1}{3}N_{}^{\gamma }{}_{\gamma }{}^{}N_{\alpha \beta }+\frac{1}{3}_\alpha A_\beta `$ (23)
$``$ $`\frac{2}{3}_\alpha r_\beta \frac{1}{3}ϵ_{}^{\gamma \delta }{}_{\alpha }{}^{}\left(_\gamma 2A_\gamma \right)N_{\beta \delta }`$
$`_tQ_\alpha `$ $`=`$ $`2(q1)Q_\alpha \mathrm{\Sigma }_{\alpha \beta }Q^\beta _\alpha P^\beta \mathrm{\Pi }_{\alpha \beta }`$ (24)
$`+`$ $`(P\mathrm{\Omega })r_\alpha +\mathrm{\Pi }_{\alpha \beta }(3A^\beta +r^\beta )`$
$`+`$ $`ϵ_{\alpha \beta \gamma }N^{\beta \delta }\mathrm{\Pi }_{\delta }^{}{}_{}{}^{\gamma }`$
In addition to the evolution equations, the variables satisfy constraint equations as follows:
$`0`$ $`=`$ $`(𝒞_{\mathrm{com}})^{\gamma i}ϵ^{\alpha \beta \gamma }\left(_\alpha E_{\beta }^{}{}_{}{}^{i}[r_\alpha +A_\alpha ]E_{\beta }^{}{}_{}{}^{i}\right)`$ (25)
$``$ $`N^{\gamma \alpha }E_{\alpha }^{}{}_{}{}^{i}`$
$`0`$ $`=`$ $`𝒞_\mathrm{G}1+\frac{1}{3}(2_\alpha 2r_\alpha 3A_\alpha )A^\alpha \frac{1}{6}N_{\alpha \beta }N^{\alpha \beta }`$ (26)
$`+`$ $`\frac{1}{12}(N_{}^{\alpha }{}_{\alpha }{}^{})^2\frac{1}{6}\mathrm{\Sigma }_{\alpha \beta }\mathrm{\Sigma }^{\alpha \beta }\mathrm{\Omega }`$
$`0`$ $`=`$ $`(𝒞_\mathrm{C})^\alpha _\beta \mathrm{\Sigma }^{\alpha \beta }+2r^\alpha \mathrm{\Sigma }_{}^{\alpha }{}_{\beta }{}^{}r^\beta 3A_\beta \mathrm{\Sigma }^{\alpha \beta }`$ (27)
$``$ $`ϵ^{\alpha \beta \gamma }N_{\beta \delta }\mathrm{\Sigma }_{\gamma }^{}{}_{}{}^{\delta }+3Q^\alpha `$
$`0`$ $`=`$ $`𝒞_qq\frac{1}{3}\mathrm{\Sigma }^{\alpha \beta }\mathrm{\Sigma }_{\alpha \beta }+\frac{1}{3}_\alpha r^\alpha \frac{2}{3}A_\alpha r^\alpha `$ (28)
$``$ $`\frac{1}{2}(\mathrm{\Omega }+3P)`$
$`0`$ $`=`$ $`(𝒞_\mathrm{J})^\alpha _\beta N^{\alpha \beta }(r_\beta +2A_\beta )N^{\alpha \beta }`$ (29)
$`+`$ $`ϵ^{\alpha \beta \gamma }(_\beta A_\gamma r_\beta A_\gamma )`$
$`0`$ $`=`$ $`(𝒞_\mathrm{W})^\alpha =ϵ^{\alpha \beta \gamma }\left(_\beta r_\gamma A_\beta r_\gamma \right)N^{\alpha \beta }r_\beta `$ (30)
We want a class of initial data satisfying these constraints that is general enough for our purposes but simple enough to find numerically. Recall that on an initial data surface, the spatial metric $`h_{ij}`$ and extrinsic curvature $`K^{ij}`$ must satisfy the constraint equations
$`D_i(K^{ij}Kh^{ij})=q^j`$ (31)
$`R+K^2K^{ij}K_{ij}=2\mu `$ (32)
Here $`D_i`$ and $`R`$ are respectively the derivative operator and scalar curvature associated with $`h_{ij}`$ and $`\mu `$ and $`q_i`$ are the components of the stress-energy tensor given in equation (3). We use the York methodjimmy which begins by defining the quantities $`\overline{h}_{ij}`$ and $`\overline{A}^{ij}`$ by
$`h_{ij}=\psi ^4\overline{h}_{ij}`$ (33)
$`K^{ij}\frac{1}{3}Kh^{ij}=\psi ^{10}\overline{A}^{ij}`$ (34)
We choose $`K`$ to be constant, $`\overline{h}_{ij}`$ to be the flat metric $`\delta _{ij}`$ and $`q^i`$ to vanish. With these choices, equations (31) and (32) become
$`_i\overline{A}^{ij}`$ $`=`$ $`0`$ (35)
$`_i^i\psi `$ $`=`$ $`\left(\frac{1}{12}K^2\frac{1}{4}\mu \right)\psi ^5\frac{1}{8}\overline{A}^{ij}\overline{A}_{ij}\psi ^7`$ (36)
Here $`_i`$ is the ordinary derivative with respect to Cartesian coordinates and indicies are raised and lowered with $`\delta _{ij}`$.
We choose space to have topology $`T^3`$ with the Cartesian coordinates $`x,y`$ and $`z`$ each going from $`0`$ to $`2\pi `$. We choose the following solution of equation (35)
$`\overline{A}^{11}`$ $`=`$ $`a_2\mathrm{cos}y+a_3\mathrm{cos}z+b_2+b_3`$
$`\overline{A}^{22}`$ $`=`$ $`a_1\mathrm{cos}xa_3\mathrm{cos}z+b_1b_3`$
$`\overline{A}^{33}`$ $`=`$ $`a_1\mathrm{cos}xa_2\mathrm{cos}yb_1b_2`$ (37)
with the off-diagonal components of $`\overline{A}^{ij}`$ vanishing. Here the $`a_i`$ and $`b_i`$ are constants. Note that due to the periodicity of the coordinates and the linearity of equation (35) the general solution of equation (35) is a Fourier series. The solution that we choose is then essentially the simplest solution of equation (35) without symmetries. The quantity $`\mu `$ can be freely specified and we choose it to be
$$\mu =c_0+c_1\mathrm{cos}x+c_2\mathrm{cos}y+c_3\mathrm{cos}z$$
(38)
where the $`c_i`$ are constants. With these choices for $`\overline{A}^{ij}`$ and $`\mu `$, equation (36) is solved numerically (in a manner to be described later) to yield $`\psi `$ and therefore $`h_{ij}`$ and $`K^{ij}`$.
From this initial data, we must then produce the initial values of the scale invariant variables. From equation (1) it follows that $`H=K/3`$ and since $`H`$ is constant it then follows that $`r_\alpha `$ vanishes. Since the initial spatial metric is conformally flat, we choose the initial spatial tetrad vectors by
$$E_{\alpha }^{}{}_{}{}^{i}=H^1\psi ^2\delta _{\alpha }^{}{}_{}{}^{i}$$
(39)
It then follows from equation (2) that $`N_{\alpha \beta }`$ vanishes and that
$$A_\alpha =2\psi ^1_\alpha \psi $$
(40)
From equation (1) it then follows that
$$\mathrm{\Sigma }^{\alpha \beta }=H^1\psi ^6\delta _{}^{\alpha }{}_{i}{}^{}\delta _{}^{\beta }{}_{j}{}^{}\overline{A}^{ij}$$
(41)
while $`\mathrm{\Omega }`$ is given by equation (8) and $`q`$ by the vanishing of equation (28).
The numerical method used is as follows: each spatial direction corresponds to $`n+2`$ grid points with spacing $`dx=2\pi /n`$. The variables on grid points $`2`$ to $`n+1`$ are evolved using the evolution equations, while at points $`1`$ and $`n+2`$ periodic boundary conditions are imposed. The initial data is determined once equation (36) is solved. This is done iteratively as follows:jim Define
$$S(\psi )2\psi +\left(\frac{1}{12}K^2\frac{1}{4}\mu \right)\psi ^5\frac{1}{8}\overline{A}^{ij}\overline{A}_{ij}\psi ^7$$
(42)
Then equation (36) takes the form $`^i_i\psi 2\psi =S(\psi )`$. We make an initial guess $`\psi ^0`$ for $`\psi `$ and solve using the conjugate gradient method Saul the equation
$$^i_i\psi ^{k+1}2\psi ^{k+1}=S(\psi ^k)$$
(43)
iterating until $`\psi ^k`$ satisfies equation (36) to within a set tolerance.
The evolution proceeds using equations (13-20) with the exception that the term $`(52q)𝒞_q`$ is added to the right hand side of equation (20) to prevent the growth of constraint violating modes and the term $`0.6(𝒞_\mathrm{C})^\alpha `$ is added to the right hand side of equation (15) to make the system well posed.meandcarsten Spatial derivatives are evaluated using centered differences, and the evolution is done using a three step iterated Crank-Nicholson methodICN (a type of predictor-corrector method). In equation (20) the highest spatial derivative term is $`\frac{1}{3}^\alpha _\alpha q`$ which gives this equation the form of a diffusion equation. Note that diffusion equations can only be evolved in one direction in time, in this case the negative direction which corresponds to the approach to the singularity. Stability of numerical evolution of diffusion equations generally requires a time step proportional to the square of the spatial step. However, the constant of proportionality depends on the coefficient of the second spatial derivative. To ensure stability, we define $`E_{\mathrm{max}}`$ to be the maximum value of $`|E_{\alpha }^{}{}_{}{}^{i}|`$ (over all space and over all $`\alpha `$ and $`i`$) and then define $`dt_1\frac{1}{4}(dx/E_{\mathrm{max}})^2`$ and $`dt_2\frac{1}{8}dx`$. The time step $`dt`$ is then chosen to be whichever of $`dt_1`$ and $`dt_2`$ has the smaller magnitude.
Before presenting numerical results, it is helpful to consider what behavior to expect as the singularity is approached (that is as $`t\mathrm{}`$). First denote the eigenvalues of $`\mathrm{\Sigma }_{}^{\alpha }{}_{\beta }{}^{}`$ by ($`\mathrm{\Sigma }_1,\mathrm{\Sigma }_2,\mathrm{\Sigma }_3`$). Then suppose that at sufficiently early times the time averages of $`q\mathrm{\Sigma }_i`$ are all positive. Then the time averages of the eigenvalues of $`F_{}^{\alpha }{}_{\beta }{}^{}`$ are all positive. Since we are evolving in the negative time direction, this should lead (through equation (13)) to an exponential decrease in $`E_{\alpha }^{}{}_{}{}^{i}`$. However, since all spatial derivatives appear in the equations in the form $`_\alpha =E_{\alpha }^{}{}_{}{}^{i}_i`$ we would expect the spatial derivatives to become negligible. That is, the approach to the singularity is local. Furthermore, this positivity of the time averages of the eigenvalues of $`F_{}^{\alpha }{}_{\beta }{}^{}`$ should lead (through equations (14-15)) to exponential decrease in $`r_\alpha `$ and $`A^\alpha `$. A similar argument applied to equations (18) and (24) and using equations (26) and (28) indicates that as the singularity is approached $`Q_\alpha `$ should become negligible, but $`\mathrm{\Omega }`$ should not, and therefore that $`v_\alpha `$ should become negligible. Thus, as the singularity is approached, the system should be well described by a simplified set of evolution and constraint equations where spatial derivatives as well as $`r_\alpha ,A^\alpha `$ and $`v_\alpha `$ are negligible. Note that the fact that spatial derivatives are becoming negligible does $`\mathrm{𝑛𝑜𝑡}`$ mean that the spacetime is becoming homogeneous. Rather the considerable spatial variation is becoming a negligible part of the equations of motion since all spatial derivatives appear multiplied by $`E_{\alpha }^{}{}_{}{}^{i}`$ which is becoming negligible. We now write down this simplified system of evolution and constraint equations where all these terms are neglected. In this limit equation (27) implies that the matricies $`\mathrm{\Sigma }_{}^{\alpha }{}_{\beta }{}^{}`$ and $`N_{}^{\alpha }{}_{\beta }{}^{}`$ commute and therefore have the same eigenvalues. The evolution equations for $`\mathrm{\Sigma }^{\alpha \beta }`$ and $`N^{\alpha \beta }`$ can then be written as evolution equations for their eigenvalues. The non-trivial evolution and constraint equations then become in this limit
$`_t\mathrm{\Sigma }_i`$ $`=`$ $`(q2)\mathrm{\Sigma }_i2N_i^2+\left({\displaystyle \underset{k=1}{\overset{3}{}}}N_k\right)N_i+\frac{2}{3}Y`$ (44)
$`_tN_i`$ $`=`$ $`(q+2\mathrm{\Sigma }_i)N_i`$ (45)
$`_t\mathrm{\Omega }`$ $`=`$ $`2(q2)\mathrm{\Omega }`$ (46)
$`0`$ $`=`$ $`Y+{\displaystyle \underset{k=1}{\overset{3}{}}}\mathrm{\Sigma }_k^2+6\mathrm{\Omega }6`$ (47)
$`0`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{3}{}}}\mathrm{\Sigma }_k^2+6\mathrm{\Omega }3q`$ (48)
Here $`\mathrm{\Sigma }_i`$ and $`N_i`$ are the eigenvalues of $`\mathrm{\Sigma }_{}^{\alpha }{}_{\beta }{}^{}`$ and $`N_{}^{\alpha }{}_{\beta }{}^{}`$ respectively and $`Y`$ is given by
$$Y=\underset{k=1}{\overset{3}{}}N_k^2\frac{1}{2}\left(\underset{k=1}{\overset{3}{}}N_k\right)^2$$
(49)
and indicies are not summed over unless explicitly indicated.
Suppose that the dynamics is in a period (called a Kasner epoch) when all the $`N_i`$ are negligibly small. Then it follows from equations (47) and (48) that $`q=2`$ and therefore, from equations (44) and (46) that $`\mathrm{\Omega }`$ and the $`\mathrm{\Sigma }_i`$ are constant. From equations (47) and (45) it follows that there are two possibilities for a Kasner epoch: (i) all the $`\mathrm{\Sigma }_i`$ are $`1`$ in which case the $`N_i`$ remain negligible and the Kasner epoch lasts all the way to the singularity. (ii) one of the $`\mathrm{\Sigma }_i`$ is $`<1`$ in which case the corresponding $`N_i`$ grows until it is large enough to bring the Kasner epoch to an end. We now look in more detail at possibility (ii). Let $`\mathrm{\Sigma }_1`$ be the $`\mathrm{\Sigma }_i`$ that is $`<1`$. Then $`N_1`$ is the $`N_i`$ that is growing. We are therefore led to examine equations (44-48) neglecting $`N_2`$ and $`N_3`$ but not $`N_1`$. In this regime equations (44) and (45) become
$`_t\mathrm{\Sigma }_1=S^2(\mathrm{\Sigma }_1+4)`$ (50)
$`_t\mathrm{\Sigma }_2=S^2(\mathrm{\Sigma }_22)`$ (51)
$`_t\mathrm{\Sigma }_3=S^2(\mathrm{\Sigma }_32)`$ (52)
$`_t\mathrm{\Omega }=2S^2\mathrm{\Omega }`$ (53)
where $`SN_1/\sqrt{6}`$. Define $`Z4+\mathrm{\Sigma }_1`$. Then equation (50) becomes $`_tZ=S^2Z`$ while from equations (51) and (52) it follows that there are constants $`c_2`$ and $`c_3`$ with $`c_2+c_3=1`$ such that $`\mathrm{\Sigma }_2=2+c_2Z`$ and $`\mathrm{\Sigma }_3=2+c_3Z`$. Similarly, it follows from equation (53) that there is a constant $`c_4`$ such that $`\mathrm{\Omega }=c_4Z^2`$. Finally, it then follows from equation (47) that $`Z`$ satisfies the equation of motion
$$_tZ=\left[\frac{1}{6}(4\alpha ^2)Z^24Z+6\right]Z$$
(54)
where $`\alpha ^2=1[12c_4+(c_2c_3)^2]`$. Note that the quantity in square brackets vanishes at $`Z_+`$ and $`Z_{}`$ where $`Z_\pm =6/(2\alpha )`$. Therefore the dynamics is a “bounce” from a Kasner epoch corresponding to $`Z_{}`$ to one corresponding to $`Z_+`$. Use a minus subscript to denote a quantity before the bounce and a plus subscript to denote a quantity after the bounce. We then have $`\mathrm{\Omega }_+/\mathrm{\Omega }_{}=(Z_+/Z_{})^2`$. However, from the definition of $`Z`$ it follows that $`Z_{}=4+\mathrm{\Sigma }_1`$ while from the definition of $`Z_\pm `$ it follows that $`(1/Z_+)+(1/Z_{})=2/3`$. We then find the following “bounce rule” relating a quantity after the bounce to quantities before the bounce.
$$\mathrm{\Omega }_+=\mathrm{\Omega }_{}\left(\frac{3}{5+2\mathrm{\Sigma }_1}\right)^2$$
(55)
Note that from the bounce rule it follows that $`\mathrm{\Omega }`$ increases at each bounce. Furthermore, it follows from equation (47) (and from the fact that $`\mathrm{\Sigma }_{\alpha \beta }`$ is trace-free) that the minimum possible value for a $`\mathrm{\Sigma }_i`$ during a Kasner epoch is $`2\sqrt{1\mathrm{\Omega }}`$ and therefore that no further bounces can happen once $`\mathrm{\Omega }>3/4`$ (though bounces may cease at lower values of $`\mathrm{\Omega }`$). Thus we expect that at each spatial point there is a last bounce followed by a Kasner epoch that lasts all the way to the singularity. In other words, we expect the approach to the singularity to be of the Andersson and Rendall class.
## III results
All runs were done in double precision on a SunBlade 2000 with $`n=50`$ (except for a convergence test which also used $`n=25`$). The equations were evolved from $`t=0`$ to $`t=90`$. For the initial data, the trace of the extrinsic curvature was $`1`$ corresponding to an initial value of $`1/3`$ for $`H`$. The constants $`a_i,b_i`$ and $`c_i`$ characterizing the initial data were
$`a_i=(0.2,0.1,0.04)`$ (56)
$`b_i=(1.7,0.1,0)`$ (57)
$`c_i=(0.005,0.005,0.005)`$ (58)
and the constant $`c_0`$ was $`0.02`$.
We would like to know whether $`E_{\alpha }^{}{}_{}{}^{i},r_\alpha ,A^\alpha `$ and $`v^\alpha `$ become negligible as the singularity is approached. In figure 1 are plotted the maximum values (over all space, $`\alpha `$ and $`i`$) of $`\mathrm{ln}|E_{\alpha }^{}{}_{}{}^{i}|,\mathrm{ln}|r_\alpha |,\mathrm{ln}|A^\alpha |`$ and $`\mathrm{ln}|v^\alpha |`$. Note that after a certain amount of evolution, all these quantities steadily decrease. This indicates that after a certain amount of time spatial derivatives become negligible in the eqautions of motion and that the approximation considered at the end of the previous section becomes valid.
It then follows that the interesting part of the dynamics is the development of the variables at a single point as a function of time. We now present data of that form. The behavior at the spatial point chosen is typical.
The results of a convergence test are plotted in figure 2. Here what is plotted is $`4𝒞_q`$ for the $`n=50`$ run (solid line) and $`𝒞_q`$ for the $`n=25`$ run (dotted line). Both quantities are plotted vs $`t`$. Note that the two curves roughly agree in magnitude but become out of sync in time. This indicates second order convergence but with the system having sensitive dependence on initial conditions. Similar results were obtained for the other constraints.
Figures 3 and 4 show respectively the diagonal components of $`\mathrm{\Sigma }_{\alpha \beta }`$ and $`N_{\alpha \beta }`$ in the asymptotic frame, i.e. the frame of the eigenvectors that $`\mathrm{\Sigma }_{}^{\alpha }{}_{\beta }{}^{}`$ has at the end of the evolution. For that part of the evolution where the approximation made at the end of the previous section is valid, these diagonal components are the eigenvalues of $`\mathrm{\Sigma }_{}^{\alpha }{}_{\beta }{}^{}`$ and $`N_{}^{\alpha }{}_{\beta }{}^{}`$ respectively. Note that the dynamics of the eigenvalues of $`\mathrm{\Sigma }_{}^{\alpha }{}_{\beta }{}^{}`$ consists of epochs where they are apporoximately constant (Kasner epochs) punctuated by short bounces. Furthermore the components of $`N_{\alpha \beta }`$ are negligible except at the bounces. Also note that there is a last bounce and that this coincides with the most negative eigenvalue of $`\mathrm{\Sigma }_{}^{\alpha }{}_{\beta }{}^{}`$ becoming greater than $`1`$.
In figure 5 is plotted $`\mathrm{\Omega }`$ vs $`t`$. Note that the behavior of $`\mathrm{\Omega }`$ is a sequence of constant values that are punctuated by short bounces and that the bounces in $`\mathrm{\Omega }`$ occur at the same times as the bounces in $`\mathrm{\Sigma }_{\alpha \beta }`$. The sequence of the values of $`\mathrm{\Omega }`$ is 0.05956, 0.1607, 0.4139, 0.6231, while the corresponding values of the most negative eigenvalues of $`\mathrm{\Sigma }_{}^{\alpha }{}_{\beta }{}^{}`$ are -1.583, -1.565, -1.277, -0.6596. These values obey the “bounce rule” of equation (55).
## IV conclusions
These simulations support the expected picture for the approach to the generic singularity in a spacetime where the matter is a stiff fluid. As the singularity is approached the terms in the equations of motion involving spatial derivatives become negligible. The dynamics at each spatial point consists of a sequence of Kasner epochs punctuated by short bounces. The sequence of values of $`\mathrm{\Omega }`$ obeys the expected bounce rule. There is a last bounce, after which the dynamics is described by a single Kasner epoch all the way to the singularity, thus yielding a spacetime in the class of reference larsandalan .
What remains to be done is to treat the approach to the singularity for non-stiff fluids. Here the behavior that is expected is quite different. The BKL conjecture is that the matter will become negligible and the dynamics of the gravitational variables as the singularity is approached will be that of vacuum spacetimes. The formalism of jw includes a class of non-stiff fluids, and the resulting equations are similar to those of the stiff fluid case. Nonetheless, the numerical methods of this paper are not adequate to treat the case of non-stiff fluids. That is because in non-stiff fluids shock waves form, while the numerical methods of the present paper are appropriate for smooth solutions. A shock capturing method would be appropriate to treat the non-stiff fluid case.
## V acknowledgements
This work was partially supported by a grant from the National Science and Engineering Research Council of Canada and by NSF grant PHY-0456655 to Oakland University.
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# Hard X-ray diffuse emission from the Galactic Center seen by INTEGRAL.
## 1 Introduction
Recent observations of the Galactic Center (GC) in 1–10 keV (Baganoff et al. (2001)), 20-100 keV (Bélanger et al. (2004)), 10 MeV–10 GeV (Mayer-Hasselwander et al. (1998)) and 1–10 TeV (Aharonian et al. (2004)) reveal the high-energy activity of the nucleus of the Milky Way. This activity is powered by a supermassive black hole (BH) of mass $`M_{BH}3\times 10^6M_{}`$ (Genzel et al. (2000); Ghez et al. (2000); Schödel et al. (2002)) and has a number of puzzling properties such as an extremely low luminosity, $`L_{BH}10^{36}`$ erg/s (eight orders of magnitude less than the Eddington luminosity). A number of theoretical models of supermassive BHs accreting at low rates were put forward to explain the data (Narayan et al. (2002); Melia and Falcke (2001); Aharonian and Neronov (2005)) but the “broad-band” picture of activity of the GC is still missing.
Emission from the GC has two major contributions. The first one is the variable emission from the direct vicinity of the central BH (the source Sgr A\*) detected in X-rays (Baganoff et al. (2001); Porquet et al. (2003); Goldwurm et al. (2003)) and in the infrared (Genzel et al. (2003)). The typical variability time scales $`T1`$ ksec are close to the light-crossing time of a compact region of the size of about ten gravitational radii of the central supermassive BH, $`R10R_{grav}10^{13}`$ cm. Apart from the highly variable emission from a compact object, there is a strong diffuse component which extends over tens of parsecs around the compact source. This diffuse X-ray emission has complex spatial and spectral properties (Muno et al. (2004)). In particular, it contains an extended “hard” component which is tentatively explained by the presence of hot plasma with temperature $`T8`$ keV in the emission region. However, such an explanation faces serious problems because it is difficult to find objects which would produce such a hot plasma and, moreover, this plasma would not be gravitationally bound in the GC. The origin of the hard component of diffuse X-ray emission can be clarified using the data on diffuse emission in higher energy bands (above 10 keV). If the hard diffuse emission is really of thermal origin, one expects to see a sharp cut-off in the spectrum above 10 keV.
The GC was recently detected in the 20-100 keV energy band by the INTEGRAL satellite (Bélanger et al. (2004)). The source IGR J1745.6-2901 is found to be coincident with Sgr A\* to within $`1`$ arcmin. However, the wide point spread function of ISGRI imager on board of INTEGRAL encircles the whole region of size $`20`$ pc around the BH which does not allow to separate the contributions of Sgr A\* itself and of the possible extended emission around Sgr A\*. The only possibility of separation of the two contributions is to study the time variability of the signal. Indeed, one expects to see a highly variable signal, if the 20-100 keV emission from the GC is mostly from the Sgr A\*. Moreover, a hint of the variability in this energy band (a possible 40-min flare) was reported in (Bélanger et al. (2004)). At the same time, the data sample analyzed in (Bélanger et al. (2004)) was too small to exclude that the “flare” is a statistical fluctuation of the signal (the analysis presented below shows that the latter is the case). The question of variability of the INTEGRAL GC source was addressed recently by (Goldwurm et al. (2004)), where the absence of significant variability of the source at the time scales of 1 day and 1 year was reported and the previous detection of the “flare” by (Bélanger et al. (2004)) was attributed to a “background feature”. It is clear that in order to find whether the signal detected by INTEGRAL is from the supermassive black hole itself or from an extended region around the GC, it is important to study systematically the variability of the source at the most relevant time scale of $`1`$ ksec (the dynamical time scale near the black hole horizon). Such a systematic analysis should allow to separate the “artificial” or “background” variability which is specific to the coded mask instruments (see Section 2) from the true variability of the source.
In what follows we develop a systematic approach to the detection of variability of INTEGRAL sources in the “crowded” fields, which contain many sources in the field of view (Section 2). We apply the method to the $`1^{}\times 2^{}`$ sky region around the GC (Section 3). This region contains six hard X-ray sources previously detected by INTEGRAL (Bélanger et al. (2004)). We find that the source IGR J1745.6-2901 coincident with the Galactic Center is not variable in the 20-100 keV energy band. In Section 4 we analyze the spectrum of non-variable hard X-ray emission from the GC in more details, using the INTEGRAL and XMM-Newton data. We find that the normalization of the XMM-Newton spectrum in 1-10 keV band matches the one of the INTEGRAL spectrum only if the XMM-Newton spectrum is extracted from an extended region with the radius of $`6^{}`$ around the GC position. This indicates that the size of the hard X-ray emission region is $`1020`$ pc. In Section 5 we study the physical mechanism behind the non-variable hard X-ray emission from the GC. Taking into account very special physical conditions in the GC region, we show that the only viable mechanism is the synchrotron emission from electrons of energies of 10-100 TeV. We find that inverse Compton and bremsstrahlung emission from such electrons are at the level of the detected TeV flux from the GC.
## 2 Method for the detection of variability
In principle, the variability of INTEGRAL sources can be analyzed in a standard way studying (in)consistency of the detected signal with the one expected from a non-variable source. The INTEGRAL data are naturally organized by pointings of duration of $`13`$ ksec in average (Science Windows, ScW). The simplest way to detect the variability of a source on the ksec time scale is therefore to analyze the evolution of the flux from the source on ScW-by-ScW basis. If the flux from a source in the $`i`$-th ScW is $`F_i`$ and the flux variance is $`\sigma _i`$, the $`\chi ^2`$ of the fit by a constant is $`\chi ^2=_{i=1}^N(F_i\overline{F})^2/\sigma _i^2`$ (here $`\overline{F}=(_{i=1}^NF_i/\sigma _i^2)/(_{i=1}^N1/\sigma _i^2)`$ is the weighted average flux and $`N`$ is the total number of ScWs). The probability of obtaining a given $`\chi ^2`$ under the assumption that the source is non-variable is $`𝒫=1P(N1,\chi ^2)`$ (here $`P(a,x)`$ is the incomplete gamma-function).
However, in reality, if one applies the above method to the analysis of of “crowded” fields which contain many sources, one would obtain a surprising (wrong) result that all the detected sources are highly variable! The reason for this is that INTEGRAL is a coded mask instrument. In the coded mask instruments each source casts a shadow of the mask on the detector plane. Knowing the position of the shadow one can reconstruct the position of the source on the sky. If there are several sources in the field of view, each of them produces a shadow which is spread over the whole detector plane. Some detector pixels are illuminated by more than one source. If the signal in a detector pixel is variable, one can tell only with certain probability, which of the sources illuminating this pixel is responsible for the variable signal. Thus, in a coded mask instrument, the presence of bright variable sources in the field of view introduces artificial variability for almost all the other sources in the field of view. Moreover, since the overlap between the shadowgrams of the bright variable source and of the sources at different positions on the sky varies with the position on the sky, one can not know “in advance” what is the level of “artificial” variability in a given region of the deconvolved sky image.
In order to overcome this difficulty, one has to measure the variability of the flux not only directly in the sky pixels at the position of the source of interest, but also in the background pixels around the source. Obviously, the “artificial” variability introduced by the nearby bright sources is the same in the background pixels and in the pixel(s) at the source position. Practically, one can calculate the $`\chi ^2`$ for each pixel of the sky image and compare the values of $`\chi ^2`$ at the position of the source of interest to the mean values of $`\chi ^2`$ in the adjacent background pixels. The variable sources should be visible as local excesses in the $`\chi ^2`$ map of the region of interest.
If a source can be localized in the variability image one can estimate the ratio between the average flux, $`\overline{F}`$ and the amplitude of the flux variations, $`\mathrm{\Delta }F`$. One can assume, for the sake of the argument, that the variations are normally distributed around the average value with the typical variance $`\mathrm{\Sigma }_s^2`$. Assuming that the variations of the background and of the source flux are independent one can ascribe the excess scatter of the flux data points at the location of the source to the source variability. If the typical variance of the background in each ScW is $`\mathrm{\Sigma }_b^2`$, the resulting variance in the image pixel containing the source is $`\mathrm{\Sigma }^2`$(Ra<sub>s</sub>, Dec<sub>s</sub>)$`=\chi _{red}^2`$(Ra<sub>s</sub>, Dec$`{}_{s}{}^{})\mathrm{\Sigma }_b^2/N=(\mathrm{\Sigma }_b^2+\mathrm{\Sigma }_s^2)/N`$. One can estimate the source variability as $`V=\mathrm{\Delta }F/\overline{F}=\sqrt{(\chi _{red}^2(Ra_s,Dec_s)1)\mathrm{\Sigma }_b^2}/\overline{F}`$. Note, that if a source is not detected in the variability image, this means that typical variations of its flux are $`\mathrm{\Delta }F\sqrt{\mathrm{\Delta }\chi ^2\mathrm{\Sigma }_b^2}`$ where $`\mathrm{\Delta }\chi ^2`$ is the typical scatter of the values of $`\chi ^2`$ for the background pixels. Taking into account that for the large number of ScWs (large $`N`$) $`\mathrm{\Delta }\chi ^2N^{1/2}`$, one can estimate the variable contribution to the flux as $`\mathrm{\Delta }F\sqrt{\mathrm{\Sigma }_b^2/N}`$.
Note that the above upper limit on the variable contribution to the flux is of the order of the square root of the variance of the mean value $`\sqrt{\mathrm{\Sigma }^2}`$. It is known that with the current version of OSA (4.2), the systematic effects start to contribute significantly to the variance when the exposure time is more than $`10^5`$ s (Bodaghee et al. (2004)). This leads to a slower than $`N^{1/2}`$ decrease of the variance with the increasing number of ScWs. The same should apply for the variability upper limit $`\mathrm{\Delta }F`$.
## 3 Variability of the sources in GC region
We have applied the above method to the $`1^{}\times 2^{}`$ sky region centered on the GC. Using the publicly available data (see http://isdc.unige.ch) with effective exposure time at the position of the GC of more than $`1`$ Msec we have produced the hard X-ray images of the region of interest with the Offline Software Analysis (OSA) package (Courvoisier et al. (2004)), version 4.2. Fig. 1a shows the 20-60 keV intensity map of the field around the GC. The size of each pixel in the images is $`3^{}\times 3^{}`$. The value of each pixel of the image is the weighted average $`\overline{F}`$ of the flux from the sky direction corresponding to the center of the pixel. Fig. 1b is the $`\chi _{red}^2`$ (variability) map of the same region in the same energy band.
One can see that from all sources visible on the intensity map, Fig. 1a, only 3 appear at the variability map, Fig. 1b. Among the variable sources are the well-known X-ray binaries, 1E 1740.7-2942, A 1742-294 and KS 1741-293. The variability indexes $`V=\mathrm{\Delta }F/F`$ for these sources are $`0.07`$, $`0.22`$ and $`0.59`$, respectively. One can see that in spite of large excess $`\chi _{red}^2`$ at the location of the brightest source, the X-ray binary 1E 1740.7-2942, the source flux varies at the level of only about 7%. This implies an important limitation for the possible applications of the method of variability search described above: the X-ray binaries which produce much lower flux (e.g. at the level of KS 1741-293) but have similar variability properties as 1E 1740.7-2942, would not be localized in the variability map, because the excess of $`\chi _{red}^2`$ produced by such sources would be less than the statistical scatter of the $`\chi _{red}^2`$ values over the image pixels.
In this respect the source 1E 1743.1-2843 is an interesting example. As one can see, this source is not detected in the variability image. Calculating the upper limit on the variable fraction of the source flux, $`\mathrm{\Delta }F`$, along the lines explained above, one gets an upper limit on possible variability at the level of $`8`$%. 1E 1743.1-2843 was discovered by the Einstein Observatory (Watson et al. (1981)). Its spectral characteristics suggest an interpretation as a neutron star LMXB (Cremonesi et al. (1999)) while its luminosity suggests that the source should belong to the Atoll-class sources and should exhibit Type-I X-ray bursts. The lack of substantial variations of the X-ray flux from this source during the last 20 yr of observations questions such an interpretation (Porquet et al. 2003a ). The lack of variability in 20-60 keV energy band also argues against the neutron-star X-ray binary interpretation for this source, although a definite conclusion would be possible only after a systematic study of variability properties of LMXBs detected by the INTEGRAL.
## 4 Non-variability of the GC source.
The GC itself (the source IGR J1745.6-2901) also does not appear in the variability map. The upper limit on the variable contribution to the flux can be obtained, as it is explained above, from the typical value of the variance in a single ScW of duration of $`2`$ ksec in 20-60 keV band (Bodaghee et al. (2004)), $`\mathrm{\Sigma }_b^210^{10}`$ erg/(cm<sup>2</sup> s). Dividing $`\mathrm{\Sigma }_b^2`$ by $`\sqrt{N}`$ and taking into account the systematic “renormalization” by a factor $`1.4`$ which has to be applied to the image Fig. 1b (see the discussion at the end of the previous section), we find an upper limit to the variable contribution at the level of $`V=\mathrm{\Delta }F/F0.08`$, or, equivalently, $`\mathrm{\Delta }F5\times 10^{12}`$ erg/(cm<sup>2</sup> sec) in the 20 to 60 keV band.
This upper limit is obtained under certain assumption about the variability type (Gaussian fluctuations around the mean flux). At the same time, if the variable emission is of different type, the above upper limit would not apply. For example, if the source has exhibited one or two powerful flares during the whole observation period, the amplitude of the flares can not be deduced from the statistical fluctuations of the signal.
The non-detection of variability of IGR J1745.6-2901 indicates that the main contribution to the hard X-ray flux from this source possibly comes from an extended region around the supermassive BH, rather than from the direct vicinity of the BH horizon. Indeed, simple powerlaw extrapolation of the observed X-ray flux from Sgr A\* shows that during the brightest flares the 10-100 keV flux would be roughly at the level detected by the INTEGRAL. However, in this case the hard X-ray flux should vary by $`100`$% at the ksec time scale, which is well beyond the upper limit found above. At the same time, the flux in 1-10 keV energy band is dominated by the diffuse emission and a simple powerlaw extrapolation of the diffuse X-ray signal to hard X-ray band is in good agreement with the INTEGRAL data (see below).
As it is discussed in the Introduction, the diffuse X-ray emission from the region around Sgr A\* contains a hard component of unclear nature. The 8 keV “temperature” of this hard component is at the upper end of the energy band probed by the X-ray telescopes and in order to prove the thermal nature of this component one has to observe the exponential cut-off in the spectrum above 10 keV. The analysis of INTEGRAL spectrum of IGR J1745.6-2901 shows that such a sharp cut-off is not observed.
The 20-200 keV spectrum of IGR J1745.6-2901 is shown on Fig. 2 together with XMM-Newton spectrum in 3-10 keV energy band extracted from a circular region of radius $`R=6^{}`$ around the GC (chosen to match the angular resolution of ISGRI imager on board of INTEGRAL).
The XMM-Newton Observation Data Files (ODFs), obs\_id 0111350101, were obtained from the online Science Archive http://xmm.vilspa.esa.es/external/xmm\_data\_acc/xsa/index.shtml. The data were processed and the event-lists filtered using xmmselect within the Science Analysis Software (sas) v6.0.1. The 20-200 keV spectrum can be fit by a simple power law model with the photon spectral index $`\mathrm{\Gamma }_{20200}=3.0\pm 0.1`$ and flux $`F_{20200}=(4\pm 2)\times 10^{11}`$ erg/(cm<sup>2</sup> s). However, a broken power-law provides a better fit to the INTEGRAL data. The broken power law with a break energy at $`E_{break}=26\pm 1`$ keV is characterized by a low energy spectral index $`\mathrm{\Gamma }_{low}=1.85_{0.06}^{+0.02}`$, and a high-energy spectral index $`\mathrm{\Gamma }_{high}=3.3\pm 0.1`$. It provides the best simultaneous fit to the XMM-Newton and INTEGRAL data. The intercalibration factor between the two instruments is just $`1.2`$ which indicates that 20-200 keV flux detected by INTEGRAL matches well the 1-10 keV flux collected from the region of the size about the point spread function of ISGRI.
It is important to note that because of the presence of diffuse emission component around the Galactic Center, the normalization of the XMM-Newton flux depends strongly on the size of the region from which the spectrum is extracted. For example, if one collects the flux from a disk-like region of the radius $`r`$ centered on the Galactic Center, the flux grows as $`r^2`$ for $`r<1^{}`$, see Fig. 3. From Fig. 3 one can see that the mismatch between the normalization of the INTEGRAL spectrum and XMM-Newton spectrum collected from the $`r1^{}`$ disk would be a factor of $`10`$. For the disk radii $`r>1^{}`$ the flux grows proportionally to $`r`$. The fact that the XMM-Newton spectrum matches the INTEGRAL spectrum when the disk radius reaches $`r6^{}`$ has an important implication for the physics of diffuse hard X-ray emission.
## 5 The nature of extended non-thermal hard X-ray emission around the GC.
We have seen in the previous sections that several complementary arguments indicate that the hard X-ray emission from the GC detected by INTEGRAL originates from an extended region around the supermassive black hole. The inter-calibration factor of order of $`1`$ between XMM-Newton and INTEGRAL spectra can be achieved only if the XMM-Newton flux is collected from a circular region of the radius $`r6^{}`$. This means that the size of the hard X-ray extended emission region is about $`D20`$ pc.
In order to understand the mechanism of diffuse nonthermal 10-100 keV emission from the inner $`10`$ pc of the Galaxy, it is useful to recall that this region is characterized by quite special astrophysical conditions. It is quite densely populated with giant molecular clouds and the typical gas/dust density throughout the region is $`n10^4`$ cm<sup>-3</sup> (see (Metzger et al. (1996)) for a review). Most of the estimates of magnetic field are in the range $`B10^410^3`$ G (Yusef-Zadeh et al. (1996)), much stronger than the typical Galactic magnetic field. Besides, the density of the infrared-optical background in this region, $`U_{iro}10^210^3`$ cm<sup>-3</sup> (Cox and Laureijs (1988)), is more than two orders of magnitude higher than the density of cosmic microwave background radiation. Such “exotic” conditions should be taken into account in the modeling of physical processes leading to the emission of diffuse non-thermal X-rays.
Among the possible emission mechanisms, synchrotron radiation from electrons of energies $`E_e10^{13}10^{14}`$ eV is the most plausible candidate. The typical energy of synchrotron photons produced by such electrons is $`E_{syn}50(B/10^4`$ G$`)(E_e/10^{14}`$ eV$`)^2`$ keV. The synchrotron cooling time of electrons emitting at 10 keV is $`t_{syn}16(B/10^4`$ G$`)^{3/2}(ϵ/10`$ keV$`)^{1/2}`$ yr. One can see that the cooling time is too short for electrons accelerated near the supermassive black hole to spread over the emission region of the size $`D20`$ pc if one assumes diffusion of electrons injected by a central source. Naively, to overcome this difficulty, one would assume a lower magnetic field strength in the emission region. However, in this case the inverse Compton (IC) flux from electrons which emit synchrotron radiation at $`10`$ keV would be stronger than the TeV flux from the GC detected by HESS. Thus, within the synchrotron scenario one has to assume that either magnetic field in the emission region is mostly ordered, so that electrons escape along the magnetic field lines, rather than diffuse in the random magnetic field, or that 10-100 TeV electrons are injected not by a point source at the location of the GC, but throughout the whole extended region of the size $`1020`$ pc. Possible mechanism leading for such “extended injection” is e.g. cascading of the 1-100 TeV gamma quanta on dense infrared photon background in the central 20 pc of the Galaxy (Neronov et al. (2002); Aharonian and Neronov (2005)). Otherwise, electrons can be accelerated in the shell of supernova remnant Sgr A East, whose size is about 5 pc.
Large density of the infrared-optical photon background and of the molecular gas in the emission region lead to significant IC and bremsstrahlung emission from the $`10^{13}10^{14}`$ eV electrons. The results of calculation of the broad-band spectrum synchrotron-IC-bremsstrahlung emission from the inner 20 pc of the Galaxy are shown in Fig. 4. One can see that the expected IC and bremsstrahlung fluxes are strong enough to match the observed level of the TeV emission from the GC. If the TeV flux detected by HESS is produced via the above mechanism, it should be not variable, since the IC and bremsstrahlung flux also come from the extended region of the size of several tens of kiloparsecs. The (non)variability of the TeV signal from the GC can be tested with future HESS observations.
Since the synchrotron emission results in efficient cooling of the high-energy electrons, all the power supplied by the supermassive black hole is dissipated with almost 100% efficiency in the form of the hard non-thermal X-ray emission. The required power of the supermassive black hole in the “synchrotron-IC-bremsstrahlung” scenario is just $`P_{syn}10^{36}`$ erg/s, about the total power of Sgr A\* observed in infrared. Thus, the above scenario is the most “economic” one.
Other possible mechanisms of non-thermal 10-100 keV emission appear to be much less efficient. For example, if one tries to explain the observed hard X-ray flux with the bremsstrahlung, one finds immediately that for moderately relativistic electrons which can emit bremsstrahlung radiation at $`2050`$ keV, the bremsstrahlung cooling rate is some 5 orders of magnitude less than the Coulomb loss rate, which means that the bremsstrahlung is very energetically inefficient mechanism of powering the nonthermal hard X-ray emission. The supermassive black hole should produce the power at the level of $`L_{BH}10^5L_{Xray}10^{41}`$ erg/s in this scenario. Although such a luminosity is still much below the Eddington luminosity of a $`4\times 10^6M_{}`$ black hole, it is orders of magnitude higher than the observed luminosity of Sgr A\* in the wide photon energy range from radio to very-high-energy gamma-rays.
The possibility that the observed non-thermal emission is produced via IC scattering of the dense infrared photon background is also unsatisfactory. The IC cooling time for electrons emitting in the 10-100 keV band, $`t_{IC}3\times 10^6(U_{iro}/(100\text{ eV/cm}^3))^1(E_e/1`$ GeV$`)^1`$ yr, is much larger than the bremsstrahlung cooling time, $`t_{br}4\times 10^3(n/10^4`$ cm$`{}_{}{}^{3})^1`$ yr. This means that (1) the supplied power should be 2-3 orders of magnitude higher than the observed hard X-ray luminosity (depending on the assumptions about the dust density) and (2) the hard X-ray IC emission should be accompanied by much stronger 100 MeV bremsstrahlung emission. In fact, the predicted bremsstrahlung flux is at least an order of magnitude higher than the GC flux found in the 100 MeV-GeV band by EGRET (see Fig.4).
Thus, the synchrotron radiation is the most probable mechanism of the diffuse hard X-ray emission from the inner 20 pc of the Galaxy. In order to test the above synchrotron-IC-bremsstrahlung scenario, one has to study the correlation of the spatial distributions of 20-100 keV and 1-10 TeV fluxes. Although it is quite difficult to analyze extended sources with the coded mask instruments, like INTEGRAL, it would be interesting to extract the information on the extension of the source directly from the 10-100 keV data (not from the matching with the lower-energy spectrum, like it is done above). We leave this question for future study.
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# 1 The quark spin fraction ΔΣ of the nucleon in dependence of the pion mass 𝑚_𝜋²["GeV"]². The filled circles with error bars represent the predictions for ΔΣ by the LHPC group, corresponding respectively to 𝑚_𝜋={744,831}, and 897"MeV", while the cross with error bar stands for the empirical value corresponding to the physical pion mass 𝑚_𝜋=138"MeV". The predictions of the CQSM are shown by the open squares for four values of 𝑚_𝜋, i.e., 𝑚_𝜋={0,200,400}, and 600"MeV". Also shown by the filled triangles are the predictions of the CQSM scaled by the factor 0.7.
OU-HET-531
Moments of generalized parton distribution functions
and the nucleon spin contents
M. Wakamatsu<sup>1</sup><sup>1</sup>1Email : [email protected]
Department of Physics, Faculty of Science,
Osaka University,
Toyonaka, Osaka 560-0004, JAPAN
Abstract
It is shown that, based only on two empirically known facts besides two reasonable theoretical postulates, we are inevitably led to a conclusion that the quark orbital angular momentum carries nearly half of the total nucleon spin. We also perform a model analysis to find that the quark spin fraction $`\mathrm{\Delta }\mathrm{\Sigma }`$ is extremely sensitive to the pion mass, which may resolve the discrepancy between the observation and the prediction of the recent lattice QCD simulation carried out in the heavy pion region.
The so-called “nucleon spin puzzle” raised more than 15 years ago is still an unsolved fundamental puzzle in hadron physics ,. If intrinsic quark spin carries little of the total nucleon spin, what carries the rest of the nucleon spin ? It is the question to be answered. Admitting that the QCD is a correct theory of strong interaction, the answer must naturally be sought for in the following three ; the quark orbital angular momentum (OAM), the gluon polarization, and the gluon orbital angular momentum.
Roughly speaking, there exist two contrasting or opposing standpoints to try to answer the above question. The chiral soliton picture of the nucleon emphasizes the importance of the quark orbital angular momentum ,. On the other hand, the possible importance of the gluon polarization was stressed by several authors in relation with the axial anomaly of QCD ,,. Later, the role of QCD anomaly was understood more clearly within the framework of the perturbative QCD, especially in view of the factorization-scheme dependence of parton distribution functions . Nonetheless, the serious problem is that no one can give any reliable theoretical prediction for the actual magnitude of $`\mathrm{\Delta }g`$.
An important remark here is that it is meaningless to talk about the nucleon spin contents without reference to the energy scale of observation. In fact, it is a widely known fact that the gluon polarization grows rapidly as $`Q^2`$ increases, even if it is small at low energy . In contrast, the gluon orbital angular momentum decreases rapidly to partially compensate the increase of $`\mathrm{\Delta }g`$. (Strictly speaking, these statements are gauge dependent, since it is known that there is no gauge invariant decomposition of gluon angular momentum into spin and orbital angular momentum .) Consequently, when we talk about the nucleon spin contents naively, we should implicitly understand that we are thinking of it at low energy scale of nonperturbative QCD. At this low energy, the CQSM predicts that ,,
$$\mathrm{\Delta }\mathrm{\Sigma }0.35,2L_q0.65,$$
(1)
which means that the quark OAM dominates over the contribution of quark intrinsic spin.
We repeat the question,“Which carries the rest of the nucleon spin, $`L_q`$ or $`\mathrm{\Delta }g`$?” Naturally, only experiments can answer it. A direct measurement of $`\mathrm{\Delta }g`$ via photon-gluon fusion processes is one of the most promising direction of study. For instance, the Compass group recently extracted the value of $`\mathrm{\Delta }g/g`$ through the analysis of the asymmetry of high $`p_T`$ hadron pairs. Their first result for $`\mathrm{\Delta }g/g`$ has turned out to be fairly small ,
$$\mathrm{\Delta }g/g0.06\pm 0.31(stat.)\pm 0.06(syst.),$$
(2)
although it is premature to draw any decisive conclusion only from this analysis. On the other hand, the key quantity for the direct measurement of $`J_q`$ or $`L_q`$ is the generalized parton distributions appearing in the cross sections of deeply virtual Compton scattering and/or deeply virtual meson productions . What plays the central role here is Ji’s quark angular momentum sum rule.
Here, we start our argument with the familiar definition of the generalized form factors $`A_{20}(t)`$ and $`B_{20}(t)`$ of the nucleon, which is given as a nonforward matrix element of QCD energy momentum tensor $`T_{q,g}^{\mu \nu }`$ :
$$P^{}|T_{q,g}^{\mu \nu }|P=\overline{U}(P^{})\left[A_{20}^{q,g}(t)\gamma ^{(\mu }P^{\nu )}+B_{20}^{q,g}(t)\frac{P^{(\mu }i\sigma ^{\nu )\alpha }\mathrm{\Delta }_\alpha }{2M}\right]U(P)+\mathrm{}.$$
(3)
According to Ji’s sum rule, the total angular momentum carried by quark and gluon fields in the nucleon is related to the forward ($`t=0`$) limit of these generalized form factors as
$`J^{u+d}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[A_{20}^{u+d}(0)+B_{20}^{u+d}(0)],`$ (4)
$`J^g`$ $`=`$ $`{\displaystyle \frac{1}{2}}[A_{20}^g(0)+B_{20}^g(0)].`$ (5)
Remembering the fact that the above generalized form factors $`A_{20}^{u+d}(0)`$ and $`A_{20}^g(0)`$ are related to the second moments of the unpolarized generalized parton distribution functions of quarks and gluons, which reduce to the familiar unpolarized distributions for quarks and gluons in the forward limit, they just represent the the total momentum fraction of quarks and gluons in the nucleon as
$`A_{20}^{u+d}(0)`$ $`=`$ $`{\displaystyle _0^1}x[u(x)+\overline{u}(x)+d(x)+\overline{d}(x)]𝑑xx^{u+d},`$ (6)
$`A_{20}^g(0)`$ $`=`$ $`{\displaystyle _0^1}xg(x)x^g.`$ (7)
On the other hand, the second $`B`$ parts are sometimes called the anomalous gravitomagnetic moments (AGM) of the constituents of the nucleon . From the conservation of total momentum and angular momentum, it follows that
$`A_{20}^{u+d}(0)+A_{20}^g(0)=1,`$ (8)
$`A_{20}^{u+d}(0)+B_{20}^{u+d}(0)+A_{20}^g(0)+B_{20}^g(0)=1,`$ (9)
which in turn dictates a nontrivial identity :
$$B_{20}^{u+d}(0)+B_{20}^g(0)=0.$$
(10)
To proceed further, we must distinguish three possibilities below :
1. $`B_{20}^{u+d}(0)=B_{20}^g(0)0`$,
2. $`B_{20}^{u+d}(0)=B_{20}^g(0)=0`$,
3. $`B_{20}^u(0)=B_{20}^d(0)=B_{20}^g(0)=0`$ .
The recent lattice QCD simulation by LHPC Collaboration gives a strong support to the second possibility that the total quark contribution to the nucleon AGM vanishes ,. This happens as a cancellation of the $`u`$\- and $`d`$-quark contributions, i.e., $`B_{20}^u(0)`$ and $`B_{20}^d(0)`$, which have sizable magnitudes with opposite signs. Noteworthy here is the fact that both of $`B_{20}^u(0)`$ and $`B_{20}^d(0)`$ have fairly strong dependence on the pion mass but their sum is almost independent on it. In any case, this lattice analysis seems to deny the third possibility indicated in on the basis of the equivalence principle, but strongly supports the second possibility, which is the basis of the following argument. In fact, once we accept this postulate, we are led to a surprisingly simple result that the total quark angular momentum is just a half of the total quark angular momentum fraction ,,, :
$$J^{u+d}=\frac{1}{2}x^{u+d}.$$
(11)
Now, we can go further. First, let us recall an empirically well-accepted understanding that, even at low energy scale like $`Q^2(600\text{MeV})^2`$, the gluon field seems to carry about $`(2030)\%`$ of the total nucleon momentum. (The widespread belief that the quark and gluon fields share equal amounts of nucleon momentum applies to the asymptotic case of large $`Q^2`$). For instance, one may consult the well-established GRV fit of the unpolarized parton densities . (See also .) Their next-to-leading order fit of the gluon density is given at $`Q^2=\mu _{NLO}^2=0.40\text{GeV}^2(630\text{MeV})^2`$ as
$$xg(x,\mu _{NLO}^2)=20.8x^{1.6}(1x)^{4.1}.$$
(12)
This turns out to give
$$x^g_0^1xg(x,\mu _{NLO}^2)𝑑x0.30.$$
(13)
Conversely saying, we can say that, at low energy, the quark field carries at least $`(7080)\%`$ of nucleon momentum, which in turn must be equal to the total quark angular momentum fraction, according to the aforementioned argument, such that
$$2J^{u+d}=x^{u+d}=(0.70.8).$$
(14)
On the other hand, through the analysis of polarized deep-inelastic scatterings, we already know that the intrinsic quark polarization $`\mathrm{\Delta }\mathrm{\Sigma }`$ is about $`(2035)\%`$ (see, for instance, the recent review ) :
$$\mathrm{\Delta }\mathrm{\Sigma }(0.20.35).$$
(15)
Putting these two observations (14) and (15) together, we find that the quark orbital angular momentum fraction is nearly $`50\%`$,
$$2L^{u+d}=2J^{u+d}\mathrm{\Delta }\mathrm{\Sigma }0.5.$$
(16)
That is, once admitting that the isosinglet combination of the quark contribution to the nucleon AGM vanishes, we are inevitably led to a surprising conclusion that the quark OAM carries nearly half of the nucleon spin, only with use of the empirically known information.
One might wonder why our conclusion is entirely different from that obtained by the LHPC Collaboration ,, who claims that the quark OAM is negligibly small, in spite that our argument above is based on a result of the LHPC group, i.e., $`B_{20}^{u+d}(0)=0`$. The rest of the present report is devoted to clarifying this point. The reason can easily be traced back to the fact that, instead of using the empirical value of $`\mathrm{\Delta }\mathrm{\Sigma }`$, they used their theoretical predictions for it,
$$\mathrm{\Delta }\mathrm{\Sigma }(\text{LHPC})0.682,$$
(17)
which is fairly large and clearly contradicts the observation. Why does their analysis give very large $`\mathrm{\Delta }\mathrm{\Sigma }`$, then? This is probably because their simulation was performed with quite large pion mass around $`m_\pi (700900)\text{MeV}`$, which is far from our realistic world close to the chiral limit. As we shall discuss below, the strong sensitivity of $`\mathrm{\Delta }\mathrm{\Sigma }`$ on the pion mass seems to be a likely solution to the above-mentioned discrepancy.
Now, we shall show it on the basis of the CQSM ,. Within the framework of the CQSM, we first solve the mean-field equation of motion self-consistently for several values of $`m_\pi `$. The model is defined with a physical cutoff. Here we use the Pauli-Villars regularization scheme with double substraction terms . The relevant regularization parameters are all fixed uniquely from reasonable physical requirements. How to introduce finite pion mass into the whole scheme is explained in . Here, we tried to find a self-consistent soliton profile with the fixed value of the dynamical quark mass $`M=400\text{MeV}`$. This is repeated for several values of pion mass, i.e., $`m_\pi =0,200,400`$, and $`600\text{MeV}`$. In this analysis, no stable solution was found for $`m_\pi >620\text{MeV}`$. We then evaluate $`\mathrm{\Delta }\mathrm{\Sigma }`$ for each soliton solution with different value of $`m_\pi `$. The results are shown in Fig.1 together with the predictions of the LHPC Collaboration ,. The filled circles with error bars represent the predictions for $`\mathrm{\Delta }\mathrm{\Sigma }`$ by the LHPC group, corresponding respectively to $`m_\pi =744,831`$, and $`897\text{MeV}`$, while the cross with error bar stands for the empirical value corresponding to the physical pion mass $`m_\pi =138\text{MeV}`$. The predictions of the CQSM are shown by the open squares for four values of $`m_\pi `$, i.e., $`m_\pi =0,200,400`$, and $`600\text{MeV}`$. One clearly sees that $`\mathrm{\Delta }\mathrm{\Sigma }`$ is very sensitive to the value of $`m_\pi `$, especially when approaching the chiral limit $`m_\pi 0`$. Inspired by the indication of the GRSV fit, which dictates that the quark fields carries only $`70\%`$ of the total nucleon momentum and also the total spin, one may tentatively renormalize the predictions of the CQSM by multiplying a factor of 0.7. The results are shown by the filled triangles. It is interesting to see that these points appears to be smoothly connected to the lattice predictions given in the large $`m_\pi `$ domain. We hope that the lattice simulation in the near future will be extended to the region of smaller $`m_\pi `$ and that it will confirm the strong $`m_\pi `$ dependence of $`\mathrm{\Delta }\mathrm{\Sigma }`$ predicted by the CQSM, although the lattice QCD would still need a help of some other theoretical technique like the chiral perturbation theory to explore the region of $`m_\pi `$ very close to the chiral limit .
Summarizing our arguments, we have shown that, with use of the two empirical knowledge alone, aside from the two reasonable theoretical postulates, we are inevitably led to a drastic conclusion that the quark orbital angular momentum carries nearly half of the total nucleon spin. The two theoretical postulates here are
* Ji’s angular momentum sum rule : $`J^{u+d}=\frac{1}{2}[x^{u+d}+B_{20}^{u+d}(0)]`$ ,
* absence of the net quark contribution to the anomalous gravitomagnetic moment of the nucleon : $`B_{20}^{u+d}(0)=0`$ .
On the other hand, the two empirically known facts we have used are
* the fraction of the quark momentum and angular momentum of the nucleon at low energy scale, $`Q^2(600\text{MeV})^2`$ : $`x^{u+d}=2J^{u+d}(0.70.8)`$ ,
* the quark spin fraction from polarized DIS analyses : $`\mathrm{\Delta }\mathrm{\Sigma }(0.20.35)`$ .
Although there remains some room concerning how to define the angular momentum of the constituents of the nucleon , it is reasonable to stick to Ji’s definition, which leads to the above sum rule. Otherwise, we would lose a only clue to experimentally access the quark angular momentum in the nucleon. Thus, only one factor, which might potentially alter our conclusion, is the second postulate, i.e., $`B_{20}^{u+d}(0)=0`$. Although it is strongly supported by the lattice simulation by the LHPC Collaboration, an independent check is highly desirable. Also desirable is an analytical proof of it within the framework of (nonperturbative) QCD.
We have also shown that the above-mentioned conclusion, obtained independently of any models, is qualitatively consistent with the predictions of the CQSM. The CQSM predicts very strong dependence of the quark spin fraction $`\mathrm{\Delta }\mathrm{\Sigma }`$ on the pion mass : it reproduces small $`\mathrm{\Delta }\mathrm{\Sigma }`$ in the domain close to the chiral limit, it also smoothly matches the predictions of the LHPC Collaborations obtained in the heavy pion region. It is hoped that this behavior of $`\mathrm{\Delta }\mathrm{\Sigma }`$ will be confirmed by the lattice simulation in the near future.
Acknowledgement
This work is supported in part by a Grant-in-Aid for Scientific Research for Ministry of Education, Culture, Sports, Science and Technology, Japan (No. C-16540253)
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# Unstable quasi g-modes in rotating main-sequence stars
## 1 Introduction
With the introduction of new opacity data by Iglesias et al. (iglesias92 (1992)) and a later update by Iglesias & Rogers (iglesias96 (1996)) it was suggested that the origin of $`\beta `$ Cephei pulsations in massive main-sequence stars and of the slowly pulsating B (SPB) stars (Waelkens waelkens91 (1991)) of lower mass is caused by the classical opacity valve mechanism (Eddington eddington (1930)) operating on non radial g-mode or p-mode oscillations. Early stability analyses with these new opacity results in non-rotating stars have been published by e.g. Kiriakidis et al. (kiria92 (1992)), Gautschy and Saio (gautschy93 (1993)) and Dziembowski et al. (dziembowski93 (1993)).
Observations by Balona (balona94 (1994)) and Balona and Koen (balonaK94 (1994)) showed an unexplainable lack of SPB stars in two open clusters which do contain many pulsating $`\beta `$ Cephei stars. It was suggested that the lack of SPB stars in these clusters might be correlated with the rapid rotation of the cluster stars. Ushomirsky and Bildsten (usho98 (1998)), in an attempt to explain Balona and Koen’s observations in terms of rotational stabilization of g-modes, applied the traditional approximation combined with a quasi-adiabatic stability analysis to study the effect of rotation on the $`\kappa `$ instability mechanism. They did not find a decisive solution: rotation seemed indeed to stabilize some pulsation modes but to destabilize other g-modes. This work was followed up by Lee (lee01 (2001)) who used a fully non-adiabatic analysis with a truncated expansion of spherical harmonics to approximate the low-frequency g-mode oscillations in rotating stars. Lee found some rotational stabilization caused by the Coriolis coupling with higher degree spherical harmonic components, especially for retrograde g-modes in the inertial regime, but this seems not sufficient to explain the lack of SPB stars in clusters with rapidly rotating B-stars. Townsend (town03 (2003)), also applying the traditional approximation, tried to explain the lack of observed rapidly rotating SPB stars by the well known effect of the Coriolis force on g-mode pulsations: the confinement of the oscillation to the equatorial regions ($`|\mu |<\overline{\sigma }/(2\mathrm{\Omega }_s)`$) of the star, where $`\mu =\mathrm{cos}\vartheta `$ and $`\overline{\sigma }`$ is the oscillation frequency in the corotating frame (e.g. Savonije et al. sav95 (1995)). The lack of rapidly rotating SPB stars would then be caused by a selection effect, not by rotational stabilization of g-modes. However, more recent observations indicate that the two open clusters studied by Balona and Koen may after all contain a few SPB stars (Stankov et al. stank02 (2002)), so that the problem may be non-existent.
There is, however, another complication: the existence of unstable rotational ‘q-modes’ in rotating main-sequence stars whose instability has been hitherto neglected. These q-modes exhibit hardly any rotational confinement towards the equatorial region, in contrast to normal g-modes. In this paper we will study the stability of these quasi g-modes in comparison with normal (retrograde) g-modes and check how well unstable q-modes can explain the oscillation spectra of rotating B-stars.
## 2 Basic oscillation equations
We consider uniformly rotating main sequence stars with mass $`M_\mathrm{s}`$ between 3-8 $`\mathrm{M}_{}`$ and denote their radius by $`R_\mathrm{s}`$. We wish to study the oscillatory stability of these uniformly rotating B-stars by subjecting them to a perturbing periodic force to determine the resonant behaviour. We assume the star’s angular velocity of rotation $`\mathrm{\Omega }_\mathrm{s}`$ to be much smaller than its break-up speed, i.e. $`(\mathrm{\Omega }_\mathrm{s}/\mathrm{\Omega }_\mathrm{c})^21`$, with $`\mathrm{\Omega }_\mathrm{c}^2=GM_\mathrm{s}/R_\mathrm{s}^3`$, so that effects of centrifugal distortion ($`\mathrm{\Omega }_\mathrm{s}^2`$) may be neglected to a first approximation. The Coriolis acceleration is proportional to $`\mathrm{\Omega }_\mathrm{s}`$ and we consider its effect on the induced oscillatory motions in the star. We use nonrotating spherical coordinates ($`r,\vartheta ,\phi )`$, with the origin at the star’s centre, whereby $`\vartheta =0`$ corresponds to its rotation axis.
Let us denote the displacement vector in the star by $`\xi `$ and perturbed Eulerian quantities like pressure $`P^{}`$, density $`\rho ^{}`$, temperature $`T^{}`$ and energy flux $`F^{}`$ with a prime. The linearized hydrodynamic equations governing the non-adiabatic response of the uniformly rotating star to the perturbing potential $`\mathrm{\Phi }_\mathrm{T}`$ may then be written as
$`\left[\left({\displaystyle \frac{}{t}}+\mathrm{\Omega }_\mathrm{s}{\displaystyle \frac{}{\phi }}\right)v_i^{}\right]e_i+2\mathrm{\Omega }_\mathrm{s}\times v^{}=`$
$`{\displaystyle \frac{1}{\rho }}P^{}+{\displaystyle \frac{\rho ^{}}{\rho ^2}}P\mathrm{\Phi }_\mathrm{T},`$ (1)
$$\left(\frac{}{t}+\mathrm{\Omega }_\mathrm{s}\frac{}{\phi }\right)\rho ^{}+\left(\rho v^{}\right)=0,$$
(2)
$$\left(\frac{}{t}+\mathrm{\Omega }_\mathrm{s}\frac{}{\phi }\right)\left[S^{}+\xi S\right]=\frac{1}{\rho T}F^{},$$
(3)
$$\frac{F^{}}{F}=\left(\frac{\mathrm{d}T}{\mathrm{d}r}\right)^1\left[\left(\frac{3T^{}}{T}\frac{\rho ^{}}{\rho }\frac{\kappa ^{}}{\kappa }\right)T+T^{}\right]$$
(4)
where $`\kappa `$ denotes the opacity of stellar material and $`S`$ its specific entropy. These perturbation equations represent, respectively, conservation of momentum, conservation of mass and conservation of energy, while the last equation describes the radiative diffusion of the perturbed energy flux. For simplicity we adopt the Cowling (1941) approximation, i.e. we neglect perturbations to the gravitational potential caused by the star’s oscillatory distortion. For the higher radial order oscillation modes studied here, this approximation should be adequate. As usual, we also neglect perturbations of the nuclear energy sources (not important) and of convection (frozen convective flux approximation). In main-sequence (MS)stars more massive than about 6 $`M_{}`$ in which the opacity Z-bump driving region is convectively unstable this latter approximation introduces some uncertainty with respect to the detected oscillatory instabilities.
For the periodic forcing we apply a force equal to the (real part of) the gradient of
$$\mathrm{\Phi }_\mathrm{T}(r,\vartheta ,\phi ,t)=fr^lP_l^m(\mu )e^{\mathrm{i}(\sigma tm\phi )}$$
(5)
where $`\sigma `$ is the forcing frequency in the inertial frame, $`\mu =\mathrm{cos}\vartheta `$, $`P_l^m(\mu )`$ is an associated Legendre polynomial and $`f`$ is an arbitrary constant. Since we are interested in periodic solutions of the forced oscillations we assume the perturbations have the same time and $`\phi `$ dependence as the forcing (although with a certain phase lag due to radiative damping). After factoring out the known time and $`\phi `$-dependence, Eqs. (14) form a 2D problem in the ($`r,\vartheta `$) meridional plane of the perturbed star. The 2D problem requires a large amount of computer memory and CPU time to solve, see e.g. Savonije et al. (sav95 (1995)). However, it is possible to obtain solutions adequate for initial analyses of the stability of rotating stars by applying the so called ‘traditional approximation’.
## 3 The traditional approximation
For a rotating star the solutions of Eqs. (14) are no longer separable into $`r`$-,$`\vartheta `$\- and $`\phi `$\- factors and cannot be described by a single spherical harmonic $`(l,m)`$ due to the action of the Coriolis force on oscillating stellar matter. However, for a uniformly rotating star in the traditional approximation separability is retained by neglecting the $`\vartheta `$\- component of the angular velocity vector (e.g. Unno et al. unno89 (1989)). This is done because the radial motions are expected to be small in the stably stratified layers of the star, especially for low-frequency modes (Savonije et al. sav95 (1995)). In this way the problem is reduced to two 1-dimensional problems (in $`\vartheta `$ and $`r`$) with a prescribed azimuthal harmonic dependence $`e^{im\phi }`$ since we assume the unperturbed star to be axisymmetric. When only the radial component of $`\mathrm{\Omega }_\mathrm{s}`$ is retained, the $`\vartheta `$\- and $`\phi `$-components of the equation of motion can be written:
$$\overline{\sigma }^2\xi _\vartheta 2\mathrm{i}\mathrm{\Omega }_\mathrm{s}\overline{\sigma }\mathrm{cos}\vartheta \xi _\phi =\frac{1}{r\rho }\frac{P^{}}{\vartheta }\frac{1}{r}\frac{\mathrm{\Phi }_\mathrm{T}}{\vartheta },$$
(6)
$$\overline{\sigma }^2\xi _\phi +2\mathrm{i}\mathrm{\Omega }_\mathrm{s}\overline{\sigma }\mathrm{cos}\vartheta \xi _\vartheta =\frac{\mathrm{i}m}{r\rho \mathrm{sin}\vartheta }P^{}+\frac{\mathrm{i}m}{r\mathrm{sin}\vartheta }\mathrm{\Phi }_\mathrm{T}.$$
(7)
We have expressed the velocity perturbations in terms of the displacement vector by the relation $`v^{}=\mathrm{i}\overline{\sigma }\xi `$, with $`\overline{\sigma }=\sigma m\mathrm{\Omega }_\mathrm{s}`$ the oscillation frequency in the corotating stellar frame (negative if the oscillation is retrograde, i.e. propagates in the direction counter to the stellar rotation). We will only consider positive $`m`$-values in this paper, except when explicitely stated otherwise (in Tables 2-3).
### 3.1 Separation of variables
A separation of variables can be performed by writing, see Papaloizou & Savonije (pap97 (1997)):
$$\xi _\vartheta =\underset{n=0}{\overset{\mathrm{}}{}}_n(\vartheta )D_n(r)\text{ }\xi _\phi =\underset{n=0}{\overset{\mathrm{}}{}}𝒢_n(\vartheta )D_n(r)$$
$$P^{}=\underset{n=0}{\overset{\mathrm{}}{}}𝒳_n(\vartheta )W_n(r),$$
(8)
with $`\xi _r`$, $`T^{}`$ and $`\rho ^{}`$ having expansions of the same form as that for $`P^{}`$. Here we leave out the known factor $`e^{im\phi }`$ for the azimuthal variation of the perturbed quantities, while $`_n,𝒢_n`$ and $`𝒳_n`$ are functions of $`\vartheta `$ chosen to obey the relations
$$\overline{\sigma }^2_n2\mathrm{i}\mathrm{\Omega }_\mathrm{s}\overline{\sigma }\mathrm{cos}\vartheta 𝒢_n=\frac{\mathrm{d}𝒳_n}{\mathrm{d}\vartheta }$$
and
$$\overline{\sigma }^2𝒢_n+2\mathrm{i}\mathrm{\Omega }_\mathrm{s}\overline{\sigma }\mathrm{cos}\vartheta _n=\frac{\mathrm{i}m}{\mathrm{sin}\vartheta }𝒳_n.$$
We obtain an equation for $`𝒳_n(\vartheta )`$ by imposing the constraint
$$\frac{1}{\mathrm{sin}\vartheta }\frac{\mathrm{d}\left(\mathrm{sin}\vartheta _n\right)}{\mathrm{d}\vartheta }\mathrm{i}m\frac{𝒢_n}{\mathrm{sin}\vartheta }=\frac{\lambda _n}{\overline{\sigma }^2}𝒳_n,$$
where $`\lambda _n`$ is a constant. Then $`𝒳_n`$ must satisfy the second-order equation obtained from
$$\frac{1}{\mathrm{sin}\vartheta }\frac{\mathrm{d}𝒬_n}{\mathrm{d}\vartheta }+\frac{mx\mathrm{cos}\vartheta }{\mathrm{sin}^2\vartheta }𝒬_n=𝒳_n\left(\frac{m^2}{\mathrm{sin}^2\vartheta }\lambda _n\right)$$
(9)
with
$$𝒬_n=\frac{\mathrm{sin}\vartheta }{(1x^2\mathrm{cos}^2\vartheta )}\left(\frac{\mathrm{d}𝒳_n}{\mathrm{d}\vartheta }\frac{mx\mathrm{cos}\vartheta }{\mathrm{sin}\vartheta }𝒳_n\right)$$
(10)
whereby $`𝒳_n(\vartheta )`$ is an eigenfuction with $`\lambda _n`$ the associated eigenvalue. The solution depends on the rotation parameter $`x=2\mathrm{\Omega }_\mathrm{s}/\overline{\sigma }`$. For $`\mathrm{\Omega }_\mathrm{s}=0`$ the functions $`𝒳_n(\vartheta )`$ become the associated Legendre functions $`P_{m+n}^m(\mathrm{cos}\vartheta )`$ with corresponding eigenvalues $`\lambda _n=(m+n)(m+n+1).`$ Normal modes of the rotating star correspond to normal )modes of the non-rotating star with $`(m+n)(m+n+1)`$ replaced by any permissible $`\lambda _n.`$
Different $`𝒳_n(\vartheta )`$ are orthogonal on integration with respect to $`\mu =\mathrm{cos}\vartheta `$ over the interval $`(1,1)`$. If the perturbing potential is expanded in terms of the $`𝒳_n`$ such that
$$\mathrm{\Phi }_\mathrm{T}(r,\vartheta )=\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{\Psi }_n(r)𝒳_n(\vartheta ),$$
(11)
we find from (6) and (7) that
$$D_n(r)=\frac{1}{r\rho }W_n(r)+\frac{1}{r}\mathrm{\Psi }_n(r).$$
(12)
This equation is exactly the same as in the non-rotating case. The same is true for the equation of continuity, except that $`(m+n)(m+n+1)`$ is replaced by $`\lambda _n`$. In this way the adiabatic stellar response will consist of a superposition of responses appropriate to non-rotating stars with $`(m+n)(m+n+1)`$ replaced by $`\lambda _n`$ obtained from equations (9-10). However, we are interested in the oscillatory stability of rotating stars and for that we need to include non-adiabatic effects, i.e. we have to consider the energy equation (3). We can write the divergence term on its right hand side as:
$$\frac{F^{}}{F_r}=\frac{\left(r^2\frac{F_r^{}}{F_r}\right)}{r^2r}+\frac{1}{r\mathrm{sin}\vartheta }\frac{\left(\mathrm{sin}\vartheta \frac{F_\vartheta ^{}}{F_r}\right)}{\vartheta }\frac{\mathrm{i}m}{r\mathrm{sin}\vartheta }\left(\frac{F_\phi ^{}}{F_r}\right)$$
By adopting $`T^{}(r,\vartheta )=_nT_n^{}(r)𝒳_n(\vartheta )`$ and substitution in the flux equation (4) the radial part of the above divergence can be written in the required separated form $`𝒟_r(r,\vartheta )=H(r)𝒳_n(\vartheta )`$. However, the two angular terms cannot be expressed in this way. For one can write, after applying equation (4) and equations (9-10), the angular part of the above divergence as:
$$𝒟_a(r,\vartheta )=\frac{1}{r^2}\left(\frac{\mathrm{d}\mathrm{log}T}{\mathrm{d}r}\right)^1\frac{T^{}(r)}{T(r)}𝒴_n(\vartheta )𝒳_n(\vartheta )$$
with
$$𝒴_n(\vartheta )=2\mu x^2\frac{𝒬_n(\mu )}{𝒳_n(\mu )}+\left(x^2\mu ^2mx\lambda _n\right)$$
which shows that the energy equation renders the system of equations inseparable, unless $`x0`$. Nevertheless, we wish to apply the traditional approximation and take advantage of its simple treatment of rotational effects. Therefore we will approximate the energy equation by a separable equation by averaging the angular part of the divergence over $`\vartheta `$ and using the averaged value $`ϵ_n`$ for $`𝒴_n`$:
$$ϵ_n=\frac{𝒴_n𝒳_nd\mu }{𝒳_nd\mu }$$
so that
$$𝒟_a(r,\vartheta )\frac{1}{r^2}\left(\frac{\mathrm{d}\mathrm{log}T}{\mathrm{d}r}\right)^1\frac{T^{}(r)}{T(r)}ϵ_n𝒳_n(\vartheta )$$
(13)
For small values of the rotation-parameter $`x0`$ the exact result $`ϵ_n\lambda _n=(m+n)(m+n+1)`$ for non-rotating stars is approached.
Expression (13) takes the radiative diffusion in the horizontal direction at least qualitatively into account. Below we will check our results by comparing with other linear stability calculations of rotating stars.
### 3.2 The radial part of the oscillation equations
Once we have solved the above angular eigenvalue problem for $`\lambda _n`$ we can factor out the common, now known, factor $`𝒳_n(\vartheta )e^{\mathrm{i}(\sigma tm\phi )}`$ from all perturbed quantities. The linearized equations describing the non-adiabatic forced oscillations, for the $`n`$-th component in expansion (11), then form a 1-dimensional (radial) problem. The remaining radial part of the oscillation equations can be expressed as:
$$\overline{\sigma }^2\rho \xi _r\frac{\mathrm{d}P^{}}{\mathrm{d}r}+\frac{\mathrm{d}P}{\mathrm{d}r}\left(\frac{\rho ^{}}{\rho }\right)=l\rho fr^{\left(l1\right)},$$
(14)
$$\frac{1}{\rho r^2}\frac{\mathrm{d}\left(\rho r^2\xi _r\right)}{\mathrm{d}r}=\frac{\rho ^{}}{\rho }+\frac{\lambda _n}{\overline{\sigma }^2r^2}\left[\frac{P}{\rho }\left(\frac{P^{}}{P}\right)+fr^l\right],$$
(15)
$$\frac{P^{}}{P}\mathrm{\Gamma }_1\frac{\rho ^{}}{\rho }+𝒜\xi _r=\mathrm{i}\eta \left[\frac{1}{Fr^2}\frac{\mathrm{d}(r^2F_r^{})}{\mathrm{d}r}+ϵ_n\frac{\mathrm{\Lambda }}{r^2}\left(\frac{T^{}}{T}\right)\right],$$
(16)
$$\frac{F_r^{}}{F}=\mathrm{\Lambda }\frac{\mathrm{d}}{\mathrm{d}r}\left(\frac{T^{}}{T}\right)+(4\kappa _T)\frac{T^{}}{T}(1+\kappa _\rho )\frac{\rho ^{}}{\rho }$$
(17)
where $`\mathrm{\Gamma }`$’s represent Chandrasekhar’s adiabatic coefficients,
$$\mathrm{\Lambda }=\left(\frac{\mathrm{d}\mathrm{log}T}{\mathrm{d}r}\right)^1;\text{ }𝒜=\frac{\mathrm{d}\mathrm{log}P}{\mathrm{d}r}\mathrm{\Gamma }_1\frac{\mathrm{d}\mathrm{log}\rho }{\mathrm{d}r},$$
the opacity derivatives are given by:
$$\kappa _T=\left(\frac{\mathrm{log}\kappa }{\mathrm{log}T}\right)_\rho ;\text{ }\kappa _\rho =\left(\frac{\mathrm{log}\kappa }{\mathrm{log}\rho }\right)_T;$$
$`\eta `$ is a characteristic radiative diffusion length:
$$\eta =(\mathrm{\Gamma }_31)F/\left(\overline{\sigma }P\right),\text{ with }F=\frac{4acT^3}{3\kappa \rho }\frac{\mathrm{d}T}{\mathrm{d}r}$$
being the unperturbed (radial) radiative energy flux, $`ϵ_n`$ is related to the horizontal radiative diffusion and is defined by (13). All other constants have their usual meaning. Note that the radiative diffusion introduces a factor $`\mathrm{i}`$, so that the radial parts of the perturbations are complex-valued. This expresses the induced phase-lags with respect to the external forcing caused by non-adiabatic effects due to radiative diffusion.
The oscillation equations are complemented by the linearized equation of state
$$\frac{P^{}}{P}=\left(\frac{\mathrm{log}P}{\mathrm{log}T}\right)\frac{T^{}}{T}+\left(\frac{\mathrm{log}P}{\mathrm{log}\rho }\right)\frac{\rho ^{}}{\rho }+\left(\frac{\mathrm{log}P}{\mathrm{log}\mu _a}\right)\frac{\mu _a^{}}{\mu _a}$$
(18)
with
$$\frac{\mu _a^{}}{\mu _a}=\frac{\mathrm{d}\mathrm{log}\mu _a}{\mathrm{d}r}\xi _r$$
where $`\mu _a`$ is the mean atomic weight of the stellar gas.
Finally, we prescribe the usual boundary conditions $`\xi _r=F^{}=0`$ at the stellar centre and require that the Lagrangian perturbations at the stellar surface obey: $`\delta P=0`$ and $`\delta F/F=4\delta T/T`$ (Stefan’s law).
### 3.3 Stellar input models for the oscillation code
We constructed the unperturbed stellar models for the main-sequence stars with a recent version (Pols et al. pols95 (1995)) of the stellar evolution code developed by Eggleton (egg72 (1972)). The models represent spherical main-sequence stars with masses $`48M_{}`$ and chemical composition given by various values of the central hydrogen abundance $`X_\mathrm{c}`$ and $`Z=0.02`$ or $`Z=0.03`$. These models were constructed by using the OPAL opacities (Iglesias & Rogers iglesias96 (1996)). Table 1 lists the effective temperatures of the various stellar input models characterized by their mass, core hydrogen mass fraction $`X_c`$ and metal content $`Z`$, used in our calculations.
## 4 Solution method
### 4.1 Determination of the eigenvalue $`\lambda _n(x)`$
For a given stellar rotation frequency $`\mathrm{\Omega }_\mathrm{s}`$, forcing frequency in corotating frame $`\overline{\sigma }`$, azimuthal index $`m`$ and angular order $`n`$, the eigenvalues $`\lambda _n`$ can be determined by numerically integrating the two first order differential Eqs. (9) and (10). We have used a shooting method with fourth order Runge-Kutta integration with variable stepsize (Press et al. press92 (1992)) to obtain numerical solutions. It is convenient to rewrite equations (9-10) as
$$\frac{\mathrm{d}𝒳_n}{\mathrm{d}\mu }=\frac{mx\mu }{1\mu ^2}𝒳_n\left(\frac{1x^2\mu ^2}{1\mu ^2}\right)𝒬_n,$$
(19)
$$\frac{\mathrm{d}𝒬_n}{\mathrm{d}\mu }=\left(\frac{m^2}{1\mu ^2}\lambda _n\right)𝒳_n+\frac{mx\mu }{1\mu ^2}𝒬_n$$
(20)
where the factor $`x=2\mathrm{\Omega }_\mathrm{s}/\overline{\sigma }`$ expresses the importance of stellar rotation.
To enable integration away from $`\mu =1`$ it is convenient to write $`𝒳_n(\mu )=(1\mu ^2)^{\frac{m}{2}}Y_n(\mu )`$ and expand $`Y_n(\mu )`$ in a power series which can be substituted in Eqs. (19) and (20) to determine the coefficients. $`𝒬_n`$ can then be expressed in terms of $`Y_n`$ as
$$𝒬_n=\frac{(1\mu ^2)^{\frac{m}{2}}}{1x^2\mu ^2}\left[m\mu (1x)Y_n(1\mu ^2)\frac{\mathrm{d}Y_n}{\mathrm{d}\mu }\right].$$
For the even $`n`$ (including $`n`$=0) solutions the boundary conditions are $`𝒳_n=0`$ at $`\mu =1`$ and $`𝒬_n=0`$ at $`\mu =0`$, while the odd solutions are defined by $`𝒳_n=0`$ at $`\mu =0`$. We can now integrate equations (19)-(20) from $`\mu =1\delta `$ (with $`\delta =10^4`$) to $`\mu =0`$ with an estimated value for the eigenvalue $`\lambda _n`$. We iterate by adjusting $`\lambda _n`$ until the integrated value for either $`𝒬_n`$ (even solutions), or $`𝒳_n`$ (odd solutions) is sufficiently close to zero for $`\mu =0`$.
### 4.2 Eigenvalues $`\lambda _{nm}`$ and mode classification ($`n`$,$`m`$)
Eigenvalues $`\lambda _{nm}`$ of equations (9-10) (or its equivalent) have recently been calculated by Bildsten et al. (bild96 (1996)), Papaloizou & Savonije (pap97 (1997)), Lee & Saio (lee97 (1997)) and Savonije & Witte (sav02 (2002)). It appears that in the traditional approximation the angular eigenvalue $`\lambda _{nm}`$ can take both positive and negative values, including 0. Here we will follow Lee & Saio (1997, see their Fig. 1) and allow $`n`$ to assume positive and negative values (or 0) to discriminate between these cases. The positive values $`n=0,1,2,3`$,… correspond to gravity (g)-modes (or p-modes) for which $`\lambda >0`$ for any $`x`$-value, whereby even values of $`n`$ (or 0) correspond to modes with even symmetry about the stellar equator and odd values to modes with odd symmetry. For these modes $`n`$ gives the number of nodes in $`𝒳_n(\mu )`$ over the interval $`1<\mu <1`$, except when $`1<x^1<0`$ (where $`x^1=\overline{\sigma }/2\mathrm{\Omega }_s`$) when there appears an extra set of nodes. This is related to the existence of quasi g-modes with $`n=1`$, see below.
For $`n<0`$ and $`x<0`$ the eigenvalues are positive for $`|x^1|0`$, and negative for large values of $`|x^1|`$. In this case the eigenvalue $`\lambda _{nm}`$ changes sign for $`x`$-value:
$$x_r^1=\frac{m}{(m+|n|)(m+|n+1|)}$$
(21)
where $`m>0`$. The oscillation modes corresponding to these roots $`x_r`$ are the purely toroidal r-modes (Papaloizou & Pringle pap78 (1978)). The $`n<0`$ branches thus have positive, although small, eigenvalues $`\lambda _{nm}`$ for $`|x^1|<|x_r^1|`$ and these correspond to quasi-toroidal q-modes (Savonije & Witte, 2002). Fig. 1 shows the ($`\lambda >0`$) q-mode branches with $`n=1`$ and $`n=2`$ for $`m`$=2 together with the normal g-mode $`n`$=0 branch. The special case $`n=1`$ is interesting as this corresponds to a branch of q-modes where $`\lambda `$ does not remain small when $`|x^1|`$ decreases, see Fig. 1. For $`|x^1|<|x_r^1|`$ these $`n=1`$ modes have values of $`\lambda _{nm}`$ that remain an order of magnitude smaller than those of the lowest order ($`n=0`$) g-modes. This means that for these special ’quasi g-modes’ the confinement to the equatorial regions of the rotating star at lower frequencies $`\overline{\sigma }`$ is almost absent. We will see that for sufficient rotation speeds $`n=1`$ quasi g-mode oscillations can be destabilized by the opacity bump due to the heavy element ionisation zone, like normal g-modes.
The $`n<0`$ modes with $`\lambda _{nm}<0`$ (for both positive and negative $`x`$-values) correspond to rotationally stabilized g<sup>-</sup>-modes or oscillatory convection (and exists only) in the regions where the square of the Brunt-Väisälä frequency $`𝒩^2=g𝒜/\mathrm{\Gamma }_10`$. In this paper only modes with $`\lambda _{nm}>0`$ are explored.
### 4.3 Solution of radial oscillation equations
For a given stellar rotation speed $`\mathrm{\Omega }_s`$ and stellar forcing frequency $`\overline{\sigma }`$ (in fact for a given $`x^1=\overline{\sigma }/2\mathrm{\Omega }_s`$) we can for each possible combination of $`(n,m)`$ determine the eigenvalue $`\lambda _n`$ and corresponding eigenfunction $`𝒳_n(\vartheta )`$ (and auxilliary function $`𝒬_n`$ from which $`ϵ_n`$ is calculated), substitute $`\lambda _n`$ and $`ϵ_n`$ in the radial oscillation equations (15-16) and solve the set of radial equations (14-17). The solution of these differential equations is found by transforming them into algebraic equations by means of finite differences on a staggered spatial mesh and applying matrix inversion similar to standard Henyey schemes for stellar evolution (e.g. Savonije & Papaloizou 1983). Hereby the forcing potential $`\mathrm{\Phi }_T(l,m)`$ is chosen to have the correct symmetry, i.e. for given $`(n,m)`$, we adopt $`l`$ accordingly: for g-modes $`l=m+n`$ and for q-modes $`l=m+|n2|`$. By scanning through a range of (initially) real forcing frequencies $`\overline{\sigma }`$ and solving the radial oscillation equations we find the possible oscillation modes (with a different number of radial nodes $`k`$) when the response becomes resonant for certain values $`\overline{\sigma }_k(\mathrm{\Omega }_s,n,m)`$. At resonance (that is what we are interested in) the stellar response becomes that of the free oscillation mode $`(k,n,m)`$ and appears virtually independent (apart from a fixed overall scaling factor) of the value adopted for $`l`$ in the forcing potential, as long as the equatorial symmetry of the forcing is consistent.
### 4.4 Search for unstable oscillation modes
The resonances with free oscillation modes are found by varying the real valued $`\overline{\sigma }`$ and maximizing the absolute value of the tidal torque integral (interpreting the forcing potential as that of a spherical harmonic component of an external point mass companion as in Savonije & Witte 2002):
$$𝒯_{nm}(\overline{\sigma },\mathrm{\Omega }_\mathrm{s})=\pi \zeta _{nm}f_0^{R_\mathrm{s}}m\left[\rho _{nm}^{}(r)\right]r^{l+2}dr$$
(22)
where $`\pi `$ results after the integration over $`\phi `$, $`f`$ is the constant in the tidal potential defined by Eq. (5), $`\mathrm{Im}`$ stands for imaginary part, and
$$\zeta _{nm}=\frac{\left[_1^1P_l^m(\mu )𝒳_n(\mu )d\mu \right]^2}{_1^1𝒳_n^2(\mu )d\mu }$$
(The tidal torque follows by multiplying $`𝒯_{nm}`$ with $`m`$). This procedure shows in fact immediately whether the resonance corresponds to a stable oscillation mode (if $`𝒯_{nm}`$ has same sign as $`\overline{\sigma }`$) or to an unstable one (when $`𝒯_{nm}`$ has opposite sign as $`\overline{\sigma }`$). This can be understood physically by considering a stable oscillation mode: when the forced oscillation is prograde ( $`\overline{\sigma }>0`$) the torque should lead to a spin up of the star and must be positive. For an unstable mode the phase difference between forcing and star has the opposite sign.
For a check of the stability properties of a mode and in order to estimate its growth or damping rate, we forced the star with a complex frequency $`\sigma `$ and searched with a numerical algorithm for that value of the imaginary part $`m(\sigma )`$ for which the integral over the star $`𝒥=_0^R\xi _r(r)\xi _r^{}(r)dr`$ becomes maximized. The search interval for $`m(\sigma )`$ was chosen symmetric about the zero value. The resonant forcing of stable oscillation modes becomes enhanced when the radiative damping is neutralized, so that $`𝒥`$ becomes maximal for a positive value of $`m(\sigma )`$ (identified as the ‘damping rate’), in analogy with the standard theory for a 1D forced linear harmonic oscillator with damping. For an unstable mode with ‘negative damping’, the enhanced resonant response is found for $`m(\sigma )<0`$ (the absolute value of which is identified as the ‘growth rate’). Except for modes near the boundaries of the instability bands, the integral $`𝒥`$ of the squared oscillation amplitude becomes many orders of magnitude larger when applying complex forcing frequencies to search for the maximum resonant response.
Please note that in this paper we consider $`\sigma `$ to be the real part of the forcing frequency, unless its complex character is made explicit, like in $`m(\sigma )`$, while $`\overline{\sigma }`$ always refers to the real (part of the) oscillation frequency (in the corotating frame).
### 4.5 $`\kappa `$-mechanism
The instabilities found in this paper are (checked to be) caused by the well known $`\kappa `$ valve-mechanism by which a fraction of the outflowing stellar thermal energy-flux is tapped by an increase of the radiative opacity during the compression phase of the oscillations and heat is released during the expansion phase. According to a semi-adiabatic analysis assuming constant radiative luminosity the opacity $`\kappa `$ has the correct absorbing behaviour in regions of the star where (e.g. Unno et al. unno89 (1989)):
$$\frac{\mathrm{d}}{\mathrm{dr}}\left(\kappa _T+\frac{\kappa _\rho }{\mathrm{\Gamma }_31}\right)>0$$
(23)
Regions where inequality (23) is fulfilled tend to drive the oscillation when a sufficient amount of thermal energy can be absorbed during the compression phase. For this the oscillation period should be comparable to the thermal timescale in the driving zone. Let us define a local ’adiabatic frequency’ $`\nu _{ad}=2\pi /\tau _{leak}`$ and use the leaking time $`\tau _{leak}(R_sr)^2\rho \kappa /c`$ as a crude estimate for the thermal timescale $`\tau _{th}`$ at radius $`r`$ based on random walk of photons towards the stellar surface. When the surface is approached the leaking time becomes ever shorter and $`\nu _{ad}`$ rises steeply, see e.g. Fig. 4. This figure shows the various frequencies as a function of the temperature in a 4 $`M_{}`$ main-sequence star with Z=0.02.
Note that from now on we normalize all frequencies on the star’s critical rotation frequency $`\mathrm{\Omega }_c=\sqrt{GM_s/R_s^3}`$ also when this is not indicated. In the surface regions where the oscillation frequency becomes much smaller than the ‘adiabatic frequency’ $`\nu _{ad}`$ the oscillations evidently are very non-adiabatic and no driving can occur when (23) attains positive values in these regions. For instability the rise of $`\nu _{ad}`$, where its curve in Fig. 4 intersects the actual (constant) $`\overline{\sigma }`$ should occur well into the driving region where (23) attains a significant positive value, i.e. in the opacity bump region near $`T2\times 10^5`$K, or rise steeply early in the damping region adjacent to it (so that the extra damping has no effect due to the already strong non-adiabaticity).
It is known (Berthomieu et al. berth78 (1978), Savonije et al. sav95 (1995)) that the traditional approximation yields qualitatively correct results for g-modes if the oscillation frequency $`\overline{\sigma }`$ is small compared to the Brunt-Väisälä frequency (thus outside the convective core), so that the horizontal oscillatory motion dominates. This is also a reasonable approximation for the q-modes with weak radial motion studied here. Fig. 4 shows that this requirement is amply satisfied for the 4$`M_{}`$ star.
### 4.6 Comparison with other stability analyses
We applied our calculation method to a non-rotating star and compared the results with those given in Gautschy and Saio (gautschy96 (1996)) and found an instability region for a 5$`M_{}`$ MS star very similar to the one shown in their Figure 2. We also compared our results for rotating stars with those of Lee (lee01 (2001)) who used a truncated series expansion in spherical harmonics to describe the non-adiabatic non-radial oscillations in a rotating star. Lee applied this scheme to a $`4M_{}`$ ZAMS star with $`X_c`$=0.7 and $`Z`$=0.02. We will compare the results for rotation speeds $`\mathrm{\Omega }_s`$=0.1 (Table 2) and $`\mathrm{\Omega }_s=0.3`$ (Table 3). Our 4$`M_{}`$ ZAMS stellar input model has $`{}_{}{}^{10}\mathrm{log}(L/L_{})`$=2.366 and $`{}_{}{}^{10}\mathrm{log}T_{\mathrm{eff}}`$=4.170 and thus is slightly hotter than Lee’s input model ($`{}_{}{}^{10}\mathrm{log}T_{\mathrm{eff}}`$=4.164).
The results presented in Tables 2-3 show that the two methods yield very similar instability intervals for a 4$`M_{}`$ ZAMS star. Only for the highest rotation rate we find slighly narrower instability domains for prograde g-modes on one hand and slighly more extended (to lower frequencies) intervals for retrograde g-modes on the other. This may be partly due to the slightly different input model.
It makes practically no difference to these results whether we adopt $`ϵ_n`$=0 or $`ϵ_n0`$, in the former case the instability interval is, for the case $`\mathrm{\Omega }_s`$=0.3, sometimes extended whereby $`k_{min}`$ decreases by one. The great benefit of the present calculation method is that there is no truncation problem in the solution of the angular eigenfunctions (although by neglecting the horizontal component of $`\mathrm{\Omega }_s`$, which seems reasonable for the considered modes). The horizontal part of the energy flux perturbation is treated only approximately, see equation (13). But, since the approach in this paper greatly simplifies the analysis of the stability of oscillation modes in rotating stars and provides us with an at least qualitatively correct picture of the linear stability properties of these modes, it seems a good initial approach. Townsend (town05 (2005)) used a similar approach (but neglected the horizontal diffusion altogether by adopting $`ϵ_n=0`$) to study the stability of g-modes in rotating SPB-stars in more detail than is done here. Where the results overlap they seem qualitatively consistent with each other, although Townsend finds unstable $`n`$=0 g-modes even in ZAMS stars somewhat less massive than 3$`M_{}`$ while the $`n`$=0 g-modes appear all stable in our 3$`M_{}`$ models, see below. This seems related to the higher effective temperatures of the lower mass stellar models used by Townsend as compared to the models used here.
## 5 Results
### 5.1 Unstable $`n=1`$ q-modes in a $`3M_{}`$ MS star (Z=0.02)
Before exploring the stability of more massive MS stars let us search for the lower mass limit for which overstable q-modes can exist. We studied the stability of two MS models with masses $`3M_{}`$ and $`2M_{}`$. For the $`2M_{}`$ MS star no q-mode instabilities (nor any unstable $`n`$=0 g-modes) could be found, so that the lower boundary of the q-mode instability region lies between 2$`M_{}`$ and 3$`M_{}`$, coinciding by the way with the lower boundary for SPB stars. It can be seen in Fig. 2 that for these lower stellar masses the Z-bump in the opacity near $`{}_{}{}^{10}\mathrm{log}T5.3`$ is rather modest and because it lies quite deep inside the star, where the local thermal timescale $`\tau _{th}`$ is relatively long, only very low frequency oscillation modes can be sufficiently non-adiabatic in this potential driving zone to give rise to overstability. It appears that the very low frequency quasi g-modes (or $`n`$=-1 ‘q-modes’) can indeed be excited in the 3$`M_{}`$ star, see Table 4 for a list of unstable modes. No unstable g-modes could be detected in a 3$`M_{}`$ MS star: the required low frequencies can only be reached for high radial order modes for which internal damping connected with short wavelengths is too strong.
Fig. 3 shows the oscillation frequencies in the corotating frame: $`\overline{\sigma }_k/(2\mathrm{\Omega }_s)`$ of $`n`$=-1 q-modes for several rotation rates of the $`X_c=0.6`$ and $`X_c`$=0.4 model. The (absolute value) of the frequencies in the corotating frame $`|\overline{\sigma }_k|`$ are for $`m>1`$ significantly lower than those of unstable normal g-modes. For given $`\mathrm{\Omega }_s`$ the mode frequencies $`|\overline{\sigma }_k|`$ are seen to decrease with increasing $`m`$-value, while for given $`m`$ the frequencies $`|\overline{\sigma }_k|`$ increase with rotation speed $`\mathrm{\Omega }_s`$. That $`|\overline{\sigma }_k|`$ increases with $`m`$ follows from equation (21) which tells us that the roots $`|x_r^1|`$ from where $`\lambda _{nm}>0`$ for $`n=1`$ modes shifts towards 0 (respectively $`x_r^1=1/2`$,$`1/3`$ and $`1/4`$ for $`m`$=1, 2 and 3), so that the lowest radial order mode frequency $`|\overline{\sigma }_0|`$ for a given $`m`$ is larger than that for $`m+1`$ and consequently $`|\overline{\sigma }_k(m)|>|\overline{\sigma }_k(m+1)|`$ for not too large $`k`$.
For $`m>1`$ the mode frequencies of adjacent radial orders $`k`$ are very densely packed, so that the oscillation periods $`P_k`$ of all q-modes in an instability interval $`(k_{min},k_{max})`$ can only be distinguished if the observations have good time coverage, see Table 4. This property of q-modes tends to give a much simpler observable oscillation spectrum in rotating stars than g-modes. Note also that although in the corotating frame the q-modes are retrograde modes, in the observer’s frame they are always prograde.
The location and extent of the instability intervals are very sensitive to the thermal timescale $`\tau _{th}`$ in the opacity Z-bump region. In the more evolved MS models, with cooler surface layers, this region lies deeper inside the star so that $`\tau _{th}`$ is longer and the steep rise of the ’adiabatic frequency’ curve takes place a bit further outwards in the star, at lower $`T`$, see dot-dashed curve in Fig. 2. Instabilities in the less evolved $`X_c`$=0.6 model with the shorter $`\tau _{th}`$ require somewhat higher oscillation frequencies for instability and can thus occur at higher rotation rates than in the more evolved $`X_c`$=0.4 model, see Table 4. For rotation speed $`\mathrm{\Omega }_s`$=0.3 instabilities can only be found in the $`X_c`$=0.4 model for $`m3`$ (i.e. for the lowest mode frequencies), while for still higher rotation rates (and thus higher mode frequencies) no q-mode instabilities could be found. We see in Table 4 that the modes with higher $`m`$ values (i.e. with lower frequencies $`|\overline{\sigma }_k|`$) become stable when $`\mathrm{\Omega }_s`$ decreases as the oscillations become too slow and hence too non-adiabatic.
### 5.2 Unstable modes in a $`4M_{}`$ evolved MS-star
#### 5.2.1 Unstable retrograde $`n`$=0 g-modes
Let us now turn to 4 $`M_{}`$ MS star models of decreasing central hydrogen abundance $`X_c`$. Fig. 4 shows the (modified) Brunt-Väisälä frequency, the ‘adiabatic frequency’ and the driving/damping zones, all as a function of temperature, for the slightly evolved 4$`M_{}`$ MS model with $`X_c=0.60`$. The opacity Z-bump near $`{}_{}{}^{10}\mathrm{log}T5.3`$ coincides with the region where non-adiabicity of low-frequency, $`|\overline{\sigma }/(2\mathrm{\Omega }_s)|<1`$, oscillations becomes noticeable. A little further out the driving is, according to inequality (23), replaced by damping but the extra thermal energy exchange for relevant frequencies is small because of existing strong non-adiabaticity. By comparing the dashed and dot-dashed curves it is seen that the thermal timescale $`\tau _{th}`$ in the driving zone increases significantly when the star evolves further from the ZAMS and its surface layers become cooler.
We determined the unstable n=0 g-modes with $`m`$=1 to $`m`$=4 in a 4$`M_{}`$ main-sequence star rotating at various angular speeds, for two evolution phases: $`X_c`$=0.6 and $`X_c`$=0.2, see Table 5. Note that the (retrograde) g-mode frequency $`|\overline{\sigma }_k|`$ increases with $`\mathrm{\Omega }_s`$ as the increasing Coriolis force renders the star ‘stiffer’ against oscillatory motions. In Fig. 5 it can be seen that, for a given rotation rate $`\mathrm{\Omega }_s`$, the $`n`$=0 g-mode frequencies $`|\overline{\sigma }_k|`$ increase with increasing $`m`$-values, contrary to the behaviour of $`n=1`$ q-modes in Fig. 3. This is the usual g-mode behaviour, remember that $`n`$=0 means $`m`$=$`l`$ in the non-rotating star terminology. The results in Table 5 show that in the little evolved MS model all combinations of $`\mathrm{\Omega }_s`$ and $`m`$ yield unstable g-mode intervals. However, for the evolved model with $`X_c=0.2`$ the opacity bump region has moved into the star with a corresponding longer local value for $`\tau _{th}`$. As a result the parameter combinations that correspond with higher mode frequencies, i.e. the larger $`m`$ values and higher rotation rates, no longer yield unstable g-modes because the oscillations have become too adiabatic for destabilization by the $`\kappa `$-mechanism. For $`\mathrm{\Omega }_s`$=0.5 all considered g-modes appear stable.
#### 5.2.2 Unstable $`n=1`$ q-modes
Let us now consider the behaviour of the $`n=1`$ q-modes in a 4$`M_{}`$ main-sequence star. We have extended the calculations to a rather high angular rotation speed of $`\mathrm{\Omega }_s=0.5`$ for which the neglect of centrifugal distortion is questionable just to see the effect of stronger Coriolis forces on the oscillation.
Fig. 6 shows the eigenfunction $`𝒳_n(\mu )`$ for the first and last unstable $`m`$=2 $`n`$=0 g-mode and $`m`$=2 $`n=1`$ q-mode, bracketing their instability domains, as listed in Tables 5 and 6 for $`\mathrm{\Omega }_s=0.3`$. The two bracketing unstable g-modes are at the boundary, respectively well inside the inertial range: with $`1/x1.09`$ and -0.60, respectively. For the latter low frequency g<sub>30</sub>-mode rotational effects are thus significant and its eigenfunction $`𝒳_n`$ exhibits the typical equatorial confinement of these modes: the eigenfunction has significant amplitude only for $`|\mu |<1/x`$. The two q-modes, by definition inside the inertial regime, have respectively the values $`1/x0.302`$ and $`0.271`$, but are seen to exhibit hardly any confinement toward low values of $`\mu `$ (these q-modes are odd and vanish at the equator).
The q-modes, in spite of their weak compressions also unstable due to the $`\kappa `$-mechanism, extent the instability region down to lower frequencies than that of the normal g-modes of Table 5. The q-modes have significantly smaller $`\lambda `$-values and their angular eigenfunctions do not exhibit the strong rotational confinement to the stellar equator for small values of $`|1/x|`$ shown by normal retrograde g-modes.
In Table 6 we list the unstable q-modes with $`n=1`$ and $`m`$=1 to 4 at various evolutionary stages on the main-sequence and for angular rotation speeds up to $`\mathrm{\Omega }_s`$= 0.5. The r-mode frequencies are proportional to $`\mathrm{\Omega }_s`$, and the same remains (roughly) true for the q-modes, although non-toroidal effects begin to become significant with increasing $`k`$.
In order to understand the shifting boundaries of the instability regions with $`m`$ and $`\mathrm{\Omega }_s`$ one should again realize that for a given rotation speed $`\mathrm{\Omega }_s`$ the absolute value of the oscillation frequency $`|\overline{\sigma }_k|`$ for $`n=0`$ g-modes increases with $`m(=l)`$, see Fig. 5, while Fig. 7 shows that for $`n=1`$ q-modes the reverse is true (for an explanation see section 5.1). It can be seen in Fig. 7, e.g. the case with $`m=1`$ and $`\mathrm{\Omega }_s`$=0.5, that sometimes not all radial orders in an instability domain are unstable according to our criteria (the mode amplitude should increase when the search for the resonant forcing frequency is made in the complex plane). In some cases the resonances remain too weak.
Fig. 8 shows the imaginary parts of the complex forcing frequency $`\sigma `$ for which the resonant response in a 4$`M_{}`$ model is maximized for both the unstable $`n=1`$ q-modes and $`n`$=0 g-modes with m-values in the range 1-3. As expected, for the quasi g-modes the highest growth rates are attained by the $`m`$=1 modes which have the largest $`\lambda `$ values and highest oscillation frequencies $`\overline{\sigma }`$ (see Fig. 7) and are thus closest in character to the g-modes. But in this little evolved stellar model ($`X_c=0.6`$) the driving layer is relatively close to the stellar surface and the (low-frequency) modes have rather small growth rates.
In the more evolved stellar model with $`X_c`$=0.4 and for the lowest rotation speed $`\mathrm{\Omega }_s`$=0.1 the $`m`$=2 and $`m`$=3 q-modes have oscillation frequencies too small for the $`\kappa `$ instability to be effective, the oscillations are too non-adiabatic in the $`\kappa `$-bump region. In the most evolved stellar model ($`X_c=0.1`$), on the other hand, the sudden increase of $`\tau _{ad}`$ occurs further out in the star, see Fig. 4. For $`\mathrm{\Omega }_s=0.5`$ this results in stable $`m`$=1 q-modes as the mode frequencies are in this case too high for efficient driving in the opacity bump region.
Note again that for $`m2`$ the subsequent radial orders of the $`n=1`$ q-modes in an instability interval are densely spaced in period. For normal g-modes similar close spacing of radial orders $`k`$ generally occurs for much larger values of $`k`$ for which the oscillations are significantly more confined to the equatorial plane of the rotating star and for which the strong radiative damping prevents instability.
### 5.3 Unstable $`n=2`$ q-modes in 3-4 $`M_{}`$ MS stars
It appears that in the evolved 3-4$`M_{}`$ models even unstable higher order $`n=2`$ quasi-toroidal modes are excited by the $`\kappa `$-mechanism. These modes have very small eigenvalues $`\lambda _2`$ and very low oscillation frequencies, see Tables 7-8 and Fig. 10. G-modes of similarly low oscillation frequencies would be of much higher radial degree and have much larger $`\lambda `$-values and be stable due to strong non-adiabaticity. In less evolved stellar models with $`X_c0.4`$ the $`n=2`$ q-modes appear all stable.
The angular eigenfunction $`𝒳_n(\vartheta )`$ for three (unstable) $`n=2`$ q<sub>17</sub>-modes with $`m`$=2, 3 and 4, all with frequencies $`|1/x|<0.115`$ (see Fig. 10), are depicted in Fig. 9. The eigenfunction $`𝒳_n(\vartheta )`$ determines the latitudinal distribution of temperature, density (and small radial velocity) over the stellar surface. The shape varies little as a function of $`m`$.
Figures 2 and 4 show that in the evolved stars the ‘adiabatic’ frequency rises steeply only near the outer edge of the opacity bump region. Low frequency q-mode oscillations apparently can still absorb sufficient heat during compression in the opacity bump zone, at least for rotation speeds $`\mathrm{\Omega }_s0.3`$. When the opacity derivatives $`\kappa _T`$ and $`\kappa _\rho `$ are put to zero the instabilities disappear.
For the higher stellar masses on the main sequence the opacity bump is closer to the surface and the non-adiabaticity too strong for destabilization of the very low frequency $`n=2`$ q-modes.
Fig. 10 shows that the absolute value of the $`n=2`$ mode frequencies (in corotating frame) $`|\overline{\sigma }_k|`$ as a function of azimuthal index $`m`$ behave differently from the $`n=1`$ modes shown in Fig. 7: the absolute value of the frequencies of e.g. $`m`$=3 modes are higher than those of $`m`$=2, while for $`k16`$ the same holds for $`m`$=4 with respect to $`m`$=3. This can again be explained by considering equation (21). For $`n=2`$ the roots $`x_r^1`$ have values $`1/6`$, $`1/6`$, $`3/20`$ and $`1/10`$ for $`m`$=1, 2, 3 and 4, respectively, i.e. lie closely together so that already for relatively small $`k`$ we get the behaviour stated above.
Fig. 11 shows the imaginary parts of the complex forcing frequency $`\sigma `$ for which the resonant response is maximized for the unstable $`n=2`$ q-modes with m-values in the range 2-4. The real parts of the oscillation frequencies are shown in Fig. 10. The very low frequency $`n=2`$ q-modes with small $`\lambda `$-values have significantly smaller growth rates than the $`n=1`$ q-modes. The $`m`$=3 modes attain the largest maximum growth rate.
### 5.4 Unstable modes in a $`6M_{}`$ MS star
#### 5.4.1 Unstable $`n`$=0 retrograde g-modes for Z=0.02
Fig. 12 shows the frequencies of unstable $`n`$=0 g-modes in the 6$`M_{}`$ model with $`X_c`$=0.6, rotating at speeds $`\mathrm{\Omega }_s`$=0.1 and 0.3. By comparing with Fig. 5 which depicts the same for a 4$`M_{}`$ stellar model it appears that in this more massive star the instability intervals are less extended and terminate at larger $`k`$-values ($`\mathrm{\Delta }k10`$), while the high frequency edge is shifted to slightly lower $`k`$ with $`\mathrm{\Delta }k2`$. This is caused by the shorter thermal timescale in the driving zone of the more massive star. Apart from the slight shift in $`k`$, the frequencies $`\overline{\sigma }_k`$ of the unstable g-modes are very similar for the 4$`M_{}`$ and 6$`M_{}`$ models.
Table 9 lists the g-mode instability intervals for various rotation rates and evolution stages of the 6$`M_{}`$ MS star. The table shows that we expect to observe many unstable g-modes with a retrograde character in observer’s frame (negative periods). By writing $`\sigma =\overline{\sigma }+m\mathrm{\Omega }_s`$ we can understand that the modes become prograde in the observer’s frame for larger $`m`$-values and or higher stellar rotation rates as listed in table 9. The unstable g-modes, both prograde and retrograde, occur in many wide instability intervals which is, however, not observed in SPB stars. We list here only the $`n`$=0 g-modes, but there exist also many higher order ($`n>0`$) unstable g-modes. It is a serious puzzle why all these apparently unstable g-modes are not observed.
#### 5.4.2 Unstable $`n=1`$ q-modes for Z=0.02
The frequencies of unstable $`n=1`$ q-modes detected for the 6$`M_{}`$ model with $`X_c`$=0.2 are shown in Fig. 14. The instability regions ($`k_{min}`$,$`k_{max}`$) are smaller than for the g-modes: due to the lower oscillation frequency of q-modes the low-frequency (high k) end is less extended as the opacity bump region becomes too non-adiabatic. Note that all unstable q-modes in Fig. 14 have $`|1/x|<0.45`$, while many of the unstable g-modes in the same stellar model (except the $`m`$=1 modes) lie outside the inertial range, see Fig. 12. In Table 10 we list the unstable $`n=1`$ modes with $`m`$=1, 2 or 3 for a $`6M_{}`$ main-sequence star. The frequency of the $`m`$=4 q-modes appears too low for destabilization by the $`\kappa `$ mechanism. With increasing stellar mass the opacity bump driving region for the instabilities moves closer to the stellar surface and the steep rise of $`\tau _{ad}`$ occurs in the 6$`M_{}`$ model at higher temperature, see Fig. 13. For the unevolved 6 $`M_{}`$ models (with $`X_c>0.4`$) the thermal timescale in most of the opacity bump region is therefore shorter than the oscillation periods (in corotating frame) of the q-modes, except for the fastest rotation rates considered. For this reason we find only $`m`$=1 instabilities and those only for $`\mathrm{\Omega }_s0.4`$. When the star evolves away from the ZAMS the effective temperature decreases and the driving zone moves inwards, so that the local thermal timescale increases. For the models with $`X_c<0.4`$ we therefore find more unstable q-modes and also instabilities at lower rotation rates, although not for $`\mathrm{\Omega }_s<0.3`$ as in a 4 $`M_{}`$ star: the frequencies are in that case too low.
#### 5.4.3 Unstable $`n=1`$ q-modes for Z=0.03
When the metal fraction is increased to Z=0.03 the opacity bump effect gets stronger and more q-modes become unstable, as can be seen by comparing Table 10 with Table 11. With the enhanced opacity bump the higher $`m`$ modes can also be destabilized, for high rotation rates $`\mathrm{\Omega }_s>`$0.3 even $`m`$=4 in the evolved $`X_c`$=0.2 stellar model.
Fig. 15 shows the imaginary parts of the complex forcing frequency $`\sigma `$ for which the resonant response is maximized for both a set of $`n=1`$ q-modes and $`n`$=0 g-modes for various m-values. Again the highest frequency $`m`$=1 q-modes show the largest growth rates.
It can be seen that the maximum growth rate of $`m`$=1 q-modes becomes almost comparable ($`20\%`$) to that of $`n`$=0 g-modes in this more evolved ($`X_c=0.2`$) main-sequence model.
### 5.5 Unstable $`n=1`$ modes in a $`8M_{}`$ MS star with Z=0.03 $``$ LPV in Be stars?
Fig. 16 shows the driving- and damping regions in a 8 $`M_{}`$ MS star with $`X_c=0.2`$ and Z=0.03. It can be seen that for this more massive star the opacity bump region is convectively unstable. Nevertheless, by assuming the convective flux is not heavily perturbed by the oscillation, we have sought for unstable modes, see Table 12. We have extended the calculations up to completely illegal rotation speeds for which the (neglected) rotational deformation of the star becomes strong in order to compare our results with observations of rapidly rotating Be-stars.
It is tempting to identify the observed line profile variability (LPV) in the well observed Be star $`\mu `$ Cen (Rivinius et al. riv01 (2001)) in terms of $`n=1`$ q-mode oscillations in a rapidly rotating B-star. The observed variations are indeed all prograde in the observer’s frame and fall in two narrow period intervals of about 0.50 days ($`m`$=2) and 0.28 days ($`m`$=3). Qualitatively, this is exactly what is found for the unstable $`n=`$-1 q-modes calculated here (see Table 12), although the exact periods do not match. But that is not surprising in view of the over-simplistic modelling of extremely rapidly (differentially) rotating Be-stars.
## 6 Summary and discussion
In this paper we study the stability of quasi g-modes (’q-modes’), a branch of oscillation modes that exists in between the normal retrograde g-modes and the r-modes in rotating stars. It appears that these retrograde (in corotating frame) low frequency oscillations can be destabilized by the $`\kappa `$-mechanism in the driving region connected with the opacity bump close to the stellar surface. From relatively low angular rotation speeds upwards the unstable q-modes with $`m>1`$ occur in a few narrow period bands (for different $`m`$-values) in which the different radial orders lie densely packed together. For $`m`$=1 the periods of q-modes are typically larger than a day, up to many days for low radial orders. In the observer’s frame all q-modes are prograde. The predicted observable q-mode spectrum is thus in sharp contrast to the predicted spectrum of unstable g-modes which occur in much wider period intervals, whereby modes of different radial order $`k`$ should be separately observable with even modest time resolution. The unstable q-modes seem to fit the observed rather scarce frequency spectrum in SPB stars and $`\beta `$ Cephei stars (e.g. Stankov et al. stank02 (2002)) better than normal g-modes. The observed closely spaced frequencies in e.g. the SPB star HD 160124 (Waelkens, waelkens91 (1991)) cannot be explained by unstable g-modes, but is typical for the q-mode spectrum. It is, however, doubtful whether the rotational confinement of g-modes towards the stellar equator in rotating stars can explain away the many predicted unstable g-modes (see Townsend town05 (2005)). It remains mysterious why we don’t observe the many predicted unstable g-modes and would instead observe the (weaker) unstable q-modes. It helps in any case that the unstable q-modes hardly have a tendency to focus on the equatorial region of rotating stars.
Although the more massive $`\beta `$ Cephei pulsators seem to show clustered periods similar to SPB stars (but with smaller periods) they cannot be explained straightforwardly in terms of q-modes because in stars more massive than 8$`M_{}`$ the thermal timescale in the the Z-bump driving region becomes too short to enable destabilization of the q-modes. Q-mode instability for these more massive $`\beta `$ Cephei pulsators would require the Z-bump region to be deeper inside the star than predicted by the current stellar evolution models.
Nevertheless, the observed LPV in Be stars like $`\mu `$ Cen is (qualitatively) consistent with the calculated properties of q-modes in a rapidly rotating 8M B-star: two narrow period intervals corresponding to two unstable series (for $`m`$=2 and $`m`$=3) of closely spaced q-modes, see previous section. One may speculate about the cause of the related outbursts in this Be-star in terms of q-modes. Excitation of the (negative angular momentum) retrograde q-modes in the star will accelerate the rotation of the stellar material in the driving region just beneath the stellar surface, perhaps bringing the surface layers up to ejection velocities. A good understanding of the Be phenomenon requires a much more detailed study with differential rotation including the centrifugal force and non-linear effects.
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# NEWS FROM HERWIG
## 1 Introduction
HERWIG$`^\mathrm{?}`$ is a Monte Carlo event generator for simulation of hadronic final states in lepton–lepton, lepton–hadron and hadron–hadron collisions. It incorporates important colour coherence effects in the final state$`^\mathrm{?}`$ and initial state$`^\mathrm{?}`$ parton showers, as well as in heavy quark processes$`^\mathrm{?}`$ and the hard process generation$`^\mathrm{?}`$. It uses the cluster$`^\mathrm{?}`$ hadronization model and a cluster-based simulation of the underlying event$`^\mathrm{?}`$. While earlier versions$`^\mathrm{?}`$ concentrated on QCD and a few other SM processes, recent versions contain a vast library of MSSM$`^\mathrm{?}`$ and other BSM processes. A review of current Monte Carlo event generators including HERWIG can be found in $`^\mathrm{?}`$.
We are currently in a period of intense activity, finalizing the HERWIG program and writing a completely new event generator, HERWIG++. In this very short contribution, I can do little more than mention the areas of progress and provide references to sources of more details.
## 2 HERWIG version 6.5
HERWIG version 6.5 was released$`^\mathrm{?}`$ in October 2002. Its main new features were an interface to the Les Houches Accord event format$`^\mathrm{?}`$, the hooks needed by the MC@NLO package$`^\mathrm{?}`$ and various bug fixes and minor improvements. It was advertised as the final fortran version of HERWIG before work switched to HERWIG++.
Despite this, the period since then has seen intense development with several new subversion releases and new features, most notably version 6.505, which featured an improved interface to the Jimmy generator for multiparton interactions, which I will discuss in more detail shortly. The most recent version is 6.507, which can be obtained from the HERWIG web site$`^\mathrm{?}`$.
Development of fortran HERWIG is now slowing, and the only new feature still foreseen is the implementation of matrix element corrections to the production of Higgs bosons, both SM and MSSM, preliminary versions of which have been discussed in $`^\mathrm{?}`$. Beyond this, the HERWIG collaboration has made a commitment to all running (and ceased) experiments to support their use of HERWIG throughout their lifetimes. Due to lack of manpower, making the same promise to the LHC experiments would divert too much effort away from support of HERWIG++, and we will only support their use of HERWIG until we believe that HERWIG++ is a stable alternative for production running.
## 3 Jimmy
Early versions of the Jimmy model$`^\mathrm{?}`$ generated jet events in photoproduction using a multiparton interaction picture. The recent update$`^\mathrm{?}`$ enables it to work efficiently as a generator of underlying events in high $`E_T`$ jet events and other hard processes in hadron–hadron collisions for the first time. For a given pdf set, the main adjustable parameters are PTJIM, the minimum transverse momentum of partonic scattering, and JMRAD(73), related to the effective proton radius. Varying these one is able to get a good description of the CDF data$`^\mathrm{?}`$ and other data held in the JetWeb database$`^\mathrm{?}`$ that are sensitive to underlying event effects in hard process events. However, a poor description of minimum bias data in which there is no hard scale is still obtained. This is probably due to the fact that PTJIM is a hard cutoff and there is no soft component below it; preliminary attempts to rectify this are encouraging$`^\mathrm{?}`$. It is interesting to note that with tunings that give equally good descriptions of current data, Jimmy predicts twice as much underlying event activity as PYTHIA at the LHC.
## 4 HERWIG++
The HERWIG program is now more than ten times the size it was when it was designed and is maintained by a collaboration of about ten authors. Its structure has become too unwieldy to maintain reliably and is too rigid to incorporate many of the physics improvements that have occurred to us recently. We therefore took the decision to write a completely new event generator, HERWIG++$`^\mathrm{?}`$, retaining the main features of HERWIG, angular ordering and cluster hadronization, but with a completely new design offering more flexible and scalable development, with the aim to have a reliable product throughout the lifetime of the LHC experiments.
### Use of ThePEG
At around the same time, the developers of PYTHIA took a similar decision$`^\mathrm{?}`$ and started work$`^\mathrm{?}`$ on PYTHIA 7, a replacement for PYTHIA. As part of this project, an extremely powerful framework for the administration of event generation was developed$`^\mathrm{?}`$. In order for this to be used by other Monte Carlo packages, it has been separated off as an independent Toolkit for high energy Physics Event Generation, ThePEG$`^\mathrm{?}`$. HERWIG++ is based on ThePEG, which offers a number of advantages: the administrative overhead is shared, while retaining completely independent physics implementations; users only need to learn to use one framework to use several different event generators; and can (with care!) mix modules from different event generators simply by selecting the appropriate components from the toolkit.
### Hard Interactions
A small library of basic $`22`$ processes is built in to ThePEG, a few more will be implemented in HERWIG++ using a new HELAS-like structure that has already been implemented, but we do not foresee ever developing a hard process library to the extent that we did in HERWIG. Instead, the plan is to provide a clean interface to external codes to generate the majority of hard processes: interfaces to AMEGIC++ and a Les Houches accord file reader already exist. Of course MC@NLO- and CKKW-style matrix element+parton shower matchings are also planned.
### Parton Showers
A completely new parton shower algorithm has been designed$`^\mathrm{?}`$, based on the quasi-collinear limit$`^\mathrm{?}`$, which gives a smooth suppression of forward radiation from massive partons, rather than a sharp dead-cone as in HERWIG. It also has advantages for matrix element matching, as it gives a smooth coverage of the soft limit, with the emission region from colour-connected jet pairs just touching, with no overlap or missing region. The new algorithm has been used for an interesting study of the theoretical uncertainties in Sudakov form factors$`^\mathrm{?}`$.
### Cluster Hadronization
The hadronization model is largely a rewrite, but with an improved treatment of the baryon sector, inspired by, but slightly different from, the one of Kupco$`^\mathrm{?}`$.
### Phenomenology for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}`$ Annihilation
The above components, together with a treatment of secondary decays that for the moment is an exact copy of HERWIG’s, are sufficient to give full simulation of $`\mathrm{e}^+\mathrm{e}^{}`$ annihilation events. The first phenomenological study was made in $`^{\mathrm{?},\mathrm{?}}`$. Although not a complete parameter tune, a first attempt yields a similar overall description of $`\mathrm{Z}^0`$ decays to HERWIG, with significant improvements in the yields of identified baryons. Particularly significant is the description of B meson production – in HERWIG heavy and light quark events could not be simultaneously well described, always yielding tension in parameter tunings, but in HERWIG++ the perturbative cutoff largely determines the B fragmentation function and its best-fit value also gives a good description of light-quark event shapes.
### Current Developments
Work is under way on the extension to hadron collisions. A simple model of the underlying event exists and a new model based on Jimmy plus a soft component is planned. The initial state parton shower exists, but needs further development and testing. A complete description of simple processes like Drell–Yan is anticipated for this summer.
A complete rewrite of the secondary decays is under way, with more sophisticated treatments of almost all decay modes: a general treatment of the spin structure and spin correlations; interference between hadronic resonances and non-resonant diagrams; and specialized decayers for important special cases. The decay tables themselves are stored in an external xml database and a direct link with the CEDAR project$`^\mathrm{?}`$ is planned. At present 448 particles with 2607 decay modes have been incorporated, loosely based on the particle data group tables. The new database allows the many massages that are needed, for example to make branching fractions add up to 100%, to be documented, with a star-rating system for how trustworthy they are.
### Future Outlook
We plan to complete the implementation and testing of initial state showers this summer and release a preliminary version to the LHC experiments. Next will come 2-jet production in hadron collisions and the first serious comparisons with Tevatron data. Development will continue on secondary decays and start on a multiparton interaction model of the underlying event and CKKW-type matching to mutijet matrix elements. HERWIG++ will be a stable alternative to HERWIG for first LHC analyses, allowing us a platform from which to further develop the theoretical framework.
## Acknowledgments
I am grateful to all the authors of Jimmy, HERWIG, HERWIG++ and ThePEG for their fruitful collaboration. Stefan Gieseke’s help in preparing the transparencies for this presentation is particularly appreciated.
## References
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# Non-linear ripple dynamics on amorphous surfaces patterned by ion-beam sputtering
## Abstract
Erosion by ion-beam sputtering (IBS) of amorphous targets at off-normal incidence frequently produces a (nanometric) rippled surface pattern, strongly resembling macroscopic ripples on aeolian sand dunes. Suitable generalization of continuum descriptions of the latter allows us to describe theoretically for the first time the main nonlinear features of ripple dynamics by IBS, namely, wavelength coarsening and non-uniform translation velocity, that agree with similar results in experiments and discrete models. These properties are seen to be the anisotropic counterparts of in-plane ordering and (interrupted) pattern coarsening in IBS experiments on rotating substrates and at normal incidence.
Ever since their earliest observation navez , the production of ripples on the surfaces of amorphous targets subject to ion-beam sputtering (IBS) at intermediate energies, has been fascinating due to the similarities with macroscopic ripples, like those produced underwater vortex\_exp , or on the surface of aeolian sand dunes dunas . Beyond the morphological resemblance, IBS ripples share many other properties with e.g. aeolian ripples, such as wavelength coarsening and pattern translation with time carter ; habenicht . Remarkably, while typical wavelengths of the latter are above 1 cm, the periodicity of IBS ripples is in the 100 nm range valbusa , these patterns having gained increased interest recently for applications in Nanotechnology, ranging from optoelectronic to catalytic applic . IBS ripples are produced on a wide class of substrates, from amorphous or amorphizable (silica, Si, GaAs, InP) to metallic targets (Cu, Au, Ag) valbusa . In view moreover of their implied loss of in-plane symmetry, IBS ripples provide interesting instances of systems hosting a competition between pattern forming and disordering mechanisms patt .
A successful description of the main features of IBS ripples was provided by Bradley and Harper (BH) BH , based on Sigmund’s linear cascade approximation of sputtering processes in amorphous or polycrystalline targets Sigmund . The linear equation derived by BH describes satisfactorily some properties of IBS ripples, such as their alignment with the ion beam as a function of the incidence angle to target normal $`\theta `$ \[wave vector parallel (perpendicular) to the projection of the ion beam for $`\theta <\theta _c`$ ($`\theta >\theta _c`$), for some threshold $`\theta _c`$\]. Other features, such as ripple stabilization or wavelength dependence with ion energy or flux, required non-linear extensions of BH’s approach Makeev ; park , leading to an anisotropic generalization of the well known Kuramoto-Sivashinsky (KS) equation patt ; Makeev . However, a notable limitation of the anisotropic KS (AKS) equation is its inability to predict ripples that coarsen with time, contradicting observations in many experiments and/or discrete models of IBS (see habenicht ; yewande and Refs. therein).
In this Letter, we introduce a “hydrodynamic” model dunas ; Aste ; prl for IBS ripple production at off-normal incidence. Time scale separation between microscopic processes (collision cascades, surface diffusion) and pattern evolution allows us to derive an improved equation for the surface height. The new non-linear terms appearing allow for ripple coarsening and pattern translation with non-uniform velocity, as seen in experiments and discrete models. Our theory has both the AKS equation Makeev and the normal incidence hydrodynamic theory prl ; normal as particular limits, and enables analysis of the important case of rotating substrates frost . In addition, our model may be important also to the context of ripples on aeolian sand dunes, where the standard 1D approximation requires validation, fully anisotropic 2D models being scarce dunas , as incidentally occurs in many other contexts within Pattern Formation bar\_nepomnyaschy2 .
During IBS of amorphous or semiconductor substrates, in which the ions amorphize the subsurface layer, incident ions lose their energy through random collision cascades in the bulk Sigmund . Near-surface atoms receiving enough energy and momentum to break their bonds are in principle sputtered away, although they may join the current of surface adatoms that are available to other relaxation mechanisms, such as surface diffusion, before incorporating back to the solid bulk. Within the so-called “hydrodynamic” approach to aeolian sand dunes dunas and ion-sputtered surfaces Aste ; prl , we consider two coupled fields, namely, the thickness of the mobile surface adatoms layer, $`R(𝐱,t)`$, (related with the density of mobile adatoms through the atomic volume) and the height of the bombarded surface above a reference plane, $`h(𝐱,t)`$. Their time evolution is provided by
$`_tR`$ $`=(1\varphi )\mathrm{\Gamma }_{ex}\mathrm{\Gamma }_{ad}+D^2R,`$ (1)
$`_th`$ $`=\mathrm{\Gamma }_{ex}+\mathrm{\Gamma }_{ad},`$ (2)
where $`\mathrm{\Gamma }_{ex}`$ and $`\mathrm{\Gamma }_{ad}`$ are, respectively, rates of atom excavation from and addition to the immobile bulk, $`(1\varphi )=\overline{\varphi }`$ measures the fraction of eroded atoms that become mobile, and the third term in Eq. (1) describes motion of mobile atoms along the surface as due to isotropic thermal diffusion ($`D`$ being a constant for amorphous materials). Even if all eroded atoms are sputtered away ($`\varphi =1`$), we assume a non-zero average fraction of mobile atoms, $`R_{eq}`$.
Under the assumption that nucleation events are more likely in surface protrusions, in analogy to the Gibbs-Thompson relation we have
$$\mathrm{\Gamma }_{ad}=\gamma _0\left[R(1+\gamma _{2x}_x^2h+\gamma _{2y}_y^2h)R_{eq}\right],$$
(3)
where $`\gamma _0`$ is the mean nucleation rate for a flat surface (on the $`xy`$ plane) and $`\gamma _{2x}`$, $`\gamma _{2y}0`$ describe variation of the nucleation rate with surface curvatures. Note that, in (3), the full thickness of the mobile adatoms layer is affected by the shape of the surface.
The rate at which material is sputtered from the bulk depends on the angle of incidence, ion and substrate species, ion flux, energy, and many other experimental parameters. If the beam direction is in the $`xz`$ plane, we have, following Makeev ; prl
$$\begin{array}{c}\mathrm{\Gamma }_{ex}=\alpha _0[1+\alpha _{1x}_xh+\alpha _{2x}_x^2h+\alpha _{2y}_y^2h\hfill \\ \hfill +\alpha _{3x}(_xh)^2+\alpha _{3y}(_yh)^2(_xh)(\alpha _{4x}_x^2h+\alpha _{4y}_y^2h)],\end{array}$$
(4)
where parameters reflect the dependence of $`\mathrm{\Gamma }_{ex}`$ on the local shape of the surface detalles , as described by more microscopic derivations such as BH or generalizations thereof Makeev ; Feix . Analogous, but not equal, geometrical couplings to the driving occur in aeolian sand dunes dunas , or in growth onto amorphous substrates Raible . Note the loss of reflection symmetry in the $`x`$ direction, but not in the $`y`$ direction. For a planar surface, atoms are sputtered from the bulk in a typical time of order $`\alpha _0^1`$.
The main difference between former models BH ; Makeev ; park and our present model, Eqs. (1)-(4), is that, in the latter, eroded material is allowed to redeposit locally, and there is an implicit viscous flow viscousflow in the amorphized layer through the evolution of $`R`$. These additional mechanisms are seen below to induce richer pattern dynamics than in BH ; Makeev ; park .
The linearized Eqs. (1)-(4) have solutions $`R^l=\widehat{R^l}e^{i𝐤𝐱+\omega _𝐤t}`$, $`h^l=\widehat{h^l}e^{i𝐤𝐱+\omega _𝐤t}`$, with a dispersion relation, $`\omega _𝐤`$, given in the long wavelength limit by detalles
$`e(\omega _𝐤)`$ $`=`$ $`ϵ\varphi \gamma _0(\alpha _{2x}k_x^2+\alpha _{2y}k_y^2)ϵ^2\varphi \overline{\varphi }\gamma _0\alpha _{1x}^2k_x^2`$
$`R_{eq}D(k_x^2+k_y^2)(\gamma _{2x}k_x^2+\gamma _{2y}k_y^2)+𝒪(ϵk^4),`$ (5)
$`m(\omega _𝐤)`$ $`=`$ $`ϵ\varphi \gamma _0\alpha _{1x}k_x+𝒪(k^3),`$ (6)
where $`ϵ\alpha _0/\gamma _0`$ is a dimensionless parameter; the erosion rate being much smaller than the nucleation rate, $`ϵ1`$ for typical experiments. Eq. (6) is a simple consequence of the asymmetry in the $`x`$ direction, induced by the incoming flux.
The surface morphology is dominated by the periodic pattern with wave vector $`𝐤^{max}`$ making Eq. (5) a positive maximum. It can be shown detalles that $`𝐤^{max}`$ is oriented along the $`\widehat{𝐱}`$ or $`\widehat{𝐲}`$ directions, as observed experimentally valbusa . Close to the instability threshold, before nonlinear terms are no longer negligible, one has $`k_{x,y}^{max}ϵ^{1/2}`$. Substituting this into Eqs. (5) and (6) provides us with estimations of the typical time and length scales of the pattern, that we employ to rescale $`X=ϵ^{1/2}x`$, $`Y=ϵ^{1/2}y`$, $`T_2=ϵ^2t`$ and $`T_1=ϵ^{3/2}t`$, and perform a multiple scale expansion of Eqs. (1), (2), in which $`R`$ can be adiabatically eliminated. To lowest non-linear order, we get detalles ; nota\_no\_redep
$`_th`$ $`=\gamma _xh+{\displaystyle \underset{i=x,y}{}}\{\nu _i_i^2h+\lambda _i^{(1)}(_ih)^2+\mathrm{\Omega }_i_i^2_xh`$
$`+\xi _i`$ $`(_xh)(_i^2h)\}+{\displaystyle }_{i,j=x,y}\{𝒦_{ij}_i^2_j^2h+\lambda _{ij}^{(2)}_i^2(_jh)^2\};`$
$`\gamma `$ $`=ϵ\varphi \gamma _0\alpha _{1x},\nu _x=ϵ\varphi \gamma _0\alpha _{2x}ϵ^2\overline{\varphi }\varphi \gamma _0\alpha _{1x}^2,`$
$`\nu _y`$ $`=ϵ\varphi \gamma _0\alpha _{2y},\lambda _i^{(1)}=ϵ\varphi \gamma _0\alpha _{3i},`$
$`\mathrm{\Omega }_i`$ $`=ϵ(\overline{\varphi }D\varphi R_{eq}\gamma _0\gamma _{2i})\alpha _{1x},\xi _i=ϵ\varphi \gamma _0\alpha _{4i},`$
$`𝒦_{ij}`$ $`=DR_{eq}\gamma _{2i}+ϵ\left[D\overline{\varphi }(\gamma _{2i}\alpha _{2j})+\varphi \gamma _0R_{eq}\gamma _{2i}\alpha _{2j}\right],`$
$`\lambda _{ij}^{(2)}`$ $`=ϵ\left[\overline{\varphi }D\varphi R_{eq}\gamma _0\gamma _{2i}\right]\alpha _{3j}.`$ (7)
Eq. (7) generalizes the AKS type equations obtained BH ; Makeev ; park within BH approach to IBS. While sharing the same reflection properties in the $`x`$ and $`y`$ directions and most of the terms on the rhs, both equations differ crucially by the presence here of the $`\lambda _{ij}^{(2)}`$ nonlinearities. Moreover, in the absence of redeposition ($`\varphi =1`$), $`\lambda _i^{(1)}`$ and $`\lambda _{ij}^{(2)}`$ have the same signs, making Eq. (7) nonlinearly unstable, as in the BH case prl ; Kim ; comments . Note, the linear dispersion relation of (7) matches Eqs. (5) and (6) above. Under normal incidence, parameters are isotropic and $`\alpha _{1x}=\alpha _{4x}=\alpha _{4y}=0`$, Eq. (7) reducing to that obtained in prl ; Kim , and in Raible . As a final remark, let us quote the form of Eq. (7) for sample rotation around the $`z`$ axis during bombardment (see e.g. Bradley ; frost ). Dynamics of $`h`$ are given by a different isotropic limit of Eq. (7), namely,
$`_th`$ $`=\nu _\mathrm{r}^2h𝒦_\mathrm{r}^4h+\lambda _\mathrm{r}^{(1)}(h)^2+\lambda _\mathrm{r}^{(2)}^2(h)^2`$
$`+\lambda _\mathrm{r}^{(3)}[(^2h)h]];`$ (8)
$`\nu _\mathrm{r}`$ $`=(\nu _x+\nu _y)/2,\lambda _\mathrm{r}^{(1)}=(\lambda _x^{(1)}+\lambda _y^{(1)})/2,`$
$`\lambda _\mathrm{r}^{(2)}`$ $`={\displaystyle \frac{1}{4}}{\displaystyle \underset{i,j=x,y}{}}\lambda _{i,j}^{(2)},\lambda _\mathrm{r}^{(3)}={\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j=x,y}{}}\lambda _{i,j}^{(2)}\delta _{i,j}\lambda _\mathrm{r}^{(2)},`$
$`𝒦_\mathrm{r}`$ $`=(3𝒦_{x,x}+3𝒦_{y,y}+𝒦_{x,y}+𝒦_{y,x})/8,`$
with parameters (susbscript “r” denotes “rotating”).
To the best of our knowledge, Eq. (7) is new, and indeed has a rich parameter space. Numerical integration (not shown) within linear regime retrieves all features of the ripple structure as predicted by the BH theory, i.e., dependence of the ripple wavelength with linear terms, and ripple orientation as a function of $`\theta `$. Entering the nonlinear regime, and as occurs in the AKS equation and its generalizations Makeev ; park , nonlinearities $`\lambda _i^{(1)}`$ lead to saturation of the pattern with constant wavelength and amplitude. In absence of these terms, the ripple wavelength grows indefinitely as $`\mathrm{}(t)t^n`$ with $`n=1/2`$ Raible\_b ; nota\_exponente until a single ripple remains in a finite simulation domain. As an anisotropic generalization of the ordering process observed for normal incidence prl ; normal , pattern coarsening requires the presence of $`\lambda _{i,j}^{(2)}`$, whose magnitude and mathematically correct sign comments are due to describing redeposition by means of the additional field $`R`$. When the values of these coefficients increase relative to $`\lambda _i^{(1)}`$, coarsening stops later, and the amplitude and wavelength of the pattern also increase. The coarsening exponent $`n`$ will take an effective value that will be larger the later coarsening stops, and may depend on simulation parameter values. For instance, we show in Fig. 1 snapshots of a numerical integration of Eq. (7) for relatively large $`\lambda _{i,j}^{(2)}`$. The apparent coarsening is quantified in the plot of $`\mathrm{}(t)`$ shown in Fig. 2$`(a)`$, compatible (after transient effects analogous to those in Raible\_b ) with $`n=0.19`$. Note the saturation of ripple wavelength at long times, together with saturation in amplitude, as shown in the plot of the surface roughness (rms width) $`W(t)`$ in Fig. 2$`(b)`$. These results are similar to those obtained experimentally for IBS of Si in carter . For a different experiment on Si, precise measurements of the coarsening exponent habenicht yield $`n=0.50(4)`$, no saturation having been observed in this system, as we would expect. Here, the dispersion velocity of the pattern was also measured, finding that it decays with ripple wavelength or, equivalently, with time. We have also observed the same trend in the dispersion velocity of the ripples shown in Fig. 1. Experiments also exist, e.g. for IBS of Si, in which ripple coarsening is absent or residual erlebacher , that would correspond to smaller $`\lambda _{i,j}^{(2)}`$ values in (7), see e.g. an example in Fig. 3, where $`\mathrm{}(t)\mathrm{log}t`$ approximately. Ripple coarsening has been also observed in a Monte Carlo model of IBS yewande , in which rules implement Sigmund’s theory. To our purpose, the main conclusion of this study is the correlation between an increasing $`\mathrm{}(t)`$ and non-uniform dispersion velocity, and on the parameter-dependent values of the coarsening exponent. Let us note that, on the experimental side, values of $`n`$ show a large scatter in the literature (see Refs. in valbusa ; yewande ).
Additional non-linear effects can be described by Eq. (7). A first one is production of grooves (as opposed to ripples), due to the loss of up-down symmetry induced by quadratic nonlinearities. Indeed, by changing the sign of $`\lambda _i^{(1)}`$, grooves replace ripples, see Fig. 1$`(d)`$; this calls for systematic experimental exploration. A second effect is related with cancellation modes (CM), known in the AKS equation rost\_krug ; park and other models of IBS Aste ; prl ; Kim . These are linearly unstable modes for which nonlinearities cancel one another exactly. While for equations like the one considered in prl ; Raible CM affect well-posedness, anisotropic systems rost\_krug ; park may remain better defined in the presence of CM. In our case, if reflection symmetry breaking terms are neglected in Eq. (7) and BH parameters are used, the same CM occur as in the AKS equation. Additional CM ensue between the $`\lambda _i^{(1)}`$ and $`\lambda _{i,j}^{(2)}`$ terms for appropriate relative signs Raible ; comments . We have verified numerically that the full Eq. (7) breaks down for the latter CM, but can support AKS-type CM. That these solutions are physically realizable or are artifacts of the small slope approximation made, remains to be assessed. Useful information on this issue might come from field experiments in aeolian sand dunes.
In closing, we mention IBS of metals as an immediate experimental domain to which the above results may be relevant valbusa , albeit differing in the degree of universality. There, however, the correct extension of BH theory is not yet clear, nor is its importance relative to anisotropic surface diffusion. We have taken preliminar steps in this context Feix , and expect to make progress in this direction soon.
###### Acknowledgements.
We thank L. Vázquez and R. Gago for discussions. This work has been partially supported by MECD (Spain), through Grants Nos. BFM2003-07749-C05, -01 and -05, and an FPU fellowship (J.M.-G.).
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# Three lines proof of the lower bound for the matrix rigidity
## 1 Introduction
### 1.1 Problem
Changing some entries of a complex matrix can reduce its rank. The rigidity of a matrix $`M`$ is the function $`R_M(r)`$, which for a given rank $`r`$, gives the minimum number of entries of $`M`$ which one has to change in order to reduce $`M`$’s rank to $`r`$ or less. More formally,
$$R_M(r)=\underset{rank(\stackrel{~}{M})r}{\mathrm{min}}\{weight(M\stackrel{~}{M})\},$$
where $`weight`$ denotes the number of non-zero entries. In other words, large rigidity shows that the matrix’s rank is stable under perturbations. It is easy to see that $`R_M(r)nr`$ for any full rank matrix $`M`$, because change of one entry reduces the rank by at most 1.
### 1.2 History
In this section I survey all known results for the matrix rigidity over infinite fields up to my best knowledge. There have been done a lot of work on rigidity on finite fields and on restricted and generalized versions of the rigidity problem as well .
The rigidity was defined by Valiant ; a similar notion independently was proposed by Grigoriev . The main motivation to study rigidity is that good lower bounds on rigidity would give important complexity results in other computational models, like linear algebraic circuits and communication complexity. For communication complexity $`(0,1)`$-matrices are especially important. Valiant showed $`R_M(r)(nr)^2`$ for ”almost all” matrices $`M`$. Pudlak and Rodl showed a similar result for $`(0,1)`$-matrices. However, to show a good lower bound for an explicit matrix still remains unsolved task.
The most interesting matrix probably is Hadamard matrix. Pudlak and Savicky showed that for any Hadamard matrix $`H`$, $`R_H(r)=\mathrm{\Omega }(\frac{n^2}{r^4\mathrm{log}^2r})`$, Razborov improved their result to $`R_H(r)=\mathrm{\Omega }(\frac{n^2}{r^3\mathrm{log}r})`$. Grigoriev and Nisan independently observed an easy method to get lower bound for any *totally non-singular* matrix $`M`$ (i.e. a matrix in which all submatrices are non-singular) $`R_M(r)=\mathrm{\Omega }(\frac{n^2}{r})`$. Similar strategy was used by Alon to improve rigidity of Hadamard matrix $`R_H(r)=\mathrm{\Omega }(\frac{n^2}{r^2})`$, Lokam to give an alternate proof of the same result, Kashin and Razborov to prove $`R_H(r)\frac{n^2}{256r}`$ and de Wolf to use quantum information theoretical arguments to give a neat proof of $`R_H(r)\frac{n^2}{4r}`$ (the last result holds for any orthogonal matrix where all entries have the same magnitude, including Discrete Fourier Transform).
There are results for other matrices as well. Razborov showed $`R_M(r)=\mathrm{\Omega }(\frac{n^2}{r})`$ if $`M`$ is the generalized Fourier Transform matrix or the inverse of the Vandermonde matrix. Kimmel and Settle gave the lower of the rigidity of the triangular matrix $`T`$, their result in simplified form looks like $`R_T(r)\mathrm{\Omega }(\frac{n^2}{r})`$. Independently, Pudlak and Vavrin determined the exact value of $`T`$, particulary, for a large $`n`$ but small $`r`$ it is like $`R_T(r)\frac{n^2}{4r}`$. Pudlak showed that $`R_M(r)=\mathrm{\Omega }(\frac{n^2}{r})`$ if $`M`$ belongs to a class of matrices called Densely Regular, that includes triangular matrix, Vandermonde matrices, shifters and parity shifters. Shokrollahi et al. showed that $`R_C(r)=\mathrm{\Omega }(\frac{n^2}{r}\mathrm{log}\frac{n}{r})`$ for a Cauchy matrix $`C`$. Codenotti et al. studied the rigidity of some matrices under combinatorial assumptions.
Lokam gives some quadratic lower bounds for ”less explicit” matrices. Landsberg et al. gave geometrical interpretation of matrix rigidity.
However, these results do not give a superlinear rigidity for an explicit matrix when $`r=O(n)`$. Lokam observes that a method used in all those results (and this paper as well) by getting ”candidate” matrices that are close to full rank does not give a results like $`R(r)=\omega (\frac{n^2}{r}\mathrm{log}\frac{n}{r})`$.
Codenotti gives a survey paper on the matrix rigidity problem as well as some interesting problems.
In this paper we give a simple proof in ”three lines” of $`R_S(r)\frac{n^2}{4r}`$ for any Sylvester matrix $`S`$ (special case of Hadamard matrix). The same proof works for other ”well behaved” matrices, like Discrete Fourier Transform. However, our main contribution is the simplicity of the proof.
### 1.3 Matrices
If $`A=(a_{ij})`$ and $`B=(b_{kl})`$ are matrices of size $`m\times n`$ and $`p\times q`$ respectively, the *Kronecker product* $`AB`$ is the $`mp\times nq`$ matrix made up of $`p\times q`$ blocks, where the $`(k,l)`$ block is $`b_{kl}A`$.
*Sylvester matrix* $`S(n)`$ of order $`n:=2^k`$ is $`n\times n`$ matrix made by iterating Kronecker product of $`k`$ copies of the following matrix
$$S(2)=\left(\begin{array}{cc}+& +\\ +& \end{array}\right)$$
where $`+`$ and $``$ denotes $`+1`$ and $`1`$ respectively.
For example,
$$S(4)=\left(\begin{array}{cccc}+& +& +& +\\ +& & +& \\ +& +& & \\ +& & & +\end{array}\right)$$
Sylvester matrices are special case of Hadamard matrices. A real valued matrix $`H`$ is called *Hadamard matrix* iff $`HH^T=nI`$.
*Discrete Fourier Transform* is $`n\times n`$ matrix $`FN(n)=(f_{jk})`$ defined by $`f_{jk}:=\omega ^{(j1)(k1)}`$, where $`\omega :=e^{\frac{2\pi i}{n}}`$ and $`i:=\sqrt{1}`$.
## 2 Proof
###### Theorem 1
If $`S(n)`$ is a Sylvester matrix and $`rn/2`$ is a power of $`2`$ then
$$R_{S(n)}(r)\frac{n^2}{4r}.$$
In other words, for any $`n\times n`$ matrix $`\stackrel{~}{S}`$ such that $`rank(\stackrel{~}{S})r`$ holds
$$weight(S(n)\stackrel{~}{S})\frac{n^2}{4r}.$$
###### Proof
Assume the opposite, $`weight(S(n)\stackrel{~}{S})<\frac{n^2}{4r}`$. Let uniformly divide $`\stackrel{~}{S}`$ in $`(\frac{n}{2r})^2`$ submatrices $`\stackrel{~}{S}_{ij}`$ of size $`2r\times 2r`$. By a counting argument, there exists $`i,j`$ s.t. $`weight(S(2r)\stackrel{~}{S}_{ij})<r`$. Thus, $`rank(\stackrel{~}{S})rank(\stackrel{~}{S}_{ij})>2rr=r`$. Contradiction. ∎
The same proof works for $`R_{FT(n)}(r)\frac{n^2}{4r}`$, where $`FT(n)`$ denotes $`n\times n`$ Discrete Fourier Transform matrix, because DFT matrix where columns with even indexes are written first is represented as a matrix
$$\left(\begin{array}{cc}FT(n/2)& \omega ^jFT(n/2)\\ FT(n/2)& \omega ^{jn/2}FT(n/2)\end{array}\right)$$
where $`j`$ denote the index of a row. Since rows of a $`FT(n)`$ are orthogonal and multiplication by some constant does not change this property, each submatrix can be recursively divided again and again, by getting full rank submatrices. This is the only property of matrices we need in the proof.
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# Jet quenching
## 1 Introduction
Recent results at high transverse momentum from RHIC$`^\mathrm{?}`$ on the inclusive particle suppression and the absence of away-side two particle correlations in central heavy ion collisions, together with the negative-effect results from the reference deuteron-gold run, suggest the formation of a very dense partonic medium with which the triggered particles interact strongly. The nature of this medium is still unknown and the study of its properties is the main goal of the experimental program of high-energy heavy ion collisions. More differential measurements of particles with high transverse momentum will give very valuable information as explained in the next sections.
In heavy ion collisions, the particles produced perturbatively at high transverse momentum are expected to be uncorrelated from the small transverse momentum bulk. At the same time, this high-multiplicity state is expected to form a deconfined and thermalized medium which modifies the properties of the parton shower developed by the high-$`p_t`$ particles. The formalism to compute the medium-induced gluon radiation has been developed using several techniques and different approximations $`^\mathrm{?}`$. Apart from details, most of the main results depend on coherence effects which suppress gluon radiation at small transverse momentum and/or large energies $`k_t^2\widehat{q}L`$, $`\omega \widehat{q}L^2`$, where the transport coefficient $`\widehat{q}`$ depends on properties of the medium as the density. This results in the well-known quadratic dependence of the radiative energy loss with the traverse length of the medium. Moreover, the medium-modification of the transverse momentum spectrum of radiated gluons translates into a jet broadening $`k_t\mathrm{\Delta }EL/\alpha _s`$ $`^\mathrm{?}`$. One of the main predictions of this formalism is, then, the broadening of the associated gluon radiation when compared with the evolution in the vacuum.
## 2 Energy loss and jet quenching
The medium-induced gluon radiation spectrum $`\omega dI/d\omega `$ depends on the length of the medium and the transport coefficient $`\widehat{q}`$. In the absence of a more elaborated formalism, taking into account interference effects on multiple gluon radiation, the independent gluon emission approximation is usually taken. In this way, the probability that an additional energy $`\mathrm{\Delta }E`$ is radiated by medium effect is given by the quenching weights $`^{\mathrm{?},\mathrm{?}}`$ $`P(\mathrm{\Delta }E)`$ and the medium-modified fragmentation functions are modeled by the convolution $`D_{ih}^{\mathrm{med}}=P(\mathrm{\Delta }E)D_{ih}`$. These medium-modified fragmentation functions can be used to compute the cross section for the production of a hadron $`h`$ through the perturbative expression \[for precise definitions of the convolutions see e.g. $`^\mathrm{?}`$\]
$$\frac{d\sigma ^{AAh}}{dp_t}f_i^A(x_1,Q^2)f_j^A(x_2,Q^2)\sigma ^{ijk}D_{ih}^{\mathrm{med}}(z,\mu _F)$$
(1)
The strategy is then to fit the best value of $`\widehat{q}`$ that reproduces the experimental suppression measured by the ratio
$$R_{AA}(p_t)=\frac{dN^{AA}/dp_t}{N_{\mathrm{coll}}dN^{pp}/dp_t}.$$
(2)
The transport coefficient $`\widehat{q}`$ is proportional to the medium density. Thus, comparing the different $`\widehat{q}`$ obtained by applying this procedure to different systems, information about the density of the media is obtained. One limitation, however, appears when the medium is very dense and the suppression is so strong that the effect is dominated by surface emission $`^{\mathrm{?},\mathrm{?}}`$. In this case, the measure gives only a lower limit for the transport coefficient $`^\mathrm{?}`$. In order to improve the determination of the medium properties, one possible solution is to study the case of heavy quarks $`^\mathrm{?}`$. In this case, the radiation is suppressed by mass terms $`^\mathrm{?}`$ and, hence, the effect is smaller. In Fig.1 the prediction $`^\mathrm{?}`$ for the suppression of electrons from the decay of charm quarks at RHIC is presented using the transport coefficient obtained from the light meson case.
## 3 Medium-modified jet shapes
The structure of the jets is expected to be strongly modified when developed in a medium. The larger emission angle of the medium-induced spectrum translates into a broadening of the jet shapes. Although the broadening in energy could remain small for moderate transport coefficients, the intrajet multiplicity distribution is expected to present a harder spectrum in the transverse momentum with respect to the jet axis $`^\mathrm{?}`$ (see Fig. 2). This situation would be ideal for the study of medium-modified jet shapes at the LHC as i) it would allow for a good calibration of the jet energy (essential in order to study the jet properties) for the moderate values of the jet cone ($`R0.3`$) to be measured in the high-multiplicity environment of a heavy ion collision; and ii) the broadening produced in the intrajet multiplicities would be sizable enough to measure with high precision the medium effects.
A more interesting situation is, however, when the high-$`p_t`$ particle travels through a flowing medium. In this case, the flow introduces a preferred direction in the medium-induced gluon radiation, producing asymmetric jet shapes $`^\mathrm{?}`$ (see Fig. 3). These asymmetries are, in this way, a measurement of the flow field in a medium. Interestingly, preliminary results $`^\mathrm{?}`$ from RHIC on high-$`p_t`$ two particle correlations show a strong elongation of the jet-like signal in the longitudinal direction for central gold-gold collisions. In the spirit of the effects shown in Fig. 3, this elongation is produced by the strong longitudinal flow present in these collisions.
## References
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# An obstruction for the mean curvature of a conformal immersion 𝑆^𝑛→ℝ^{⋉⊉⊮}
(Date: June 2005)
## Abstract.
We prove a Pohozaev type identity for non-linear eigenvalue equations of the Dirac operator on Riemannian spin manifolds with boundary. As an application, we obtain that the mean curvature $`H`$ of a conformal immersion $`S^n\mathrm{}^{\mathrm{}\mathrm{}\mathrm{}}`$ satisfies $`_XH=0`$ where $`X`$ is a conformal vector field on $`S^n`$ and where the integration is carried out with respect to the Euclidean volume measure of the image. This identity is analogous to the Kazdan-Warner obstruction that appears in the problem of prescribing the scalar curvature on $`S^n`$ inside the standard conformal class.
MSC 2000: 53A27, 53A30, 35J60
Let $`(M,g)`$ be a compact Riemannian manifold with a conformal vector field $`X`$. Given a function $`s`$ on $`M`$, then it is a classical question to ask whether $`s`$ is the scalar curvature of a metric $`\stackrel{~}{g}`$ conformal to $`g`$. The determination of the set of all such functions $`s`$ is still open, although several partial results are known, in particular there are necessary conditions that $`s`$ has to satisfy in order to be a scalar curvature.
On the one hand there are topological obstructions. If for example $`M`$ is spin and has non-vanishing $`\widehat{A}`$ genus, then the scalar curvature of any metric on $`M`$ has either to be negative somewhere or the Ricci curvature vanishes everywhere on $`M`$.
However, if one fixes the conformal class $`[g]`$ as described above, there are further obstructions that arise from conformal vector fields. For example if $`M`$ is $`S^n`$ with the standard conformal structure, Kazdan and Warner \[KW75\] derived a famous obstruction. A slightly stronger version of this obstruction due to Bourguignon and Ezin \[BE87\] is described in the following theorem.
###### Theorem 1.
Let $`X`$ be a conformal vector field on the compact manifold $`(M,g)`$. If $`s`$ is the scalar curvature of a metric $`\stackrel{~}{g}=u^{4/(n2)}g`$, then
$$_M_Xsdv_{\stackrel{~}{g}}=0$$
where $`dv_{\stackrel{~}{g}}=u^{\frac{2n}{n2}}dv_g`$ is the volume measure associated to $`\stackrel{~}{g}`$.
Tightly related to the Kazdan-Warner obstruction is the Pohozaev identity. Let $`\mathrm{\Omega }`$ be a star-shaped open set of $`\mathrm{}^{\mathrm{}}`$ ($`n\mathrm{}`$) with smooth boundary. We denote by $`\mathrm{\Delta }=_{i=1}^n_{ii}`$ the Laplacian on $`\mathrm{}^{\mathrm{}}`$. Let $`uC^2(\overline{\mathrm{\Omega }})`$ be a positive solution of $`\mathrm{\Delta }u=u^{p1}`$ on $`\mathrm{\Omega }`$ with $`u_{|\mathrm{\Omega }}0`$. The vector field $`X=_{i=1}^nx^i_i`$ is conformal. If one uses similar methods as in the proof of the Kazdan-Warner obstruction, then one obtains the Pohozaev identity (\[Po65\]) which asserts that:
$$\left(1\frac{n}{2}+\frac{n}{p}\right)_\mathrm{\Omega }u^p=\frac{1}{2}_\mathrm{\Omega }\nu ,X(_\nu u)^2$$
(1)
where $`\nu `$ resp. $`_\nu `$ is the outer normal vector resp. the outer normal derivative on $`\mathrm{\Omega }`$. One among many important consequences of this inequality is that no non-trivial solutions exist if $`p\frac{2n}{n2}`$. Another application is an alternative proof of the Kazdan-Warner obstructions in the case that $`(M,g)`$ is the sphere with the standard conformal structure \[DR99\].
In the present short article, we establish a similar identity for the classical Dirac operator $`D`$. We derive this identity on manifolds with boundary in order to admit future Pohozaev type applications. Then, we will specialize to compact manifolds without boundary, where we will derive a Kazdan-Warner type obstruction for the mean curvature of a conformal immersion $`S^2\mathrm{}^{\mathrm{}}`$.
Our main theorem is:
###### Theorem 2.
Let $`(M,g,\chi )`$ be a compact Riemannian spin manifold of dimension $`n`$ with boundary $`M`$ (possibly equal to $`\mathrm{}`$) and with Dirac operator $`D`$. We assume that there exists a smooth spinor field $`\psi `$ which satisfies for some $`p>1`$,
$$D\psi =H|\psi |^{p2}\psi ,HC^{\mathrm{}}(M).$$
(2)
Furthermore, we assume that $`X`$ is a conformal vector field on $`M`$. Then, we have the following Pohozaev type identity:
$$_M\nu _X\psi ,\psi =\frac{p2}{p}_MH|\psi |^pg(X,\nu )+\left(1\frac{p2}{p}n\right)_MH\beta |\psi |^p+\frac{2}{p}_M(_XH)|\psi |^p$$
where $`\nu `$ denotes the outward pointing normal vector along $`M`$, and where $`,`$ denotes the real scalar product on spinors.
Proof: The flow associated to the conformal vector field $`X`$ will be denoted as $`\alpha ^t`$. If $`p`$ is in the interior of $`M`$, then $`\alpha ^t(p)`$ exists for times $`t`$ close to $`0`$. For any $`t\mathrm{}`$ let $`f^t`$ be the conformal scaling function of $`\alpha ^t`$, i.e. $`(d\alpha ^t)_p`$ is $`f^t(p)`$ times an isometry from $`T_pM`$ to $`T_{\alpha ^t(p)}M`$. Let $`\alpha _{}^t:\mathrm{\Sigma }_pM\mathrm{\Sigma }_{\alpha ^t(p)}M`$ be the spinor identification map as constructed in \[Ht74, Hi86, BG92\]. In particular, this map has the pointwise properties that
$$|\alpha _{}^t(\psi )|=|\psi |$$
and the following transformation formula for conformal changes of the metric. Let $`\phi \mathrm{\Gamma }(\mathrm{\Sigma }M)`$ be a spinor field. For $`t`$ close to $`0`$, we then define the map $`\alpha _\mathrm{\#}^t:\mathrm{\Gamma }(\mathrm{\Sigma }M)\mathrm{\Gamma }(\mathrm{\Sigma }\stackrel{~}{M})`$, $`\alpha _\mathrm{\#}^t(\phi ):=\alpha _{}^t\phi \alpha ^t`$, where $`\stackrel{~}{M}`$ is $`M`$ without an open neighborhood of the boundary.
Then
$$D\alpha _\mathrm{\#}^t\left((f^t)^{\frac{n1}{2}}\psi \right)=\alpha _\mathrm{\#}^t((f^t)^{\frac{n+1}{2}}D\psi ).$$
Now we assume that $`\psi `$ satisfies (2), and we obtain
$$D\alpha _\mathrm{\#}^t\left((f^t)^{\frac{n1}{2}}\psi \right)=\alpha _\mathrm{\#}^t\left((f^t)^{\frac{n+1}{2}}H|\psi |^{p2}\psi \right).$$
Deriving with respect to $`t`$ at $`t=0`$ yields
$`{\displaystyle \frac{n1}{2}}D\beta \psi +D{\displaystyle \frac{d}{dt}}|_{t=0}\alpha _\mathrm{\#}^t\psi `$ $`=`$ $`{\displaystyle \frac{n+1}{2}}H\beta |\psi |^{p2}\psi +H|\psi |^{p2}{\displaystyle \frac{d}{dt}}|_{t=0}\alpha _\mathrm{\#}^t\psi `$ (4)
$`+(p2)H{\displaystyle \frac{d}{dt}}|_{t=0}\alpha _\mathrm{\#}^t\psi ,\psi |\psi |^{p4}\psi (_XH)|\psi |^{p2}\psi .`$
where $`\beta :=\frac{d}{dt}|_{t=0}f^t`$. We reformulate using definition of the *Lie derivative of spinor fields in the direction $`X`$* \[BG92\], i.e.
$$_X(\psi )=\frac{d}{dt}|_{t=0}\alpha _\mathrm{\#}^t(\psi ).$$
(5)
Together with $`D\beta \psi =\beta D\psi +\beta \psi `$ and (2) one then concludes that
$`{\displaystyle \frac{n1}{2}}\beta \psi +D_X\psi `$ $`=`$ $`H\beta |\psi |^{p2}\psi +H|\psi |^{p2}_X\psi `$ (7)
$`+(p2)H_X\psi ,\psi |\psi |^{p4}\psi +(_XH)|\psi |^{p2}\psi .`$
After multiplication with $`\psi `$, the $`\beta \psi `$-term vanishes, and we obtain
$$D_X\psi ,\psi =(p1)H|\psi |^{p2}_X\psi ,\psi +H\beta |\psi |^p+(_XH)|\psi |^p.$$
The product rule for the Lie derivative tells us that
$$|\psi |^{p2}_X\psi ,\psi =\frac{1}{2}|\psi |^{p2}_X|\psi |^2=|\psi |^{p1}_X|\psi |=\frac{1}{p}_X|\psi |^p.$$
Hence, we obtain
$$D_X\psi ,\psi =\frac{p1}{p}H_X|\psi |^p+H\beta |\psi |^p+(_XH)|\psi |^p.$$
Strictly speaking, this equation is valid in the interior, but it extends to the boundary by continuity. Now, we integrate over $`M`$. With partial integration for the Dirac operator one obtains
$$_MD_X\psi ,\psi =_M_X\psi ,D\psi +_M\nu _X\psi ,\psi =_MH\underset{=\frac{1}{p}_X|\psi |^p}{\underset{}{_X\psi ,|\psi |^{p2}\psi }}+_M\nu _X\psi ,\psi .$$
This yields
$$_M\nu _X\psi ,\psi =\frac{p2}{p}_MH_X|\psi |^p+_MH\beta |\psi |^p+_M(_XH)|\psi |^p.$$
Using $`div(H|\psi |^pX)=(_XH)|\psi |^p+H_X|\psi |^p+H|\psi |^pdivX`$ and $`divX=n\beta `$ we obtain
$$_M\nu _X\psi ,\psi =\frac{p2}{p}_MH|\psi |^pg(X,\nu )+\left(1\frac{p2}{p}n\right)_MH\beta |\psi |^p+\frac{2}{p}_M(_XH)|\psi |^p.$$
###### Examples 3.
1.) Let $`\mathrm{\Omega }`$ be domain in $`\mathrm{}^{\mathrm{}}`$ with smooth boundary, let $`X=r_r=x^i_i`$, and we will assume that $`H=\lambda `$ is constant. Then $`\beta 1`$ and we obtain
$$_\mathrm{\Omega }\nu _X\psi ,\psi =\lambda \frac{p2}{p}_\mathrm{\Omega }X,\nu |\psi |^p+\lambda \left(1\frac{p2}{p}n\right)_\mathrm{\Omega }|\psi |^p.$$
This inequality bears many analogies to equation (1). In particular, the constant $`1\frac{p2}{p}`$ before the integral over $`\mathrm{\Omega }`$ vanishes if $`p`$ takes the value $`p=2n/(n1)`$. This value plays the same role in non-linear Dirac equations as the value $`p=2n/(n2)`$ does for the Laplace operator.
2.) If $`M`$ is a closed manifold and $`X`$ is a conformal vector field, then for $`p=2n/(n1)`$ we obtain
$$_M(_XH)|\psi |^p=0.$$
###### Corollary 4.
\[Kazdan-Warner type obstructions\] Let $`f:S^n\mathrm{}^{\mathrm{}\mathrm{}\mathrm{}}`$, $`n2`$, be a conformal immersion (possibly with branching points of even order in the case $`n=2`$). We denote by $`H:S^n\mathrm{}`$ the mean curvature of this immersion. Then, for any conformal vector field $`X`$ we have
$$_{S^n}(_XH)f^{}(d\mu )=0$$
where $`d\mu `$ is the volume element on $`f(S^n)`$ induced from the euclidean metric on $`\mathrm{}^{\mathrm{}}`$. In particular, $`_XH`$ changes sign.
The corollary is particularly interesting in dimension $`n=2`$. If $`f:S^2\mathrm{}^{\mathrm{}}`$ is any immersion, then after possibly composing with a diffeomorphism $`S^2S^2`$, we can assume that $`f`$ is conformal.
The corollary is analogous to results in \[KW75\], \[BE87\] and \[DR99\].
Proof: Let $`\psi `$ be parallel spinor on $`\mathrm{}^{\mathrm{}\mathrm{}\mathrm{}}`$. Then, as proven in \[KS96, Ba98, Fr98\], the restriction of $`\psi `$ on $`\mathrm{\Sigma }`$ satisfies equation (2) with $`p=2n/(n1)`$, and $`|\psi |^pd\nu =f^{}(d\mu )`$ where $`d\nu `$ is the standard volume element on $`S^n`$. Since this equation is conformally invariant we obtain a solution of (2) on $`S^n`$ equipped with the standard metric. The corollary then immediately follows from example (2) above.
###### Example 5.
Let $`x_3:S^2\mathrm{}`$ be the third component of the standard inclusion. One shows that $`X:=gradx_3`$ is a conformal vector field on $`S^2`$, where the gradient is taken with respect to the standard metric on $`S^2`$. Then for any $`C\mathrm{}`$ one has $`_X(x_3+C)=g(gradx_3,gradx_3)0`$. Hence $`x_3+C:S^2\mathrm{}`$ is not the mean curvature of a conformal immersion.
Authors’ addresses:
Bernd Ammann and Emmanuel Humbert,
Institut Élie Cartan BP 239
Université de Nancy 1
54506 Vandoeuvre-lès -Nancy Cedex
France
M. Ould Ahmedou
Mathematisches Institut
der Universität Tübingen
Auf der Morgenstelle 10
72076 Tübingen
Germany
E-Mail:
ammann at iecn.u-nancy.fr, humbert at iecn.u-nancy.fr and ahmedou at analysis.mathematik.uni-tuebingen.de
WWW:
http://www.berndammann.de/publications
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# A minimal classical sequent calculus free of structural rules
## 1 Introduction
The following Gentzen-Schütte-Tait \[Gen39, Sch50, Tai68\] system, denoted GS1p in \[TS96\], is a standard right-sided formulation of the propositional fragment of Gentzen’s classical sequent calculus LK:
System GS1p tensy $`P,\neg P`$ tensy $`\mathrm{\Gamma },A\mathrm{\Gamma },B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Gamma },A_i`$ $`_i`$ $`\mathrm{\Gamma },A_1A_2`$ tensy $`\mathrm{\Gamma }`$ $`𝖶`$ $`\mathrm{\Gamma },A`$ tensy $`\mathrm{\Gamma },A,A`$ $`𝖢`$ $`\mathrm{\Gamma },A`$
Here $`P`$ ranges over propositional variables, $`A,A_i,B`$ range over formulas, $`\mathrm{\Gamma }`$ ranges over disjoint unions of formulas, and comma denotes disjoint union.<sup>1</sup><sup>1</sup>1We label the conjunction and disjunction rules with $`\&`$ and $``$ for reasons which will become apparent later. By defining a sequent as a disjoint union of formulas, rather than an ordered list, we avoid an exchange/permutation rule (*cf.* \[TS96, §1.1\]). Negation is primitive on propositional variables $`P`$, and extends to compound formulas by de Morgan duality.<sup>2</sup><sup>2</sup>2$`\neg (AB)=(\neg A)(\neg B)`$ and $`\neg (AB)=(\neg A)(\neg B)`$.
The structural rules, weakening $`𝖶`$ and contraction $`𝖢`$, are absorbed in the following variant, a right-sided formulation of the propositional part of the calculus of \[Ket44\], called GS3p in \[TS96\].<sup>3</sup><sup>3</sup>3We label the disjunction rule as $`\&`$ to distinguish it from the disjunction rule $``$ of GS1p. The notation is derived from linear logic \[Gir87\].
System GS3p tensy $`\mathrm{\Gamma },P,\neg P`$ tensy $`\mathrm{\Gamma },A\mathrm{\Gamma },B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Gamma },A,B`$ $`\&`$ $`\mathrm{\Gamma },AB`$
The new axiom $`\mathrm{\Gamma },P,\neg P`$ amounts to the original axiom $`P,\neg P`$ followed immediately by weakenings. This paper presents a propositional classical sequent calculus Mp which is also free of structural rules:
System Mp tensy $`P,\neg P`$ tensy $`\mathrm{\Gamma },\mathrm{\Delta },A\mathrm{\Gamma },\mathrm{\Sigma },B`$ $``$ $`\mathrm{\Gamma },\mathrm{\Delta },\mathrm{\Sigma },AB`$ tensy $`\mathrm{\Gamma },A,B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Gamma },A_i`$ $`_i`$ $`\mathrm{\Gamma },A_1A_2`$
A distinguishing feature of Mp is the *blended conjunction rule*<sup>4</sup><sup>4</sup>4By analogy with GS3 and GS3p in \[TS96\], we reserve the symbol M for a full system with quantifiers, and use Mp to denote the propositional system. Following \[TS96\], we treat cut separately. To maximise emphasis on the blended conjunction rule, we omit quantifiers and cut in this paper.
tensy
Γ,Δ,AΓ,Σ,B
Γ,Δ,Σ,AB
𝑡𝑒𝑛𝑠𝑦
Γ,Δ,AΓ,Σ,B
Γ,Δ,Σ,AB
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,\Delta,A\hskip 17.22217pt\Gamma,\Sigma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=76.42003pt\hbox{\kern 3.00003pt$\;\wedge$}}}\hbox{\kern 5.06952pt\hbox{$\displaystyle\;\Gamma,\Delta,\Sigma,A\wedge B\;$}}}}
which combines the standard context-sharing and context-splitting conjunction rules:
tensy $`\mathrm{\Gamma },A\mathrm{\Gamma },B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Delta },A\mathrm{\Sigma },B`$ $``$ $`\mathrm{\Delta },\mathrm{\Sigma },AB`$
We refer to Mp as (cut-free propositional) *minimal sequent calculus*. In contrast to GS3p, contraction and weakening are absorbed symmetrically: both are absorbed into the conjunction rule, leaving the axiom rule intact.
Mp is evidently sound, since each of its rules can be derived (encoded) in GS1p. Theorem 1 (page 1) is completeness for formulas: a formula is valid iff it is derivable in Mp.<sup>5</sup><sup>5</sup>5Completeness here refers specifically to formulas, not to sequents. Section 6 discusses completeness for sequents.
### 1.1 Minimality
The blended conjunction rule $``$ is critical for the liberation from structural rules: Proposition 2 (page 2) shows that relaxing it to the union of the the two standard conjunction rules $`\&`$ and $``$ breaks completeness.<sup>6</sup><sup>6</sup>6In other words, if we remove the $``$ rule and add both the $`\&`$ and the $``$ rules, the resulting system fails to be complete. The formula $`\left((PQ)(\overline{Q}P)\right)\overline{P}`$ becomes underivable (see the proof of Proposition 2, page 2). The main theorem of the paper (page 2) formalises the sense in which Mp is a minimal complete core of classical sequent calculus:
Theorem 2: Minimality A standard sequent calculus $`S`$ is complete iff $`S\text{Mp}`$.
Here $`ST`$ (“$`S`$ contains $`T`$”) iff every rule of $`T`$ is derivable in $`S`$, and a *standard sequent calculus* is any propositional sequent calculus with the axiom $`\overline{P,\neg P}`$ and any subset of the following *standard rules*:
tensy $`\mathrm{\Gamma },A\mathrm{\Gamma },B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Gamma },A,B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Gamma }`$ $`𝖶`$ $`\mathrm{\Gamma },A`$
tensy $`\mathrm{\Delta },A\mathrm{\Sigma },B`$ $``$ $`\mathrm{\Delta },\mathrm{\Sigma },AB`$ tensy $`\mathrm{\Gamma },A_i`$ $`_i`$ $`\mathrm{\Gamma },A_1A_2`$ tensy $`\mathrm{\Gamma },A,A`$ $`𝖢`$ $`\mathrm{\Gamma },A`$
## 2 Notation and terminology
Formulas are built from literals (propositional variables $`P,Q,R\mathrm{}`$ and their formal complements $`\overline{P},\overline{Q},\overline{R},\mathrm{}`$) by the binary connectives and $``$ and or $``$. Define negation or not $`\neg `$ as an operation on formulas (rather than as a connective): $`\neg P=\overline{P}`$ and $`\neg \overline{P}=P`$ for all propositional variables $`P`$, with $`\neg (AB)=(\neg A)(\neg B)`$ and $`\neg (AB)=(\neg A)(\neg B)`$.
We identify a formula with its parse tree, a tree labelled with literals at the leaves and connectives at the internal vertices. A sequent is a non-empty disjoint union of formulas.<sup>7</sup><sup>7</sup>7Thus a sequent is a particular kind of labelled forest. This foundational treatment of formulas and sequents as labelled trees and forests sidesteps the common problem of “formulas” versus “formula occurrences”: disjoint unions of graphs are well understood in graph theory \[Bol02\]. Comma denotes disjoint union. Throughout the document, $`P,Q,\mathrm{}`$ range over propositional variables, $`A,B,\mathrm{}`$ over formulas, and $`\mathrm{\Gamma },\mathrm{\Delta },\mathrm{}`$ over (possibly empty) disjoint unions of formulas.
A formula $`A`$ is valid if it evaluates to $`1`$ under all possible $`0/1`$-assignments of its propositional variables (with the usual interpretation of $``$ and $``$ on $`\{0,1\}`$). A sequent $`A_1,\mathrm{},A_n`$ is valid iff the formula $`A_1(A_2(\mathrm{}(A_{n1}A_n)\mathrm{}))`$ is valid. A subsequent of a sequent $`\mathrm{\Gamma }`$ is any result of deleting zero or more formulas from $`\mathrm{\Gamma }`$; if at least one formula is deleted, the result is a proper subsequent.
## 3 Completeness
###### Theorem 1 (Completeness)
Every valid formula is derivable in Mp.
The proof is via the following auxiliary definitions and lemmas.
A sequent is minimally valid, or simply minimal, if it is valid while no proper subsequent is valid. For example, the sequents $`P,\neg P`$ and $`PQ,\overline{Q}P,\overline{P}`$ are minimal, while $`P,\neg P,Q`$ is not.
###### Lemma 1
Every valid sequent contains a minimal subsequent.
###### Proof.
Immediate from the definition of minimality. ∎
###### Lemma 2
Suppose a sequent $`\mathrm{\Gamma }`$ is a disjoint union of literals (*i.e.*, $`\mathrm{\Gamma }`$ contains no $``$ or $``$). Then $`\mathrm{\Gamma }`$ is minimal iff $`\mathrm{\Gamma }=P,\neg P`$ for some propositional variable $`P`$.
###### Proof.
By definition of validity in terms of valuations, $`\mathrm{\Gamma }`$ is valid iff it contains a complementary pair of literals, *i.e.*, iff $`\mathrm{\Gamma }=P,\neg P,\mathrm{\Delta }`$ with $`\mathrm{\Delta }`$ a disjoint union of zero or more literals. Since $`P,\neg P`$ is valid, $`\mathrm{\Gamma }`$ is minimal iff $`\mathrm{\Delta }`$ is empty. ∎
Suppose $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ are each disjoint unions of formulas (so each is either a sequent or empty). Write $`\mathrm{\Gamma }\mathrm{\Delta }`$ if $`\mathrm{\Gamma }`$ results from deleting zero or more formulas from $`\mathrm{\Delta }`$.
###### Lemma 3
Suppose $`\mathrm{\Gamma },A_1A_2`$ is minimal. Choose $`\mathrm{\Gamma }_1\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_2\mathrm{\Gamma }`$ such that $`\mathrm{\Gamma }_1,A_1`$ and $`\mathrm{\Gamma }_2,A_2`$ are minimal (existing by Lemma 1, since $`\mathrm{\Gamma },A_1`$ and $`\mathrm{\Gamma },A_2`$ are valid). Then every formula of $`\mathrm{\Gamma }`$ is in at least one of the $`\mathrm{\Gamma }_i`$.
###### Proof.
Suppose the formula $`B`$ of $`\mathrm{\Gamma }`$ is in neither $`\mathrm{\Gamma }_i`$. Let $`\mathrm{\Gamma }^{}`$ be the result of deleting $`B`$ from $`\mathrm{\Gamma }`$. Then $`\mathrm{\Gamma }^{},A_1A_2`$ is a valid proper subsequent of $`\mathrm{\Gamma },A_1A_2`$, contradicting minimality. (The sequent $`\mathrm{\Gamma }^{},A_1A_2`$ is valid since $`\mathrm{\Gamma }_1,A_1`$ and $`\mathrm{\Gamma }_2,A_2`$ are valid.) ∎
###### Lemma 4
Suppose $`\mathrm{\Gamma },AB`$ is minimal and $`\mathrm{\Gamma },A`$ is valid. Then $`\mathrm{\Gamma },A`$ is minimal.
###### Proof.
If not, some proper subsequent $`\mathrm{\Delta }`$ of $`\mathrm{\Gamma },A`$ is valid. If $`\mathrm{\Delta }`$ does not contain $`A`$, then it is also a proper subsequent of $`\mathrm{\Gamma },AB`$, contradicting minimality. Otherwise let $`\mathrm{\Delta }^{}`$ be the result of replacing $`A`$ in $`\mathrm{\Delta }`$ by $`AB`$. Since $`\mathrm{\Delta }`$ is valid, so also is $`\mathrm{\Delta }^{}`$. Thus $`\mathrm{\Delta }^{}`$ is a valid proper subsequent of $`\mathrm{\Gamma },AB`$, contradicting minimality. ∎
###### Lemma 5
Suppose $`\mathrm{\Gamma },AB`$ is minimal and neither $`\mathrm{\Gamma },A`$ nor $`\mathrm{\Gamma },B`$ is valid. Then $`\mathrm{\Gamma },A,B`$ is minimal.
###### Proof.
Suppose $`\mathrm{\Gamma },A,B`$ had a valid proper subsequent $`\mathrm{\Delta }`$. Since neither $`\mathrm{\Gamma },A`$ nor $`\mathrm{\Gamma },B`$ is valid, $`\mathrm{\Delta }`$ must contain both $`A`$ and $`B`$. Let $`\mathrm{\Delta }^{}`$ result from replacing $`A,B`$ by $`AB`$ in $`\mathrm{\Delta }`$. Then $`\mathrm{\Delta }^{}`$ is a valid proper subsequent of $`\mathrm{\Gamma },AB`$, contradicting minimality. ∎
Since a formula (viewed as a singleton sequent) is a minimal sequent, the Completeness Theorem (Theorem 1) is a special case of:
###### Proposition 1
Every minimal sequent is derivable in Mp.
###### Proof.
Suppose $`\mathrm{\Gamma }`$ is a minimal sequent. We proceed by induction on the number of connectives in $`\mathrm{\Gamma }`$.
* *Induction base (no connective).* Since $`\mathrm{\Gamma }`$ is minimal, Lemma 2 implies $`\mathrm{\Gamma }=P,\neg P`$, the conclusion of the axiom rule $`\overline{P,\neg P}`$.
* *Induction step (at least one connective).*
1. *Case: $`\mathrm{\Gamma }=\mathrm{\Delta },A_1A_2`$.* By Lemma 3, $`\mathrm{\Gamma }=\mathrm{\Sigma },\mathrm{\Delta }_1,\mathrm{\Delta }_2,A_1A_2`$ for $`\mathrm{\Sigma },\mathrm{\Delta }_1,A_1`$ and $`\mathrm{\Sigma },\mathrm{\Delta }_2,A_2`$ minimal. Write down the conjunction rule
tensy
Σ,Δ1,A1Σ,Δ2,A2
Σ,Δ1,Δ2,A1A2
𝑡𝑒𝑛𝑠𝑦
Σ,Δ1,A1Σ,Δ2,A2
Σ,Δ1,Δ2,A1A2
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Sigma,\Delta_{1},A_{1}\hskip 17.22217pt\Sigma,\Delta_{2},A_{2}$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=92.31104pt\hbox{\kern 3.00003pt$\;\wedge$}}}\hbox{\kern 7.22227pt\hbox{$\displaystyle\;\Sigma,\Delta_{1},\Delta_{2},A_{1}\wedge A_{2}\;$}}}}
and appeal to induction with the two hypothesis sequents.
2. *Case: $`\mathrm{\Gamma }=\mathrm{\Delta },A_1A_2`$.*
1. *Case: $`\mathrm{\Delta },A_i`$ is valid for some $`i\{1,2\}`$.* Write down the disjunction rule
tensy
Δ,Ai
i
Δ,A1A2
𝑡𝑒𝑛𝑠𝑦
Δ,Ai
i
Δ,A1A2
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 12.80751pt\hbox{$\displaystyle\penalty 1\Delta,A_{i}$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=47.82208pt\hbox{\kern 3.00003pt$\;\oplus_{i}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Delta,A_{1}\vee A_{2}\;$}}}}
then appeal to induction with $`\mathrm{\Delta },A_i`$, which is minimal by Lemma 4.
2. *Case: $`\mathrm{\Delta },A_i`$ is not valid for each $`i\{1,2\}`$.* Thus $`\mathrm{\Delta },A_1,A_2`$ is minimal, by Lemma 5. Write down the disjunction rule
tensy
Δ,A1,A2
&
Δ,A1A2
𝑡𝑒𝑛𝑠𝑦
Δ,A1,A2
&
Δ,A1A2
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 4.99994pt\hbox{$\displaystyle\penalty 1\Delta,A_{1},A_{2}$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=47.8221pt\hbox{\kern 3.00003pt$\raisebox{0.0pt}{\raisebox{6.94444pt}{\rotatebox{180.0}{$\&$}}}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Delta,A_{1}\vee A_{2}\;$}}}}
then appeal to induction with $`\mathrm{\Delta },A_1,A_2`$.
($`\mathrm{\Gamma }`$ may match both 1 and 2 in the inductive step, permitting some choice in the construction of the derivation. There is choice in case 2(a) if both $`\mathrm{\Delta },A_1`$ and $`\mathrm{\Delta },A_2`$ are valid.) ∎
Note that completeness does not hold for arbitrary valid sequents. For example, the sequent $`P,\neg P,Q`$ is valid but not derivable in Mp. A sequent is valid iff some some subsequent is derivable in Mp. Thus Mp is complete for sequents modulo final weakenings. In this sense, Mp is akin to system GS5p of \[TS96, §7.4\] (related to resolution). (See also Section 6.)
## 4 The Minimality Theorem
Relaxing blended conjunction to the pair of standard conjunction rules (context-sharing $`\&`$ and context-splitting $``$) breaks completeness. Let Mp<sup>-</sup> be the following subsystem of Mp:<sup>8</sup><sup>8</sup>8This precursor of Mp is (cut-free) multiplicative-additive linear logic \[Gir87\] with *tensor* $``$ and *with* $`\&`$ collapsed to $``$, and *plus* $``$ and *par* $`\&`$ collapsed to $``$.
System Mp<sup>-</sup> tensy $`P,\neg P`$ tensy $`\mathrm{\Gamma },A\mathrm{\Gamma },B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Gamma },A,B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Delta },A\mathrm{\Sigma },B`$ $``$ $`\mathrm{\Delta },\mathrm{\Sigma },AB`$ tensy $`\mathrm{\Gamma },A_i`$ $`_i`$ $`\mathrm{\Gamma },A_1A_2`$
###### Proposition 2
System Mp<sup>-</sup> is incomplete.
###### Proof.
We show that the valid formula $`A=\left((PQ)(\overline{Q}P)\right)\overline{P}`$ is not derivable in Mp<sup>-</sup>. The placement of the two outermost $``$ connectives forces the last two rules of a potential derivation to be disjunction rules. Since $`PQ,\overline{Q}P,\overline{P}`$ is minimal (no proper subsequent is valid), the two disjunction rules must be $`\&`$ rather than $``$:
tensy
tensy
PQ,¯QP,¯P
&
(PQ)(¯QP),¯P
&
((PQ)(¯QP))¯P
𝑡𝑒𝑛𝑠𝑦
tensy
PQ,¯QP,¯P
&
(PQ)(¯QP),¯P
&
((PQ)(¯QP))¯P
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 10.83334pt\hbox{$\displaystyle\penalty 1P\wedge Q,\,\overline{Q}\wedge P,\,\overline{P}$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=86.30118pt\hbox{\kern 3.00003pt$\raisebox{0.0pt}{\raisebox{6.94444pt}{\rotatebox{180.0}{$\&$}}}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle(P\wedge Q)\vee(\overline{Q}\wedge P),\,\overline{P}$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=128.2978pt\hbox{\kern 3.00003pt$\raisebox{0.0pt}{\raisebox{6.94444pt}{\rotatebox{180.0}{$\&$}}}$}}}\hbox{\kern 13.2761pt\hbox{$\displaystyle\;\big{(}(P\wedge Q)\vee(\overline{Q}\wedge P)\big{)}\vee\overline{P}\;$}}}}
It remains to show that $`PQ,\overline{Q}P,\overline{P}`$ is not derivable in Mp<sup>-</sup>. There are only two connectives, both $``$, so the last rule must be a conjunction.
1. *Case: the last rule is a context-sharing $`\&`$-rule.*
1. *Case: The last rule introduces $`PQ`$.*
tensy
P,¯QP,¯PQ,¯QP,¯P
&
PQ,¯QP,¯P
𝑡𝑒𝑛𝑠𝑦
P,¯QP,¯PQ,¯QP,¯P
&
PQ,¯QP,¯P
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1P,\,\overline{Q}\wedge P,\overline{P}\hskip 17.22217ptQ,\,\overline{Q}\wedge P,\,\overline{P}$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=104.11015pt\hbox{\kern 3.00003pt$\;\&$}}}\hbox{\kern 16.96011pt\hbox{$\displaystyle\;P\wedge Q,\,\overline{Q}\wedge P,\,\overline{P}\;$}}}}
The left hypothesis $`P,\overline{Q}P,\overline{P}`$ cannot be derived in Mp<sup>-</sup>, since there is no $`Q`$ to match the $`\overline{Q}`$ (and no weakening).
2. *Case: The last rule introduces $`\overline{Q}P`$.* The same as the previous case, by symmetry, and exchanging $`Q\overline{Q}`$.
2. *Case: the last rule is a context-splitting $``$-rule.*
1. *Case: The last rule introduces $`PQ`$.*
tensy
P,ΓQ,Δ
PQ,¯QP,¯P
𝑡𝑒𝑛𝑠𝑦
P,ΓQ,Δ
PQ,¯QP,¯P
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 5.77942pt\hbox{$\displaystyle\penalty 1P,\,\Gamma\hskip 17.22217ptQ,\,\Delta$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=70.18993pt\hbox{\kern 3.00003pt$\;\otimes$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;P\wedge Q,\,\overline{Q}\wedge P,\,\overline{P}\;$}}}}
We must allocate each of $`\overline{Q}P`$ and $`\overline{P}`$ either to $`\mathrm{\Gamma }`$ or to $`\mathrm{\Delta }`$. If $`\overline{Q}P`$ is in $`\mathrm{\Gamma }`$, then $`P,\mathrm{\Gamma }`$ is not derivable in Mp<sup>-</sup>, since it contains no $`Q`$ to match the $`\overline{Q}`$. So $`\overline{Q}P`$ is in $`\mathrm{\Delta }`$. But then the $`\overline{P}`$ is required in both $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$.
2. *Case: The last rule introduces $`\overline{Q}P`$.* The same as the previous case, by symmetry, and exchanging $`Q\overline{Q}`$.
A standard system is any propositional sequent calculus containing the axiom $`\overline{P,\neg P}`$ and any of the following standard rules:
tensy $`\mathrm{\Gamma },A\mathrm{\Gamma },B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Gamma },A,B`$ $`\&`$ $`\mathrm{\Gamma },AB`$ tensy $`\mathrm{\Gamma }`$ $`𝖶`$ $`\mathrm{\Gamma },A`$
tensy $`\mathrm{\Delta },A\mathrm{\Sigma },B`$ $``$ $`\mathrm{\Delta },\mathrm{\Sigma },AB`$ tensy $`\mathrm{\Gamma },A_i`$ $`_i`$ $`\mathrm{\Gamma },A_1A_2`$ tensy $`\mathrm{\Gamma },A,A`$ $`𝖢`$ $`\mathrm{\Gamma },A`$
Thus there are $`2^6=64`$ such systems (many of which will not be complete).
System $`S`$ contains system $`T`$, denoted $`ST`$, if each rule of $`T`$ is a derived rule of $`S`$. For example, system GS1p (page 1) contains Mp since the blended conjunction rule $``$ and the disjunction rule $`\&`$ of Mp can be derived in GS1p:
tensy
Γ,Δ,AΓ,Σ,B
Γ,Δ,Σ,AB
tensy
tensy
Γ,Δ,A
W
Γ,Δ,Σ,A
tensy
Γ,Σ,B
W
Γ,Δ,Σ,B
&
Γ,Δ,Σ,AB
tensy
Γ,A,B
&
Γ,AB
tensy
tensy
tensy
Γ,A,B
2
Γ,A,AB
1
Γ,AB,AB
C
Γ,AB
𝑡𝑒𝑛𝑠𝑦
Γ,Δ,AΓ,Σ,B
Γ,Δ,Σ,AB
𝑡𝑒𝑛𝑠𝑦
tensy
Γ,Δ,A
W
Γ,Δ,Σ,A
tensy
Γ,Σ,B
W
Γ,Δ,Σ,B
&
Γ,Δ,Σ,AB
missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpression𝑡𝑒𝑛𝑠𝑦
Γ,A,B
&
Γ,AB
𝑡𝑒𝑛𝑠𝑦
tensy
tensy
Γ,A,B
2
Γ,A,AB
1
Γ,AB,AB
C
Γ,AB
\begin{array}[]{ccc}{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,\Delta,A\hskip 17.22217pt\Gamma,\Sigma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=76.42003pt\hbox{\kern 3.00003pt$\;\wedge$}}}\hbox{\kern 5.06952pt\hbox{$\displaystyle\;\Gamma,\Delta,\Sigma,A\wedge B\;$}}}}&\hskip 21.52771pt\longleftarrow&{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 5.83333pt\hbox{$\displaystyle\penalty 1\Gamma,\Delta,A$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=41.52776pt\hbox{\kern 3.00003pt$\;\mathsf{W}^{\ast}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\Gamma,\Delta,\Sigma,A$}}}}\hskip 5.0pt plus 1.0fil\penalty 2\hskip 17.22217pt\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 6.38889pt\hbox{$\displaystyle\penalty 1\Gamma,\Sigma,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=42.11453pt\hbox{\kern 3.00003pt$\;\mathsf{W}^{\ast}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\Gamma,\Delta,\Sigma,B$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=201.01314pt\hbox{\kern 3.00003pt$\;\&$}}}\hbox{\kern 67.36607pt\hbox{$\displaystyle\;\Gamma,\Delta,\Sigma,A\wedge B\;$}}}}\\
\\
\\
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 6.11102pt\hbox{$\displaystyle\penalty 1\Gamma,A,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=41.8366pt\hbox{\kern 3.00003pt$\raisebox{0.0pt}{\raisebox{6.94444pt}{\rotatebox{180.0}{$\&$}}}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A\vee B\;$}}}}&\hskip 21.52771pt\longleftarrow&{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 10.69438pt\hbox{$\displaystyle\penalty 1\Gamma,A,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=51.00331pt\hbox{\kern 3.00003pt$\;\oplus_{2}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A,A\vee B$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=100.79982pt\hbox{\kern 3.00003pt$\;\oplus_{1}$}}}\hbox{\kern 16.6882pt\hbox{$\displaystyle\Gamma,A\vee B,A\vee B$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=150.59633pt\hbox{\kern 3.00003pt$\;\mathsf{C}$}}}\hbox{\kern 54.37987pt\hbox{$\displaystyle\;\Gamma,A\vee B\;$}}}}\end{array}
where $`𝖶^{}`$ denotes a sequence of zero or more weakenings.
###### Theorem 2 (Minimality Theorem)
A standard system is complete iff it contains Mp.
### 4.1 Proof of the Minimality Theorem
Two systems are equivalent if each contains the other. For example, it is well known that GS1p (page 1) is equivalent to:<sup>9</sup><sup>9</sup>9This system is multiplicative linear logic \[Gir87\] plus contraction and weakening (with the connectives denoted $``$ and $``$ instead of $``$ and $`\&`$ ).
tensy
P,¬P
tensy
Δ,AΣ,B
Δ,Σ,AB
tensy
Γ,A1,A2
&
Γ,A1A2
tensy
Γ
W
Γ,A
tensy
Γ,A,A
C
Γ,A
𝑡𝑒𝑛𝑠𝑦
P,¬P
𝑡𝑒𝑛𝑠𝑦
Δ,AΣ,B
Δ,Σ,AB
𝑡𝑒𝑛𝑠𝑦
Γ,A1,A2
&
Γ,A1A2
𝑡𝑒𝑛𝑠𝑦
Γ
W
Γ,A
𝑡𝑒𝑛𝑠𝑦
Γ,A,A
C
Γ,A
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 16.42004pt\hbox{$\displaystyle\penalty 1\rule{0.0pt}{5.59721pt}$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=32.84009pt\hbox{\kern 3.00003pt$$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;P,\neg P\;$}}}}\hskip 21.52771pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 3.54153pt\hbox{$\displaystyle\penalty 1\Delta,A\hskip 12.91663pt\Sigma,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=60.03093pt\hbox{\kern 3.00003pt$\;\otimes$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Delta,\,\Sigma,\,A\wedge B\;$}}}}\hskip 21.52771pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 5.83325pt\hbox{$\displaystyle\penalty 1\Gamma,A_{1},A_{2}$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=46.29427pt\hbox{\kern 3.00003pt$\raisebox{0.0pt}{\raisebox{6.94444pt}{\rotatebox{180.0}{$\&$}}}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,\,A_{1}\vee A_{2}\;$}}}}\hskip 21.52771pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 6.80551pt\hbox{$\displaystyle\penalty 1\Gamma$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=19.86104pt\hbox{\kern 3.00003pt$\;\mathsf{W}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A$}}}}\hskip 21.52771pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,A,A$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=29.02776pt\hbox{\kern 3.00003pt$\;\mathsf{C}$}}}\hbox{\kern 3.1945pt\hbox{$\displaystyle\;\Gamma,A\;$}}}}
via the following rule derivations:
tensy
Δ,AΣ,B
Δ,Σ,AB
tensy
tensy
Δ,A
W
Δ,Σ,A
tensy
Σ,B
W
Δ,Σ,B
&
Δ,Σ,AB
tensy
Γ,Ai
i
Γ,A1A2
tensy
tensy
Γ,Ai
W
Γ,A1,A2
&
Γ,A1A2
tensy
Γ,A,B
&
Γ,AB
tensy
tensy
tensy
Γ,A,B
2
Γ,A,AB
1
Γ,AB,AB
C
Γ,AB
tensy
Γ,AΓ,B
&
Γ,AB
tensy
tensy
Γ,AΓ,B
Γ,Γ,AB
C
Γ,AB
𝑡𝑒𝑛𝑠𝑦
Δ,AΣ,B
Δ,Σ,AB
𝑡𝑒𝑛𝑠𝑦
tensy
Δ,A
W
Δ,Σ,A
tensy
Σ,B
W
Δ,Σ,B
&
Δ,Σ,AB
missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpression𝑡𝑒𝑛𝑠𝑦
Γ,Ai
i
Γ,A1A2
𝑡𝑒𝑛𝑠𝑦
tensy
Γ,Ai
W
Γ,A1,A2
&
Γ,A1A2
missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpression𝑡𝑒𝑛𝑠𝑦
Γ,A,B
&
Γ,AB
𝑡𝑒𝑛𝑠𝑦
tensy
tensy
Γ,A,B
2
Γ,A,AB
1
Γ,AB,AB
C
Γ,AB
missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpression𝑡𝑒𝑛𝑠𝑦
Γ,AΓ,B
&
Γ,AB
𝑡𝑒𝑛𝑠𝑦
tensy
Γ,AΓ,B
Γ,Γ,AB
C
Γ,AB
\begin{array}[]{ccc}{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Delta,A\hskip 17.22217pt\Sigma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=57.2534pt\hbox{\kern 3.00003pt$\;\otimes$}}}\hbox{\kern 0.27786pt\hbox{$\displaystyle\;\Delta,\Sigma,A\wedge B\;$}}}}&\hskip 21.52771pt\longleftarrow&{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 5.83333pt\hbox{$\displaystyle\penalty 1\Delta,A$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=31.94444pt\hbox{\kern 3.00003pt$\;\mathsf{W}^{\ast}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\Delta,\Sigma,A$}}}}\hskip 5.0pt plus 1.0fil\penalty 2\hskip 17.22217pt\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 6.38889pt\hbox{$\displaystyle\penalty 1\Sigma,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=32.53122pt\hbox{\kern 3.00003pt$\;\mathsf{W}^{\ast}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\Delta,\Sigma,B$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=181.84651pt\hbox{\kern 3.00003pt$\;\&$}}}\hbox{\kern 62.57442pt\hbox{$\displaystyle\;\Delta,\Sigma,A\wedge B\;$}}}}\\
\\
\\
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 12.80751pt\hbox{$\displaystyle\penalty 1\Gamma,A_{i}$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=44.62762pt\hbox{\kern 3.00003pt$\;\oplus_{i}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A_{1}\vee A_{2}\;$}}}}&\hskip 21.52771pt\longleftarrow&{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 9.19643pt\hbox{$\displaystyle\penalty 1\Gamma,A_{i}$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=37.40546pt\hbox{\kern 3.00003pt$\;\mathsf{W}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A_{1},A_{2}$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=84.6798pt\hbox{\kern 3.00003pt$\raisebox{0.0pt}{\raisebox{6.94444pt}{\rotatebox{180.0}{$\&$}}}$}}}\hbox{\kern 20.02608pt\hbox{$\displaystyle\;\Gamma,A_{1}\vee A_{2}\;$}}}}\\
\\
\\
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 6.11102pt\hbox{$\displaystyle\penalty 1\Gamma,A,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=41.8366pt\hbox{\kern 3.00003pt$\raisebox{0.0pt}{\raisebox{6.94444pt}{\rotatebox{180.0}{$\&$}}}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A\vee B\;$}}}}&\hskip 21.52771pt\longleftarrow&{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 10.69438pt\hbox{$\displaystyle\penalty 1\Gamma,A,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=51.00331pt\hbox{\kern 3.00003pt$\;\oplus_{2}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A,A\vee B$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=100.79982pt\hbox{\kern 3.00003pt$\;\oplus_{1}$}}}\hbox{\kern 16.6882pt\hbox{$\displaystyle\Gamma,A\vee B,A\vee B$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=150.59633pt\hbox{\kern 3.00003pt$\;\mathsf{C}$}}}\hbox{\kern 54.37987pt\hbox{$\displaystyle\;\Gamma,A\vee B\;$}}}}\\
\\
\\
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,A\hskip 17.22217pt\Gamma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=51.97562pt\hbox{\kern 3.00003pt$\;\&$}}}\hbox{\kern 5.06952pt\hbox{$\displaystyle\;\Gamma,A\wedge B\;$}}}}&\hskip 21.52771pt\longleftarrow&{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,A\hskip 17.22217pt\Gamma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=51.97562pt\hbox{\kern 3.00003pt$\;\otimes$}}}\hbox{\kern 0.27786pt\hbox{$\displaystyle\;\Gamma,\Gamma,A\wedge B\;$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=98.97212pt\hbox{\kern 3.00003pt$\;\mathsf{C}^{\ast}$}}}\hbox{\kern 28.56776pt\hbox{$\displaystyle\;\Gamma,A\wedge B\;$}}}}\end{array}
We shall abbreviate these four rule derivations as follows, and write analogous abbreviations for other rule derivations.<sup>10</sup><sup>10</sup>10“Context-splitting conjunction $``$ is derivable from context-sharing conjunction $`\&`$ and weakening $`𝖶`$”, *etc.*
$$\begin{array}{cccccc}& & \&𝖶& \&& & 𝖢\\ & & \&𝖶& \&& & 𝖢\end{array}$$
#### 4.1.1 The three complete standard systems
As a stepping stone towards the Minimality Theorem, we shall prove that, up to equivalence, there are only three complete standard systems.
We abbreviate a system by listing its non-axiom rules. For example, $`\text{GS1p}=(\&,,𝖶,𝖢)`$ and $`\text{Mp}=(,,\&)`$. Besides GS1p, we shall pay particular attention to the systems
$$\begin{array}{ccc}\hfill \text{Pp}=& (,,𝖢)\hfill & \text{Positive calculus}\hfill \\ \hfill \text{Np}=& (\&,\&,𝖶)\hfill & \text{Negative calculus}\hfill \end{array}$$
(Our terminology comes from polarity of connectives in linear logic \[Gir87\]: *tensor* $``$ and *plus* $``$ are positive, and *with* $`\&`$ and *par* $`\&`$ are negative.)
###### Proposition 3
Up to equivalence:
* $`\text{GS1p}=(\&,,𝖢,𝖶)`$ is the only complete standard system with both contraction $`𝖢`$ and weakening $`𝖶`$;
* $`\text{Pp}=(,,𝖢)`$ is the only complete standard system without weakening $`𝖶`$;
* $`\text{Np}=(\&,\&,𝖶)`$ is the only complete standard system without contraction $`𝖢`$.
The proof is via the following lemmas.
###### Lemma 6
$`\text{Mp}=(,,\&)`$ is contained in each of $`\text{Pp}=(,,𝖢)`$, $`\text{Np}=(\&,\&,𝖶)`$ and $`\text{GS1p}=(\&,,𝖢,𝖶)`$.
###### Proof.
Pp contains Mp since $`𝖢`$ ,
tensy
Γ,Δ,AΓ,Σ,B
Γ,Δ,Σ,AB
tensy
tensy
Γ,Δ,AΓ,Σ,B
Γ,Γ,Δ,Σ,AB
C
Γ,Δ,Σ,AB
𝑡𝑒𝑛𝑠𝑦
Γ,Δ,AΓ,Σ,B
Γ,Δ,Σ,AB
𝑡𝑒𝑛𝑠𝑦
tensy
Γ,Δ,AΓ,Σ,B
Γ,Γ,Δ,Σ,AB
C
Γ,Δ,Σ,AB
\begin{array}[]{ccc}{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,\Delta,A\hskip 17.22217pt\Gamma,\Sigma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=76.42003pt\hbox{\kern 3.00003pt$\;\wedge$}}}\hbox{\kern 5.06952pt\hbox{$\displaystyle\;\Gamma,\Delta,\Sigma,A\wedge B\;$}}}}&\hskip 21.52771pt\longleftarrow&{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,\Delta,A\hskip 17.22217pt\Gamma,\Sigma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=77.53114pt\hbox{\kern 3.00003pt$\;\otimes$}}}\hbox{\kern 0.83342pt\hbox{$\displaystyle\;\Gamma,\Gamma,\Delta,\Sigma,A\wedge B\;$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=124.52765pt\hbox{\kern 3.00003pt$\;\mathsf{C}^{\ast}$}}}\hbox{\kern 29.12334pt\hbox{$\displaystyle\;\Gamma,\Delta,\Sigma,A\wedge B\;$}}}}\end{array}
(where $`𝖢^{}`$ denotes zero or more consecutive contractions) and $`\&𝖢`$ :
tensy
Γ,A,B
&
Γ,AB
tensy
tensy
tensy
Γ,A,B
2
Γ,A,AB
1
Γ,AB,AB
C
Γ,AB
𝑡𝑒𝑛𝑠𝑦
Γ,A,B
&
Γ,AB
𝑡𝑒𝑛𝑠𝑦
tensy
tensy
Γ,A,B
2
Γ,A,AB
1
Γ,AB,AB
C
Γ,AB
\begin{array}[]{ccc}{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 6.11102pt\hbox{$\displaystyle\penalty 1\Gamma,A,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=41.8366pt\hbox{\kern 3.00003pt$\raisebox{0.0pt}{\raisebox{6.94444pt}{\rotatebox{180.0}{$\&$}}}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A\vee B\;$}}}}&\hskip 21.52771pt\longleftarrow&{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 10.69438pt\hbox{$\displaystyle\penalty 1\Gamma,A,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=51.00331pt\hbox{\kern 3.00003pt$\;\oplus_{2}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A,A\vee B$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=100.79982pt\hbox{\kern 3.00003pt$\;\oplus_{1}$}}}\hbox{\kern 16.6882pt\hbox{$\displaystyle\Gamma,A\vee B,A\vee B$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=150.59633pt\hbox{\kern 3.00003pt$\;\mathsf{C}$}}}\hbox{\kern 54.37987pt\hbox{$\displaystyle\;\Gamma,A\vee B\;$}}}}\end{array}
Np contains Mp since $`𝖶\&`$ , and $`𝖶\&`$ (see page 4.1). $`\text{GS1p}=(\&,,𝖢,𝖶)`$ is equivalent to $`(,\&,,\&,𝖢,𝖶)`$ since $``$ and $`\&`$ are derivable. Thus GS1p contains Pp (and Np), hence Mp. ∎
###### Lemma 7
$`\text{Pp}=(𝖢,,)`$ and $`\text{Np}=(\&,\&,𝖶)`$ are complete. <sup>11</sup><sup>11</sup>11Recall that completeness refers to formulas, not sequents in general.
###### Proof.
Each contains Mp by Lemma 6, which is complete (Theorem 1). ∎
###### Lemma 8
Up to equivalence, system $`\text{GS1p}=(\&,,𝖢,𝖶)`$ is the only complete standard system with both contraction $`𝖢`$ and weakening $`𝖶`$.
###### Proof.
GS1p is complete (see *e.g.* \[TS96\], or by the fact that GS1p contains Mp which is complete). Any complete system must have a conjunction rule ($``$ or $`\&`$) and a disjunction rule ($``$ or $`\&`$ ). In the presence of $`𝖢`$ and $`𝖶`$, the two conjunctions are derivable from one other, as are the two disjunctions (see page 4.1). ∎
###### Lemma 9
A complete standard system without weakening $`𝖶`$ must contain $`\text{Pp}=(,,𝖢)`$.
###### Proof.
System $`\text{Mp}\text{-}=(,,\&,\&)`$, with both conjunction rules and both disjunction rules, is incomplete (Proposition 2, page 2), therefore we must have contraction $`𝖢`$.
Without the $``$ rule, the valid formula $`(P\overline{P})Q`$ is not derivable: the last rule must be $`\&`$ , leaving us to derive $`P\overline{P},Q`$, which is impossible without weakening $`𝖶`$ (*i.e.*, with at most $`\&`$ , $`\&`$, $``$ and $`𝖢`$ available), since, after a necessary axiom $`\overline{P,\overline{P}}`$ at the top of the derivation, there is no way to introduce the formula $`Q`$.
Without the context-splitting $``$ rule, the valid formula $`P(Q(\overline{P}\overline{Q}))`$ is not derivable. The last two rules must be $`\&`$ , for if we use a $``$ we will not be able to match complementary literals in the axioms at the top of the derivation. Thus we are left to derive $`P,Q,\overline{P}\overline{Q}`$, using $`\&`$ and $`𝖢`$. The derivation must contain an axiom rule $`\overline{P,\overline{P}}`$. The next rule can only be a $`\&`$ (since $`P,\overline{P}`$ cannot be the hypothesis sequent of a contraction $`𝖢`$ rule). Since the only $``$-formula in the final concluding sequent $`P,Q,\overline{P}\overline{Q}`$ is $`\overline{P}\overline{Q}`$, and the $`\&`$ rule is context sharing, the $`\&`$-rule must be
tensy
tensy
P,¯P
P,¯Q
&
P,¯P¯Q
𝑡𝑒𝑛𝑠𝑦
tensy
P,¯P
P,¯Q
&
P,¯P¯Q
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\hskip 5.0pt plus 1.0fil{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 10.01556pt\hbox{$\displaystyle\penalty 1$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=20.03113pt\hbox{}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;P,\overline{P}\;$}}}}\hskip 5.0pt plus 1.0fil\penalty 2\hskip 17.22217pt\hskip 5.0pt plus 1.0fil{\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 8.62671pt\hbox{$\displaystyle\penalty 1$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 7.23782pt$\vbox to20.6665pt{\leaders\vbox to5.16663pt{\vss\hbox{$\cdot$}\vss}\vfill}$\hbox{}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;P,\overline{Q}$}}}}\hskip 5.0pt plus 1.0fil\penalty 2$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=98.72557pt\hbox{\kern 3.00003pt$\;\&$}}}\hbox{\kern 33.51385pt\hbox{$\displaystyle\;P,\overline{P}\wedge\overline{Q}\;$}}}}
but $`P,\overline{Q}`$ is not derivable. ∎
###### Lemma 10
Up to equivalence, $`\text{Pp}=(,,𝖢)`$ is the only complete standard system without weakening $`𝖶`$.
###### Proof.
By Lemma 7, Pp is complete. By Lemma 9, every $`𝖶`$-free complete standard system contains Pp. All other $`𝖶`$-free standard systems containing Pp are equivalent to Pp, since the standard rule derivations $`\&𝖢`$ and $`\&𝖢`$ yield $`\&`$ and $`\&`$ (see page 4.1). ∎
###### Lemma 11
A complete standard system without contraction $`𝖢`$ must contain $`\text{Np}=(\&,\&,𝖶)`$.
###### Proof.
System $`\text{Mp}\text{-}=(,,\&,\&)`$, with both conjunction rules and both disjunction rules, is incomplete (Proposition 2, page 2), therefore we must have weakening $`𝖶`$.
Without the $`\&`$ rule, the valid formula $`P\overline{P}`$ would not be derivable.
Without the $`\&`$ rule the valid formula $`P(\overline{P}\overline{P})`$ would not be derivable. The last rule must be a $`\&`$ (rather than a $``$, otherwise we lack either $`P`$ or $`\overline{P}`$), so we are left to derive $`P,\overline{P}\overline{P}`$. The last rule cannot be $`\&`$ or $``$, as the only connective is $``$. It cannot be $`𝖶`$, or else we lack either $`P`$ or $`\overline{P}`$. It cannot be $``$, as one of the two hypotheses will be the single formula $`\overline{P}`$. ∎
###### Lemma 12
Up to equivalence, $`\text{Np}=(\&,\&,𝖶)`$ is the only complete standard system without contraction $`𝖢`$.
###### Proof.
By Lemma 7, Np is complete. By Lemma 11, every $`𝖢`$-free complete standard system contains Np. All other $`𝖢`$-free standard systems containing Np are equivalent to Np, since the standard rule derivations $`\&𝖶`$ and $`\&𝖶`$ yield $``$ and $``$ (see page 4.1). ∎
* Proof of Proposition 3. Parts (1), (2) and (3) are Lemmas 8, 10 and 12, respectively. $`\mathrm{}`$
###### Lemma 13
Every standard complete system has contraction $`𝖢`$ or weakening $`𝖶`$.
###### Proof.
Otherwise it is contained in $`\text{Mp}\text{-}=(,,\&,\&)`$, which is incomplete (Prop. 2). ∎
###### Theorem 3
Up to equivalence, there are only three complete standard systems:
1. The Gentzen-Schütte-Tait system $`\text{GS1p}=(\&,,𝖢,𝖶)`$.
2. Positive calculus $`\text{Pp}=(,,𝖢)`$.
3. Negative calculus $`\text{Np}=(\&,\&,𝖶)`$.
###### Proof.
Proposition 3 and Lemma 13. ∎
* Proof of Minimality Theorem (Theorem 2). Each of the three complete standard systems contains Mp (Lemma 6). $`\mathrm{}`$
The three inequivalent complete standard systems GS1p, Pp and Np, together with propositional minimal sequent calculus Mp, sit in the following Hasse diagram of containments:
Containments of complete inequivalent systems
Thus we can view propositional minimal sequent calculus Mp as a minimal complete core of GS1p, hence of (propositional) Gentzen’s LK.
## 5 Extended Minimality Theorem
Define an extended system as one containing the axiom rule $`\overline{P,\overline{P}}`$ and any of the following rules. (We have extended the definition of *standard system* by making blended conjunction available.)
Extended system rules tensy
Δ,AΣ,B
Δ,Σ,AB
tensy
Γ,Δ,AΓ,Σ,B
Γ,Δ,Σ,AB
tensy
Γ,AΓ,B
&
Γ,AB
tensy
Γ,Ai
i
Γ,A1A2
tensy
Γ,A,B
&
Γ,AB
tensy
Γ,A,A
C
Γ,A
tensy
Γ
W
Γ,A
missing-subexpression𝑡𝑒𝑛𝑠𝑦
Δ,AΣ,B
Δ,Σ,AB
𝑡𝑒𝑛𝑠𝑦
Γ,Δ,AΓ,Σ,B
Γ,Δ,Σ,AB
𝑡𝑒𝑛𝑠𝑦
Γ,AΓ,B
&
Γ,AB
missing-subexpressionmissing-subexpression𝑡𝑒𝑛𝑠𝑦
Γ,Ai
i
Γ,A1A2
𝑡𝑒𝑛𝑠𝑦
Γ,A,B
&
Γ,AB
missing-subexpressionmissing-subexpression𝑡𝑒𝑛𝑠𝑦
Γ,A,A
C
Γ,A
𝑡𝑒𝑛𝑠𝑦
Γ
W
Γ,A
\begin{array}[]{c}\\[-4.30554pt]
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Delta,A\hskip 17.22217pt\Sigma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=57.2534pt\hbox{\kern 3.00003pt$\;\otimes$}}}\hbox{\kern 0.27786pt\hbox{$\displaystyle\;\Delta,\Sigma,A\wedge B\;$}}}}\hskip 25.83325pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,\Delta,A\hskip 17.22217pt\Gamma,\Sigma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=76.42003pt\hbox{\kern 3.00003pt$\;\wedge$}}}\hbox{\kern 5.06952pt\hbox{$\displaystyle\;\Gamma,\Delta,\Sigma,A\wedge B\;$}}}}\hskip 25.83325pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,A\hskip 17.22217pt\Gamma,B$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=51.97562pt\hbox{\kern 3.00003pt$\;\&$}}}\hbox{\kern 5.06952pt\hbox{$\displaystyle\;\Gamma,A\wedge B\;$}}}}\\
\\
\\[-4.30554pt]
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 12.80751pt\hbox{$\displaystyle\penalty 1\Gamma,A_{i}$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=44.62762pt\hbox{\kern 3.00003pt$\;\oplus_{i}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A_{1}\vee A_{2}\;$}}}}\hskip 51.6665pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 6.11102pt\hbox{$\displaystyle\penalty 1\Gamma,A,B$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=41.8366pt\hbox{\kern 3.00003pt$\raisebox{0.0pt}{\raisebox{6.94444pt}{\rotatebox{180.0}{$\&$}}}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A\vee B\;$}}}}\\
\\
\\[-4.30554pt]
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,A,A$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=29.02776pt\hbox{\kern 3.00003pt$\;\mathsf{C}$}}}\hbox{\kern 3.1945pt\hbox{$\displaystyle\;\Gamma,A\;$}}}}\hskip 60.27759pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{\kern 6.80551pt\hbox{$\displaystyle\penalty 1\Gamma$}}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=19.86104pt\hbox{\kern 3.00003pt$\;\mathsf{W}$}}}\hbox{\kern 0.0pt\hbox{$\displaystyle\;\Gamma,A$}}}}\\
\\
\end{array}
The Minimality Theorem (Theorem 2, page 2) extends as follows.
###### Theorem 4 (Extended Minimality Theorem)
An extended system is complete iff it contains propositional minimal sequent calculus Mp.
To prove this theorem, we require two lemmas.
###### Lemma 14
Suppose $`S`$ is a complete extended system with the blended conjunction rule $``$, and with at least one of contraction $`𝖢`$ or weakening $`𝖶`$. Then $`S`$ is equivalent to a standard system.
###### Proof.
If $`S`$ has weakening $`𝖶`$, let $`S^{}`$ be the result of replacing the blended conjunction rule $``$ in $`S`$ by context-sharing conjunction $`\&`$; otherwise $`S`$ has contraction, and let $`S^{}`$ result from replacing $``$ by context-splitting $``$. Then $`S^{}`$ is equivalent to $`S`$, since $`𝖢`$ (page 4.1.1) and $`\&𝖶`$ (page 4). ∎
###### Lemma 15
Suppose $`S`$ is a complete extended system with neither contraction $`𝖢`$ nor weakening $`𝖶`$. Then $`S`$ is equivalent to propositional minimal sequent calculus Mp.
###### Proof.
Since $`\text{Mp}\text{-}=(,\&,,\&)`$ is incomplete (Proposition 2, page 2), $`S`$ must have the blended conjunction rule $``$ either directly or as a derived rule. Since $`S`$ is complete, it must have a disjunction rule, therefore it could only fail to be equivalent to $`\text{Mp}=(,,\&)`$ if (a) it has $``$ and $`\&`$ is not derivable, *i.e.*, $`S`$ is equivalent to $`(,)`$, or (b) it has $`\&`$ and $``$ is not derivable, *i.e.*, $`S`$ is equivalent to $`(,\&)`$. In case (a), the valid formula $`P\overline{P}`$ would not be derivable, and in case (b) the valid formula $`(P\overline{P})Q`$ would not be derivable, either way contradicting the completeness of $`S`$. ∎
* Proof of the Extended Minimality Theorem (Theorem 4). Suppose $`S`$ is a complete extended system. If $`S`$ has contraction $`𝖢`$ or weakening $`𝖶`$ then it is equivalent to a standard system by Lemma 14, hence contains Mp by the original Minimality Theorem. Otherwise $`S`$ is equivalent to Mp by Lemma 15, hence in particular contains Mp.
Conversely, suppose $`S`$ is an extended system containing Mp. Then $`S`$ is complete since Mp is complete. $`\mathrm{}`$
We also have the following extension of Theorem 3 (page 3), which stated that, up to equivalence, there are only three complete standard systems, GS1p, Pp and Np.
###### Theorem 5
Up to equivalence, there are only four complete extended systems:
1. The Gentzen-Schütte-Tait system $`\text{GS1p}=(\&,,𝖢,𝖶)`$.
2. Positive calculus $`\text{Pp}=(,,𝖢)`$.
3. Negative calculus $`\text{Np}=(\&,\&,𝖶)`$.
4. Propositional minimal sequent calculus $`\text{Mp}=(,,\&)`$.
###### Proof.
Theorem 3 together with Lemmas 14 and 15. ∎
## 6 Degrees of completeness
We defined a system as *complete* if every valid formula (singleton sequent) is derivable. To avoid ambiguity with forthcoming definitions, let us refer to this default notion of completeness as formula-completeness. Define a system as minimal-complete if every minimal<sup>12</sup><sup>12</sup>12Recall that a valid sequent is minimal if no proper subsequent is valid. sequent is derivable, and sequent-complete if every valid sequent is derivable. (Thus *sequent-complete* implies *minimal-complete* implies *formula-complete*.)
For a minimal-complete system $`S`$, a sequent $`\mathrm{\Gamma }`$ is valid iff a subsequent of $`\mathrm{\Gamma }`$ is derivable in $`S`$. Thus a minimal-complete system $`S`$ can be viewed as sequent-complete, modulo final weakenings. (*Cf.* system GS5p of \[TS96, §7.4\] (related to resolution).)
###### Proposition 4
$`\text{Pp}=(,,𝖢)`$ and $`\text{Mp}=(,,\&)`$ are formula-complete and minimal-complete, but not sequent-complete.
###### Proof.
We have already proved that Mp (hence also Pp, by containment) is minimal-complete (Proposition 1).
We show that the valid (non-minimal) sequent $`P,\overline{P},Q`$ is not derivable in Pp (hence also in Mp). A derivation must contain an axiom rule $`\overline{P,\overline{P}}`$. This cannot be followed by a $``$ or $``$ rule, otherwise we introduce a connective $``$ or $``$ which cannot subsequently be removed by any other rule before the concluding sequent $`P,\overline{P},Q`$. Neither can it be followed by contraction $`𝖢`$, since there is nothing to contract. ∎
###### Proposition 5
$`\text{Np}=(\&,\&,𝖶)`$ is formula-, minimal- and sequent-complete.
###### Proof.
Np is minimal-complete since it contains Mp. Suppose $`\mathrm{\Gamma }`$ is a valid but not minimal sequent. Choose a minimal subsequent $`\mathrm{\Delta }`$ of $`\mathrm{\Gamma }`$ (see Lemma 1, page 1). By minimal-completeness, $`\mathrm{\Delta }`$ has a derivation. Follow this with weakenings to obtain $`\mathrm{\Gamma }`$. ∎
Below we have annotated our Hasse diagram with completeness strengths.
## 7 Possible future work
1. *Cut.* Chapter 4 of \[TS96\] gives a detailed analysis of cut for Gentzen systems. One could pursue an analogous analysis of cut for minimal sequent calculus. Aside from context-splitting and context-sharing cut rules
tensy $`\mathrm{\Delta },A\mathrm{\Sigma },\neg A`$ $`\mathrm{𝖼𝗎𝗍}_{}`$ $`\mathrm{\Delta },\mathrm{\Sigma }`$ tensy $`\mathrm{\Gamma },A\mathrm{\Gamma },\neg A`$ $`\mathrm{𝖼𝗎𝗍}_\&`$ $`\mathrm{\Gamma }`$
one might also investigate a blended cut rule:
tensy $`\mathrm{\Gamma },\mathrm{\Delta },A\mathrm{\Gamma },\mathrm{\Sigma },\neg A`$ $`\mathrm{𝖼𝗎𝗍}`$ $`\mathrm{\Gamma },\mathrm{\Delta },\mathrm{\Sigma }`$
2. *Quantifiers.* Explore the various ways of adding quantifiers to Mp, for a full first-order system M.
3. *Mix (nullary multicut).* Gentzen’s multicut rule
tensy $`\mathrm{\Delta },A_1,\mathrm{},A_m\mathrm{\Sigma },\neg A_1,\mathrm{},\neg A_n`$ $`\mathrm{\Delta },\mathrm{\Sigma }`$
in the nullary case $`m=n=0`$ has been of particular interest to linear logicians \[Gir87\], who call it the *mix* rule. One could investigate context-splitting, context-sharing and blended incarnations:
tensy
ΔΣ
mix
Δ,Σ
tensy
ΓΓ
mix&
Γ
tensy
Γ,ΔΓ,Σ
mix
Γ,Δ,Σ
𝑡𝑒𝑛𝑠𝑦
ΔΣ
mix
Δ,Σ
𝑡𝑒𝑛𝑠𝑦
ΓΓ
mix&
Γ
𝑡𝑒𝑛𝑠𝑦
Γ,ΔΓ,Σ
mix
Γ,Δ,Σ
{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Delta\hskip 17.22217pt\Sigma$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=32.77776pt\hbox{\kern 3.00003pt$\mathsf{\mkern 1.0mumix}_{\otimes}$}}}\hbox{\kern 3.61116pt\hbox{$\displaystyle\;\Delta,\Sigma\;$}}}}\hskip 34.44434pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma\hskip 17.22217pt\Gamma$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=29.7222pt\hbox{\kern 3.00003pt$\mathsf{\mkern 1.0mumix}_{\&}$}}}\hbox{\kern 8.95839pt\hbox{$\displaystyle\;\Gamma\;$}}}}\hskip 34.44434pt{tensy\vbox{\hbox spread0.0pt{\hskip 0.0pt plus 0.0001fil\hbox{$\displaystyle\penalty 1\Gamma,\Delta\hskip 17.22217pt\Gamma,\Sigma$}\hskip 0.0pt plus 0.0001fil}\hbox{\hbox{\kern 0.0pt\vrule height=0.40001pt,depth=0.40001pt,width=51.94438pt\hbox{\kern 3.00003pt$\mathsf{\mkern 1.0mumix}$}}}\hbox{\kern 8.40282pt\hbox{$\displaystyle\;\Gamma,\Delta,\Sigma\;$}}}}
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# Adhesive interactions of viscoelastic spheres
## 1 INTRODUCTION
Numerous phenomena observed in granular systems, ranging from sand and powders to granular gases in planetary rings or protoplanetary discs, are direct consequences of the specific particle interactions. Besides elastic forces, common for molecular or atomic materials (solids, liquids, and gases), colliding granular particles exert also dissipative forces. These forces acting between contacting grains give rise to unusual properties of granular matter. Hence, the use of an appropriate model of the dissipative interactions is necessary for the adequate description of granular systems. In real granular systems the particles may have a complicated non-spherical shape, differ in size, mass and material properties. In what follows, however, we assume that granular particles are smooth spheres of the same material. We also assume that particles interact exclusively via pairwise mechanical contact.
There are three different types of forces acting between contacting particles: (i) repelling elastic forces, due to the compression of particles, (ii) attractive adhesive forces which appear when particles share a common surface and (iii) dissipative forces, acting against the relative motion of the particles. The dependence on the material parameters and on the quantities which characterize the relative position and motion of particles is known for all of these forces, e.g. \[Brilliantov and Pöschel (2004a)\]. The total force acting between the particles is, however, not just the sum of the above three components. Instead, a more careful analysis, sketched below, reveals that the total force contains an additional cross-term which depends on both dissipative and adhesive parameters. Using the obtained expression for the total force acting between adhesive viscoelastic spheres, we estimate the effect of the adhesive force on the coefficient of restitution and analyze the range of validity of the frequently used viscoelastic interaction model \[Brilliantov et al. (1996), Brilliantov and Pöschel (2004b)\].
## 2 FORCES OF GRANULAR PARTICLES
Elastic force. Consider a static contact of two spheres of radii $`R_1`$ and $`R_2`$. When the spheres are squeezed, the material in the bulk is deformed (Fig. 1).
The displacement field $`\stackrel{}{u}\left(\stackrel{}{r}\right)`$ causes the deformation field
$$u_{ij}(\stackrel{}{r})=\frac{1}{2}\left(\frac{u_i}{x_j}+\frac{u_j}{x_i}\right),i,j=\{x,y,z\}.$$
(1)
In the elastic regime $`u_{ij}(\stackrel{}{r})`$ is proportional to the stress tensor $`\sigma ^{ij}(\stackrel{}{r})`$, which gives the $`i`$-component of the force, acting on a unit surface normal to the direction $`j`$ \[Landau and Lifshitz (1965)\]:
$$\sigma ^{ij}(\stackrel{}{r})=E_1(u_{ij}(\stackrel{}{r})\frac{1}{3}\delta _{ij}u_{ll}(\stackrel{}{r}))+E_2\delta _{ij}u_{ll}(\stackrel{}{r}).$$
(2)
Repeated indices imply summation and the elastic coefficients $`E_1`$ and $`E_2`$ are related to the Young modulus $`Y`$ and the Poisson ratio $`\nu `$ by
$$E_1=\frac{Y}{(1+\nu )},E_2=\frac{Y}{3(12\nu )}.$$
(3)
The contact problem for elastic spheres has been solved by Hertz (?): the circular contact of radius $`a`$ gives rise to the elastic force, which read in terms of the compression $`\xi R_1+R_2\left|\stackrel{}{r}_1\stackrel{}{r}_2\right|`$,
$$a^2=R_{\mathrm{eff}}\xi ,F_H=\frac{\sqrt{R_{\mathrm{eff}}}}{D}\xi ^{3/2}$$
(4)
where
$$R_{\mathrm{eff}}\frac{R_1R_2}{R_1+R_2},D\frac{3}{2}\frac{(1\nu ^2)}{Y}.$$
(5)
The normal pressure $`P_H(x,y)\sigma ^{zz}(x,y,z=0)`$, which acts between the compressed bodies in the plane of contact $`z=0`$ reads \[Hertz (1882)\]
$$P_H(x,y)=\frac{3F_H}{2\pi ab}\sqrt{1\frac{x^2}{a^2}\frac{y^2}{b^2}}.$$
(6)
Adhesive force. The Hertz theory for the elastic contact of spheres was extended to adhesive contact by Johnson, Kendall, and Roberts (JKR) (?). They found that the contact area is enlarged owing to the adhesive force and, thus, introduced an effective Hertz load $`F_H`$ which would cause this enlarged area. The contact area of radius $`a`$ corresponds then to the compression $`\xi _H`$ for the Hertz load $`F_H`$. In reality, however, this contact radius occurs at the compression $`\xi <\xi _H`$. The difference between the Hertz compression $`\xi _H`$ and the actual one, $`\xi `$, was attributed to the additional stress
$$P_B(x,y)\frac{F_B}{2\pi a^2}\left(1\frac{x^2}{a^2}\frac{y^2}{a^2}\right)^{1/2},$$
(7)
which is the solution of the classical Boussinesq problem \[Timoshenko (1970)\]: This distribution of the normal surface traction gives rise to a constant displacement over a circular region of an elastic body. The displacement $`\xi _B`$ corresponding to the contact radius $`a`$ and the total load $`F_B`$ are related by
$$\xi _B=\frac{2DF_B}{3a},D\frac{3}{2}\frac{(1\nu ^2)}{Y}.$$
(8)
The quantity $`F_B`$ stands for the additional attractive adhesive force, which acts against $`F_H`$. In the JKR theory it is expressed by the adhesion coefficient $`\gamma `$, which is twice the surface free energy per unit area of a solid in vacuum:
$$F_B=2\pi a^2\sqrt{\frac{3\gamma }{2\pi Da}}.$$
(9)
Thus, the total force $`F=F_HF_B`$ and the actual compression, $`\xi =\xi _H\xi _B`$, read
$`\xi (a)={\displaystyle \frac{a^2}{R_{\mathrm{eff}}}}\sqrt{{\displaystyle \frac{8\pi \gamma Da}{3}}}`$ (10)
$`F(a)={\displaystyle \frac{a^3}{DR_{\mathrm{eff}}}}\sqrt{{\displaystyle \frac{6\pi \gamma }{D}}}a^{3/2},`$ (11)
implying a finite contact radius for $`F=0`$:
$$a_0^3=6\pi D\gamma R_{\mathrm{eff}}^2.$$
(12)
For decreasing force ($`F<0`$) the contact radius decreases until the minimal value
$$a_{\mathrm{sep}}^3=\frac{1}{4}a_0^3$$
(13)
corresponding to maximal (in absolute value) negative force. At this point the particles separate.
Viscous force. If the deformation of contacting spheres changes with time, an additional dissipative force arises and the stress tensor contains an additional dissipative component $`\sigma _{\mathrm{dis}}^{ij}`$. For small deformation rate $`\dot{u}_{ij}(\stackrel{}{r})`$ it reads \[Landau and Lifshitz (1965)\],
$$\sigma _{\mathrm{dis}}^{ij}(t)=\eta _1\left[\dot{u}_{ij}(t)\frac{1}{3}\delta _{ij}\dot{u}_{ll}(t)\right]+\eta _2\delta _{ij}\dot{u}_{ll}(t),$$
(14)
where $`\eta _1`$ and $`\eta _2`$ are the viscous constants.
If the impact velocity of the colliding bodies is much smaller than the speed of sound in the particle material and if the characteristic relaxation time of the viscous processes in the bulk of the material is much smaller than the duration of the collision, one can apply the quasistatic approximation \[Brilliantov et al. (1996)\]. In this approximation the displacement field $`\stackrel{}{u}(\stackrel{}{r})`$ coincides with that for the static case $`\stackrel{}{u}_{\mathrm{el}}(\stackrel{}{r})`$, which is the solution of the corresponding elastic problem. The field $`\stackrel{}{u}_{\mathrm{el}}(\stackrel{}{r})`$ in its turn is completely determined by the time-dependent compression $`\xi `$, i.e., $`\stackrel{}{u}_{\mathrm{el}}=\stackrel{}{u}_{\mathrm{el}}(\stackrel{}{r},\xi )`$ \[Brilliantov et al. (1996)\]. Relating $`\xi `$ and $`a`$ via Eq. (10) and neglecting a small hysteresis in the very beginning of the contact \[Smith et al. (1989)\], we obtain for adhesive interaction
$$\dot{\stackrel{}{u}}(\stackrel{}{r},t)\dot{\xi }\frac{}{\xi }\stackrel{}{u}_{\mathrm{el}}(\stackrel{}{r},\xi )=\dot{a}\frac{}{a}\stackrel{}{u}_{\mathrm{el}}(\stackrel{}{r},\xi (a)).$$
(15)
The dissipative stress tensor reads, respectively
$$\sigma _{\mathrm{dis}}^{ij}=\dot{a}\frac{}{a}\left[\eta _1\left(u_{ij}^{\mathrm{el}}\frac{1}{3}\delta _{ij}u_{ll}^{\mathrm{el}}\right)+\eta _2\delta _{ij}u_{ll}^{\mathrm{el}}\right].$$
(16)
From Eqs. (16) and (2) follows the relation between the elastic and dissipative stress tensors in quasistatic approximation \[Brilliantov (2005)\],
$$\sigma _{\mathrm{dis}}^{ij}=\dot{a}\frac{}{a}\sigma _{\mathrm{el}}^{ij}(E_1\eta _1,E_2\eta _2),$$
(17)
meaning that the dissipative tensor is obtained from the corresponding elastic tensor by substituting the elastic constants by the viscous constants and applying the operator $`\dot{a}/a`$. In particular the normal component of the stress tensor at the plane $`z=0`$ for an adhesive contact reads
$$\sigma _{\mathrm{el}}^{zz}(x,y,z=0)=P_H(x,y)P_B(x,y),$$
(18)
with $`P_H(x,y)`$ and $`P_B(x,y)`$ given by Eqs. (6,7,9).
From Eqs. (17,18) we find the dissipative stress at the contact plane and, integrating it over the contact area, the dissipative force. Referring for detail to \[Brilliantov (2005)\], we present here the final result:
$`F_{\mathrm{dis}}=\dot{a}\left(A{\displaystyle \frac{3a^2}{DR_{\mathrm{eff}}}}+{\displaystyle \frac{3}{2}}B\sqrt{{\displaystyle \frac{6\pi \gamma }{D}}}\sqrt{a}\right)`$ (19)
$`A{\displaystyle \frac{1}{3}}{\displaystyle \frac{\left(3\eta _2\eta _1\right)^2}{\left(3\eta _2+2\eta _1\right)}}\left[{\displaystyle \frac{\left(1\nu ^2\right)(12\nu )}{Y\nu ^2}}\right]`$ (20)
$`B{\displaystyle \frac{(3\eta _2\eta _1)(1+\nu )(12\nu )}{3Y\nu }}.`$ (21)
The first term in the rhs of Eq. (19) corresponds to the dissipative force in the absence of adhesion. The second term is the corresponding cross-term, which depends on both adhesive and dissipative parameters.
Remind that the above relations for the dissipative force have been obtained within the simple JKR theory, which is quite accurate in the range of parameters of practical interest \[Attard and Parker (1992)\].
## 3 COEFFICIENT OF RESTITUTION
An important characteristic of rarefied systems (granular gases \[Brilliantov and Pöschel (2004b)\]) is the coefficient of restitution, which quantifies the loss of mechanical energy for pairwise collisions. It relates the pre-collision relative velocity, $`gv_1v_2`$, to that after the collision, $`g^{}=v_1^{}v_2^{}`$:
$$\epsilon (g)g^{}/g.$$
(22)
The coefficient of restitution may be evaluated solving the two-body collision problem with given interaction forces, yielding $`g^{}`$ as a function of $`g`$.
For small velocities, when the kinetic energy of the relative motion of colliding particles is close to the surface energy of the contact, the adhesive forces may change the coefficient of restitution qualitatively. Indeed, adhesive particles may stay compressed in contact even if the external load vanishes. That is, a tensile force must be applied to separate the particles. The work against this tensile force at the very end of the collision reduces the kinetic energy of the relative motion after the impact, that is, it reduces the coefficient of restitution. For small impact velocity the kinetic energy of the relative motion may be too small to overcome the attractive barrier, i.e., the particles stick together after the collision corresponding to $`\epsilon =0`$.
Consider first a pure viscoelastic collision with the impact velocity $`g`$ and coefficient of restitution $`\epsilon _\mathrm{v}`$. Using the definition (22), the energy balance reads
$$\frac{1}{2}m^{\mathrm{eff}}g^2\frac{1}{2}m^{\mathrm{eff}}\epsilon _\mathrm{v}^2(g)g^2=W_{\mathrm{dis}}.$$
(23)
where
$$m^{\mathrm{eff}}\frac{m_1m_2}{m_1+m_2}.$$
(24)
The work $`W_{\mathrm{dis}}`$ results from the dissipative force Eq. (19) with $`\gamma =0`$. The corresponding coefficient of restitution $`\epsilon _\mathrm{v}`$ is analytically known \[Schwager and Pöschel (1998), Ramírez et al. (1999)\]. Turn now to viscoelastic impacts with adhesion, characterized by the coefficient of restitution $`\epsilon _{\mathrm{ad}}`$. The energy balance reads
$$\frac{1}{2}m^{\mathrm{eff}}g^2\frac{1}{2}m^{\mathrm{eff}}\epsilon _{\mathrm{ad}}^2(g)g^2=W_{\mathrm{dis}\mathrm{A}}+W_{\mathrm{dis}\mathrm{B}}+W_{\mathrm{ad}},$$
(25)
where $`W_{\mathrm{dis}\mathrm{A}}`$ is the work of the dissipative force $`F_{\mathrm{dis}}`$ due to the first term in Eq. (19) and $`W_{\mathrm{dis}\mathrm{B}}`$ is the work due to the second term. Finally, $`W_{\mathrm{ad}}`$ is the work due to adhesion, i.e., it results from the adhesive force in the region where the total force Eq. (11) is negative, that is, in the region where the contact radius varies from $`a_0`$ to $`a_{\mathrm{sep}}`$ \[Brilliantov and Pöschel (2004a)\]:
$$W_{\mathrm{ad}}=_{\xi (a_0)}^{\xi (a_{\mathrm{sep}})}F(\xi )𝑑\xi =_{a_0}^{a_{\mathrm{sep}}}F(a)\frac{d\xi }{da}𝑑a.$$
(26)
Using the approximation $`W_{\mathrm{dis}\mathrm{A}}W_{\mathrm{dis}}`$ from Eqs. (23,25) we obtain the coefficient of restitution for viscoelastic collisions with adhesion:
$$\epsilon _{\mathrm{ad}}^2(g)=\epsilon _\mathrm{v}^2(g)\frac{2(W_{\mathrm{ad}}+W_{\mathrm{dis}\mathrm{B}})}{m^{\mathrm{eff}}g^2}.$$
(27)
$`W_{\mathrm{ad}}`$ may be found employing Eq. (11) for the total force $`F(a)`$, Eq. (10) for the compression, which allows to obtain $`d\xi /da`$, and Eqs. (12,13) for $`a_0`$ and $`a_{\mathrm{sep}}`$:
$$W_{\mathrm{ad}}=q_0\left(\pi ^5\gamma ^5D^2R_{\mathrm{eff}}^4\right)^{1/3},$$
(28)
with the analytically known pure number $`q_00.578`$ \[Brilliantov and Pöschel (2004a)\]. For small dissipation and small adhesion, $`W_{\mathrm{dis}\mathrm{B}}`$ may be roughly estimated as \[Brilliantov (2005)\]:
$$W_{\mathrm{dis}\mathrm{B}}=q_1B\gamma ^{1/2}D^{1/5}g^{8/5}\left(m^{\mathrm{eff}}R_{\mathrm{eff}}^2\right)^{3/10},$$
(29)
where $`q_17.93`$.
From Eq. (27) we obtain the condition for the validity of the viscoelastic collision model:
$$g^2\frac{2}{m^{\mathrm{eff}}}\left(W_{\mathrm{ad}}+W_{\mathrm{dis}\mathrm{B}}\right),$$
(30)
which with Eqs. (28,29) may be written as
$`g^2g_c^2=\left(\gamma ^5D^2R_{\mathrm{eff}}^4\right)^{1/3}/m^{\mathrm{eff}}`$ (31)
$`B\left(D^2/\gamma R_{\mathrm{eff}}^2\right)^{1/6}\left(m^{\mathrm{eff}}\right)^{1/2}.`$ (32)
With the above condition for $`B`$ one can neglect $`W_{\mathrm{dis}\mathrm{B}}`$ as compared to $`W_{\mathrm{ad}}`$ for the impact velocity $`g`$ being of the order of $`g_c`$ ($`gg_c`$). Then we obtain respectively the condition for sticking collision when $`\epsilon _{\mathrm{ad}}(g_{\mathrm{st}})=0`$:
$$\frac{1}{2}m^{\mathrm{eff}}\epsilon _\mathrm{v}^2(g_{\mathrm{st}})g_{\mathrm{st}}^2=W_{\mathrm{ad}}.$$
(33)
Hence, if $`g<g_{\mathrm{st}}`$ for head-on collisions (vanishing tangential component of the impact velocity), the colliding particles stick together and form a joint particle of mass $`m_1+m_2`$.
## 4 CONCLUSION
The collision of adhesive viscoelastic spheres is characterized by (i) the elastic Hertz force, (ii) the dissipative force originating from viscoelastic bulk deformation, and (iii) the adhesive force. We use the continuum model of adhesive contact by Johnson, Kendall, and Roberts (?) which is adequate in the range of parameters of practical interest. The total force was derived under the approximation of quasistatic deformation, that is, the impact velocity is assumed to be much smaller than the speed of sound in the material and the viscous relaxation time is much smaller than the duration of the collision. This force is not only the superposition of its three components (i-iii), but there appears an additional cross-term, which depends on both viscous and adhesive parameters of the material.
Using this force we estimated the contribution of adhesive forces to the normal coefficient of restitution as well as the range of validity of the viscoelastic collision model \[Brilliantov et al. (1996)\] and the condition for sticking impact of head-on collisions.
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# Localizations at infinity and essential spectrum of quantum Hamiltonians: I. General theory
## 1 Introduction
In this paper we continue the investigation of the spectral properties of quantum Hamiltonians with $`C^{}`$-algebra methods on the lines of our previous work \[GI4\]. More precisely, our aim is to study the essential spectrum of general classes of (unbounded) operators in $`L^2(X)`$, where $`X`$ is a locally compact non-compact abelian group, by using crossed product techniques. For some historical remarks and comparison with other recently obtained results, see Subsections 1.7 and 1.16.
1.1. We set $`B(X)=B(L^2(X))`$ and we denote by $`U_x`$ the operator of translation by $`xX`$ and by $`V_k`$ the operator of multiplication by the character $`kX^{}`$ (our notations, although rather standard, are summarized in Section 2). We define <sup>2</sup><sup>2</sup>2 We make the following convention: if a symbol like $`T^{()}`$ appears in a relation, then the relation must hold for the operator $`T`$ and for its adjoint $`T^{}`$.
(1.1)
$$C(X)=\{TB(X)\underset{k0}{lim}[T,V_k]=0\text{ and }\underset{x0}{lim}(U_x1)T^{()}=0\}$$
which is clearly a $`C^{}`$-algebra of operators on $`L^2(X)`$ (without unit if $`X`$ is not discrete). Besides the norm topology on $`C(X)`$ we shall also consider on it the topology defined by the family of seminorms $`S_\theta =S\theta (Q)+\theta (Q)S`$ with $`\theta 𝒞_0(X)`$ and we shall denote $`C_s(X)`$ the corresponding topological space (see Remark 5.7). Here $`\theta (Q)`$ is the operator of multiplication by $`\theta `$ in $`L^2(X)`$.
Our main result is a description of the essential spectrum of the operators $`TC`$ in terms of their “localizations at infinity”. We denote by $`\delta X`$ the set of all ultrafilters on $`X`$ finer than the Fréchet filter (cf. page 2.9). If $`A_i`$ are subsets of a topological space we denote $`\overline{}_{iI}A_i`$ the closure of their union.
###### Theorem 1.1
If $`TC(X)`$ is a normal operator, then for each $`\varkappa \delta X`$ the limit $`lim_{x\varkappa }U_xTU_x^{}=\varkappa .T`$ exists in $`C_s(X)`$ and
(1.2)
$$\sigma _{\mathrm{ess}}(T)=\overline{}_\varkappa \sigma (\varkappa .T).$$
Note that $`x\varkappa `$ should be read “$`x`$ tends to infinity along the filter $`\varkappa `$”. The limit operator $`\varkappa .T`$ will be called *localization at $`\varkappa `$ of $`T`$*. Since an ultrafilter finer than the Fréchet filter can be thought as a point on an ideal boundary at infinity of $`X`$, the operators $`\varkappa .T`$ will also be called *localizations at infinity of $`T`$*.
We are mainly interested in the essential spectrum of unbounded self-adjoint operators $`H`$ “affiliated” to $`C(X)`$, but the corresponding result is an immediate consequence of Theorem 1.1. We say that $`H`$ is affiliated to some $`C^{}`$-algebra $`A`$ of operators on $`L^2(X)`$ if $`\phi (H)A`$ for all $`\phi 𝒞_0()`$ (for this it suffices to have $`(Hz)^1A`$ for one $`z\rho (H)`$). For technical reasons we have to consider self-adjoint operators which are not necessarily densely defined and, in order to avoid confusions with the standard terminology, we shall call these more general objects *observables*. A more detailed presentation of this notion can be found in Subsection 2. For the moment we note only that an observable $`H`$ is affiliated to $`C(X)`$ if and only if
(1.3)
$$\underset{k0}{lim}[V_k,(Hz)^1]=0\text{and}\underset{x0}{lim}(U_x1)(Hz)^1=0$$
for some $`z\rho (H)`$. This follows from the fact that if $`TB(X)`$ is normal, then
$$\underset{x0}{lim}(U_x1)T=0\underset{x0}{lim}(U_x1)T^{}=0.$$
###### Theorem 1.2
Let $`H`$ be an observable on $`L^2(X)`$ affiliated to $`C(X)`$. Then for each $`\varkappa \delta X`$ the limit $`\varkappa .H:=lim_{x\varkappa }x.H`$ exists in the following sense: there is an observable $`\varkappa .H`$ affiliated to $`C(X)`$ such that $`lim_{x\varkappa }U_x\phi (H)U_x^{}=\phi (\varkappa .H)`$ in $`C_s(X)`$ for all $`\phi 𝒞_0()`$. Moreover, we have
(1.4)
$$\sigma _{\mathrm{ess}}(H)=\overline{}_\varkappa \sigma (\varkappa .H).$$
Practically we are interested only in the case when $`H`$ is a self-adjoint operator in the standard sense. However, even in this case $`\varkappa .H`$ could be not densely defined and quite often we have $`\varkappa .H=\mathrm{}`$ (i.e. the domain of $`\varkappa .H`$ is $`\{0\}`$). For example, if $`H`$ has purely discrete spectrum, then $`\phi (H)`$ is a compact operator and we clearly get $`\varkappa .H=\mathrm{}`$ for all $`\varkappa `$. Since $`\sigma (\mathrm{})=\mathrm{}`$, we then obtain $`\sigma _{\mathrm{ess}}(H)=\mathrm{}`$, as it should be.
###### Remark 1.3
The observable $`H`$ should be thought as the Hamiltonian (energy observable) of a physical system. Thus (1.4) says that the essential spectrum of the Hamiltonian $`H`$ can be computed in terms of the spectra of its localizations at infinity $`\varkappa .H`$. We emphasize that this notion of infinity is determined by the position observable $`Q`$. In other terms, if $`H`$ satisfies (1.3) then $`\sigma _{\mathrm{ess}}(H)`$ is given by its localizations in the region $`Q=\mathrm{}`$. This property does not hold in many situations of physical interest (e.g. if magnetic fields which do not vanish at infinity are involved) because localizations at infinity with respect to other observables must be taken into account, see \[GI3\].
###### Remark 1.4
It will be clear from the proof of Theorem 1.2 (see Lemma 5.3 and Proposition 5.10) that (1.4) remains valid if $`\varkappa `$ runs over sets much smaller than $`\delta X`$: it suffices to take $`\varkappa 𝒦`$ if $`𝒦\delta X`$ has the property: if $`\phi `$ is a bounded uniformly continuous function on $`X`$ and $`lim_{x\varkappa }\phi (x+y)=0`$ for all $`yX,\varkappa 𝒦`$, then $`\phi 𝒞_0(X)`$.
###### Remark 1.5
We mention the following immediate consequence of (1.4): *if two observables affiliated to $`C(X)`$ have the same localizations at infinity, then they have the same essential spectrum*. If the difference of the resolvents is a compact operator, then clearly they have the same localizations at infinity, but the converse is far from being true (e.g. see the example on page 531 from \[GI4\], where the essential spectrum is independent of the details of the shape of the function $`\omega `$). On the other hand, one may find in \[GG2\] criteria which ensure the compactness of the difference of the resolvents of two self-adjoint operators under rather weak conditions, e.g. an example from \[GG2, p. 26 \] is a general version of \[LaS, Proposition 4.1\].
###### Remark 1.6
The following remark is useful in applications: if $`H`$ is an observable affiliated to $`C(X)`$ and if $`\theta :\sigma (H)`$ is a proper continuous function, then $`\theta (H)`$ is affiliated to $`C(X)`$ and we have $`\varkappa .\theta (H)=\theta (\varkappa .H)`$ for all $`\varkappa \delta X`$ (see page 2).
###### Remark 1.7
As explained in \[GI4, p. 520\], all our results extend trivially to the case when $`L^2(X)`$ is replaced with the space of $`L^2`$ functions with values in a Hilbert space $`𝐄`$: it suffices to replace the algebra $`𝒜`$ with $`𝒜K(𝐄)`$. For example, Theorem 1.2 remains valid without any change if $`L^2(X)`$ is replaced by $`L^2(X;𝐄)`$, where $`𝐄`$ is finite dimensional, and $`C(X)`$ is defined exactly as before. Thus in applications we can consider differential operators with matrix valued coefficients, like Dirac operators.
1.2. We give here the simplest applications of Theorems 1.1 and 1.2, a more detailed study and more general examples can be found in Section 4.
Assume first that $`X`$ is discrete. Note that in the particularly important case $`X=^n`$ Theorem 1.1 has been proved in \[RRS1\] (with a slightly different formulation and with quite different methods). Now we have
(1.5)
$$C(X)=\{TB(X)\underset{k0}{lim}[T,V_k]=0\}.$$
Since $`V_k^{}U_xV_k=k(x)U_x`$ we see that each operator of the form $`T=_{aX}\phi _a(Q)U_a`$, with $`\phi _a\mathrm{}_{\mathrm{}}(X)`$ and $`\phi _a0`$ only for a finite number of $`a`$, belongs to $`C(X)`$. Clearly, we have $`\varkappa .T=_{aX}(\varkappa .\phi _a)(Q)U_a`$ where the function $`\varkappa .\phi _a\mathrm{}_{\mathrm{}}(X)`$ is defined by $`(\varkappa .\phi _a)(y)=lim_{x\varkappa }\phi (x+y)`$. The Jacobi and CMV operators considered in \[LaS\] are particular cases of such operators $`T`$.
Now we give three examples in the case $`X=^n`$. We start with the Schrödinger operator. We denote by $`^s`$ the Sobolev space of order $`s`$ associated to $`L^2(^n)`$. Note that $`\mathrm{\Delta }`$ is the positive Laplacian. From Proposition 4.12 we get:
###### Proposition 1.8
Let $`W`$ be a continuous symmetric sesquilinear form on $`^1`$ such that:
(1) $`W\mu \mathrm{\Delta }\nu `$ as forms on $`^1`$ for some numbers $`\mu <1`$ and $`\nu >0`$,
(2) $`lim_{k0}[V_k,W]_{^1^1}=0`$.
Let $`H_0`$ be the self-adjoint operator associated to the form sum $`\mathrm{\Delta }+W`$ and let $`V`$ be a real function in $`L_{\mathrm{loc}}^1(^n)`$ such that its negative part is relatively bounded with respect to $`H_0`$ with relative bound $`<1`$. Then the self-adjoint operator $`H=H_0+V(Q)`$ (form sum) is affiliated to $`C(^n)`$, hence the conclusions of Theorem 1.2 hold for it.
This can be extended to a general class of hypoelliptic operators, cf. Proposition 4.16. We present below a very particular case.
###### Proposition 1.9
Let $`h:^n`$ be of class $`C^m`$ for some $`m1`$ and such that:
(1) $`lim_k\mathrm{}h(k)=+\mathrm{}`$,
(2) the derivatives of order $`m`$ of $`h`$ are bounded,
(3) $`_{|\alpha |m}|h^{(\alpha )}(k)|C(1+|h(k)|)`$.
Let $`𝒢=D(|h(P)|^{1/2})`$ be the form domain of the operator $`h(P)`$ and assume that $`W`$ is a symmetric continuous form on $`𝒢`$ such that:
(4) $`W\mu h(P)\nu `$ as forms on $`𝒢`$ for some numbers $`\mu <1`$ and $`\nu >0`$,
(5) $`lim_{k0}[V_k,W]_{𝒢𝒢^{}}=0`$.
Let $`H_0=h(P)+W`$ (form sum) and let $`VL_{\mathrm{loc}}^1(^n)`$ real such that its negative part is relatively bounded with respect to $`H_0`$ with relative bound $`<1`$. Then the self-adjoint operator $`H=H_0+V(Q)`$ (form sum) is affiliated to $`C(^n)`$, hence the conclusions of Theorem 1.2 hold for it.
###### Remark 1.10
If $`X`$ is an arbitrary group, $`h:X`$ is continuous and satisfies $`|h(k)|\mathrm{}`$ as $`k\mathrm{}`$, and if $`VL^{\mathrm{}}(X)`$, then obviously $`h(P)+V(Q)`$ is affiliated to $`C(X)`$ and so we can apply Theorem 1.2. In order to cover unbounded $`V`$ without much effort a quite weak regularity condition on $`h`$ is sufficient, see Proposition 4.16 and especially relation (4.8). We shall not try to optimize on this here.
Finally, we consider a Dirac operator $`D`$. Let $`=L^2(^n;𝐄)`$ for some finite dimensional Hilbert space $`𝐄`$. We only need to know that $`D`$ is a symmetric first order differential operator with constant coefficients acting on $`𝐄`$-valued functions and which is realized as a self-adjoint operator on $``$ such that the domain of $`|D|^{1/2}`$ is the Sobolev space $`^{1/2}`$. Now from Corollary 4.8 we get:
###### Proposition 1.11
Let $`W`$ be a continuous symmetric form on $`^{1/2}`$ such that:
(1) $`\pm W\mu |D|+\nu `$ as forms on $`^{1/2}`$ for some numbers $`\mu <1`$ and $`\nu >0`$,
(2) $`lim_{k0}[V_k,W]_{^{1/2}^{1/2}}=0`$.
Then the self-adjoint operator $`H=D+W`$, defined as explained on page 4.7, is affiliated to $`C(^n)`$, hence the conclusions of Theorem 1.2 hold for it.
Observe that condition (2) is trivially satisfied if $`W`$ is the operator of multiplication by an operator valued function $`W:^nB(𝐄)`$.
###### Remark 1.12
We emphasize that the conditions on the perturbation $`W`$ in Propositions 1.8-1.11 is such that $`W`$ can contain terms of the same order as $`\mathrm{\Delta },h(P)`$ or $`D`$ respectively. For example, operators of the form
$$_{j,k}_ja_{jk}_k+\text{ singular lower order terms }$$
with $`a_{jk}L^{\mathrm{}}`$ such that the matrix $`(a_{jk}(x))`$ is bounded from below by a strictly positive constant are already covered by Proposition 1.8. See Example 4.13 for much more general results. These examples may be combined with the Remark 1.6 to cover functions of operators, e.g. $`\sqrt{H}`$ if $`H0`$.
1.3. Crossed products of $`C^{}`$-algebras by the action of $`X`$ play a fundamental rôle in our proof of Theorem 1.1 but we have to stress that they are important for two distinct reasons. First, they are in a natural sense $`C^{}`$-algebras of *energy*<sup>3</sup><sup>3</sup>3 We emphasize “energy” because algebras of observables and crossed products were frequently used in various domains of the quantum theory in the last 50 years, but with different meanings and scopes than here. observables (or quantum Hamiltonians), and hence they allow one to organize the Hamiltonians in classes each having some specific properties, e.g. the essential spectrum of the operators in a class is given by a “canonical” formula specific to that class (see (1.6)). On the other hand, crossed products are very efficient at a technical level, their use allows one to solve a non-abelian problem by abelian means: the problem of computing the quotient of a non-commutative $`C^{}`$-algebra $`AB(X)`$ with respect to the ideal $`K(X)K(L^2(X))`$ is reduced to that of computing $`𝒜/𝒞_0(X)`$ where $`𝒜`$ is a $`C^{}`$-algebra of bounded uniformly continuous functions on $`X`$.
The first reason mentioned above will be clarified by the later developments, but one may observe already now that the decomposition (1.2) is far from efficient. Indeed, its extreme redundancy becomes clear when we realize that many $`\varkappa `$ give the same $`\varkappa .H`$ (e.g. if the filters $`\varkappa `$ and $`\chi `$ have the same envelope then $`\chi .T=\varkappa .T`$, see page 2) and many more give the same $`\sigma (\varkappa .H)`$ (e.g. $`\chi .T=U_x\varkappa .TU_x^{}`$ if $`\chi `$ is the translation by $`xX`$ of $`\varkappa `$).
Thus at a qualitative level (1.2) is not very significant, it does not say much about $`\sigma _{\mathrm{ess}}(H)`$, at least when compared with the $`N`$-body situation where the HVZ theorem has such a nice physical interpretation that you can predict it and believe it without proof.
In order to partially remediate this drawback we consider smaller classes of Hamiltonians. The following framework, introduced in \[GI4\], gives us more specific information about $`\sigma _{\mathrm{ess}}(H)`$. Let $`𝒞(X)`$ be the $`C^{}`$-algebra of all bounded uniformly continuous functions on $`X`$ and $`𝒞_{\mathrm{}}(X)`$ that of continuous functions which have a limit at infinity (in the usual sense).
###### Definition 1.13
An *algebra of interactions $`𝒜`$* on $`X`$ is a $`C^{}`$-subalgebra of $`𝒞(X)`$ which is stable under translations and which contains $`𝒞_{\mathrm{}}(X)`$. The *$`C^{}`$-algebra of quantum Hamiltonians of type $`𝒜`$* is the norm closed linear space $`A𝒜XB(X)`$ generated by the operators of the form $`\phi (Q)\psi (P)`$ with $`\phi 𝒜`$ and $`\psi 𝒞_0(X^{})`$.
We have denoted $`\phi (Q)`$ the operator of multiplication by $`\phi `$ in $`L^2(X)`$ and $`\psi (P)`$ becomes multiplication by $`\psi `$ after a Fourier transformation. The Propositions 3.3 and 3.4 explain why we think of $`A`$ as a $`C^{}`$-algebra of Hamiltonians. For example, if $`X=^n`$, the self-adjoint operators of the form $`\mathrm{\Delta }+_{k=1}^na_k(x)_k+a_0(x)`$ with $`a_j𝒜^{\mathrm{}}`$ (functions in $`𝒜`$ with all derivatives in $`𝒜`$) generate $`A`$. It turns out that $`A`$ is canonically isomorphic with the crossed product of $`𝒜`$ by the natural action of $`X`$, which explains the notation $`𝒜X`$ and the relevance of crossed products in our context.
###### Remark 1.14
Note that the definition and the quoted propositions tend to give the impression that the algebra $`A`$ is rather small. But this is wrong, $`A`$ *is much larger than expected*. For example, $`𝒞(X)X=C(X)`$ and we shall see in Section 4 that the set of self-adjoint operators affiliated to $`C(X)`$ is very large. Other examples are the $`N`$-body algebra and the “bumps” algebras. In fact, we may summarize our approach as follows: we first isolate a class of *elementary Hamiltonians*, these being the simplest operators we would like to study, but our results concern all the operators affiliated to the $`C^{}`$-algebra they generate, which happens to be a crossed product and is very rich.
In order to state the next consequence of Theorem 1.1 we have to introduce some new notations. Let $`\sigma (𝒜)`$ be the space of characters of the abelian $`C^{}`$-algebra $`𝒜`$. Then $`\sigma (𝒜)`$ is a compact topological space which contains $`X`$ as an open dense subset, so $`\delta (𝒜)=\sigma (𝒜)X`$ is a compact space. We shall adopt the following abbreviation: $`H^{}A`$ means that $`H`$ is either a normal element of the algebra $`A`$ or an observable affiliated to $`A`$. If $`H`$ is an observable affiliated to $`A`$ then $`U_xHU_x^{}`$ is also an observable affiliated to $`A`$ and we have $`\phi (U_xHU_x^{})=U_x\phi (H)U_x^{}`$ for $`\phi 𝒞_0()`$. By “continuity” of a map $`\sigma (𝒜)\varkappa \varkappa .H^{}C_s(X)`$ whose values are observables we mean that $`\sigma (𝒜)\varkappa \phi (\varkappa .H)C_s(X)`$ is continuous for all $`\phi 𝒞_0()`$.
###### Theorem 1.15
If $`H^{}A`$ then the map $`XxU_xHU_x^{}^{}A`$ extends to a continuous map $`\sigma (𝒜)\varkappa \varkappa .H^{}C_s(X)`$ and we have
(1.6)
$$\sigma _{\mathrm{ess}}(H)=\overline{}_{\varkappa \delta (𝒜)}\sigma (\varkappa .H).$$
Remark: To see the connection between this and Theorems 1.1 and 1.2, we recall that an ultrafilter finer than the Fréchet filter is the same thing as a character $`\varkappa `$ of the algebra of all bounded functions on $`X`$ such that $`\varkappa (\phi )=0`$ if $`\phi 𝒞_0(X)`$ (see Subsection 2). Moreover, if $`\chi `$ is a second such ultrafilter and $`\varkappa (\phi )=\chi (\phi )`$ for all $`\phi 𝒞(X)`$, then $`\varkappa .H=\chi .H`$ for all $`H^{}C(X)`$, thus the union in (1.2) and (1.4) may be taken in fact over $`\varkappa \delta (𝒞(X))`$. We emphasize that, although Theorem 1.15 seems stronger than Theorems 1.1 and 1.2, it is in fact an immediate consequence of Theorem 1.1 (just “abstract nonsense”, see Subsection 5.14 for details). Note also that (1.6) is a canonical decomposition of the essential spectrum of $`H`$, all the objects in the formula being canonically associated to $`𝒜`$. The representation (1.6) is further discussed in Subsection 5.14, see page 5.7.
###### Remark 1.16
We mention that, by using a more involved algebraic formalism as in \[GI4\], one can obtain partial, but often relevant, information concerning the essential spectrum of $`H`$ as follows. Let $`𝒥`$ be an $`X`$-ideal such that $`𝒞_0(X)𝒥𝒜`$ and let $`J=𝒥X`$ (we use here notations and results from \[GI4\]). Then $`K(X)JA`$ and $`J`$ is an ideal in $`A`$, so the image $`H_𝒥`$ of $`H`$ is well defined as observable affiliated to the quotient algebra $`A/J`$. By using the natural surjection $`A/K(X)A/J`$ we clearly get $`\sigma (H_𝒥)\sigma _{\mathrm{ess}}(H)`$. In this argument $`J`$ need not be a crossed product, but if it is, we can use $`A/J(𝒜/𝒥)X`$ to get a concrete representation of $`H_𝒥`$.
1.4. This subsection is devoted to some historical comments and a discussion of some related results from the literature.
Theorem 1.1 was announced in the preprints \[Ift, GI2\], see Theorems 1.3 and 4.2 in \[Ift\] and Theorem 4.1 and Corollary 4.2 in \[GI2\]. In fact, the theorem was stated in a stronger form, namely we assert that the union in (1.2) is already closed. Moreover, some nontrivial applications are stated at page 149 of \[GI5\]. The closedness of the union in (1.2) as well as more explicit applications of Theorem 1.1 will be discussed in the second part of this paper. However we show here that the union in (1.6) is closed for some special algebras $`𝒜`$ when the result is far from obvious (Section 6).
The main idea of the proof of Theorem 1.1 we had in mind at that moment is presented at \[GI2, p. 30–31\] and it has to be combined with the two main points of the algebraic approach we used in that paper, namely:
* If $``$ is a Hilbert space then the quotient algebra $`B()/K()`$ is a $`C^{}`$-algebra and, if $`\widehat{T}`$ is the projection of $`TB()`$ in the quotient, then $`\sigma _{\mathrm{ess}}(T)=\sigma (\widehat{T})`$.
* We have $`K(L^2(X))=𝒞_0(X)X`$ and if $`𝒜`$ is an algebra of interactions then
(1.7)
$$\left(𝒜X\right)/\left(𝒞_0(X)X\right)\left(𝒜/𝒞_0(X)\right)X.$$
If $`T𝒜X`$ the isomorphism (1.7) allows us to reduce the computation of $`\widehat{T}`$ to an abelian problem and hence to deduce $`\widehat{T}(\varkappa .T)_{\varkappa \delta (𝒜)}_{\varkappa \delta (𝒜)}C(X)`$.
The preceding strategy requires a lot of abstract machinery and is not adapted to a purely Hilbert space setting. For example, the isomorphism (1.7) is a consequence of the fact that the functor $`𝒜𝒜X`$ transforms short exact sequences in short exact sequences, an assertion which does not even make sense if we fix the Hilbert space on which the algebras are realized.
Instead, in the present paper we decided to avoid step (2) of this strategy and to base our arguments on a beautiful theorem due to M. B. Landstad \[Lan\] which gives an intrinsic characterization of crossed products. We feel that this makes the argument more elementary and gives deeper insight into the matters treated here. In fact, we could now avoid completely going out from the purely Hilbert space setting (in particular, forget about the step (1) above), but this does not seem to us a natural attitude and we finally decided to adopt a median approach.
It is remarkable that $`C(X)`$ as defined in (1.1) is precisely the crossed product $`𝒞(X)X`$. Initially, this fact was proved by direct methods in the case $`X=^n`$ in \[DG2\] (because of this Corollary 4.2 from \[GI2\] was stated only for $`X=^n`$). The general case follows in fact immediately from Landstad’s theorem.
We make now some comments concerning other papers with goals similar to ours. We note first that, in the particular case $`X=^n`$, Rabinovich, Roch and Silberman \[RRS1\] discovered Theorem 1.1 before us and proved it with no $`C^{}`$-algebra techniques (in Remark 3.18 we explain why (1.5) is just their algebra of “band dominated operators”). It seems that they realized the fact that their algebra in the case $`X=`$ is a crossed product only in \[RRR\] (this fact is a particular case of \[GI4, Theorem 4.1\]). In \[RRS1\] and in subsequent works \[Rab, RRS1, RRS2\] (see also \[RRS2\] for references to earlier papers) these authors use a discretization technique in order to treat perturbations of pseudo-differential operators in $`L^2(^n)`$. They get relations like (1.4) and show that in some situations the union is already closed. Moreover, in Chapter 7 of \[RRS2\] they present an abstract version of their approach (in particular they consider groups more general than $`^n`$) which seems to us complementary to our approach and relevant in contexts like that of \[Gol\]. We learned about these works quite recently thanks to a correspondence with Barry Simon who sent us a copy of the paper \[Rab\]; this explains why the above references were not included in our previous works on this topic.
We discuss now the relation between our paper and the article \[Ro3\] (this reference was pointed out to us by one of the referees). We shall do it in some detail because $`C^{}`$-algebra techniques are emphasized in \[Ro3\]. The purpose of Roe is to extend the results of Rabinovich, Roch and Silberman to nonabelian groups. He considers a finitely generated discrete (nonabelian) group $`\mathrm{\Gamma }`$ and defines $`A`$ as the $`C^{}`$-algebra of operators on $`\mathrm{}_2(\mathrm{\Gamma })`$ generated by $`\mathrm{}_{\mathrm{}}(\mathrm{\Gamma })`$ and by the right translation operators $`R_\gamma `$ (this is a natural extension of the procedure introduced in \[RRS1\]). Then, denoting $`L_\gamma `$ the left translation operators, he shows that for each $`TA`$ the map $`\gamma L_\gamma TL_\gamma ^{}`$ extends to a $``$-strongly continuous map $`\beta \mathrm{\Gamma }A`$, where $`\beta \mathrm{\Gamma }`$ is the Stone-Čech compactification of $`\mathrm{\Gamma }`$ (the space of characters of $`\mathrm{}_{\mathrm{}}(\mathrm{\Gamma })`$). The restriction to $`\delta \mathrm{\Gamma }=\beta \mathrm{\Gamma }\mathrm{\Gamma }`$ of this map is the *symbol* of $`T`$ and the main result of \[Ro3\] is that for exact groups (in the $`C^{}`$-algebra sense) an operator $`TA`$ has symbol equal to zero if and only if $`T`$ is compact.
On \[GI2, p. 30-31\], where we describe the main ideas of the proof of Theorem 1.2, we introduce the notion of regular operator on $`L^2(X)`$ for $`X`$ an abelian locally compact group (and in a more general context in the footnote on \[GI2, p. 31\]): we say that a bounded operator $`T`$ on $`L^2(X)`$ is regular if $`\{U_xT^{()}U_x^{}xX\}`$ are strongly relatively compact sets. Then we note that for such operators the map $`xU_xTU_x^{}`$ extends to a strongly continuous map $`\beta X\varkappa T_\varkappa B(L^2(X))`$ (this time the Stone-Čech compactification $`\beta X`$ involves the topology of $`X`$) and call the values $`T_\varkappa `$ with $`\varkappa \delta X`$ localizations at infinity of $`T`$. We show that the elements of $`𝒞(X)X`$ are regular and from the arguments on page 31 it is rather obvious that their localizations at infinity belong to the same algebra $`𝒞(X)X`$. This is more explicitly stated and proved in \[GI5, Lemma 3.10\] (which is Lemma 3.9 in the preprint version and Lemma 5.8 here). All this can also be done at the level of the algebra $`𝒞(X)`$ and at the bottom of \[GI2, p. 31\] we say that if $`\phi 𝒞(X)`$ and all its localizations at infinity are zero, then $`\phi 𝒞_0(X)`$ (this is easy to prove, cf. Lemma 5.3 here) and finish by saying that this remains true after taking crossed products (which is not obvious but can be deduced from \[GI2, Theorem 3.4\] or (1.7) here; as we said before, in this paper we prefer to use Landstad’s theorem at this last step).
We emphasize that although the starting points of \[GI2\] and \[Ro3\] (in particular the relevance of the Stone-Čech compactification) are similar, the proofs of the main fact (that the kernel of the symbol map, in the terminology of \[RRS1\], is just the compacts) are of a quite different nature. Indeed, Roe mentions that $`A`$ is the reduced crossed product $`L^{\mathrm{}}(\mathrm{\Gamma })_r\mathrm{\Gamma }`$ but never uses this fact, cf. the proof of \[Ro3, Proposition 3.3\]. On the other hand, the crossed product structure and relations like (1.7) are the heart of our approach (and we expect that (1.7) is also true under Roe’s conditions).
The methods used by Roe also seem relevant for the solution of a problem left open (but not explicitly stated) in \[Gol\]. The space $`\mathrm{\Gamma }`$ considered there is a tree, which is a finitely generated *monoïd*. The natural object in this case is the $`C^{}`$-algebra generated by the right translations and by $`\mathrm{}_{\mathrm{}}(\mathrm{\Gamma })`$, the localizations at infinity being given by left translations. Due to obvious technical difficulties the algebra considered in \[Gol\] is much smaller: it is generated by the Laplacian (which is a certain polynomial in the right translations) and by the functions in $`\mathrm{}_{\mathrm{}}(\mathrm{\Gamma })`$ which extend continuously to the hyperbolic compactification of $`\mathrm{\Gamma }`$ (see Sections 3.4 and 3.5 and Theorem 5.1 and its proof in \[Gol\]; the references are to the preprint version). A larger algebra, associated to the analogue of the slowly oscillating functions on $`\mathrm{\Gamma }`$, is considered in \[GG1\], where the problem is treated by very different techniques. It would be interesting to know if the techniques from \[Ro3, Section 3\] can be adapted to solve the most general situation.
Y. Last and B. Simon obtained in \[LaS\] relations like (1.4) for large classes of Schrödinger operators on $`^n`$ and their discrete versions (Jacobi or CMV operators). Their proofs involve “classical” geometrical methods (localization with the help of a partition of unity).
We have to emphasize that many people working on pseudo-differential operators have been led to consider $`C^{}`$-algebras generated by such operators and to describe their quotients with respect to the ideal of compact operators: in fact, this is one of the most efficient ways to define the symbol of an operator (see \[CMS\] for example). Much more specific and relevant with respect to our goals is the work of H. O. Cordes (see \[Cor\] for a review). For example, the $`C^{}`$-algebra generated by a hypoelliptic operator and by the algebra of slowly oscillating functions and the computation of its quotient with respect to the compacts seem to have been considered for the first time in M. Taylor’s thesis (see \[Tay, Theorem 1\]). For more recent work on these lines, we refer to \[Nis\].
A rather different class of “$`C^{}`$-algebras of Hamiltonians” appears in the work of J. Bellissard on solid state physics \[Be1, Be2\]: he fixes a Hamiltonian $`H`$ and considers the $`C^{}`$-algebra generated by its translates. These algebras do not contain compact operators in general, so the techniques we use do not seem relevant in his setting. A more detailed discussion of the connection between the approach of Bellissard and ours can be found in \[GI4\].
The origin of our approach can be traced back to the algebraic treatment of the $`N`$-body problem from \[BG1, BG2\] (where the HVZ theorem and the Mourre estimate are proved in an abstract graded $`C^{}`$-algebra framework for a very general class of $`N`$-body Hamiltonians). The rôle of the crossed products was pointed out in \[GI2, GI3, GI4\] and a treatment of the $`N`$-body problem along these lines is presented in \[DG1, DG2, DG3\]. Various applications and extensions of the crossed product technique can be found in \[AMP, Man, MPR, Ric, Rod\] and references therein.
Our interest in localizations at infinity of a Hamiltonian was initially motivated by our desire to go beyond the $`N`$-body problem and to consider general (phase space) anisotropic systems \[GI1, Ift\]. Indeed, in the $`N`$-body case there is a lot of supplementary structure which makes the theory simple and beautiful (cf. Subsection 6.14), but this structure has no analogue in other types of anisotropy. We first found that the $`C^{}`$-algebra techniques are quite well adapted to the study of Hamiltonians with Klaus type potentials, see \[GI2, GI4\] and also Subsection 6.6 here for a treatment in the spirit of Theorem 1.1. We finally realized that the relation (1.7), which is the main point of the algebraic approach that we used, predicts in fact the description (1.4) of the essential spectrum of $`H`$ in terms of its localizations at infinity.
The paper \[HeM\] played an important rôle in our understanding of this fact. Indeed, B. Helffer and A. Mohamed prove there that the essential spectrum of a magnetic Hamiltonian $`(PA)^2+V`$ is the closure of the union of the spectra of some limit Schrödinger operators. Their proof is based on hypoellipticity techniques and the result is already interesting if the magnetic field is not present. The class of potentials they consider is quite large, but the function $`V`$ has to be bounded from below and to satisfy some regularity conditions. These assumptions imply that the limit operators have only polynomial electric and magnetic potentials, which is easily explained in our framework, see \[GI5, Proposition 3.13\].
1.5. Plan of the paper. Our purpose being to emphasize not only the power but also the simplicity of the $`C^{}`$-algebra techniques, we made an effort to make the paper essentially self-contained and easy to read by people working in the spectral theory of quantum Hamiltonians and with little background in $`C^{}`$-algebras.
We could have written a much shorter paper but which would have been accessible mostly to people with no interest in spectral theory. Instead, we have chosen to present in some detail most of the tools which are not standard among those interested in the subject. In particular we give in an Appendix a simple and self-contained proof of Landstad’s theorem (Theorem 3.7) which plays an important rôle in our arguments.
In Section 2 we introduce our notations and make a résumé of what we need concerning (ultra)filters and their relation with the characters of some abelian $`C^{}`$-algebras. In Section 3 we introduce crossed products in the version we need and we point out several useful consequences of Landstad’s theorem. This replaces the much more abstract arguments from \[GI2, GI4\], since we remain in a purely Hilbert space setting, but also gives stronger and more explicit results in applications. Section 4 is devoted to criteria of affiliation to the algebra $`C(X)`$, we show there that this algebra is much larger than one would think at first sight.
In Section 5 we prove our main result, Theorem 5.11. Finally, in Section 6 we consider three algebras of quantum Hamiltonians, those which seem the most interesting to us. The first one $`V(X)`$ is generated by slowly oscillating potentials and is the simplest non trivial algebra of Hamiltonians since it is defined by the property that if $`H`$ is affiliated to $`V(X)`$ then all its localizations at infinity are free Hamiltonians (i.e. functions of the momentum). The second one is the algebra associated to a *sparse set* and it is remarkable because the localizations at infinity of the Hamiltonians affiliated to it are two-body Hamiltonians and thus their essential spectrum has a quite interesting structure. The third one is, of course, the $`N`$-body algebra, or rather a more general and geometrically natural algebra that we call *Grassmann algebra*, an object of a remarkable simplicity, richness and beauty. The final subsections are devoted to some remarks of a different nature on the localizations at infinity of Hamiltonians of the form $`h(P)+v(Q)`$ with $`v(Q)`$ relatively bounded with respect to $`h(P)`$.
Acknowledgments: It is a pleasure to thank Barry Simon for stimulating correspondence and for sending us a copy of the paper \[Rab\]. We are also indebted to Françoise Piquard and George Skandalis for helpful discussions and to Steffen Roch for pointing out to us an erroneous assertion in the first version of this paper, cf. Remark 5.12. Finally, we are grateful to the referees, their comments and critics allowed us to eliminate several errors form the first version of this paper and to significantly improve the general presentation.
## 2 Preliminaries
In this section we describe our notations and recall facts needed in the rest of the paper.
2.1. If $`X`$ is a locally compact topological space then $`𝒞_{\mathrm{}}(X)`$ is the $`C^{}`$-algebra of continuous functions which have a limit at infinity and $`𝒞_0(X)`$ is the subalgebra of functions which converge to zero at infinity; thus $`𝒞_{\mathrm{}}(X)=+𝒞_0(X)`$. Let $`𝒞_\text{c}(X)`$ be the subalgebra of functions with compact support. If $`A`$ is a $`C^{}`$-algebra then we similarly define $`𝒞_0(X;A)`$ for example, which is also a $`C^{}`$-algebra. If $`X`$ is compact we set $`𝒞(X;A)=𝒞_0(X;A)`$ and $`𝒞(X)=𝒞(X;)`$, which does not conflict with the notation (2.3) because the continuous functions on $`X`$ are uniformly continuous. The characteristic function of a set $`SX`$ is denoted $`\text{1}_S`$.
In order to facilitate the reading of the paper we tried to respect as much as possible the following notational conventions. For abelian algebras (abstract as well as concrete ones, like function algebras) we use “mathcal” fonts, like $`𝒜,𝒞`$. For nonabelian algebras we use “mathscr” fonts, like $`A,C`$. Moreover, the crossed product of an abelian algebra $`𝒜`$ by the action of some group is denoted $`A`$. For other mathematical objects we use either greek letters or “mathcal” fonts with one exception: filters are often denoted by small gothic letters like $`𝔣,𝔤`$. However, ultrafilters are generally denoted $`\varkappa `$ because we think of them as “points at infinity” of the space $`X`$ whose points are denoted $`x`$.
2.2. If $``$ is a Hilbert space then $`B()`$ and $`K()`$ are the $`C^{}`$-algebras of bounded and compact operators on $``$ respectively. The resolvent set, the spectrum and the essential spectrum of an operator $`S`$ are denoted $`\rho (S)`$, $`\sigma (S)`$ and $`\sigma _{\mathrm{ess}}(S)`$ respectively.
By morphism between two $`C^{}`$-algebras we understand $``$-homomorphism. An ideal in a $`C^{}`$-algebra is assumed to be closed and two-sided.
An *observable* is a linear operator $`H:D(H)`$ such that $`HD(H)𝒦`$, where $`𝒦`$ is the closure of $`D(H)`$ in $``$, and such that $`H`$ when considered as operator in $`𝒦`$ is self-adjoint in the usual sense. A trivial observable which, however, is quite important, is the unique observable whose domain is equal to $`\{0\}`$; we shall denote it $`\mathrm{}`$. One has to think that $`H`$ is equal to $`\mathrm{}`$ on $`𝒦^{}`$ and for this reason we set $`\phi (H)=0`$ on $`𝒦^{}`$ if $`\phi 𝒞_0()`$. Note that we keep the notation $`(Hz)^1`$ for the resolvent of $`H`$ in $``$ but $`(Hz)^1=0`$ on $`𝒦^{}`$.
If $`CB()`$ is any $`C^{}`$-subalgebra then an observable $`H`$ is said to be affiliated to $`C`$ if $`(Hz)^1C`$ for some $`z\rho (H)`$. Then $`\phi (H)C`$ for all $`\phi 𝒞_0()`$.
It is theoretically much more convenient to define an observable affiliated to $`C`$ as a morphism $`H:𝒞_0()C`$ and then to set $`H(\phi )\phi (H)`$. We refer to \[GI4, p. 522–523\] for a résumé of what we need and also to \[DG3\] for comments on this notion which should not be confused with that introduced by S. Baaj and S. L. Woronowicz (in \[ABG, Sec. 8.1\] one can find a systematic presentation of this point of view).
We recall two definitions which make the transition from Theorem 1.1 to Theorem 1.2 trivial. The *spectrum* of the observable $`H`$ is the set
(2.1)
$$\sigma (H)=\{\lambda \phi 𝒞_0()\text{ and }\phi (\lambda )0\phi (H)0\}.$$
Let $`K=CK()`$, this is an ideal in $`C`$. Then the *essential spectrum* of $`H`$ is the set
(2.2)
$$\sigma _{\mathrm{ess}}(H)=\{\lambda \phi 𝒞_0()\text{ and }\phi (\lambda )0\phi (H)K\}.$$
We also note that any morphism $`\pi :CB`$ between two $`C^{}`$-algebras extends in a trivial way to a map between observables affiliated to $`C`$ to observables affiliated to $`B`$. Indeed, it suffices to define $`\pi (H)`$ by the condition $`\phi (\pi (H))=\pi (\phi (H))`$. For example, if $`\pi :CC/K`$ is the canonical morphism of $`C`$ onto the quotient algebra $`C/K`$, we have $`\sigma _{\mathrm{ess}}(H)=\sigma (\pi (H))`$.
Finally, we mention one more immediate consequence of the definition of an observable in terms of morphisms, cf. \[ABG, p. 370\]. We shall use the easily proven fact that $`\phi (H)`$ depends only on the restriction of $`\phi `$ to the closed real set $`\sigma (H)`$. Let $`\theta :\sigma (H)`$ be continuous and proper (i.e. $`|\theta (\lambda )|\mathrm{}`$ if $`|\lambda |\mathrm{}`$). Then the observable $`\theta (H)`$ is well defined by the rule $`\phi (\theta (H))=(\phi \theta )(H)`$ for $`\phi 𝒞_0()`$ (if $`H`$ is a self-adjoint operator then $`\theta (H)`$ is just the operator defined by the usual functional calculus). Clearly: *if $`H`$ is affiliated to $`C`$, the observable $`\theta (H)`$ is also affiliated to $`C`$*.
2.3. We describe now objects and notations from the harmonic analysis on groups. Everything we need can be found in \[Fol\] or \[FeD\]; see also \[Wei\].
Let $`X`$ be an abelian locally compact group (with the operation denoted additively) equipped with a Haar measure $`\text{d}x`$. We abbreviate $`B(X)=B(L^2(X)),K(X)=K(L^2(X))`$ and note that these are $`C^{}`$-algebras depending on $`X`$ and not on the choice of the Haar measure. Other such $`C^{}`$-algebras are $`L^{\mathrm{}}(X)`$, $`𝒞_{\mathrm{}}(X)`$, $`𝒞_0(X)`$ and the $`C^{}`$-algebra of bounded uniformly continuous functions on $`X`$, which plays the most important rôle in what follows:
(2.3)
$$𝒞(X)=\{\phi :X\phi \text{ is bounded and uniformly continuous }\}.$$
In order to avoid ambiguities, if $`\phi `$ is a measurable function on $`X`$ then we denote $`\phi (Q)`$ the operator of multiplication by $`\phi `$ in $`L^2(X)`$ (the symbol $`Q`$ has no operator meaning). By using this map we identify the algebra $`L^{\mathrm{}}(X)`$ and its $`C^{}`$-subalgebras with $`C^{}`$-subalgebras of $`B(X)`$, in particular we always embed
(2.4)
$$𝒞_0(X)𝒞_{\mathrm{}}(X)𝒞(X)B(X).$$
Note that the $`C^{}`$-algebra $`\mathrm{}_{\mathrm{}}(X)`$ of all bounded functions on $`X`$ cannot be embedded in $`B(X)`$ (neither can the $`C^{}`$-algebra $`(X)`$ of bounded Borel functions).
Let $`X^{}`$ be the set of characters of $`X`$ (continuous homomorphisms $`k:X`$ with $`|k(x)|=1`$) equipped with the locally compact group structure defined by the operation of multiplication and the topology of uniform convergence on compact sets. We denote the operation in $`X^{}`$ additively and its neutral element by $`0`$, as in \[Wei, ch.II, §5\] (this convention looks rather strange if $`X=^n`$, for example). If $`X`$ is a real finite dimensional vector space then $`X^{}`$ is identified with the vector space dual to $`X`$ as follows: let $`,:X\times X^{}`$ be the canonical bilinear map and take $`k(x)=\text{e}^{ix,k}`$. In fact, the field of real numbers can be replaced here by an arbitrary non-discrete locally compact field, see \[Fol, page 91\]) and \[Wei, ch.II, §5\]. We recall that the dual group $`(X^{})^{}`$ of $`X^{}`$ is identified with $`X`$, each $`xX`$ being seen as a character of $`X^{}`$ through the formula $`x(k)=k(x)`$.
The Fourier transform of $`uL^1(X)`$ is the function $`u\widehat{u}:X^{}`$ given by $`\widehat{u}(k)=_X\overline{k(x)}u(x)\text{d}x`$. We equip $`X^{}`$ with the unique Haar measure $`\text{d}k`$ such that $``$ induces a unitary map $`:L^2(X)L^2(X^{})`$. From $`^1=^{}`$ we get $`(^1v)(x)=_X^{}k(x)v(k)\text{d}k`$ for $`vL^2(X^{})`$. By taking into account the identification $`X^{}=X`$, the Fourier transform of $`\psi L^1(X^{})`$ and the Fourier inversion formula are
(2.5)
$$\widehat{\psi }(x)=_X^{}\overline{k(x)}\psi (k)\text{d}k\text{and}\psi (k)=_Xk(x)\widehat{\psi }(x)\text{d}x.$$
For each measurable $`\psi :X^{}`$ we define the operator $`\psi (P)`$ on $`L^2(X)`$ by $`\psi (P)=^{}M_\psi `$, where $`M_\psi `$ is the operator of multiplication by $`\psi `$ in $`L^2(X^{})`$. In particular, the restriction to $`L^{\mathrm{}}(X^{})`$ of the map $`\psi \psi (P)`$ is injective and gives us $`C^{}`$-subalgebras
(2.6)
$$𝒞_0(X^{})L^{\mathrm{}}(X^{})B(X).$$
Let $`\{U_x\}_{xX}`$ and $`\{V_k\}_{kX^{}}`$ be the strongly continuous unitary representations of $`X`$ and $`X^{}`$ in $`L^2(X)`$ defined by $`(U_xf)(y)=f(x+y)`$ and $`(V_kf)(y)=k(y)f(y)`$ respectively. Note that $`U_x`$ and $`V_k`$ satisfy the canonical commutation relations
(2.7)
$$U_xV_k=k(x)V_kU_x.$$
Observe that we have $`U_x=x(P)`$ if $`xX`$ is identified with the function $`kk(x)`$ and similarly $`V_k=k(Q)`$. Also, we have, cf. (2.5):
(2.8)
$$\psi (P)=_XU_x\widehat{\psi }(x)\text{d}x\text{if}\widehat{\psi }L^1(X).$$
2.4. We summarize here some facts we need concerning filters, cf. \[Bou, HiS, Sam\].
A *filter* on $`X`$ is a family $`𝔣`$ of subsets of $`X`$ which does not contain the empty set, is stable under finite intersections, and has the property: $`GF𝔣G𝔣`$ (the empty set is a filter!). If $`Y`$ is a topological space and $`\theta :XY`$ is any map, then $`lim_𝔣\theta =y`$ means that $`\theta ^1(V)𝔣`$ if $`V`$ is a neighborhood of $`y`$. We shall often write $`lim_{x𝔣}\theta (x)`$ instead of $`lim_𝔣\theta `$ for reasons which will become clear later on.
If $`𝔣,𝔤`$ are filters and $`𝔣𝔤`$ then $`𝔤`$ is said to be *finer* than $`𝔣`$. An *ultrafilter* is a maximal element in the set of all filters on $`X`$ for this order relation. If $`xX`$ then the family of sets which contain $`x`$ is the ultrafilter determined by $`x`$. A filter $`𝔣`$ is an ultrafilter if and only if for each $`AX`$ one has $`A𝔣`$ or $`A^cXA𝔣`$.
Ultrafilters are important because of the following property: if $`𝔣`$ is an ultrafilter and $`\theta :XY`$ is an arbitrary map with values in a compact space $`Y`$, then $`lim_𝔣\theta `$ exists. This fact will become clear after the explanations in Subsection 6.6.
The space $`\gamma X`$ of all ultrafilters on $`X`$ is a compact space for the topology defined as follows: the map $`𝔣\{\varkappa \gamma X\varkappa 𝔣\}`$ is a bijection from the set of all filters on $`X`$ onto the set of all closed subsets of $`\gamma X`$. Thus one should think that *a filter is a closed subset of* $`\gamma X`$. Another description of this topology will be given later on. The compact topological space $`\gamma X`$ is the *discrete Stone-Čech compactification of $`X`$* and it is characterized by the following universal property: *if $`Y`$ is a compact space then each map $`\theta :XY`$ has a unique extension to a continuous map $`\theta :\gamma XY`$*. Since this property is important for us, we shall further discuss it in Subsection 6.6, see page 6.6.
The set $`X`$ is identified with an open dense subset of $`\gamma X`$ (to $`xX`$ one associates the ultrafilter determined by $`x`$) and the topology induced by $`\gamma X`$ on $`X`$ is the discrete topology. However, the space $`\gamma XX`$ is much too large for our purposes, the only ultrafilters of interest to us belong to the compact subset of $`\gamma X`$ defined by
(2.9)
$$\delta X=\{\varkappa \varkappa \text{ is an ultrafilter finer than the Fréchet filter}\}.$$
We call *Fréchet filter* the filter consisting of the sets with relatively compact complement (this is not quite standard). This filter depends on the locally compact non compact topology given on $`X`$. In view of the standard meaning of the notation $`lim_x\mathrm{}`$ it is natural to denote by $`\mathrm{}`$ the Fréchet filter. As explained above, one should think of $`\mathrm{}`$ as a certain compact subset of $`\gamma X`$ and then we have in fact $`\mathrm{}=\delta X`$.
2.5. We shall explain now the relation between filters and characters of certain abelian $`C^{}`$-algebras. If $`𝒜`$ is such an algebra we denote $`\sigma (𝒜)`$ the space of characters of $`𝒜`$ (a character is a non zero morphism $`𝒜`$) equipped with the weak topology. This is a locally compact topological space which is compact if and only if $`𝒜`$ is unital.
Let $``$ be a unital abelian $`C^{}`$-algebra and let $`𝒜`$ be a $`C^{}`$-subalgebra which contains the unit of $``$. Then each character of $``$ restricts to a character of $`𝒜`$ and each character of $`𝒜`$ is obtained in this way. This gives a canonical map $`\pi :\sigma ()\sigma (𝒜)`$ which is continuous and surjective and if we define in $`\sigma ()`$ an equivalence relation $`\varkappa \chi `$ by the condition $`\varkappa (S)=\chi (S)S𝒜`$, the compact topological space $`\sigma (𝒜)`$ is just the quotient of $`\sigma ()`$ with respect to this relation.
In particular, *a map $`f:\sigma (𝒜)Y`$ is continuous if and only if $`f\pi :\sigma ()Y`$is continuous, where $`Y`$ is an arbitrary topological space*.
Let $`\mathrm{}_{\mathrm{}}(X)`$ be the $`C^{}`$-algebra of all bounded functions on $`X`$. Then the space of all characters of $`\mathrm{}_{\mathrm{}}(X)`$ can be identified with the space $`\gamma X`$ of all ultrafilters on $`X`$:
(2.10)
$$\gamma X=\sigma (\mathrm{}_{\mathrm{}}(X)).$$
Indeed, the map which associates to an ultrafilter $`𝔣`$ the character $`\phi lim_𝔣\phi `$ is a homeomorphism and the inverse map associates to the character $`\varkappa `$ the ultrafilter $`𝔣=\{FX\varkappa (\text{1}_F)=1\}`$. From now on we shall identify $`𝔣`$ and $`\varkappa `$, so *an ultrafilter is the same thing as a character of $`\mathrm{}_{\mathrm{}}(X)`$*, and we shall work with the interpretation which is most suited to the context. We also set $`\varkappa (F)=\varkappa (\text{1}_F)`$ for $`FX`$. Then
(2.11)
$$\delta X=\{\varkappa \gamma X\varkappa (K)=0KX\text{ compact }\}.$$
The algebras $`𝒜`$ that we consider are unital subalgebras of $`\mathrm{}_{\mathrm{}}(X)`$, thus their character spaces $`\sigma (𝒜)`$ are quotients of $`\gamma X`$. In other terms, we can view the characters of $`𝒜`$ as equivalence classes of ultrafilters: if $`\varkappa `$ is a character of $`𝒜`$, then there is an ultrafilter $`𝔣`$ such that $`\varkappa (\phi )=lim_𝔣\phi `$ for all $`\phi 𝒜`$, and in fact there are many such ultrafilters.
For the algebras which are of interest for us we always have
(2.12)
$$𝒞_{\mathrm{}}(X)𝒜𝒞(X)\mathrm{}_{\mathrm{}}(X)$$
Then $`X`$ is identified with an open dense subset of $`\sigma (𝒜)`$ and the topology induced by $`\sigma (𝒜)`$ on $`X`$ coincides with the initial topology, so $`\sigma (𝒜)`$ is a compactification of the locally compact space $`X`$. Thus
(2.13)
$$\delta (𝒜)=\sigma (𝒜)X=\{\varkappa \sigma (𝒜)\varkappa (\phi )=0\phi 𝒞_0(X)\}$$
is a compact subset of $`\sigma (𝒜)`$, the boundary of $`X`$ in the compactification $`\sigma (𝒜)`$. The *uniform compactification* $`\beta _\mathrm{u}X`$ of $`X`$ is defined by the largest algebra $`𝒞(X)`$:
(2.14)
$$\beta _\mathrm{u}X=\sigma (𝒞(X)),\delta _\mathrm{u}X=\beta _\mathrm{u}XX=\delta (𝒞(X)).$$
Later on we shall explicitly describe the equivalence relation in $`\gamma X`$ which defines $`\beta _\mathrm{u}`$.
We are interested only in the boundary $`\delta (𝒜)`$ of $`X`$ in $`\sigma (𝒜)`$. We show now that this is a quotient of $`\delta X`$.
###### Lemma 2.1
Let $`𝔣`$ be an ultrafilter on $`X`$ and let $`\varkappa `$ be the character of $`𝒜`$ defined by $`\varkappa (\phi )=lim_𝔣\phi `$. Then $`\varkappa \delta (𝒜)`$ if and only if $`𝔣\delta X`$.
Proof: If $`𝔣`$ is an ultrafilter and $`YX`$ then there are only two possibilities: either $`Y𝔣`$, and then $`XY𝔣`$ hence $`YZ=\mathrm{}`$ for all $`Z𝔣`$, or $`Y𝔣`$, and then the sets $`YZ`$ with $`Z𝔣`$ form an ultrafilter on $`Y`$. If $`𝔣`$ is not finer than the Fréchet filter then there is a set with compact complement $`Y`$ which does not belong to $`𝔣`$, and so $`Y𝔣`$. Since any ultrafilter on a compact set is convergent, we see that there is $`yY`$ such that $`𝔣`$ contains the filter of neighborhoods of $`y`$. But then clearly $`lim_𝔣\phi =\phi (y)`$ for any continuous function $`\phi `$, hence the character $`\varkappa (\phi )=lim_𝔣\phi `$ is just $`y`$ and does not belong to $`\delta (𝒜)`$. On the other hand, if $`𝔣\delta X`$ then clearly $`\varkappa \delta (𝒜)`$.
Thus the *characters $`\varkappa \delta (𝒜)`$ are equivalence classes of ultrafilters $`𝔣\delta X`$*. In general, we do not distinguish between a character and the elements of the equivalence class of ultrafilters which define it. However, when needed for the clarity of the argument, we shall use the map $`\delta `$ which sends an element into its equivalence class. More precisely, from (2.12) we see that there are canonical surjections
(2.15)
$$\delta X\delta _uX\delta (𝒜)\{\mathrm{}\}$$
and all of them (and their compositions) will be denoted $`\delta `$. Here $`\mathrm{}`$ is the Fréchet filter and we have $`\delta (𝒞_{\mathrm{}}(X))=\{\mathrm{}\}`$.
2.6. The space $`\beta _\mathrm{u}X`$ is the quotient of $`\gamma X`$ given by an equivalence relation that we describe now (see \[Sam, p. 121\]). If $`𝔣`$ is a filter then its *envelope* is the filter $`𝔣^{}`$ generated by the sets $`F+V`$ where $`F𝔣`$ and $`V`$ belongs to the filter of neighborhoods of the origin (observe that the sets $`F+V`$, with $`V`$ an open neighborhood of the origin, are open and form a basis of $`𝔣^{}`$). Note that $`𝔣𝔣^{}`$ and $`(𝔣^{})^{}=𝔣^{}`$. Two filters are *u-equivalent* (uniformly equivalent) if they have the same envelope.
*The quotient of $`\gamma X`$ with respect to this relation is $`\beta _\mathrm{u}X`$*. We shall give a complete proof of this assertion since in \[Sam\] the $`C^{}`$-algebra point of view is not explicitly considered. The following simple fact will be useful for other purposes too.
###### Lemma 2.2
Let $`\phi :X`$ be uniformly continuous and let $`𝔣`$ be a filter on $`X`$.
(1) $`lim_𝔣\phi `$ exists if and only if $`lim_𝔣^{}\phi `$ exists and in this case they are equal.
(2) If $`lim_{x𝔣}\phi (x+y)\xi (y)`$ exists for each $`yX`$ then the limit exists locally uniformly in $`y`$ and $`\xi `$ is a uniformly continuous function.
Proof: To prove (1) it suffices to show that $`lim_𝔣^{}\phi =0`$ if $`lim_𝔣\phi =0`$. For $`\epsilon >0`$ let $`F_\epsilon `$ be the set of points where $`|\phi (x)|<\epsilon `$. We have $`F_\epsilon 𝔣`$ and if we choose a neighborhood $`V`$ of the origin such that $`|\phi (x)\phi (y)|<\epsilon `$ if $`xyV`$, then for $`xF_\epsilon +V`$ we have $`|\phi (x)|<2\epsilon `$ hence $`F_{2\epsilon }𝔣^{}`$.
Now we prove (2). Set $`\omega _V(\phi )=sup_{yzV}|\phi (y)\phi (z)|`$ if $`V`$ is a neighborhood of the origin. Then $`\phi `$ is uniformly continuous if and only if for each $`\epsilon >0`$ there is $`V`$ such that $`\omega _V(\phi )<\epsilon `$. Clearly $`\omega _V(\xi )\omega _V(\phi )`$, so $`\xi `$ is uniformly continuous. Now let $`V`$ be open and let $`K`$ be a compact set. Then $`K`$ is covered by the open sets $`x+V`$, $`xK`$, hence there is $`ZK`$ finite such that $`K_{zZ}(z+V)`$. Thus for each $`yK`$ there is $`zZ`$ such that $`yz+V`$ and then $`|\phi (x+y)\phi (x+z)|\omega _V(\phi )`$ for all $`xX`$. Then we have:
$`|\phi (x+y)\xi (y)|`$ $``$ $`|\phi (x+y)\phi (x+z)|+|\phi (x+z)\xi (z)|+|\xi (z)\xi (y)|`$
$``$ $`\omega _V(\phi )+|\phi (x+z)\xi (z)|+\omega _V(\xi )`$
$``$ $`2\omega _V(\phi )+|\phi (x+z)\xi (z)|.`$
We choose $`V`$ such that $`\omega _V(\phi )\epsilon /3`$ and then we fix $`Z`$ as above. Since $`Z`$ is finite, there is $`F𝔣`$ such that $`|\phi (x+z)\xi (z)|\epsilon /3`$ for all $`xF`$ and $`zZ`$. Finally, we get $`|\phi (x+y)\xi (y)|\epsilon `$ for all $`xF`$ and $`yK`$.
###### Lemma 2.3
Assume $`𝔣=𝔣^{}`$ where $`𝔣`$ is a filter on $`X`$. Then for each $`F𝔣`$ there is an open subset $`G𝔣`$ of $`F`$ and a function $`\theta 𝒞(X)`$ such that $`\theta =1`$ on $`G`$ and $`\theta =0`$ on $`F^\mathrm{c}XF`$.
Proof: Note first that the open sets from $`𝔣`$ form a basis of $`𝔣`$. Clearly there is an open $`G𝔣`$ and an open, relatively compact neighborhood of the origin $`U`$ such that $`G+(UU)F`$, so denoting $`A=GU`$ we shall have $`A+UF`$. We then set $`\theta =|U|^1\text{1}_A\text{1}_U`$, so for each $`xX`$ we have $`\theta (x)=|U|^1|A(xU)|`$. For $`xG`$, $`xUA`$ thus $`\theta (x)=|U|^1|xU|=1`$, and for $`xA+U`$, $`A(xU)=\mathrm{}`$ hence $`\theta (x)=0`$. But $`A+UF`$, thus $`\theta =0`$ on $`F^\mathrm{c}`$ too. Finally, from
$$uv_L^{\mathrm{}}u_{L^1}v_L^{\mathrm{}}\text{and}x.(uv)uv_L^{\mathrm{}}x.uu_{L^1}v_L^{\mathrm{}},$$
where $`(x.u)(y)=u(y+x)`$, we get $`L^1(X)L^{\mathrm{}}(X)𝒞(X)`$, hence $`\theta 𝒞(X)`$.
###### Proposition 2.4
Let $`\varkappa ,\chi `$ be ultrafilters on $`X`$. Then $`\varkappa ^{}=\chi ^{}`$ if and only if $`\varkappa (\phi )=\chi (\phi )`$ for each $`\phi 𝒞(X)`$.
Proof: The “only if” part follows from $`\varkappa (\phi )=lim_\varkappa \phi =lim_\varkappa ^{}\phi =lim_\chi ^{}\phi =lim_\chi \phi =\chi (\phi ),`$ the second and the fourth equality being consequences of Lemma 2.2. Conversely, let $`\varkappa ^{}\chi ^{}`$. Then there is $`F\varkappa ^{}`$ such that $`F\chi ^{}\chi `$, hence $`F^\mathrm{c}\chi `$ because $`\chi `$ is an ultrafilter. Let now $`G`$ and $`\theta `$ be as in Lemma 2.3. Since $`G\varkappa ^{}\varkappa `$ we have $`\varkappa (\text{1}_G)=1`$, thus $`\varkappa (\text{1}_{G^\mathrm{c}})=0`$. Hence $`\varkappa (\theta )=\varkappa (\theta \text{1}_G)+\varkappa (\theta \text{1}_{G^\mathrm{c}})=1+\varkappa (\theta )\varkappa (\text{1}_{G^\mathrm{c}})=1`$. On the other hand, $`F^\mathrm{c}\chi `$ implies $`\chi (\text{1}_{F^\mathrm{c}})=1`$ and $`\theta \text{1}_{F^\mathrm{c}}=0`$, thus $`0=\chi (\theta \text{1}_{F^\mathrm{c}})=\chi (\theta )\chi (\text{1}_{F^\mathrm{c}})=\chi (\theta )`$.
## 3 Crossed products
In this section we recall some facts concerning crossed products and point out some properties important for our later arguments. A locally compact non compact abelian group $`X`$ is fixed in what follows.
We shall say that a $`C^{}`$-algebra $`𝒜`$ is an $`X`$-algebra if a homomorphism $`\alpha :x\alpha _x`$ of $`X`$ into the group of automorphisms of $`𝒜`$ is given, such that for each $`A𝒜`$ the map $`x\alpha _x(A)`$ is norm continuous <sup>2</sup><sup>2</sup>2 The terminology “$`C^{}`$-dynamical system” used by some $`C^{}`$-algebra theorists seems to us extremely confusing in our context, even if $`X`$ is $``$ or $``$, so we shall not use it.. An *$`X`$-subalgebra* of $`𝒜`$ is a $`C^{}`$-subalgebra that is left invariant by all the automorphisms $`\alpha _x`$. An *$`X`$-ideal* is an ideal stable under the $`\alpha _x`$. If $`(𝒜,\alpha )`$ and $`(,\beta )`$ are two $`X`$-algebras, a morphism $`\varphi :𝒜`$ is called $`X`$-morphism if $`\varphi [\alpha _x(A)]=\beta _x[\varphi (A)]`$ for all $`xX`$ and $`A𝒜`$.
We shall not need the abstract definition of the crossed product $`𝒜X`$ of an $`X`$-algebra $`𝒜`$ by the action of $`X`$. We mention only that $`𝒜X`$ is a $`C^{}`$-algebra uniquely defined modulo a canonical isomorphism by a certain universal property (see \[Rae\] for example) and that the correspondence $`𝒜𝒜X`$ has certain functorial properties (see \[GI5\]) which play an important rôle in \[GI4\] but will not be used here. On the other hand, the following concrete realization of $`𝒜X`$ for certain $`𝒜`$ will be important.
There is a natural action of $`X`$ on $`L^{\mathrm{}}(X)`$ by translations $`(\tau _x\phi )(y)=\phi (y+x)`$ and it is clear that $`x\tau _x\phi L^{\mathrm{}}(X)`$ is norm continuous if and only if $`\phi 𝒞(X)`$. Thus $`𝒞(X)`$ becomes an $`X`$-algebra and we will be interested only in crossed products $`𝒜X`$ with $`𝒜`$ an $`X`$-subalgebra of $`𝒞(X)`$, i.e. a $`C^{}`$-subalgebra stable under translations. In many cases we shall slightly simplify the writing and set $`x.\phi =\tau _x\phi `$. Note that if $`\phi 𝒞(X)L^2(X)`$ we have $`x.\phi =U_x\phi `$ but $`(x.\phi )(Q)=U_x\phi (Q)U_x^{}`$. More generally, we shall use the notations:
(3.1)
$$xX,TB(X)x.T\tau _x(T)=U_xTU_x^{}.$$
The next definition describes $`𝒜X`$ in what we could call the *pseudo-differential operator representation*, or $`\mathrm{\Psi }`$DO-representation.
###### Definition 3.1
If $`𝒜`$ is an $`X`$-subalgebra of $`𝒞(X)`$, the *crossed product* $`𝒜XA`$ is the norm closed linear subspace of $`B(X)`$ generated by the operators of the form $`\phi (Q)\psi (P)`$ with $`\phi 𝒜`$ and $`\psi 𝒞_0(X^{})`$.
The fact that $`A`$ is a $`C^{}`$-algebra follows from:
###### Lemma 3.2
If $`\phi 𝒞(X)`$ and $`\psi 𝒞_0(X^{})`$ then for each number $`\epsilon >0`$ there are elements $`x_1,\mathrm{},x_nX`$ and functions $`\psi _1,\mathrm{},\psi _n𝒞_0(X^{})`$ such that:
(3.2)
$$\psi (P)\phi (Q)_k\phi (Q+x_k)\psi _k(P)<\epsilon .$$
For the proof, first approximate $`\psi `$ by functions such that $`\widehat{\psi }L^1(X)`$ and then adapt the proof of \[DG1, Lemma 2.1\]. We mention two results which explain why we think of $`A`$ as a $`C^{}`$-algebra of quantum Hamiltonians. The first one is \[GI4, Proposition 4.1\].
###### Proposition 3.3
Let $`𝒜`$ be an $`X`$-subalgebra of $`𝒞(X)`$ which contains the constants. Let $`h:X^{}`$ be a continuous non-constant function such that $`lim_k\mathrm{}|h(k)|=\mathrm{}`$. Then $`𝒜X`$ is the $`C^{}`$-algebra generated <sup>3</sup><sup>3</sup>3 If $`𝒮`$ is a family of self-adjoint operators then the $`C^{}`$-algebra generated by $`𝒮`$ is the smallest $`C^{}`$-algebra of operators on $``$ to which is affiliated each $`H𝒮`$. by the self-adjoint operators of the form $`h(P+k)+v(Q)`$, with $`kX^{}`$ and $`v𝒜`$ real.
The second one is \[DG1, Corollary 2.4\]. Here we assume $`X=^n`$ and denote $`𝒜^{\mathrm{}}`$ the set of functions in $`𝒜`$ such that all their derivatives exist and belong to $`𝒜`$.
###### Proposition 3.4
Let $`h`$ be a real elliptic polynomial of order $`m`$ on $`X`$ and let $`𝒜`$ be as in Proposition 3.3. Then $`𝒜X`$ is the $`C^{}`$-algebra generated by the self-adjoint operators of the form $`h(P)+S`$, where $`S`$ runs over the set of symmetric differential operators of order $`<m`$ with coefficients in $`𝒜^{\mathrm{}}`$.
###### Examples 3.5
We shall point out now the simplest crossed products. The smallest crossed product $`\{0\}=\{0\}X`$ is, of course, of no interest.
* The largest crossed product is $`C(X)=𝒞(X)X`$, see Theorem 3.10.
* The $`C_0`$ functions of momentum: $`𝒞_0(X^{})=X`$.
* The algebra of compact operators: $`K(X)=𝒞_0(X)X`$.
* The *two-body algebra*: $`T(X):=𝒞_{\mathrm{}}(X)X=𝒞_0(X^{})+K(X)`$.
The name of the fourth algebra is justified by Propositions 3.3 and 3.4. Indeed, if $`X=^n`$ then $`T(X)`$ is the $`C^{}`$-algebra generated by the self-adjoint operators of the form $`(P+k)^2+v(Q)`$ with $`kX`$ and $`v𝒞_\text{c}^{\mathrm{}}(X)`$ is real, or by those of the form $`\mathrm{\Delta }+_{j=1}^na_j_j+a_0`$ where $`a_j`$ are $`C^{\mathrm{}}`$ functions constant outside a compact.
###### Remark 3.6
Note that the only abelian crossed products are $`\{0\}`$ and $`𝒞_0(X^{})`$.
We have defined a map $`𝒜𝒜X`$ from the set of all $`X`$-subalgebras of $`𝒞(X)`$ into the set of $`C^{}`$-subalgebras of $`B(X)`$ which is obviously increasing. The following theorem, which is an immediate consequence of a more general abstract result due to M.B. Landstad, cf. \[Lan, Theorem 4\] or \[Ped\], says that this map is injective and describes its range.
###### Theorem 3.7
A $`C^{}`$-subalgebra $`A`$ of $`B(X)`$ is a crossed product if and only if for each $`AA`$ the following two conditions are satisfied:
* If $`kX^{}`$ then $`V_k^{}AV_kA`$ and $`lim_{k0}V_k^{}AV_kA=0`$,
* If $`xX`$ then $`U_xAA`$ and $`lim_{x0}(U_x1)A=0`$.
In this case, there is a unique $`X`$-subalgebra $`𝒜𝒞(X)`$ such that $`A=𝒜X`$, and this algebra is given by
(3.3)
$$𝒜=A_{\mathrm{}}:=\{\phi 𝒞(X)\phi (Q)^{()}\psi (P)A,\psi 𝒞_0(X^{})\}.$$
Note that, since $`A`$ is stable under taking adjoints, if we replace $`U_xA`$ by $`AU_x`$ and $`(U_x1)A`$ by $`A(U_x1)`$ in the second condition above we get an equivalent condition. If each element $`A`$ of a $`C^{}`$-subalgebra $`AB(X)`$ verifies the two conditions of the theorem, we shall say that $`A`$ *satisfies Landstad’s conditions*.
The following reformulation of the second Landstad condition is useful.
###### Lemma 3.8
If $`TB(X)`$ then the next three assertions are equivalent:
* $`lim_{x0}(U_x1)T=0`$,
* $`T=\psi (P)T_0`$ for some $`\psi 𝒞_0(X^{})`$ and $`T_0B(X)`$,
* $`\epsilon >0`$ $`FX^{}`$ with $`X^{}F`$ compact and $`\text{1}_F(P)T<\epsilon `$.
Proof: It suffices to consider only the first two conditions. If $`T=\psi (P)T_0`$ then
$$(U_x1)T(U_x1)\psi (P)T_0T_0\underset{k}{sup}|(k(x)1)\psi (k)|0\text{as }x0.$$
To prove the converse assertion, let $`B_0=\{TBlim_{x0}(U_x1)T=0\}`$. This is clearly a closed subspace of $`B`$ such that $`\psi (P)B_0B_0`$ if $`\psi 𝒞_0(X^{})`$. By taking $`\widehat{\psi }(k)=|K|^1\text{1}_K`$ in (2.8), where $`K`$ runs over the family of compact neighborhoods of the origin in $`X^{}`$, we easily see that each $`TB_0`$ is a norm limit of operators of the form $`\psi (P)T`$. Now the Cohen-Hewitt factorization theorem \[FeD, Th. V.9.2\] shows that each $`TB_0`$ can be written as $`T=\psi (P)T_0`$ with $`\psi 𝒞_0(X^{})`$ and $`T_0B_0`$.
###### Corollary 3.9
If $`A`$ is a crossed product then each $`AA`$ can be factorized as $`A=A_1\psi _1(P)=\psi _2(P)A_2`$ with $`A_iA`$ and $`\psi _i𝒞_0(X^{})`$. In particular, if $`AA`$ and $`\psi `$ is a bounded continuous function on $`X^{}`$ then $`A\psi (P)`$ and $`\psi (P)A`$ belong to $`A`$.
Theorem 3.7 allows us to give an intrinsic description of some crossed products. By “intrinsic” we mean a description which makes no reference to the crossed product operation. Examples may be found in Section 6, here we give the description of the largest crossed product $`C(X)`$ which makes the connection with the definition (1.1).
###### Theorem 3.10
The crossed product $`C(X)=𝒞(X)X`$ is given by (1.1).
For the proof, it suffices to note that the right hand side of (1.1) is a $`C^{}`$-algebra and to apply Theorem 3.7. It is useful to view the last condition in (1.1) from the perspective of Lemma 3.8: this gives a precise meaning to the fact that the operators from $`C(X)`$ tend to zero as $`P\mathrm{}`$.
###### Remark 3.11
If $`X=^n`$ we see that $`C(X)`$ is the norm closed linear subspace of $`B(X)`$ generated by the operators $`\phi (Q)\psi (P)`$ with $`\phi `$ in the space of $`C^{\mathrm{}}`$ functions which are bounded together with all their derivatives and $`\psi `$ in the space of $`C^{\mathrm{}}`$ functions with compact support. So $`C(X)`$ is generated by a rather restricted class of pseudo-differential operators. In particular, $`C(X)`$ is the norm closure of the set of pseudo-differential operators with symbols of class $`S^m`$ if $`m<0`$ (see \[Hör, Definition 18.1.1\] and use \[Hör, Theorem 18.1.6\]). From Proposition 3.4 it also follows that $`C(X)`$ is generated by a rather small class of elliptic operators.
As a consequence, we get an intrinsic description of the algebras of quantum Hamiltonians, in the sense of Definition 1.13.
###### Proposition 3.12
A $`C^{}`$-subalgebra $`AB(X)`$ is a $`C^{}`$-algebra of quantum Hamiltonians if and only if $`AT(X)`$ and
* $`xX,kX^{},AAV_k^{}AV_k\text{ and }U_xA\text{ belong to }A,`$
* $`lim_{k0}[A,V_k]=lim_{x0}(U_x1)A=0.`$
###### Remark 3.13
Observe that the classical Riesz-Kolmogorov compactness criterion
$`K(X)`$ $`=`$ $`\{TB(X)\underset{k0}{lim}(V_k1)T=0\text{ and }\underset{x0}{lim}(U_x1)T=0\}`$
$`=`$ $`\{TB(X)T=\phi (Q)S=\psi (P)R\text{ with }\phi 𝒞_0(X),\psi 𝒞_0(X^{})`$
$`\text{and }S,RB(X)\}`$
is also an intrinsic characterization of a crossed product and follows easily from Theorem 3.7 and Lemma 3.8 together with a similar fact with the group $`U_x`$ replaced by $`V_k`$. In more intuitive terms, the compact operators are characterized by the fact that they vanish when $`P\mathrm{}`$ and $`Q\mathrm{}`$.
Now we show that the set of $`C^{}`$-subalgebras of $`B(X)`$ which are crossed products is stable under arbitrary intersections and that the $`C^{}`$-algebra generated by an arbitrary family of crossed products is again a crossed product. We denote by $`C^{}\left(_\lambda _\lambda \right)`$ the $`C^{}`$-subalgebra generated by a family of $`C^{}`$-subalgebras $`_\lambda `$.
###### Theorem 3.14
If $`(𝒜_\lambda )`$ is an arbitrary family of $`X`$-subalgebras of $`𝒞(X)`$ then $`_\lambda 𝒜_\lambda `$ and $`C^{}\left(_\lambda 𝒜_\lambda \right)`$ are $`X`$-subalgebras and:
(3.4) $`_\lambda (𝒜_\lambda X)`$ $`=`$ $`(_\lambda 𝒜_\lambda )X,`$
(3.5) $`C^{}\left(_\lambda (𝒜_\lambda X)\right)`$ $`=`$ $`C^{}\left(_\lambda 𝒜_\lambda \right)X.`$
Proof: The fact that $`_\lambda 𝒜_\lambda `$ and $`C^{}\left(_\lambda 𝒜_\lambda \right)`$ are $`X`$-subalgebras is easy to prove and the inclusions $``$ in (3.4) and $``$ in (3.5) are obvious. The proof of $``$ in (3.5) is elementary. Indeed, it suffices to show that $`\phi (Q)\psi (P)`$ belongs to the left hand side of (3.5) if $`\phi C^{}\left(_\lambda 𝒜_\lambda \right)`$. Then we may assume that $`\phi =\phi _L=_{\lambda L}\phi _\lambda `$ with $`\phi _\lambda 𝒜_\lambda `$ and $`L`$ a finite set. Let $`\lambda L`$ and $`M=L\{\lambda \}`$. Then Corollary 3.9 applied to $`\phi _\lambda (Q)\psi (P)𝒜_\lambda X`$ gives
$$\phi _L(Q)\psi (P)=\phi _M(Q)\phi _\lambda (Q)\psi (P)=\phi _M(Q)\psi _\lambda (P)A_\lambda $$
for some $`\psi _\lambda 𝒞_0(X^{})`$ and $`A_\lambda 𝒜_\lambda X`$. Repeating the argument with $`\phi _L`$ replaced by $`\phi _M`$ we see that $`\phi _L(Q)\psi (P)`$ can be written as a product of elements of $`𝒜_\lambda X`$ with $`\lambda L`$. This proves (3.5).
The inclusion $``$ in (3.4) is a deeper fact, it depends on Theorem 3.7. Let $`A_\lambda =𝒜_\lambda X`$ and $`A=_\lambda A_\lambda `$. It is easy to check that $`A`$ satisfies the two conditions of Theorem 3.7, so $`A=𝒜X`$ where $`𝒜`$ is defined by (3.3). If $`\phi 𝒞(X)`$ has the property $`\phi (Q)^{()}\psi (P)A`$ for all $`\psi 𝒞_0(X^{})`$ then we also have $`\phi (Q)^{()}\psi (P)A_\lambda `$ for all such $`\psi `$, hence $`\phi (A_\lambda )_{\mathrm{}}=𝒜_\lambda `$ for each $`\lambda `$. Thus $`\phi _\lambda 𝒜_\lambda `$, hence $`𝒜_\lambda 𝒜_\lambda `$.
###### Proposition 3.15
If $`𝒜,𝒥`$ are $`X`$-subalgebras then $`𝒥`$ is an ideal of $`𝒜`$ if and only if $`𝒥X`$ is an ideal of $`𝒜X`$.
Proof: The fact that “$`𝒥𝒜`$ ideal $`𝒥X𝒜X`$ ideal” follows easily from Lemma 3.2. For the converse it suffices to show that if $`J,A`$ are crossed products and if $`J`$ is an ideal of $`A`$, then $`J_{\mathrm{}}`$ is an ideal of $`A_{\mathrm{}}`$. Let $`\xi J_{\mathrm{}}`$ and $`\phi A_{\mathrm{}}`$, then by Corollary 3.9, for each $`\psi 𝒞_0(X^{})`$ we can factorize $`\phi (Q)\psi (P)=\psi _0(P)S`$ for some $`\psi _0𝒞_0(X^{})`$ and $`SA`$. Thus $`(\xi \phi )(Q)\psi (P)=\xi (Q)\psi _0(P)S𝒥`$ because $`\xi (Q)\psi _0(P)J`$ and $`J`$ is an ideal of $`A`$, hence $`\xi \phi J_{\mathrm{}}`$.
###### Proposition 3.16
Assume that $`𝒜,,𝒥`$ are $`X`$-subalgebras of $`𝒞(X)`$ such that $`𝒜=+𝒥`$ and that $`𝒥`$ is an ideal in $`𝒜`$. Then $`𝒥X`$ is an ideal in $`𝒜X`$ and $`𝒜X=X+𝒥X`$. If $`𝒜=+𝒥`$ is a linear direct sum, then $`𝒜X=X+𝒥X`$ is a linear direct sum.
Proof: We know that $`𝒥X`$ is an ideal in $`𝒜X`$ and that $`X𝒜X`$ is a $`C^{}`$-subalgebra. From \[Dix, Corollary 1.8.4\] we see that $`X+𝒥X`$ is closed in $`𝒜X`$, and since it is clearly dense in $`𝒜X`$, we have $`𝒜X=X+𝒥X`$. Finally,
$$(X)(𝒥X)=\left(𝒥\right)X$$
because of (3.4), and this is $`\{0\}`$ if $`𝒥=\{0\}`$. .
We mention a fact which is useful in the explicit computations of $`A_{\mathrm{}}`$.
###### Remark 3.17
It is clear that in (3.3) it suffices to consider only $`\psi 𝒞_\text{c}(X^{})`$. Since, by Corollary 3.9, a crossed product is a $`𝒞_0(X^{})`$-bimodule, we get the following simpler description of $`𝒜`$: *if there is $`\xi 𝒞_0(X^{})`$ such that $`\xi (k)0`$ for all $`kX^{}`$, then*
(3.6)
$$𝒜=\{\phi 𝒞(X)\phi (Q)^{()}\xi (P)A\}.$$
Such a $`\xi `$ exists if and only if $`X^{}`$ is $`\sigma `$-compact (i.e. a countable union of compact sets).
###### Remark 3.18
The following comment on the first Landstad condition is of some interest, although it does not play any rôle in our arguments. Let $`C^\mathrm{u}(Q)`$ be the set of $`SB(X)`$ which verify the first Landstad condition; this is clearly a $`C^{}`$-algebra. Let us say that an operator $`SB(X)`$ is of *finite range* (not rank!) if there is a compact neighborhood $`\mathrm{\Lambda }`$ of the origin such that $`S\text{1}_K(Q)=\text{1}_{K+\mathrm{\Lambda }}(Q)S\text{1}_K(Q)`$ for any Borel set $`K`$. Clearly, the set of finite range operators is a $``$-subalgebra of $`B(X)`$ and it can be shown that the set of finite range operators which belong to $`C^\mathrm{u}(Q)`$ is dense in $`C^\mathrm{u}(Q)`$. Moreover, under quite general conditions on $`X`$ it can be shown that a finite range operator belongs to $`C^\mathrm{u}(Q)`$ (this is probably always true). Thus, if $`X=^n`$ or if $`X`$ is a discrete group for example, then $`C^\mathrm{u}(Q)`$ is exactly the norm closure of the set of finite range operators. These questions are treated in \[GG2, Propositions 4.11 and 4.12\].
## 4 Affiliation to C $`(X)`$
Theorem 1.2 shows that the essential spectrum of the operators affiliated to $`C(X)`$ is determined by their localizations at infinity, so it is important to show that the class of operators affiliated to $`C(X)`$ is large. We show in this section that this is indeed the case: singular perturbations of hypoelliptic self-adjoint pseudo-differential operators are affiliated to $`C(X)`$. If one thinks of $`C(X)`$ as the $`C^{}`$-algebra generated by the operators of the form $`\phi (Q)\psi (P)`$ with $`\phi 𝒞(X),\psi 𝒞_0(X)`$, this is far from obvious.
In the rest of the section we fix a finite dimensional Hilbert space $`𝐄`$, we set $`=L^2(X;𝐄)`$ and define $`C=C(X)`$ as in (1.1). Since the adjoint space <sup>2</sup><sup>2</sup>2 The adjoint space (space of antilinear continuous forms) of a Hilbert space $`𝒢`$ is denoted $`𝒢^{}`$ and if $`u𝒢`$ and $`v𝒢^{}`$ then we set $`v(u)=u,v`$. $`^{}`$ is identified with $``$ by using the Riesz isomorphism, if $`𝒢`$ is a Hilbert space with $`𝒢`$ continuously and densely then we get a similar embedding $`𝒢^{}`$.
Let $`H`$ be a self-adjoint operator on $``$ and let $`z\rho (H)`$. As we saw in (1.3), $`H`$ is affiliated to $`C`$ if and only if
(4.1)
$$\underset{x0}{lim}(U_x1)(Hz)^1=0\text{ and }\underset{k0}{lim}[V_k,(Hz)^1]=0.$$
In the next subsection we make an abstract analysis of these relations and in Subsection 4.9 we give concrete examples.
4.1. A function $`\theta :X^{}`$ such that $`lim_k\mathrm{}\theta (k)=+\mathrm{}`$ will be called divergent. Lemma 3.8 and an interpolation argument give:
###### Lemma 4.1
The first condition in (4.1) is fulfilled if and only there are $`s>0`$ and a continuous divergent function $`\theta `$ such that $`D(|H|^s)D(\theta (P))`$. And then this property holds for all real numbers $`s>0`$.
Let $`S(𝐄)`$ be the space of symmetric operators on $`𝐄`$. If $`h:X^{}S(𝐄)`$ is Borel, then $`h(P)`$ is the self-adjoint operator on $``$ such that $`h(P)^{}`$ is the operator of multiplication by $`h`$ in $`L^2(X^{};𝐄)`$. If $`lim_k\mathrm{}\mathrm{dist}(0,\sigma (h(k)))=\mathrm{}`$ then we write $`lim_k\mathrm{}h(k)=\mathrm{}`$. This property is equivalent to $`lim_k\mathrm{}(h(k)+i)^1=0`$ and implies $`lim_k\mathrm{}\phi (h(k))=0`$ for all $`\phi 𝒞_0()`$. If $`𝐄=`$ this means $`lim_k\mathrm{}|h(k)|=\mathrm{}`$.
###### Corollary 4.2
If $`h:X^{}S(𝐄)`$ is a continuous function on $`X^{}`$ then $`h(P)`$ is affiliated to $`C`$ if and only if $`lim_k\mathrm{}h(k)=\mathrm{}`$.
In particular, if $`X=^2`$ then the operator $`H=_1^2_2^2`$ is not affiliated to $`C`$. A second interesting operator not affiliated to $`C`$ is $`H=(_1+ix_2)^2+(_2+ix_1)^2`$.
We now give the simplest affiliation criterion.
###### Proposition 4.3
Assume that $`H_0`$ is a self-adjoint operator affiliated to $`C`$ and that $`V`$ is a bounded symmetric operator such that $`lim_{k0}[V_k,V]=0`$. Then $`H=H_0+V`$ is a self-adjoint operator affiliated to $`C`$.
Proof: Let $`R=(H+i)^1`$ and $`R_0=(H_0+i)^1`$. Since $`H`$ and $`H_0`$ have the same domain and $`R[1+VR_0]=R_0`$, the operator $`1+VR_0`$ is invertible. On the other hand, $`1+VR_0`$ clearly satisfies the second condition in (4.1), hence its inverse verifies it too. From $`R=R_0[1+VR_0]^1`$ we see that both conditions in (4.1) are satisfied.
From now on we consider only situations when $`V`$ is not bounded.
###### Proposition 4.4
Let $`H`$ be a self-adjoint operator such that $`V_kD(H)D(H)`$ for all $`k`$. Then $`H`$ is affiliated to $`C`$ if and only if $`D(H)D(\theta (P))`$ for some continuous divergent function $`\theta `$ and
(4.2)
$$\underset{k0}{lim}[V_k,H]_{D(H)D(H)^{}}=0.$$
Proof: It is clear that $`V_kD(H)D(H)`$ for all $`k`$ if and only if $`V_k`$ extends to a continuous map $`D(H)^{}D(H)^{}`$ for each $`k`$, and then we have in $`B()`$:
(4.3)
$$[V_k,(Hz)^1]=(Hz)^1[H,V_k](Hz)^1.$$
The operator $`[H,V_k]`$ belongs to $`B(D(H),)`$ and so we can consider it as a map $`D(H)D(H)^{}`$. But $`(Hz)^1`$ is an isomorphism $`D(H)`$ and $`D(H)^{}`$. To end the proof it suffices to use Lemma 4.1.
We shall give below three perturbative criteria of affiliation: we add to an operator affiliated to $`C`$ an operator which is not necessarily affiliated to it. Note that functions of $`Q`$ are never affiliated to $`C`$. First we consider operator bounded perturbations.
###### Corollary 4.5
Let $`H_0`$ be a self-adjoint operator affiliated to $`C`$ such that $`V_kD(H_0)D(H_0)`$ for all $`k`$. Let $`V`$ be a symmetric operator with domain $`D(H_0)`$ and such that $`H=H_0+V`$ is self-adjoint. Then $`H`$ is affiliated to $`C`$ if and only if
(4.4)
$$\underset{k0}{lim}[V_k,V]_{D(H_0)D(H_0)^{}}=0.$$
Now we want to consider form bounded perturbations in a generalized sense (in order to cover not semibounded operators). Let $`H`$ be a self-adjoint operator on $``$. We say that a Hilbert space $`𝒢`$ is adapted to $`H`$ if $`D(H)𝒢`$ continuously and densely and $`Hz`$ extends to an isomorphism $`𝒢𝒢^{}`$ for some (hence for all) $`z`$ outside the spectrum of $`H`$. Then $`H`$ extends to a continuous operator $`𝒢𝒢^{}`$ and we keep the notation $`H`$ for the extended map. It is not difficult to show that *if $`H`$ is a semibounded operator then $`𝒢`$ is adapted to $`H`$ if and only if $`𝒢=D(|H|^{1/2})`$ as topological vector spaces*, see \[GG2, page 47\]. But in general, for example in the case of Dirac operators, this is not the case. Observe that
$$D(H)𝒢𝒢^{}D(H)^{}$$
continuously and densely, in particular $`B(𝒢,𝒢^{})B(D(H),D(H)^{})`$. It is then clear that one has $`V_k𝒢𝒢`$ for all $`k`$ if and only if $`V_k`$ extends to a continuous map $`𝒢^{}𝒢^{}`$ for each $`k`$, and in this case the identity (4.3) is valid in $`B(𝒢^{},𝒢)`$. The operator $`[H,V_k]`$ belongs to $`B(𝒢,𝒢^{})`$ and so we can consider it as a map $`D(H)D(H)^{}`$. But $`(Hz)^1`$ is an isomorphism $`D(H)`$ and $`D(H)^{}`$. Thus:
###### Proposition 4.6
Let $`H`$ be a self-adjoint operator on $``$ such that $`D(H)D(\theta (P))`$ for some continuous divergent function $`\theta `$. Assume that $`𝒢`$ is a Hilbert space adapted to $`H`$ and that $`V_k𝒢𝒢`$ for all $`k`$. Then $`H`$ is affiliated to $`C`$ if and only if
(4.5)
$$\underset{k0}{lim}[V_k,H]_{D(H)D(H)^{}}=0.$$
In many situations of interest in quantum mechanics the domain of the Hamiltonian is difficult to determine while its form domain is quite explicit. For this reason the following condition stronger than (4.5) is often more convenient:
(4.6)
$$\underset{k0}{lim}[V_k,H]_{𝒢𝒢^{}}=0.$$
We shall use this in the following context.
###### Definition 4.7
Let $`H_0`$ be a self-adjoint operator on $``$ and let $`𝒢`$ be a Hilbert space adapted to it. We say that $`V`$ is a standard form perturbation of $`H_0`$ if $`V`$ is a continuous symmetric sesquilinear form on $`𝒢`$ and there are numbers $`\mu [0,1)`$ and $`\nu 0`$ such that one of the following conditions is satisfied:
(1) $`\pm V\mu |H_0|+\nu `$ as forms on $`𝒢`$
(2) $`H_0`$ is bounded from below and $`V\mu H_0\nu `$ as forms on $`𝒢`$.
Then the operator $`H=H_0+V:𝒢𝒢^{}`$ is such that its restriction to $`D(H)=\{u𝒢Hu\}`$ is a self-adjoint operator on $``$ (and will also be denoted $`H`$) and $`𝒢`$ is adapted to $`H`$ too (see \[DG3\]). Note that $`V`$ is seen as a continuous operator $`𝒢𝒢^{}`$.
###### Corollary 4.8
Let $`H_0`$ and $`V`$ as above. We assume that $`𝒢D(\theta (P))`$ for some continuous divergent function $`\theta `$, that $`V_k𝒢𝒢`$ for all $`k`$, and $`lim_{k0}[V_k,H]_{𝒢𝒢^{}}`$. Then $`H`$ is affiliated to $`C`$.
The next result covers perturbations of $`H_0`$ which are not dominated by $`H_0`$.
###### Proposition 4.9
Let $`H_1,H_2`$ be bounded from below self-adjoint operators and let us denote $`𝒢_i=D(|H_i|^{1/2})`$. Assume that $`𝒢𝒢_1𝒢_2`$ is dense in $``$ and let $`H=H_1+H_2`$, the sum being defined in form sense. Let us suppose that $`𝒢D(\theta (P))`$ for some continuous divergent function $`\theta `$ and that for $`i=1,2`$ we have $`V_k𝒢_i𝒢_i`$ and $`lim_{k0}[V_k,H_i]_{B(𝒢_i,𝒢_i^{})}=0`$. Then $`H`$ is affiliated to $`C`$.
Proof: Let us recall that the form sum $`H=H_1+H_2`$ is defined as the unique self-adjoint operator such that $`D(|H|^{1/2})=𝒢`$ and $`u,Hu=u,H_1u+u,H_2u`$ for all $`u𝒢`$. The topology of $`𝒢`$ is the intersection topology of $`𝒢_1`$ and $`𝒢_2`$, so thinking in terms of sesquilinear forms we see that
$$[V_k,H]_{B(𝒢,𝒢^{})}C[V_k,H_1]_{B(𝒢_1,𝒢_1^{})}+C[V_k,H_2]_{B(𝒢_2,𝒢_2^{})}$$
for some constant $`C`$. Hence (4.6) is satisfied.
4.2. If $`w`$ is a continuous divergent function on $`X^{}`$ let $`^w^w(X)=D(w(P))`$ equipped with the graph norm. We saw in Lemma 4.1 that if $`H`$ is affiliated to $`C`$ then $`D(|H|^{1/2})^w`$ for such a $`w`$. We consider now operators whose form domain is equal to some $`^w`$.
We say that $`w`$ is a weight<sup>2</sup><sup>2</sup>2 The terminology is suggested by that from \[Hör, Section 10.1\], cf. the remark after Theorem 10.1.5. if $`w:X^{}]0,\mathrm{}[`$ is continuous and $`w(k+p)\omega (k)w(p)`$ for some function $`\omega `$ and all $`k,pX^{}`$. If $`\omega `$ is the smallest function satisfying such an estimate, then $`\omega (k+p)\omega (k)\omega (p)`$. From now on we shall assume that $`\omega `$ satisfies this submultiplicativity condition. We also say $`\omega `$-weight if we need to be more specific. If $`X=^n`$ then a standard choice is $`w(k)=k^s`$ for some real $`s`$.
###### Lemma 4.10
A continuous divergent function $`w`$ on $`X^{}`$ is an $`\omega `$-weight if and only if $`V_k^w^w`$ and $`V_k_{B(^w)}\omega (k)`$ for all $`k`$.
Proof: We may take $`w(P)u`$ as norm on $`^w`$. From $`V_k^{}w(P)V_k=w(P+k)`$ we see that we have $`V_k^w^w`$ and $`V_k_{B(^w)}\omega (k)`$ if and only if $`w(P+k)u\omega (k)w(P)u`$ for all $`u`$, which is equivalent to $`w(k+p)\omega (k)w(p)`$ for all $`k,p`$.
###### Proposition 4.11
A self-adjoint operator on $``$ with $`D(|H|^{1/2})=^w`$ for some divergent weight $`w`$ and such that $`lim_{k0}[V_k,H]_{B(^w,^w)}=0`$ is affiliated to $`C`$.
This is an immediate consequence of Proposition 4.6.
###### Proposition 4.12
Let $`H`$ be as in Proposition 4.11 and bounded from below. Let $`VL_{\mathrm{loc}}^1(X)`$ be a real function whose negative part is form bounded with respect to $`H`$ with relative bound strictly less than 1. Then the self-adjoint operator $`H+V(Q)`$ (form sum) is affiliated to $`C`$.
Proof: Let $`V_+,V_{}`$ be the positive and negative parts of $`V`$, then we define the sum as $`(HV_{})+V_+`$ and apply successively Propositions 4.6 and 4.9.
###### Example 4.13
The most common situation is $`X=^n`$ and $`w(k)=k^s`$ for some real $`s>0`$. Then $`^w`$ is the usual Sobolev space $`^s`$ and typical operators satisfying the conditions of the Proposition 4.11 are the uniformly elliptic operators of order $`2s`$. For example, let $`s=m1`$ integer and
$$L=_{|\alpha |,|\beta |m}P^\alpha a_{\alpha \beta }P^\beta $$
for some measurable functions $`a_{\alpha \beta }:XB(𝐄)`$ such that the operator of multiplication by $`a_{\alpha \beta }`$ is a continuous map $`^{m|\beta |}^{|\alpha |m}`$ (this is a very general assumption which allows one to give a meaning to the differential expression $`L`$). Then $`L:^m^m`$ is a continuous map and $`V_k^{}LV_k`$ is a polynomial in $`k`$. If $`u,Lu\mu u_^m^2\nu u^2`$ for some $`\mu ,\nu >0`$, then $`L`$ induces a self-adjoint operator in $``$ which is affiliated to $`C`$.
###### Example 4.14
We give an explicit example of physical interest in the case $`s=1`$. Let
(4.7)
$$H=_{i,j}(P_iA_i)G_{ij}(P_jA_j)+V(PA)G(PA)+V$$
where $`G_{ij},A_i,V`$ are (the operators of multiplication by) locally integrable real functions having the following properties ($`_1`$ is the norm of $`^1`$):
(1) $`G_{ij}L^{\mathrm{}}(X)`$, the matrix $`G(x)=(G_{ij}(x))`$ is symmetric and $`G(x)\nu >0`$,
(2) for each $`\epsilon >0`$ there is $`\delta `$ such that $`A_ju\epsilon u_1+\delta u`$ for all $`u^1`$,
(3) if $`V_{}`$ is the negative part of $`V`$ then for each $`\epsilon >0`$ there is a real number $`\delta `$ such that $`u,V_{}u\epsilon u_1^2+\delta u^2`$ for all $`u^1`$.
Note that the conditions on $`A_j`$ and $`V_{}`$ are satisfied if there is $`s<1`$ such that $`A_juCu_s`$ and $`u,V_{}uCu_s^2`$. Then $`H`$ is affiliated to $`C`$. Indeed, observe first that $`H_0(PA)G(PA)`$ is a self-adjoint operator with form domain equal to $`^1`$, because there is $`\delta `$ such that:
$$u,H_0u\nu (PA)u^2\frac{\nu }{2}Pu^2\nu Au^2\frac{\nu }{4}Pu^2\delta u^2$$
Hence, according to Proposition 4.12, it suffices to prove that $`H_0`$ is affiliated to $`C`$. But
$`V_k^{}H_0V_k`$ $`=`$ $`(PA+k)G(PA+k)`$
$`=`$ $`H_0+kG(PA)+(PA)Gk+kGk.`$
Thus $`V_k^{}H_0V_kH_0_{B(^1,^1)}C(|k|+|k|^2)`$ so we can use Proposition 4.11.
###### Remark 4.15
Let us consider the operator $`H_0`$ under the more general condition $`A_jL_{\mathrm{loc}}^2(X)`$. More precisely, $`H_0`$ is the positive self-adjoint operator associated to the closed quadratic form $`(PA)u^2`$ whose domain is the set $`𝒢`$ of $`u`$ such that the distributions $`(P_jA_j)u`$ belong to $``$. The preceding computation shows that $`V_k𝒢𝒢`$ and that (4.6) is satisfied. Hence $`H_0`$ is affiliated to $`C`$ if and only if $`𝒢\theta (P)`$ for some continuous divergent function $`\theta `$. But this cannot be true without some boundedness conditions on $`A`$ at infinity.
As a final example we consider singular perturbations of $`h(P)`$, where $`h:X^{}`$ is a continuous divergent function and $`X`$ is an arbitrary group. Let $`𝒢=D(|h(P)|^{1/2})`$. Two functions $`u,v`$ on a neighborhood of infinity will be called equivalent if they satisfy $`c_1|u(k)||v(k)|c_2|u(k)|`$ for all large $`k`$ and some constants $`c_1,c_2>0`$. It is clear that $`𝒢=^w`$ if and only if $`h`$ is equivalent to $`w^2`$. Then Proposition 4.12 implies:
###### Proposition 4.16
Let $`h:X^{}`$ be a divergent function equivalent to a weight and such that
(4.8)
$$\underset{k0}{lim}\underset{p}{sup}\frac{|h(p+k)h(p)|}{1+|h(p)|}=0.$$
Let $`W`$ be a standard form perturbation of $`h(P)`$ with $`lim_{k0}[V_k,W]_{B(𝒢,𝒢^{})}=0`$ and define $`H_0=h(P)+W`$ as a form sum. Let $`VL_{\mathrm{loc}}^1(X)`$ real and such that $`V_{}\mu H_0+\nu `$ on $`𝒢`$ for some $`\mu <1,\nu >0`$. Then the form sum $`H=H_0+V(Q)`$ is a self-adjoint operator affiliated to $`C`$.
###### Example 4.17
Let $`X=^n`$ and assume that $`h`$ is of class $`C^1`$ and satisfies $`|h^{}(k)|C(1+|h(k)|)`$. Then (4.8) is fulfilled because
$$|h(p+k)h(p)|\underset{0<\theta <1}{sup}|h^{}(p+\theta k)||k|C\left(1+\underset{0<\theta <1}{sup}|h(p+\theta k)|\right)|k|$$
which is $`C^{}(1+|h(p)|)|k|`$ if $`|k|1`$ because $`h`$ is a equivalent to a weight. On the other hand, assume that $`h`$ is of class $`C^m`$ for some integer $`m1`$ and that we have: (1) $`lim_k\mathrm{}h(k)=+\mathrm{}`$, (2) the derivatives of order $`m`$ of $`h`$ are bounded, (3) $`_{|\alpha |m}|h^{(\alpha )}(k)|C(1+|h(k)|)`$. Then from \[ABG, p. 342–343\] we get that $`h`$ is equivalent to a weight. *Any real hypoelliptic polynomial satisfies all these conditions*, see Definition 11.1.2 and Theorem 11.1.3 in \[Hör\].
## 5 Localizations at infinity
In this section we prove our main result, Theorem 5.11, and some easy consequences.
5.1. We define first the localizations at infinity for functions in $`𝒞(X)`$. We denote $`𝒞_s(X)`$ the space $`𝒞(X)`$ equipped with the topology given by the seminorms $`\phi _\theta =\phi \theta `$ with $`\theta 𝒞_0(X)`$ (this is the strict topology associated to the essential ideal $`𝒞_0(X)`$).
###### Lemma 5.1
If $`\phi 𝒞(X)`$ and $`\varkappa \delta X`$ then $`\varkappa .\phi (y):=lim_{x\varkappa }\phi (x+y)`$ exists locally uniformly in $`yX`$. Equivalently, we have $`x.\phi \varkappa .\phi `$ in $`𝒞_s(X)`$ if $`x\varkappa `$ in $`\gamma X`$. The function $`\varkappa .\phi `$ belongs to $`𝒞(X)`$ and we have $`(\varkappa .\phi )(y)=\varkappa (y.\phi )`$.
Proof: Since $`\phi `$ is a bounded function, we have
$$\underset{x\varkappa }{lim}\phi (x+y)=\underset{x\varkappa }{lim}(y.\phi )(x)=\varkappa (y.\phi )$$
by taking into account the two interpretations of $`\varkappa `$. Then we use Lemma 2.2.
Thus $`\varkappa .\phi 𝒞(X)`$ is well defined for all $`\varkappa \gamma X`$ (if $`\varkappa =xX`$, see page 3) and all $`\phi 𝒞(X)`$. The next lemma is a slight improvement of Lemma 5.1, it will allow us to give a completely elementary proof of Theorem 5.16 (see the remark after the proof of the theorem). Note that the relation $`(\varkappa .\phi )(y)=\varkappa (y.\phi )`$ remains true for all $`\varkappa \gamma X`$ if we interpret $`xX`$ as a character of $`\mathrm{}_{\mathrm{}}(X)`$. Since $`y.\phi 𝒞(X)`$ we see that $`\varkappa .\phi `$ depends in fact only on the class of $`\varkappa `$ in $`\beta _\mathrm{u}X`$, cf. Subsection 2. We shall keep the notation $`\varkappa .\phi `$ even if $`\varkappa \beta _\mathrm{u}X`$. Recall that $`X\beta _\mathrm{u}X`$ is an open dense subset.
###### Lemma 5.2
Let $`\phi 𝒞(X)`$. Then $`Xxx.\phi 𝒞(X)`$ extends to a continuous function $`\beta _\mathrm{u}X\varkappa \varkappa .\phi 𝒞_s(X)`$. We have $`\varkappa .\phi (y)=\varkappa (y.\phi )`$ for all $`yX`$.
Proof: For $`\varkappa \beta _\mathrm{u}X=\sigma (𝒞(X))`$ the function $`\varkappa .\phi `$ is given by $`\varkappa .\phi (y)=\varkappa (y.\phi )`$, $`yX`$. It is easy to check directly that $`\varkappa .\phi `$ so defined belongs to $`𝒞(X)`$: we have $`|\varkappa (y.\phi )|y.\phi =\phi `$ and
$$|\varkappa (y.\phi )\varkappa (z.\phi )|=|\varkappa (y.\phi z.\phi )|y.\phi z.\phi =(yz).\phi \phi .$$
It remains to prove that $`\varkappa \varkappa .\phi \theta 𝒞_0(X)`$ is continuous for any $`\theta 𝒞_0(X)`$, i.e. that for each $`\chi \beta _\mathrm{u}X`$, each $`\epsilon >0`$ and each $`\theta 𝒞_0(X)`$ there is a neighborhood $`V`$ of $`\chi `$ in $`\beta _\mathrm{u}X`$ such that $`(\varkappa .\phi \chi .\phi )\theta <\epsilon `$ if $`\varkappa V`$. Since $`\theta `$ is $`C_0`$, it will suffice to prove that for each $`\chi `$ and $`\epsilon `$ as before and each compact set $`KX`$ there is a neighborhood $`V`$ of $`\chi `$ such that $`\varkappa V`$ implies $`|\varkappa (y.\phi )\chi (y.\phi )|<\epsilon `$ for $`yK`$. But the map $`yy.\phi 𝒞(X)`$ is norm continuous, thus $`\{y.\phi yK\}`$ is a compact subset of $`𝒞(X)`$. Hence there is a finite subset $`Z`$ of $`K`$ such that for each $`yK`$ we have $`\mathrm{min}_{zZ}y.\phi z.\phi <\epsilon `$. Thus for each $`yK`$ and $`zZ`$ we have
$`|\varkappa (y.\phi )\chi (y.\phi )|`$ $`=`$ $`|\varkappa (y.\phi z.\phi )+\varkappa (z.\phi )\chi (z.\phi )+\chi (z.\phi y.\phi )|`$
$`<`$ $`2\epsilon +|\varkappa (z.\phi )\chi (z.\phi )|.`$
Now, if we take $`V=\{\varkappa \beta _\mathrm{u}Xsup_{zZ}|\varkappa (z.\phi )\chi (z.\phi )|<\epsilon \}`$, then $`V`$ is a neighborhood of $`\chi `$ in $`\beta _\mathrm{u}X`$ because $`Z`$ is a finite set, and for each $`\varkappa V`$ and each $`yK`$ we have $`|\varkappa (y.\phi )\chi (y.\phi )|<3\epsilon `$.
###### Lemma 5.3
If $`\phi 𝒞(X)`$ then $`\varkappa .\phi =0`$ for all $`\varkappa \delta X`$ if and only if $`\phi 𝒞_0(X)`$.
Proof: If $`\varkappa .\phi =0`$ for all $`\varkappa \delta X`$ then $`\varkappa (\phi )=(\varkappa .\phi )(0)=0`$ for such $`\varkappa `$. If $`\phi 𝒞_0(X)`$ then there is a number $`a>0`$ such that the set $`U=\{x|\phi (x)|>a\}`$ is not relatively compact. Since $`UV\mathrm{}`$ for each $`V`$ with relatively compact complement, we see that the family of sets $`UV`$ is a filter basis and the filter $`𝔣`$ it generates is finer than Fréchet and contains $`U`$. Let $`\varkappa `$ be any ultrafilter finer than $`𝔣`$, then $`\varkappa \delta X`$ and $`\varkappa (\text{1}_U)=1`$. Finally, from $`|\phi |a\text{1}_F`$ we get $`|\varkappa (\phi )|^2=\varkappa (|\phi |^2)a^2\varkappa (\text{1}_V)=a^2`$, so we cannot have $`\varkappa (\phi )=0`$.
###### Definition 5.4
If $`\phi 𝒞(X)`$ and $`\varkappa \delta X`$ then the function $`\varkappa .\phi 𝒞(X)`$ is the *localization of $`\phi `$ at $`\varkappa `$*. And $`\mathrm{}(\phi ):=\{\varkappa .\phi \varkappa \delta X\}𝒞(X)`$ is the set of *localizations of $`\phi `$ at infinity*.
For each $`\varkappa \delta X`$ let $`\tau _\varkappa :𝒞(X)𝒞(X)`$ be given by $`\tau _\varkappa (\phi )=\varkappa .\phi `$. Clearly this is a unital morphism and, since the property $`x.(\varkappa .\phi )=\varkappa .(x.\phi )`$ is easy to check, $`\tau _\varkappa `$ is in fact an $`X`$-morphism. By Lemma 5.3 we have
(5.1)
$$_{\varkappa \delta X}\mathrm{ker}\tau _\varkappa =𝒞_0(X).$$
Note that *$`\mathrm{ker}\tau _\varkappa `$ is the maximal $`X`$-ideal included in the maximal ideal $`\mathrm{ker}\varkappa `$ of $`𝒞(X)`$*.
###### Remark 5.5
In general $`\tau _\varkappa \tau _\chi \tau _\chi \tau _\varkappa `$.
5.2. In this subsection we extend the notion of localization to operators in $`C(X)`$.
###### Definition 5.6
Let $`C_s(X)`$ be the space $`C(X)`$ equipped with the topology defined by the family of seminorms $`T_\theta =T\theta (Q)+\theta (Q)T`$ with $`\theta 𝒞_0(X)`$.
Note that if $`X`$ is $`\sigma `$-compact then there is $`\theta 𝒞_0(X)`$ with $`\theta (x)>0`$ for all $`xX`$ and then $`_\theta `$ is a norm on $`C(X)`$ which induces on bounded subsets of $`C(X)`$ the topology of $`C_s(X)`$. In any case, the topology of $`C_s(X)`$ is finer than the strong operator topology induced by $`B(X)`$. Note also that the topology of $`C_s(X)`$ does not depend on any Hilbert space realization of $`C(X)`$ because $`C(X)`$ is a $`𝒞(X)`$-bimodule and $`𝒞_0(X)`$ is an ideal of $`𝒞(X)`$. Finally, observe that we could consider on $`C(X)`$ the (intrinsically defined) strict topology associated to the ideal $`K(X)`$; this is weaker than that of $`C_s(X)`$ and finer than the strong operator topology (but coincides with it on bounded sets).
###### Remark 5.7
That this is the natural topology in our context should have been clear for us a long time ago, since it is induced by the strict topology of $`𝒞(X)`$, cf. \[GI2, p. 31\] and \[GI5, p. 148\]. However, we did not realize it until B. Simon, in a private communication, emphasized its importance, in relation with Proposition 3.11 and Theorem 4.5 from \[LaS\]. We are indebted to him for this remark. On the other hand, note that this topology does not play any rôle in our paper, the strong operator topology on $`C(X)`$ (used in \[GI2, GI5\]) suffices.
We now describe some topological properties of $`C_s(X)`$.
###### Lemma 5.8
The map $`TT^{}`$ is continuous on $`C_s(X)`$ and the operation of multiplication is continuous on bounded sets. If $`TC(X)`$ the map $`xU_xTU_x^{}C(X)`$ is norm continuous and the set $`\{U_xTU_x^{}xX\}`$ is relatively compact in $`C_s(X)`$.
Proof: The first assertion is obvious. To prove the second one, note first that *if $`SC(X)`$ and $`\theta 𝒞_0(X)`$ then the operators $`S\theta (Q)`$ and $`\theta (Q)S`$ are compact*. Indeed, it suffices to show this for $`S`$ of the form $`\phi (Q)\psi (P)`$ and then the assertion is obvious. In particular, from the Remark 3.13 it follows that there are $`KK(X)`$ and $`\theta ^{}𝒞_0(X)`$ such that $`S\theta (Q)=\theta ^{}(Q)K`$, and similarly for $`\theta (Q)S`$. Thus for $`A,B,S,TC(X)`$ we have
$$(BATS)\theta (Q)B(AS)\theta (Q)+(BT)\theta ^{}(Q)K$$
from which the continuity of multiplication follows. The norm continuity of $`xU_xTU_x^{}`$ is obvious by (1.1). Finally, the last assertion of the lemma says that $`xU_xTU_x^{}\theta (Q)`$ has relatively compact range and similarly when $`\theta `$ is on the left side. Clearly it suffices to take $`T=\phi (Q)\psi (P)`$ and then $`U_xTU_x^{}\theta (Q)=\phi (Q+x)\psi (P)\theta (Q)`$ and $`\psi (P)\theta (Q)`$ is a compact operator. Now the assertion follows from the Riesz-Kolmogorov criterion (Remark 3.13) which clearly implies: if $`K`$ is a compact operator and $`\phi 𝒞(X)`$ then $`\phi (Q+x)K`$ is a norm relatively compact family of operators.
###### Proposition 5.9
If $`TC(X)`$ and $`\varkappa \delta X`$ then $`\varkappa .T:=lim_{x\varkappa }U_xTU_x^{}`$ exists in the topological space $`C_s(X)`$. The map $`\tau _\varkappa :C(X)C(X)`$ defined by $`\tau _\varkappa (T)=\varkappa .T`$ is a morphism uniquely determined by the property:
(5.2)
$$\phi 𝒞(X),\psi 𝒞_0(X^{})\tau _\varkappa \left(\phi (Q)\psi (P)\right)=(\varkappa .\phi )(Q)\psi (P).$$
If $`TC(X)`$ and $`\psi :X^{}`$ is a bounded continuous function, then
(5.3)
$$\tau _\varkappa \left(T\psi (P)\right)=\tau _\varkappa (T)\psi (P)\text{and}\tau _\varkappa \left(\psi (P)T\right)=\psi (P)\tau _\varkappa (T).$$
For each $`kX^{}`$ we have $`\tau _\varkappa \left(V_k^{}TV_k\right)=V_k^{}\tau _\varkappa (T)V_k`$.
Proof: We must show that there is an operator $`\varkappa .TC(X)`$ such that
$$\underset{x\varkappa }{lim}(U_xTU_x^{}\varkappa .T)\theta (Q)=\underset{x\varkappa }{lim}\theta (Q)(U_xTU_x^{}\varkappa .T)=0$$
for all $`\theta 𝒞_0(X)`$. It is clearly sufficient to consider $`T=\phi (Q)\psi (P)`$ with $`\phi 𝒞(X)`$ and $`\psi 𝒞_0(X^{})`$. Then we have
$$U_xTU_x^{}\theta (Q)=\phi (Q+x)\psi (P)\theta (Q)=\phi (Q+x)\theta ^{}(Q)K$$
for some $`\theta ^{}𝒞_0(X)`$ and $`KK(X)`$. Indeed, $`\psi (P)\theta (Q)`$ is a compact operator and so we can use the Remark 3.13. Now it suffices to use Lemma 5.1. The argument for $`\theta (Q)U_xTU_x^{}`$ is even simpler. The other assertions are easy to prove, for example the last assertion follows from $`V_k^{}U_xTU_x^{}V_k=U_xV_k^{}TV_kU_x^{}`$.
###### Proposition 5.10
Let $`TC(X)`$. Then $`\varkappa .T=0`$ for each $`\varkappa \delta X`$ if and only if $`TK(X)`$.
Proof: In order to prove that $`\varkappa .T=0`$ if $`TK(X)`$ it suffices to consider $`T=\phi (Q)\psi (P)`$ with $`\phi 𝒞_0(X)`$. Then $`\varkappa .T=(\varkappa .\phi )(Q)\psi (P)`$ and $`\varkappa .\phi =0`$ if $`\phi 𝒞_0(X)`$. Reciprocally, let $`J=\{TC(X)\varkappa .T=0,\varkappa \delta X\}`$ and notice that $`J`$ is a $`C^{}`$-algebra and, moreover, it is a crossed product because of the last assertions of Proposition 5.9. Also, for each $`SC(X)`$ we have $`\varkappa .(ST)=(\varkappa .S)(\varkappa .T)=0`$ so $`J`$ is an ideal. Thus, by Proposition 3.15, there is an ideal $`𝒥`$ in $`𝒞(X)`$ such that $`J=𝒥X`$. Let us show that $`𝒥=𝒞_0(X)`$. This will finish the proof, because then $`J=𝒞_0(X)X=K(X)`$. From (3.3) we get
$$𝒥=\{\phi 𝒞(X)\varkappa .(\phi (Q)\psi (P))=0\psi 𝒞_0(X^{})\text{ and }\varkappa \delta X\}.$$
But $`\varkappa .(\phi (Q)\psi (P))=(\varkappa .\phi )(Q)\psi (P)`$. On the other hand, if $`\theta 𝒞(X)`$ is such that $`\theta (Q)\psi (P)=0`$ $`\psi 𝒞_0(X^{})`$, then $`\theta =0`$. Indeed, $`V_k^{}\theta (Q)\psi (P)V_k=\theta (Q)\psi (P+k)`$, so if $`\psi L^1(X^{})`$ then we have in the weak operator topology
$$0=_X^{}V_k^{}\theta (Q)\psi (P)V_k\text{d}k=\theta (Q)_X^{}\psi (P+k)\text{d}k=\theta (Q)_X^{}\psi \text{d}k.$$
Thus it suffices to take $`\psi `$ such that $`_X^{}\psi \text{d}k0`$. So we finally see that $`𝒥`$ is the set of $`\phi 𝒞(X)`$ such that $`\varkappa .\phi =0`$ for all $`\varkappa \delta X`$, i.e. $`𝒥=𝒞_0(X)`$ by Lemma 5.3.
The next result follows easily from Propositions 5.9 and 5.10.
###### Theorem 5.11
The map $`T(\varkappa .T)_{\varkappa \delta X}`$ is a morphism $`C(X)_{\varkappa \delta X}C(X)`$ with $`K(X)`$ as kernel, so we have a canonical embedding
(5.4)
$$C(X)/K(X)_{\varkappa \delta X}C(X).$$
Theorem 1.1 is an immediate consequence. As explained in Subsection 2, the morphism $`\tau _\varkappa `$ extends to observables affiliated to $`A`$ and Theorem 1.2 follows easily.
###### Remark 5.12
It has been brought to our attention by Steffen Roch that it is not possible to deduce Theorem 1.1 for not normal operators from Theorem 5.11, as we stated in an earlier version of this paper, because the spectrum of a general element of an infinite product of $`C^{}`$-algebras is not so simply related to the spectra of its components. We could have stated a version of Theorem 1.1 valid for not normal operators in the spirit of \[RRS2, Theorem 2.2.1\] but we did not do it because the only applications we have in mind refer to quantum Hamiltonians, which are self-adjoint operators. We mention, however, that for some algebras $`𝒜`$ the Theorem 1.15 remains true (without closure) for non normal operators, see Sections 2.4 and 2.5 in \[RRS2\].
###### Definition 5.13
If $`H`$ is an observable affiliated to $`C(X)`$ and if $`\varkappa \delta X`$ then the observable $`\varkappa .H`$ affiliated to $`C(X)`$ is called *localization of $`H`$ at $`\varkappa `$*. The set of operators $`\mathrm{}(H):=\{\varkappa .H\varkappa \delta X\}`$ is the set of *localizations of $`H`$ at infinity*.
Then we can write the relation (1.4) as follows:
(5.5)
$$\sigma _{\mathrm{ess}}(H)=\overline{}_{\varkappa \delta X}\sigma (\varkappa .H)=\overline{}_{K\mathrm{}(H)}\sigma (K).$$
###### Remark 5.14
By using the universal property of the Stone-Čech compactification $`\gamma X`$ (cf. page 2) we see that for $`TB(X)`$ the following two assertions are equivalent:
(1) the set $`\{x.TxX\}`$ is strongly relatively compact in $`B(X)`$;
(2) $`Xxx.T`$ extends to a strongly continuous map $`\gamma X\varkappa \varkappa .TB(X)`$.
The set of operators having these properties is a norm closed subalgebra of $`B(X)`$ (quite large, it contains $`C(X),L^{\mathrm{}}(X),L^{\mathrm{}}(X^{})`$ and much more). It is easy to check that $`\sigma (\varkappa .T)\sigma _{\mathrm{ess}}(T)`$ if $`\varkappa \delta X`$, but in most cases the operators $`\varkappa .T`$ do not suffice to determine the essential spectrum of $`T`$. This fact extends to observables affiliated to this algebra. For example, if $`H`$ is the Hamiltonian of a particle in 2 dimensions in a constant non-zero magnetic field, then $`T=\phi (H)`$ has the property (1) and $`\varkappa .T=0`$ if $`\phi 𝒞_0()`$, i.e. $`\varkappa .H=\mathrm{}`$ for all $`\varkappa \delta X`$. But $`\sigma _{\mathrm{ess}}(H)\mathrm{}`$.
5.3. We fix now an algebra of interactions $`𝒜`$ on $`X`$, and set $`A=𝒜XC`$. Theorem 5.11 gives a description of $`A/K(X)`$ but we can make it more precise because many ultrafilters give the same character of $`𝒜`$.
###### Definition 5.15
If $`\varkappa \delta X`$ the $`C^{}`$-algebras $`𝒜_\varkappa =\tau _\varkappa (𝒜)`$ and $`A_\varkappa =\tau _\varkappa (A)=𝒜_\varkappa X`$ are the *localizations at $`\varkappa `$* of the algebras $`𝒜`$ and $`A`$ respectively.
As explained in Subsection 2, and taking into account the relation $`(\varkappa .\phi )(y)=\varkappa (y.\phi )`$ (see Lemma 5.1) and Lemma 2.1, we see that $`𝒜_\varkappa `$ and $`A_\varkappa `$ depend only on the restriction to $`𝒜`$ of the character $`\varkappa `$. In other terms, we have for example $`𝒜_\varkappa =𝒜_\chi `$ if $`\delta (\varkappa )=\delta (\chi )`$, where $`\delta :\delta X\delta (𝒜)`$ is the canonical surjection, cf. (2.15). According to the convention made in Subsection 2 (see page 2.1) we shall use the same notations $`𝒜_\varkappa `$ and $`A_\varkappa `$ if $`\varkappa \delta (𝒜)`$.
In the statement of the next theorem we use the canonical identification of $`X`$ (as topological space) with an open dense subset of $`\sigma (𝒜)`$
###### Theorem 5.16
If $`TA`$ the norm continuous map $`Xxx.TAC`$ extends to a continuous map $`\sigma (𝒜)\varkappa \varkappa .TC_s(X)`$. For each $`\varkappa \delta (𝒜)`$ the map $`\tau _\varkappa :AC`$ defined by $`\tau _\varkappa (T)=\varkappa .T`$ is a morphism with $`A_\varkappa `$ as range. One has $`\varkappa .T=0`$ for all $`\varkappa \delta (A)`$ if and only if $`TK(X)`$ which gives a canonical embedding
(5.6)
$$A/K(X)_{\varkappa \delta (𝒜)}A_\varkappa .$$
Proof: Consider for each $`TA`$ the map $`F_T:\gamma XC_s(X)`$ defined by $`F_T(\varkappa )=\varkappa .T`$. From Lemma 5.2 it follows that $`F_T`$ is continuous: indeed, it suffices to assume that $`T=\phi (Q)\psi (P)`$ and to argue as in the proof of Lemma 5.8. Notice that if the characters $`\varkappa ,\chi \gamma X`$ are equal on $`𝒜`$, then $`F_T(\varkappa )=F_T(\chi )`$. Indeed, for $`T`$ as above we have $`\varkappa .T=(\varkappa .\phi )(Q)\psi (P)=(\chi .\phi )(Q)\psi (P)=\chi .T`$. Thus, as explained on page 2, if $`\pi :\gamma X\sigma (A)`$ is the canonical surjection, we shall have $`F_T=f_T\pi `$, where $`f_T:\sigma (A)C_s(X)`$ is continuous. If $`xX`$ then $`\pi (x)=x`$ so $`f_T(x)=F_T(x)=x.T`$. We have both $`X\sigma (A)`$ and $`X\gamma X`$ and since the restriction of $`\pi `$ to $`X`$ is the identity mapping, $`\pi `$ acts non-trivially only on the boundary. Let $`\delta `$ be the restriction of the map $`\pi `$ to $`\delta X`$, hence $`\delta :\delta X\delta (A)`$ is a canonical surjection. Thus $`f_T(\varkappa )=0`$ for all $`\varkappa \delta (A)`$ is equivalent to $`F_T(\varkappa )=0`$ for all $`\varkappa \delta X`$ which means that $`TK(X)`$.
Remark: By using the last assertion of Lemma 5.8 and the universal property of the space $`\gamma X`$, cf. page 2, one may avoid the use of Lemma 5.2.
###### Remark 5.17
In nice situations, the localization at infinity $`𝒜_\varkappa `$ is simpler than $`𝒜`$, and $`(𝒜_\varkappa )_\chi `$ is still simpler, and so on, but this is not always the case. Note also that in general $`𝒜_\varkappa 𝒜`$. If, however, this holds for each $`\varkappa \delta (𝒜)`$, then it is natural to ask whether we have $`\tau _\varkappa \tau _\chi \phi =\tau _\chi \tau _\varkappa \phi `$ for all $`\phi 𝒜`$ and all $`\varkappa ,\chi \delta (𝒜)`$. Although this is not true if $`𝒜=𝒞(X)`$, in several non-trivial and physically interesting situations this property is satisfied. See Examples 5.18 and 5.19 and Section 6.
###### Example 5.18
We shall consider here the localizations at infinity of the simplest algebras. If $`𝒜=𝒞_{\mathrm{}}(X)`$ then $`\sigma (𝒜)=X\{\mathrm{}\}`$ is the Alexandroff compactification of $`X`$, we have $`\delta (𝒜)=\{\mathrm{}\}`$, and the localization of $`\phi 𝒜`$ at $`\mathrm{}`$ is the constant function which takes the value $`\phi (\mathrm{})=lim_x\mathrm{}\phi (x)`$. If $`X=`$ and $`𝒜`$ is the set of bounded continuous functions which have limits as $`x\pm \mathrm{}`$ then $`\sigma (𝒜)=[\mathrm{},+\mathrm{}]`$, $`\delta (𝒜)=\{\mathrm{},+\mathrm{}\}`$, and the localization of $`\phi 𝒜`$ at $`+\mathrm{}`$ is again the constant function which takes the value $`\phi (+\mathrm{})=lim_{x+\mathrm{}}\phi (x)`$ and similarly for the localization at $`\mathrm{}`$. Thus in both examples we have $`𝒜_\varkappa =`$ for all $`\varkappa \delta (𝒜)`$. In Subsection 6.1 we shall describe explicitly the largest $`X`$-subalgebra $`𝒜𝒞(X)`$ such that $`𝒜_\varkappa =`$ for all $`\varkappa \delta (𝒜)`$.
###### Example 5.19
The next example is due to Gilles Godefroy (we thank him for answering to our questions) and is relevant in the context of Remark 5.17. Let $`X=\times `$ and let $`𝒜`$ be the set of $`\phi \mathrm{}_{\mathrm{}}(X)`$ such that $`lim_k\mathrm{}\phi (j,k)=0`$ for all $`j`$. Let $`\theta \mathrm{}_{\mathrm{}}()`$ and set $`\phi (j,k)=\theta (k)`$ if $`|k|j`$ and $`=0`$ otherwise. Then $`\phi 𝒜`$ and $`lim_{a+\mathrm{}}\phi (a+j,k)=\theta (k)`$ for each $`j,k`$. It is clear now that we may construct an ultrafilter $`\varkappa \delta X`$ such that $`\varkappa .\phi =1\theta `$ so $`\varkappa .\phi 𝒜`$ in general.
Theorem 1.15 is a corollary of Theorem 5.16. Thus, if $`H`$ is a normal element of $`A`$ or an observable affiliated to $`A`$ and if we set $`\varkappa .H=\tau _\varkappa (H)`$, then
(5.7)
$$\sigma _{\mathrm{ess}}(H)=\overline{}_{\varkappa \delta (𝒜)}\sigma (\varkappa .H).$$
This representation of the essential spectrum of $`H`$, although more precise than (5.5), is still quite redundant, cf. page 1.12, and can be improved in many situations (the most interesting one being the $`N`$-body case). To explain this, for $`\varkappa \delta (𝒜)`$ let us denote
(5.8)
$$𝒥_\varkappa =\mathrm{ker}\tau _\varkappa =\{\phi 𝒜\varkappa (x.\phi )=0xX\}.$$
This is is the maximal $`X`$-ideal included in the maximal ideal $`\mathrm{ker}\varkappa `$ of $`𝒜`$. Although the ideals $`\mathrm{ker}\varkappa `$ for different $`\varkappa `$ are not comparable, it often happens that the $`𝒥_\varkappa `$ are comparable, i.e. we may have $`𝒥_\varkappa 𝒥_\chi `$ for $`\varkappa \chi `$.
###### Lemma 5.20
If $`𝒥_\varkappa 𝒥_\chi `$ then $`\sigma (H_\chi )\sigma (H_\varkappa )`$. In particular, (5.7) remains true if we restrict the union to the $`\varkappa `$ such that the ideal $`𝒥_\varkappa `$ is minimal in $`\{𝒥_\varkappa \varkappa \delta (𝒜)\}`$.
Proof: Here we use more abstract algebraic tools, as in \[GI2, GI4\]. The morphism $`\tau _\varkappa :𝒜𝒜_\varkappa `$ is surjective and has $`𝒥_\varkappa `$ as kernel, hence induces an isomorphism $`𝒜/𝒥_\varkappa 𝒜_\varkappa `$. If $`T𝒜`$ and if $`T/𝒥_\varkappa `$ is its projection in the quotient $`𝒜/𝒥_\varkappa `$, then $`T/𝒥_\varkappa `$ is sent by this isomorphism into $`\varkappa .T`$, hence $`\sigma (T/𝒥_\varkappa )=\sigma (\varkappa .T)`$. From $`𝒥_\varkappa 𝒥_\chi `$ we get a canonical surjective morphism $`𝒜/𝒥_\varkappa 𝒜/𝒥_\chi `$ which sends $`T/𝒥_\varkappa `$ into $`T/𝒥_\chi `$. Finally, we recall that if $`\mathrm{\Phi }`$ is a morphism then $`\sigma (\mathrm{\Phi }(S))\sigma (S)`$.
###### Example 5.21
If, for $`xX`$ and $`\varkappa \delta (𝒜)`$, we denote $`x+\varkappa `$ the character $`\varkappa \tau _x`$, then clearly $`𝒥_{x+\varkappa }=𝒥_\varkappa `$ , hence $`\sigma ((x+\varkappa ).H)=\sigma (\varkappa .H)`$. However, this case is trivial because clearly $`(x+\varkappa ).H=U_x(\varkappa .H)U_x^{}`$.
One further simplification may be obtained as follows.
###### Lemma 5.22
Let $`𝒦\delta (𝒜)`$ such that: if $`\phi 𝒜`$ and $`\varkappa (x.\phi )=0`$ for all $`\varkappa 𝒦`$ and $`xX`$, then $`\phi 𝒞_0(X)`$. Then (5.7) remains valid if $`\delta (𝒜)`$ is replaced by $`𝒦`$.
Proof: This is a consequence of the proof of Theorem 5.16 but can also be proved directly as follows. One first notices that the condition on $`𝒦`$ is equivalent to the density in $`\delta (𝒜)=\sigma (𝒜/𝒞_0(X))`$ of the set of characters of the form $`\varkappa \tau _x`$, with $`\varkappa 𝒦`$ and $`xX`$. Then one can use the following easily proven fact: if $`S_\alpha `$ is a net of operators such that $`S_\alpha ^{()}S^{()}`$ strongly, then $`\sigma (S)`$ is included in the closure of $`_\alpha \sigma (S_\alpha )`$.
## 6 Applications
After some preliminaries, we describe here three classes of $`C^{}`$-algebras of Hamiltonians which seem to us particularly relevant and treat some more explicit examples.
6.1. Algebras associated to translation invariant filters. In this preliminary subsection we give an intrinsic description of a class of crossed products introduced in \[GI2, GI4\]. Recall that a filter $`𝔣`$ is *translation invariant* if: $`xX,F𝔣x+F𝔣`$. Note that $`𝔣^{}`$ will also be translation invariant. If $`𝔣`$ is a translation invariant filter let
(6.1)
$$𝒥(𝔣)=\{\phi 𝒞(X)lim_𝔣\phi =0\}.$$
This is clearly an $`X`$-ideal in $`𝒞(X)`$ and from Lemma 2.2 we get:
(6.2)
$$𝒥(𝔣)=𝒥(𝔣^{}).$$
Then $`𝒞(𝔣)=+𝒥(𝔣)`$ is the $`X`$-algebra consisting of the bounded uniformly continuous functions $`\phi `$ such that $`lim_𝔣\phi `$ exists. Observe that if $`𝔣`$ is the Fréchet filter then $`𝒥(𝔣)=𝒞_0(X)`$ and $`𝒞(𝔣)=𝒞_{\mathrm{}}(X)`$.
Below we shall consider nets indexed by the filter $`𝔣`$ equipped with the order relation $`FGFG`$. For example, $`lim_{F𝔣}\text{1}_F(Q)T=0`$ means that for each $`\epsilon >0`$ there is a Borel set $`F𝔣`$ such that $`\text{1}_F(Q)T<\epsilon `$.
###### Proposition 6.1
$`𝒥(𝔣)X=\{TC(X)lim_{F𝔣}\text{1}_F(Q)T^{()}=0\}`$.
Proof: Each $`T𝒥(𝔣)X`$ has the property $`lim_{F𝔣}\text{1}_F(Q)T=0`$. Indeed, it suffices to consider operators of the form $`T=\phi (Q)\psi (P)`$ with $`\phi 𝒥(𝔣),\psi 𝒞_0(X^{})`$. But then the set $`F`$ of points $`x`$ such that $`|\phi (x)|<\epsilon `$ is open and belongs to $`𝔣`$, and so we have $`\text{1}_F(Q)\phi (Q)\epsilon `$, which is more than needed.
Conversely, let $`J`$ be the set of $`TC(X)`$ such that $`lim_{F𝔣}\text{1}_F(Q)T^{()}=0`$. This is clearly a $`C^{}`$-subalgebra of $`C(X)`$ which is stable under the morphisms $`TV_k^{}TV_k`$. By Theorem 3.7 we have $`J=𝒥X`$ for a unique $`X`$-algebra $`𝒥`$, namely the set of $`\phi 𝒞(X)`$ such that $`lim_{F𝔣}\text{1}_F(Q)\phi (Q)^{()}\psi (P)=0`$ for all $`\psi 𝒞_0(X^{})`$.
Thus it remains to prove the following assertion: if $`\phi 𝒞(X)`$ has the property $`lim_{F𝔣}\text{1}_F(Q)\phi (Q)\psi (P)=0`$ for $`\psi 𝒞_0(X^{})`$, then $`lim_{F𝔣}\text{1}_F(Q)\phi (Q)=0`$. Observe that, due to (6.2) we may assume $`𝔣=𝔣^{}`$.
Fix $`fL^2(X),\psi 𝒞_0(X^{})`$ and let us set $`\theta =\psi (P)f`$ and $`\theta _a(x)=(U_a^{}\theta )(x)=\theta (xa)`$. Clearly $`lim_{F𝔣}\text{1}_F(Q)\phi (Q)U_a^{}\theta =0`$ uniformly in $`aX`$. Thus, for any $`\epsilon >0`$, there is $`F𝔣`$ Borel such that $`\text{1}_F\phi \theta _a<\epsilon `$ for all $`a`$, hence
$$|\phi (a)|\text{1}_F\theta _a\text{1}_F(\phi (a)\phi )\theta _a+\text{1}_F\phi \theta _a\text{1}_F(\phi (a)\phi )\theta _a+\epsilon .$$
Since $`𝔣=𝔣^{}`$ we may assume that $`F=G+V`$ where $`G𝔣`$ and $`V`$ is a a compact neighborhood of the origin. Moreover, since $`\phi `$ is uniformly continuous and since we may choose $`V`$ as small as we wish, we may assume that $`|\phi (x)\phi (a)|<\epsilon `$ if $`xaV`$. It is possible to choose $`f,\psi `$ such that $`\text{supp}\theta V`$ and $`\theta =1`$. Indeed, $`\theta `$ is equal to the convolution product $`\stackrel{~}{\psi }f`$ where $`\stackrel{~}{\psi }(x)=\widehat{\psi }(x)`$ and it suffices to choose $`f,\stackrel{~}{\psi }`$ continuous, positive and not zero and such that $`\text{supp}f+\text{supp}\stackrel{~}{\psi }V`$. Then for $`aG`$ we clearly have $`\text{supp}\theta _aF`$ hence
$$|\phi (a)|=|\phi (a)|\theta _a=|\phi (a)|\text{1}_F\theta _a(\phi (a)\phi )\theta _a+\epsilon 2\epsilon .$$
This proves that $`lim_𝔣\phi =0`$.
From Proposition 6.1 we easily get:
$$𝒞(𝔣)X=\{TC(X)S𝒞_0(X^{})\text{ such that }\underset{F𝔣}{lim}\text{1}_F(Q)(TS)^{()}=0\}.$$
The $`X`$-algebras of the form $`_\lambda 𝒞(𝔣_\lambda )`$ are of some physical interest \[Ric\]. Indeed, one should think of a filter finer than the Fréchet filter as the set of traces on $`X`$ of the filter of neighborhoods of some closed part of the boundary of $`X`$ in a compactification of $`X`$. This explains the interest of the algebras $`_\lambda 𝒞(𝔣_\lambda )`$ in the present context: they consist of “potentials” which have limits at infinity when going in certain directions. One may easily deduce from Theorem 3.14 and Proposition 6.1 an intrinsic description of the crossed products $`_\lambda 𝒞(𝔣_\lambda )X`$.
6.2. The $`V(X)`$ algebra. We shall consider now the simplest non-trivial functions in $`𝒞(X)`$, those all of whose localizations at infinity are constants. Our purpose is to give a simple characterization of the $`X`$-algebra $`𝒜`$ defined by the condition $`𝒜_\varkappa =`$ for all $`\varkappa \delta X`$ and of the associated crossed product. So we introduce the $`X`$-algebra:
(6.3)
$$𝒱(X):=\{\phi 𝒞(X)\varkappa .\phi ,\varkappa \delta X\}$$
Observe that the relation $`\varkappa .\phi `$ is equivalent to $`\varkappa .\phi =\varkappa (\phi )`$.
###### Lemma 6.2
We have $`\phi 𝒱(X)`$ if and only if $`\phi 𝒞(X)`$ and
(6.4)
$$\underset{x\mathrm{}}{lim}(\phi (x+y)\phi (x))=0,yX.$$
Proof: The condition (6.4) is equivalent to $`y.\phi \phi 𝒞_0(X)`$ for all $`yX`$ and, by (5.1), this is equivalent to $`\varkappa (y.\phi \phi )=0`$ for all $`\varkappa \delta X`$ and all $`y`$, hence to $`\varkappa .\phi (y)=\varkappa (\phi )`$ for all $`\varkappa ,y`$, which means $`\phi 𝒱(X)`$.
It is easily shown that $`\phi 𝒞(X)`$ satisfies (6.4) if and only if $`\phi `$ is a bounded continuous function such that $`lim_x\mathrm{}(\phi (x+y)\phi (x))=0`$ uniformly in $`y`$ when $`y`$ runs over a compact neighborhood of the origin. Thus the functions from $`𝒱(X)`$ are of *vanishing oscillation at infinity* or *slowly oscillating*, and their rôle in the theory of pseudo-differential operators was noticed a long time ago due especially to a well known theorem of H. Cordes concerning the compactness of the commutators $`[\phi (Q),\psi (P)]`$ (see \[ABG, p. 176–177\] for a short presentation of the main ideas). If $`X=^n`$ then $`𝒱(X)`$ is just the norm closure of the set of bounded functions of class $`C^1`$ whose derivative tends to zero at infinity. Thus results of the same nature as the embedding (6.6) may be found already in \[Tay\].
The algebra $`𝒱(X)`$ was systematically considered in the works \[Rab, RRR, RRS1, RRS2\]; see especially \[RRS2\] where one may find references to other earlier papers. Although the authors emphasize the case $`X=^n`$, it is clear for us that their methods extend to many other groups. On the other hand, since they allow the functions $`\phi `$ to be Banach space valued, the applications of their theory cover directly the case of operators on $`L^2(^n)`$ for example (this involves a certain discretization technique). In particular, Theorems 2.4.2 and Corollary 2.4.28 from \[RRS2\] are much stronger than our next Proposition 6.3 in the case $`X=^n`$. Taking into account the wealth of informations and applications in connection to these question which may be found in \[RRS2, Chapters 2,4,5\], we decided to keep this section to a minimum, just to point out the special role of the algebra $`𝒱(X)`$ in the crossed product formalism.
More recently, the relevance of $`𝒱(X)`$ in questions related to the computation of the essential spectrum has been independently noticed in \[LaS, Man\].
We mention that the compactification $`\sigma (𝒱(X))`$ and the boundary $`\upsilon X=\delta (𝒱(X))`$ are called *Higson compactification* and *Higson corona* of $`X`$ and play an important rôle in recent questions of topology, $`C^{}`$-algebras, $`K`$-theory, etc. \[Ro1, Ro2\].
Finally, we note that a non-abelian version of $`𝒱(X)`$ appears in a natural way in the spectral analysis of Schrödinger operators on a tree $`X`$, see \[GG1\].
We now give an intrinsic description of the crossed product $`V(X)=𝒱(X)X`$ and a more specific decomposition of the essential spectrum.
###### Proposition 6.3
We have
(6.5)
$$V(X)=\{TC(X)\varkappa .T𝒞_0(X^{}),\varkappa \delta X\}.$$
If $`TV(X)`$ then the map $`\varkappa \varkappa .T𝒞_0(X^{})`$ is norm continuous, hence (5.6) takes the more precise form
(6.6)
$$V(X)/K(X)𝒞(\upsilon X;𝒞_0(X^{})).$$
In particular, $`\mathrm{}(T)=\{\varkappa .T\varkappa \upsilon X\}𝒞_0(X^{})`$ is a compact set. If $`H`$ is a normal element of $`V(X)`$ or is an observable affiliated to $`V(X)`$ then:
(6.7)
$$\sigma _{\mathrm{ess}}(H)=_{\varkappa \upsilon X}\sigma (\varkappa .H)=_{K\mathrm{}(H)}\sigma (K).$$
Proof: To show the inclusion $``$ in (6.5) and the norm continuity of the map $`\varkappa \varkappa .T𝒞_0(X^{})`$ it suffices to consider $`T=\phi (Q)\psi (P)`$ with $`\phi 𝒱(X)`$ and $`\psi 𝒞_0(X^{})`$. But then $`\varkappa .T=\varkappa (\phi )\psi (P)`$ and these facts become obvious. Note that the compactness of the set $`\mathrm{}(T)`$ implies that the union $`_{T\mathrm{}(T)}\sigma (T)`$ is closed, hence (6.7) is true. It remains to show the inclusion $``$ in (6.5). Since $`\varkappa .(V_k^{}TV_k)=V_k^{}(\varkappa .T)V_k`$ and $`(\varkappa .T)U_x=\varkappa .(TU_x)`$ and since $`𝒞_0(X^{})`$ is stable under the automorphism generated by $`V_k`$ and under multiplication by $`U_x`$, it is clear that the right hand side of (6.5) satisfies Landstad’s conditions. Hence Theorem 3.7 shows that it suffices to prove that if $`\phi 𝒞(X)`$ has the property $`(\varkappa .\phi )(Q)\psi (P)𝒞_0(X^{})`$ for all $`\psi 𝒞_0(X^{})`$ and all $`\varkappa \delta X`$, then $`\phi 𝒱(X)`$. Thus it suffices to show that if $`\xi 𝒞(X)`$ and $`\xi (Q)\psi (P)𝒞_0(X^{})`$ for all $`\psi 𝒞_0(X^{})`$, then $`\xi `$ is a constant. But we have
$$\xi (Q)\psi (P)=U_x\xi (Q)\psi (P)U_x^{}=\xi (x+Q)\psi (P)$$
hence $`(\xi (Q)\xi (x+Q))\psi (P)=0`$ for all $`\psi `$, so $`\xi (Q)=\xi (x+Q)`$ for all $`x`$.
Remark: If the reader has any difficulty in proving that the union in (6.7) is closed, he should look at the proof of \[DG2, Theorem 2.10\].
###### Remark 6.4
$`V(X)`$ is the largest crossed product $`A`$ such that $`A/K(X)`$ is abelian. Indeed, $`A/K(X)_{\varkappa \delta (𝒜)}A_\varkappa `$ by (5.6) and the $`A_\varkappa `$ are crossed products, so $`A_\varkappa `$ is abelian if and only if $`A_\varkappa =\{0\}`$ or $`A_\varkappa =𝒞_0(X^{})`$.
###### Remark 6.5
The observables affiliated to $`𝒞_0(X^{})`$ are functions of momentum, so that it is natural to call them *free Hamiltonians*. Then we may describe in physical terms $`V(X)`$ as the largest $`C^{}`$-algebra of energy observables such that if $`H`$ is affiliated to it then all its localizations at infinity are free Hamiltonians.
###### Remark 6.6
We reconsider here the question of Remark 5.17 for $`𝒜=𝒱(X)`$. If $`\varkappa \delta X`$ then $`\tau _\varkappa :𝒱(X)`$ is just the character associated to $`\varkappa `$ and so if $`\chi \delta X`$ then $`\tau _\chi \tau _\varkappa \phi =\tau _\varkappa \phi \tau _\chi \phi =\tau _\varkappa \tau _\chi \phi `$ in general.
6.3. More remarks on filters. The following general remarks will be useful in the next subsections. Let $`Y`$ be a closed subspace of $`X`$ (thus $`KY`$ is compact for each compact $`KX`$). If $`𝔣`$ is a filter on $`Y`$ then $`𝔣`$ can be seen as a filter basis on $`X`$ and we shall denote (just for a moment) by $`𝔣^X`$ the filter on $`X`$ that it generates (this is the set of subsets of $`X`$ which contain a set from $`𝔣`$). The map $`𝔣𝔣^X`$ is an injective map from the set of filters on $`Y`$ onto the set of filters on $`X`$ which contain $`Y`$. Indeed, we have $`𝔣=\{FYF𝔣^X\}`$. It is also clear that if $`\varkappa `$ is an ultrafilter on $`Y`$ then $`𝔣^X`$ is also an ultrafilter. Finally, if $`𝔣`$ is finer than Fréchet on $`Y`$ then $`𝔣^X`$ is finer than Fréchet on $`X`$.
Since $`Y𝔣^X`$, if $`T:XZ`$ then $`lim_{𝔣^X}T`$ exists if and only if $`lim_𝔣T|_Y`$ exists and then they are equal.
From now on we shall not distinguish $`𝔣^X`$ from $`𝔣`$, so we use the same notation $`𝔣`$ for both. In particular, we get natural embeddings
(6.8)
$$\gamma Y\gamma X\text{and}\delta Y\delta X.$$
It is convenient to understand this when the ultrafilters are interpreted as characters. We have an obvious embedding $`\mathrm{}_{\mathrm{}}(Y)\mathrm{}_{\mathrm{}}(X)`$ so each character of $`\mathrm{}_{\mathrm{}}(X)`$ gives a character of $`\mathrm{}_{\mathrm{}}(Y)`$ by restriction, and reciprocally, each character of $`\mathrm{}_{\mathrm{}}(Y)`$ has a canonical extension to a character of $`\mathrm{}_{\mathrm{}}(X)`$, namely $`\varkappa (\phi ):=\varkappa (\phi \text{1}_Y)`$. Thus:
$$\gamma Y=\{\varkappa \gamma X\varkappa (Y)=1\}\text{and}\delta Y=\gamma Y\delta X.$$
It is easy to see now that $`\gamma Y`$ is a *clopen* subset of $`\gamma X`$, equal to the closure of $`Y`$ in $`\gamma X`$.
One says that a filter on a topological space is *convergent* to some point $`x`$ if it is finer than the filter of neighborhoods of $`x`$. *Any ultrafilter on a compact space is convergent*. This is easily seen to be equivalent to any of the usual definitions of compactness \[Bou, Chapter 1, §9\].
It is easy now to understand the universal property of $`\gamma X`$, cf. page 2. We first observe that $`\gamma `$ should be considered as a functor from the category of sets into the category of compact spaces. Indeed, if $`X,Y`$ are sets and $`\theta :XY`$ then it is obvious how to define $`\gamma \theta :\gamma X\gamma Y`$ if ultrafilters are thought as characters: note first that $`\phi \phi \theta `$ is a morphism $`\theta ^{}:\mathrm{}_{\mathrm{}}(Y)\mathrm{}_{\mathrm{}}(X)`$ and then if $`\varkappa \gamma X`$ define $`\gamma \theta (\varkappa )`$ as the character of $`\mathrm{}_{\mathrm{}}(Y)`$ given by $`\gamma \theta (\varkappa )=\varkappa \theta ^{}`$. The continuity of $`\gamma \theta `$ is clear.
Now assume $`Y`$ is a compact topological space. The only thing we need to accept is that $`\sigma (𝒞(Y))=Y`$, this is not difficult to prove directly. Then we have a natural continuous map $`\gamma Y\chi \chi _{\mathrm{}}Y`$ which associates to a character $`\chi `$ of $`\mathrm{}_{\mathrm{}}(Y)`$ its restriction to $`𝒞(Y)`$. In fact, the ultrafilter $`\chi `$ is convergent and $`\chi _{\mathrm{}}`$ is just its limit. Finally, $`\varkappa \gamma \theta (\varkappa )_{\mathrm{}}`$ is the unique extension of $`\theta `$ to a continuous map $`\gamma XY`$.
6.4. Sparse sets. From the point of view of the complexity of the interactions, the algebra of interactions that one should consider next is
(6.9)
$$𝒜=\{\phi 𝒞(X)\varkappa .\phi 𝒞_{\mathrm{}}(X),\varkappa \delta X\}.$$
The corresponding algebra of energy observables is
(6.10)
$$A=𝒜X=\{TC(X)\varkappa .TT(X),\varkappa \delta X\}.$$
Thus $`A`$ is the largest $`C^{}`$-algebra of energy observables such that all the localizations at infinity of a Hamiltonian $`H`$ affiliated to it are two-body Hamiltonians. We shall leave for the second part of our work the study of the algebra (6.10) and we shall consider here only subalgebras corresponding to Klaus type potentials.
###### Remark 6.7
The algebra $`𝒜`$ defined by (6.9) is characterized by $`𝒜_\varkappa =𝒞_{\mathrm{}}(X)`$ for each $`\varkappa `$, hence contains $`𝒞_{\mathrm{}}(X)`$ and is stable under all the morphisms $`\tau _\varkappa `$. It is also clear that $`\tau _\chi \tau _\varkappa \phi `$ and is distinct from $`\tau _\varkappa \tau _\chi \phi `$ in general, cf. Remark 5.17.
M. Klaus discovered in \[Kla\] the following class of Hamiltonians with nontrivial essential spectrum. Let $`L`$ be a discrete set such that the distance between two successive points of $`L`$ tends to infinity when we approach infinity. For each $`lL`$ let $`V_lL^1()`$ real such that $`V_l_{L^1}A`$ and $`\text{supp}V_l[A,A]`$ for a fixed finite $`A`$. Denote $`H=P^2+_lV_l(Ql)`$ and $`H_l=P^2+V_l(Q)`$. Then the description of $`\sigma _{\mathrm{ess}}(H)`$ given in \[Kla\] is equivalent to:
(6.11)
$$\sigma _{\mathrm{ess}}(H)=_{F(L)}\overline{}_{lF^c}\sigma (H_l)$$
where $`(L)`$ is the set of finite subsets of $`L`$ and $`F^c=LF`$. One of the main examples in \[GI2, GI4\] consisted in an algebraic treatment of this example, treatment based on the construction of a $`C^{}`$-algebra to which operators like $`H`$ are affiliated. We recall below the definition of this type of algebras and then we shall give a description of $`\sigma _{\mathrm{ess}}(H)`$ for the operators affiliated to them which is more in the spirit of Theorem 1.1 (description which also appears in \[GI2, GI4\] but which is deduced there by very different means).
If $`L,\mathrm{\Lambda }`$ are subsets of $`X`$ we denote $`L_\mathrm{\Lambda }=L+\mathrm{\Lambda }`$ and $`L_\mathrm{\Lambda }^c=XL_\mathrm{\Lambda }`$. If $`L`$ has the property $`L_\mathrm{\Lambda }X`$ for each compact $`\mathrm{\Lambda }`$ then we associate to it the filter
(6.12)
$$𝔣_L=\{AXAL_\mathrm{\Lambda }^c\text{ for some compact }\mathrm{\Lambda }X\}.$$
This is clearly a translation invariant filter finer than the Fréchet filter and such that $`𝔣_L^{}=𝔣_L`$. Thus
(6.13)
$$𝒞_L(X)=\{\phi 𝒞(X)lim_{𝔣_L}\phi \text{ exists }\}$$
is an algebra of interactions on $`X`$. An intrinsic description of the corresponding algebra of Hamiltonians $`C_L(X)`$ follows immediately from the results of Subsection 6. Let
(6.14)
$$\delta _LX=\delta (𝒞_L(X))=\sigma (𝒞_L(X))X$$
be the boundary of $`X`$ in the compactification associated to $`𝒞_L(X)`$. We recall that $`\delta _LX`$ is a quotient of $`\delta X`$. We set
(6.15)
$$\mathrm{}_L=\{\varkappa \gamma X\varkappa 𝔣_L\}=\{\varkappa \gamma XL_\mathrm{\Lambda }^c\varkappa \text{ if }\mathrm{\Lambda }X\text{ is compact }\}.$$
This is a compact subset of $`\delta X`$ and if $`\varkappa \mathrm{}_L`$ then $`\varkappa (\phi )=lim_{𝔣_L}\phi `$, so that $`\mathrm{}_L`$ gives just a point in $`\delta _LX`$. The problem that remains to be solved is the description of the other points of $`\delta _LX`$.
In this subsection we consider only the case when $`L`$ is a *sparse* set, in the following sense: $`L`$ is locally finite and for each compact set $`\mathrm{\Lambda }`$ there is a co-finite set $`ML`$ (i.e. such that $`LM`$ is finite) with the following property: if $`mM`$ and $`lL`$, $`lm`$, then $`(m+\mathrm{\Lambda })(l+\mathrm{\Lambda })=\mathrm{}`$.
With the conventions made in Subsection 6.6, we have $`\delta L\delta X`$, more explicitly for $`\varkappa \delta L`$ and $`\phi \mathrm{}_{\mathrm{}}(X)`$ we set
$$\varkappa (\phi )\varkappa (\phi \text{1}_L)=\underset{l\varkappa }{lim}\phi (l).$$
Below we use the symbol $``$ to denote disjoint union of sets.
###### Lemma 6.8
Let $`\theta :X\times \delta L\delta _LX`$ be defined by $`\theta (x,\varkappa )=\varkappa \tau _x`$. Then $`\theta `$ is injective and its range is $`\delta _LX\{\mathrm{}_L\}`$, which gives us an identification
(6.16)
$$\delta _LX(X\times \delta L)\{\mathrm{}_L\}.$$
Proof: We set $`\theta (x,\varkappa )=\theta _{x,\varkappa }`$ and note the more explicit formula
$$\theta _{x,\varkappa }(\phi )=\underset{l\varkappa }{lim}\phi (l+x).$$
We first prove that $`\theta `$ is injective. It is clearly sufficient to show that if $`xX`$ and $`\varkappa ,\chi \delta L`$ are such that $`\varkappa (x.\phi )=\chi (\phi )`$ for all $`\phi 𝒞_L`$, then $`x=0`$ and $`\varkappa =\chi `$. Let $`ML`$ such that $`\varkappa (M)=1`$. Since $`\varkappa `$ is finer than the Fréchet filter, $`M`$ is infinite and $`\varkappa (N)=1`$ if $`N`$ is a co-finite subset of $`M`$. Let $`\mathrm{\Lambda }X`$ be compact and such that $`0,x\mathrm{\Lambda }`$. Eliminating if needed a finite number of points from $`M`$, we may assume that $`(LM)M_\mathrm{\Lambda }=\mathrm{}`$ and $`M_\mathrm{\Lambda }=_{lM}(l+\mathrm{\Lambda })`$. Choose $`\phi 𝒞_0(X)`$ with $`\text{supp}\phi \mathrm{\Lambda }`$ and let us define $`\phi _M=_{lM}\tau _l\phi `$. Then:
$$(\text{1}_Lx.\phi _M)(y)=_{lM}\text{1}_M(y)(\tau _{xl}\phi )(y)=_{lM}\text{1}_M(y)\phi (x+yl).$$
In the sum from the right hand side the terms are zero unless $`l,yM`$ and $`x+yl+\mathrm{\Lambda }`$; but this implies $`l=y`$ because $`x\mathrm{\Lambda }`$. We get $`\text{1}_Lx.\phi _M=\text{1}_M\phi (x)`$ and so, by choosing $`\phi `$ such that $`\phi (x)0`$, we see that
$$\varkappa (x.\phi _M)=\varkappa (\text{1}_Lx.\phi _M)=\varkappa (\text{1}_M)\phi (x)=\phi (x)0.$$
Similarly $`\text{1}_L\phi _M=\text{1}_M\phi (0)`$ and so $`\chi (\phi _M)=\chi (\text{1}_M\phi _M)=\chi (\text{1}_M)\phi (0)`$. If $`x0`$ we may choose $`\phi `$ such that $`\phi (0)=0`$ and we see that $`\varkappa (x.\phi )\chi (\phi )`$ for some $`\phi 𝒞_L`$. If $`x=0`$ but $`\varkappa \chi `$ then $`M`$ can be chosen such that $`\chi (M)=0`$ (because $`\varkappa `$ and $`\chi `$ are distinct ultrafilters) hence again $`\varkappa (x.\phi )\chi (\phi )`$ for some $`\phi 𝒞_L`$. This proves the injectivity of the map $`\theta `$.
Now we show that for any $`\chi \delta X`$ such that $`\chi \mathrm{}_L`$ there is $`(x,\varkappa )X\times \delta L`$ such that $`\chi (\phi )=\varkappa (x.\phi )`$ for all $`\phi 𝒞_L`$. Since $`\chi `$ is not finer than $`𝔣_L`$, there is a compact set $`\mathrm{\Lambda }X`$ such that $`L_\mathrm{\Lambda }^c\chi `$. But $`\chi `$ is an ultrafilter, so $`L_\mathrm{\Lambda }\chi `$. Since $`\chi `$ is finer than the Fréchet filter, there is $`ML`$ such that $`\chi (M_\mathrm{\Lambda })=1`$ and
$$M_\mathrm{\Lambda }=_{lM}(l+\mathrm{\Lambda })M\times \mathrm{\Lambda }.$$
The sets $`FM_\mathrm{\Lambda }`$ with $`\chi (F)=1`$ form a basis for $`\chi `$ and each such $`F`$ can be uniquely written as a disjoint union $`F=_{lN}(l+F(l))`$ with $`NM`$ and $`F(l)\mathrm{\Lambda }`$ non empty sets. We define surjective maps $`\pi _M:M_\mathrm{\Lambda }M`$ and $`\pi _\mathrm{\Lambda }:M_\mathrm{\Lambda }\mathrm{\Lambda }`$ with the help of the identification $`M_\mathrm{\Lambda }M\times \mathrm{\Lambda }`$. The image $`\varkappa =\pi _M(\chi )`$, i.e. the filter of subsets of $`M`$ generated by the $`\pi _M(F)`$ with $`F\chi `$, is obviously an ultrafilter on $`M`$, hence on $`L`$, finer then the Fréchet filter. Similarly, $`\pi _\mathrm{\Lambda }(\chi )`$ is an ultrafilter on $`\mathrm{\Lambda }`$, which is a compact space, hence $`\pi _\mathrm{\Lambda }(\chi )`$ converges to some point $`x\mathrm{\Lambda }`$. If $`F`$ is as above then $`\pi _M(F)=N`$ and $`\pi _\mathrm{\Lambda }(F)=_{lN}F(l)`$ and the families of these sets are bases for the filters $`\varkappa `$ and $`\pi _\mathrm{\Lambda }(\chi )`$ respectively. In particular, since $`\pi _\mathrm{\Lambda }(\chi )`$ is finer than the filter of neighborhoods of $`x`$, for each neighborhood $`V`$ of $`x`$ there is $`F`$ such that $`_{lN}F(l)V`$.
We prove now that $`\chi (\phi )=\varkappa (x.\phi )`$ if $`\phi 𝒞_L`$. We have $`\chi (\phi )=lim_\chi \phi `$, thus for each $`\epsilon >0`$ there is $`F\chi `$ as above such that $`|\phi (y)\chi (\phi )|<\epsilon `$ for all $`yF`$. Thus $`|\phi (l+\lambda )\chi (\phi )|<\epsilon `$ for all $`lN`$ and $`\lambda F(l)`$. On the other hand, $`\phi `$ being uniformly continuous, there is a neighborhood $`V`$ of $`x`$ such that $`|\phi (l+\lambda )\phi (l+x)|<\epsilon `$ for all $`lN`$ and $`\lambda V`$. By what we said above, the preceding $`F`$ may be chosen such that $`_{lN}F(l)V`$. Hence we get $`|\phi (l+x)\chi (\phi )|<2\epsilon `$ for all $`lN`$. Since $`N\varkappa `$ and $`\epsilon >0`$ is arbitrary, this shows that $`lim_{l\varkappa }\phi (l+x)=\chi (\phi )`$.
Remark: It is easy to show that $`\delta _LX`$ is homeomorphic with $`(X\times \delta L)\{\mathrm{}_L\}`$, thought as the one point compactification of $`X\times \delta L`$, but we do not need this.
In the next lemma we use the notation of Definition 5.15. Let $`\varkappa `$ be a point in $`\delta X`$.
###### Lemma 6.9
If $`\varkappa \mathrm{}_L`$ then $`𝒞_L(X)_\varkappa =`$. If $`\varkappa \mathrm{}_L`$, then $`𝒞_L(X)_\varkappa =𝒞_{\mathrm{}}(X)`$.
Proof: If $`\varkappa \mathrm{}_L`$ then $`\varkappa .\phi (x)=\varkappa (x.\phi )=lim_{𝔣_L}x.\phi =lim_{𝔣_L}\phi `$ because $`𝔣_L`$ is translation invariant. Thus $`\varkappa .\phi `$ in this case. Now let $`\chi \mathrm{}_L`$. It suffices then to show that $`\chi .\phi 𝒞_0(X)`$ if $`lim_{𝔣_L}\phi =0`$ and by an easy density argument we see that it suffices to assume that $`\text{supp}\phi L_K`$ for a compact subset $`K`$ of $`X`$. If $`\varkappa ,x`$ are such that $`\theta (x,\varkappa )=\chi `$ then
$$\chi .\phi (y)=\chi (y.\phi )=\varkappa (x.(y.\phi ))=\varkappa ((x+y).\phi )=\underset{l\varkappa }{lim}\phi (l+x+y).$$
But if $`zK`$ then there is $`ML`$ co-finite such that $`l+zL_K`$ if $`lM`$, and then $`\phi (l+z)=0`$ for all such $`l`$, and so $`lim_{l\varkappa }\phi (l+z)=0`$. Hence $`\text{supp}\chi .\phi Kx`$.
To finish the proof it remains to show that if $`\chi \mathrm{}_L`$ and $`\xi 𝒞_{\mathrm{}}(X)`$, then there is $`\phi 𝒞_L`$ such that $`\chi .\phi =\xi `$. It suffices to show this under the assumption that $`\xi `$ has compact support. Then it suffices to take $`\phi =\xi _L=_{lL}\tau _l\xi `$
###### Lemma 6.10
If $`\phi 𝒞_L(X)`$ the map $`\delta L\varkappa \varkappa .\phi 𝒞_{\mathrm{}}(X)`$ is norm continuous.
Proof: By a density argument, it suffices to show this for $`\text{supp}\phi M_\mathrm{\Lambda }`$, where $`ML`$ is a co-finite set and $`\mathrm{\Lambda }X`$ is a compact set such that $`M_\mathrm{\Lambda }=_{lM}(l+\mathrm{\Lambda })`$. If $`lM`$ let $`\phi _l`$ be the function defined by $`\phi _l(x)=\phi (l+x)`$ for $`x\mathrm{\Lambda }`$ and $`\phi _l(x)=0`$ otherwise. Then clearly $`\phi _l𝒞_0(X)`$, $`\text{supp}\phi _l\mathrm{\Lambda }`$, and the family $`\{\phi _l\}_{lM}`$ is equicontinuous. Thus the set $`\{\phi _llM\}`$ is relatively compact in $`𝒞_0(X)`$. From the universal property of $`\gamma M`$, cf. page 2, there is a unique continuous map $`\gamma M\varkappa \phi _\varkappa 𝒞_0(X)`$ such that $`\phi _\varkappa =\phi _l`$ for all $`lM`$. Since $`\delta M=\delta L`$, it suffices to show that $`\phi _\varkappa =\varkappa .\phi `$ if $`\varkappa \delta M`$. But we have
$$\varkappa .\phi (x)=\underset{l\varkappa }{lim}\phi (l+x)=\underset{l\varkappa }{lim}\phi _l(x)=\phi _\varkappa (x)$$
because $`\gamma M\varkappa \phi _\varkappa (x)`$ is continuous.
Putting all this together we obtain, for the algebra $`A=C_L(X)`$, an improvement of Theorem 5.16. If $`TC_L(X)`$ then, according to that theorem, we have a continuous map $`\sigma (𝒞_L(X))\varkappa \varkappa .TC_s(X)`$ which induces an embedding
(6.17)
$$C_L(X)/K(X)_{\varkappa \delta _LX}C_L(X)_\varkappa .$$
From Lemma 6.9 we see that the localization $`C_L(X)_\varkappa =𝒞_L(X)_\varkappa X`$ at $`\varkappa `$ is
(6.18)
$$C_L(X)_\varkappa =\{\begin{array}{cc}𝒞_0(X^{})\text{if}\varkappa =\mathrm{}_L,\hfill & \\ T(X)\text{if}\varkappa \mathrm{}_L.\hfill & \end{array}$$
Here $`\delta _LX`$ is represented as in (6.16). We identify $`\delta L\{0\}\times \delta LX\times \delta L`$ and we simplify the relation (6.17) by taking into account the discussion made on page 5.7. First, since $`(\varkappa \tau _x).T=U_xTU_x^{}`$, it suffices to restrict the product to the set $`\delta L\{\mathrm{}_L\}`$. Second, we note that the contribution of the point $`\mathrm{}_L`$ is already covered by the other ones. Indeed, this follows from the easy to check relation <sup>2</sup><sup>2</sup>2 Note that this is related to Remark 5.17: we have $`\tau _\varkappa \tau _\chi =\tau _\chi \tau _\varkappa =\mathrm{}_L`$ on $`𝒞_L`$.
$$\mathrm{}_L(\phi )=\mathrm{}_L(\tau _\varkappa (\phi ))\text{for all}\phi 𝒞_L(X),$$
which implies $`𝒥_\varkappa 𝒥_\mathrm{}_L`$, and from Lemma 5.20. Finally, we get:
###### Theorem 6.11
If $`TC_L(X)`$ and $`\varkappa \delta L`$ then $`lim_{l\varkappa }U_lTU_l^{}=\varkappa .T\tau _\varkappa (T)`$ exists in $`C_s(X)`$ and belongs to $`T(X)`$. The map $`\delta L\varkappa \varkappa .TT(X)`$ is norm continuous. The maps $`\tau _\varkappa :C_L(X)T(X)`$ are surjective morphisms and the intersection of their kernels is $`K(X)`$, which gives us a canonical embedding
(6.19)
$$C_L(X)/K(X)𝒞(\delta L;T(X)).$$
If $`H`$ a normal operator in $`C_L(X)`$ or an observable affiliated to $`C_L(X)`$, then
(6.20)
$$\sigma _{\mathrm{ess}}(H)=_{\varkappa \delta L}\sigma (\varkappa .H).$$
The last assertion follows from the norm continuity of the map $`\varkappa \varkappa .T`$.
###### Remark 6.12
Theorem 6.11 has been obtained by rather different methods in \[GI2, GI4\], see for example Theorems 5.5 and 5.6 in \[GI4\]. The point is that in these references the quotient $`C_L(X)/K(X)`$ was computed directly and the notion of localization at infinity did not play a rôle. Our purpose here was only to show that Theorem 1.1 can be effectively used even in some rather complicated situations. Our arguments in this subsection may, in fact, serve as a model for other computations.
###### Remark 6.13
That the class of operators affiliated to $`C_L(X)`$ is quite large can be seen from the following result \[GI4, Theorem 6.1\]. Let $`X=^n`$ and denote $`^t`$ the Sobolev space of order $`t`$ and $`_t`$ the norm in $`B(^t,^t)`$. Let $`h:X`$ be a continuous function such that $`c^1(1+|k|)^{2s}|h(k)|c(1+|k|)^{2s}`$ for some constant $`c`$ and all large $`k`$, and denote $`H_0=h(P)`$. Let $`0t<s`$ reals, choose a sparse set $`LX`$ and let $`\{W_l\}_{lL}`$ be a family of symmetric operators $`W_l:^t^t`$ with the property $`sup_{lL}(1+|Q|)^\lambda W_l_t<\mathrm{}`$ for some $`\lambda >2n`$. Then the series $`_{lL}U_l^{}W_lU_lW`$ converges in the strong topology of $`B(^t,^t)`$. Let $`H=H_0+W`$, $`H_l=H_0+W_l`$ be the self-adjoint operators in $`L^2(X)`$ defined as form sums. Then $`H`$ is affiliated to $`C_L(X)`$. If $`\varkappa \delta L`$ then we also have $`\varkappa .H=lim_{l\varkappa }H_l`$ in norm resolvent sense.
###### Remark 6.14
The preceding arguments can be simplified and everything becomes an elementary exercise for the subalgebras of $`C_L(X)`$ corresponding to a finite number of types of bumps \[GI2, p. 548\]. The case of just one type is already interesting. More precisely, let $``$ be a finite partition of $`L`$ consisting of $`n`$ infinite sets and let $`𝒞_{}`$ be the set of $`\phi 𝒞_L`$ such that $`lim_{al\mathrm{}}\phi (l+x)a.\phi (x)`$ exists for each $`x`$ and for each $`a`$. Then $`\delta L`$ is replaced by the finite set $``$ and the Hamiltonians affiliated to $`C_{}`$ have (modulo translations) exactly $`n+1`$ localizations at infinity: a free one $`H_0^{}𝒞_0(X^{})`$ and a two body one $`a.H^{}T(X)`$ for each $`a`$. And $`\sigma _{\mathrm{ess}}(H)=_a\sigma (a.H)`$.
We make a final comment on the algebra $`𝒜`$ defined in (6.9). We saw that for any sparse set $`L`$ we have $`𝒞_L(X)𝒜`$. On the other hand, if $`M`$ is a second sparse set, then $`LM`$ is not sparse in general. However, the $`C^{}`$-algebra $`𝒞_{L,M}`$ generated by $`𝒞_L𝒞_M`$ is still included in $`𝒜`$. Note that to each Hamiltonian affiliated to $`C_L`$ one may associate in a canonical way a free Hamiltonian, this is the localization of $`H`$ at the point $`\mathrm{}_L`$. But this is not the case for Hamiltonians affiliated to $`A`$.
6.5. Grassmann algebras. We shall construct here $`C^{}`$-algebras canonically associated to finite dimensional vector spaces and which allow one to consider a very general version of $`N`$-body Hamiltonians. This algebras have first been pointed out in \[DG1\] and the spectral theory of the operators affiliated to them (essential spectrum and the Mourre estimate) has been studied in detail in \[DG2\]. Our approach here is rather different, the graded algebra structure so important in the quoted works does not play a big rôle anymore.
If $`Y`$ is a closed subgroup of a locally compact abelian group $`X`$ then $`X/Y`$ is also a locally compact abelian group and we have a continuous surjective group morphism $`\pi _Y:XX/Y`$. Then the map defined by $`\phi \phi \pi _Y`$ gives us a natural embedding $`𝒞(X/Y)𝒞(X)`$. In fact
(6.21)
$$𝒞(X/Y)=\{\phi 𝒞(X)y.\phi =\phi yY\}.$$
Note that we shall denote just $`0`$ the group $`\{0\}`$ and then $`𝒞(0)=𝒞_0(0)=`$, hence $`𝒞(X/X)=𝒞_0(X/X)=`$. On the other hand, if $`0YZX`$ are closed subgroups then $`X/Z(X/Y)/(Y/Z)`$ and we have natural maps
(6.22)
$$XX/YX/Z0$$
hence we get embeddings
(6.23)
$$𝒞(X/Z)𝒞(X/Y)𝒞(X).$$
In the rest of this subsection *we shall consider only finite dimensional real vector spaces*, although much of the theory can be extended to more general groups. We shall consider the algebra generated by the $`C^{}`$-subalgebras $`𝒞_0(X/Y)𝒞(X/Y)𝒞(X)`$. We recall that the *Grassmannian* $`𝔾(X)`$ is the set of all vector subspaces of $`X`$ and the *projective space* $`(X)`$ is the set of all one dimensional subspaces of $`X`$.
###### Definition 6.15
The *(classical) Grassmann algebra* of the vector space $`X`$ is the $`X`$-subalgebra $`𝒢(X)𝒞(X)`$ defined by
(6.24)
$$𝒢(X)=\text{norm closure of }_Y𝒞_0(X/Y)$$
where $`Y`$ runs over $`𝔾(X)`$. The *quantum Grassmann algebra* of $`X`$ is the $`C^{}`$-algebra $`G(X)B(X)`$ defined by
(6.25)
$$G(X)=𝒢(X)X=\text{norm closure of }_Y𝒞_0(X/Y)X.$$
###### Remarks 6.16
The fact that $`𝒢(X)`$ is a $`C^{}`$-algebra follows from the obvious relation
(6.26)
$$𝒞_0(X/Y)𝒞_0(X/Z)𝒞_0(X/(YZ)).$$
The second equality from (6.25) follows from Theorem 3.14.
Let $`G(X)`$ be the set of finite unions of strict vector subspaces of $`X`$:
$$G(X)=\{LX𝔽𝔾(X)\{X\}\text{ finite such that }L=_{Y𝔽}Y\}.$$
If $`L`$ is as above and $`\mathrm{\Lambda }X`$ is compact then $`L_\mathrm{\Lambda }=L+\mathrm{\Lambda }=_{Y𝔽}(Y+\mathrm{\Lambda })`$. Thus $`L_\mathrm{\Lambda }`$ is a closed set, $`L_\mathrm{\Lambda }X`$, and we have $`L_\mathrm{\Lambda }M_\mathrm{\Lambda }=(LM)_\mathrm{\Lambda }`$ and $`L_\mathrm{\Lambda }^{}L_{\mathrm{\Lambda }^{\prime \prime }}`$ if $`\mathrm{\Lambda }^{}\mathrm{\Lambda }^{\prime \prime }`$. This justifies the next definition.
###### Definition 6.17
The *Grassmann filter* $`𝔤=𝔤_X`$ on $`X`$ is the filter generated by the family of open sets $`L_\mathrm{\Lambda }^c=XL_\mathrm{\Lambda }`$ where $`L`$ runs over $`G(X)`$ and $`\mathrm{\Lambda }`$ over the set of compact subsets of $`X`$. If $`Y`$ is a subspace of $`X`$, then we denote also by $`𝔤_Y`$ the filter on $`X`$ generated by the Grassmann filter of $`Y`$.
Clearly, $`𝔤_X`$ is translation invariant, finer than the Fréchet filter, and $`𝔤_X^{}=𝔤_X`$. If $`X`$ is one dimensional, then $`𝔤_X`$ is just the Fréchet filter.
###### Remark 6.18
For $`LG(X)`$ we may consider the filter $`𝔣_L`$ defined as in (6.12). Then $`𝔤_X`$ is just the filter generated by $`_L𝔣_L`$. This can be expressed in other terms as follows: (1) $`𝔤_X`$ is the upper bound of the set of filters $`𝔣_L`$; (2) when seen as compact subsets of $`\delta X`$ (cf. page 2), $`𝔤_X`$ is the intersection of the compact sets $`𝔣_L`$.
###### Remark 6.19
If we equip $`X`$ with an Euclidean norm $`||`$ and denote $`\pi _Y`$ the orthogonal projection onto $`Y^{}X/Y`$, then $`\delta _L(x)\text{dist}(x,L)=\mathrm{min}_{Y𝔽}|\pi _Yx|`$ (with $`L,𝔽`$ as before). Then the sets $`L_r^c=\{xX\delta _L(x)>r\}`$, with $`LG(X)`$ and $`r>0`$ real, form a basis of the filter $`𝔤_X`$. Note that $`L`$ has empty interior and if $`x`$ is outside it then $`\delta _L(tx)=|t|\delta _L(x)\mathrm{}`$ as $`t\mathrm{}`$.
If $`𝔣`$ is a filter on a set $`S`$ and $`\pi `$ is a map from $`S`$ to a locally compact space $`T`$ then $`lim_𝔣\pi =\mathrm{}`$ means: for each compact $`KT`$ there is $`F𝔣`$ such that $`\pi (F)K=\mathrm{}`$.
###### Lemma 6.20
Let $`Y,Z𝔾(X)`$. If $`YZ`$ then $`lim_{𝔤_Y}\pi _Z=0`$. If $`YZ`$ then $`lim_{𝔤_Y}\pi _Z=\mathrm{}`$.
Proof: Since $`Y𝔤_Y`$ the above limits involve only the restriction of $`\pi _Z`$ to $`Y`$. In the first case, if $`yY`$ then $`\pi _Z(y)=0`$, so the assertion is clear. If $`YZ`$ then $`E=YZ`$ is a strict subspace of $`Y`$. Let $`E^{}`$ be a supplementary subspace for $`E`$ in $`Y`$. Then $`\pi _Z:E^{}X/Z`$ is injective, hence if $`KX/Z`$ is compact then the set $`\mathrm{\Lambda }`$ of $`yE^{}`$ such that $`\pi _ZyK`$ is a compact in $`E^{}`$ and thus in $`Y`$. If $`yYE_\mathrm{\Lambda }𝔤_Y`$ then $`y=e+e^{}`$ with $`eE`$ and $`e^{}E^{}\mathrm{\Lambda }`$ so $`\pi _Zy=\pi _Ze^{}K`$.
###### Corollary 6.21
If $`\phi 𝒢(X)`$ and $`Y𝔾(X)`$ then $`lim_{𝔤_Y}\phi `$ exists. If $`\phi =_Z\phi ^Z`$ is a finite sum of $`\phi ^Z𝒞_0(X/Z)`$ , then $`lim_{𝔤_Y}\phi =_{ZY}\phi ^Z(0)`$.
We see that each filter $`𝔤_Y`$ defines a character of $`𝒢(X)`$ and we could proceed as in the proof of Lemma 6.8 and describe $`\delta _𝔾X\delta (𝒢(X))`$ in terms of couples $`(Y,y)`$ with $`yX/Y`$. We shall not do it explicitly, but this is hidden in what follows. We only note that the rôle of $`\mathrm{}_L`$ is now played by $`𝔤_X`$.
###### Proposition 6.22
For each $`\phi 𝒢(X)`$ the limit
(6.27)
$$\tau _Y\phi =\underset{y𝔤_Y}{lim}y.\phi $$
exists locally uniformly on $`X`$. If $`\phi `$ is a finite sum $`\phi =_Z\phi ^Z`$ with $`\phi ^Z𝒞_0(X/Z)`$, then $`\tau _Y\phi =_{ZY}\phi ^Z`$. If $`Y(X)`$, $`yY0`$, then $`\tau _Y\phi (x)=lim_t\mathrm{}\phi (x+ty)`$.
Proof: We have to show that $`lim_{y𝔤_Y}\phi (x+y)`$ exists locally uniformly in $`x`$. But this is an immediate consequence of the Corollary 6.21 and Lemma 2.2.
According to the conventions we made at the beginning of this subsection, we have $`𝒢(X/Y)𝒞(X/Y)𝒞(X)`$ if $`Y`$ is a subspace of $`X`$.
###### Proposition 6.23
We have $`𝒢(X/Y)𝒢(X)`$. Moreover, there is a unique morphism $`\tau _Y:𝒢(X)𝒢(X/Y)`$ such that $`\tau _Y`$ is a projection (in the sense of linear spaces). The map $`\tau _Y`$ is given by (6.27). If $`YZX`$ then
(6.28)
$$=𝒢(0)𝒢(X/Z)𝒢(X/Y)𝒢(X)$$
and $`\tau _Y\tau _Z=\tau _Z\tau _Y=\tau _Z`$. More generally, for any $`Y,Z𝔾(X)`$ we have
(6.29)
$$𝒢(X/Y)𝒢(X/Z)=𝒢(X/(Y+Z))$$
and $`\tau _Y\tau _Z=\tau _Z\tau _Y=\tau _{Y+Z}`$.
Proof: The algebra $`𝒢(X/Y)`$ is generated the $`𝒞_0((X/Y)/E)`$ with $`EX/Y`$ subspace. If $`Z=\pi _Y^1(E)`$ then $`YZX`$, $`E=Z/Y`$ and $`(X/Y)/EX/Z`$ allows us to identify $`𝒞_0((X/Y)/E)=𝒞_0(X/Z)`$ and thus to get the first assertion of the proposition. Observe that
$`𝒢(X/Y)`$ $`=`$ $`\text{norm closure of }_{ZY}𝒞_0(X/Z)`$
$`=`$ $`\{\phi 𝒢(X)y.\phi =\phi yY\}.`$
The other assertions of the proposition are easy to check.
###### Proposition 6.24
If $`\phi 𝒢(X)`$ and $`\tau _Y\phi =0`$ for all $`Y(X)`$, then $`\phi 𝒞_0(X)`$.
Proof: This follows from Theorem 3.2 and Lemma 4.1 of \[DG1\], but we give a self-contained proof here. Consider first a finite set $`𝔽𝔾(X)`$ which is stable under intersections and such that $`0𝔽`$ and let $`𝒜=_{Y𝔽}𝒞_0(X/Y)`$. Then $`𝒜`$ is a $``$-algebra because of (6.26) and $`𝒞_0(X)𝒜`$. Clearly $`\tau _Y\phi \phi `$ for all $`Y𝔽,\phi 𝒜`$. Let us write $`\phi =_Y\phi ^Y`$ with $`\phi ^Y𝒞_0(X/Y)`$. From Proposition 6.22 we get $`\tau _Y\phi =_{ZY}\phi ^Z`$, so if $`Y`$ is a maximal element of $`𝔽`$ then $`\tau _Y\phi =\phi ^Y`$, hence $`\phi ^Y\phi `$. By induction,we easily see that there is a constant $`c`$ such that
$$\phi ^Yc\phi \text{for all}Y𝔽,\phi 𝒜.$$
This clearly implies that $`𝒜`$ is a $`C^{}`$-algebra and that $`_{Y𝔽}𝒞_0(X/Y)`$ is a topological direct sum. If $`\phi `$ is as above and $`\tau _Y\phi =0`$ for all $`Y0`$ then $`_{ZY}\phi ^Z=0`$ if $`Y0`$ hence, the sum being direct, we get $`\phi ^Z=0`$ for all $`Z0`$, thus $`\phi 𝒞_0(X)`$.
It follows that the map $`\phi (\tau _Y\phi )_{Y0}`$ is a morphism from $`𝒜`$ into $`_{Y0}𝒢(X/Y)`$ with kernel equal to $`𝒞_0(X)`$. In particular, the induced map $`𝒜/𝒞_0(x)_{Y0}𝒢(X/Y)`$ is an isometry, so that if $`\psi 𝒜`$ is such that $`\tau _Y\psi \epsilon `$ for all $`Y0`$ then there is $`\psi _0𝒞_0(X)`$ such that $`\psi \psi _02\epsilon `$ (just by definition of the quotient norm).
Let now $`\phi 𝒢(X)`$ such that $`\tau _Y\phi =0`$ for all $`Y(X)`$. From Proposition 6.23 it follows that this property remains true for all $`Y𝔾(X),Y0`$. From the definition (6.24) it follows easily that for each $`\epsilon >0`$ there is $`𝒜`$ as above and there is $`\psi 𝒜`$ such that $`\phi \psi \epsilon `$. Then clearly we have $`\tau _Y\psi \epsilon `$ for all $`Y0`$, so by what we proved above there is $`\psi _0𝒞_0(X)`$ such that $`\psi \psi _02\epsilon `$, and hence $`\phi \psi _03\epsilon `$. This clearly implies $`\phi 𝒞_0(X)`$.
The next theorem is now an immediate consequence of Theorem 5.16, Propositions 6.23 and 6.24, and of Lemma 5.22. We denote
(6.31)
$$G_Y(X)=G(X/Y)X=\text{norm closure of }_{ZY}𝒞_0(X/Z)X.$$
We mention that we have non canonical isomorphisms $`G_Y(X)G(X/Y)𝒞_0(Y^{})`$.
###### Theorem 6.25
If $`TG(X)`$ and $`Y𝔾(X)`$ then $`\tau _YT=lim_{y𝔤_Y}U_xTU_x^{}`$ exists in $`C_s(X)`$ and belongs to $`G_Y(X)`$. The map $`\tau _Y:G(X)G_Y(X)`$ is a morphism and a linear projection and is uniquely characterized by these properties. We have $`\tau _Y\tau _Z=\tau _Z\tau _Y=\tau _{Y+Z}`$. If $`Y(X)`$ and $`yY,y0`$ then $`\tau _YT=lim_t\mathrm{}U_{ty}TU_{ty}^{}`$. We have $`TK(X)`$ if and only if $`\tau _YT=0`$ for all $`Y(X)`$, which gives us
(6.32)
$$G(X)/K(X)_{Y(X)}G_Y(X).$$
From (6.32) we get that the the essential spectrum of an observable $`H`$ affiliated to $`G(X)`$ is equal to the closure of the union $`\sigma (\tau _YH)`$ with $`Y(X)`$. But now we can prove more: as in the situations considered in Theorem 6.11 and Proposition 6.3, the union is already closed (although it is not finite, as in the usual $`N`$-body problem).
###### Theorem 6.26
If $`TG(X)`$ then $`\{\tau _YTY(X)\}`$ is a compact set in $`G(X)`$. In particular, if $`H`$ a normal operator in $`G(X)`$ or an observable affiliated to $`G(X)`$, then
(6.33)
$$\sigma _{\mathrm{ess}}(H)=_{Y(X)}\sigma (\tau _YH).$$
One should note that the map $`Y\tau _YT`$ is not continuous: if $`T𝒞_0(X/Z)X`$ then $`\tau _YT=T`$ if $`YZ`$ and $`\tau _YT=0`$ if $`YZ`$.
Theorem 6.26 is a corollary of Theorem 4.2 and Proposition 5.4 from \[DG2\]. We shall give below a slightly improved proof. Note that only some general properties of the lattice $`𝔾(X)`$ and of the graded algebra structure of $`G(X)`$ are really needed. The next two lemmas imply the first assertion of Theorem 6.26 (hence the second).
###### Lemma 6.27
If $`TG(X)`$ then for each $`Z𝔾(X),Z0`$ there is $`Y(X)`$ such that $`\tau _ZT=\tau _YT`$.
Proof: Let $`𝔼𝔾(X)`$ be countable. Then $`\{EZE𝔼,EZZ\}`$ is a countable set of strict subspaces of $`Z`$, so its union is not $`Z`$. Let $`Y(Z)`$ such that $`YEZ=0`$ if $`E`$ is in the preceding set. Then from $`E𝔼`$ and $`EY`$ we get $`EZ`$. Now if $`T_𝔼`$ is a finite sum $`_{E𝔼}T^E`$ with $`T^E𝒞_0(X/E)X`$ then clearly $`\tau _ZT=\tau _YT`$. Finally, if $`T`$ is arbitrary, then there is $`𝔼`$ as above such that $`T`$ be a norm limit of operators of the form $`T_𝔼`$, so we have $`\tau _ZT=\tau _YT`$.
###### Lemma 6.28
Let $`\{Y_n\}_{n0}`$ be a sequence of linear subspaces of $`X`$ and let us define $`Y=_{n0}_{mn}Y_m`$. If $`k`$ is the dimension of $`Y`$, then there is $`N`$ such that for all $`nN`$ and all $`TG(X)`$:
$$(\tau _{Y_n}\tau _Y)Tk\underset{mn}{sup}(\tau _{Y_n}\tau _{Y_m})T.$$
Proof: Since a decreasing sequence of subspaces is eventually constant, there is $`N`$ such that $`Y=_{mn}Y_m`$ for all $`nN`$. The dimension of $`Y`$ being $`k`$, for each $`nN`$ there are $`n<n_1<\mathrm{}<n_k`$ such that $`Y=Y_n+Y_{n_1}+\mathrm{}+Y_{n_k}`$. From Theorem 6.25 we get $`\tau _Y=\tau _{Y_n}\tau _{Y_{n_1}}\mathrm{}\tau _{Y_{n_k}}`$. Let $`𝒫=\tau _{Y_n}`$, $`𝒫_i=\tau _{Y_{n_i}}`$, and $`𝒫_i^{}=1𝒫_i`$. Then:
$$𝒫\tau _Y=𝒫[1𝒫_1\mathrm{}𝒫_k]=\underset{i=1}{\overset{k1}{}}𝒫𝒫_i^{}𝒫_{i+1}\mathrm{}𝒫_k+𝒫𝒫_k^{}.$$
Since the morphisms $`𝒫_i`$ commute, we get $`(𝒫\tau _Y)T_{i=1}^k𝒫𝒫_i^{}T`$. Now it suffices to note that $`𝒫𝒫_i^{}=𝒫(𝒫𝒫_i)`$.
Proof of Theorem 6.26: If $`\{\tau _{Y_n}T\}`$ is a norm Cauchy sequence and $`Y`$ is as in the Lemma 6.28 then $`(\tau _{Y_n}\tau _Y)T0`$. Observe that we do not have $`k=0`$ because this would imply $`Y_n=0`$ for large $`n`$. Thus we can use Lemma 6.27 and find $`E(X)`$ such that $`\tau _YT=\tau _ET`$, which proves the first assertion of the theorem.
###### Remark 6.29
The usual form of the HVZ theorem for $`N`$-body Hamiltonians follows easily from Theorem 6.26. Indeed, in the Agmon-Froese-Herbst formalism \[ABG\] one is given a finite lattice $``$ and an injective map $`aX_a𝔾(X)`$ such that $`X_{ab}=X_aX_b`$, $`X_{\mathrm{max}}=X`$ and $`X_{\mathrm{min}}=0`$. The $`N`$-body Hamiltonians are observables $`H`$ affiliated to the $`C^{}`$-algebra $`C=_a𝒞_0(X^a)XG(X)`$, where $`X^a=X/X_a`$. Let $`\tau _a=\tau _{X_a}`$, then $`\tau _a`$ is a morphism and a linear projection of $`C`$ onto the $`C^{}`$-subalgebra $`C_a=_{ba}𝒞_0(X^b)X`$. Let us set $`H_a=\tau _aH`$. Then (a generalized version of) the HVZ theorem says that
(6.34)
$$\sigma _{\mathrm{ess}}(H)=_a\sigma (H_a),$$
where $``$ is the set of atoms of $``$. To get this from (6.33), note that for each $`Y(X)`$ there is a smallest $`b`$ in $``$ such that $`YX_b`$, so we have $`YX_c`$ if and only if $`bc`$. Then for $`TC`$ we have $`\tau _YT=\tau _bT`$. On the other hand, there is an atom $`a`$ such that $`ab`$, and then $`\tau _bT=\tau _b\tau _aT`$. Thus $`\sigma (\tau _YT)\sigma (\tau _aT)`$. Reciprocally, if $`Z(X)`$ and $`ZX_a`$ then $`\tau _aT=\tau _a\tau _ZT`$ and so $`\sigma (\tau _aT)\sigma (\tau _ZT)`$.
###### Example 6.30
The simplest application of Theorem 6.26 is obtained by taking $`X=^n`$ and $`H=\mathrm{\Delta }+V(x)`$ where $`V𝒢(X)`$. Although simple, this situation is, however, not trivial because the union in (6.33) contains an infinite number of distinct terms in general. For example, the construction of $`V`$ may involve an infinite number of subspaces $`Y`$ whose union is dense in $`X`$.
###### Example 6.31
We show here that in an $`N`$-body type situation (i.e. involving only a finite number of subspaces $`Y`$) the class of Hamiltonians for which (6.34) applies is very large. We use the setting of Remark 6.29 and, to simplify notations, we equip $`X`$ with an Euclidean structure, so that $`X`$ is identified with $`X^{}`$ and $`X^a=X_a^{}`$. For real $`s`$ let $`^s`$ be the usual Sobolev spaces, set $`=^0=L^2(X)`$, and embed as usual $`^s^s`$ if $`s>0`$. Fix $`s>0`$ and denote $`_s`$ the norm in $`B(^s,^s)`$. Let $`h:X`$ be continuous and such that $`c^{}(1+|k|)^{2s}h(k)c^{\prime \prime }(1+|k|)^{2s}`$ outside a compact, for some constants $`c^{},c^{\prime \prime }`$. Then $`H(\mathrm{max}):=h(P)`$ is a self-adjoint operator with domain $`^{2s}`$ and form domain $`^s`$. Then for each $`a\mathrm{max}`$ let $`H(a):^s^s`$ be a symmetric continuous operator such that the following properties hold:
(1) $`U_xH(a)U_x^{}=H(a)`$ if $`xX_a`$,
(2) $`V_kH(a)V_k^{}H_a_s0`$ as $`k0`$ in $`X`$,
(3) $`(V_k1)H(a)_s0`$ as $`k0`$ in $`X^a`$,
(4) $`H_a:=_{ba}H(b)\mu h(P)\nu `$ as forms on $`^s`$, for some $`\mu ,\nu >0`$.
Then $`HH(\mathrm{min})`$ is affiliated to $`C`$ and $`\tau _aH=H_a`$, so (6.34) holds. See \[DG2, Theorem 4.6\] for the details of the computation and for more general results.
6.6. On the operators $`\varkappa .H`$. We observed after Theorem 1.2 that if $`H`$ is a self-adjoint operator affiliated to $`C(X)`$ then its localizations at infinity $`\varkappa .H`$ are not necessarily densely defined. We shall make in this subsection some comments on this question and we shall give conditions which allow one to compute $`\varkappa .H`$ directly in terms of $`x.H`$ and so to avoid considering the resolvent of $`H`$. This is not possible for $`H=h(P)+v(Q)`$ if the operator $`V=v(Q)`$ is not relatively bounded with respect to $`h(P)`$, so we shall consider here only more elementary situations which are of some physical interest.
To fix the ideas we consider here only the case $`X=^n`$ and take $`=L^2(X;𝐄)`$, where $`𝐄`$ a finite dimensional Hilbert space, cf. Section 4. For simplicity we consider only operators whose form domain is a Sobolev space $`^s`$ with $`s>0`$ (everything extends with no difficulty to hypoelliptic operators). Set $`k=(1+|k|^2)^{1/2}`$ and denote $`S(𝐄)`$ the space of symmetric operators on $`𝐄`$ and $`|S|`$ the absolute value of $`SS(𝐄)`$.
Let $`h:XS(𝐄)`$ be locally Lipschitz and such that $`c^{}k^{2s}|h(k)|c^{\prime \prime }k^{2s}`$ and $`|h^{}(k)|ck^{2s}`$ outside a compact, where $`c,c^{},c^{\prime \prime }`$ are constants. We set $`H_0=h(P)`$ and observe that $`D(|H_0|^{1/2})=^s`$.
Let $`v:XS(𝐄)`$ be a locally integrable function such that the operator $`V=v(Q)`$ satisfies $`V^s^s`$ and $`\pm V\mu |H_0|+\nu `$ for some real numbers $`\mu ,\nu `$ with $`\mu <1`$. Then the self-adjoint operator $`H=H_0+V`$ (form sum) is affiliated to $`C(X)`$, cf. Corollary 4.8. Note that $`x.V=U_xVU_x^{}=v(x+Q)`$ satisfies the same estimates as $`V`$ and that $`x.H=H_0+x.V`$. We mention that the next lemma is valid under the much more general conditions of Definition 4.7.
###### Lemma 6.32
Let us assume that for each $`𝒞^{\mathrm{}}`$ function with compact support $`f`$ the set $`\{x.VfxX\}`$ is relatively compact in $`^s`$. Then for each $`\varkappa \delta X`$ the limit $`lim_{x\varkappa }x.V=\varkappa .V`$ exists in the strong operator topology in $`B(^s,^s)`$, we have $`\pm \varkappa .V\mu |H_0|+\nu `$ as forms on $`^s`$, and we have $`\varkappa .H=H_0+\varkappa .V`$ if $`\varkappa .H`$ is defined as in Theorem 1.2.
Proof: Let $`z\rho (H)`$ and $`R=(Hz)^1C(X)`$. Then $`\varkappa .H`$ is defined by the operator $`\varkappa .R=lim_{x\varkappa }x.R`$ where the limit exists in $`C_s(X)`$. Note that we know that the limit exists but we do not yet know whether $`\varkappa .R`$ is injective or not. On the other hand, the existence of $`\varkappa .V`$ follows from the fact that the set of operators $`x.V`$ is bounded in $`B(^s,^s)`$: thus it suffices to show the existence of the limit $`lim_{x\varkappa }x.Vf`$ in $`^s`$ for $`f`$ a $`𝒞^{\mathrm{}}`$ function with compact support, and this is obvious by the universal property of the Stone-Čech compactification $`\gamma X`$ of the discrete space $`X`$ and our assumption. Note that $`\varkappa .V`$ is the operator of multiplication by a distribution which could not be a function, but clearly the estimate verified by $`V`$ remains valid in the limit. Hence if *we define* $`\varkappa .H=H_0+\varkappa .V`$ as form sum, we get a densely defined self-adjoint operator such that $`\varkappa .Hz`$ extends to an isomorphism $`^s^s`$. Now it suffices to prove that $`\varkappa .R=(\varkappa .Hz)^1`$. Since $`x.Hz:^s^s`$ is also an isomorphism, one can easily justify the equality
$$(x.Hz)^1(\varkappa .Hz)^1=(x.Hz)^1(\varkappa .Vx.V)(\varkappa .Hz)^1$$
in $`B(^s,^s)`$. Then for $`f^s`$ we have
$$[(x.Hz)^1(\varkappa .Hz)^1]f_^sC(\varkappa .Vx.V)(\varkappa .Hz)^1f_^s$$
where
$$C=(Hz)^1_{B(^s,^s)}=(x.Hz)^1_{B(^s,^s)}.$$
This clearly finishes the proof.
This lemma gives a rather concrete method of computing $`\varkappa .H`$ and also shows that this operator is densely defined. The most elementary way of checking the relative compacity assumption from the lemma is described below.
###### Proposition 6.33
Assume that for each $`\mu >0`$ there is $`\nu `$ such that $`|V|\mu P^{2s}+\nu `$. Then $`lim_{x\varkappa }x.V=\varkappa .V`$ exists strongly in $`B(^s,^s)`$ for each $`\varkappa \delta X`$, for each $`\mu >0`$ there is $`\nu `$ such that $`\pm \varkappa .V\mu |H_0|+\nu `$ as forms on $`^s`$, and $`\varkappa .H=H_0+\varkappa .V`$ if $`\varkappa .H`$ is defined as in Theorem 1.2. In particular, we have
$$\sigma _{\mathrm{ess}}(H)=\overline{}_{\varkappa \delta X}\sigma (\varkappa .H).$$
Proof: We only have to show that the set $`\{x.VfxX\}`$ is relatively compact in $`^s`$ if $`f𝒞_c^{\mathrm{}}(X)`$, i.e. if $`f`$ is a $`𝒞^{\mathrm{}}`$ function with compact support. This is equivalent to the relative compactness in $``$ of the set $`\{P^sx.VfxX\}`$. Let $`\psi ,\xi 𝒞_c^{\mathrm{}}(X)`$ with $`\xi (x)=1`$ on $`\text{supp}f`$ and let $`S=P^s\xi (Q)P^s`$ and $`T=P^sVP^s`$. Then:
$$\psi (P)P^sx.Vf=\psi (P)P^s\xi (Q)x.Vf=\psi (P)SU_xTU_x^{}P^sf\psi (P)Sf_x.$$
The set $`\{f_xxX\}`$ is bounded in $``$ and the operator $`\psi (P)S`$ is compact in $``$, so the set $`\{\psi (P)P^sx.VfxX\}`$ is relatively compact in $``$. Thus it suffices to prove the following assertion: for each $`\epsilon >0`$ there is $`\psi 𝒞_c^{\mathrm{}}(X)`$ such that $`\psi (P)^{}P^sx.Vf\epsilon `$ for all $`xX`$, where $`\psi (P)^{}=1\psi (P)`$. Let $`V_\pm `$be the positive and negative parts of $`V`$, so that $`V=V_+V_{}`$ and $`|V|=V_++V_{}`$, then it is clearly sufficient to prove this assertion with $`V`$ replaced by $`V_\pm `$. If $`T_\pm =P^sV_\pm P^s`$ then
$$\psi (P)^{}P^sx.V_\pm f=\psi (P)^{}U_xT_\pm U_x^{}P^sf\psi (P)^{}T_\pm P^sf$$
and if we set $`C_\pm =T_\pm ^{1/2}`$ then
$$\psi (P)^{}T_\pm C_\pm \psi (P)^{}T_\pm ^{1/2}=C_\pm \psi (P)^{}T_\pm \psi (P)^{}^{1/2}.$$
On the other hand, from $`|V|\mu P^{2s}+\nu `$ we get $`V_\pm \mu P^{2s}+\nu `$ and then $`T_\pm \mu +\nu P^{2s}`$ and so
$$\psi (P)^{}T_\pm \psi (P)^{}[\mu +\nu P^{2s}](1\psi (P))^2.$$
Since $`\mu `$ can be chosen as small as we wish, it is clear that the right hand side above can be made $`\epsilon `$ for any $`\epsilon >0`$ by choosing $`\psi `$ conveniently. Since the left hand side is positive we then get $`\psi (P)^{}T_\pm \psi (P)^{}\epsilon `$.
###### Corollary 6.34
If the conditions of Proposition 6.33 are satisfied and if for each $`C^{\mathrm{}}`$ function $`f`$ with support in the unit ball we have $`lim_x\mathrm{}x.Vf_^s=0`$, then the essential spectrum of $`H`$ is given by $`\sigma _{\mathrm{ess}}(H)=\sigma (H_0)`$.
###### Example 6.35
There are at least three physically interesting situations covered by Proposition 6.33:
(1) The Schrödinger operator $`H=P^2+V=\mathrm{\Delta }+V(x)`$. Then $`s=1`$ and the assumptions of the proposition are satisfied if $`V`$ is of Kato class, so we get Theorem 4.5 from \[LaS\]. Corollary 6.34 is similar to \[LaS, Theorem 4.3\].
(2) The relativistic spin zero operator $`H=(P^2+m^2)^{1/2}+V(x)`$, then $`s=1/2`$. Here $`m`$ is any real number.
(3) The Dirac operator $`H=D+V(x)`$. Here $`D`$ is the free Dirac operator of mass $`m0`$, $`s=1/2`$, $`𝐄=^N`$ is not trivial, and $`V(x)`$ is matrix valued. The last two situations are also considered in \[Rab\].
6.7. Cocompact subgroups. We consider now a situation similar to that from \[LaS, Section 5\]. In a $`C^{}`$-algebra setting such examples and generalizations appear in \[Man\].
Throughout this subsection $`X`$ is an abelian locally compact non compact group and $`Y`$ a closed subgroup such that $`X/Y`$ is compact. Since $`Y`$ is fixed, we shall abbreviate $`\pi _Y=\pi `$. We embed $`𝒞(X/Y)𝒞(X)`$ as explained on page 6.14, so we think of $`𝒞(X/Y)`$ as a translation invariant $`C^{}`$-subalgebra of $`𝒞(X)`$ containing the constants, explicitly described by (6.21). More generally, we identify any function $`v`$ defined on $`X/Y`$ with the function $`v\pi `$ defined on $`X`$.
The full justification of the class of functions introduced below will become clear later on, cf. Lemma 6.41.
###### Definition 6.36
If $`\theta :XX`$ then write $`\theta (a+x)a+\theta (x)`$ if
(6.35)
$$\underset{x\mathrm{}}{lim}[\theta (a+x)a\theta (x)]=0aX.$$
If $`\theta `$ is uniformly continuous then this is equivalent to having $`\theta (x)=x+\xi (x)`$ where $`\xi :XX`$ is uniformly continuous and slowly oscillating in a sense similar to (6.4).
The next proposition is completely elementary but we state it separately because the main idea of the proof is very clear in this context.
###### Proposition 6.37
Let $`h:X^{}`$ be a continuous function such that $`|h(k)|\mathrm{}`$ as $`k\mathrm{}`$. Assume that $`\theta :XX`$ is uniformly continuous and $`\theta (a+x)a+\theta (x)`$. If $`v:X/Y`$ is continuous, $`H=h(P)+v(Q)`$ and $`H_\theta =h(P)+v\theta (Q)`$, then $`\sigma _{\mathrm{ess}}(H_\theta )=\sigma (H)`$.
Proof: This will be a consequence of Theorem 1.2. The self-adjoint operator $`H_\theta `$ is affiliated to $`C(X)`$ because of Proposition 4.3. It remains to compute the localizations $`\varkappa .H`$ for $`\varkappa \delta X`$. The image $`\pi \theta (\varkappa )`$ is an ultrafilter on the compact space $`X/Y`$, hence it converges to some unique point $`\widehat{\varkappa }X/Y`$. Let $`zX`$ such that $`\widehat{\varkappa }=\pi (z)`$. Since $`vv\pi `$ we may define the translated function $`\widehat{\varkappa }.v(s)=v(\widehat{\varkappa }+s)=v\pi (z+x)`$ for $`s=\pi (x)X/Y`$. We shall prove that $`\varkappa .H_\theta =\widehat{\varkappa }.H`$ where $`\widehat{\varkappa }.H=h(P)+\widehat{\varkappa }.v(Q)`$. Note that $`\widehat{\varkappa }.H=U_zHU_z^{}`$, so $`\sigma (\widehat{\varkappa }.H)=\sigma (H)`$, which finishes the proof.
Observe that $`D(H)=D(h(P))`$ is stable under translations, so it suffices to prove the much stronger fact: $`\text{s\hspace{0.17em}-}lim_{x\varkappa }x.H_\theta f=\widehat{\varkappa }.Hf`$ if $`fD(H)`$. But this follows from $`\text{s\hspace{0.17em}-}lim_{x\varkappa }x.v\theta (Q)=\widehat{\varkappa }.v(Q)`$. This means that for each $`fL^2(X)`$ we have
(6.36)
$$\underset{x\varkappa }{lim}_X|v\pi (\theta (x+y))v\pi (z+y)|^2|f(y)|^2\text{d}y=0$$
where $`z`$ is as above. Now for large $`x`$ we have $`\pi (\theta (x+y))\pi (\theta (x))+\pi (y)`$ and $`lim_{x\varkappa }\pi (\theta (x))=\widehat{\varkappa }=\pi (z)`$ and since $`v`$ is uniformly continuous we have
$$\underset{x\varkappa }{lim}\underset{yK}{sup}|v\pi (\theta (x+y))v\pi (z+y)|=0$$
for each compact $`KX`$. This clearly implies (6.36).
The extension of Proposition 6.37 to bounded measurable functions $`v`$ seems to require further conditions. Indeed, one could say that it suffices to use the dominated convergence theorem in (6.36). But this requires some care because $`\varkappa `$ is a filter, we did not assume $`X`$ separable, and the dominated convergence theorem is not true if sequences are replaced by nets. We indicate below two situations where these problems can be avoided.
###### Proposition 6.38
If $`X/Y`$ is separable or if $`\theta `$ is such that $`\theta ^1(N)`$ is of measure zero whenever $`NX`$ is of measure zero, then Proposition 6.37 remains valid if $`v`$ is a bounded measurable function.
Proof: If $`X/Y`$ is separable then the point $`\widehat{\varkappa }`$ has a countable fundamental system of neighborhoods $`\{G_n\}`$. For each $`n`$ choose $`F_n\varkappa `$ such that $`\pi (\theta (F_n))G_n`$ and then choose points $`x_nX`$ such that $`x_nF_n`$. Clearly we shall have $`\pi (\theta (x_n))\widehat{\varkappa }`$ and $`\text{s\hspace{0.17em}-}lim_n\mathrm{}x_n.H_\theta =\varkappa .H_\theta `$ if the left hand side exists. Now the rest of the proof of Proposition 6.37 works after replacing $`x`$ by $`x_n`$ and $`x\varkappa `$ by $`n\mathrm{}`$, this time we can use the dominated convergence theorem directly in (6.36).
If $`\theta `$ has the property $`|N|=0|\theta ^1(N)|=0`$, we argue as follows. Since $`v`$ is bounded, it is sufficient to prove (6.36) for $`f`$ the characteristic function of a compact set. Then we approximate $`v`$ in $`L^2(X/Y)`$ by functions $`w𝒞(X/Y)`$, for such $`w`$ the relation (6.36) being obvious. The only problem which appears is to estimate the term
$$_K|v\eta (x+y)w\eta (x+y)|^2\text{d}y$$
where $`KX`$ is a compact and $`\eta =\pi \theta `$. The map $`\eta :XX/Y`$ is continuous and has the property $`|N|=0|\eta ^1(N)|=0`$, by hypothesis and \[Fol, Theorem 2.9\]. But this implies that there is an integrable function $`g0`$ on $`X/Y`$ such that the preceding integral be $`_{X/Y}|vw|^2g\text{d}s`$, and this can be made as small as we wish.
In the case $`X=^n`$ and under stronger assumptions on $`\theta `$ we may extend Proposition 6.37 to unbounded functions $`v`$, in particular we may recover Theorem 5.1 of the revised version of \[LaS\]. In order to be specific, *we assume that $`h`$ is as in Subsection 6.31 and that $`s1`$*. In particular our assumptions below imply those of Proposition 6.33. Then we easily obtain:
###### Proposition 6.39
Let $`Y`$ be the additive subgroup of $`X=^n`$ generated by $`n`$ linearly independent vectors. Let $`\theta :XX`$ be a homeomorphism such that $`\theta `$ and $`\theta ^1`$ are Lipschitz and such that $`\theta (a+x)a+\theta (x)`$. Assume that $`v:XB(𝐄)`$ is a locally integrable $`Y`$-periodic function and that $`V=v(Q)`$ has the property: for each $`\mu >0`$ there is $`\nu `$ such that $`|V|\mu P^{2s}+\nu `$. Then the operator $`V_\theta =v\theta (Q)`$ has the same property and if we set $`H=h(P)+V`$ and $`H_\theta =h(P)+V_\theta `$ then $`\sigma _{\mathrm{ess}}(H_\theta )=\sigma (H)`$.
###### Example 6.40
We give an elementary example. Let $`X=,Y=`$ and let $`v`$ be a real periodic locally integrable function on $``$. Then the form sum $`H=\frac{d^2}{dx^2}+v(x)`$ is a self-adjoint operator on $`L^2()`$ and its spectrum is purely absolutely continuous. Let $`\theta :`$ be of class $`C^1`$ with $`\theta ^{}(x)>0`$ for all $`x`$ and such that $`\theta ^{}(x)1`$ as $`|x|\mathrm{}`$. Then the form sum $`H_\theta =\frac{d^2}{dx^2}+v(\theta (x))`$ is a self-adjoint operator and its essential spectrum is equal to the spectrum of $`H`$.
We shall now consider the questions treated above in this subsection from the point of view of Theorem 1.15. As in the proof of Proposition 6.3, we get from (6.21) and from Theorem 3.7:
$$𝒞(X/Y)X=\{TC(X)y.T=TyY\}.$$
Clearly $`𝒞_0(X)𝒞(X/Y)=\{0\}`$, from which it follows easily that $`𝒜=𝒞_0(X)+𝒞(X/Y)`$ is an algebra of interactions and that we are in the conditions of Proposition 3.16, hence we have a topological direct sum
$`A`$ $``$ $`𝒜X=K(X)+𝒞(X/Y)X`$
$`=`$ $`\{TC(X)y.TTK(X)yY\}.`$
This is a rather trivial algebra but things become less trivial when we look at the image of $`A`$ under an automorphism of $`C(X)`$.
If $`\theta :XX`$ is a uniformly continuous homeomorphism then $`\theta ^{}:𝒞(X)𝒞(X)`$ is the injective morphism defined by $`\theta ^{}\phi =\phi \theta `$. Clearly $`\theta ^{}𝒞_0(X)=𝒞_0(X)`$ but in nontrivial situations the image through $`\theta ^{}`$ of an $`X`$-algebra is not an $`X`$-algebra. However, we are interested only in algebras of interactions (which contain $`𝒞_0(X)`$), and the property of $`\theta `$ isolated in Definition 6.36 will be sufficient. The next lemma and its corollary are obvious.
###### Lemma 6.41
If $`\theta :XX`$ is uniformly continuous and $`\theta (a+x)a+\theta (x)`$, then for each $`aX`$ the map $`\tau _a\theta ^{}\theta ^{}\tau _a`$ sends $`𝒞(X)`$ into $`𝒞_0(X)`$.
###### Corollary 6.42
Let $`\theta :XX`$ be a uniformly continuous homeomorphism such that $`\theta (a+x)a+\theta (x)`$. Then, if $`𝒜`$ is an algebra of interactions, $`𝒜^\theta :=\theta ^{}𝒜`$ is also an algebra of interactions. Moreover, $`\theta ^{}`$ leaves $`𝒞_0(X)`$ invariant and so it induces a canonical isomorphism of $`X`$-algebras $`𝒜/𝒞_0(X)𝒜^\theta /𝒞_0(X).`$
We apply this to the situation (6). Since $`X/Y`$ is compact, we have
$$\delta (𝒜)=\sigma (𝒜/𝒞_0(X))=\sigma (𝒞(X/Y))=X/Y$$
and thus we get $`\delta (𝒜^\theta )X/Y`$. Since $`X`$ acts transitively on $`X/Y`$ we see that, modulo a unitary equivalence, we have only one localization at infinity for an observable affiliated to $`A_\theta :=𝒜^\theta X`$. This assertion can be made more precise as follows.
###### Proposition 6.43
If $`\theta :XX`$ is a uniformly continuous homeomorphism such that $`\theta (a+x)a+\theta (x)`$, then there is a unique morphism $`𝒫:A_\theta 𝒞(X/Y)X`$ such that
(6.38)
$$𝒫(\phi \theta (Q)\psi (P))=\{\begin{array}{cc}0\text{if}\phi 𝒞_0(X),\hfill & \\ \phi (Q)\psi (P)\text{if}\phi 𝒞(X/Y).\hfill & \end{array}$$
This morphism is surjective and has $`K(X)`$ as kernel. If $`\varkappa `$ is a filter on $`X`$ finer than the Fréchet filter and such that $`lim_\varkappa \pi \theta =0`$, then $`𝒫(T)=lim_{x\varkappa }U_xTU_x^{}`$, where the limit exists in $`C_s(X)`$.
Proof: The uniqueness of $`𝒫`$ follows from the fact that the operators $`\phi \theta (Q)\psi (P)`$ with $`\phi 𝒜`$ generate $`A_\theta `$, and the surjectivity holds for a similar reason. A filter $`\varkappa `$ as in the statement of the proposition exists because $`\pi \theta :XX/Y`$ is surjective. If $`T=\phi \theta (Q)\psi (P)`$ then $`U_xTU_x^{}=\phi \theta (x+Q)\psi (P)`$ and $`\phi \theta (x+Q)\psi (P)\xi (Q)`$ converges strongly to zero or to $`\phi (Q)\psi (P)\xi (Q)`$ as $`x\varkappa `$ if $`\phi 𝒞_0(X)`$ or $`\phi 𝒞(X/Y)`$ respectively. Here $`\xi 𝒞_0(X)`$ and Remark 3.13 is used. Thus we can define $`𝒫(T)`$ by the last assertion of the proposition. If $`𝒫(T)=0`$ then the beginning of the proof of Proposition 6.37 shows that $`\tau _\chi T=0`$ for all $`\chi \delta X`$, so $`T`$ is compact.
###### Remark 6.44
Thus, if $`H`$ is an element of $`A_\theta `$ or an observable affiliated to $`A_\theta `$, then
(6.39)
$$\sigma _{\mathrm{ess}}(H)=\sigma (𝒫(H)).$$
One may get a large class of Hamiltonians $`H`$ affiliated to $`A_\theta `$ by using Theorem 2.8 and Lemma 2.9 from\[DG3\]. For example, let $`H_00`$ be self-adjoint operator strictly<sup>2</sup><sup>2</sup>2 Strictly means $`(1+\epsilon H_0)^1TT0`$ as $`\epsilon 0`$ for all $`TA_\theta `$. For example, it suffices that $`H_0=h(P)`$ where $`h`$ is a positive continuous function on $`X^{}`$ which diverges at infinity. affiliated to $`A_\theta `$. Let $`V`$ be a quadratic form with $`\mu H_0\nu V\nu (H_0+1)`$ for some $`0<\mu <1`$ and $`\nu >0`$ and such that $`(H_0+1)^\alpha V(H_0+1)^{1/2}A_\theta `$ for some $`\alpha >0`$. Then the form sum $`H=H_0+V`$ is a self-adjoint operator affiliated to $`A_\theta `$. For example, the last condition is satisfied if $`VA^\theta `$ and then one gets singular functions $`V`$ as limits of sequences $`V_n`$ such that $`(H_0+1)^\alpha V_n(H_0+1)^{1/2}A_\theta `$ is norm convergent (this gives a class of $`V`$ larger than that from Proposition 6.39).
## Appendix A Appendix
A.1. We give here a detailed proof of Theorem 3.7. We follow rather closely Landstad’s arguments, but we use the characterization of $`𝒜X`$ taken as Definition 3.1, which makes the proof more transparent. We mention that the space $`B_2^{}(X)`$ is suggested by Kato’s theory of smooth operators, cf. \[RS\]. We shall not discuss the uniqueness of $`𝒜`$ because the proof of \[Lan, Lemma 3.1\] can hardly be simplified (if $`X`$ is discrete we have $`𝒜=I(A)`$ so uniqueness is trivial, see Remark A.8).
We begin with some heuristic comments which will make the rigorous proof quite natural. The first question is, given $`A`$, how to determine $`𝒜`$. Observe that if we know $`A\psi (P+k)`$ for all $`k`$ then we can recuperate the operator $`A`$ by integrating over $`k`$, because this operation will give $`A\psi `$ with $`\psi :=_X^{}\psi (k)\text{d}k`$. On the other hand, $`\psi (P+k)=V_k^{}\psi (P)V_k`$ so that if $`A`$ commutes with $`V_k`$ then we get $`A\psi =_X^{}V_k^{}A\psi (P)V_k\text{d}k`$. Thus if $`T=_j\phi _j(Q)\psi _j(P)`$ then
$$\underset{j}{}\phi _j(Q)\psi _j=_X^{}V_k^{}TV_k\text{d}k=:I(T)$$
If the group $`X`$ is discrete, so that $`X^{}`$ is compact, this formal argument can easily be made rigorous, the map $`I`$ is well defined on all $`B(X)`$ and we have $`𝒜=I(A)`$ (we strongly advise the reader to first prove Landstad’s theorem for discrete $`X`$; this is a really pleasant exercise). In general, one can give a meaning to $`I(T)`$ for a sufficiently large class of $`T`$ for the rest of the proof to work. Anyway, the preceding formula shows how to extract the part in $`𝒜`$ of the operator $`TA`$. The second point is that one can reconstruct $`T`$ from such quantities by using the formally obvious relation
$$T=_XI(TU_x^{})U_x\text{d}x$$
which is just the Fourier inversion formula, see (A.11). But the right hand side here is, again formally, in $`𝒜X`$.
To make all this rigorous demands some preliminary constructions that we expose in Subsection A in a more general context. We set $`=L^2(X)`$ and we abbreviate $`B=B(X)=B()`$. We recall that we have unitary representations $`U_x`$ and $`V_k`$ of $`X`$ and $`X^{}`$ in $``$ which satisfy the canonical commutation relations
(A.1)
$$U_xV_k=k(x)V_kU_x.$$
Most of the next arguments do not depend on the explicit form of the operators $`U_x,V_k`$.
A.2. We first introduce the space of “smooth operators” with respect to the unitary representation $`V_k`$:
(A.2)
$$B_2^{}:=\{TB_X^{}TV_kf^2\text{d}k<\mathrm{},f\}.$$
###### Lemma A.1
$`B_2^{}`$ is a left ideal (not closed in general) in $`B`$. If $`TB_2^{}`$ then
(A.3)
$$|T|:=\underset{f=1}{sup}\left(_X^{}TV_kf^2\text{d}k\right)^{1/2}<\mathrm{}$$
and $`(B_2^{},||||||)`$ is a Banach space such that $`|ST|S|T|`$ for all $`SB`$. Finally, if $`xX`$ and $`TB_2^{}`$ then $`TU_xB_2^{}`$ and $`|TU_x|=|T|`$.
For the proof of (A.3) we have only to remark that the map which sends $`f`$ into $`(TV_kf)_{kX^{}}L^2(X^{};)`$ is clearly closed and linear, hence it is continuous. The last assertion of the lemma follows from (A.1).
The map $`xTU_xB_2`$ is not norm continuous in general. For this reason it will be convenient to consider the following left ideal in $`B`$ and closed subspace of $`B_2^{}`$
(A.4)
$$B_2:=\{TB_2^{}\underset{x0}{lim}|TU_xT|=0\}.$$
The following property of $`B_2`$ will be important in what follows: if $`S^{}ST_j^{}T_j`$ for some $`T_jB_2`$, then $`SB_2`$.
###### Lemma A.2
If $`\psi L^{\mathrm{}}(X^{})`$ then $`\psi (P)B_2^{}`$ if and only if $`\psi L^2(X^{})`$. In this case we have $`\psi (P)B_2`$ and $`|\psi (P)|=\psi _{L^2}`$.
Proof: We have
$`{\displaystyle _X^{}}\psi (P)V_kf^2\text{d}k`$ $`=`$ $`{\displaystyle _X^{}}V_k^{}\psi (P)V_kf^2\text{d}k={\displaystyle _X^{}}\psi (P+k)f^2\text{d}k`$
$`=`$ $`{\displaystyle _X^{}}\text{d}k{\displaystyle _X^{}}|\psi (p+k)|^2|\widehat{f}(p)|^2\text{d}p`$
$`=`$ $`\psi _{L^2(X^{})}\widehat{f}_{L^2(X^{})}=\psi _{L^2(X^{})}f_{L^2(X)}.`$
Then $`|\psi (P)U_x\psi (P)|=x\psi \psi _{L^2}`$ where $`x`$ is identified with the map $`kk(x)`$ on $`X^{}`$, we clearly get $`\psi B_2`$.
###### Definition A.3
$`B_1`$ is the linear subspace of $`B`$ generated by the operators of the form $`S^{}T`$ with $`S,TB_2`$.
The polarization identity
(A.5)
$$4S^{}T=\underset{m=0}{\overset{3}{}}i^m(i^mS+T)^{}(i^mS+T)$$
shows that $`B_1`$ is linearly generated by the operators of the form $`S^{}S`$ with $`SB_2`$.
We recall that a subset $`CB`$ is called hereditary if: $`0STCSC`$.
###### Lemma A.4
$`B_1`$ is an hereditary $``$-subalgebra of $`B`$. If $`S0`$ then $`SB_1`$ if and only if $`\sqrt{S}B_2`$. If $`SB_1`$ then $`U_xSU_yB_1`$ for all $`x,yX`$.
Proof: The fact that $`B_1`$ is a linear space and that $`S^{}B_1`$ if $`SB_1`$ is obvious. $`B_1`$ is stable under multiplication because for $`S,TB_2`$ we have $`S^{}ST^{}T=S^{}ST^{}TB_1`$ the space $`B_2`$ being a left ideal.
We prove now that if $`S0`$ and $`SB_1`$ then $`\sqrt{S}B_2`$ (the reverse implication being obvious). Since $`SB_1`$ we have $`S=_{j=1}^n\lambda _jS_j^{}S_j`$ with $`\lambda _j`$ and $`S_jB_2`$. If $`S=S^{}`$ then by taking the real parts we may assume that $`\lambda _j`$. Then
$$S=\left(\underset{\lambda _j>0}{}+\underset{\lambda _j<0}{}\right)\lambda _jS_j^{}S_j\underset{\lambda _j>0}{}\lambda _jS_j^{}S_j,$$
which implies $`\sqrt{S}B_2`$ by the property mentioned after (A.4).
Finally, if $`0STB_1`$ then $`\sqrt{T}B_2`$, so $`SB_1`$ by the same property.
Let $`TB_1`$ and let us write $`T=_{j=1}^nS_j^{}T_j`$ with $`S_j,T_jB_2`$. Then if $`f,g`$:
$`{\displaystyle _X^{}}|V_kf,TV_kg|\text{d}k`$ $``$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}\left({\displaystyle _X^{}}S_jV_kf^2\text{d}k\right)^{1/2}\left({\displaystyle _X^{}}T_jV_kg^2\text{d}k\right)^{1/2}`$
$``$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}|S_j||T_j|fg<\mathrm{}.`$
From the operator version of the Riesz lemma it follows that there is a unique operator $`I(T)B(X)`$ such that
$$f,I(T)g=_X^{}V_kf,TV_kg\text{d}k\text{ for all }f,g.$$
In other terms, we see that the strongly continuous map $`kV_k^{}TV_k`$ is such that the integral
(A.6)
$$I(T)=_X^{}V_k^{}TV_k\text{d}k$$
exists in the weak operator topology of $`B(X)`$. It is clear that for all $`TB_2`$ we have
(A.7)
$$I(T^{}T)^{1/2}=|T|.$$
Moreover, the computation done above gives for $`S,TB_2`$:
(A.8)
$$I(S^{}T)|S||T|.$$
###### Example A.5
If $`SB(X)`$ and $`\xi ,\eta L^{\mathrm{}}(X^{})L^2(X^{})`$ then $`\xi (P)S\eta (P)B_1`$ and
(A.9)
$$I\left(\xi (P)S\eta (P)\right)S\xi _{L^2(X^{})}\eta _{L^2(X^{})}.$$
Indeed, we write $`\xi (P)S\eta (P)=\left(S^{}\overline{\xi }(P)\right)^{}\eta (P)`$ and use (A.8) and Lemma A.2.
###### Lemma A.6
If $`TB_1`$ then $`(x,y)I(U_xTU_y)`$ is bounded and norm continuous.
Proof: Due to (A.8) it suffices to assume that $`T=S^{}S`$ for some $`SB_2`$ and to show continuity at $`x=y=0`$ of the map $`(x,y)I(S_x^{}S_y)`$ with $`S_x=SU_x`$. Then
$`I(S_x^{}S_y)I(S^{}S)`$ $``$ $`I\left((S_xS)^{}S_y\right)+I\left(S^{}(S_yS)\right)`$
$``$ $`|S_xS||S|+|S||S_yS|`$
because of the estimate (A.8).
###### Proposition A.7
If $`TB_1`$ then $`I(T)𝒞(X)`$ and $`U_x^{}I(T)U_x=I(U_x^{}TU_x)`$ for all $`xX`$. The map $`I:B_1𝒞(X)`$ is linear and positive.
Proof: We clearly have $`V_kI(T)V_k^{}=I(T)`$ for all $`kX^{}`$. Since the Von Neumann algebra generated by $`\{V_k\}_{kX^{}}`$ is just $`L^{\mathrm{}}(X)`$, we get $`\phi (Q)I(T)=I(T)\phi (Q)`$ for all $`\phi L^{\mathrm{}}(X)`$. But $`L^{\mathrm{}}(X)`$ is maximal abelian in $`B`$, thus $`I(T)L^{\mathrm{}}(X)`$. From (A.1) we get $`U_x^{}V_k^{}TV_kU_x=V_k^{}U_x^{}TU_xV_k`$ hence $`I(U_x^{}TU_x)=U_x^{}I(T)U_x`$. Since $`\phi 𝒞(X)`$ if and only if $`\phi L^{\mathrm{}}(X)`$ and $`xU_x^{}\phi (Q)U_x`$ is norm continuous, we get $`I(T)𝒞(X)`$. The last assertion of the proposition is obvious.
###### Remark A.8
In a similar way we can associate an hereditary $``$-subalgebra $`B_1^{}`$ to $`B_2^{}`$ and define an extension of $`I`$ to it, but then we only have $`I:B_1^{}L^{\mathrm{}}(X)`$. If $`X`$ is a discrete group, then $`B_1=B`$ and $`I`$ is a conditional expectation.
For $`TB_1`$ and for $`xX`$ we set $`\stackrel{~}{T}(x):=I(TU_x^{})`$, so that we associate to $`T`$ a function $`\stackrel{~}{T}:X𝒞(X)`$ which, by Lemma A.6, is bounded and norm continuous. Let $`\widehat{T}`$ be the Fourier transform of the function $`kV_k^{}TV_k`$, more precisely
$$\widehat{T}(x)=_X^{}\overline{k(x)}V_k^{}TV_k\text{d}k$$
where the integral exists in the weak operator topology. From (A.1) we get
(A.10)
$$\stackrel{~}{T}(x)=\widehat{T}(x)U_x^{}.$$
so that $`\stackrel{~}{T}`$ is a kind of twisted Fourier transform. Now the inversion formula for the Fourier transform gives us a *formal relation*
(A.11)
$$T=_X\stackrel{~}{T}(x)U_x\text{d}x$$
whose rigorous meaning is given below.
###### Lemma A.9
For each $`TB_1`$ and $`\theta L^1(X)`$ we have
(A.12)
$$_X\stackrel{~}{T}(x)U_x\theta (x)\text{d}x=_X^{}V_k^{}TV_k\widehat{\theta }(k)\text{d}k$$
where both integrals exist in the weak operator topology.
Proof: For each $`f`$,
$`{\displaystyle _X}f,\stackrel{~}{T}(x)U_xf\theta (x)\text{d}x`$ $`=`$ $`{\displaystyle _X}\left({\displaystyle _X^{}}V_kf,T\overline{k(x)}V_kf\text{d}k\right)\theta (x)\text{d}x`$
$`=`$ $`{\displaystyle _X}\left({\displaystyle _X^{}}\overline{k(x)}V_kf,TV_kf\text{d}k\right)\theta (x)\text{d}x.`$
Since $`\theta L^1(X)`$ and the function $`kV_kf,TV_kf`$ is in $`L^1(X^{})`$, we can apply the Fubini theorem and get thus get (A.12).
Let us remark that the l.h.s. of the identity (A.12) always exists in the strong operator topology, and the same is true for the r.h.s. if $`\widehat{\theta }L^1(X^{})`$. We recall the following result (see e.g. \[Fol, Lemma 4.19\]).
###### Lemma A.10
Let $`\mathrm{\Lambda }X^{}`$ be a neighborhood of the neutral element in $`X^{}`$ and let $`\epsilon >0`$. Then there is $`\theta 𝒞_\text{c}(X)`$ such that $`\widehat{\theta }0`$, $`\widehat{\theta }L^1(X^{})`$, $`_X^{}\widehat{\theta }(k)\text{d}k=1`$, and $`_{X^{}\mathrm{\Lambda }}\widehat{\theta }(k)\text{d}k\epsilon `$.
The next version of the Fourier inversion formula is an easy consequence of Lemmas A.9 and A.10:
###### Proposition A.11
If $`TB_1`$ and $`kV_k^{}TV_k`$ is norm continuous, then $`T`$ belongs to the norm closure of the set of operators of the form $`T_\theta =_X\stackrel{~}{T}(x)U_x\theta (x)\text{d}x`$ with $`\theta 𝒞_\text{c}(X)`$ and $`\widehat{\theta }L^1(X^{})`$.
###### Lemma A.12
Let $`\psi 𝒞_0(X^{})`$ and $`TB_1`$. Then if $`\theta L^1(X)`$ and $`\widehat{\theta }L^1(X^{})`$ the integral $`_X\stackrel{~}{T}(x)U_x\psi (P)\theta (x)\text{d}x`$ exists in the norm operator topology and $`I(T)\psi (P)`$ is a norm limit of such integrals.
Proof: The map $`xU_x\psi (P)`$ is norm continuous if $`\psi 𝒞_0(X^{})`$, hence the integrand is norm continuous. The last assertion follows by choosing $`\theta `$ as in Lemma A.10 but with the rôles of $`X`$ and $`X^{}`$ inverted.
A.3. We are now ready to prove Landstad’s theorem (Theorem 3.7). From now on, $`A`$ and $`𝒜`$ are as in that theorem.
###### Lemma A.13
$`A`$ is a non-degenerate $`𝒞_0(X^{})`$-bimodule. More precisely, if $`AA`$ and $`\psi 𝒞_0(X^{})`$ then $`A\psi (P)A`$, $`\psi (P)AA`$ and $`A`$ is a limit of operators of the form $`A\psi (P)`$ and of operators of the form $`\psi (P)A`$.
Proof: It is clearly sufficient to consider only the right action and, since each $`\psi 𝒞_0(X^{})`$ is limit in the sup norm of functions whose Fourier transform is integrable, we may assume $`\widehat{\psi }L^1(X)`$. Then $`A\psi (P)=_XAU_x\widehat{\psi }(x)\text{d}x`$, we have $`AU_xA`$ and the integral converges in norm by the second assumption of Theorem 3.7, so $`A\psi (P)A`$. By taking $`\widehat{\psi }=|K|^1\text{1}_K`$, where $`K`$ runs over the set of compact neighborhoods of the origin in $`X`$, and by taking into account the norm continuity of the map $`xAU_x`$, we see that $`A`$ is a norm limit of operators of the form $`A\psi (P)`$.
###### Lemma A.14
$`𝒜`$ is an $`X`$-subalgebra of $`𝒞(X)`$.
Proof: It is clear that $`𝒜`$ is a norm closed subspace of $`𝒞(X)`$ stable under conjugation and stable under translations (note that $`(\tau _x\phi )(Q)=U_x\phi (Q)U_x^{}`$). To show that it is stable under multiplication, let $`\phi _1,\phi _2𝒜`$ and $`\psi 𝒞_0(X^{})`$. Since $`\phi _2(Q)\psi (P)A`$ we can write it as a norm limit of operators of the form $`\xi (P)A`$ with $`\xi 𝒞_0(X^{})`$ and $`AA`$, so that $`\phi _1(Q)\phi _2(Q)\psi (P)`$ is a norm limit of operators of the form $`\phi _1(Q)\xi (P)A`$ which belong to $`A`$.
Now we may consider the crossed product $`𝒜X`$, this is the norm closed subspace of $`B(X)`$ generated by the operators $`\phi (Q)\psi (P)`$ with $`\phi 𝒜`$ and $`\psi 𝒞_0(X^{})`$. We clearly have $`𝒜XA`$ and it remains to prove the reverse inclusion.
###### Lemma A.15
If $`TAB_1`$ then $`I(T)𝒜`$
Proof: Due to Proposition A.7 it suffices to show that $`I(T)\psi (P)A`$ if $`\psi 𝒞_0(X^{})`$. Because of Lemma A.12, it is enough to prove that $`_X\stackrel{~}{T}(x)U_x\psi (P)\theta (x)\text{d}xA`$ if $`\theta L^1(X)`$ and $`\widehat{\theta }L^1(X^{})`$. But (A.12) implies:
$`{\displaystyle _X}\stackrel{~}{T}(x)U_x\psi (P)\theta (x)\text{d}x={\displaystyle _X}V_k^{}TV_k\psi (P)\widehat{\theta }(k)\text{d}k.`$
Since $`V_k^{}TV_k\psi (P)A`$ and is a norm continuous function of $`k`$ the last integral belongs to $`A`$.
###### Lemma A.16
If $`TAB_1`$ then $`T\psi (P)𝒜X`$ for all $`\psi 𝒞_0(X^{})`$.
Proof: We shall have $`TU_x^{}AB_1`$, hence $`\stackrel{~}{T}(x)=I(TU_x^{})𝒜`$, and thus the map $`\stackrel{~}{T}:X𝒜`$ is bounded and norm continuous. On the other hand, Proposition A.11 shows that for each $`\psi 𝒞_0(X^{})`$ the operator $`T\psi (P)`$ is a norm limit of integrals $`_X\stackrel{~}{T}(x)U_x\psi (P)\theta (x)\text{d}x`$. But $`U_x\psi (P)𝒞_0(X^{})`$ and the map $`xU_x\psi (P)`$ is norm continuous, thus the preceding integral converges in norm. Also, we have $`\stackrel{~}{T}(x)U_x\psi (P)𝒜X`$ for each $`x`$, thus the integral belongs to $`𝒜X`$.
Now we prove $`A𝒜X`$. For this it suffices to find a dense subset of $`A`$ which is included in $`𝒜X`$. The Example A.5 and Lemma A.13 imply that $`AB_1`$ is a dense subspace of $`A`$. Thus *it suffices to show that $`AB_1𝒜X`$*. But this follows from Lemma A.16 because each $`TA`$ is a norm limit of operators of the form $`T\psi (P)`$ with $`\psi 𝒞_0(X^{})`$.
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# I Introduction
## I Introduction
Considerable activity in search of possible New Physics beyond the Standard Model has recently been devoted to the measurements of time-dependent CP asymmetries in neutral $`B`$ meson decays into final CP eigenstates defined by
$`{\displaystyle \frac{\mathrm{\Gamma }\left(\overline{B}\left(t\right)f\right)\mathrm{\Gamma }\left(B\left(t\right)f\right)}{\mathrm{\Gamma }\left(\overline{B}\left(t\right)f\right)+\mathrm{\Gamma }\left(B\left(t\right)f\right)}}=𝒮_f\mathrm{sin}\left(\mathrm{\Delta }mt\right)+𝒜_f\mathrm{cos}\left(\mathrm{\Delta }mt\right),`$ (1)
where $`\mathrm{\Delta }m`$ is the mass difference of the two neutral $`B`$ eigenstates, $`S_f`$ monitors mixing-induced CP asymmetry and $`𝒜_f`$ measures direct CP violation (in the BaBar notation, $`𝒞_f=𝒜_f`$). The time-dependent CP asymmetries in the $`bsq\overline{q}`$ penguin-induced two-body decays such as $`B^0(\varphi ,\omega ,\pi ^0,\eta ^{},f_0)K_S`$ measured by BaBar BaBarSf1 ; BaBarSf2 and Belle BelleSf1 ; BelleSf2 ; BelleSf3 show some indications of sizable deviations from the expectation of the SM where CP asymmetry in all above-mentioned modes should be equal to $`S_{J/\psi K_S}=0.687\pm 0.032`$ HFAG with a small deviation at most $`𝒪\left(0.1\right)`$ LS ; Browder . Based on the framework of QCD factorization BBNS , the mixing-induced CP violation parameter $`S_f`$ in the seven 2-body modes $`(\varphi ,\omega ,\rho ^0,\eta ^{},\eta ,\pi ^0,f_0)K_S`$ has recently been quantitatively studied in CCSsin2beta and Buchalla ; Beneke . It is found that the sign of $`\mathrm{\Delta }S_f\eta _fS_fS_{J/\psi K_S}`$ ($`\eta _f`$ being the CP eigenvalue of the final state $`f`$) at short distances is positive except for the channel $`\rho ^0K_S`$. After including final-state rescattering effects, the central values of $`\mathrm{\Delta }S_f`$ become positive for all the modes under consideration, but they tend to be rather small compared to the theoretical uncertainties involved so that it is difficult to make reliable statements on the sign at present CCSsin2beta .
Time-dependent CP asymmetries in the $`bsq\overline{q}`$ induced three-body decays $`B^0K^+K^{}K_S`$ and $`K_SK_SK_S`$ have also been measured by $`B`$ factories BaBarSf2 ; BaBarKKK ; BaBarKKKL ; BaBar3Ks ; BelleSf2 ; BelleSf3 ; BelleKKK (see Table 1). Three-body modes such as these were first discussed by Gershon and Hazumi Gershon . While $`K_SK_SK_S`$ has fixed $`CP`$-parity, $`K^+K^{}K_S`$ is an admixture of $`CP`$-even and $`CP`$-odd components, rendering its CP analysis more complicated. By excluding the major $`CP`$-odd contribution from $`\varphi K_S`$, the 3-body $`K^+K^{}K_S`$ final state is primarily $`CP`$-even. A measurement of the CP-even fraction $`f_+`$ in the $`B^0K^+K^{}K_S`$ decay yields $`f_+=0.89\pm 0.08\pm 0.06`$ by BaBar BaBarSf2 and $`0.93\pm 0.09\pm 0.05`$ by Belle BelleSf3 , while the CP-odd fraction in $`K^+K^{}K_L`$ is measured to be $`f_{}=0.92\pm 0.33_{0.14}^{+0.13}\pm 0.10`$ by BaBar BaBarKKKL . Hence, while $`\eta _f=1`$ for the $`K_SK_SK_S`$ mode, $`\eta _f=2f_+1`$ for $`K^+K^{}K_S`$ and $`\eta _f=\left(2f_{}1\right)`$ for $`K^+K^{}K_L`$. It is convenient to define an effective $`\mathrm{sin}2\beta `$ via $`S_f\eta _f\mathrm{sin}2\beta _{\mathrm{eff}}`$. The results of $`\mathrm{sin}2\beta _{\mathrm{eff}}`$ for $`K^+K^{}K_S`$ obtained from the measurements of $`S_{K^+K^{}K_S}`$ and $`f_+`$ are also shown in Table 1.
In order to see if the current measurements of the deviation of $`\mathrm{sin}2\beta _{\mathrm{eff}}`$ in $`KKK`$ modes from $`\mathrm{sin}2\beta _{J/\psi K_S}`$ signal New Physics in $`bs`$ penguin-induced modes, it is of great importance to examine and estimate how much of the deviation of $`\mathrm{sin}2\beta _{\mathrm{eff}}`$ is allowed in the SM. One of the major uncertainties in the dynamic calculations lies in the hadronic matrix elements which are nonperturbative in nature. One way to circumvent this difficulty is to impose SU(3) flavor symmetry Grossman03 ; Engelhard or isospin and U-spin symmetries Gronau to constrain the relevant hadronic matrix elements. While this approach is model independent in the symmetry limit, deviations from that limit can only be computed in a model dependent fashion. In addition, it may have some weakness as discussed in Engelhard .
We shall apply the factorization approach in this work as it seems to work even in the case of three-body $`B`$ decays DKK . By using factorization and kaon time-like form factors extracted from the $`e^+e^{}K\overline{K}`$ process, the predicted $`\overline{B}{}_{}{}^{0}D^{()+}K^{}K^0`$ rate agrees well with the data DKK . In general, three-body $`B`$ decays are more complicated than the two-body case as they receive resonant and nonresonant contributions and involve 3-body matrix elements. Nonresonant charmless three-body $`B`$ decays have been studied extensively Deshpande ; Fajfer1 ; Fajfer2 ; Deandrea1 ; Deandrea ; Fajfer3 based on heavy meson chiral perturbation theory (HMChPT) Yan ; Wise ; Burdman . However, the predicted decay rates are in general unexpectedly large. For example, the branching ratio of the nonresonant decay $`B^{}\pi ^+\pi ^{}\pi ^{}`$ is predicted to be of order $`10^5`$ in Deshpande and Fajfer1 , which is too large compared to the BaBar’s preliminary result $`\left(0.68\pm 0.41\right)\times 10^6`$ BaBar3pi . The issue has to do with the applicability of HMChPT. In order to apply this approach, two of the final-state pseudoscalars have to be soft. The momentum of the soft pseudoscalar should be smaller than the chiral symmetry breaking scale $`\mathrm{\Lambda }_\chi 830`$ MeV. For 3-body charmless $`B`$ decays, the available phase space where chiral perturbation theory is applicable is only a small fraction of the whole Dalitz plot. Therefore, it is not justified to apply chiral and heavy quark symmetries to a certain kinematic region and then generalize it to the region beyond its validity. In order to have a reliable prediction for the total rate of direct 3-body decays, one should try to utilize chiral symmetry to a minimum. Therefore, we will apply HMChPT only to the strong vertex and use the form factors to describe the weak vertex Cheng:2002qu . Moreover, we shall introduce a form factor to take care of the off-shell effect.
As shown in CCSsin2beta , among the aforementioned seven neutral $`PK_S`$ modes, only the $`\omega K_S`$ and $`\rho ^0K_S`$ modes are expected to have a sizable deviation of the mixing-induced CP asymmetry $`S_f`$ from $`S_{J/\psi K_S}`$. More precisely, it is found $`\mathrm{\Delta }S_{\omega K_S}=0.12_{0.06}^{+0.05}`$ and $`\mathrm{\Delta }S_{\rho ^0K_S}=0.09_{0.07}^{+0.03}`$ <sup>1</sup><sup>1</sup>1Note that since $`K^+K^{}K_S`$ is not a pure CP eigenstate, we define $`\mathrm{\Delta }\mathrm{sin}2\beta _{\mathrm{eff}}\mathrm{sin}2\beta _{\mathrm{eff}}\mathrm{sin}2\beta _{J/\psi K}`$ with $`\mathrm{sin}2\beta _{\mathrm{eff}}=S_f/\eta _f`$. In general, the relation $`\mathrm{\Delta }S_f=\mathrm{\Delta }\mathrm{sin}2\beta _f^{\mathrm{eff}}`$ holds for the final state with fixed CP parity. in the absence of final-state interactions CCSsin2beta . Although the tree contribution in these two modes is color suppressed, the large cancellation between $`a_4`$ and $`a_6`$ penguin terms renders the tree pollution relatively significant. Unlike the above-mentioned case for two-body decays, the tree contribution to the 3-body decay $`B^0K^+K^{}K_S`$ is color-allowed and hence it has the potential for producing a large deviation from $`\mathrm{sin}2\beta `$ measured in $`BJ/\psi K_S`$. We shall see in this work that it is indeed the case. In contrast, the absence of tree pollution in $`K_SK_SK_S`$ renders it theoretically very clean in our picture.
The layout of the present paper is as follows. In Sec. II we apply the factorization approach to study $`B^0K^+K^{}K_S`$ and $`K_SK_SK_S`$ decays and discuss resonant and nonresonant contributions in Sec. II. Numerical results for decay rates and CP-violating parameters $`S_f`$ and $`A_f`$ and discussions are presented in Sec. III. Sec. IV contains our conclusions.
## II Formalism for charmless 3-body $`B`$ decays
In the factorization approach, the matrix element of the $`\overline{B}\overline{K}\overline{K}K`$ decay amplitude is given by
$`\overline{K}\overline{K}K\left|_{\mathrm{eff}}\right|\overline{B}={\displaystyle \frac{G_F}{\sqrt{2}}}{\displaystyle \underset{p=u,c}{}}\lambda _p\overline{K}\overline{K}K\left|T_p\right|\overline{B},`$ (2)
where $`\lambda _pV_{pb}V_{ps}^{}`$ and BBNS
$`T_p`$ $`=`$ $`a_1\delta _{pu}\left(\overline{u}b\right)_{VA}\left(\overline{s}u\right)_{VA}+a_2\delta _{pu}\left(\overline{s}b\right)_{VA}\left(\overline{u}u\right)_{VA}+a_3\left(\overline{s}b\right)_{VA}{\displaystyle \underset{q}{}}\left(\overline{q}q\right)_{VA}`$ (3)
$`+a_4^p{\displaystyle \underset{q}{}}\left(\overline{q}b\right)_{VA}\left(\overline{s}q\right)_{VA}+a_5\left(\overline{s}b\right)_{VA}{\displaystyle \underset{q}{}}\left(\overline{q}q\right)_{V+A}`$
$`2a_6^p{\displaystyle \underset{q}{}}\left(\overline{q}b\right)_{SP}\left(\overline{s}q\right)_{S+P}+a_7\left(\overline{s}b\right)_{VA}{\displaystyle \underset{q}{}}{\displaystyle \frac{3}{2}}e_q\left(\overline{q}q\right)_{V+A}`$
$`2a_8^p{\displaystyle \underset{q}{}}\left(\overline{q}b\right)_{SP}{\displaystyle \frac{3}{2}}e_q\left(\overline{s}q\right)_{S+P}+a_9\left(\overline{s}b\right)_{VA}{\displaystyle \underset{q}{}}{\displaystyle \frac{3}{2}}e_q\left(\overline{q}q\right)_{VA}`$
$`+a_{10}^p{\displaystyle \underset{q}{}}\left(\overline{q}b\right)_{VA}{\displaystyle \frac{3}{2}}e_q\left(\overline{s}q\right)_{VA},`$
with $`\left(\overline{q}q^{}\right)_{V\pm A}\overline{q}\gamma _\mu \left(1\pm \gamma _5\right)q^{}`$, $`\left(\overline{q}q^{}\right)_{S\pm P}\overline{q}\left(1\pm \gamma _5\right)q^{}`$ and a summation over $`q=u,d,s`$ being implied. The matrix element $`\overline{K}\overline{K}K\left|jj^{}\right|\overline{B}`$ corresponds to $`\overline{K}K\left|j\right|\overline{B}\overline{K}\left|j^{}\right|0`$, $`\overline{K}\left|j\right|\overline{B}\overline{K}K\left|j^{}\right|0`$ or $`0\left|j\right|\overline{B}\overline{K}\overline{K}K\left|j^{}\right|0`$, as appropriate, and $`a_i`$ are the NLO effective Wilson coefficients. In this work, we take
$`a_10.99\pm 0.37i,a_20.190.11i,a_30.002+0.004i,a_50.00540.005i,`$
$`a_4^u0.030.02i,a_4^c0.040.008i,a_6^u0.060.02i,a_6^c0.060.006i,`$
$`a_70.54\times 10^4i,a_8^u\left(4.50.5i\right)\times 10^4,a_8^c\left(4.40.3i\right)\times 10^4,`$ (4)
$`a_90.0100.0002i,a_{10}^u\left(58.3+86.1i\right)\times 10^5,a_{10}^c\left(60.3+88.8i\right)\times 10^5,`$
for typical $`a_i`$ at the renormalization scale $`\mu =m_b/2=2.1`$ GeV which we are working on.
Applying Eqs. (2), (3) and the equation of motion, we obtain the $`\overline{B}{}_{}{}^{0}K^+K^{}\overline{K}^0`$ decay amplitude as
$`\overline{K}{}_{}{}^{0}K_{}^{+}K^{}\left|T_p\right|\overline{B}`$ $`=`$ $`K^+\overline{K}{}_{}{}^{0}|\left(\overline{u}b\right)_{VA}|\overline{B}{}_{}{}^{0}K^{}\left|\left(\overline{s}u\right)_{VA}\right|0[a_1\delta _{pu}+a_4^p+a_{10}^p(a_6^p+a_8^p)r_\chi ]`$ (5)
$`+\overline{K}{}_{}{}^{0}|\left(\overline{s}b\right)_{VA}|\overline{B}{}_{}{}^{0}K^+K^{}\left|\left(\overline{u}u\right)_{VA}\right|0(a_2\delta _{pu}+a_3+a_5+a_7+a_9)`$
$`+\overline{K}{}_{}{}^{0}|\left(\overline{s}b\right)_{VA}|\overline{B}{}_{}{}^{0}K^+K^{}\left|\left(\overline{d}d\right)_{VA}\right|0[a_3+a_5{\displaystyle \frac{1}{2}}(a_7+a_9)]`$
$`+\overline{K}{}_{}{}^{0}|\left(\overline{s}b\right)_{VA}|\overline{B}{}_{}{}^{0}K^+K^{}\left|\left(\overline{s}s\right)_{VA}\right|0[a_3+a_4^p+a_5{\displaystyle \frac{1}{2}}(a_7+a_9+a_{10}^p)]`$
$`+\overline{K}{}_{}{}^{0}|\overline{s}b|\overline{B}{}_{}{}^{0}K^+K^{}\left|\overline{s}s\right|0(2a_6^p+a_8^p)`$
$`+K^+K^{}\overline{K}{}_{}{}^{0}|\left(\overline{s}d\right)_{VA}|00\left|\left(\overline{d}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}(a_4^p{\displaystyle \frac{1}{2}}a_{10}^p)`$
$`+K^+K^{}\overline{K}{}_{}{}^{0}|\overline{s}\gamma _5d|00\left|\overline{d}\gamma _5b\right|\overline{B}{}_{}{}^{0}(2a_6^p+a_8^p),`$
with $`r_\chi =2m_K^2/\left(m_bm_s\right)`$. In the factorization terms, the $`K\overline{K}`$ pair can be produced through a transition from the $`\overline{B}`$ meson or can be created from vacuum through $`V`$ and $`S`$ operators. There exist two weak annihilation contributions, where the $`\overline{B}`$ meson is annihilated and a final state with three kaons is created. Note that the OZI suppressed matrix element $`K^+K^{}\left|\left(\overline{d}d\right)_{VA}\right|0`$ is included in the factorization amplitude since it could be enhanced through the long-distance pole contributions via the intermediate vector mesons such as $`\rho ^0`$ and $`\omega `$.
To evaluate the above amplitude, we need to consider the $`\overline{B}K\overline{K}`$, $`0K\overline{K}`$ and $`0\overline{K}\overline{K}K`$ matrix elements, the so-called two-meson transition, two-meson and tree-meson creation matrix elements in addition to the usual one-meson transition and creation ones. The two-meson transition matrix element $`\overline{K}{}_{}{}^{0}K_{}^{+}\left|\left(\overline{u}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}`$ has the general expression LLW
$`\overline{K}{}_{}{}^{0}\left(p_1\right)K^+\left(p_2\right)\left|\left(\overline{u}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}`$ $`=`$ $`ir\left(p_Bp_1p_2\right)_\mu +i\omega _+\left(p_2+p_1\right)_\mu +i\omega _{}\left(p_2p_1\right)_\mu `$ (6)
$`+hϵ_{\mu \nu \alpha \beta }p_B^\nu \left(p_2+p_1\right)^\alpha \left(p_2p_1\right)^\beta .`$
This leads to
$`K^{}\left(p_3\right)\left|\left(\overline{s}u\right)_{VA}\right|0\overline{K}{}_{}{}^{0}\left(p_1\right)K^+\left(p_2\right)\left|\left(\overline{u}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}`$
$`={\displaystyle \frac{f_K}{2}}\left[2m_3^2r+\left(m_B^2s_{12}m_3^2\right)\omega _++\left(s_{23}s_{13}m_2^2+m_1^2\right)\omega _{}\right],`$ (7)
where $`s_{ij}\left(p_i+p_j\right)^2`$. A pole model calculation of the $`\overline{B}^0\overline{K}^0K^+`$ transition matrix element amounts to considering the strong interaction $`\overline{B}{}_{}{}^{0}\overline{K}{}_{}{}^{0}\overline{B}_{s}^{}`$ followed by the weak transition $`\overline{B}_s^{}K^+`$ and the result is Cheng:2002qu
$`\left[K^{}\left(p_3\right)\left|\left(\overline{s}u\right)_{VA}\right|0\overline{K}{}_{}{}^{0}\left(p_1\right)K^+\left(p_2\right)\left|\left(\overline{u}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}\right]_{\mathrm{pole}}`$
$`={\displaystyle \frac{f_K}{f_\pi }}{\displaystyle \frac{g\sqrt{m_Bm_{B_s^{}}}}{s_{23}m_{B_s^{}}^2}}F(s_{23},m_{B_s^{}})F_1^{B_sK}\left(m_3^2\right)\left[m_B+{\displaystyle \frac{s_{23}}{m_B}}m_B{\displaystyle \frac{m_B^2s_{23}}{m_3^2}}\left(1{\displaystyle \frac{F_0^{B_sK}\left(m_3^2\right)}{F_1^{B_sK}\left(m_3^2\right)}}\right)\right]`$
$`\times \left[m_1^2+m_3^2s_{13}+{\displaystyle \frac{\left(s_{23}m_2^2+m_3^2\right)\left(m_B^2s_{23}m_1^2\right)}{2m_{B_s^{}}^2}}\right],`$ (8)
where $`g`$ is a heavy-flavor independent strong coupling which can be extracted from the recent CLEO measurement of the $`D^+`$ decay width, $`g=0.59\pm 0.01\pm 0.07`$ CLEOg , and $`F_{0,1}^{B_sK}`$ are the $`B_sK`$ weak transition from factors in the standard convention BSW . Since $`B_s^{}`$ can be far from the mass shell, it is necessary to introduce a form factor $`F(s_{23},m_{B_s^{}})`$ to take into account the off-shell effect of the $`B_s^{}`$ pole. Following CCS , it is parameterized as $`F(s_{23},m_{B_s^{}})=\left(\mathrm{\Lambda }^2m_{B_s^{}}^2\right)/\left(\mathrm{\Lambda }^2s_{23}\right)`$ with the cut-off parameter $`\mathrm{\Lambda }`$ chosen to be $`\mathrm{\Lambda }=m_{B_s^{}}+\mathrm{\Lambda }_{\mathrm{QCD}}`$.
It is worth making a digression for a moment. In principle, one can apply HMChPT twice to evaluate the form factors $`r,\omega _+`$ and $`\omega _{}`$ LLW . However, this will lead to too large decay rates in disagreement with experiment Cheng:2002qu . This is because the use of HMChPT is reliable only in the kinematic region where $`K^+`$ and $`\overline{K}^0`$ are soft. Therefore, the available phase space where chiral perturbation theory is applicable is very limited. If the soft meson result is assumed to be applicable to the whole Dalitz plot, the decay rate will be greatly overestimated. Therefore, we employ the pole model to evaluate the aforementioned form factors. We shall apply HMChPT only once to the $`\overline{B}{}_{}{}^{0}K{}_{}{}^{0}B_{s}^{}`$ strong vertex and introduce a form factor to take care of the momentum dependence of the strong coupling.
The resonant pole contributions to the form factors $`r`$, $`\omega _\pm `$ and $`h`$ can be worked out from Eq. (II). In principle, there are also nonresonant contributions to these form factors. It turns out that the leading nonresonant contribution can be determined as follows. We notice that the same $`\overline{B}K\overline{K}`$ two-meson transition matrix element also appears in the decay $`B^{}D^0K^0K^{}`$ under factorization DKK . The data favors a $`1^{}`$ configuration in the $`K^0K^{}`$ pair DKKdata . The corresponding two-meson transition matrix element is dominated by the $`\omega _{}`$ term. Following DKK we shall include a nonresonant contribution to $`\omega _{}`$ parametrized as
$`\omega _{}^{NR}=\kappa {\displaystyle \frac{2p_Bp_2}{s_{12}^2}},`$ (9)
and employ the $`B^{}D^0K^0K^{}`$ data and apply isospin symmetry to the $`\overline{B}K\overline{K}`$ matrix elements to determine the unknown parameter $`\kappa `$. The denominator in the above parametrization is inspired by the QCD counting rule which gives rise to a $`1/s_{12}^2`$ asymptotic behavior,<sup>2</sup><sup>2</sup>2As explained in DKK , at least two hard gluon exchanges are needed: one creating the $`s\overline{s}`$ pair in $`\overline{K}^0K^+`$, the other kicking the spectator to catch up with the energetic $`s`$ quark to form the $`K`$ meson. This gives rise to a $`1/s_{12}^2`$ asymptotic behavior. while the numerator $`p_Bp_2=m_BE_{K^+}`$ is motivated by the observation that $`K^+`$ contains an energetic $`u`$ quark coming from the $`bu`$ transition.
The matrix elements involving 3-kaon creation are given by Cheng:2002qu
$`\overline{K}{}_{}{}^{0}\left(p_1\right)K^+\left(p_2\right)K^{}\left(p_3\right)\left|\left(\overline{s}d\right)_{VA}\right|00\left|\left(\overline{d}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}0,`$ (10)
$`\overline{K}{}_{}{}^{0}\left(p_1\right)K^+\left(p_2\right)K^{}\left(p_3\right)\left|\overline{s}\gamma _5d\right|00\left|\overline{d}\gamma _5b\right|\overline{B}{}_{}{}^{0}=v{\displaystyle \frac{f_Bm_B^2}{f_\pi m_b}}(1{\displaystyle \frac{s_{13}m_1^2m_3^2}{m_B^2m_K^2}})F^{KKK}\left(m_B^2\right),`$
where
$`v={\displaystyle \frac{m_{K^+}^2}{m_u+m_s}}={\displaystyle \frac{m_K^2m_\pi ^2}{m_sm_d}},`$ (11)
characterizes the quark-order parameter $`\overline{q}q`$ which spontaneously breaks the chiral symmetry. Both relations in Eq. (10) are originally derived in the chiral limit Cheng:2002qu and hence the quark masses appearing in Eq. (11) are referred to the scale $``$ 1 GeV . The first relation reflects helicity suppression which is expected to be even more effective for energetic kaons. For the second relation, we introduce the form factor $`F^{KKK}`$ to extrapolate the chiral result to the physical region. Following Cheng:2002qu we shall take $`F^{KKK}\left(q^2\right)=1/\left[1\left(q^2/\mathrm{\Lambda }_\chi ^2\right)\right]`$ with $`\mathrm{\Lambda }_\chi =0.83`$ GeV being a chiral symmetry breaking scale.
We now turn to the 2-kaon creation matrix element which can be expressed in terms of time-like kaon current form factors as
$`K^+\left(p_{K^+}\right)K^{}\left(p_K^{}\right)\left|\overline{q}\gamma _\mu q\right|0`$ $`=`$ $`\left(p_{K^+}p_K^{}\right)_\mu F_q^{K^+K^{}},`$
$`K^0\left(p_{K^0}\right)\overline{K}^0\left(p_{\overline{K}^0}\right)\left|\overline{q}\gamma _\mu q\right|0`$ $`=`$ $`\left(p_{K^0}p_{\overline{K}^0}\right)_\mu F_q^{K^0\overline{K}^0}.`$ (12)
The weak vector form factors $`F_q^{K^+K^{}}`$ and $`F_q^{K^0\overline{K}^0}`$ can be related to the kaon electromagnetic (e.m.) form factors $`F_{em}^{K^+K^{}}`$ and $`F_{em}^{K^0\overline{K}^0}`$ for the charged and neutral kaons, respectively. Phenomenologically, the e.m. form factors receive resonant and nonresonant contributions and can be expressed by
$`F_{em}^{K^+K^{}}=F_\rho +F_\omega +F_\varphi +F_{NR},F_{em}^{K^0\overline{K}^0}=F_\rho +F_\omega +F_\varphi +F_{NR}^{}.`$ (13)
It follows from Eqs. (II) and (13) that
$`F_u^{K^+K^{}}`$ $`=`$ $`F_d^{K^0\overline{K}^0}=F_\rho +3F_\omega +{\displaystyle \frac{1}{3}}\left(3F_{NR}F_{NR}^{}\right),`$
$`F_d^{K^+K^{}}`$ $`=`$ $`F_u^{K^0\overline{K}^0}=F_\rho +3F_\omega ,`$
$`F_s^{K^+K^{}}`$ $`=`$ $`F_s^{K^0\overline{K}^0}=3F_\varphi {\displaystyle \frac{1}{3}}\left(3F_{NR}+2F_{NR}^{}\right),`$ (14)
where use of isospin symmetry has been made.
The resonant and nonresonant terms in Eq. (13) can be parametrized as
$`F_h\left(s_{23}\right)={\displaystyle \frac{c_h}{m_h^2s_{23}im_h\mathrm{\Gamma }_h}},F_{NR}^{()}\left(s_{23}\right)=\left({\displaystyle \frac{x_1^{()}}{s_{23}}}+{\displaystyle \frac{x_2^{()}}{s_{23}^2}}\right)\left[\mathrm{ln}\left({\displaystyle \frac{s_{23}}{\stackrel{~}{\mathrm{\Lambda }}^2}}\right)\right]^1,`$ (15)
with $`\stackrel{~}{\mathrm{\Lambda }}0.3`$ GeV. The expression for the nonresonant form factor is motivated by the asymptotic constraint from pQCD, namely, $`F\left(t\right)\left(1/t\right)\left[\mathrm{ln}\left(t/\stackrel{~}{\mathrm{\Lambda }}^2\right)\right]^1`$ in the large $`t`$ limit Brodsky . The unknown parameters $`c_h`$, $`x_i`$ and $`x_i^{}`$ are fitted from the kaon e.m. data, giving the best fit values (in units of GeV<sup>2</sup> for $`c_h`$) DKK :
$$\begin{array}{ccc}c_\rho =3c_\omega =c_\varphi =0.363,\hfill & c_{\rho \left(1450\right)}=7.98\times 10^3,\hfill & c_{\rho \left(1700\right)}=1.71\times 10^3,\hfill \\ c_{\omega \left(1420\right)}=7.64\times 10^2,\hfill & c_{\omega \left(1650\right)}=0.116,\hfill & c_{\varphi \left(1680\right)}=2.0\times 10^2,\hfill \end{array}$$
(16)
and
$`x_1=3.26\mathrm{GeV}^2,x_2=5.02\mathrm{GeV}^4,x_1^{}=0.47\mathrm{GeV}^2,x_2^{}=0.`$ (17)
Note that the form factors $`F_{\rho ,\omega ,\varphi }`$ in Eqs. (13) and (14) include the contributions from the vector mesons $`\rho \left(770\right),\rho \left(1450\right),\rho \left(1700\right)`$, $`\omega \left(782\right),\omega \left(1420\right),\omega \left(1650\right),`$ $`\varphi \left(1020\right)`$ and $`\varphi \left(1680\right)`$. It is interesting to note that (i) the fitted values of $`c_V`$ are very close to the vector meson dominance expression $`g_{_{V\gamma }}g_{VKK}`$ for $`V=\rho ,\omega ,\varphi `$ DM2 ; PDG , where $`g_{_{V\gamma }}`$ is the e.m. coupling of the vector meson defined by $`V\left|j_{em}\right|0=g_{V\gamma }ϵ_V^{}`$ <sup>3</sup><sup>3</sup>3The vector meson e.m. couplings are given by $`g_{\varphi \gamma }=e_sm_\varphi f_\varphi `$, $`g_{\rho \gamma }=[(e_ue_d)/\sqrt{2}]m_\rho f_\rho `$ and $`g_{\omega \gamma }=[(e_u+e_d)/\sqrt{2}]m_\omega f_\omega `$ where $`e_q`$ is the quark’s charge and $`f_V`$ is the vector decay constant. and $`g_{VKK}`$ is the $`VKK`$ strong coupling with, $`g_{\varphi K^+K^{}}g_{\rho K^+K^{}}/\sqrt{2}=g_{\omega K^+K^{}}/\sqrt{2}3.03`$, and (ii) the vector-meson pole contributions alone yield $`F_{u,s}^{K^+K^{}}\left(0\right)1,1`$ and $`F_d^{K^+K^{}}\left(0\right)0`$ as the charged kaon does not contain the valence $`d`$ quark. <sup>4</sup><sup>4</sup>4The sign convention is fixed by using $`M(q\overline{q}^{},p)\overline{M}(q\overline{q}^{},p^{})|\overline{q}\gamma _\mu q|0=M(q\overline{q}^{},p)|\overline{q}\gamma _\mu q|M(q\overline{q}^{},p^{})=(pp^{})_\mu |F_q^{MM}|`$ in the case of a real $`F_q^{MM}`$. The matrix element in the decay amplitude relevant for our purpose then has the expression
$`\overline{K}{}_{}{}^{0}\left(p_1\right)\left|\left(\overline{s}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}K^+\left(p_2\right)K^{}\left(p_3\right)\left|\left(\overline{q}q\right)_{VA}\right|0=(s_{12}s_{13})F_1^{BK}\left(s_{23}\right)F_q^{K^+K^{}}\left(s_{23}\right).`$ (18)
We also need to specify the 2-body matrix element $`K^+K^{}\left|\overline{s}s\right|0`$ induced from the scalar density. It receives resonant and non-resonant contributions:
$`K^+\left(p_2\right)K^{}\left(p_3\right)\left|\overline{s}s\right|0`$ $``$ $`f_s^{K^+K^{}}\left(s_{23}\right)={\displaystyle \underset{i}{}}{\displaystyle \frac{m_i\overline{f}_ig^{iKK}}{m_i^2s_{23}im_i\mathrm{\Gamma }_i}}+f_s^{NR},`$
$`f_s^{NR}`$ $`=`$ $`{\displaystyle \frac{v}{3}}\left(3F_{NR}+2F_{NR}^{}\right)+v{\displaystyle \frac{\sigma }{s_{23}^2}}\left[\mathrm{ln}\left({\displaystyle \frac{s_{23}}{\stackrel{~}{\mathrm{\Lambda }}^2}}\right)\right]^1,`$ (19)
where the scalar decay constant $`\stackrel{~}{f}_i`$ is defined in $`i\left|\overline{s}s\right|0=m_i\overline{f}_i`$, $`g^{iKK}`$ is the $`iKK`$ strong coupling, and the nonresonant terms are related to those in $`F_s^{K^+K^{}}`$ through the equation of motion.<sup>5</sup><sup>5</sup>5The use of equations of motion also leads to $`f_s^{K^+K^{}}=vF_s^{K^+K^{}}.`$ (20) Note that the pole contribution to $`F_s^{K^+K^{}}`$ should be dropped in the above relation as it applies only to nonresonant contributions. The main scalar meson pole contributions are those that have dominant $`s\overline{s}`$ content and large coupling to $`K\overline{K}`$. It is found in ANS that among the $`f_0`$ mesons, only $`f_0\left(980\right)`$ and $`f_0\left(1530\right)`$ have the largest couplings with the $`K\overline{K}`$ pair. Note that $`f_0\left(1530\right)`$ is a very broad state with the width of order 1 GeV ANS . To proceed with the numerical calculations, we use $`g^{f_0\left(980\right)KK}=1.5`$ GeV, $`g^{f_0\left(1530\right)KK}=3.18`$ GeV, $`\mathrm{\Gamma }_{f_0\left(980\right)}=80`$ MeV, $`\mathrm{\Gamma }_{f_0\left(1530\right)}=1.160`$ GeV ANS , $`\overline{f}_{f_0\left(980\right)}\left(\mu =m_b/2\right)0.39`$ GeV Cheng:2005ye and $`\overline{f}_{f_0\left(1530\right)}\overline{f}_{f_0\left(980\right)}`$. The sign of the resonant terms is fixed by $`f_s^{K^+K^{}}\left(0\right)=v`$ from a chiral perturbation theory calculation (see, for example, Cheng:1988va ). It should be stressed that although the nonresonant contributions to $`f_s^{KK}`$ and $`F_s^{KK}`$ are related through the equation of motion, the resonant ones are different and not related a priori. To apply the equation of motion, the form factors should be away from the resonant region. In the large $`s_{23}`$-region, the nonresonant contribution dominated by the $`1/s_{23}`$ term is far away from the resonant one. In contrast, the $`1/s_{23}^2`$ term dominates in the low $`s_{23}`$-region where resonant contributions cannot be ignored. The $`1/s_{23}^2`$ term in $`F_s`$ is not necessarily conveyed to $`f_S`$ through the equation of motion. Hence, the $`1/s_{23}^2`$ term in Eq. (II) is undetermined and a new parameter $`\sigma `$, which is expected to be of similar size as $`x_2`$, is assigned and will be determined later by fitting to the data. The corresponding matrix element is now given by
$`\overline{K}{}_{}{}^{0}\left(p_1\right)\left|\overline{s}b\right|\overline{B}{}_{}{}^{0}K^+\left(p_2\right)K^{}\left(p_3\right)\left|\overline{s}s\right|0={\displaystyle \frac{m_B^2m_K^2}{m_bm_s}}F_0^{BK}\left(s_{23}\right)f_s^{K^+K^{}}\left(s_{23}\right).`$ (21)
Collecting all the relevant matrix elements evaluated above, we are ready to compute the amplitude $`A\left(\overline{B}{}_{}{}^{0}K_{S\left(L\right)}K^+K^{}\right)=\pm A\left(\overline{B}{}_{}{}^{0}\overline{K}{}_{}{}^{0}K_{}^{+}K^{}\right)/\sqrt{2}`$. Since under CP-conjugation we have $`K_S\left(\stackrel{}{p}_1\right)K_S\left(\stackrel{}{p}_1\right)`$, $`K^+\left(\stackrel{}{p}_2\right)K^{}\left(\stackrel{}{p}_2\right)`$ and $`K^{}\left(\stackrel{}{p}_3\right)K^+\left(\stackrel{}{p}_3\right)`$, the $`\overline{B}{}_{}{}^{0}K_SK^+K^{}`$ amplitude can be decomposed into CP-odd and CP-even components
$`A\left[\overline{B}{}_{}{}^{0}K_S\left(p_1\right)K^+\left(p_2\right)K^{}\left(p_3\right)\right]=A(s_{12},s_{13},s_{23})=A_{CP}+A_{CP+},`$
$`A_{CP\pm }={\displaystyle \frac{1}{2}}\left[A(s_{12},s_{13},s_{23})\pm A(s_{13},s_{12},s_{23})\right].`$ (22)
Correspondingly, we have
$`\mathrm{\Gamma }`$ $`=`$ $`\mathrm{\Gamma }_{CP+}+\mathrm{\Gamma }_{CP},`$
$`\mathrm{\Gamma }_{CP\pm }`$ $`=`$ $`{\displaystyle \frac{1}{\left(2\pi \right)^3}}{\displaystyle \frac{1}{32m_B^3}}{\displaystyle \left|A_{CP\pm }\right|^2𝑑s_{12}𝑑s_{13}}={\displaystyle \frac{1}{\left(2\pi \right)^3}}{\displaystyle \frac{1}{32m_B^3}}{\displaystyle \left|A_{CP\pm }\right|^2𝑑s_{12}𝑑s_{23}}.`$ (23)
The vanishing cross terms due to the interference between CP-odd and CP-even components can be easily seen from the (anti)symmetric properties of the amplitude and the integration variables under the interchange of $`s_{12}s_{13}`$. Similar relations hold for the conjugated $`B^0`$ decay rate $`\overline{\mathrm{\Gamma }}`$. The $`CP`$-even fraction $`f_+`$ is defined by
$`f_+{\displaystyle \frac{\mathrm{\Gamma }_{CP+}+\overline{\mathrm{\Gamma }}_{CP+}}{\mathrm{\Gamma }+\overline{\mathrm{\Gamma }}}}|_{\varphi K_S\mathrm{excluded}.}`$ (24)
Note that results for the $`K^+K^{}K_L`$ mode are identical to the $`K^+K^{}K_S`$ ones with the CP eigenstates interchanged. For example, results for $`\left(K^+K^{}K_L\right)_{CP+}`$ are the same as those for $`\left(K^+K^{}K_S\right)_{CP}`$ and hence $`f_+`$ in $`K^+K^{}K_S`$ corresponds to $`f_{}`$ in $`K^+K^{}K_L`$.
We next turn to the $`\overline{B}{}_{}{}^{0}K_SK_SK_S,K_SK_SK_L`$ decays. The decay amplitudes are given by
$`A\left[\overline{B}{}_{}{}^{0}K_S\left(p_1\right)K_S\left(p_2\right)K_{S,L}\left(p_3\right)\right]`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}\right)^{3/2}\{\pm A\left[\overline{B}{}_{}{}^{0}K^0\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_2\right)\overline{K}{}_{}{}^{0}\left(p_3\right)\right]`$ (25)
$`\pm A\left[\overline{B}{}_{}{}^{0}K^0\left(p_2\right)\overline{K}{}_{}{}^{0}\left(p_3\right)\overline{K}{}_{}{}^{0}\left(p_1\right)\right]`$
$`+A\left[\overline{B}{}_{}{}^{0}K^0\left(p_3\right)\overline{K}{}_{}{}^{0}\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_2\right)\right]\},`$
with
$`A\left[\overline{B}{}_{}{}^{0}K^0\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_2\right)\overline{K}{}_{}{}^{0}\left(p_3\right)\right]`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}{\displaystyle \underset{p=u,c}{}}\lambda _p\{[K^0\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_2\right)\left|\left(\overline{d}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}\overline{K}{}_{}{}^{0}\left(p_3\right)\left|\left(\overline{s}d\right)_{VA}\right|0`$ (26)
$`+K^0\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_3\right)\left|\left(\overline{d}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}\overline{K}{}_{}{}^{0}\left(p_2\right)\left|\left(\overline{s}d\right)_{VA}\right|0]`$
$`\times \left(a_4^p+{\displaystyle \frac{1}{2}}a_{10}^p\left(a_6^p{\displaystyle \frac{1}{2}}a_8^p\right)r_\chi \right)`$
$`+[\overline{K}{}_{}{}^{0}\left(p_2\right)\left|\overline{s}b\right|\overline{B}{}_{}{}^{0}K^0\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_3\right)\left|\overline{s}s\right|0`$
$`+\overline{K}{}_{}{}^{0}\left(p_3\right)\left|\overline{s}b\right|\overline{B}{}_{}{}^{0}K^0\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_2\right)\left|\overline{s}s\right|0\left]\right(2a_6^p+a_8^p)`$
$`+K^0\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_2\right)\overline{K}{}_{}{}^{0}\left(p_3\right)\left|\overline{s}\gamma _5d\right|00\left|\overline{d}\gamma _5b\right|\overline{B}{}_{}{}^{0}(2a_6^p+a_8^p)`$
$`+[\overline{K}{}_{}{}^{0}\left(p_2\right)\left|\left(\overline{s}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}K^0\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_3\right)\left|\left(\overline{s}s\right)_{VA}\right|0`$
$`+\overline{K}{}_{}{}^{0}\left(p_3\right)\left|\left(\overline{s}b\right)_{VA}\right|\overline{B}{}_{}{}^{0}K^0\left(p_1\right)\overline{K}{}_{}{}^{0}\left(p_2\right)\left|\left(\overline{s}s\right)_{VA}\right|0]`$
$`\times [a_3+a_4^p+a_5{\displaystyle \frac{1}{2}}(a_7+a_9+a_{10})]\},`$
where the last term will not contribute to the purely CP-even decay $`\overline{B}{}_{}{}^{0}K_SK_SK_S`$. Decay rates for the $`K_SK_SK_S`$ and $`K_SK_SK_L`$ modes can be obtained from Eq. (23) with an additional factor of $`1/3!`$ and $`1/2!`$, respectively, for identical particles in the final state.
We now consider the CP asymmetries for $`\overline{B}{}_{}{}^{0}K^+K^{}K_{S\left(L\right)},K_SK_SK_{S\left(L\right)}`$ decays. The direct CP asymmetry and the mixing induced CP violation are defined by
$`𝒜_{KKK}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }\overline{\mathrm{\Gamma }}}{\mathrm{\Gamma }+\overline{\mathrm{\Gamma }}}}`$
$`=`$ $`{\displaystyle \frac{\left|A\right|^2𝑑s_{12}𝑑s_{23}\left|\overline{A}\right|^2𝑑s_{12}𝑑s_{23}}{\left|A\right|^2𝑑s_{12}𝑑s_{23}+\left|\overline{A}\right|^2𝑑s_{12}𝑑s_{23}}},`$
$`𝒮_{KKK,CP\pm }`$ $`=`$ $`{\displaystyle \frac{2\mathrm{Im}\left(e^{2i\beta }A_{CP\pm }\overline{A}_{CP\pm }^{}\right)𝑑s_{12}𝑑s_{23}}{\left|A_{CP\pm }\right|^2𝑑s_{12}𝑑s_{23}+\left|\overline{A}_{CP\pm }\right|^2𝑑s_{12}𝑑s_{23}}},`$
$`𝒮_{KKK}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{Im}\left(e^{2i\beta }A\overline{A}^{}\right)𝑑s_{12}𝑑s_{23}}{\left|A\right|^2𝑑s_{12}𝑑s_{23}+\left|\overline{A}\right|^2𝑑s_{12}𝑑s_{23}}}`$ (27)
$`=`$ $`f_+S_{KKK,CP+}+\left(1f_+\right)S_{KKK,CP},`$
where $`\overline{A}`$ is the decay amplitude of $`B^0K^+K^{}K_{S\left(L\right)}`$ or $`K_SK_SK_{S\left(L\right)}`$. For the $`K^+K^{}K_S`$ mode, it is understood that the contribution from $`\varphi K_S`$ is excluded. It is expected in the SM that $`𝒮_{KKK,CP+}\mathrm{sin}2\beta _{\mathrm{eff}}\mathrm{sin}2\beta `$, $`𝒮_{KKK,CP}\mathrm{sin}2\beta `$ and hence $`𝒮_{KKK}\left(2f_+1\right)\mathrm{sin}2\beta `$.<sup>6</sup><sup>6</sup>6Writing the CP-conjugated decay amplitude as $`\overline{A}=\overline{A}_{CP+}+\overline{A}_{CP}`$, we have $`\overline{A}_{CP\pm }=\pm A_{CP\pm }`$ with $`\lambda _p\lambda _p^{}`$. This leads to $`𝒮_{KKK,CP}𝒮_{KKK,CP+}`$.
## III Numerical results and discussions
To proceed with the numerical calculations, we need to specify the input parameters. For the CKM matrix elements, we use the Wolfenstein parameters $`A=0.825`$, $`\lambda =0.22622`$, $`\overline{\rho }=0.207`$ and $`\overline{\eta }=0.340`$, corresponding to $`\left(\mathrm{sin}2\beta \right)_{CKM}=0.724`$ CKMfitter . For $`BK`$ form factors we shall use those derived in the covariant light-front quark model CCH with the assigned error to be $`0.03`$, namely, $`F_{0,1}^{BK}\left(0\right)=0.35\pm 0.03`$. The parameter $`\kappa `$ in Eq. (9) is determined from the $`B^{}D^0K^0K^{}`$ data. From the measured branching ratio $`\left(B^{}D^0K^0K^{}\right)=\left(5.5\pm 1.4\pm 0.8\right)\times 10^4`$ DKKdata , we obtain $`\kappa =3.1_{1.8}^{+5.1}`$ GeV where use of $`a_1^{DKK}=0.935`$ and $`a_2^{DKK}\left(a_2^{D\rho }\right)=0.4\pm 0.2`$ has been made DKK . For the quark masses and the unitarity angle $`\gamma `$, we shall use $`m_b\left(m_b\right)=4.2`$ GeV, $`m_s\left(m_b/2\right)=80\pm 20`$ MeV and $`\gamma =\left(58.6\pm 7\right)^{}`$ CKMfitter . The $`K_SK_SK_S`$ rate sensitive to the parameter $`\sigma `$ in Eq. (II) is used to determine $`\sigma =\left(10.4_{4.8}^{+5.4}\right)`$GeV<sup>4</sup>, where the errors include the uncertainties in the $`K_SK_SK_S`$ decay rate, the strange quark mass and the $`F_0^{BK}`$ form factor.
Results for the decay rates and CP asymmetries in $`\overline{B}{}_{}{}^{0}K^+K^{}K_{S\left(L\right)},K_SK_SK_{S\left(L\right)}`$ are exhibited in Table 2 and Table 3, respectively. The theoretical errors shown there are from the uncertainties in (i) the parameter $`\kappa `$ which governs the nonresonant contribution to the form factor $`\omega _{}`$ \[see Eq. (9)\], (ii) the strange quark mass $`m_s`$, the form factor $`F_0^{BK}`$ and $`\sigma `$ \[see Eq. (II)\] constrained from the $`K_SK_SK_S`$ rate, and (iii) the unitarity angle $`\gamma `$. To compute the CP-even fraction $`f_+`$ and $`\mathrm{sin}2\beta _{\mathrm{eff}}`$ for $`K^+K^{}K_S`$, we need to turn off the coefficient $`c_\varphi `$ in Eq. (14). As one can see from Table 2, the predicted rates for $`\overline{B}{}_{}{}^{0}K^+K^{}K_{S\left(L\right)}`$ decays and the $`CP`$-even (odd) ratio $`f_{+()}`$ are in accordance with the data within errors, though the theoretical central values on rates are somewhat smaller than the experimental ones. Theoretical errors on the branching ratios are dominated by the sizable error in $`\kappa `$ and the uncertainty in the strange quark mass as the penguin term $`a_6r_\chi `$ and the parameter $`v`$ are very sensitive to $`m_s`$. Note that the second error in rates (including the contribution from the uncertainty in $`\sigma `$) are constrained from the $`K_SK_SK_S`$ rate and hence are reduced significantly. For the first error, we note that the larger the value of $`\left|\kappa \right|`$ we have, the larger rate on CP-odd $`K^+K^{}K_S`$ is obtained, leading to a smaller value of $`f_+\left(K^+K^{}K_S\right)`$. Since the central value of our $`f_+\left(K^+K^{}K_S\right)`$ agrees well with data, $`\kappa `$ is preferred to be around its central value.
The $`K^+K^{}`$ mass spectra of the $`\overline{B}{}_{}{}^{0}K^+K^{}K_S`$ decay from $`CP`$-even and $`CP`$-odd contributions are shown in Fig. 1. In the spectra, there are peaks at the threshold and a milder one in the large $`m_{K^+K^{}}`$ region. For the $`CP`$-even part, the threshold enhancement arises from the $`f_0\left(980\right)K_S`$ and the nonresonant $`f_S^{K^+K^{}}`$ contributions \[see Eq. (II)\], while the peak at large $`m_{K^+K^{}}`$ comes from the nonresonant two-meson transition $`\overline{B}{}_{}{}^{0}K^+K_S`$ followed by a current produced $`K^{}`$. Since the nonresonant term \[Eq. (9)\] favors a small $`m_{K^+K_s}`$ region, the spectrum should peak at the large $`m_{K^+K^{}}`$ end. For the $`CP`$-odd spectrum the bump at the large $`m_{K^+K^{}}`$ end originates from the same two-meson transition term, while the peak on the lower end corresponds to the $`\varphi K_s`$ contribution, which is also shown in the insert. The full $`K^+K^{}K_S`$ spectrum is basically the sum of the $`CP`$-even and the $`CP`$-odd parts. Note that although we include $`f_0\left(1530\right)K_S`$ contribution, its effect is not as prominent as one may expect from the $`K^{}K^+K^{}`$ spectrum where a large $`f_X\left(1500\right)K^{}`$ contribution is found Garmash:2004wa .
For the mixing-induced CP asymmetry in the $`K^+K^{}K_S`$ mode, we compute the effective $`\mathrm{sin}2\beta `$ in two different ways: In one way, we calculate $`S`$ with $`\varphi K_S`$ excluded in $`K^+K^{}K_S`$ and then apply the relation $`S=\left(2f_+1\right)\mathrm{sin}2\beta _{\mathrm{eff}}`$ and the theoretical value of $`f_+`$ to obtain $`\mathrm{sin}2\beta _{\mathrm{eff}}`$. This procedure follows closely the BaBar and Belle method of measuring the effective $`\mathrm{sin}2\beta `$. In the other way, we calculate $`S`$ directly for the CP-even $`K^+K^{}K_S`$ and identify $`S_{KKK,CP+}`$ with $`\mathrm{sin}2\beta _{\mathrm{eff}}`$. As for the $`K_SK_SK_S`$ mode, there is no such ambiguity as it is a purely CP-even state. As shown in Table 3 and Fig. 2, the resulting $`\mathrm{sin}2\beta _{\mathrm{eff}}`$ is slightly different in these two different approaches.
The deviation of the mixing-induced CP asymmetry in $`B^0K^+K^{}K_S`$ and $`K_SK_SK_S`$ from that measured in $`BJ/\psi K_S`$ (or the fitted CKM’s $`\mathrm{sin}2\beta `$ CKMfitter ), namely, $`\mathrm{\Delta }\mathrm{sin}2\beta _{\mathrm{eff}}\mathrm{sin}2\beta _{\mathrm{eff}}\mathrm{sin}2\beta _{J/\psi K_S\left(CKM\right)}`$, is calculated from Table 3 to be
$`\mathrm{\Delta }\mathrm{sin}2\beta _{K^+K^{}K_S}=0.06_{0.02}^{+0.08}\left(0.02_{0.02}^{+0.08}\right),\mathrm{\Delta }\mathrm{sin}2\beta _{K_SK_SK_S}=0.06_{0.00}^{+0.00}\left(0.02_{0.00}^{+0.00}\right).`$ (28)
Note that part of the deviation comes from that between the measured $`\mathrm{sin}2\beta _{J/\psi K_S}`$ and the fitted CKM’s $`\mathrm{sin}2\beta `$. The $`K^+K^{}K_S`$ has a potentially sizable $`\mathrm{\Delta }\mathrm{sin}2\beta `$, as this penguin-dominated mode is subject to a tree pollution due to the presence of color-allowed tree contributions. For the $`K_SK_SK_S`$ mode, the central value and the error on $`\mathrm{\Delta }\mathrm{sin}2\beta `$ are small.
It is instructive to see the dependence of $`\mathrm{sin}2\beta _{\mathrm{eff}}`$ on the $`K^+K^{}`$ invariant mass, $`m_{K^+K^{}}m_{23}=\sqrt{s_{23}}`$. For the phase space integration in Eq. (II), for a given $`s_{23}`$, the upper and lower bounds of $`s_{12}`$ are fixed. The invariant mass $`m_{23}`$ is integrated from $`m_{23}^{}=m_2+m_3`$ to $`m_{23}^+=m_Bm_1`$. When the variable $`s_{23}`$ or $`m_{23}`$ is integrated from $`m_{23}^{}`$ to a fixed $`m_{23}^{\mathrm{max}}`$ (of course, $`m_{23}^{}<m_{23}^{\mathrm{max}}m_{23}^+`$), the effective $`\mathrm{sin}2\beta `$ thus obtained is designated as $`\mathrm{sin}2\beta _{\mathrm{eff}}\left(m_{23}^{\mathrm{max}}\right)`$. Fig. 2 shows the plot of $`\mathrm{sin}2\beta _{\mathrm{eff}}\left(m_{K^+K^{}}^{\mathrm{max}}\right)`$ versus $`m_{K^+K^{}}^{\mathrm{max}}`$ for $`K^+K^{}K_S`$. Since there are two different methods for the determination of $`\mathrm{sin}2\beta _{\mathrm{eff}}`$, the results are depicted in two different curves. It is interesting that $`\mathrm{sin}2\beta \left(m_{23}^{\mathrm{max}}\right)`$ is slightly below $`\mathrm{sin}2\beta _{CKM}`$ at the bulk of the $`m_{K^+K^{}}`$ region and gradually increases and becomes slightly larger than $`\mathrm{sin}2\beta _{CKM}`$ when the phase space is getting saturated. The deviation $`\mathrm{\Delta }\mathrm{sin}2\beta _{K^+K^{}K_S}`$ arises mainly from the large $`m_{K^+K^{}}`$ region.
Direct CP violation is found to be very small in both $`K^+K^{}K_S`$ and $`K_SK_SK_S`$ modes. It is interesting to notice that direct CP asymmetry in the CP-even $`K^+K^{}K_S`$ mode is only of order $`10^3`$, but it becomes $`0.2\times 10^2`$ in $`K^+K^{}K_S`$ with $`\varphi K_S`$ excluded. Since these direct CP asymmetries are so small they can be used as approximate null tests of the SM.
Since direct CP violation in charmless $`B`$ decays can be significantly affected by final-state rescattering CCS , we have studied to what extent indications of possibly large deviations of the mixing-induced CP violation seen in the penguin-induced two-body decay modes from $`\mathrm{sin}2\beta `$ determined from $`BJ/\psi K_S`$ can be accounted for by final-state interactions CCSsin2beta . It is natural to extend the study of final-state rescattering effect on time-dependent CP asymmetries to $`BKKK_S`$ decays. Final-state interactions in three-body decays are expected to be much more complicated than the two-body case. For example, the color allowed tree decay $`\overline{B}{}_{}{}^{0}D_s^{()+}D^{()}`$ can rescatter into a $`K^+K^{}K_S`$ final state, where we have $`D_s^{()+}K^+\overline{D}^0`$, $`D^{()}K_SD_s^{}`$ followed by a $`\overline{D}^0D_s^{}K^{}`$ fusion. These diagrams are too complicated and will not be included in this study.<sup>7</sup><sup>7</sup>7In passing we note that these diagrams could have the effect of increasing somewhat our predictions for the rates of 3$`K`$ final states. Although these contributions carry negligible CP-odd (weak) phases, they also contribute to the strong phases and hence will tend to dilute our prediction on $`\mathrm{\Delta }𝒮`$ but not necessarily on direct CP asymmetries. Nevertheless, we attempt to incorporate final state rescattering effects in a simple way by including resonance contributions to the corresponding kaon pairs in the final state hewett-res . We note that another attempt in this direction has recently been made by Furman et al. Furman . They considered rescattering of $`\pi \pi `$ and $`K\overline{K}`$ pairs in the $`\pi \pi `$ effective mass range from threshold to 1.1 GeV. While their predicted direct CP asymmetry is very small, the parameter $`𝒮`$ is found to be $`0.64`$ or $`0.77`$, depending on the set of penguin amplitudes. However, due to the limitation on phase space, the calculated branching ratios of order $`1\times 10^6`$ for $`K^+K^{}K_S`$ and $`K_SK_SK_S`$ are only small portions of the total experimental rates (see Table 1) and, consequently, the predictions of $`S`$ may be affected when the whole phase space is taken into consideration.
## IV Conclusions
In the present work we have studied the decay rates and time-dependent CP asymmetries in the decays $`B^0K^+K^{}K_{S\left(L\right)}`$ and $`K_SK_SK_{S\left(L\right)}`$ within the framework of factorization. Our main results are as follows:
1. Resonant and nonresonant contributions to the hadronic matrix elements are carefully investigated. We incorporate final state rescattering effects in a simple way by including resonance contributions to the corresponding kaon pairs in the final state. Instead of applying heavy meson chiral perturbation theory to the matrix element for $`BKK`$, which is valid only for a small kinematic region, we consider the resonant contribution from the $`B_s^{}`$ pole and nonresonant contributions constrained by QCD counting rules.
2. Using the $`K_SK_SK_S`$ decay rate as an input, the predicted branching ratio of $`K^+K^{}K_{S\left(L\right)}`$ modes and the CP-even (-odd) fraction of $`B^0K^+K^{}K_{S\left(L\right)}`$ are consistent with the data within the theoretical and experimental errors, though the theoretical central values on rates are somewhat smaller than the experimental ones.
3. Owing to the presence of color-allowed tree contributions in $`B^0K^+K^{}K_{S\left(L\right)}`$, this penguin-dominated mode is subject to a potentially significant tree pollution and the deviation of the mixing-induced CP asymmetry from that measured in $`BJ/\psi K_S`$, namely, $`\mathrm{\Delta }\mathrm{sin}2\beta _{K^+K^{}K_{S\left(L\right)}}\mathrm{sin}2\beta _{K^+K^{}K_{S\left(L\right)}}\mathrm{sin}2\beta _{J/\psi K_S}`$, can be as large as $`𝒪\left(0.10\right)`$. The deviation $`\mathrm{\Delta }\mathrm{sin}2\beta _{K^+K^{}K_{S\left(L\right)}}`$ arises mainly from the large $`m_{K^+K^{}}`$ region.
4. The $`K_SK_SK_{S\left(L\right)}`$ mode appears theoretically very clean in our picture: The uncertainties in $`\mathrm{\Delta }\mathrm{sin}2\beta _{\mathrm{eff}}`$ are negligible.
5. Direct CP asymmetries are very small in both $`K^+K^{}K_{S\left(L\right)}`$ and $`K_SK_SK_{S\left(L\right)}`$ modes.
###### Acknowledgements.
We wish to thank Kai-Feng Chen for discussion. This research was supported in part by the National Science Council of R.O.C. under Grant Nos. NSC93-2112-M-001-043, NSC93-2112-M-001-053 and by the U.S. DOE contract No. DE-AC02-98CH10886(BNL).
Note added: After the paper was submitted for publication, BaBar has presented a Dalitz plot study of $`B^0K^+K^{}K_S^0`$ decays BaBarKKKsEPS . The BaBar results constrain the tree contribution (incorporated via Eq. (2.18) in the present work) in rates and, as a result, a small $`\mathrm{\Delta }\mathrm{sin}2\beta _{K^+K^{}K_S}`$ is preferable.
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# Primarily quasilocal fields and 1-dimensional abstract local class field theory
## 1. Introduction and statements of the main results
This paper is concerned with finite abelian extensions of primarily quasilocal (abbr, PQL) fields, and can be viewed as a continuation of , I. When $`E`$ is a strictly PQL-field, it shows that these extensions and their norm groups are related as in the fundamental correspondence of the classical local class field theory (see , page 101). The paper proves that they are subject to an exact analogue to the local reciprocity law (formulated, e.g., in , Ch. 6, Theorem 8, and , Theorem 7.1), and to a partial analogue to the local Hasse symbol (as characterized in , Ch. 2, see also , and , Theorems 6.9 and 6.10) of form determined by invariants of the Brauer group Br$`(E)`$. It takes a step towards characterizing the fields whose finite abelian extensions have the above properties. When $`E`$ belongs to some special classes of traditional interest, the present research enables one to achieve this aim and to find a fairly complete description of the norm groups of arbitrary finite separable extensions of $`E`$ (see Section 3, and the references in , Remark 3.9).
The basic notation, terminology and conventions kept in this paper are standard and virtually the same as in , I, and . Throughout, $``$ denotes the set of prime numbers and every algebra $`A`$ is understood to be associative with a unit lying in the considered subalgebras of $`A`$. Simple algebras are supposed to be finite-dimensional over their centres, Brauer and value groups are written additively, Galois groups are viewed as profinite with respect to the Krull topology, and the considered profinite group homomorphisms are continuous. For each field $`E`$, $`E^{}`$ denotes its multiplicative group, $`E_{\mathrm{sep}}`$ a separable closure of $`E`$, $`𝒢_E=𝒢(E_{\mathrm{sep}}/E)`$ is the absolute Galois group of $`E`$, $`s(E)`$ stands for the class of central simple $`E`$-algebras, $`[B]`$ denotes the similarity class of any $`Bs(E)`$, $`d(E)`$ is the subclass of division algebras $`Ds(E)`$, and $`P(E)=\{p:E(p)E\}`$, where $`E(p)`$ is the maximal $`p`$-extension of $`E`$ in $`E_{\mathrm{sep}}`$. As usual, $`E`$ is said to be formally real, if $`1`$ is not presentable over $`E`$ as a finite sum of squares; $`E`$ is called nonreal, otherwise. We say that $`E`$ is Pythagorean, if it is formally real and the set $`E^2=\{\lambda ^2:\lambda E^{}\}`$ is additively closed. For any field extension $`F/E`$, $`\rho _{E/F}`$ denotes the scalar extension map of Br$`(E)`$ into Br$`(F)`$, Br$`(F/E)`$ the relative Brauer group of $`F/E`$, and $`I(F/E)`$ is the set of intermediate fields of $`F/E`$. When $`F/E`$ is finite and separable, Cor<sub>F/E</sub> stands for the corestriction map of Br$`(F)`$ into Br$`(E)`$ (see ). We say that $`E`$ is $`p`$-quasilocal, for some $`p`$, if it satisfies one of the following conditions: (i) every cyclic degree $`p`$ extension of $`E`$ is embeddable as a subalgebra in each $`\mathrm{\Delta }d(E)`$ of (Schur) index $`p`$; (ii) the $`p`$-component Br$`(E)_p`$ of Br$`(E)`$ is trivial or $`pP(E)`$. The field $`E`$ is called PQL, if it is $`p`$-quasilocal for every $`p`$; when this holds, $`E`$ is said to be strictly PQL in case Br$`(E)_p\{0\}`$, $`pP(E)`$. We say that $`E`$ is quasilocal, if its finite extensions are PQL.
Local fields and $`p`$-adically closed fields are strictly quasilocal (abbr, SQL), i.e. their finite extensions are strictly PQL (cf. , Ch. XIII, Sect. 3, , Theorem 3.1 and Lemma 2.9, and , Sect. 3). The strictly PQL-property has been fully characterized in the following two classes: (i) algebraic extensions of global fields ; (ii) Henselian discrete valued fields , Sect. 2. These facts extend the arithmetic basis of this research which is further motivated by results of , I, on Brauer groups of PQL-fields and on absolute Galois groups of quasilocal fields. For reasons clarified in the sequel, our approach to the main topic of this paper is purely algebraic. Our starting point is the fact that residue fields of Henselian valued stable fields are PQL in the case of totally indivisible value groups, in the sense of (2.1) (i). Specifically, a Henselian discrete valued field $`(K,v)`$ with a perfect residue field $`\widehat{K}`$ is stable if and only if $`\widehat{K}`$ is stable and PQL , Proposition 2 (cf. also , I, and , Proposition 2.3). These and other related results show, with their proofs, that PQL-fields partly resemble local fields in a number of respects (see Propositions 2.2-2.3, , I, Lemma 4.3 and , II, Lemma 2.3). For example, by , I, Lemma 4.3, if $`\widehat{K}`$ is $`p`$-quasilocal and $`\stackrel{~}{L}_1`$ and $`\stackrel{~}{L}_2`$ are different extensions of $`\widehat{K}`$ in $`\widehat{K}(p)`$ of degree $`p`$, then the inner product $`N(\stackrel{~}{L}_1/\widehat{K})N(\stackrel{~}{L}_2/\widehat{K})`$ of the norm groups $`N(\stackrel{~}{L}_i/\widehat{K})`$, $`i=2`$, is equal to $`\widehat{K}^{}`$. This result and its key role in the proof of , I, Theorem 4.1, attract interest in the study of the PQL-property along the lines of the classical local class field theory, with the notion of a local field extended as follows:
###### Definition 1.
Let $`E`$ be a field, Nr$`(E)`$ the set of norm groups of $`E`$, and $`\mathrm{\Omega }(E)`$ the set of finite abelian extensions of $`E`$ in $`E_{\mathrm{sep}}`$. We say that $`E`$ admits $`1`$-dimensional local class field theory (abbr, LCFT), if the natural mapping of $`\mathrm{\Omega }(E)`$ into Nr$`(E)`$ (by the rule $`MN(M/E)`$, $`M\mathrm{\Omega }(E)`$) is injective and the following condition holds, for each $`M_1`$, $`M_2\mathrm{\Omega }(E)`$:
(1.1) The norm group (over $`E`$) of the compositum $`M_1M_2`$ equals the intersection $`N(M_1/E)N(M_2/E)`$, and $`N(M_1M_2/E)=N(M_1/E)N(M_2/E)`$.
We say that $`E`$ is a field with $`1`$-dimensional local $`p`$-class field theory (abbr, local $`p`$-CFT), for a given $`p`$, if the fields from the set $`\mathrm{\Omega }_p(E)=\{L\mathrm{\Omega }(E):LE(p)\}`$ are uniquely determined by their norm groups and satisfy condition (1.1). When this is the case and $`pP(E)`$, we have Br$`(E)_p\{0\}`$ (see Proposition 3.3). Observe that $`E`$ admits LCFT if and only if it admits local $`p`$-CFT, for every $`pP(E)`$. This follows from Lemma 2.1 and shows that PQL-fields with LCFT are strictly PQL.
The main purpose of this paper is to shed light on the place of strictly PQL-fields in LCFT by proving the following:
###### Theorem 1.1.
Strictly PQL-fields admit LCFT. Conversely, a field $`E`$ admitting such a theory is strictly PQL, if each $`Dd(E)`$ of prime exponent $`p`$ is similar to a tensor product of cyclic division $`E`$-algebras of index $`p`$.
###### Theorem 1.2.
Let $`E`$ and $`M`$ be fields, such that $`E`$ is strictly PQL, $`P(E)\varphi `$ and $`M\mathrm{\Omega }(E)`$. For each $`pP(E)`$, let $`{}_{p}{}^{}\mathrm{Br}(E)=\{b_p\mathrm{Br}(E):pb_p=0\}`$, $`I_p`$ be a basis and $`d(p)`$ the dimension of $`{}_{p}{}^{}\mathrm{Br}(E)`$ as a vector space over the field $`𝔽_p`$ with $`p`$ elements, $`𝒢(M/E)_p`$ the Sylow $`p`$-subgroup of the Galois group $`𝒢(M/E)`$, and $`𝒢(M/E)_p^{d(p)}`$ the direct product, indexed by $`I_p`$, of isomorphic copies of $`𝒢(M/E)_p`$. Then the direct product $`𝒢(M/E)^{\mathrm{Br}(E)}=_{pP(E)}𝒢(M/E)_p^{d(p)}`$ and the quotient group $`E^{}/N(M/E)`$ are isomorphic.
Before stating our third main result, recall that a field $`F`$ is Pythagorean with $`F(2)=F(\sqrt{1})`$ if and only if it is formally real and $`2`$-quasilocal , I, Lemma 3.5. Note also that, by , Theorem 2, if $`pP(F)`$, then $`F(p)`$ contains as a subfield a $`_p`$-extension $`U_p`$ of $`F`$ (i.e. $`U_p/F`$ is Galois with $`𝒢(U_p/F)`$ isomorphic to the additive group $`_p`$ of $`p`$-adic integers) unless $`p=2`$ and $`F`$ is Pythagorean. We retain notation as in Theorem 1.2.
###### Theorem 1.3.
Let $`E`$ be a strictly PQL-field, such that $`P(E)\varphi `$, and let $`E_{\mathrm{}}E_{\mathrm{sep}}`$ be the compositum of fields $`E_p`$, $`pP(E)`$, where $`E_p/E`$ is a $`_p`$-extension, if $`p>2`$ or $`E`$ is nonreal, and $`E_2=E(2)`$ when $`E`$ is formally real. Then there exists a set $`H_E=\{(,M/E):E^{}𝒢(M/E)^{\mathrm{Br}(E)}`$, $`M\mathrm{\Omega }(E)\}`$ of surjective group homomorphisms with the following properties:
(i) The kernel of $`(,M/E)`$ coincides with $`N(M/E)`$, for each $`M\mathrm{\Omega }(E)`$;
(ii) If $`M\mathrm{\Omega }(E)`$ and $`K`$ is an intermediate field of $`M/E`$, then $`(,K/E)`$ equals the composition $`\pi _{M/K}(,M/E)`$, where $`\pi _{M/K}:𝒢(M/E)^{\mathrm{Br}(E)}𝒢(K/E)^{\mathrm{Br}(E)}`$ is the homomorphism mapping the $`i_p`$-th component of
$`𝒢(M/E)_p^{d(p)}`$ on the $`i_p`$-th component of $`𝒢(K/E)_p^{d(p)}`$ as the natural projection $`𝒢(M/E)_p𝒢(K/E)_p`$, for each pair $`(p,i_p)P(E)\times I_p`$;
(iii) The set $`H_E`$ is uniquely determined by the mappings $`(,\mathrm{\Gamma }/E)`$, where $`\mathrm{\Gamma }`$ runs through the set of finite extensions of $`E`$ in $`E_{\mathrm{}}`$ of primary degrees.
Theorems 1.1 and 1.2 show the strong influence of Br$`(E)`$ on a number of algebraic properties of a PQL-field $`E`$. They are obtained from similar results on finite abelian $`p`$-extensions of $`p`$-quasilocal fields, stated as Theorem 3.1. This approach uses several properties of $`p`$-quasilocal fields without generally valid analogues for PQL-fields (see Propositions 2.2 and 2.3, , I, Corollary 8.5, and , Proposition 6.3). As shown in and , II, Sect. 3, it enables one to describe the isomorphism classes of Brauer groups of major types of PQL-fields, and of the reduced parts of Brauer groups of equicharacteristic Henselian valued absolutely stable fields with totally indivisible value groups. Thus it turns out that usually powerful methods of valuation theory are virtually inapplicable to many PQL and most presently known quasilocal fields (see (2.4) and (2.5) (iii), and compare (2.1) with (2.2) and Remark 2.4). At the same time, the proofs in show at crucial points that the study of the PQL-property can effectively rely on constructive methods based on properties (established in the 1990’s, see e.g. , Proposition 2.6, and ) of relative Brauer groups of extensions obtained as transfers of function fields of Brauer-Severi and other varieties. This makes it possible to apply Theorem 3.1 and other results about $`p`$-quasilocal fields to Brauer groups of arbitrary fields, and so leads to a better understanding of the relations between Galois groups and norm groups of finite Galois extensions of quasilocal fields (see Remark 4.3 and ).
When $`E`$ is a local field, the former assertion of Theorem 1.1 yields the fundamental correspondence of the classical local class field theory. If $`E`$ is merely strictly PQL with $`d(p)=1`$, for all $`pP(E)`$ (i.e. with Br$`(E)`$ embeddable in $`/`$, the quotient group of the additive group of rational numbers by the subgroup of integers), then Theorem 1.2 states that $`E^{}/N(M/E)𝒢(M/E)`$. This holds, for instance, in the following cases:
(1.2) (i) $`E`$ is an algebraic strictly PQL-extension of a global field $`E_0`$; when Br$`(E)_p\{0\}`$, by , Theorem 2.1, $`pP(E)`$ and Br$`(E)_p`$ is isomorphic to the quasicyclic $`p`$-group $`(p^{\mathrm{}})`$ unless $`p=2`$ and $`E`$ is formally real (big families of such $`E`$ can be constructed by applying , Theorem 2.2).
(ii) The triple $`\mathrm{AT}(E)=(𝒢_E,\{𝒢_F:F\mathrm{\Sigma }\},E_{\mathrm{sep}}^{})`$, $`\mathrm{\Sigma }`$ being the set of finite extensions of $`E`$ in $`E_{\mathrm{sep}}`$, is an Artin-Tate class formation (see , Ch. 14).
It is known that (1.2) (ii) holds, if $`E`$ is $`p`$-adically closed or has a Henselian discrete valuation with a quasifinite residue field. Note further that the statement of Theorem 1.2 coincides in case (1.2) (ii) with the local reciprocity law of the Artin-Tate abstract class field theory , Ch. 14, Sect. 5. As to Theorem 1.3, it can be viewed as a partial analogue to the local Hasse symbol (compare with Theorem 3.2, Remark 7.1 and , Theorem 3). The question of whether fields $`E`$ with LCFT are strictly PQL is open. By Theorem 1.1, its answer depends on the solution to one of the leading unsolved problems in Brauer group theory (see Remark 3.4, , Sect. 16, and , Sect. 5). This allows us to prove in Section 3 that SQL-fields are those whose finite extensions admit LCFT. Note also that when $`{}_{p}{}^{}\mathrm{Br}(E)`$ is finite, for each $`pP(E)`$, the answer is positive if and only if $`E^{}/N(M/E)`$ $`𝒢(M/E)^{\mathrm{Br}(E)}`$, $`M\mathrm{\Omega }(E)`$ (apply Theorem 1.2 and , Sect. 15.1, Proposition b).
The paper is organized as follows: Section 2 includes generalities about PQL-fields. The former and the latter assertions of Theorem 1.1 are proved in Sections 4 and 3, respectively. Our main results on local $`p`$-CFT are stated in Section 3. Theorems 1.2, 1.3 and these results are proved in Sections 4, 6 and 7. Section 5 contains an interpretation of a part of local $`p`$-CFT in terms of Galois cohomology. It generalizes , Theorem 1, as well as known relations between local fields and Demushkin groups (cf. , Ch. II, Theorem 4), and the sufficiency part of the main results of .
## 2. Preliminaries on PQL-fields
The present research is based on the possibility to reduce the study of norm groups of finite abelian extensions to the special case of $`p`$-extensions. The reduction is obtained by applying the following lemma (which can be deduced from Galois theory (see e.g., , Ch. VIII) and , II, Lemma 2.2).
###### Lemma 2.1.
Let $`E`$ be a field, $`M\mathrm{\Omega }(E)`$, $`ME`$, $`\mathrm{\Pi }`$ the set of prime divisors of $`[M:E]`$, and $`M_p=ME(p)`$, for each $`p\mathrm{\Pi }`$. Then $`M`$ coincides with the compositum of the fields $`M_p:p\mathrm{\Pi }`$, $`N(M/E)=_{p\mathrm{\Pi }}N(M_p/E)`$ and $`E^{}/N(M/E)`$ is isomorphic to the direct product $`_{p\mathrm{\Pi }}E^{}/N(M_p/E)`$.
The main results of , I, used in this paper can be stated as follows:
###### Proposition 2.2.
Let $`E`$ be a $`p`$-quasilocal field with Br$`(E)_p\{0\}`$, for some $`pP(E)`$. Suppose further that $`R`$ is a finite extension of $`E`$ in $`E(p)`$ and $`Dd(E)`$ is an algebra of $`p`$-primary dimension. Then:
(i) $`R`$ is $`p`$-quasilocal and $`D/E`$ is a cyclic of exponent exp$`(D)=\mathrm{ind}(D)`$;
(ii) Br$`(R)_p`$ is a divisible group unless $`p=2`$, $`R=E`$ and $`E`$ is formally real; when $`E`$ is formally real, $`E(2)=E(\sqrt{1})`$ and Br$`(E)_2`$ is of order $`2`$;
(iii) $`\rho _{E/R}`$ maps Br$`(E)_p`$ surjectively on Br$`(R)_p`$ and Cor<sub>R/E</sub> maps Br$`(R)_p`$ injectively in Br$`(E)_p`$; every $`E`$-automorphism $`\psi `$ of the field $`R`$ is extendable to a ring automorphism on each $`D_Rs(R)`$ of $`p`$-primary index;
(iv) $`R`$ embeds in $`D`$ as an $`E`$-subalgebra if and only if the degree $`[R:E]`$ divides ind$`(D)`$; $`R`$ is a splitting field of $`D`$ if and only if $`\mathrm{ind}(D)[R:E]`$.
Our next result gives an equivalent form of Proposition 2.2 (iv), for PQL-fields, and shows its optimality in the class of algebraic strictly PQL-extensions of the field $``$ of rational numbers. We refer the reader to , II, for a proof of this result, which demonstrates the applicability of the arithmetic method of constructing such extensions, based on , Theorem 2.2.
###### Proposition 2.3.
(i) Let $`E`$ be a PQL-field, $`M/E`$ a finite Galois extension and $`R`$ an intermediate field of $`M/E`$. If $`𝒢(M/E)`$ is nilpotent, then $`R`$ embeds as an $`E`$-subalgebra in each $`Dd(E)`$ of index divisible by $`[R:E]`$.
(ii) For each nonnilpotent finite group $`G`$, there exists a strictly PQL-field $`\mathrm{\Psi }=\mathrm{\Psi }(G)`$, such that $`E/`$ is an algebraic extension, Br$`(\mathrm{\Psi })/`$ and there is a Galois extension $`\mathrm{\Psi }^{}`$ of $`\mathrm{\Psi }`$ with $`𝒢(\mathrm{\Psi }^{}/\mathrm{\Psi })G`$, which does not embed as a $`\mathrm{\Psi }`$-subalgebra in any $`\mathrm{\Delta }d(\mathrm{\Psi })`$ of index ind$`(\mathrm{\Delta })=[\mathrm{\Psi }^{}:\mathrm{\Psi }]`$.
The main results of , I, and , Sect. 3, show that Br$`(E)`$ and $`𝒢_E`$ are the main algebraic structures associated with any quasilocal field $`E`$. Therefore, we would like to point out that $`𝒢_E`$ is prosolvable and Br$`(E)`$ is embeddable in $`/`$ in the following two cases:
(2.1) (i) $`E`$ is an algebraic extension of a quasilocal field $`K`$ with a Henselian valuation $`v`$, such that $`v(K)`$ is totally indivisible (i.e. $`v(K)pv(K)`$, for every $`p`$); then Br$`(E)`$ is divisible with Br$`(E)_p^{}=\{0\}`$, $`p^{}P(E)`$ (concerning $`𝒢_E`$ and Br$`(E)`$, see , I, Lemma 1.2 and Proposition 3.1, and , Corollary 5.3 as well as , I, (1.3), respectively).
(ii) $`E`$ is formally real and quasilocal; then Br$`(E)`$ is of order $`2`$ and $`𝒢_E`$ has an abelian open subgroup of index $`2`$, which is either procyclic or $`2`$-generated as a topological group (see , Sect. 3, and , Proposition 3.4).
Condition (2.1) (i) holds, for example, when $`(K,v)`$ is Henselian discrete valued quasilocal or $`E`$ is an algebraic quasilocal nonreal extension of a global field, such that Br$`(E)\{0\}`$ (cf. , Sect. 3). The existence of a PQL-field $`F`$ of essentially nonarithmetic nature, i.e. with Br$`(F)`$ not embeddable in $`/`$, has been established in , Sect. 3, by observing that every abelian torsion group $`A`$ embeds in Br$`(F(A))`$, for a suitably chosen strictly PQL-field $`F(A)`$. This result and Proposition 2.2 (ii) are complemented by the following statement which describes, in conjunction with , Theorem 23.1, the isomorphism classes of Brauer groups of nonreal PQL-fields (see , Theorem 1.2 (i)-(ii), and for the formally real case, , Proposition 6.4). This statement also shows that the absolute Galois groups of nonreal SQL-fields need not be prosolvable and may have a very complex structure:
(2.2) An abelian torsion group $`T`$ is isomorphic to Br$`(E)`$, for some nonreal PQL-field $`E=E(T)`$ if and only if $`T`$ is divisible. If $`T`$ is divisible, $`S`$ is a set of finite groups, $`E_0`$ is an arbitrary field, and $`T_0`$ is a subgroup of Br$`(E_0)`$ embeddable in $`T`$, then $`E`$ can be chosen so as to satisfy the following:
(i) $`E`$ is quasilocal, $`P(E)=`$ and $`\rho _{E/L}`$ surjective, for any finite extension $`L/E`$ (apply also the Albert-Hochschild theorem (cf. , Ch. II, 2.2)).
(ii) $`E/E_0`$ is an extension, such that $`E_0`$ is algebraically closed in $`E`$, $`T_0\mathrm{Br}(E/E_0)=\{0\}`$ and every $`GS`$ is realizable as a Galois group over $`E`$.
When the set $`{}_{p}{}^{}T=\{tT:pt=0\}`$ is infinite, for each $`p`$, (2.2) implies the following result (combined with , Remark 6.6, it points up the need for the restrictions on the ground fields considered in and ):
(2.3) In order that $`A,Bd(E)`$ have a common set of splitting fields among the intermediate fields of $`E_{\mathrm{sep}}/E`$ it is necessary and sufficient that ind$`(A)=\mathrm{ind}(B)`$ , I, Corollary 8.5. For each $`n`$, $`d(E)`$ contains infinitely many nonisomorphic $`E`$-algebras of index $`n`$.
Statements (2.1) and (2.2) can be supplemented by the following partial conversion of the Koenigsmann-Neukirch theorem , Theorem B:
(2.4) If $`(E,v)`$ is a Henselian quasilocal field, such that all finite groups are isomorphic to subquotients of $`𝒢_E`$ (by open normal subgroups), then $`v(E)`$ is divisible and every $`Dd(E)`$ is inertial relative to $`v`$ , Proposition 6.3.
###### Remark 2.4.
Let $`T`$ be a divisible abelian torsion group and Supp$`(T)`$ the set of those $`p`$ for which the $`p`$-component $`T_p`$ of $`T`$ is nontrivial. Fix a field $`E`$ with Br$`(E)T`$ as in (2.2) (i) and denote by $`E_{\mathrm{sol}}`$ the maximal Galois extension of $`E`$ in $`E_{\mathrm{sep}}`$ with a prosolvable Galois group. Then $`E_{\mathrm{sol}}/E`$ possesses an intermediate field $`E^{}`$ that is strictly PQL with Br$`(E^{})T`$ (take as $`E^{}`$, e.g., the fixed field of some Hall pro-$`\mathrm{Supp}(T)`$-subgroup of $`𝒢(E_{\mathrm{sol}}/E)`$). Note also that if Supp$`(T)=`$, then $`E`$ is SQL.
Remark 2.4 and statements (1.2) (i) and (2.1)$`÷`$(2.4) draw interest in the following open questions:
(2.5) (i) Describe the isomorphism classes of Brauer groups of SQL-fields.
(ii) Let $`E`$ be an SQL-field with $`𝒢_E`$ prosolvable. Find whether the nonzero groups Br$`(E)_p`$ except, possibly, one are isomorphic to the quasicyclic $`p`$-group $`(p^{\mathrm{}})`$. Prove whether the Sylow pro-$`p`$-subgroups of $`𝒢_E`$ are of rank $`r_p2`$, for all $`p`$, with at most $`2`$ exceptions.
(iii) Let $`(F,v)`$ be a Henselian quasilocal field with Br$`(F)_\pi `$ and Br$`(E)_\pi ^{}`$ not embeddable in $`/`$, for some $`\pi ,\pi ^{}`$. Is $`v(F)`$ divisible?
The questions in (2.5) (ii) are related by Galois cohomology (cf. , page 265, , (11.5), , Ch. I, 4.2, and , Lemma 7) and the results of Section 5. Also, by , (1.4), the assumptions of (2.1) (i) require that $`r_p2`$ whenever $`p\mathrm{char}(\widehat{E})`$. In view of , Theorem B, one can prove that a negative answer to the latter question in (2.5) (ii) would imply the existence of a field $`F`$ with Br$`(F)\{0\}`$ whose nontrivial valuations have Henselizations equal to $`F_{\mathrm{sep}}`$. Moreover, if $`3r_pr_p^{}`$, for some $`p,p^{}`$, $`p<p^{}`$, this would solve affirmatively Problem 11.5.9 (b) in .
###### Remark 2.5.
It is known (see ) that if $`E`$ is a field, then the triple AT$`(E)`$ defined in (1.2) (ii) is an Artin-Tate class formation if and only if $`E`$ is SQL, Br$`(E)`$ is embeddable in $`/`$, and $`\rho _{E/R}`$ is surjective, for every finite extension $`R`$ of $`E`$ in $`E_{\mathrm{sep}}`$. Observing that (2.2) provides access to the richest presently known sources of such fields, we add to the examples given after the statement of (1.2) that AT$`(E)`$ is Artin-Tate when $`E`$ is an algebraic SQL-extension of a global field (see and , Sect. 3).
## 3. Statements of the main results of local $`p`$-class field theory
Our main results on local $`p`$-CFT are stated as the following two theorems:
###### Theorem 3.1.
Let $`E`$ be a $`p`$-quasilocal field with Br$`(E)_p\{0\}`$, for some $`pP(E)`$, and let $`\mathrm{\Omega }_p(E)`$ and $`d(p)`$ be as in the Introduction. Then $`E`$ admits local $`p`$-CFT and, for each $`M\mathrm{\Omega }_p(E)`$, $`E^{}/N(M/E)`$ is isomorphic to the group $`𝒢(M/E)^{d(p)}`$ defined in Theorem 1.2.
Theorem 3.1 plays a major role in the proof of the embeddability of Br$`(E)`$ into $`/`$ in case (2.1) (i), which in turn enables one to characterize the quasilocal property in the class of Henselian valued fields with totally indivisible value groups (see , Sect. 6, and , II). Note also that Theorem 3.1 and Lemma 2.1 imply Theorem 1.2 and the former assertion of Theorem 1.1. In addition, Theorem 3.1 contains exact analogues to the fundamental correspondence and the local reciprocity law. Similarly, our next result can be viewed as such an analogue to Hasse’s symbol for local fields (see , Ch. 2, as well as Remark 7.1 and Corollary 7.3 below).
###### Theorem 3.2.
Under the hypotheses of Theorem 3.1, let $`E_{\mathrm{}}`$ be a $`_p`$-extension of $`E`$ in $`E(p)`$, and for any $`n`$, let $`\mathrm{\Gamma }_n`$ be the intermediate field of $`E_{\mathrm{}}/E`$ of degree $`[\mathrm{\Gamma }_n:E]=p^n`$. Then there exist sets
$`H_p(E^{})=\{(,M^{}/E^{}):E^{}𝒢(M^{}/E^{})^{d(p)},M^{}\mathrm{\Omega }_p(E^{})\}`$, $`E^{}\mathrm{\Omega }_p(E)`$,
of surjective group homomorphisms satisfying the following:
(i) The kernel of $`(,M^{}/E^{})`$ is equal to $`N(M^{}/E^{})`$, for each $`E^{}\mathrm{\Omega }_p(E)`$, $`M^{}\mathrm{\Omega }_p(E^{})`$;
(ii) If $`M\mathrm{\Omega }_p(E)`$ and $`K`$ is an intermediate field of $`M/E`$, then $`(,K/E)=\pi _{M/K}(,M/E)`$, where $`\pi _{M/K}:𝒢(M/E)^{d(p)}𝒢(K/E)^{d(p)}`$ is the homomorphism acting componentwise as the natural mapping of $`𝒢(M/E)`$ on $`𝒢(K/E)`$;
(iii) In the setting of (ii), $`(\lambda ,M/K)=(N_E^K(\lambda ),M/E)`$, for each $`\lambda K^{}`$;
(iv) The maps $`(,\mathrm{\Gamma }_n/E)`$, $`n`$, determine the sets $`H_p(E^{}),E^{}\mathrm{\Omega }_p(E)`$.
The rest of this Section concerns the open question of whether fields with LCFT are strictly PQL. Lemma 2.1 and first result in this direction imply the latter assertion of Theorem 1.1.
###### Proposition 3.3.
Let $`E`$ be a field admitting local $`p`$-CFT, for some $`pP(E)`$, and let $`L`$ be a degree $`p`$ extension of $`E`$ in $`E(p)`$. Then Br$`(L/E)\{0\}`$ and Br$`(L/E)`$ does not depend on the choice of $`L`$. Furthermore, if each $`Dd(E)`$ of exponent $`p`$ is similar to tensor products of cyclic division $`E`$-algebras of index $`p`$, then $`E`$ is $`p`$-quasilocal.
###### Proof.
Our assumptions ensure that $`L/E`$ is cyclic, whence Br$`(L/E)E^{}/N(L/E)`$ (cf. , Ch. I, Sect. 6, and , Sect. 15.1, Proposition b). As $`E`$ admits local $`p`$-CFT and $`LE`$, this yields $`N(L/E)N(E/E)=E^{}`$ and Br$`(L/E)\{0\}`$. In order to complete our proof, it suffices to show that $`L`$ embeds as an $`E`$-subalgebra in each cyclic division $`E`$-algebra of index $`p`$. If $`𝒢(E(p)/E)`$ is procyclic, this is evident, so we assume that $`\mathrm{\Omega }_p(E)`$ contains a field $`FL`$, such that $`[F:E]=p`$. It follows from Galois theory that $`LF\mathrm{\Omega }(E)`$, $`[LF:E]=p^2`$ and $`𝒢(LF/E)`$ is noncyclic. Therefore, $`LF/E`$ possesses $`p+1`$ intermediate fields of degree $`p`$ over $`E`$. Let $`F^{}`$ be such a field different from $`L`$ and $`F`$. Then $`N(F/E)N(F^{}/E)=E^{}`$ and $`LF^{}=LF`$. Moreover, the norm maps $`N_E^F`$ and $`N_E^F^{}`$ are induced by $`N_L^{LF}`$, so it turns out that $`E^{}N(LF/L)`$. Hence, by , Sect. 15.1, Proposition b, $`L`$ embeds over $`E`$ in each $`\mathrm{\Delta }d(E)`$ split by $`F`$, which proves Proposition 3.3. ∎
###### Remark 3.4.
Let $`E`$ be a field and $`pP(E)`$.
(i) If $`𝒢(E(p)/E)_p`$ and $`L`$ is the unique degree $`p`$ extension of $`E`$ in $`E(p)`$, then $`E`$ admits local $`p`$-CFT if and only if Br$`(L/E)\{0\}`$ (apply Proposition 3.3 and , Sect. 15.1, Corollary b).
(ii) When Br$`(E)_p\{0\}`$, the concluding condition of Proposition 3.3 is satisfied in the following cases: (a) $`E`$ is an algebraic extension of a global or local field $`E_0`$; (b) $`E`$ contains a primitive $`p`$-th root of unity or char$`(E)=p`$; (c) $`p=3,5`$. In cases (b) and (c), this follows from the Merkur’ev-Suslin theorem , (16.1), , Ch. VII, Theorem 28, , Sect. 4, Corollary, and . In case (a), by class field theory, every $`A_0s(E_0)`$ is cyclic (cf. , Ch. 10, Theorem 5), which implies the same, for all $`As(E)`$.
(iii) It is not known whether $`E`$ is strictly PQL, if it admits LCFT and has a Henselian discrete valuation (see , Sect. 2).
###### Corollary 3.5.
A finite purely inseparable extension $`K`$ of a field $`E`$ of characteristic $`p>0`$ admits local $`p`$-CFT if and only if so does $`E`$.
###### Proof.
By , I, Proposition 4.4, $`E`$ is $`p`$-quasilocal if and only if $`K`$ is $`p`$-quasilocal, and by the Albert-Hochschild theorem, $`\rho _{E/K}`$ is surjective. Since the exponent of Br$`(K/E)`$ divides $`[K:E]`$ (see , Sects. 13.4 and 14.4), and by Witt’s theorem (cf. , page 110), Br$`(E)_p`$ and Br$`(K)_p`$ are divisible, this implies that Br$`(E)_p\{0\}\mathrm{Br}(K)_p\{0\}`$, so Corollary 3.5 follows from Theorem 3.1 and Remark 3.4 (ii). ∎
The concluding result of this Section clarifies the role of SQL-fields in LCFT. In view of (2.2), Remark 2.5 and the main results of , it also determines the place in the study of strictly PQL-fields of the Neukirch-Perlis variant of the theory , built upon (1.2) (ii).
###### Proposition 3.6.
A field $`E`$ is SQL if and only if its finite extensions admit LCFT. When occurs, Br$`(E_p)\{0\}`$, provided that $`E_p`$ is the fixed field of a Sylow pro-$`p`$-subgroup of $`𝒢_E`$, where $`p`$ is chosen so that $`𝒢_E`$ is of cohomological $`p`$-dimension cd$`{}_{p}{}^{}(𝒢_E)0`$.
###### Proof.
The left-to-right implication is contained in Theorem 1.1, so we prove the converse one and the nontriviality of Br$`(E_p)_p`$, for any $`p`$ satisfying cd$`{}_{p}{}^{}(𝒢_E)0`$. Suppose that every finite extension $`L`$ of $`E`$ admits LCFT, and $`B_p`$ is the extension of $`E`$ generated by the $`p`$-th roots of unity in $`E_{\mathrm{sep}}`$. It is known (cf. , Ch. VIII, Sect. 3) that $`B_pE_p`$, i.e. $`E_p`$ contains a primitive $`p`$-th root of unity unless $`p=\mathrm{char}(E)`$. Also, Proposition 3.3 and Remark 3.4 (ii) indicate that $`L`$ is $`p`$-quasilocal, provided that $`B_pL`$. As $`p[L^{}:E]`$, for any finite extension $`L^{}`$ of $`E`$ in $`E_p`$, this enables one to obtain from general properties of scalar extension maps and Schur indices (cf. , Sect. 13.4, and , I, (1.3)) that $`E_p`$ is $`p`$-quasilocal. Moreover, it follows from , I, Lemma 8.3, and the choice of $`p`$ that $`E`$ is quasilocal. Since, by Proposition 3.3, PQL-fields with LCFT are strictly PQL, the obtained result shows that $`E`$ is SQL. Note further that the inequality cd$`{}_{p}{}^{}(𝒢_E)0`$ ensures the existence of finite Galois extensions of $`E`$ in $`E_{\mathrm{sep}}`$ of degrees divisible by $`p`$. Therefore, by Sylow’s theorems and Galois theory, there is a finite extension $`L_0`$ of $`E`$ in $`E_p`$, such that $`pP(L_0)`$. As in the proof of the statement that $`E_p`$ is $`p`$-quasilocal, these observations imply $`pP(E_p)`$ and Br$`(E_p)\{0\}`$. ∎
Let us note that the abstract approach to LCFT dates back to the early 1950’s (cf. , Sect. 14). It fits the character of the related class formation theory, see , and accounts for the fact that the Neukirch-Perlis variant of LCFT goes substantially beyond the limits of (2.1). The fact itself is established by summing up Proposition 3.6, Remark 2.5, statements (2.1), (2.2) and (2.4), and the results on (2.5) (ii) mentioned in Section 2.
## 4. Proof of Theorem 3.1
Let $`E`$ be a $`p`$-quasilocal field, for a given $`pP(E)`$. By Proposition 2.2 (iv), then Br$`(M/E)=\{b\mathrm{Br}(\mathrm{E}):[M:E]b=0\}`$, for every cyclic $`p`$-extension $`M/E`$. Hence, by the structure of divisible abelian torsion groups, and by the fact that Br$`(M/E)E^{}/N(M/E)`$ (cf. , Theorem 23.1, and , Sect. 15.1, Proposition b), $`E^{}/N(M/E)𝒢(M/E)^{d(p)}`$. To prove the obtained isomorphism, for any $`M\mathrm{\Omega }_p(E)`$, we need the following lemmas.
###### Lemma 4.1.
Let $`E`$ be a $`p`$-quasilocal field, for some $`pP(E)`$, and let $`E_1,\mathrm{}E_t`$ be cyclic extensions of $`E`$ in $`E(p)`$, for a given integer $`t2`$. Assume that the compositum $`E^{}`$ of the fields $`E_j:j=1,\mathrm{},t`$, satisfies the equality $`[E^{}:E]=_{j=1}^t[E_j:E]`$. Then $`N(E^{}/E)=_{j=1}^tN(E_j/E)`$.
###### Proof.
The inclusion $`N(E^{}/E)_{j=1}^tN(E_j/E)`$ follows from the transitivity of norm maps in towers of finite separable extensions (cf. , Ch. VIII, Sect. 5). Conversely, let $`c_{j=1}^tN(E_j/E)`$ and $`\beta E_1^{}`$ be of norm $`N_E^{E_1}(\beta )=c`$. The equality $`[E^{}:E]=_{j=1}^t[E_j:E]`$, Galois theory and , Lemma 4.2, imply $`[E^{}:E_1]=_{i=2}^t[(E_1E_i):E_1]`$ and $`\beta _{i=2}^tN(E_1E_i/E_1)`$. This proves Lemma 4.1 in the case where $`t=2`$. Since $`(E_1E_i)/E_1`$ is cyclic, for each $`i2`$, and by Proposition 2.2 (i), $`E_1`$ is $`p`$-quasilocal, the obtained result makes it easy to complete our proof by induction on $`t`$. ∎
###### Lemma 4.2.
Assume that $`E`$, $`F`$ and $`L`$ are fields, such that $`E`$ is $`p`$-quasilocal, $`L\mathrm{\Omega }_p(E)`$, $`EFL`$ and $`F/E`$ is cyclic. Then $`\psi (\alpha )\alpha ^1N(L/F)`$, for each $`\alpha F^{}`$ and $`\psi 𝒢(F/E)`$.
###### Proof.
As $`F`$ is $`p`$-quasilocal and $`L\mathrm{\Omega }_p(F)`$, whence $`𝒢(L/F)`$ decomposes into a direct product of cyclic groups, Galois theory and Lemma 4.1 allow one to consider only the case in which $`L/F`$ is cyclic. Let $`\psi ^{}`$ be an automorphism of $`L`$ extending $`\psi `$. Fix a generator $`\sigma `$ of $`𝒢(L/F)`$, denote by $`A_\alpha `$ the cyclic $`F`$-algebra $`(L/F,\sigma ,\alpha )`$, for an arbitrary $`\alpha F^{}`$, put $`m=[L:F]`$, and take an invertible element $`\eta A_\alpha `$ so that $`\eta ^m=\alpha `$ and $`\eta \lambda \eta ^1=\sigma (\lambda )`$, for every $`\lambda L`$. Then Proposition 2.2 (iii) and the Skolem-Noether theorem (cf. , Sect. 12.6) imply that $`A_\alpha `$ has a ring automorphism $`\stackrel{~}{\psi }`$, such that $`\stackrel{~}{\psi }(\lambda )`$ $`=\psi ^{}(\lambda )`$, for any $`\lambda L`$. Since $`\sigma \psi ^{}=\psi ^{}\sigma `$, this ensures that the element $`\eta ^1\stackrel{~}{\psi }(\eta ):=\mu `$ lies in the centralizer of $`L`$ in $`A_\alpha `$. Thereby, we have $`\mu L^{}`$ and $`\stackrel{~}{\psi }(\eta )^m=\psi (\alpha )=\eta ^mN_F^L(\mu )=\alpha N_F^L(\mu )`$, which proves Lemma 4.2. ∎
Assuming that $`E`$ is a $`p`$-quasilocal field, $`M_u\mathrm{\Omega }_p(E)`$, $`[M_u:E]=p^{\mu _u}:u=1,2`$, and putting $`M^{}=M_1M_2`$, we prove the following assertions:
(4.1) (i) If $`N(M_1/E)=N(M_2/E)`$ and $`M_1M_2`$, then $`M_1=M_2`$;
(ii) If $`M_1M_2=E`$, then $`N(M_1/E)N(M_2/E)=N(M^{}/E)`$,
$`N(M_1/E)N(M_2/E)=E^{}`$, each $`\sigma 𝒢(M_2/E)`$ has a unique prolongation $`\sigma ^{}𝒢(M^{}/M_1)`$, and the mapping of $`𝒢(M_2/E)`$ on $`𝒢(M^{}/M_1)`$ by the rule $`\sigma \sigma ^{}`$ is an isomorphism;
(iii) Under the hypotheses of (ii), if $`M_2/E`$ is cyclic and $`𝒢(M_2/E)=\sigma `$, then $`𝒢(M^{}/M_1)=\sigma ^{}`$ and Cor$`_{M_1/E}`$ maps Br$`(M^{}/M_1)`$ into Br$`(M_2/E)`$ by the formula $`[(M^{}/M_1,\sigma ^{},\theta )][(M_2/E,\sigma ,N_E^{M_1}(\theta ))]`$, $`\theta M_1^{}`$; when Br$`(E)_p`$ is divisible, Cor$`_{M_1/E}`$ induces isomorphisms Br$`(M_1)_p\mathrm{Br}(E)_p`$ and Br$`(M^{}/M_1)\mathrm{Br}(M_2/E)`$.
It is sufficient to consider the special case where $`M_uE:u=1,2`$, and $`M_1M_2`$. In view of Proposition 2.2 (ii), then Br$`(E)_p`$ is divisible. At the same time, Galois theory and the assumptions on $`M_1/E`$ imply the existence of a cyclic degree $`p`$ extension $`M_0`$ of $`E`$ in $`M_1`$. Suppose that $`M_1M_2`$ and $`dN(M_u/E):u=1,2`$, fix a generator $`\psi `$ of $`𝒢(M_0/E)`$ and elements $`\eta _jM_j^{}`$, $`j=1,2`$, so that $`N_E^{M_1}(\eta _1)=N_E^{M_2}(\eta _2)=d`$. It follows from Hilbert’s Theorem 90 that $`N_{M_0}^{M_1}(\eta _1)=N_{M_0}^{M_2}(\eta _2)\psi (\beta )\beta ^1`$, for some $`\beta M_0^{}`$. Hence, by Lemma 4.2, $`N(M_1/E)=N(M_2/E)`$ if and only if $`N(M_1/M_0)=N(M_2/M_0)`$. Observing also that $`M_0`$ is $`p`$-quasilocal and $`[M_1:M_0]=p^{\mu _11}`$, and proceeding by induction on $`\mu _1`$, one proves (4.1) (i).
Assume now that $`M_1M_2=E`$. Then the Galois-theoretic parts of (4.1) (ii) and (iii) are contained in , Ch. VIII, Theorem 4. The rest of the former assertion of (4.1) (iii) is implied by Proposition 2.2 (iii), the basic restriction-corestriction (abbr, RC) formula (cf. , Theorem 2.5) and the lemmas in , Sect. 4. The equality $`N(M^{}/E)=N(M_1/E)N(M_2/E)`$ follows from the presentability of $`M_1`$ and $`M_2`$ as compositums of cyclic extensions of $`E`$ satisfying the conditions of Lemma 4.1. It remains for us to show that $`N(M_1/E)N(M_2/E)=E^{}`$ and to prove the latter part of (4.1) (iii). Suppose first that $`M_2/E`$ is cyclic and the automorphisms $`\sigma ,\sigma ^{}`$ are determined as in (4.1) (iii), fix an element $`cE^{}`$ out of $`N(M_2/E)`$, and put $`A_c=(M_2/E,\sigma ,c)`$. It is clear from , Sect. 15.1, Proposition b, and the choice of $`c`$ that ind$`(A_c)p^{\mu _2}`$ and ind$`(A_c)>1`$. As Br$`(E)_p`$ is divisible, there exists $`\mathrm{\Delta }_cd(E)`$, such that $`p^{\mu _1}[\mathrm{\Delta }_c]=[A_c]`$. In addition, it follows from Proposition 2.2 (i) that ind$`(\mathrm{\Delta }_c)=p^{\mu _1}.\mathrm{ind}(A_c)`$. Observing also that $`[M^{}:E]=p^\mu `$ and ind$`(\mathrm{\Delta }_c)p^\mu `$, where $`\mu =\mu _1+\mu _2`$, one obtains from Proposition 2.2 (iv) that $`[\mathrm{\Delta }_c]\mathrm{Br}(M^{}/E)`$. Since $`\mathrm{\Delta }_c_EM^{}(\mathrm{\Delta }_c_EM_1)_{M_1}M^{}`$ as $`M^{}`$-algebras (cf. , Sect. 9.4, Corollary a), this means that $`[\mathrm{\Delta }_c_EM_1]\mathrm{Br}(M^{}/M_1)`$, or equivalently, that $`[\mathrm{\Delta }_c_EM_1]=[(M^{}/M_1,\sigma ^{},\alpha )]`$, for some $`\alpha M_1^{}`$. Therefore, by Proposition 2.2 (iii), the RC-formula and the former statement of (4.1) (iii), $`[A_c]=[(M_2/E,\sigma ,N_E^{M_1}(\alpha ))]`$. As $`[A_c:E]=[(M_2/E,\sigma ,N_E^{M_1}(\alpha )):E]`$ $`=[M_2:E]^2`$, this proves that $`A_c(M_2/E,\sigma ,N_E^{M_1}(\alpha ))`$ over $`E`$. Hence, by , Sect. 15.1, Proposition b, $`c.N_E^{M_1}(\alpha )^1N(M_2/E)`$, which yields $`N(M_1/E)N(M_2/E)=E^{}`$ in case $`M_2/E`$ is cyclic. Our argument, combined with Proposition 2.2 (ii)-(iii) and the RC-formula, also proves the latter part of (4.1) (iii). Henceforth, we assume that $`M_2/E`$ is noncyclic, i.e. $`𝒢(M_2/E)`$ is an abelian $`p`$-group of rank $`r2`$. Then it follows from Galois theory and the structure of finite abelian groups that there exist cyclic extensions $`F_1`$ and $`F_2`$ of $`E`$ in $`M_2`$ such that $`F_1F_2=E`$ and the $`p`$-groups $`𝒢(M_2/F_1)`$ and $`𝒢(M_2/F_2)`$ are of rank $`r1`$. As $`M_1M_2=E`$, one also sees that $`(M_1F_u)M_2=F_u:u=1,2`$. Taking now into account that $`F_1`$ and $`F_2`$ are $`p`$-quasilocal fields, and arguing by induction on $`r`$, one concludes that it suffices to deduce the equality $`N(M_1/E)N(M_2/E)=E^{}`$ under the extra hypothesis that $`N(M_1F_u/F_u)N(M_2/F_u)=F_u^{}:u=1,2`$. Then, by norm transitivity in towers of finite extensions, $`N(F_1/E)N(F_2/E)N(M_1/E)N(M_2/E)`$, and since $`N(F_1/E)N(F_2/E)=E^{}`$, this completes the proof of (4.1).
###### Remark 4.3.
Let $`E`$ be a field and $`M_1`$, $`M_2`$ be finite extensions of $`E`$ in $`E_{\mathrm{sep}}`$, such that $`M_2/E`$ is cyclic and $`M_1M_2=E`$. It is known that then the former statement of (4.1) (iii) remains valid. The assertions on $`𝒢(M_2/E)`$ and $`𝒢(M^{}/M_1)`$ follow from , Ch. VIII, Theorem 4, and the formula for the action of Cor$`_{M_1/E}`$ on Br$`(M^{}/M_1)`$ can be proved by a group-cohomological technique (cf. , Proposition 4.3.7). It is therefore worth noting that (4.1) (iii) plays a role in the proof not only of Theorems 3.1 and 3.2 but also of Proposition 2.2 (i) (see , I, Sect. 7). This, combined with , Theorem 1.2 (i)-(ii), enables one to find alternative field-theoretic proofs of the formula in (4.1) (iii) and of other results on Cor<sub>F/E</sub>, for any finite extension $`F`$ of $`E`$ in $`E_{\mathrm{sep}}`$ (by reduction to the setting of (2.2) (i), see ).
We are now in a position to prove Theorem 3.1 (and thereby, Theorem 1.2 and the former part of Theorem 1.1 as well). Assuming as above that $`M_u\mathrm{\Omega }_p(E):u=1,2`$, and $`M_1M_2=M^{}`$, put $`L^{}=M_1M_2`$. We first show that $`N(L^{}/E)=N(M_1/E)N(M_2/E)`$ and $`N(M^{}/E)=N(M_1/E)N(M_2/E)`$. In view of (4.1) (ii), it suffices to consider the case where $`L^{}E`$. As $`L^{}`$ is $`p`$-quasilocal, (4.1) (ii) yields $`L^{}=N(M_1/L^{})N(M_2/L^{})`$, so the equality $`N(L^{}/E)=N(M_1/E)N(M_2/E)`$ is obtained from norm transitivity in the towers $`EL^{}M_u:u=1,2`$. Suppose further that $`[L^{}:E]=p^m`$ and fix a degree $`p`$ extension $`L`$ of $`E`$ in $`L^{}`$. As in the proof of (4.1) (i), it is seen that an element $`\lambda L^{}`$ lies in $`N(M_i/L)`$, for some $`i\{1,2\}`$, if and only if $`N_E^L(\lambda )N(M_i/E)`$. Since $`[L^{}:L]=p^{m1}`$ and $`L`$ is $`p`$-quasilocal, this makes it easy to prove inductively that $`N(M^{}/E)=N(M_1/E)N(M_2/E)`$. It follows from this result and (4.1) (i) that the natural mapping $`\mathrm{\Omega }_p(E)\mathrm{Nr}(E)`$ is injective. Thus the statement that $`E`$ admits local $`p`$-CFT is proved, which allows us to deduce Theorem 1.1 from Lemma 2.1 and Proposition 3.3.
To finish the proof of Theorem 3.1 (and Theorem 1.2) we show that $`E^{}/N(M/E)𝒢(M/E)^{d(p)}`$, provided that $`M\mathrm{\Omega }_p(E)`$ and $`𝒢(M/E)`$ is noncyclic. Then $`𝒢(M/E)`$ is an abelian $`p`$-group of rank $`r(M)2`$, so it follows from Galois theory that $`M/E`$ has intermediate fields $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$, such that $`\mathrm{\Phi }_1\mathrm{\Phi }_2=M`$, $`\mathrm{\Phi }_1\mathrm{\Phi }_2=E`$ and $`\mathrm{\Phi }_2/E`$ is cyclic. This means that $`𝒢(M/E)`$ is isomorphic to the direct product $`𝒢(\mathrm{\Phi }_1/E)\times 𝒢(\mathrm{\Phi }_2/E)`$, and $`𝒢(\mathrm{\Phi }_1/E)`$ is of rank $`r(M)1`$ as a $`p`$-group. Since $`N(\mathrm{\Phi }_1/E)N(\mathrm{\Phi }_2/E)=E^{}`$ and $`N(\mathrm{\Phi }_1/E)N(\mathrm{\Phi }_2/E)=N(M/E)`$, the natural diagonal embedding of $`E^{}`$ into $`E^{}\times E^{}`$ induces a group isomorphism $`E^{}/N(M/E)E^{}/N(\mathrm{\Phi }_1/E)\times E^{}/N(\mathrm{\Phi }_2/E)`$. Now our proof is easily completed proceeding by induction on $`r(M)`$.
###### Remark 4.4.
Given a PQL-field $`E`$, denote by $`\mathrm{c}(M)`$ the intersection of fields $`N\mathrm{\Omega }(E)`$ with $`N(N/E)=N(M/E)`$, for each $`M\mathrm{\Omega }(E)`$. It follows from Theorem 3.1, Lemma 2.1 and , I, Lemma 4.2 (ii), that the natural map $`\nu `$ of $`\mathrm{\Omega }(E)`$ on the set $`N_{\mathrm{ab}}(E)=\{N(M/E):M\mathrm{\Omega }(E)\}`$ satisfies (1.1). Therefore, Br$`(E)_p\{0\}`$ and $`p[M:c(M)]`$ whenever $`M\mathrm{\Omega }(E)`$ and $`p`$ divides $`[\mathrm{c}(M):E]`$. Note also that $`N(\mathrm{c}(M)/E)=N(M/E)`$ and $`N(M_0/E)N(M/E)`$ in case $`M_0\mathrm{\Omega }(E)`$, $`M_0\mathrm{c}(M)`$ and $`M_0\mathrm{c}(M)`$. Since the set Cl$`(E)=\{\mathrm{c}(\mathrm{\Lambda }):\mathrm{\Lambda }\mathrm{\Omega }(E)\}`$ is closed under taking subextensions of $`E`$ and finite compositums, whence $`\nu `$ induces a bijection of Cl$`(E)`$ on $`N_{\mathrm{ab}}(E)`$, these observations allow us to view $`\mathrm{c}(M)`$ as a class field of $`N(M/E)`$.
## 5. Galois cohomological interpretation of Theorem 3.1
In this Section we consider some Galois cohomological aspects of the problem of characterizing fields with LCFT. Let $`P`$ an infinite pro-$`p`$-group, cd$`(P)`$ the cohomological dimension of $`P`$ and $`𝔽_p`$ a field with $`p`$ elements, for some $`p`$. We say that $`P`$ is a $`p`$-group of Demushkin type, if the (continuous) cohomology group homomorphism $`\phi _h:H^1(P,𝔽_p)H^2(P,𝔽_p)`$ mapping each $`gH^1(P,𝔽_p)`$ into the cup-product $`hg`$ is surjective, for every $`hH^1(P,𝔽_p)\{0\}`$. We call a degree of $`P`$ the dimension of $`H^2(P,𝔽_p)`$ as an $`𝔽_p`$-vector space. The defined groups and local $`p`$-CFT are related as follows:
###### Proposition 5.1.
Let $`E`$ be a nonreal field containing a primitive $`p`$-th root of unity, for some $`pP(E)`$. Then the following conditions are equivalent:
(i) $`E`$ admits local $`p`$-CFT;
(ii) $`𝒢(E(p)/E)`$ is a $`p`$-group of Demushkin type of degree $`d1`$;
(iii) $`E`$ is $`p`$-quasilocal with Br$`(E)_p\{0\}`$;
(iv) $`\mathrm{cd}(𝒢(E(p)/E))=2`$ and $`{}_{p}{}^{}\mathrm{Br}(E^{})`$ is a trivial module over the integral group ring $`[𝒢(E^{}/E)]`$, for every degree $`p`$ extension $`E^{}`$ of $`E`$ in $`E(p)`$.
###### Proof.
The equivalence (i)$``$(iii) is implied by Proposition 3.3, Remark 3.4 (ii) and Theorem 3.1. Note also that Br$`(E)_p=\{0\}`$ if and only if $`𝒢(E(p)/E)`$ is a free pro-$`p`$-group, i.e. a $`p`$-group of Demushkin type of degree zero (cf. , Theorem 3.1, , page 725, and , Ch. I, 4.1 and 4.2). In particular, this holds when $`𝒢(E(p)/E)`$ is of rank $`1`$ as a pro-$`p`$-group, since then $`𝒢(E(p)/E)_p`$ (see, e.g., , I, Remark 3.4 (ii)). These observations, , I, Lemma 3.8, and the end of Proposition 2.2 (ii) prove that (ii)$``$(iii). It remains to be seen that (iii)$``$(iv). As $`E`$ is nonreal and $`pP(E)`$, one obtains from Galois theory and , Theorem 2, that $`𝒢(E(p)/E)`$ possesses a closed normal subgroup $`H`$, such that $`𝒢(E(p)/E)/H_p`$. Since Br$`(E)_p\{0\}`$, this means that $`H\{1\}`$, so it follows from , I, Proposition 4.6 (ii), and Galois cohomology (see , Ch. I, 4.2 and Proposition 15) that cd$`(H)=1`$ and cd$`(𝒢(E(p)/E))=2`$. Hence, by Teichmüller’s theorem (cf. , Ch. 9, Theorem 4) and Proposition 2.2 (iii), (iii)$``$(iv). To prove that (iv)$``$(iii) we need the following results (see , (11.5), and , Proposition 3.26):
(5.1) If $`F`$ is a field containing a primitive $`p`$-th root of unity, and $`M`$ is a finite Galois extension of $`F`$ in $`F(p)`$, then there exists an isomorphism $`\kappa _M:H^2(𝒢(F(p)/M),𝔽_p)`$ $`{}_{p}{}^{}\mathrm{Br}(M)`$ as $`[𝒢(M/F)]`$-modules, and the compositions Cor$`{}_{M/F}{}^{}\kappa _M`$ and $`\kappa _F`$ $`\mathrm{cor}_{M/F}`$ coincide, where cor<sub>M/F</sub> is the corestriction map $`H^2(𝒢(F(p)/M),𝔽_p)H^2(𝒢(F(p)/F),𝔽_p)`$. Hence, Cor<sub>M/F</sub> is surjective if and only if so is cor<sub>M/F</sub>.
It suffices to show that Br$`(E^{}/E)=`$ $`{}_{p}{}^{}\mathrm{Br}(E)`$, for an arbitrary degree $`p`$ extension $`E^{}`$ of $`E`$ in $`E(p)`$. The equality cd$`(𝒢(E(p)/E))=2`$ ensures that Br$`(E)_p\{0\}`$ and it follows from (5.1) and , Proposition 3.3.8, that Cor$`_{E^{}/E}`$ maps $`{}_{p}{}^{}\mathrm{Br}(E^{})`$ surjectively on $`{}_{p}{}^{}\mathrm{Br}(E)`$. We show that $`{}_{p}{}^{}\mathrm{Br}(E^{})`$ includes the preimage $`\mathrm{\Pi }_p(E)`$ of $`{}_{p}{}^{}\mathrm{Br}(E)`$ in Br$`(E^{})_p`$ under Cor$`_{E^{}/E}`$. Our assumptions guarantee (in conjunction with Galois theory and the normality of maximal subgroups of finite $`p`$-groups, see , Ch. VIII and Ch. I, Sect. 6) that $`E^{}/E`$ is cyclic. Observe also that Br$`(E^{})_p`$ is a trivial $`[𝒢(E^{}/E)]`$-module. Indeed, let $`\sigma `$ be a generator of $`𝒢(E^{}/E)`$, and let $`b\mathrm{Br}(E^{})`$ be an element of order $`p^k`$, for some $`k`$. Proceeding by induction on $`k`$, one sees that it suffices to prove that $`\sigma (b)=b`$, provided that $`k2`$ and $`\sigma (b)=b+b_0`$, for some $`b_0_p\mathrm{Br}(E^{})`$. The assumption on $`b_0`$ indicates that Cor$`{}_{E^{}/E}{}^{}(b_0)=0`$, so it follows from (5.1) and , Corollary, that $`b_0=a\sigma (a)`$, for some $`a`$ $`{}_{p}{}^{}\mathrm{Br}(E^{})`$. This yields $`\sigma (b+a)=b+a`$, which implies $`\sigma (b)=b`$ and our assertion about Br$`(E^{})_p`$ (because $`{}_{p}{}^{}\mathrm{Br}(E^{})`$ is a trivial $`[𝒢(E^{}/E)]`$-module). By Teichmüller’s theorem, the established property of Br$`(E^{})_p`$ is equivalent to the inclusion of Br$`(E^{})_p`$ in the image of $`\rho _{E/E^{}}`$. Since Cor$`_{E^{}/E}`$ maps $`{}_{p}{}^{}\mathrm{Br}(E^{})`$ surjectively upon $`{}_{p}{}^{}\mathrm{Br}(E)`$, it is now easy to obtain from the RC-formula that $`\mathrm{\Pi }_p(E)`$ $`{}_{p}{}^{}\mathrm{Br}(E^{})`$. Moreover, it becomes clear that the implication (iv)$``$(iii) will follow, if we show that $`{}_{p}{}^{}\mathrm{Br}(E)`$ is included in the subgroup $`p\mathrm{Br}(E)_p=\{p\beta :\beta \mathrm{Br}(E)_p\}`$ of Br$`(E)_p`$. Denote by $`B`$ the extension of $`E`$ in $`E(p)`$ obtained by adjunction of a primitive $`p^2`$-th root of unity. It follows from , (16.1), , Sect. 15.1, Corollary b, and Kummer theory that $`{}_{p}{}^{}\mathrm{Br}(B)p\mathrm{Br}(B)_p`$. Hence, the surjectivity of the homomorphism $`{}_{p}{}^{}\mathrm{Br}(B)`$ $`{}_{p}{}^{}\mathrm{Br}(E)`$ induced by Cor<sub>B/E</sub> implies $`{}_{p}{}^{}\mathrm{Br}(E)p\mathrm{Br}(E)_p`$. Therefore, Br$`(E^{}/E)=`$ $`{}_{p}{}^{}\mathrm{Br}(E)`$ and (iv)$``$(iii), so Proposition 5.1 is proved. ∎
###### Corollary 5.2.
Let $`E`$ be an SQL-field, such that $`\mathrm{cd}_p(𝒢_E)0`$, for a given $`p`$. Then the Sylow pro-$`p`$-subgroups of $`𝒢_E`$ are $`p`$-groups of Demushkin type of degree $`d_p1`$, unless $`p=\mathrm{char}(E)`$ or $`E`$ is formally real and $`p=2`$.
###### Proof.
One may consider only the special case of $`E_{\mathrm{sep}}=E(p)`$ (see the proof of Proposition 3.6). Then Br$`(E)_p\{0\}`$ and our conclusion follows from Propositions 5.1, 2.2 (ii) and , Ch. II, Proposition 3. ∎
It is easily seen that the degrees of $`p`$-groups of Demushkin type are bounded by their ranks. The following conversion of this fact can be deduced from , (11.5), Proposition 5.1, , I, Theorem 8.1, and the sufficiency part of (2.2) (by specifying the cardinalities of the fields $`E_0`$ and $`E`$ in (2.2) (ii), see , Remark 5.4, for more details):
(5.2) For any system $`d\mathrm{}_0`$ and $`d_pd:p`$, of cardinal numbers, there is a field $`E`$ containing a primitive $`p`$-th root of unity, for every $`p`$, and such that $`𝒢(E(p)/E)`$ and the Sylow pro-$`p`$-subgroups of $`𝒢_E`$ are of rank $`d`$, Demushkin type and degree $`d_p`$.
By a Demushkin pro-$`p`$-group, we mean a $`p`$-group of Demushkin type of degree $`1`$. Demushkin pro-$`p`$-groups of rank $`r(P)\mathrm{}_0`$ have been classified by Demushkin, Labute and Serre (cf. and further references there). When $`r(P)=\mathrm{}_0`$, by and , $`P`$ has $`s`$-invariant zero if and only if $`P𝒢(F(p)/F)`$, for some field $`F`$. Hence, by applying Lemma 3.5 of (in a modified form adjusted to the case singled out by , Proposition 3.1 (iii)), and arguing as in the proof of (5.2), one supplements it as follows:
(5.3) For each sequence $`G_p:p`$, of Demushkin pro-$`p`$-groups of rank $`\mathrm{}_0`$ and $`s`$-invariant zero, there exists a field $`E`$ such that $`𝒢(E(p)/E)`$ and the Sylow pro-$`p`$-subgroups of $`𝒢_E`$ are isomorphic to $`G_p`$ when $`p`$ ranges over $``$.
Next we show that the inequality $`d\mathrm{}_0`$ in (5.2) is essential. This result is a special case of , Corollary 5 (see also , Sects. 1 and 2). For convenience of the reader, we present it here with a short proof based on Theorem 3.1 and Proposition 5.1.
###### Corollary 5.3.
Let $`E`$ be a field containing a primitive $`p`$-th root of unity, for some $`pP(E)`$, and with $`𝒢(E(p)/E)`$ a $`p`$-group of Demushkin type of finite rank $`r(p)`$. Then $`𝒢(E(p)/E)`$ is a Demushkin group or a free pro-$`p`$-group.
###### Proof.
In view of , Lemma 7, and Proposition 5.1 (with its proof), one may consider only the case where $`r(p)3`$ and Br$`(E)_p\{0\}`$. Take a field $`E^{}\mathrm{\Omega }(E)`$ so that $`𝒢(E^{}/E)`$ has exponent $`p`$ and order $`p^{r(p)1}`$. By Theorem 3.1 and , (11.5), then $`E^pN(E^{}/E)`$ and $`N(E^{}/E)`$ is a subgroup of $`E^{}`$ of index $`p^{(r(p)1).d_p}`$, so we have $`(r(p)1).d_p`$ $`r(p)`$, where $`d_p`$ is the degree of $`𝒢(E(p)/E)`$. In our case, this implies $`d_p=1`$, which proves Corollary 5.3. ∎
###### Remark 5.4.
Statement (5.1) and Corollary 5.3 show that the equivalence (ii)$``$(iv) in Proposition 5.1 generalizes the concluding assertion of , Theorem 1 (independently of the elementary type conjecture formulated in ).
Corollary 5.3 gives us the possibility to determine the structure of the continuous character group $`C(E(p)/E)`$ of $`𝒢(E(p)/E)`$, for a $`p`$-quasilocal nonreal field $`E`$ containing a primitive $`p`$-th root of unity. Recall that $`C(E(p)/E)`$ is an abelian torsion $`p`$-group, whence, it decomposes into the direct sum $`D(E(p)/E)R(E(p)/E)`$, where $`D(E(p)/E)`$ is the maximal divisible subgroup of $`C(E(p)/E)`$ and $`R(E(p)/E)`$ is a (reduced) subgroup of $`C(E(p)/E)`$ isomorphic to $`C(E(p)/E)/D(E(p)/E)`$ (see , Ch. 7, Sect. 5, and , Theorem 24.5). With this notation, our next result can be stated as follows:
###### Proposition 5.5.
Let $`E`$ be a $`p`$-quasilocal nonreal field, $`\mu _p(E)`$ the group of roots of unity in $`E`$ of $`p`$-primary degrees, and $`r_p(E)`$ the rank of $`𝒢(E(p)/E)`$ as a pro-$`p`$-group. Suppose that $`\mu _p(E)\{1\}`$ and $`\epsilon _pE`$ is a primitive $`p`$-th root of unity. Then:
(a) $`C(E(p)/E)=D(E(p)/E)`$ if and only if $`\mu _p(E)`$ is infinite or Br$`(E)_p=\{0\}`$; when Br$`(E)_p\{0\}`$ and $`\mu _p(E)`$ is finite of order $`p^\nu `$, the group $`R(E(p)/E)`$ is isomorphic to the maximal subgroup of Br$`(E)`$ of period $`p^\nu `$;
(b) Br$`(E)_p`$ is embeddable as a subgroup of $`D(E(p)/E)`$.
###### Proof.
(a): It follows from Kummer theory that $`C(E(p)/E)=D(E(p)/E)`$, provided that $`\mu _p(E)`$ is infinite. We show that the same equality holds in the case of Br$`(E)_p=\{0\}`$. Then it follows from , Sect. 15.1, Proposition b, that $`\epsilon _p`$ lies in the norm group $`N(L^{}/E)`$, for every cyclic extension $`L^{}`$ of $`E`$ in $`E(p)`$; hence, by Albert’s height theorem (cf. , Ch. IX, Sect. 6, and , Sect. 2), there is a cyclic extension $`L_1^{}`$ of $`E`$ in $`E(p)`$, such that $`L^{}I(L_1^{}/E)`$ and $`[L_1^{}:L^{}]=p`$. This implies $`L^{}I(L_1/E)`$, for some $`_p`$-extension $`L_1`$ of $`E`$ in $`E(p)`$, and so proves that $`C(E(p)/E)=D(E(p)/E)`$. It remains to consider the case where Br$`(E)_p\{0\}`$ and $`\mu _p(E)`$ has finite order $`p^\nu `$. As $`\mu _p(E)\{1\}`$, then we have $`E(p)E`$, so it follows from , Theorem 2, and the condition that $`E`$ is a nonreal field that $`E(p)`$ contains as a subfield a $`_p`$-extension of $`E`$. In view of Galois theory, this means that $`C(E(p)/E)`$ possesses a quasicyclic $`p`$-subgroup, which proves that $`D(E(p)/E)\{0\}`$. Note further that Br$`(E)_p`$ is a divisible group, since $`E`$ is nonreal and $`p`$-quasilocal (cf. , I, Theorem 3.1). Therefore, by , Theorem 23.1, Br$`(E)_p`$ decomposes into the direct sum $`(p^{\mathrm{}})^{d(p)}`$ of isomorphic copies of the quasicyclic $`p`$-group $`(p^{\mathrm{}})`$, indexed by a set $`I`$ of cardinality $`d(p)`$ equal to the dimension of $`{}_{p}{}^{}\mathrm{Br}(E)H^2(𝒢(E(p)/E),𝔽_p)`$ as an $`𝔽_p`$-vector space. Moreover, it becomes clear that, for each $`m`$, the maximal subgroup of Br$`(E)_p`$ of period $`p^m`$ decomposes into a direct sum of cyclic groups of order $`p^m`$, indexed by $`I`$. Thus the latter conclusion of Proposition 5.5 (a) is equivalent to the former part of the following assertions:
(5.4) (a) $`C(E(p)/E)`$ and the direct sum $`D(E(p)/E)\mu _p(E)^{d(p)}`$ are isomorphic, where $`\mu _p(E)^{d(p)}`$ is a direct sum of isomorphic copies of $`\mu _p(E)`$, indexed by a set of cardinality $`d(p)`$ (for a proof, see , II, Lemma 2.3).
(b) A cyclic extension $`M`$ of $`E`$ in $`E(p)`$ is a subfield of a $`_p`$-extension of $`E`$ in $`E(p)`$ if and only if there is $`M^{}I(E(p)/M)`$, such that $`M^{}/E`$ is cyclic and $`[M^{}:M]=p^\nu `$; this is the case if and only if $`\mu _p(E)N(M/E)`$.
The former part of (5.4) (b) is implied by (5.4) (a) and Galois theory, and the latter one follows from Albert’s height theorem.
(b): Proposition 5.5 (a) allows us to consider only the case where Br$`(E)_p\{0\}`$ and $`\mu _p(E)`$ has order $`p^\nu `$, for some $`\nu `$. Then it follows from (5.4) (a) and the nontriviality of $`D(E(p)/E)`$ that $`r_p(E)2`$. Using the observations preceding the statement of (5.4), one also sees that it is sufficient to prove the embeddability of $`{}_{p}{}^{}\mathrm{Br}(E)`$ in $`D(E(p)/E)`$. Let now $`\delta _\nu `$ be a primitive $`p^\nu `$-th root of unity, and $`M_\lambda `$ be an extension of $`E`$ generated by a $`p`$-th root $`\eta _\lambda E(p)`$ of an element $`\lambda E^{}E^p`$. Then $`M_\lambda /E`$ is cyclic, $`[M_\lambda :E]=p`$ and $`𝒢(M_\lambda /E)`$ contains a generator $`\sigma _\lambda `$, such that the cyclic $`E`$-algebra $`(M_\lambda /E,\sigma _\lambda ,\delta _\nu )`$ is isomorphic to the symbol $`E`$-algebra $`A_{\epsilon _p}(\lambda ,\delta _\nu ;E)`$. It is well-known that $`A_{\epsilon _p}(\lambda ,\delta _\nu ;E)`$ and $`A_{\epsilon _p}(\delta _\nu ,\lambda ;E)`$ are inversely-isomorphic $`E`$-algebras. This, combined with , Sect. 15.1, Proposition b, implies $`\delta _\nu N(M_\lambda /E)`$ if and only if $`\lambda N(M_{\delta _\nu }/E)`$. Hence, by (5.4) (b), the divisibility of Br$`(E)_p`$, and Theorem 23.1 of , Proposition 5.5 (b) and the latter assertion of Proposition 5.5 (a) are equivalent to the statement that $`{}_{p}{}^{}\mathrm{Br}(E)`$ embeds as a subgroup of $`N(M_{\delta _\nu }/E)/E^p`$. Since $`N(M_{\delta _\nu }/E)/E^p`$ is an abelian group of period $`p`$, this amounts to proving that $`{}_{p}{}^{}\mathrm{Br}(E)`$ is its homomorphic image. We show that $`{}_{p}{}^{}\mathrm{Br}(E)`$ is a homomorphic image of $`N(M_\mu /E)/E^p`$, for an arbitrary element $`\mu E^{}E^p`$. Fix $`\mu ^{}E^{}E^p`$ so that $`M_\mu ^{}M_\mu `$. Then $`E^{}/N(M_\mu ^{}/E)`$ $`{}_{p}{}^{}\mathrm{Br}(E)`$, by , Sect. 15.1, Proposition b (and the $`p`$-quasilocality of $`E`$), and $`N(M_\mu /E)N(M_\mu ^{}/E)=E^{}`$, by , I, Lemma 4.3. Since, by Theorem 3.1, $`N(M_\mu /E)N(M_\mu ^{}/E)=N(M_\mu M_\mu ^{}/E)`$, this yields
$`E^{}/N(M_\mu ^{}/E)N(M_\mu /E)/N(M_\mu M_\mu ^{}/E)`$, $`E^pN(M_\mu M_\mu ^{}/E)`$ and
$$N(M_\mu /E)/N(M_\mu M_\mu ^{}/E)(N(M_\mu /E)/E^p)/(N(M_\mu M_\mu ^{}/E)/E^p);$$
in particular, $`{}_{p}{}^{}\mathrm{Br}(E)`$ is a homomorphic image of $`N(M_\mu /E)/E^p`$, which completes the proof of Proposition 5.5 (b). ∎
The concluding results of this Section present applications of Proposition 5.5 and Corollary 5.3 to the study of $`p`$-primary index-exponent $`K`$-pairs, for a Henselian field $`(K,v)`$ with a $`p`$-quasilocal field $`\widehat{K}`$, for some $`p`$. Recall first that, for any field $`E`$ and each $`Dd(E)`$, exp$`(D)`$ divides ind$`(D)`$ and is divisible by any $`p`$ dividing ind$`(D)`$; in addition, $`(\mathrm{ind}(D),\mathrm{exp}(D))`$is obtained as a componentwise product of index-exponent $`l`$-pairs, where $`l`$ runs across the set of prime divisors of ind$`(D)`$ (see , Sect. 14.4). Index-exponent relations of algebras from $`d(E)`$ depend essentially on specific properties of $`E`$, and their description reduces to the special case of algebras $`D_pd(E)`$ of $`p`$-primary dimensions, for an arbitrary $`p`$. The study of index-exponent $`p`$-primary $`E`$-pairs relies on the knowledge of the Brauer $`p`$-dimension Brd$`{}_{p}{}^{}(E)`$, defined as the least integer $`b(p)0`$, for which ind$`(D_p)`$ divides exp$`(D_p)^{b(p)}`$ whenever $`D_pd(E)`$ and $`[D_p]\mathrm{Br}(E)_p`$; when no such $`b(p)`$ exists, we say that Brd$`{}_{p}{}^{}(E)`$ is infinite. Note that $`(p^k,p^n):k,n,kn`$, are index-exponent $`L`$-pairs whenever $`(L,\lambda )`$ is a Henselian field, such that the quotient group $`\lambda (L)/p\lambda (L)`$ of the value group $`\lambda (L)`$ is infinite, $`\widehat{L}`$ is a nonreal field and $`\mu _p(\widehat{L})\{1\}`$. This is demonstrated by the proof of , Corollary 4.5, which also shows that the result does not change, if the condition on $`\mu _p(\widehat{L})`$ is replaced by the one that the rank $`r_p(\widehat{L})`$ of $`𝒢(\widehat{L}(p)/\widehat{L})`$ is infinite. Therefore, we restrict our considerations to the case where $`\widehat{K}`$ is a $`p`$-quasilocal nonreal field, $`\mu _p(\widehat{K})\{1\}`$ and $`v(K)/pv(K)`$ is a nontrivial finite group. Then Brd$`{}_{p}{}^{}(K)`$ is determined by , Theorem 4.1, as follows:
(5.5) If $`v(K)/pv(K)`$ has order $`p^{\tau (p)}`$, then Brd$`{}_{p}{}^{}(K)=[(m_p(\widehat{K})+\tau (p))/2]`$, where $`m_p(\widehat{K})=\mathrm{min}\{r_p(\widehat{K}),\tau (p)\}`$.
Relying on Proposition 5.5, we first show that if $`r_p(\widehat{K})=\mathrm{}`$, then index-exponent $`p`$-primary $`K`$-pairs are fully determined by Brd$`{}_{p}{}^{}(K)`$.
###### Corollary 5.6.
Let $`(K,v)`$ be a Henselian field with $`\widehat{K}`$ satisfying the conditions of Proposition 5.5, for some $`p`$. Assume also that $`r_p(\widehat{K})<\mathrm{}`$, $`v(K)pv(K)`$ and $`v(K)/pv(K)`$ is finite of order $`p^{\tau (p)}`$. Then:
(a) $`\widehat{K}`$ has Galois extensions $`U_n`$, $`U_n^{}`$ in $`E(p)`$, $`n`$, such that $`[U_n:E]=p`$, $`[U_1\mathrm{}U_n:E]=p^n`$, $`𝒢(U_n^{}/E)_p`$ and $`U_nI(U_n^{}/E)`$, for each $`n`$;
(b) $`(p^k,p^n):k,n,nkn\tau (p)`$, are all nontrivial $`p`$-primary index-exponent pairs over $`K`$.
###### Proof.
(a): In view of (5.4) (b) and Galois theory, it suffices to prove that the maximal subgroup $`{}_{p}{}^{}D(\widehat{K}(p)/\widehat{K})`$ of $`D(\widehat{K}(p)/\widehat{K})`$ of period $`p`$ is infinite. The infinity of $`{}_{p}{}^{}D(\widehat{K}(p)/\widehat{K})`$ follows from Proposition 5.5 (b), if $`{}_{p}{}^{}\mathrm{Br}(\widehat{K})`$ is infinite. Note also that $`r_p(\widehat{K})=\mathrm{}`$ if and only if $`\widehat{K}^{}/\widehat{K}^p`$ is infinite (cf. , Ch. I, 4.1), which holds if and only if $`C(\widehat{K}(p)/\widehat{K})`$ contains infinitely many elements of order $`p`$. It is therefore clear from (5.4) (a) that if $`{}_{p}{}^{}\mathrm{Br}(\widehat{K})`$ is finite, then $`{}_{p}{}^{}D(\widehat{K}(p)/\widehat{K})`$ is infinite, as required.
(b): It follows from Corollary 5.6 (a) and Galois theory that, for each finite abelian $`p`$-group $`G`$, there exists a Galois extension $`U_G`$ of $`K`$ in $`K_{\mathrm{ur}}`$ with $`𝒢(U_G/K)G`$. When the rank of $`G`$ is $`\tau (p)`$, one obtains from , Theorem 1 (or , Example 4.3), that there is a nicely semi-ramified (abbr, NSR) algebra $`N_Gd(K)`$, in the sense of , possessing a maximal subfield $`K`$-isomorphic to $`U_G`$; this ensures that ind$`(N_G)`$ and exp$`(N_G)`$ equal the order and the period of $`G`$, respectively. Thus it becomes clear that there exist $`N_{k,n}d(K):k,n,nk\tau (p)n`$, such that $`N_{k,n}/K`$ is NSR, ind$`(N_{k,n})=p^k`$ and exp$`(N_{k,n})=p^n`$. This proves Corollary 5.6 (b), since by , Theorem 4.1, the infinity of $`r_p(\widehat{K})`$ implies $`\tau (p)=\mathrm{Brd}_p(K)`$. ∎
Our next result supplements Corollary 5.6 as follows:
###### Corollary 5.7.
Assume that $`(K,v)`$ is a Henselian field, such that $`v(K)/pv(K)`$ is finite of order $`p^{\tau (p)}>1`$, and $`\widehat{K}`$ is $`p`$-quasilocal with $`\mu _p(\widehat{K})\{1\}`$, $`r_p(\widehat{K})<\mathrm{}`$, and $`C(\widehat{K}(p)/\widehat{K})=D(\widehat{K}(p)/\widehat{K})`$, for some $`p`$. Then nontrivial $`p`$-primary index-exponent $`K`$-pairs are described by as follows:
(a) $`(p^k,p^n):k,n,nkn\mathrm{Brd}_p(K)`$, if $`\mu _p(\widehat{K})`$ is infinite;
(b) $`(p^k,p^n):k,n`$, $`nknm_p(\widehat{K})+\mathrm{min}\{\nu ,n\}([(\tau (p)m_p(\widehat{K}))/2]`$, in case $`\mu _p(\widehat{K})`$ is finite of order $`p^\nu `$, and $`m_p(\widehat{K})`$ is defined as in (5.5).
###### Proof.
The description of index-exponent $`p`$-primary $`K`$-pairs is obtained in the case of Corollary 5.7 (b) by the method of proving , Lemma 5.1 (replacing min$`\{r_p(\widehat{K})1,\tau (p)\}`$ by $`m_p(\widehat{K})`$). Henceforth, we assume that $`\mu _p(\widehat{K})`$ is infinite. Note that, for any finite abelian $`p`$-group $`G`$ of rank $`m_p(\widehat{K})`$, there exist a Galois extension $`U_G`$ of $`K`$ in $`K_{\mathrm{ur}}`$, and an algebra $`N_Gd(K)`$, such that $`𝒢(U_G/K)G`$, $`N_G/K`$ is NSR and $`U_G`$ is $`K`$-isomorphic to a maximal subfield of $`N_G`$. As in the proof of Corollary 5.6 (b), one obtains further that there exist NSR-algebras $`N_{k,n}d(K):k,n`$, $`nkm_p(\widehat{K})n`$, with ind$`(N_{k,n})=p^k`$ and exp$`(N_{k,n})=p^n`$, for each admissible pair $`k,n`$. Since Brd$`{}_{p}{}^{}(K)=[(m_p(\widehat{K})+\tau (p))/2]`$, this fact proves Corollary 5.7 in case $`r_p(\widehat{K})\tau (p)1`$. Suppose now that $`\tau (p)m_p(\widehat{K})2`$, put $`\overline{m}=\mathrm{Brd}_p(K)=[(m_p(\widehat{K})+\tau (p))/2]`$, fix a divisible hull $`v(N_G)^{}`$ of $`v(N_G)`$, and take a finite abelian $`p`$-group $`H`$ of rank $`r(H)[(\tau (p)m_p(\widehat{K}))/2]`$. Using , Theorem 1, and the natural bijection between $`I(Y/K)`$ and the set of subgroups of $`v(Y)/v(K)`$, for any totally ramified finite abelian $`p`$-extension $`Y/K`$ (cf. , Ch. 3, Theorem 2), one obtains that there is $`T_Hd(K)`$ with the following properties: $`T_H/K`$ is totally ramified, $`v(T_H)v(N_G)^{}`$ and $`v(T_H)/v(K)`$ is isomorphic to $`HH`$; $`v(N_G)v(T_H)=v(K)`$ and $`T_H_KK_{\mathrm{ur}}d(K_{\mathrm{ur}})`$; $`T_H`$ is a tensor product of $`r(H)`$ cyclic totally ramified $`K`$-algebras; ind$`(T_H)`$ and exp$`(T_H)`$ equal the order and the period of $`H/v(K)`$, respectively (see, e.g., , for more details). Note also that, by , Theorem 1, $`N_G_KT_Gd(K)`$. These observations indicate that, for any $`n`$, there exist $`\mathrm{\Delta }_nd(K)`$ and $`T_{n,\rho }d(K):\rho =1,\mathrm{}n\theta `$, where $`\theta =[(\tau (p)m_p(\widehat{K}))/2]`$, satisfying the following conditions:
(5.6) (a) $`\mathrm{\Delta }_n/K`$ is NSR, exp$`(\mathrm{\Delta }_n)=p^n`$ and ind$`(\mathrm{\Delta }_n)=p^{n.\overline{m}}`$;
(b) For each index $`\rho `$, $`T_{n,\rho }/K`$ is totally ramified, $`\mathrm{\Delta }_n_KT_{n,\rho }d(K)`$, exp$`(\mathrm{\Delta }_n_KT_{n,\rho })=p^n`$, and ind$`(T_{n,\rho })=p^\rho `$.
Statements (5.6) (b) show that ind$`(\mathrm{\Delta }_n_KT_{n,\rho })=p^{n+\rho }`$, $`\rho =1,\mathrm{},\theta `$, which completes the proof of Corollary 5.7. ∎
Summing-up Corollaries 5.6, 5.7 and , Lemma 5.1, and combining the latter with Corollary 5.3, one one fully describes $`p`$-primary index-exponent $`K`$-pairs, for a Henselian field $`(K,v)`$, such that $`v(K)pv(K)`$ and $`\widehat{K}`$ is nonreal and $`p`$-quasilocal with $`\mu _p(\widehat{K})\{1\}`$. We refer the reader to , Corollary 2.2 and Remark 4.2, for an analogous description in the case where $`(K,v)`$ is Henselian, $`p=2`$, and $`\widehat{K}`$ is formally real $`2`$-quasilocal.
## 6. Preparation for the proof of Theorems 3.2 and 1.3
Let $`E`$ be a field, $`M\mathrm{\Omega }(E)`$, $`\mathrm{\Pi }`$ the set of prime divisors of $`[M:E]`$, and $`M_p=ME(p)`$, for each $`p\mathrm{\Pi }`$. When $`ME`$, the homomorphism group $`\mathrm{Hom}(E^{},𝒢(M/E))`$ is isomorphic to the direct group product $`_{p\mathrm{\Pi }}\mathrm{Hom}(E^{},`$ $`𝒢(M_p/E))`$, so Theorem 1.3 can be deduced from Theorem 3.2, Lemma 2.1 and the primary tensor decomposition of cyclic $`E`$-algebras (cf. , Sect. 15.3). The results of this Section serve as a basis for the proof of Theorem 3.2, presented in Section 7. Our starting point are the following two lemmas.
###### Lemma 6.1.
Let $`E`$ and $`M`$ be fields, such that $`M\mathrm{\Omega }_p(E)`$, for some $`pP(E)`$. Suppose that $`𝒢(M/E)`$ has rank $`t1`$ as a $`p`$-group, and $`F`$ is an intermediate field of $`M/E`$ of degree $`[F:E]=p`$. Then there exist cyclic extensions $`E_1,\mathrm{},E_t`$ of $`E`$ in $`M`$ with $`E_1\mathrm{}E_t=M`$, $`_{i=1}^t[E_i:E]=[M:E]`$ and $`_{i=1}^t[(FE_i):F]=[M:F]`$.
###### Proof.
In view of Galois theory, this is equivalent to the following statement:
(6.1) Let $`G`$ be a finite abelian $`p`$-group of rank $`t1`$ and $`H`$ a maximal subgroup of $`G`$. Then these exist cyclic subgroups $`G_1,\mathrm{},G_t`$ of $`G`$, such that the (inner) products of $`G_i:i=1,\mathrm{},t`$, and $`HG_i:i=1,\mathrm{},t`$, are direct and equal to $`G`$ and $`H`$, respectively.
To prove (6.1) take a cyclic subgroup $`G_1G`$ of maximal order so that the order of the group $`H_1=HG_1`$ equals the exponent of $`H`$. The choice of $`G_1`$ ensures that $`G_1G_0=G`$ whenever $`G_0`$ is maximal among the subgroups of $`G`$, which trivially intersect $`G_1`$ (cf. , Sects. 15 and 27). This implies the existence of subgroups $`H_1^{}H`$ and $`G_1^{}G`$, such that $`H_1H_1^{}=G_1G_1^{}=\{1\}`$, $`H_1^{}G_1^{}`$ and the products $`H_1H_1^{}`$ and $`G_1G_1^{}`$ are equal to $`H`$ and $`G`$, respectively. Therefore, (6.1) can be proved by induction on $`t`$. ∎
###### Lemma 6.2.
Assume that $`E`$, $`M`$, $`p`$ and $`t`$ satisfy the conditions of Lemma 6.1, and let $`F`$ be a maximal subfield of $`M`$ including $`E`$. Then there exist cyclic extensions $`E_1,\mathrm{},E_t`$ of $`E`$ in $`M`$, such that $`_{i=1}^t[E_i:E]=[M:E]`$, $`_{i=1}^t[(E_iF):E]=[F:E]`$, and the compositums $`E_1\mathrm{}E_t`$ and $`(E_iF)\mathrm{}(E_tF)`$ are equal to $`M`$ and $`F`$, respectively.
###### Proof.
This is equivalent to the following statement:
(6.2) Let $`G`$ be a finite abelian $`p`$-group of rank $`t1`$, $`H`$ a subgroup of $`G`$ of order $`p`$, and $`\pi `$ the natural homomorphism of $`G`$ on $`G/H`$. Then $`G`$ contains elements $`g_1,\mathrm{},g_t`$, for which the products of the cyclic groups $`g_1,\mathrm{},g_t`$ and $`\pi (g_1),\mathrm{},\pi (g_t)`$ are direct and equal to $`G`$ and $`G/H`$, respectively.
For the proof of (6.2), consider a cyclic subgroup $`C_1G`$ of maximal order, and such that $`C_1H`$ is of minimal possible order. Then one can find $`C_1^{}G`$ so that $`C_1C_1^{}=\{1\}`$, $`C_1C_1^{}=G`$ and $`HC_1C_1^{}`$. This implies $`G/H`$ is isomorphic to the direct product $`C_1H/H\times C_1^{}H/H`$, which allows one to prove (6.2) by induction on $`t`$. ∎
Let now $`G`$ be a finite abelian $`p`$-group of rank $`t2`$. A subset $`g=\{g_1,\mathrm{},g_t\}`$ of $`G`$ is called an ordered basis of $`G`$, if $`g`$ is a basis of $`G`$ (i.e. the group product $`g_1\mathrm{}g_t`$ is direct and equal to $`G`$) and the orders of the elements of $`g`$ satisfy the inequalities $`o(g_1)\mathrm{}o(g_t)`$. It is known that the automorphism group Aut$`(G)`$ acts transitively on the set Ob$`(G)`$ of ordered bases of $`G`$. Arguing by induction on $`t`$, and using the fact that cyclic subgroups of $`G`$ of order $`o(g_1)`$ are direct summands in $`G`$ (see , Sect. 27), one obtains that Aut$`(G)`$ has the following system of generators:
(6.3) Aut$`(G)=d(k,m;h),t(i,j,s;h):h\mathrm{Ob}(G)`$, where $`i,j,k,s`$ and $`m`$ are integers with $`1k,i,jt`$, $`ij`$, $`1<m<o(h_k)`$, $`1s<o(h_i)`$, $`\mathrm{gcd}(m,o(h_k))=1`$ and max$`\{1,o(h_i)/o(h_j)\}s`$, and $`d(k,m;h)`$, $`t(i,j,s;h)`$ are defined by the data $`d(k,m;h)(h_k)=h_k^m`$, $`d(k,m;h)(h_k^{})=h_k^{}:k^{}k`$, and $`t(i,j,s;h)(h_j)=h_jh_i^s`$, $`t(i,j,s;h)`$ $`(h_j^{})=h_j^{}:j^{}j`$.
This allows us to prove the following two lemmas without serious technical difficulties.
###### Lemma 6.3.
Let $`E`$, $`M`$, $`p`$ and $`t`$ satisfy the conditions of Lemma 6.1, and let $`E`$ be a $`p`$-quasilocal field. Suppose that $`B`$ a cyclic subgroup of $`\mathrm{Br}(E)_p`$ of order $`o(B)`$ divisible by $`[M:E]`$, $`b`$ is a generator of $`B`$ and $`B(M/E)`$ is the group of those $`\beta E^{}`$, for which $`[(L/E,\tau ,\beta )]B`$ whenever $`L/E`$ is a cyclic extension, $`LM`$ and $`\tau =𝒢(L/E)`$. Assume also that $`E_1,\mathrm{},E_t`$ are cyclic extensions of $`E`$ in $`M`$, such that $`E_1\mathrm{}E_t=M`$ and $`[M:E]=_{j=1}^t[E_j:E]`$, and for each index $`jt`$, let $`E_j^{}`$ be the compositum of the fields $`E_i:ij`$, $`\tau _j`$ a generator of $`𝒢(E_j/E)`$, and $`\sigma _j`$ the unique $`E_j^{}`$-automorphism of $`M`$ extending $`\tau _j`$. Then there exist a group homomorphism $`\omega _{M/E,b}:B(M/E)𝒢(M/E)`$, and elements $`c_1,\mathrm{},c_t`$ of $`E^{}`$ with the following properties:
(i) $`c_jN(E_j^{}/E)`$ and $`[(E_j/E,\tau _j,c_j)]=[o(B)/[E_j:E]].b`$, for each index $`j`$; the co-set $`c_jN(M/E)`$ is uniquely determined by $`b`$ and $`\tau _j`$;
(ii) $`\omega _{M/E,b}`$ is the unique homomorphism of $`B(M/E)`$ on $`𝒢(M/E)`$ mapping $`c_j`$ into $`\sigma _j:j=1,\mathrm{},t`$, and with a kernel equal to $`N(M/E)`$;
(iii) $`\omega _{M/E,b}`$ does not depend on the choice of the $`t`$-tuples $`(E_1,\mathrm{},E_t)`$ and $`(\tau _1,\mathrm{},\tau _t)`$; it induces a group isomorphism $`B(M/E)/N(M/E)𝒢(M/E)`$.
###### Proof.
If $`t=1`$, our assertions can be deduced from Theorem 3.1 and the general theory of cyclic algebras (see , Sect. 15.1). Henceforth, we assume that $`t2`$. Using consecutively Galois theory and Theorem 3.1, one obtains that $`N(E_j/E)N(E_j^{}/E)=E^{}:j=1,\mathrm{},t`$, which implies the existence of elements $`c_1,\mathrm{},c_t`$ of $`E^{}`$ with the properties required by Lemma 6.3 (i). Denote by $`T(M/E)`$ the subgroup of $`E^{}`$ generated by the set $`N(M/E)\{c_1,\mathrm{},c_t\}`$. It is easily verified that $`\sigma _1,\mathrm{},\sigma _t`$ and $`c_1N(M/E),\mathrm{},c_tN(M/E)`$ form bases of the groups $`𝒢(M/E)`$ and $`T(M/E)/`$ $`N(M/E)`$, respectively. Therefore, there exists a unique homomorphism $`\omega _{M/E,b}`$ of $`T(M/E)`$ on $`𝒢(M/E)`$, such that $`\omega _{M/E,b}(c_j)=\sigma _j:j=1,\mathrm{},t`$, and $`\mathrm{Ker}(\omega _{M/E,b})=N(M/E)`$. The mapping $`\omega _{M/E,b}`$ is surjective, so it induces canonically an isomorphism of $`T(M/E)/N(M/E)`$ on $`𝒢(M/E)`$. We aim at proving that $`T(M/E)=B(M/E)`$. Assuming that $`[M:E]=p^m`$, and proceeding by induction on $`m`$, one obtains that this can be deduced from Lemma 6.2 and , Sect. 15.1, Corollary b, if $`\omega _{M/E,b}`$ has the property claimed by the former part of Lemma 6.3 (iii). In order to establish this property (the crucial point in our proof), consider another basis $`\sigma _1^{},\mathrm{},\sigma _t^{}`$ of $`𝒢(M/E)`$, and for each $`j\{1,\mathrm{},t\}`$, let $`H_j`$ be the subgroup of $`𝒢(M/E)`$ generated by the elements $`\sigma _i^{}:ij`$, $`F_j`$ the fixed field of $`H_j`$, $`\tau _j^{}`$ the $`E`$-automorphism of $`F_j`$ induced by $`\sigma _j^{}`$, and $`F_j^{}`$ the compositum of all $`F_i`$ with $`ij`$. It follows from Galois theory that $`F_1\mathrm{}F_t=M`$, $`_{i=1}^t[F_i:E]=[M:E]`$, $`F_j/E`$ is a cyclic extension, $`𝒢(F_j/E)=\tau _j^{}`$, and $`\sigma _j^{}`$ is the unique $`F_j^{}`$-automorphism of $`M`$ extending $`\tau _j^{}`$ ($`j=1,\mathrm{},t`$). Since $`\omega _{M/E,b}`$ is surjective, it maps some elements $`c_1^{},\mathrm{},c_t^{}`$ of $`T(M/E)`$ into $`\sigma _1^{},\mathrm{},\sigma _t^{}`$, respectively. Applying Lemma 6.2 and , Sect. 15.1, Corollary b, one concludes that the proof of Lemma 6.3 will be complete, if we show that $`c_j^{}N(F_j^{}/E)`$ and $`[(F_j/E,\tau _j^{},c_j^{})]=[o(B)/[F_j:E]].b`$, for every index $`j`$.
It is clearly sufficient to consider the special case in which $`\{\sigma _1,\mathrm{},\sigma _t\}`$ and $`\{\sigma _1^{},\mathrm{},\sigma _t^{}\}`$ are ordered bases of $`𝒢(M/E)`$ (which implies $`[E_\rho :E]=[F_\rho :E]`$, $`\rho =1,\mathrm{},t`$). In view of , Sect. 15.1, Corollary a, and (6.3), one may assume in addition that $`\sigma _j^{}=\sigma _j\sigma _i^s`$, $`c_j^{}=c_jc_i^s`$, $`\sigma _u^{}=\sigma _u`$ and $`c_u^{}=c_u:uj`$, for some pair $`(i,j)`$ of different indices, and some integer $`s`$ satisfying the inequalities $`1s<o(\sigma _i)`$ and divisible by max$`\{1,o(\sigma _i)/o(\sigma _j)\}`$. Then it follows from Galois theory that $`F_iF_j=E_iE_j`$, $`F_w=E_w,\tau _w^{}=\tau _w:wi`$, and $`F_u^{}=E_u^{}:uj`$. One also sees that if $`t3`$ and $`y(\{1,\mathrm{},t\}\{i,j\})`$, then $`c_j^{}N(E_y/E)`$. Since the cyclic $`E`$-algebras $`(E_j/E,\tau _j,c_j)`$ and $`(E_j/E,\tau _j,c_j^{})`$ are isomorphic (see , Sect. 15.1, Proposition b), and by Lemma 4.1, $`N(F_j^{}/E)=_{uj}N(F_u/E)`$, this reduces the proof of Lemma 6.3 (ii) to the one of the assertions that $`c_j^{}N(F_i/E)`$ and there is an $`E`$-isomorphism $`(E_i/E,\tau _i,c_i)(F_i/E,\tau _i^{},c_i)`$. We first show that $`(E_i/E,\tau _i,c_i)(F_i/E,\tau _i^{},c_i)`$. The assumptions on $`E_i,F_i`$ and $`E_j=F_j`$ guarantee that the field $`E_iE_j=F_iE_j`$ is isomorphic as an $`E`$-algebra to $`E_i_EE_j`$ and $`F_i_EE_j`$ (cf. , Sect. 14.7, Lemma b). It is therefore easy to see from the general properties of tensor products (cf. , Sect. 9.2, Proposition c) that $`(E_i/E,\tau _i,c_i)_E(E_j/E,\tau _j,c_j)(F_i/E,\tau _i^{},c_i)_E(E_j/E,\tau _j,c_jc_i^s^{})`$ as $`E`$-algebras, where $`s^{}=s.[o(\sigma _j)/o(\sigma _i)]`$. Since $`s^{}`$ and $`c_iN(E_j/E)`$, there exists an $`E`$-isomorphism $`(E_j/E,\tau _j,c_j)`$ $`(E_j/E,\tau _j,c_jc_i^s^{})`$. The obtained results prove that $`(E_i/E,\tau _i,c_i)`$ and $`(F_i/E,`$ $`\tau _i^{},c_i)`$ are similar over $`E`$. As $`[E_i:E]=[F_i:E]`$, we also have $`[(E_i/E,\tau _i,c_i):E]=[(F_i/E,\tau _i^{},c_i):E]`$ $`=[E_i:E]^2`$, so it turns out that $`(E_i/E,\tau _i,c_i)(F_i/E,`$ $`\tau _i^{},c_i)`$, as claimed.
It remains to be seen that $`c_j^{}N(F_i/E)`$. Suppose first that $`F_iE_i=E`$, take elements $`\beta _iE_j`$, $`\beta _jE_i`$, $`\delta \mathrm{Br}(E)`$ and an algebra $`Dd(E)`$ so that $`N_E^{E_j}(\beta _i)=c_i`$, $`N_E^{E_i}(\beta _j)=c_j`$, $`([E_i:E].[E_j:E])\delta =b`$, and $`o(B)\delta =[D]`$, and denote by $`\phi _i`$ and $`\phi _j`$ the automorphisms of $`E_iE_j`$ induced by $`\sigma _i`$ and $`\sigma _j`$, respectively. It is easily verified that $`[E_j:E].[D]=[(E_i/E,\tau _i,c_i)]`$ and $`𝒢(E_iE_j/E_j)=\phi _i`$. This, combined with the fact that $`E_iE_j=F_iE_j`$, $`E_iE_j=E`$ and $`\phi _i`$ extends $`\tau _i`$, enables one to deduce from (4.1) (iii), Proposition 2.2 and the RC-formula that $`D_EE_j`$ is similar to the cyclic $`E_j`$-algebra $`(E_iE_j/E_j,\phi _i,\beta _i)`$. Since $`[E_i:E].[D]=[(E_j/E,\tau _j,c_j)]`$, $`𝒢(E_iE_j/E_i)=\phi _j`$, and $`\phi _j`$ is a prolongation of $`\tau _j`$, it is analogously proved that $`[D_EE_i]=[(E_iE_j/E_i,\phi _j,\beta _j)]`$ in Br$`(E_i)`$. At the same time, it is clear from Galois theory and the condition $`F_iE_i=E`$ that $`𝒢(E_iE_j/E)=\phi _j,\phi _i^s`$ (i.e. $`\mathrm{g}.\mathrm{c}.\mathrm{d}.(s,p)=1`$), $`[E_i:E][E_j:E]`$, and $`E_iF_i=E_i\mathrm{\Phi }`$, where $`\mathrm{\Phi }`$ is the extension of $`E`$ in $`E_j`$ of degree $`[\mathrm{\Phi }:E]=[E_i:E]`$. Note also that $`\phi _j`$ extends $`\tau _i^s`$, since $`\sigma _j^{}(\lambda )=\lambda `$, for every $`\lambda F_i`$. Observing now that $`[(E_i/E,\tau _i,c_i)]=\omega .[(E_j/E,\tau _j,c_j)]`$ and $`[(E_iF_i/E_i,\psi _j,\beta _j)]=\omega .[(E_iE_j/E_i,\phi _j,\beta _j)]`$, where $`\omega =[E_j:E]/[E_i:E]=[E_iE_j:E_iF_i]`$ and $`\psi _j`$ is the automorphism of $`E_iF_i`$ induced by $`\phi _j`$ (cf. , Sect. 15.1, Corollary b), one obtains from (4.1) (iii) that $`(E_i/E,\tau _i,c_i)(\mathrm{\Phi }/E,\overline{\psi }_j,c_j)(F_i/E,\tau _i^s,c_j)`$ over $`E`$ ($`\overline{\psi }_j`$ being the automorphism of $`\mathrm{\Phi }`$ induced by $`\phi _j`$). Since $`(E_i/E,\tau _i,c_i)(F_i/E,\tau _i^{},c_i)`$, these results indicate that $`(F_i/E,\tau _i^s,c_j)(F_i/E,\tau _i^s,c_i^s)`$, which means that $`c_j^{}N(F_i/E)`$.
Let now $`E_iF_i=\stackrel{~}{E}`$, $`[\stackrel{~}{E}:E]=\mu >1`$, $`s/\mu =\stackrel{~}{s}`$ and $`E_iE_j=E_{i,j}`$. Then $`E_{i,j}\mathrm{\Omega }(\stackrel{~}{E})`$ and $`𝒢(E_{i,j}/\stackrel{~}{E})=𝒢(E_{i,j}/E_i)𝒢(E_{i,j}/F_i)`$, which implies $`\stackrel{~}{s}`$ and $`\mathrm{g}.\mathrm{c}.\mathrm{d}.(\stackrel{~}{s},p)=1`$. For each index $`ui`$, take $`\gamma _uE_u^{}`$ so that $`N_E^{E_u^{}}(\gamma _u)=c_u`$, and put $`\stackrel{~}{c}_u=N_{\stackrel{~}{E}}^{E_u^{}}(\gamma _u)`$, $`\stackrel{~}{E}_u=E_u\stackrel{~}{E}`$ and $`\stackrel{~}{F}_u=F_u\stackrel{~}{E}`$. Since $`E_iE_u=E`$, it follows from (4.1) (ii) that $`E_i\stackrel{~}{E}_u=\stackrel{~}{E}`$, $`[\stackrel{~}{E}_u:\stackrel{~}{E}]=[E_u:E]`$, $`\tau _u`$ uniquely extends to a $`\stackrel{~}{E}`$-automorphism $`\stackrel{~}{\tau }_u`$ of $`\stackrel{~}{E}_u`$, and $`𝒢(\stackrel{~}{E}_u/\stackrel{~}{E})=\stackrel{~}{\tau }_u`$. Note also that $`N_E^{E_u}(\theta _u)=N_{\stackrel{~}{E}}^{\stackrel{~}{E}_u}(\theta _u)`$, $`\theta _uE_u`$, whence $`c_iN(\stackrel{~}{E}_u/\stackrel{~}{E})`$. Fix elements $`b_\mu \mathrm{Br}(E)`$ and $`\mathrm{\Delta }_ud(E)`$ so that $`\mu .b_\mu =b`$ and $`[\mathrm{\Delta }_u]=[o(B)/[E_u:E]]b_\mu `$, and put $`b^{}=\rho _{E/\stackrel{~}{E}}(b_\mu )`$. It is clear from the double centralizer theorem (cf. , Sects. 12.7 and 13.3) that $`[(E_i/E,\tau _i,c_i)_E\stackrel{~}{E}]=[(E_i/\stackrel{~}{E},\tau _i^\mu ,c_i)]`$ in Br$`(\stackrel{~}{E})`$. Applying Proposition 2.2, (4.1) (iii) and the RC-formula, one obtains consecutively that $`o(B)`$ equals the order of $`b^{}`$ in Br$`(\stackrel{~}{E})`$, $`[(E_i/E,\tau _i,c_i)]=[o(B)/[E_i:\stackrel{~}{E}]]b_\mu `$, $`[(E_i/\stackrel{~}{E},\tau _i^\mu ,c_i)]=[o(B)/[E_i:\stackrel{~}{E}]].b^{}`$, and for each $`ui`$, $`[(E_u/E,\tau _u,c_u)]=\mu [\mathrm{\Delta }_u]`$ and $`[(\stackrel{~}{E}_u/\stackrel{~}{E},\stackrel{~}{\tau }_u,\stackrel{~}{c}_u)]=[\mathrm{\Delta }_u_E\stackrel{~}{E}]=[o(B)/[\stackrel{~}{E}_u:\stackrel{~}{E}]].b^{}`$. Consider now $`M/\stackrel{~}{E}`$, $`\stackrel{~}{s}`$, $`b^{}`$, the fields $`E_i,F_i`$, $`\stackrel{~}{E}_u,\stackrel{~}{F}_u:ui`$, and the algebras $`(E_i/\stackrel{~}{E},\tau _i^\mu ,c_i)`$, $`(F_i/\stackrel{~}{E},\tau _i^\mu ,c_i)`$, $`(\stackrel{~}{E}_j/\stackrel{~}{E},\stackrel{~}{\tau }_j,\stackrel{~}{c}_j)`$, instead of $`M/E`$, $`s`$, $`b`$, $`E_u^{},F_u^{}:u^{}=1,\mathrm{},t`$, and $`(E_i/E,\tau _i,c_i)`$, $`(F_i/E,\tau _i^{},c_i)`$, $`(E_j/E,\tau _j,`$ $`c_j)`$, respectively. As in the proof of our assertion in case $`\mathrm{g}.\mathrm{c}.\mathrm{d}.(s,p)=1`$, one concludes that $`\stackrel{~}{c}_jc_i^{\stackrel{~}{s}}N(F_i/\stackrel{~}{E})`$. Since $`N_E^{\stackrel{~}{E}}(\stackrel{~}{c}_jc_i^{\stackrel{~}{s}})=c_j^{}`$, this yields $`c_j^{}N(F_i/E)`$ (and $`T(M/E)=B(M/E)`$), which completes the proof of Lemma 6.3. ∎
###### Lemma 6.4.
In the setting of Lemma 6.3, let $`F`$ be an extension of $`E`$ in $`M`$, $`\pi _{M/F}`$ the natural projection of $`𝒢(M/E)`$ on $`𝒢(F/E)`$, $`\rho _{E/F}(b)=b^{}`$, $`B^{}=b^{}`$, and $`B_\kappa =b_\kappa `$, where $`b_\kappa =\kappa .b`$, for some $`\kappa `$ dividing $`o(B)/[M:E]`$. Then $`[M:E]o(B_\kappa )`$, $`[M:F]o(B^{})`$, and the mappings $`\omega _{M/E,b}`$, $`\omega _{M/E,b_\kappa }`$, $`\omega _{F/E,b}`$ and $`\omega _{M/F,b^{}}`$, determined as required by Lemma 6.3, are related as follows:
(i) $`\omega _{M/E,b}=\omega _{M/E,b_\kappa }`$ and $`\omega _{F/E,b}=\pi _{M/F}\omega _{M/E,b}`$;
(ii) $`\omega _{M/F,b^{}}(\gamma )=\omega _{M/E,b}(N_E^F(\gamma ))`$, for every $`\gamma F^{}`$.
###### Proof.
Clearly, $`o(B_\kappa )=o(B)/\kappa `$, and Proposition 2.2 implies that $`o(B^{})=o(B)/[F:E]`$, so the assertions that $`[M:E]o(B_\kappa )`$, $`[M:F]o(B^{})`$ and $`\omega _{M/E,b}=\omega _{M/E,b_\kappa }`$ are obvious. Note also that if $`t=1`$, then the remaining statements of the lemma can be deduced from the general properties of cyclic algebras. Suppose further that $`t2`$. First we prove that $`\omega _{F/E,b}=\pi _{M/F}\omega _{M/E,b}`$ in case $`F`$ is a maximal subfield of $`M`$. Let $`E_1,\mathrm{},E_t`$ be extensions of $`E`$ in $`M`$ with the properties described by Lemma 6.2. Then there exists an index $`j`$, such that $`E_jF`$ is a maximal subfield of $`E_j`$ and $`E_iF`$, for any other index $`i`$. This allows us to prove our assertion by applying Lemma 6.3 and , Sect. 15.1, Corollary b.
We turn to the proof of Lemma 6.4 (ii) in the special case of $`[F:E]=p`$. It is briefly presented, since we argue along the same lines as the concluding part of the proof of Lemma 6.3. Take $`E_1,E_1^{},\mathrm{},E_t,E_t^{}`$, $`\tau _1,\sigma _1,\mathrm{},\tau _t,\sigma _t`$ and $`c_1,\mathrm{},c_t`$ in accordance with Lemmas 6.1 and 6.3. Then $`FE_j`$ and $`E_uF=E:uj`$, for some index $`j`$. This implies $`FE_u^{}`$ and $`\tau _u`$ uniquely extends to an $`F`$-automorphism $`\stackrel{~}{\tau }_u`$ of the field $`E_uF:=\stackrel{~}{E}_u`$, when $`u`$ runs across the set $`W_j=\{1,\mathrm{},t\}\{j\}`$. Choose elements $`\alpha _1E_1^{},\mathrm{},\alpha _tE_t^{}`$ so that $`N_E^{E_k^{}}(\alpha _k)=c_k:k=1,\mathrm{},t`$, and put $`\stackrel{~}{E}_j=E_j`$, $`\stackrel{~}{c}_j=c_j`$, $`\stackrel{~}{\tau }_j=\tau _j^p`$, $`\stackrel{~}{\sigma }_j=\sigma _j^p`$, and $`\stackrel{~}{c}_u=N_F^{E_u^{}}(\alpha _u),\stackrel{~}{\sigma }_u=\sigma _u:uW_j`$. It is easily verified that $`N_E^F(\stackrel{~}{c}_j)=c_j^p`$ and $`N_E^F(\stackrel{~}{c}_u)=c_u`$, $`uW_j`$. Therefore, it follows from (4.1) (iii), , Sect. 14.7, Lemma a, and the equality $`o(B^{})=o(B)/p`$ that $`M/F`$, $`b^{}`$, and $`\stackrel{~}{E}_\rho `$, $`\stackrel{~}{\tau }_\rho ,\stackrel{~}{\sigma }_\rho ,\stackrel{~}{c}_\rho `$, where $`\rho `$ runs through $`W_j`$ or $`\{1,\mathrm{},t\}`$ depending on whether or not $`F=E_j`$, are related as in Lemma 6.3. Since the RC-formula and , Sect. 15.1, Proposition b, yield Cor$`{}_{F/E}{}^{}([(E_j/F,\tau _j^p,c_j)])=[(E_j/E,\tau _j,c_j^p)]`$, this proves Lemma 6.4 (ii) when $`[F:E]=p`$.
Let finally $`F`$ be any proper extension of $`E`$ in $`M`$ different from $`M`$. By Galois theory and the structure of finite abelian groups, then $`M`$ possesses subfields $`F_0`$ and $`M_0`$, such that $`EF_0FM_0`$ and $`[F_0:E]=[M:M_0]=p`$. In view of the transitivity of norm mappings, canonical projections of Galois groups and scalar extension maps of Brauer groups in towers of intermediate fields of $`M/E`$ (cf. , Sect. 9.4, Corollary a), the considered special cases of Lemma 6.4 enable one to complete inductively its proof. ∎
## 7. Existence and form of Hasse symbols
The main purpose of this Section is to prove Theorem 3.2. Fix an $`𝔽_p`$-basis $`I_p`$ of $`{}_{p}{}^{}\mathrm{Br}(E)`$ and a generator $`\phi _{\mathrm{}}`$ of $`𝒢(E_{\mathrm{}}/E)`$ as a topological group. Take a subset $`\mathrm{\Lambda }_p(E)=\{b_{i,n}(p):iI_p,n\}`$ of Br$`(E)_p`$ so that $`\{b_{i,1}(p):iI_p\}`$ is a basis of $`{}_{p}{}^{}\mathrm{Br}(E)`$ as an $`𝔽_p`$-vector space, and $`pb_{i,n}(p)=b_{i,n1}(p):iI_p,n`$ and $`n2`$ (the existence of $`\mathrm{\Lambda }_p(E)`$ is implied by Proposition 2.2 (ii)). For each $`E^{}\mathrm{\Omega }_p(E)`$, let $`\rho _{E/E^{}}(b_{i,n}(p))=b_{i,n}(E^{}):(i,n)I_p\times `$, $`E_{\mathrm{}}^{}=E_{\mathrm{}}E^{}`$, and $`[(E^{}E_{\mathrm{}}):E]=m(E^{})`$. As $`𝒢(E_{\mathrm{}}/E)_pH`$ whenever $`H`$ is an open subgroup of $`_p`$ (see , Ch. I, 1.5, 4.2), it follows from Galois theory that $`E_{\mathrm{}}^{}/E^{}`$ is a $`_p`$-extension and $`\phi _{\mathrm{}}^{m(E^{})}`$ uniquely extends to an $`E^{}`$-automorphism $`\phi _{\mathrm{}}(E^{})`$ of $`E_{\mathrm{}}^{}`$; one also sees that $`\phi _{\mathrm{}}(E^{})`$ topologically generates $`𝒢(E_{\mathrm{}}^{}/E^{})`$. For each $`n`$ and $`iI_p`$, denote by $`\phi _n(E^{})`$ the automorphism of $`\mathrm{\Gamma }_nE^{}`$ induced by $`\phi _{\mathrm{}}(E^{})`$, and by $`g_{i,n}(E^{})`$ the element of $`𝒢(\mathrm{\Gamma }_nE^{}/E^{})^{d(p)}`$ with components $`g_{i,n}(E^{})_i=\phi _n(E^{})`$, and $`g_{i,n}(E^{})_i^{}=1:i^{}I_p\{i\}`$. By Proposition 2.2 and , Sect. 15.1, Proposition a, $`E^{}`$ has a subset $`C_p(E^{})=\{c_{i,n}(E^{}):iI_p,n\}`$, such that $`[((\mathrm{\Gamma }_nE^{})/E^{},\phi _n(E^{}),c_{i,n}(E^{}))]=b_{i,n}(E^{})`$, for each $`(i,n)I_p\times `$. Observe that, for any $`n`$, there is a unique surjective homomorphism $`(,(\mathrm{\Gamma }_nE^{})/E^{}):E^{}𝒢(\mathrm{\Gamma }_nE^{}/E^{})^{d(p)}`$, whose kernel is $`N(\mathrm{\Gamma }_nE^{}/E^{})`$, and which maps $`c_{i,n}(E^{})`$ into $`g_{i,n}(E^{})`$, when $`iI_p`$. Fix some $`M^{}\mathrm{\Omega }_p(E^{})`$ and $`\mu `$ so that $`p^\mu [M^{}:E]`$, and for each $`iI_p`$, let $`B(M^{}/E^{})_i=\{\beta _iE^{}:[(L^{}/E^{},\sigma ^{},\beta _i)]b_{i,\mu }(E^{})\}`$, for every cyclic extension $`L^{}`$ of $`E^{}`$ in $`M^{}`$, where $`\sigma ^{}`$ is a generator of $`𝒢(L^{}/E^{})`$. It follows from Proposition 2.2 (ii)-(iii), the definition of $`\mathrm{\Lambda }_p(E)`$ and Lemma 6.3 that the groups $`\overline{B}(M^{}/E^{})_i:=B(M^{}/E^{})_i/N(M^{}/E^{})`$, $`iI_p`$, have the following property:
(7.1) $`\overline{B}(M^{}/E^{})𝒢(M^{}/E^{})`$, for each index $`i`$, and the inner product of $`\overline{B}(M^{}/E^{})_i:iI_p`$, is direct and coincides with $`E^{}/N(M^{}/E^{})`$.
Therefore, there exists a unique homomorphism $`(,M^{}/E^{})`$ of $`E^{}`$ into
$`𝒢(M^{}/E^{})^{d(p)}`$, mapping $`B(M^{}/E^{})_i`$ into the $`i`$-th component of $`𝒢(M^{}/E^{})^{d(p)}`$ by the formula $`\beta _i\omega _{M^{}/E^{},b_{i,\mu }(E^{})}(\beta _i)`$, for each $`iI_p`$, where $`\omega _{M^{}/E^{},b_{i,\mu }(E^{})}`$ is defined as in Lemma 6.3. Hence, $`\mathrm{Ker}(,M^{}/E^{})=N(M^{}/E^{})`$, and by Lemma 6.4, the sets $`H(E^{})=\{(,M^{}/E^{}):M^{}\mathrm{\Omega }_p(E^{})\}`$, $`E^{}\mathrm{\Omega }_p(E)`$, consist of surjections related as required by Theorem 3.2 (ii)$`÷`$(iii).
Suppose now that $`\mathrm{\Theta }_p(E^{})=\{\theta (M^{}/E^{}):E^{}𝒢(M^{}/E^{})^{d(p)},M^{}\mathrm{\Omega }_p(E^{})\}`$, $`E^{}\mathrm{\Omega }_p(E)`$, is a system of surjective homomorphisms with the same kernels and relations, and such that $`\theta (\mathrm{\Gamma }_n/E)=(,\mathrm{\Gamma }_n/E)`$, for every $`n`$. Then it follows from Proposition 2.2, the RC-formula and (4.1) (iii) that $`\theta (\mathrm{\Gamma }_nE^{}/E^{})=(,(\mathrm{\Gamma }_nE^{})/E^{})`$, for each pair $`(E^{},n)\mathrm{\Omega }_p(E)\times `$. We show that $`\mathrm{\Theta }_p(E^{})=H_p(E^{})`$, for any $`E^{}\mathrm{\Omega }_p(E)`$. This is obvious, if $`E_{\mathrm{}}=E(p)`$, so we assume further that $`E(p)E_{\mathrm{}}`$. In view of Proposition 2.2 (i) and the already established part of Theorem 3.2, it is sufficient to prove that $`\theta (M/E)=(,M/E)`$, for an arbitrary fixed field $`M\mathrm{\Omega }_p(E)`$. Note first that the compositum $`ME_{\mathrm{}}`$ possesses a subfield $`M_0\mathrm{\Omega }_p(E)`$, such that $`M_0E_{\mathrm{}}=E`$ and $`M_0E_{\mathrm{}}=ME_{\mathrm{}}`$. This follows from Galois theory and the projectivity of $`_p`$ as a profinite group (cf. , Ch. I, 5.9). In particular, $`MM_0\mathrm{\Gamma }_n`$, for every sufficiently large index $`n`$. We also have $`\theta (M/E)=\pi _{M_0\mathrm{\Gamma }_n/M}\theta (M_0\mathrm{\Gamma }_n/E)`$, which allows us to consider only the special case of $`M=M_0\mathrm{\Gamma }_\kappa `$ and $`M_0E`$, where $`\kappa `$ is chosen so that $`[M_0:E]p^\kappa `$. Let now $`t`$ be the rank of $`𝒢(M/E)`$, and $`E_1,\mathrm{},E_t`$ be cyclic extensions of $`E`$ in $`M`$, such that $`_{u=1}^t[E_u:E]=[M:E]`$, $`E_1=\mathrm{\Gamma }_\kappa `$ and $`E_2\mathrm{}E_t=M_0`$. Take $`E_1^{},\mathrm{},E_t^{}`$ as in Lemma 6.3, denote for convenience by $`\tau _1`$ the automorphism $`\phi _\kappa (E)`$ of $`E_1`$, and let $`\tau _u`$ be a generator of $`𝒢(E_u/E)`$, for every $`u\{2,\mathrm{},t\}`$. Fix an index $`xI_p`$, put $`b_{x,\kappa }(p)=b_x`$, and identifying $`𝒢(M/E)`$ with the $`x`$-th component of $`𝒢(M/E)^{d(p)}`$, consider elements $`x_1,\mathrm{},x_t`$ of $`E^{}`$ determined so that $`\theta (M/E)(x_u)`$ equals the $`E_u^{}`$-automorphism $`\sigma _u`$ of $`M`$ extending $`\tau _u`$, for each positive integer $`ut`$. The assumptions on $`\mathrm{\Theta }_p(E^{}):E^{}\mathrm{\Omega }_p(E)`$, imply that $`x_uN(E_u^{}/E):u=1,\mathrm{},t`$, and $`[(E_1/E,\tau _1,x_1)]=b_x`$. We show as in the proof of Lemma 6.3 (iii) that $`[(E_u/E,\tau _u,x_u)]=(p^\kappa /[E_u:E])b_x`$, for each $`u`$. Fix an index $`u2`$, put $`\omega _u=p^\kappa /[E_u:E]`$, $`\sigma _1^{}=\sigma _1\sigma _u`$, $`x_1^{}=x_1x_u`$, and denote by $`F_u`$ the fixed field of the subgroup of $`𝒢(M/E)`$ generated by $`\sigma _1^{}`$ and $`\sigma _u^{}:u^{}\{1,u\}`$. It is easily verified that $`F_1\mathrm{}F_t=M`$ and $`_{u^{}=1}^t[F_u^{}:E]=[M:E]`$, where $`F_u^{}=E_u^{}:u^{}u`$. As $`\mathrm{Ker}(\theta (F_u/E))=N(F_u/E)`$, $`\theta (M/E)(x_1^{})=\sigma _1^{}`$ and $`\sigma _1^{}𝒢(M/F_u)`$, the equality $`\theta (F/E)=\pi _{M/E}\theta (M/E)`$ ensures that $`x_1^{}N(F_u/E)`$. Observe also that $`E_1E_u=E_1F_u`$ and $`E_uF_u=E`$. This implies that $`(E_1/E,\tau _1,x_1)_E(E_u/E,\tau _u,x_u)`$ is $`E`$-isomorphic to $`(E_1/E,\tau _1^{},x_1x_u^{\omega _u})_E(F_u/E,f_u,x_u)`$, where $`\tau _1^{}𝒢(E_1/E)`$ and $`f_u𝒢(F_u/E)`$ are induced by $`\sigma _1^{}`$ and $`\sigma _u`$, respectively. Since $`\sigma _u𝒢(M/E_1)`$ and $`x_uN(E_1/E)`$, it thereby turns out that $`(E_u/E,\tau _u,x_u)(F_u/E,f_u,x_u)(F_u/E,f_u^1,x_1)(F_u/E,f_1,x_1)`$ over $`E`$, $`f_1`$ being the automorphism of $`F_u`$ induced by $`\sigma _1`$. Furthermore, if $`\mathrm{\Phi }_u`$ is the extension of $`E`$ in $`E_1`$ of degree $`[E_u:E]`$, then $`E_u\mathrm{\Phi }_u=E_uF_u`$ and, because $`x_1N(E_u/E)`$, it follows that $`[(E_u/E,\tau _u,x_u)]=[(F_u/E,f_1,x_1)]=\omega _u[(E_1/E,\tau _1,x_1)]=\omega _ub_x`$, as claimed. The obtained result indicates that $`\theta (M/E)(\alpha _x)=(\alpha _x,M/E)`$, for each $`\alpha _xB(M/E)_x`$. As $`x`$ is an arbitrary element of $`I_p`$, this enables one to complete the proof of Theorem 3.2 by applying (7.1).
###### Remark 7.1.
Theorem 1.1 and , Theorem 2.1, show that if $`(E,v)`$ is a Henselian discrete valued strictly PQL-field, then $`\widehat{E}(p)/\widehat{E}`$ is a $`_p`$-extension, for each $`pP(E)`$. Therefore, one can take as $`E_{\mathrm{}}`$ the compositum of all $`I\mathrm{\Omega }(E)`$ that are inertial over $`E`$. Note also that if $`E`$ is SQL, then Br$`(E)`$ is isomorphic to the direct sum $`_{pP(E)}(p^{\mathrm{}})`$ (by (2.1) (i)), and for each $`pP(E)`$, the set $`\mathrm{\Lambda }_p(E)`$ can be chosen so that one may put $`C_p(E)=\{c_n(p)=\pi :n\}`$, where $`\pi `$ is a uniformizer of $`(E,v)`$. When $`E`$ is a local field, and for each $`(p,n)P(E)\times `$, $`\phi _n(p)`$ is the Frobenius automorphism of the inertial extension of $`E`$ in $`E_{\mathrm{sep}}`$ of degree $`p^n`$, the sets $`H_p(E)`$, $`pP(E)`$, from the proof of Theorem 3.2, define the Hasse (the norm residue) symbol, in the sense of , and give rise to the Artin map (cf. , Ch. 6).
###### Corollary 7.2.
In the setting of (2.2) (i), let $`M\mathrm{\Omega }(E)`$ and $`F`$ be a finite extension of $`E`$ in $`E_{\mathrm{sep}}`$. Then $`N(MF/F)=\{\lambda F^{}:N_E^F(\lambda )N(M/E)\}`$.
###### Proof.
It suffices to prove that $`N(MF/F)`$ includes the preimage of $`N(M/E)`$ in $`F^{}`$ under $`N_E^F`$. In view of Theorem 3.1, Lemma 2.1, , I, Lemma 4.2 (ii), and of the PQL-property of $`F`$, one may consider only the case where $`M/F_0`$ is cyclic, for $`F_0=MF`$. Put $`\mu _0=N_{F_0}^F(\mu )`$, for each $`\mu F^{}`$. Theorem 3.2 (iii), , I, Lemma 4.2 (ii), and norm transitivity in towers of intermediate fields of $`MF/E`$ imply that if $`N_E^F(\mu )N(M/E)`$, then $`\mu _0N(M/F_0)`$. At the same time, it follows from Proposition 2.2 (ii), , I, Corollary 8.5, the RC-formula and the assumptions on $`E`$ that Cor$`_{F/F_0}`$ is injective. Therefore, one deduces from the former part of (4.1) (iii) (in the general form pointed out in Remark 4.3) that $`\mu N(MF/F)`$, which proves our assertion. ∎
Corollary 7.2 generalizes , Theorem 7.6. Applying the RC-formula, Propositions 2.2 (ii) and 2.3 (i), as well as Lemma 2.1, norm and corestriction transitivity, and already used known relations between norms and cyclic algebras, one obtains by the method of proving Theorem 3.2 the existence of an exact analogue to the local Hasse symbol in the following situation:
###### Corollary 7.3.
Assume that $`E`$ is a nonreal field, such that every $`L\mathrm{\Omega }(E)`$ is strictly PQL with $`\rho _{E/L}`$ surjective. Then the maps Cor$`{}_{L/E}{}^{}:L\mathrm{\Omega }(E)`$, are bijective, and there are sets $`H(E^{})=\{(,M^{}/E^{}):`$ $`E^{}𝒢(M^{}/E^{})^{\mathrm{Br}(E^{})},`$ $`M^{}\mathrm{\Omega }(E^{})\}`$, $`E^{}\mathrm{\Omega }(E)`$, of group homomorphisms satisfying the following:
(i) $`(,M^{}/E^{})`$ is surjective and its kernel equals $`N(M^{}/E^{})`$, for each $`E^{}\mathrm{\Omega }(E)`$, $`M^{}\mathrm{\Omega }(E^{})`$;
(ii) $`H_E^{}`$ has the properties required by Theorem 1.3 (ii), for every $`E^{}\mathrm{\Omega }(E)`$; furthermore, if $`M\mathrm{\Omega }(E)`$ and $`K`$ is an intermediate field of $`M/E`$, then $`(\lambda ,M/K)=(N_E^K(\lambda ),M/E)`$, for any $`\lambda K^{}`$;
(iii) The sets $`H(E^{}),E^{}\mathrm{\Omega }(E)`$, are determined by the mappings $`(,\mathrm{\Gamma }/E)`$, when $`\mathrm{\Gamma }`$ ranges over finite extensions of $`E`$ in $`E_{\mathrm{}}`$ of primary degrees.
Corollary 7.3 has a partial analogue for a formally real strictly PQL-field $`E`$ and the sets $`\mathrm{\Omega }^{}(E^{})=\{M^{}\mathrm{\Omega }(E^{}):2[M^{}:E^{}]\}`$, $`E^{}\mathrm{\Omega }^{}(E)`$. Specifically, statements (i)$`÷`$(iii) hold when every $`E^{}\mathrm{\Omega }^{}(E)`$ is $`p`$-quasilocal with Br$`(E^{})_p`$ included in the image of $`\rho _{E/E^{}}`$, for each $`p\{2\}`$; also, in this case, Cor$`_{E^{}/E}`$ induces isomorphisms Br$`(E^{})_p\mathrm{Br}(E)_p`$, $`p>2`$. This is proved in the same way as Corollary 7.3. On the other hand, it follows from the Artin-Schreier theory that if $`E^{}\mathrm{\Omega }^{}(E)`$ and $`E^{}E`$, then $`E^{}`$ is formally real and the action of $`𝒢(E^{}/E)`$ induces $`[E^{}:E]`$ orderings on $`E^{}`$. Therefore, $`E^{}`$ is not $`2`$-quasilocal, so it cannot admit LCFT.
###### Remark 7.4.
Under the hypotheses of Theorem 1.2, let $`E_{\mathrm{ab}}`$ be the compositum of all $`M\mathrm{\Omega }(E)`$, $`E_{\mathrm{ab}}(p)=E_{\mathrm{ab}}E(p)`$, for each $`pP(E)`$, and $`N_1(E)=_{M\mathrm{\Omega }(E)}N(M/E)`$. Suppose that $`d(p)`$, for all $`pP(E)`$, and Pr$`{}_{c}{}^{}(E^{})`$ is the profinite completion of $`E^{}/N_1(E)`$ (concerning its existence, see , Sect. 1.2). Using Theorems 1.1, 1.3 and Galois theory, and arguing as in the proof of implication (iii)$``$(i) of , Proposition 1.14, one obtains that $`\mathrm{Pr}_c(E^{})`$ is isomorphic to the topological group product $`_{pP(E)}𝒢(E_{\mathrm{ab}}(p)/E)^{d(p)}`$, and so generalizes a part of , Proposition 6.3.
Note finally that every abelian torsion group $`T`$ admissible by Proposition 2.2 (ii) is isomorphic to Br$`(E(T))`$, for some strictly PQL-field $`E(T)`$ satisfying the conditions of Corollary 7.3 or its analogue in the formally real case. This follows from , Corollary 6.6 and (7.2) (when $`T`$ is divisible - from Remark 2.4 as well). Hence, all forms of the local reciprocity law and Hasse symbol admissible by Theorems 1.2, 1.3 and Corollary 7.3 can be realized.
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# 1 Introduction
## 1 Introduction
It would be very powerful to use the gravity/gauge correspondence in getting non-perturbative insights of various gauge theories. This idea has been extensively studied in 5d gauged supergravity - which can be formulated as the superstring theory compactified on Ad$`S_5\times S^5`$. This theory is useful since it contains various scalar fields which correspond to the relevant and marginal operators of CFT.
The advantage of the 5d supergravity is that it can be reconciled with the brane-world proposed by Randall and Sundrum (RS) . In this brane-world, it is possible to consider 4d gravity since the graviton is localized on the brane. In order to obtain more realistic braneworld, however, it is desirable to set the braneworld in a background which corresponds, in a holographic sense, to the confining gauge theory with flavor quarks as in QCD. In such a brane-world, we expect to find mesons, the quark-bound states, which are localized on the RS brane as well as the graviton.
Recently, a new idea to introduce flavor quarks has been proposed by Karch and Katz by embedding a D7-brane as a probe which is wrapping a contractible $`S^3`$ of $`S^5`$ in the AdS$`{}_{5}{}^{}\times S^5`$ background in order to avoid the RR tadpole problem. At the price of this embedding, a tachyonic mode appears on the D7-brane, but the system is stable since the mass is within the Breitenlohner-Freedman bound of the AdS background. Meanwhile, the brane-world is set in the AdS<sub>5</sub> accompanied with the compact $`S^5`$ as the inner space. In a similar sense, the $`S^3`$ is considered as the inner space of the embedded D7-brane, and the action of the D7-brane can be reduced to 5d by integrating over the coordinates of $`S^3`$. Then we can set a system of supergravity and D7-brane which are dimensionally reduced to five dimensions.
The purpose of the present paper is to study the problem of embedding the reduced D7-brane in the 5d background of an appropriate brane-world according to the procedure given in . And we also study the localization of the flavor mesons, which are given as fluctuation modes of the D7-brane, on the RS brane. The analysis is performed for the simple background studied before in .
In Section 2, we set up our brane-world model. In Section 3, in this brane-world, the dimensionally reduced D7 brane is embedded. In Section 4, the localization of the mesons are studied, and summary is given in the final section.
## 2 Setting of a brane-world
Here we set up our model of a simple brane-world to embed a dimensionally reduced D7-brane. Consider the 10d type IIB supergravity which includes the dilaton ($`\mathrm{\Phi }`$), axion ($`\alpha `$) and five form field strength to obtain $`_5\times S^5`$ space-time with the following form of the Einstein frame metric,
$$ds_{10}^2=\frac{r^2}{L^2}A^2(r)\eta _{\mu \nu }dx^\mu dx^\nu +\frac{L^2}{r^2}\left(dr^2+r^2d\mathrm{\Omega }_5^2\right),$$
(1)
where $`A(r)`$ denotes a deviation of $`_5`$ from the exact Ad$`S_5`$. The $`_5`$ part, $`ds_5^2=\frac{r^2}{L^2}A^2(r)\eta _{\mu \nu }dx^\mu dx^\nu +\frac{L^2}{r^2}dr^2`$, is rewritten by using the new variable $`y`$, defined as $`r=e^{y/L}`$, as follows
$$ds^2=\stackrel{~}{A}^2(y)\eta _{\mu \nu }dx^\mu dx^\nu +dy^2,$$
(2)
where $`\stackrel{~}{A}^2=\frac{r^2}{L^2}A^2(r)`$. This metric is also obtained from the following 5d theory,
$$S_\mathrm{g}=d^4x^\mu 𝑑y\sqrt{g}\left\{\frac{1}{2\kappa _5^2}R\frac{1}{2}(\mathrm{\Phi })^2\frac{1}{2}e^{2\mathrm{\Phi }}(\alpha )^2V\right\},$$
(3)
where $`V=6/(L^2\kappa _5^2)`$.<sup>*</sup><sup>*</sup>* In the following, we set as $`V=3`$ for $`\kappa _5^2=2`$ and $`L=1`$. This action is obtained from 10d supergravity by integrating out the compact $`S^5`$ part. And this is also equivalent to the 5d gauged supergravity written by dilaton and axion only. In this case, the potential $`V`$ in $`S_\mathrm{g}`$ is generally given as
$$V(\varphi )=\frac{v^2}{8}\underset{i}{}\left(\frac{W}{\varphi _i}\right)^2\frac{v^2}{3}W^2,$$
(4)
where $`W`$ represents the superpotential of the theory and $`v`$ denotes the gauge coupling The gauge coupling parameter $`v`$ is fixed from the AdS<sub>5</sub> vacuum by fixing the radius of AdS as a unit length. Here we set as $`v=2`$.. In the present case, $`\varphi _i=\{\mathrm{\Phi },\alpha \}`$ and $`V=3`$.
The supersymmetric solutions of $`S_\mathrm{g}`$ are given by solving the following first order equations ,
$$\varphi _i^{}=\frac{v}{2}\frac{W}{\varphi _i},\frac{\stackrel{~}{A}^{}}{\stackrel{~}{A}}=\frac{v}{3}W,$$
(5)
where $`{}_{}{}^{}=d/dy`$. However, for some non-supersymmetric solutions, there exists $`W`$ which satisfies the above relations (4) and (5). This is the case known as the fake supersymmetry .
The solutions are obtained under the ansatzs, $`\varphi _i=\varphi _i(y)`$. Here we consider both the fake and the true supersymmetric cases. The AdS<sub>5</sub> solution is considered as a simple example for the supersymmetric case,
$$A=1,\mathrm{\Phi }=\mathrm{\Phi }_0,\alpha =0,$$
(6)
where $`\mathrm{\Phi }_0`$ is a constant and $`W=3/2.`$ As for the fake-supersymmetric case, we adopt the following solution ,
$$A=\left(1(\frac{r_0}{r})^8\right)^{1/4},e^\mathrm{\Phi }=\left(\frac{1+(r_0/r)^4}{1(r_0/r)^4}\right)^{\sqrt{3/2}},\alpha =0,$$
(7)
where $`r_0`$ is a constant and
$$W=\frac{3}{2}\mathrm{cosh}(\sqrt{\frac{8}{3}}\mathrm{\Phi }).$$
(8)
The solutions of this 5d theory can be easily lifted up to 10d. However, when the other scalars are added and they take some non-trivial configurations, it is non-trivial to lift up the solutions to 10d, and $`S^5`$ has been deformed to a complicated geometry . In this case, it is not simple to embed the D7 brane in this 10d background.
The brane-world solution is obtained by solving the following system,
$$S=S_\mathrm{g}+S_\mathrm{b},$$
(9)
where $`S_\mathrm{b}`$ is the action of the RS-brane. Here, we set $`S_\mathrm{b}`$ as follows,
$$S_b=vd^4x𝑑y\sqrt{g}W(\varphi _i)\delta (yy_b).$$
(10)
where $`y_b`$ denotes the brane position. In solving the equations of motion of $`S_g+S_b`$, the role of the brane action is to give the boundary conditions for the bulk solution of $`S_g`$. Meanwhile, these conditions for $`\mathrm{\Phi }`$ and $`A`$ are equivalent to the above equations (5) evaluated just on the point $`y=y_b`$. Then, at any point of $`y_b`$, all the solutions of (5) satisfy the boundary conditions given by $`S_b`$ of the form of (10) . So it is enough to solve equations (5) to obtain the brane-world solution which satisfies the boundary conditions at $`y_b`$. In other words, it is possible to put the brane at any point on the $`y`$axis. This situation is not changed even if the ”superpotential” $`W`$ is given by (8) with the fake-supersymmetic solution given above.
## 3 Reduced D7-brane and its embedding
In , we have given an explicit embedding of D7-brane as a probe in the 10d background (1). Rewrite the metric (1) as
$$ds_{10}^2=\frac{r^2}{L^2}A^2(r)\eta _{\mu \nu }dx^\mu dx^\nu +\frac{L^2}{r^2}\left(dr^2+r^2d\mathrm{\Omega }_5^2\right)$$
$$=\frac{r^2}{L^2}A^2(r)\eta _{\mu \nu }dx^\mu dx^\nu +\frac{L^2}{r^2}\left(\underset{m=4}{\overset{7}{}}(dX^m)^2+\underset{i=8}{\overset{9}{}}(dX^i)^2\right),$$
(11)
and the four dimensional coordinates $`\{X^m\}`$ are further rewritten in the polar coordinate as
$$\underset{m=4}{\overset{7}{}}(dX^m)^2=d\rho ^2+\rho ^2d\mathrm{\Omega }_3^2,$$
(12)
where $`\rho ^2=_{m=4}^7(X^m)^2`$. Then, the world-volume of the D7-brane is taken as $`\{x^\mu ,\rho ,\mathrm{\Omega }_3\}`$. As mentioned above, the D7-brane is wrapping on the contractible $`S^3`$, and we notive the relation $`r^2=\rho ^2+(X^8)^2+(X^9)^2`$ in this case.
The remaining part of the embedding procedure is to determine the shape of the D7-brane in its outer two dimensional space, $`(X^7,X^8)`$. This is performed by solving the equations of motion of two scalar fields in the D7-brane action, which is written as
$$S_{\mathrm{D7}}=\tau _7d^8\xi \left(e^\mathrm{\Phi }\sqrt{𝒢}+\frac{1}{8!}ϵ^{i_1\mathrm{}i_8}A_{i_1\mathrm{}i_8}\right),$$
(13)
where $`𝒢=\mathrm{det}(𝒢_{i,j})`$, $`i,j=07`$. $`𝒢_{i,j}`$ denotes the induced metric and $`\tau _7`$ is the tension of D7 brane. The eight form potential $`A_{i_1\mathrm{}i_8}`$, which is Hodge dual to the axion, couples to the D7 brane minimally. It is obtained as $`F_{(9)}=dA_{(8)}`$ in terms of the Hodge dual field strength $`F_{(9)}`$ . We take the gauge, $`\xi ^i=\{x^\mu ,\rho ,\theta ^j\}`$, where $`\theta ^j(j=13)`$ denote the angle variables of $`S^3`$, and we make an ansatz for the two scalar fields $`X^8(\xi )`$ and $`X^9(\xi )`$ as $`X^9w(\rho )`$ and $`X^8=0`$, without loss of generality. Then, the induced metric and the explicit form of the action (13) are obtained as
$$ds_8^2=e^{\mathrm{\Phi }/2}(\frac{r^2}{L^2}A(r)^2\eta _{\mu \nu }dx^\mu dx^\nu +\frac{1+(w^{})^2}{r^2}L^2d\rho ^2+\frac{\rho ^2}{r^2}L^2d\mathrm{\Omega }_3^2.)$$
(14)
$$S_{\mathrm{D7}}=\tau _7d^8\xi \sqrt{ϵ_3}\left(L_{D7}^{\mathrm{cl}}+L_{D7}^{\mathrm{fluc}}\right),$$
(15)
$$L_{D7}^{\mathrm{cl}}=\rho ^3\left(e^\mathrm{\Phi }A(r)^4\sqrt{1+(w^{})^2}+C_8(r)\right)$$
(16)
where $`C_8(r)`$ denotes the contribution of the eight form potential, and the fluctuations of D7 around its classical configuration are included in $`L_{D7}^{\mathrm{fluc}}`$.
Here $`w(\rho )`$ is solved independently of the angle variables $`\theta ^i`$ on $`S^3`$, so the form of the solution $`w`$ is preserved even if the action was reduced to 5d by integrating over $`\theta ^i`$. As a result, we can reduce the above action as
$$S_{\mathrm{D7}}^{(5d)}=2\pi ^2\tau _7d^4x^\mu 𝑑\rho \left(L_{D7}^{\mathrm{cl}}+L_{D7}^{\mathrm{fluc}}\right).$$
(17)
Thus we obtain the D7-brane action which is dimensionally reduced to the 5d. Here we notice the relation, $`r^2=\rho ^2+w^2(\rho )`$, for $`r`$ and $`\rho `$. We preserve the interpretation of $`r`$ or $`\rho `$ as the energy scale of the field theory on the brane, so we select the solution of $`w(\rho )`$ such that it could give the one to one correspondence between $`r`$ and $`\rho `$ .
We are now ready to embed the reduced D7-brane in the 5d brane-world. As done in 10d, it is embedded as a probe such that the given background configuration is not altered. This is performed by solving the equation of motion for $`w`$, which is obtained from $`S_{D7}`$ only in the brane-world background. However we impose $`Z_2`$ symmetry with respect to the coordinate $`y`$ as,
$$w(y)=w(|yy_b|).$$
(18)
Then the equation of $`w`$ is obtained as
$$\stackrel{~}{\tau }_7\left(\delta _wL_{D7}^{\mathrm{cl}}F(\rho )\frac{2w^{}G(\rho )}{(1+w^{}{}_{}{}^{2})^{3/2}}\delta (\rho \rho _b)\right)=0,$$
(19)
where $`w^{}=dw/d\rho `$, $`\delta _wL_{D7}^{\mathrm{cl}}=_wL_{D7}^{\mathrm{cl}}_\rho \left(_w^{}L_{D7}^{\mathrm{cl}}\right)`$, $`F(\rho )=\rho ^3e^\mathrm{\Phi }A^4`$ and $`G(\rho )=(\rho +ww^{})^2/\rho ^2`$. And $`\rho _b`$ denotes the position of the brane on the coordinate $`\rho `$, defined as $`e^{2y_b/L}(=r_b^2)=\rho _b^2+w^2(\rho _b)`$. The $`\delta `$-function of Eq.(19) comes from the $`Z_2`$ symmetry at $`y=y_b`$, and the factor $`G(\rho )`$ appears when the $`\rho `$-derivative is changed to the $`y`$-derivative. For $`\rho <\rho _b`$, the field equation for $`w`$, $`\delta _wL_{D7}^{\mathrm{cl}}=0`$, is written as
$$w^{\prime \prime }+(1+w^{}{}_{}{}^{2})[\frac{3}{\rho }w^{}+2K_{(1)}(\rho w^{}w)]=0$$
(20)
where $`K_{(1)}=_{r^2}\mathrm{log}(e^\mathrm{\Phi }A^4)`$, and the boundary condition,
$$w^{}(\rho _b)=0.$$
(21)
For the super-symmetric background (6), $`K_{(1)}=0`$. Then we find the solution $`w^{}(\rho )=0(w=`$const.) which preserves the super-symmetry of the system , and the condition (21) is satisfied at any $`\rho `$. While, for the non-supersymmetric background (7), $`K_{(1)}0`$, then $`w^{}(\rho )0`$ and the boundary condition is not satisfied except at $`\rho =0`$ and $`\rho =\mathrm{}`$ since $`w(\rho )`$ decreases monotonically with $`\rho `$. For finite $`\rho _b`$, then, the condition (21) selects the super-symmetric solution. And the RS brane in the non-supersymmetric background is pushed toward the boundary, $`\rho _b=\mathrm{}`$.
As another possible boundary condition, we may consider $`G(\rho _b)=0`$, and we get
$$w(\rho _b)w^{}(\rho _b)+\rho _b=0.$$
(22)
This is equivalent to $`_\rho r(\rho _b)=0`$. Meanwhile $`r(\rho )`$ must increase monotonically with $`\rho `$ since it should be a single valued function of $`\rho `$ due to the requirement of the one to one correspondence between $`r^2(=\rho ^2+w(\rho )^2)`$ and $`\rho `$. Thus, the points which satisfy the condition, $`_\rho r(\rho _b)=0`$, are again restricted to $`\rho _b=0`$ and $`\rho _b=\mathrm{}`$. Then this condition is not useful as the boundary condition. By the same reason, $`F(\rho _b)=0`$ is also rejected as a useful boundary condition. After all, we find that only the supersymmetric solution satisfies the meaningful boundary condition (21).
Here we remember the asymptotic form of $`w`$ at large $`\rho `$,
$$w(\rho )=m_q+c/\rho ^2,c=\overline{\mathrm{\Psi }}\mathrm{\Psi }$$
(23)
and the meaning of the parameters in it. Namely $`m_q`$ represents the quark-mass and $`c`$ is the vacuum expectation value of bi-linear operators of quark fields ($`\mathrm{\Psi }`$) in the dual gauge theory. Then we can see $`\overline{\mathrm{\Psi }}\mathrm{\Psi }=0`$ and $`w=m_q`$ in the supersymmetric case. As a result, the chiral symmetry is preserved.
In the above analysis, we have ignored the interaction between D7 and RS brane since D7 should be treated as a probe. However, the situation is not changed even if we solve the equation of $`w`$ by taking into account of the brane action as, $`S_{D7}+S_b`$, since $`S_b`$ does not depend on $`w`$. In other words, the shape of the D7-brane is independent of the RS brane. So we need the back-reaction from the bulk in order to change the situation of the embedding. But this is out of the present work.
## 4 Meson localization
In the next, we study the fluctuation-modes of the embedded D7 brane. Some of them are trapped on the RS brane, and they might be observed as mesons in our 4d world. Then, the trapped modes are defined as the normalizable one for the integration over $`\rho `$.
### 4.1 For finite $`\rho _b`$
In this case, only the supersymmetric embeddings are allowed. So the background is given by (6), and the quadratic parts of the fluctuations, $`\varphi _8=X^8,\varphi _9=X^9w`$ and vector, are written as
$$L_{D7}^{(2)}=\rho ^3\left\{\frac{1}{2}e^\mathrm{\Phi }\frac{L^2}{r^2}\underset{i=8,9}{}(\varphi _i)^2+F_{ab}F^{ab}+\mathrm{}\right\}$$
(24)
where $`w`$ is a constant, $`w(\rho )=w(\mathrm{})m(=m_q/2\pi )`$. The dots denote the higher order terms which can be neglected here.
For the scalar, we obtain the same field equation for $`\varphi ^8`$ and $`\varphi ^9`$, then they are denoted by $`\varphi `$ for the simplicity. The field-equation is written as
$$_\rho ^2\varphi +\frac{3}{\rho }_\rho \varphi =\frac{M^2R^4}{(m^2+\rho ^2)^2}\varphi ,$$
(25)
where $`M^2`$ is defined as $`\eta ^{\mu \nu }_\mu _\nu \varphi =M^2\varphi `$. In the present case, $`\varphi `$ is $`Z_2`$ symmetric at the brane position $`y_b`$, so we must solve the equation (25) by imposing the following boundary condition,
$$\varphi ^{}(\rho _b)=0.$$
(26)
Before studying the trapping in the brane-world by using the boundary condition (26), we consider the normalizable modes which could be observed at the boundary $`\rho _b=\mathrm{}`$. This analysis is useful for finding the localized modes on the brane since the mass eigen-values of the localized modes on the branes at different $`\rho _b`$ should be related to each other by the smooth functions of $`\rho _b`$ in all the region $`0<\rho _b<\mathrm{}`$. In this sense, the spectra, which should be observed on the boundary, are considered as the one at the starting point.
The Eq.(25) is solved as
$$\varphi =(\rho ^2+m^2)^\alpha \left(c_1F(\alpha ,\alpha +1,2;\rho ^2/m^2)+c_2\rho ^2F(\alpha 1,\alpha ,0;\rho ^2/m^2)\right),$$
(27)
where $`c_1`$ and $`c_2`$ are arbitrary constants, and $`\alpha =(1+\sqrt{1+M^2L^4/m^2})/2`$. In order to find the trapped modes, consider the normalizability condition for $`\varphi `$,
$$_0^{\rho _b}𝑑\rho \rho ^3\left(\frac{L^2}{\rho ^2+m^2}\right)^2\varphi ^2(\rho )<\mathrm{}.$$
(28)
We estimate this condition for (i) $`m0`$ and (ii) $`m=0`$ separately. For the case of (i), the above integral near $`\rho 0`$ is approximately estimated as
$$_0𝑑\rho \rho ^3\varphi ^2(\rho )<\mathrm{}.$$
(29)
Meanwhile, the solution (27) is expanded near $`\rho 0`$ as,
$$\varphi =c_1(1+c_1^1\rho ^2+\mathrm{})+c_2\rho ^2(1+c_2^1\rho ^2+\mathrm{}),$$
(30)
where $`c_1^1,c_2^1`$ are the calculable coefficients. Then we find that the solution of $`c_2=0`$ satisfies (29). However, for the case of (ii), the condition (29) is replaced by
$$_0\frac{d\rho }{\rho }\varphi ^2(\rho )<\mathrm{}.$$
(31)
Then we must take as $`c_1=c_2=0`$, namely $`\varphi =0`$, to satisfy (31). In other words, there is no localized state in the case of $`m=0`$ for the supersymmetric case. So the interesting case is restricted to the case of massive quark.
On the other hand, at large $`\rho `$, the normalizability condition (28) is approximated for any $`m`$ as,
$$^{\mathrm{}}\frac{d\rho }{\rho }\varphi ^2(\rho )<\mathrm{}.$$
(32)
And the solution of $`c_2=0`$, the first term of (27), is expanded at large $`\rho `$ as,
$$\varphi /c_1=b_0(\alpha )(1+b_1(\alpha )/\rho ^2+\mathrm{})+d_0(\alpha )\rho ^2(1+d_1(\alpha )/\rho ^2+\mathrm{}),$$
(33)
where the coefficients $`b_0(\alpha )`$ and $`d_0(\alpha )`$ are the functions of $`\alpha `$. From (32), we demand $`b_0(\alpha )=0`$. As a result, we get $`\alpha =n+1`$ with the integer n , and we thus find infinite series of discrete meson mass.
We now return to the case of the brane-world. In this case, $`\rho _b`$ is finite and we need the boundary condition (26) at $`\rho _b`$ instead of the normalizable condition at large $`\rho `$ (32) given above. In order to satisfy (26), $`b_0(\alpha )/d_0(\alpha )`$ should be fixed to an appropriate value $`f`$, which depends on $`\rho _b`$,
$$b_0(\alpha )/d_0(\alpha )=f(\rho _b),$$
(34)
and $`f`$ should satisfy the condition, $`f(\mathrm{})=0`$. From this equation, we could find the mass eigen-values. Here, we can see this statement by a numerical analysis. For different three values of $`\rho _b`$, the value of $`\varphi ^{}`$ as a function of $`\alpha `$ is shown in the Fig.1. We find that the values of $`\alpha `$ at each zero point of $`\varphi ^{}`$ shift to larger $`\alpha `$ side from the integer point, which are given by $`\alpha =n+1`$ for $`\rho _b=\mathrm{}`$, smoothly when $`\rho _b`$ decreases.
We could conclude as follows. Also on the RS brane, in the present supersymmetric braneworld, we can see the meson mass-spectra, which shift to the massive side when the brane moves to the smaller $`\rho _b`$.
### 4.2 Non-supersymmetric background
For non-supersymmetric background, we must take $`\rho _b=\mathrm{}`$. In this case, the field equations for $`\varphi ^8`$ and $`\varphi ^9`$ are complicated, but the equations for both fields have the same asymptotic solution with the one given for the supersymmetric case by the equation (30) at small $`\rho `$ and (33) at large $`\rho `$ respectively. In the present case, however, the normalizable condition has a different form from (28) and it is given as
$$_0^{\mathrm{}}𝑑\rho \rho ^3e^\mathrm{\Phi }\left(\frac{L^2}{Ar^2}\right)^2\varphi (\rho )^2<\mathrm{}$$
(35)
where $`\varphi `$ denote $`\varphi ^8`$ or $`\varphi ^9`$.
In the small $`\rho `$ region, we find the following approximated condition
$$_0𝑑\rho \rho ^3\varphi (\rho )^2<\mathrm{}.$$
(36)
It should be noticed that this condition is not changed even if $`m=0`$ because, in the limit of $`\rho 0`$, we find
$$r^2=\rho ^2+w^2(\rho )w^2(0)0$$
for the non-supersymmetric case and that $`e^\mathrm{\Phi }`$ and $`A`$ approach to the constant . This is an important point to be discriminated from the supersymmetric case. And (36) is satisfied by the first one of the two independent solutions given in (30). Then we expect to find normalizable modes or the meson spectra on the RS brane for any value of $`m`$.
At large $`\rho `$, (35) is approximated by (32), and the asymptotic solution for $`\varphi `$ is also the same form with (33). Then, we can choose the second series of (33) as the normalizable solution, and many number of normalizable modes are found for both $`\varphi ^8`$ and $`\varphi ^9`$. Especially, for the case of $`m=0`$, we find one mass-less mode for $`\varphi ^8`$ as the Nambu-Goldstone boson due to the spontaneous chiral symmetry breaking. Actually, these points have been shown in when the RS brane is absent. In the present case, from the $`Z_2`$ symmetry, we need new conditions for the modes are $`_\rho \varphi ^8|_{\rho =\mathrm{}}=_\rho \varphi ^9|_{\rho =\mathrm{}}=0`$. And all the normalizable solutions satisfy this. Then it would be straightforward to find the meson-spectra on the RS brane.
## 5 Summary
Here, we consider a class of background solutions, $`_5\times S^5`$, of IIB superstring theory in order to construct a braneworld in which mesons are trapped on the RS brane. The brane-world is set in the uncompactified 5d part $`_5`$ of these solutions. Its effective action can be obtained from the 10d theory by integrating over the inner compact space, $`S^5`$. They are also obtained from the 5d gauged supergravity with the dilaton and the axion. Any solution of this 5d theory can be easily lifted up to the one of the 10d theory on $`_5\times S^5`$. On the same footing, the action of the D7-brane is reduced to 5d by integrating over $`S^3`$, which is regarded as the inner compact space of the D7-brane here.
After setting the stable RS brane, which can be put at any position of the fifth coordinate $`\rho _b`$ we like, for the 5d background given here, the reduced D7 brane is embedded in this background as a probe in the sense that the D7 brane does not change the given 5d configuration. The embedded form is obtained by solving the scalar field equation of the D7 brane action. For finite $`\rho _b`$, only the supersymmetric embedding is allowed and the quark mass is arbitrary in this case. The chiral symmetry of the corresponding gauge theory is then preserved for the mass-less quark. When the quark mass is finite, the series of meson mass-spectra are observed as trapped states on the RS brane.
In the non-supersymmetric case, on the other hand, the brane is retained on the boundary, $`\rho _b=\mathrm{}`$. In this case, the situation is similar to the case of 10d gauge/gravity correspondence. We can then observe the spontaneous chiral symmetry breaking and infinite series of the trapped mesons on the RS brane. And we assure the existence of the Nambu-Goldstone boson, which is generated as a result of the spontaneous chiral symmetry breaking.
In any case, we could show a model of a braneworld which includes flavor quarks and their bound states on the RS brane. It would be an interesting problem to study the brane-world cosmology including hadrons in terms of the model presented here. We will discuss on this point in the near future.
Acknowledgments: The author would like to thank M. Tachibana and M. Yahiro for useful discussion. This work has been supported in part by the Grants-in-Aid for Scientific Research (13135223) of the Ministry of Education, Science, Sports, and Culture of Japan.
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# The First Detailed Abundances for M giants in Baade’s Window from Infrared Spectroscopy
## 1 Introduction
The commissioning of the 4m telescope at Cerro Tololo (now the Blanco telescope) included a new wide field grating/prism developed by Blanco (Blanco, 2001) at the prime focus, permitting large scale surveys yielding slitless spectra at classification resolution. Used in combination with red sensitive photographic plates, the grism was ideally suited to surveys of M and carbon stars; the Magellanic Clouds and Galactic bulge were prime targets. The discovery and spectral classification of hundreds of M giants in the Baade’s Window field (Blanco, McCarthy & Blanco, 1984) identified large numbers of M giants with spectral types later than M4 and as cool as M9. Such cool luminous stars are not found in globular clusters, so this characteristic of the bulge population was immediately understood at the time to contradict the dominant view of the bulge being a metal poor, population II system. Using the identifications of Blanco, McCarthy & Blanco (1984), Frogel & Whitford (1987) obtained infrared photometry from which they obtained effective temperatures and bolometric luminosities (reaching as bright as $`M_{bol}<4.5`$) of the bulge M giant population.
Even at high resolution, M giants are difficult to study using optical abundance analysis due to their severe molecular blanketing. While McWilliam & Rich (1994) and subsequent studies perform high resolution abundance analysis of of the bulge K giants, the composition of the Baade’s Window M giants which includes stars as late as M9, continues to remain as an unsolved problem. Because stars like these are a major contributor to the integrated light of elliptical galaxies, the study of these stars has broader implications.
During the 1980’s, the indications of a wide abundance spread in the bulge from the presence of both RR Lyrae and late M giant stars was confirmed from the study of the K giants (Rich, 1998). Photometry and low resolution spectroscopy of the K/M giant population indicated an average metallcity of \[M/H\]$`+0.3`$(see e.g. Rich, 1998; Terndrup, Frogel, & Whitford, 1990; Terndrup, Frogel & Whitford, 1991) and a larger spread ($``$1 dex) among K–giants compared to that of the M giant population ($``$0.3 dex). Sharples et al. (1990) found the bulge M giants to have greater TiO 8415 band line strength at constant $`JK`$ relative to the Solar vicinity stars; model atmosphere fits found a +0.5 dex \[M/H\] enhancement for the bulge giants. Recently, the use of strong Ca, Na, and CO lines in the K band at medium spectral resolution has been (Frogel et al., 2001) has been applied the bulge M giant population, finding roughly Solar metallicity and no indication of a gradient in \[M/H\] within the inner 500 pc (Ramírez et al., 2000). For these cool stars especially, one would very much like to see abundances measured from high resolution spectra as the gradient in composition remains unconstrained. When the commissioning of NIRSPEC (McLean, 1998) on Keck II made it possible to obtain high dispersion infrared spectroscopy, we decided to obtain spectra of the Baade’s Window M giants. Having established abundances for the M giants in Baade’s Window, our aim is to extend our study to include the population of old late-type giants closer to the Galactic Center.
While all low mass stars ascend the first giant branch, only the most metal rich stars might be expected to become M giants. Are the M giants the progeny of the most metal rich, or perhaps the most alpha-enhanced, bulge stars? Given their high luminosities, are the M giants be a separate, perhaps younger, stellar population? Detailed abundance analysis plays a critical role in answering these questions. Here we report the first abundance analysis from high resolution infrared spectra, for M giants in the Galactic bulge field of Baade’s Window.
## 2 Observations and abundance analysis
The Baade’s Window field is a roughly 45 diameter circle centered on the globular cluster NGC 6522 $`(l,b)=(3.93,1.02)`$ corresponding to $`550`$pc South of the Galactic Center, at 8kpc. This field lies within the photometric structure of the bulge/bar and represents well the bulge stellar population. Mould (1983) measured radial velocities for 50 of the stars, finding a dispersion of $`113\pm 11`$ km/sec, a value consistent with bulge membership and confirmed by Sharples et al. (1990) in their sample of 239 M giants. The coordinates and finding charts for the M giants are given in Blanco, McCarthy & Blanco (1984) and Blanco (1986); see Table 1).
The program stars were observed on 23 June 2001, 13-15 July 2002, and 20-21 July 2003. Typical exposure times for these bright (H=8-9) stars range from 120 to 240 sec. The fainter K giants (I-195 and I-194) required a total of 1800 sec of exposure time, as they are typically 3-4 mag fainter. We nod the star by 5-10 arcsec to take at least two separate exposures. The M giants were frequently observed in conditions of less than ideal seeing and transparency. Four M giants in the Solar vicinity were observed on 27 April 2004. We used NIRSPEC (McLean, 1998) in the echelle mode. A slit width of $`0\stackrel{}{\mathrm{.}}43`$ and a length of 8″ or 12″ (depending on the targets) giving an overall spectral resolution R=25,000, and the standard NIRSPEC-5 setting, which covers most of the 1.5-1.8 $`\mu `$m H-band have been selected.
The raw stellar spectra have been reduced using the REDSPEC IDL-based package written at the UCLA IR Laboratory. Each order has been sky subtracted by using nodding pairs and flat-field corrected. Wavelength calibration has been performed using arc lamps and a 2-nd order polynomial solution, while telluric features have been removed by using a O-star featureless spectrum. The signal to noise ratio of the final spectra is $``$50. Fig. 1 illustrates the spectrum in the 1.555$`\mu `$m region for two M9 giants in our sample.
A grid of suitable synthetic spectra of giant stars has been computed by varying the photospheric parameters and the element abundances, using an updated version of the code described in Origlia, Moorwood & Oliva (1993). By combining full spectral synthesis analysis with equivalent widths measurements of selected lines, we derive abundances for Fe and O, and with somewhat greater errors, C, Mg, Si, and Ti. The lines and analysis method are described in Origlia, Rich & Castro (2002); Origlia & Rich (2004). Here we summarize the major issues.
The code uses the LTE approximation. In the H band, most of the OH and CO molecular lines are not saturated and can be safely treated under the LTE approximation, being roto-vibrational transitions in the ground electronic state, providing accurate C and O abundances (Merrill & Ridgway, 1979; Lambert et al., 1984; Smith et al., 2000). Detailed computations of possible NLTE effects for atomic lines in the H band have been performed only for AlI lines in the Sun (see Baumueller & Gehren (1996), finding indeed negligible corrections. However, most of the near IR atomic lines are of high excitation potential, indicating that they form deep in the atmosphere, where the LTE approximation should hold even in giants of low gravity. Moreover, one of the major mechanisms which can cause a deviation from LTE, namely over-ionization by UV radiation, is less efficient in cool giants, while photon suction can have some relevance. According to NLTE computations on Fe and Mg lines (see e.g. Gratton et al., 1999; Zhao & Gehren, 2000) deviations from LTE (at a level of $``$0.1 dex) are mainly observed in stars which are significantly hotter and more metal poor than those in our program.
The code is based on the molecular blanketed model atmospheres of Johnson, Bernat & Krupp (1980) in the 3000-4000 K temperature range and the ATLAS9 models for temperatures above 4000 K. Since in the near IR the major source of continuum opacity is H<sup>-</sup> with its minimum near 1.6 $`\mu `$m, the dependence of the results on the choice of reasonable model atmospheres should not be critical. However, ss a check, we also computed synthetic spectra using the more updated NextGen model atmospheres by Hauschildt et al. (1999) and we compare them with those obtained using Johnson, Bernat & Krupp (1980) models. As an expample, Fig. 2 shows the results at solar metallicity in the range of temperatures 3200-4000 K. The average ratio between NextGen and Johnson, Bernat & Krupp (1980) spectra is between 0.98 (at T<sub>eff</sub>=3200 K) and 1.0 (at T$`{}_{eff}{}^{}`$3600 K) over the entire H band, the most discrepant values being always within 0.9 and 1.1. These discrepancies are most likely due to differences in the temperature structure of the model atmospheres which however have a minor impact on the overall inferred abundances. To test this, we analyzed the 3600K M giant BMB93 using both atmospheres, finding numerical results for \[Fe/H\] and O, Si, M, Ca, Ti, and C/Fe in agreement to within 0.02 dex. Other systematic errors are of course of much greater concern.
Three main compilations of atomic oscillator strengths are used: the Kurucz database (c.f. http://cfa-www.harward.edu/amdata/ampdata/kurucz23/sekur.html), Bièmont & Grevesse (1973) and Meléndez & Barbuy (1999). On average, log-gf values from Meléndez & Barbuy (1999) are systematically lower than Kurucz ones but in most cases the difference does not exceed 0.2 dex and the overall scatter in the derived abundances is $`<0.1`$ dex (see Origlia, Rich & Castro, 2002, for a more quantitative comparison). The log-gf determinations for the IR are not as good as for the optical and for this (and many other reasons) we obtain spectra of local disk comparison M giants.
Applying this procedure to the IR spectrum of Arcturus (Hinkle, Wallace & Livingston, 1995), rebinned at the NIRSPEC resolution, we find abundances of \[Fe/H\]=$`0.6`$, \[O/H\]=$`0.25`$, \[Mg,Si/H\]=$`0.1`$, \[Ca/H\]=$`0.45`$, \[Ti/H\]=$`0.4`$, \[C/H\]=$`0.8`$, fully consistent (on average, within $`\pm `$0.1 dex) with the values published in the literature (Gonzalez & Wallerstein, 1998; Smith et al., 2000; Zoccali et al., 2004). Reference solar abundances are from Grevesse & Sauval (1998).
In the first iteration, we estimate stellar temperature from the $`(\mathrm{J}\mathrm{K})_0`$ colors (see Table 1) and the color-temperature transformation of Montegriffo et al. (1998) specifically calibrated on globular cluster giants. Gravity has been estimated from theoretical evolutionary tracks, according to the location of the stars on the Red Giant Branch (RGB) (see Origlia et al., 1997, and references therein for a more detailed discussion). For microturbulence velocity an average value $`\xi `$=2.0 km/s has been adopted (see also Origlia et al., 1997). More stringent constraints on the stellar parameters are obtained by the simultaneous spectral fitting of the several CO and OH molecular bands, which are very sensitive to temperature, gravity and microturbulence variations (see Figs. 6,7 of Origlia, Rich & Castro (2002)).
CO and OH in particular, are extremely sensitive to $`T_{eff}`$ in the range between 4500 and 3500 K. Indeed, temperature sets the fraction of molecular versus atomic carbon and oxygen. At temperatures $``$4500 K molecules barely survive, most of the carbon and oxygen are in atomic form and the CO and OH spectral features become very weak. On the contrary, at temperatures $``$3500 K most of the carbon and oxygen are in molecular form, drastically reducing the dependence of the CO and OH band strengths and equivalent widths on the temperature itself (Origlia et al., 1997).
For our analysis we compute a grid of synthetic spectra with abundances and abundance patterns varying over a large range and the photospheric parameters around the photometric estimates. The final values of our best-fit models together with random errors are listed in Table 1.
To illustrate, Fig. 3 shows a region of the observed H band spectrum of BMB93 with (superimposed) our best-fit model and two other models with $`\mathrm{\Delta }[\mathrm{X}/\mathrm{H}]=\pm `$0.2 dex, $`\mathrm{\Delta }\mathrm{T}_{\mathrm{eff}}=\pm `$200 K, $`\mathrm{\Delta }\xi =`$0.5 km s<sup>-1</sup>, and $`\mathrm{\Delta }\mathrm{log}\mathrm{g}=\pm `$0.5 dex, with respect to the best-fit parameters. It is clearly seen that our models with $`\pm `$0.2 dex abundance or $`\pm `$200 K temperature variations give remarkably different molecular line profiles. Microturbulence variation of $`\pm `$0.5 km/s mainly affects the OH lines, while gravity mainly affects the CO lines.
As a further check of the statistical significance of our best-fit solution, we also compute synthetic spectra with $`\mathrm{\Delta }\mathrm{T}_{\mathrm{eff}}=\pm `$200 K, $`\mathrm{\Delta }\mathrm{log}\mathrm{g}=\pm `$0.5 dex and $`\mathrm{\Delta }\xi =`$0.5 km s<sup>-1</sup>, and with corresponding simultaneous variations of the C and O abundances (on average, $`\pm `$0.2 dex) to reproduce the depth of the molecular features. As a figure of merit of the statistical test we adopt the difference between the model and the observed spectrum (hereafter $`\delta `$). In order to quantify systematic discrepancies, this parameter is more powerful than the classical $`\chi ^2`$ test, which is instead equally sensitive to random and systematic errors (see also Origlia et al., 2003; Origlia & Rich, 2004). Our best fit solutions always show $`>`$90% probability to be representative of the observed spectra, while spectral fitting solutions with abundance variations of $`\pm `$0.2 dex, due to possible systematic uncertainties of $`\pm `$200 K in temperature, $`\pm `$0.5 dex in gravity or $``$0.5 km/s in microturbulence have always $``$30% probability to be statistical significant. Hence, as a conservative estimate of the systematic error in the derived best-fit abundances, due to the residual uncertainty in the adopted stellar parameters, one can assume a value of $`\pm 0.1`$ dex. However, it must be noted that since the stellar features under consideration show a similar trend with variations in the stellar parameters, although with different sensitivities, relative abundances are less dependent on the adopted stellar parameters (i.e. on the systematic errors) and their values are well constrained down to $`\pm `$0.1 dex (see also Table 1).
## 3 Results and Discussion
We find no indication of striking iron or alpha abundance differences between our M giant sample and the field and globular cluster bulge K giants previously studied. The observed composition requires no special population to be the progenitors of the M giants. We cannot rule out the possibility that alpha-enhanced stars are more likely to evolve into M giants, but our sample is too small to comment on this possibility.
By observing both the local and bulge M giants with the same instrumentation, we are able to show that the bulge M giants are alpha enhanced, when compared with stars in the Solar vicinity. Earlier work comparing Solar neighborhood and bulge M giants at low resolution would not have been able to sort out the issues connected with composition. The low resolution spectroscopy of Terndrup, Frogel, & Whitford (1990) and Sharples et al. (1990) found that the bulge M giants have stronger TiO bands than their disk counterparts of the same temperature; we propose that alpha enhancement (raising Ti and O abundances) is responsible for this.
Our survey of giant late type stars in Baade’s window find their metallicities to range in \[Fe/H\] between –0.4 and +0.0 dex, (see Fig. 4). Arp (Arp, 1965) I-194 (more likely a K giant), is in common with the McWilliam & Rich (2004) sample: very similar (within $``$0.1 dex) Fe, O, Mg and Si abundances and somewhat lower (by $``$0.3 dex) Ti and higher (by $``$0.3dex) Ca abundances have been obtained in our study. It would be useful to obtain more stars with both optical and infrared abundance analysis, but the infrared approach works best for temperatures cooler than 4000K. Oxygen abundances in the optical are based on the \[OI\] 6300.3Å line, while our oxygen abundances are derived from numerous OH lines, and comparison with the disk stars confirms the oxygen enhancement. As larger samples are developed, it will be important to compare and reconcile the optical and infrared abundance determinations, especially as one pushes toward the Galactic center where only infrared spectroscopy is possible.
The metallicity distribution of the bulge stars in our survey peaks at slightly subsolar metallicity, in agreement with the giant K star distribution (Rich & McWilliam, 2000; McWilliam & Rich, 2004) and the recent, extensive photometric survey by Zoccali et al. (2003). We note the apparent absence of M giants with iron abundances above solar, while their low metallicity spread compared to the significantly larger one measured among K giants by Rich & McWilliam (2000); McWilliam & Rich (2004) is not surprising and we expect to find higher metallicity M giants as the sample size increases. Only relatively metal rich stars reach the M spectral type, with those near the red giant branch tip reaching temperatures as cool as $``$3500 K. However, we find no correlation between spectral type and metallicity (the M9 giants are not extraordinarily metal rich). We also find among the M giants no stars well above the Solar iron abundance. One of the goals of a larger survey would be to determine whether there is a genuine deficiency of metal rich stars in the M giant population. Do such stars evolve through the M giant phase too rapidly to be found in this small sample? It is conceivable that higher metallicity results in greater mass loss, and that the the most metal rich stars evolve to the UV-bright AGB-manqué phases.
Our survey also shows a homogeneous $`\alpha `$-enhancement by $``$+0.3 dex up to solar metallicity, without significant differences among the various $`\alpha `$-elements (see Fig. 5), while the four solar neighborhood M giant stars in our control sample are consistent with solar \[$`\alpha `$/Fe\] values. As a further probe of the different degree of $`\alpha `$ enhancement between the Galactic bulge and the disk, Fig. 6 compares the spectra of a bulge and a solar neighborhood M giant star with almost identical photospheric parameters and iron abundances: the OH and Ti lines of the bulge giant are clearly deeper.
While we are dealing with a small sample of M giants, it is fair to say that if they were evolved from a younger sub population of the bulge, we would expect these stars to have Solar $`[\alpha /Fe]`$.
We have also derived carbon abundances for all of the observed stars from analysis of the CO bandheads. Fig. 7 shows the \[C/Fe\] and $`{}_{}{}^{12}\mathrm{C}/^{13}\mathrm{C}`$ abundance ratios as a function of the iron abundance, which provide major clues (the latter in particular, see Origlia et al. (2003) and references therein) to the efficiency of the mixing processes in the stellar interiors during the evolution along the RGB. We find some degree of \[C/Fe\] depletion with respect to solar values and very low <sup>12</sup>C/<sup>13</sup>C$``$10, as also measured in metal-poor halo giants both in the field and in globular clusters (see e.g. Suntzeff & Smith, 1991; Shetrone, 1996; Gratton et al., 2000, and reference therein), as well as in bulge clusters (Origlia, Rich & Castro, 2002; Shetrone, 2003; Origlia & Rich, 2004).
The classical theory (Iben, 1967; Charbonnel, 1994, and references therein) predicts a decrease of <sup>12</sup>C/<sup>13</sup>C$`40`$ after the first dredge–up, the exact amount mainly depending on the chemical composition and the extent of the convective zone. Additional mixing mechanisms due to further cool bottom processing (see e.g. Charbonnel, 1995; Denissenkov & Weiss, 1996; Cavallo, Sweigart & Bell, 1998; Boothroyd & Sackmann, 1999; Weiss, Denissenkov, & Charbonnel, 2000) can explain much lower $`{}_{}{}^{12}\mathrm{C}/^{13}\mathrm{C}`$ values, as those measured in the upper RGB stars.
R. Michael Rich acknowledges support from grant AST-0098739 from the National Science Foundation. Livia Origlia acknowledges financial support by the Agenzia Spaziale Italiana (ASI) and the Ministero dell’Istruzione, Università e Ricerca (MIUR).
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# Giambelli formulae for the equivariant quantum cohomology of the Grassmannian
## 1. Introduction
Let $`X`$ denote the Grassmannian $`Gr(p,m)`$ of subspaces of dimension $`p`$ in $`^m`$. One of the fundamental problems in the study of the equivariant quantum cohomology algebra of $`X`$ is to compute its structure constants, which are the 3-point, genus 0, equivariant Gromov-Witten invariants. The goal of this paper is to give a method for such a computation. Concretely, we realize the equivariant quantum cohomology as a ring given by generators and relations and we find polynomial representatives (i.e. Giambelli formulae) for the equivariant quantum Schubert classes (which form a module-basis for this ring).<sup>1</sup><sup>1</sup>1It is a standard fact that given a ring $`R/I`$, where $`R`$ is a polynomial ring and $`I`$ an ideal, together with some elements in $`R`$ which determine a module-basis for $`R/I`$, the structure constants for this basis can be computed using e.g. Groebner basis methods. These polynomials will be given by certain determinants which appear in the Jacobi-Trudi formulae for the factorial Schur functions (see §2 below for details).
Since the equivariant quantum cohomology ring specializes to both quantum and equivariant cohomology rings, we also obtain, as corollaries, determinantal formulae for the Schubert classes in the quantum and equivariant cohomology. In fact, in the quantum case, we recover Bertram’s quantum Giambelli formula . In the case of equivariant cohomology, we show that the factorial Schur functions represent the equivariant Schubert classes. The latter result, although not explicitly stated in the literature, seems to have been known before (cf. Remark 2 in §5).
We recall next some of the basic facts about the equivariant quantum cohomology and fix the notation. The torus $`T(^{})^m`$ acts on the Grassmannian $`X`$ by the action induced from the $`GL(m)`$action. The $`T`$equivariant cohomology of a point, denoted $`\mathrm{\Lambda }`$, is the polynomial ring $`[T_1,\mathrm{},T_m]`$ in the equivariant parameters $`T_i`$, graded by $`\mathrm{deg}T_i=1`$ (see §5.1 for a geometric interpretation of $`T_i`$). Let $`q`$ be an indeterminate of degree $`m`$. The $`T`$equivariant quantum cohomology of $`X`$, denoted $`QH_T^{}(X)`$, is a graded, commutative, $`\mathrm{\Lambda }[q]`$algebra with a $`\mathrm{\Lambda }[q]`$basis $`\{\sigma _\lambda \}`$ indexed by partitions $`\lambda =(\lambda _1,\mathrm{},\lambda _p)`$ included in the $`p\times (mp)`$ rectangle (i.e. $`\lambda _1,\mathrm{},\lambda _p`$ are integers such that $`mp\lambda _1\mathrm{}\lambda _p0`$). This basis is determined by the Schubert varieties in $`X`$, defined with respect to the standard flag, and the classes $`\sigma _\lambda `$ will be called (equivariant quantum) Schubert classes. More details are given in §3 and especially in . The equivariant quantum multiplication is denoted by $``$ and it is determined by the 3-pointed, genus $`0`$, equivariant Gromov-Witten invariants $`c_{\lambda ,\mu }^{\nu ,d}`$. In this paper we refer to these coefficients as the equivariant quantum Littlewood-Richardson coefficients, abbreviated EQLR. They have been introduced by Givental-Kim in (see also ), together with the equivariant quantum cohomology. Then, by definition,
$$\sigma _\lambda \sigma _\mu =\underset{d0}{}\underset{\nu }{}q^dc_{\lambda ,\mu }^{\nu ,d}\sigma _\nu .$$
The EQLR coefficients $`c_{\lambda ,\mu }^{\nu ,d}`$ are homogeneous polynomials in $`\mathrm{\Lambda }`$ of degree $`|\lambda |+|\mu ||\nu |md`$, where $`|\alpha |`$ denotes the sum of all parts of the partition $`\alpha `$. They are equal to the structure constants of equivariant cohomology if $`d=0`$ and to those of quantum cohomology (i.e. to the ordinary Gromov-Witten invariants) if $`|\lambda |+|\mu ||\nu |md=0`$. Equivalently, the quotient of equivariant quantum cohomology ring by the ideal generated by the equivariant parameters $`T_i`$ yields the quantum cohomology ring of $`X`$ (a $`[q]`$algebra), while the quotient by the ideal generated by $`q`$ yields the $`T`$equivariant cohomology of $`X`$ (a $`\mathrm{\Lambda }`$algebra). More about the EQLR coefficients, including a certain positivity, which generalizes the positivity enjoyed by the equivariant coefficients, can be found in .
### 1.1. Statement of the results
As we have noted before, the equivariant quantum Giambelli formula which we obtain is a “factorial” generalization of Bertram’s quantum Giambelli formula , and, as in that case, it doesn’t involve the quantum parameter $`q`$. It is closely related to a certain generalization of ordinary Schur functions, called factorial Schur functions. These are polynomials $`s_\lambda (x;t)`$ in two sets of variables: $`x=(x_1,\mathrm{},x_p)`$ and $`t=(t_i)_i`$. They play a fundamental role in the study of central elements in the universal enveloping algebra of $`𝔤𝔩(n)`$ (). One of their definitions is via a “factorial Jacobi-Trudi” determinant, and this determinant will represent the equivariant quantum Schubert classes. The basic properties of the factorial Schur functions are given in §2.
Denote by $`h_i(x;t)`$ (respectively by $`e_j(x;t)`$) the complete homogeneous (respectively, elementary) factorial Schur functions. They are equal to $`s_{(i)}(x;t)`$ (respectively $`s_{(1)^i}(x;t)`$), where $`(i)`$ (resp. $`(1)^i`$) denotes the partition $`(i,0,\mathrm{},0)`$ (resp. $`(1,\mathrm{},1,0,\mathrm{},0)`$, with $`i`$ $`1`$’s). Let $`t=(t_i)_i`$ be the sequence defined by
$$t_i=\{\begin{array}{cc}T_{mi+1}\hfill & \text{if }1im\hfill \\ 0\hfill & \text{ otherwise}\hfill \end{array}$$
where $`T_i`$ is an equivariant parameter. The next definitions are inspired from those used in the theory of factorial Schur functions. For an integer $`s`$ we define the shifted sequence $`\tau ^st`$ to be the sequence whose $`i`$th term $`(\tau ^st)_i`$ is equal to $`t_{s+i}`$. Let $`h_1,\mathrm{},h_{mp}`$ and $`e_1,\mathrm{},e_p`$ denote two sets of indeterminates (these correspond to the complete homogeneous, respectively, elementary, factorial Schur functions). Often one considers shifts $`s_\lambda (x|\tau ^st)`$ of $`s_\lambda (x|t)`$ by shifting the sequence $`(t_i)`$. Corresponding to these shifts we define the shifted indeterminates $`\tau ^sh_i`$, respectively $`\tau ^se_j`$, where $`s`$ is a nonnegative integer, as elements of $`\mathrm{\Lambda }[h_1,\mathrm{},h_{mp}]`$, respectively $`\mathrm{\Lambda }[e_1,\mathrm{},e_p]`$.
The definition of $`\tau ^sh_j`$ is given inductively as a function of $`\tau ^{s+1}h_j`$ and $`\tau ^{s+1}h_{j1}`$, and it is modelled on an equation which relates $`h_j(x|\tau ^st)`$ to $`h_j(x|\tau ^{s+1}t)`$ (see eq. (2.6) below, with $`a:=\tau ^{s+1}t`$). Concretely, $`\tau ^0h_j=h_j`$, $`\tau ^1h_j=h_j+(t_{j1+p}t_0)h_{j1}`$ and, in general,
(1.1)
$$\tau ^sh_j=\tau ^{s+1}h_j+(t_{j+ps}t_{1s})\tau ^{s+1}h_{j1}.$$
Similarly, the definition of $`\tau ^se_i`$ is modelled on an equation which relates $`e_{j+1}(x|\tau ^st)`$ to $`e_{j+1}(x|\tau ^{s1}t)`$ (see eq. (2.7) below, with $`a:=\tau ^{s1}t`$), and it is given by
(1.2)
$$\tau ^se_i=\tau ^{s1}e_i+(t_st_{pi+s+1})\tau ^{s1}e_{i1}$$
with $`\tau ^0e_i=e_i`$. By convention, $`h_0=e_0=1`$, $`h_j=0`$ if $`j<0`$ or $`j>mp`$, and $`e_i=0`$ if $`i<0`$ or $`i>p`$.
For $`\lambda `$ a partition in the $`p\times (mp)`$ rectangle define $`s_\lambda \mathrm{\Lambda }[h_1,\mathrm{},h_{mp}]`$, respectively $`\stackrel{~}{s}_\lambda \mathrm{\Lambda }[e_1,\mathrm{},e_p]`$, analogously to the definition of the factorial Schur function $`s_\lambda (x|t)`$, via the factorial Jacobi-Trudi determinants (cf. (2.4) below):
(1.3)
$$s_\lambda =det(\tau ^{1j}h_{\lambda _i+ji})_{1i,jp},$$
(1.4)
$$\stackrel{~}{s}_\lambda =det(\tau ^{j1}e_{\lambda _i^{}+ji})_{1i,jmp}.$$
Here $`\lambda ^{}=(\lambda _1^{},\mathrm{},\lambda _{mp}^{})`$ denotes the partition conjugate to $`\lambda `$, i.e. the partition in the $`(mp)\times p`$ rectangle whose $`i`$th part is equal to the length of the $`i`$th column of the Young diagram of $`\lambda `$.
By $`H_k`$, for $`mp<km`$, respectively $`E_k`$, for $`p<km`$, we denote the determinants from (1.4), respectively (1.3), above, corresponding to partitions $`(k)`$, respectively $`(1)^k`$, for the appropriate $`k`$:
(1.5)
$$H_k=det(\tau ^{j1}e_{1+ji})_{1i,jk},$$
(1.6)
$$E_k=det(\tau ^{1j}h_{1+ji})_{1i,jk}.$$
With this notation, we present the main result of this paper:
###### Theorem 1.1.
(a) There is a canonical isomorphism of $`\mathrm{\Lambda }[q]`$algebras
$$\mathrm{\Lambda }[q][h_1,\mathrm{},h_{mp}]/E_{p+1},\mathrm{},E_{m1},E_m+(1)^{mp}qQH_T^{}(X),$$
sending $`h_i`$ to $`\sigma _{(i)}`$. More generally, the image of $`s_\lambda `$ is the Schubert class $`\sigma _\lambda `$.
(b)(Dual version) There is a canonical isomorphism of $`\mathrm{\Lambda }[q]`$algebras
$$\mathrm{\Lambda }[q][e_1,\mathrm{},e_p]/H_{mp+1},\mathrm{},H_{m1},H_m+(1)^pqQH_T^{}(X),$$
sending $`e_j`$ to $`\sigma _{(1)^j}`$ and $`\stackrel{~}{s}_\lambda `$ to the Schubert class $`\sigma _\lambda `$.
Presentations by generators and relations for the equivariant quantum cohomology of the type A partial flag manifolds have also been obtained by A. Givental and B. Kim in , and, in general, for flag manifolds $`X=G/B`$ ($`G`$ a connected, semisimple, complex Lie group and $`B`$ a Borel subgroup) by B. Kim in . In the case of the Grassmannian, their presentation is given as
$$\mathrm{\Lambda }[a_1,\mathrm{},a_p,b_1,\mathrm{}.,b_{mp}]/I$$
where $`I`$ is the ideal generated by $`_{i+j=k}a_ib_j=e_k(T_1,\mathrm{},T_m)`$ and $`a_pb_{mp}=e_m(T_1,\mathrm{},T_m)+(1)^pq`$; the sum is over integers $`i,j`$ such that $`0ip`$ and $`0jmp`$; $`k`$ varies between $`1`$ and $`m1`$ and $`a_0=b_0=1`$; $`e_k`$ denotes the elementary symmetric function. We also note the results from , where a recursive relation, derived from a multiplication rule with the class $`\sigma _{(1)}`$ (a Pieri-Chevalley rule), gives another method to compute the EQLR coefficients.
### 1.2. Idea of proof
The proof of the theorem uses the theory of factorial Schur functions and a characterization of the equivariant quantum cohomology (). We prove first the “dual version” of the statement. For that we show an equivalent result, where in all the formulae $`\tau ^se_j`$ is replaced by $`e_j(x|\tau ^st)`$ (here $`x=(x_1,\mathrm{},x_p)`$ and $`t=(t_i)`$ is the sequence defined in the beginning of §1.1). A key role in this “translation” is played by a factorial version of the Jacobi-Trudi formula (cf. , pag. 56). Then we prove that the images of the polynomials $`s_\lambda (x|t)`$, for $`\lambda `$ included in the $`p\times (mp)`$ rectangle, form a $`\mathrm{\Lambda }[q]`$basis in the claimed presentation. A special multiplication formula, due to Molev and Sagan (, see also ), computes the product $`s_\lambda (x|t)s_{(1)}(x|t)`$ as a sum of $`s_\mu (x|t)`$, but with $`\mu `$ having possibly a part larger than $`mp`$. Using again the factorial Jacobi-Trudi formula, we prove that, modulo the relations ideal, this multiplication is precisely the equivariant quantum Pieri-Chevalley rule (see ). But this rule determines completely $`QH_T^{}(X)`$, and so the “dual” statement is proved. To prove the first statement, we construct a morphism from $`\mathrm{\Lambda }[q][h_1,\mathrm{},h_{mp}]`$ to the dual presentation and show that its kernel is the claimed ideal of relations.
Acknowledgements: I am indebted to S. Fomin and W. Fulton whose comments and suggestions enlightened the presentation of this paper. I am also thankful to A. L. Mare and A. Yong for some useful discussions and remarks.
## 2. Factorial Schur functions
This section presents those properties of the factorial Schur functions which are used later in the paper. The factorial Schur function $`s_\lambda (x|a)`$ is a homogeneous polynomial in two sets of variables: $`x=(x_1,\mathrm{},x_p)`$ and a doubly infinite sequence $`a=(a_i)_i`$. An initial (non-homogeneous) version of these polynomials, for $`a_i=i1`$ if $`1ip`$ and $`0`$ otherwise, was first studied by L. Biedenharn and J. Louck in , then by W. Chen and J. Louck in . The general version was considered by I. Macdonald , then studied further in and , Ch. I.3. These functions, as well as a different version of them, called shifted Schur functions (see ), play an important role in the study of the center of the universal enveloping algebra of $`𝔤𝔩(n)`$, Capelli identities and quantum immanants. In a geometric context, the factorial Schur functions appeared in , expressing the result of the localization of an equivariant Schubert class to a $`T`$fixed point.
For any variable $`y`$ and any sequence $`(a_i)`$ define the “generalized factorial power”:
$$(y|a)^k=(ya_1)\mathrm{}(ya_k).$$
Let $`\lambda `$ be a partition with at most $`p`$ parts. Following , define the factorial Schur function $`s_\lambda (x|a)`$ to be
$$s_\lambda (x|a)=\frac{det[(x_j|a)^{\lambda _i+pi}]_{1i,jp}}{det[(x_j|a)^{pi}]_{1i,jp}}.$$
Denote by $`h_k(x|a)`$ (respectively $`e_k(x|a)`$) the factorial complete homogeneous Schur functions (resp. the factorial elementary Schur functions). We adopt the usual convention that $`h_k(x|a)`$ and $`e_k(x|a)`$ are equal to zero if $`k`$ is negative. For $`k`$ larger than $`p`$, $`e_k(x|a)`$ is set equal to zero, also by convention. Let $`[a]`$ denote the ring of polynomials in variables $`a_i`$ ($`i`$-integer). We start enumerating the relevant properties of the factorial Schur functions:
(A) (Basis) The factorial Schur functions $`s_\lambda (x|a)`$, where $`\lambda `$ has at most $`p`$ parts, form a $`[a]`$basis for the ring of polynomials in $`[a][x]`$ symmetric in the $`x`$variables (, I.3, pag. 55). This is a consequence of the fact that
$$s_\lambda (x|a)=s_\lambda (x)+\text{terms of lower degree in }x\text{.}$$
Then
$$s_\lambda (x|a)s_\mu (x|a)=\underset{\nu }{}c_{\lambda \mu }^\nu (a)s_\nu (x|a)$$
where the coefficients $`c_{\lambda \mu }^\nu (a)`$ are homogeneous polynomials in $`[a]`$ of degree $`|\lambda |+|\mu ||\nu |`$. They are equal to the usual Littlewood-Richardson coefficients $`c_{\lambda \mu }^\nu `$ if $`|\lambda |+|\mu |=|\nu |`$ and are equal to $`0`$ if $`|\lambda |+|\mu |<|\nu |`$ (cf. §2).
(B) (Vanishing Theorem, , Thm. 2.1, see also ) Let $`\lambda =(\lambda _1,\mathrm{},\lambda _p)`$ and $`\rho =(\rho _1,\mathrm{},\rho _p)`$ be two partitions of length at most $`p`$. Define the sequence $`a_\rho `$ by $`a_\rho =(a_{\rho _1+p},\mathrm{},a_{\rho _p+1})`$. Then
(2.1)
$$s_\lambda (a_\rho |a)=\{\begin{array}{cc}0\hfill & \text{if }\lambda \rho \hfill \\ _{(i,j)\lambda }(a_{\lambda _i+pi+1}a_{p\lambda _j^{^{}}+j})\hfill & \text{if }\lambda =\rho \hfill \end{array}$$
where $`\lambda ^{}`$ is the conjugate partition of $`\lambda `$, and $`(i,j)\lambda `$ means that $`j`$ varies between $`1`$ and $`\lambda _i`$, if $`\lambda _i>0`$. We recall a consequence of the Vanishing Theorem:
###### Corollary 2.1.
The coefficients $`c_{\lambda \mu }^\nu (a)`$ satisfy the following properties:
(i) $`c_{\lambda \mu }^\nu (a)=0`$ if $`\lambda `$ or $`\mu `$ are not included in $`\nu `$.
(ii) If the partitions $`\mu `$ and $`\nu `$ are equal, then
$$c_{\lambda \mu }^\mu (a)=s_\lambda (a_\mu |a).$$
###### Proof.
(i) is part of the Theorem 3.1 in while (ii) follows from the proof of expression (10) in .∎
(C) (Factorial Pieri-Chevalley rule, see the proof of Prop. 3.2 in or , Thm. 9.1.) Let $`(1)`$ denote the partition $`(1,0,\mathrm{},0)`$ ($`p`$ parts) and let $`\lambda `$ be a partition with at most $`p`$ parts. Then
(2.2)
$$s_{(1)}(x|a)s_\lambda (x|a)=\underset{\mu \lambda }{}s_\mu (x|a)+c_{(1),\lambda }^\lambda (a)s_\lambda (x|a)$$
where $`\mu \lambda `$ means that $`\mu `$ contains $`\lambda `$ and has one more box than $`\lambda `$ (recall that, by definition, $`\mu `$ has at most $`p`$ parts). By Corollary 2.1, $`c_{(1),\lambda }^\lambda (a)=s_{(1)}(a_\lambda |a)`$ and the last expression turns out to be
(2.3)
$$s_{(1)}(a_\lambda |a)=\underset{i=1}{\overset{p}{}}a_{\lambda _i+p+1i}\underset{j=1}{\overset{p}{}}a_j.$$
(D) (Jacobi-Trudi identities, see , I.3, Ex. 20(c), pag. 56 or , Thm. 3.1.) Let $`\tau ^ra`$ be the sequence whose $`n^{th}`$ term is $`a_{n+r}`$ and let $`\lambda `$ be a partition with at most $`p`$ parts. Then
(2.4) $`s_\lambda (x|a)`$ $`=`$ $`det[h_{\lambda _ii+j}(x|\tau ^{1j}a)]_{1i,jp}`$
(2.5) $`=`$ $`det[e_{\lambda _i^{^{}}i+j}(x|\tau ^{j1}a)]_{1i,jmp},`$
where $`\lambda ^{}`$ is the partition conjugate of $`\lambda `$.
The following proposition gives an inductive way of computing the “shifted” polynomials $`h_k(x|\tau ^sa)`$, respectively $`e_k(x|\tau ^sa)`$, starting from the “unshifted” ones.
###### Proposition 2.2.
The following identities hold in $`[a][x]`$:
(2.6)
$$h_{i+1}(x|\tau ^1a)=h_{i+1}(x|a)+(a_{i+p}a_0)h_i(x|a)$$
(2.7)
$$e_{j+1}(x|\tau a)=e_{j+1}(x|a)+(a_1a_{pj+1})e_j(x|a).$$
###### Proof.
One uses the formulae
(2.8)
$$h_k(x|a)=\underset{1i_1\mathrm{}i_kp}{}(x_{i_1}a_{i_1})(x_{i_2}a_{i_2+1})\mathrm{}(x_{i_k}a_{i_k+k1}),$$
(2.9)
$$e_k(x|a)=\underset{1i_1<\mathrm{}<i_kp}{}(x_{i_1}a_{i_1})(x_{i_2}a_{i_21})\mathrm{}(x_{i_k}a_{i_k+k1})$$
(cf. , eqs. (1.2) and (1.3)). Then the computations are straightforward. ∎
By the Jacobi-Trudi formula, $`h_i(x|\tau ^sa)`$ (resp. $`e_j(x|\tau ^sa)`$) generate the algebra of polynomials in $`[a][x]`$ symmetric in the $`x`$variables. Then Prop. 2.2 implies:
###### Corollary 2.3.
The factorial complete homogeneous (resp. elementary) symmetric functions $`h_i(x|a)`$ for $`1i`$ (resp. $`e_j(x|a)`$ for $`1jp`$) generate the algebra of polynomials in $`[a][x]`$ symmetric in the $`x`$variables.
We will need the fact that the Jacobi-Trudi formula (2.5) generalizes to the case when $`\lambda =(1)^k`$ with $`k>p`$ (when, by convention, $`e_k(x|a)=0`$), and this is the content of the next proposition (see also equation (6.10) in ):
###### Proposition 2.4.
The following holds for any positive integer $`k>p`$:
$$det\left(h_{1+ji}(x|\tau ^{1j}a)\right)_{1i,jk}=0.$$
###### Proof.
Denote by $`E_k(x|a)`$ the determinant in question. Note that, if $`kp`$, this is equal to $`e_k(x|a)`$, by the Jacobi-Trudi formula. We will need a formula proved in , I.3, pag.56, Ex. 20(b):
(2.10)
$$\underset{r=0}{\overset{p}{}}(1)^re_r(x|a)h_{sr}(x|\tau ^{1s}a)=0$$
for any positive integer $`s`$. To prove the proposition, we use induction on $`kp+1`$. Expanding $`E_{p+1}(x|a)`$ after the last column yields:
$`E_{p+1}(x|a)`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{p}{}}}(1)^{r+1+p+1}h_{p+1r}(x|\tau ^pa)e_r(x|a)`$
$`=`$ $`(1)^{p+2}{\displaystyle \underset{r=0}{\overset{p}{}}}(1)^rh_{p+1r}(x|\tau ^pa)e_r(x|a)`$
$`=`$ $`0`$
where the last equality follows from (2.10) by taking $`s=p+1`$. Assume that $`E_k(x|a)=0`$ for all $`p<k<k_0`$. Expanding the determinant defining $`E_{k_0}(x|t)`$ after the last column, and using the induction hypothesis, yields
$`E_{k_0}(x|a)`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{p}{}}}(1)^{r+1+k_0}h_{k_0r}(x|\tau ^{1k_0}a)e_r(x|a)`$
$`=`$ $`(1)^{k_0+1}{\displaystyle \underset{r=0}{\overset{p}{}}}(1)^rh_{k_0r}(x|\tau ^{1k_0}a)e_r(x|a)`$
$`=`$ $`0`$
using again (2.10) with $`s=k_0`$. ∎
## 3. Equivariant quantum cohomology of the Grassmannian
In this section we recall some basic properties of the equivariant quantum cohomology. As before, $`X`$ denotes the Grassmannian $`Gr(p,m)`$, and $`\mathrm{\Lambda }`$ the polynomial ring $`[T_1,\mathrm{},T_m]`$. The ($`T`$)equivariant quantum cohomology of the Grassmannian, denoted by $`QH_T^{}(X)`$, is a deformation of both equivariant and quantum cohomology rings (for details on the latter cohomologies, see e.g. ). More precisely, $`QH_T^{}(X)`$ is a graded, commutative, $`\mathrm{\Lambda }[q]`$algebra, where the degree of $`q`$ is equal to $`m`$, which has a $`\mathrm{\Lambda }[q]`$basis $`\{\sigma _\lambda \}`$ indexed by the partitions $`\lambda `$ included in the $`p\times (mp)`$ rectangle. If $`\lambda =(\lambda _1,\mathrm{},\lambda _p)`$, the degree of $`\sigma _\lambda `$ is equal to $`|\lambda |=\lambda _1+\mathrm{}+\lambda _p`$. The multiplication of two basis elements $`\sigma _\lambda `$ and $`\sigma _\mu `$ is given by the equivariant quantum Littlewood-Richardson (EQLR) coefficients $`c_{\lambda ,\mu }^{\nu ,d}`$, where $`d`$ is a nonnegative integer:
$$\sigma _\lambda \sigma _\mu =\underset{d}{}\underset{\nu }{}q^dc_{\lambda ,\mu }^{\nu ,d}\sigma _\nu .$$
Recall that the EQLR coefficient $`c_{\lambda ,\mu }^{\nu ,d}`$ is a homogeneous polynomial in $`\mathrm{\Lambda }`$ of polynomial degree $`|\lambda |+|\mu ||\nu |md`$. If $`d=0`$ one recovers the structure constant $`c_{\lambda ,\mu }^\nu `$ in the equivariant cohomology of $`X`$, and if the polynomial degree is equal to $`0`$ (i.e. if $`|\lambda |+|\mu |=|\nu |+md`$) the EQLR coefficient is equal to the ordinary $`3`$pointed, genus $`0`$, Gromov-Witten invariant $`c_{\lambda ,\mu }^{\nu ,d}`$. The latter is a nonnegative integer equal to the number of rational curves in $`X`$ passing through general translates of the Schubert varieties in $`X`$ corresponding to the partitions $`\lambda ,\mu `$ and the dual of $`\nu `$.
The geometric definition of these coefficients can be found in . In fact, for the purpose of this paper, the algebraic characterization of the equivariant quantum cohomology from Proposition 3.2 below (which has a geometric proof in loc. cit.), suffices. We only remark that the equivariant quantum Schubert classes $`\sigma _\lambda `$ are determined by the equivariant Schubert classes $`\sigma _\lambda ^T`$, determined in turn by the Schubert varieties in $`X`$ defined with respect to the standard flag.<sup>2</sup><sup>2</sup>2Unlike the case of classical cohomology, in equivariant cohomology the Schubert class determined by a Schubert variety $`\mathrm{\Omega }_\lambda (F_{})`$, where $`F_{}`$ is a fixed flag in $`^m`$, depends on $`F_{}`$. The precise definition of $`\sigma _\lambda ^T`$ is not presently needed, but it is given in §5, where the equivariant cohomology ring is discussed in more detail.
From now on we specialize the sequence $`a=(a_i)_i`$ from the previous section to one, denoted $`t=(t_i)_i`$, encoding the equivariant parameters $`T_i`$:
$$t_i=\{\begin{array}{cc}T_{mi+1}\hfill & \text{if }1im\hfill \\ 0\hfill & \text{ otherwise.}\hfill \end{array}$$
Using this sequence, we recall next the equivariant quantum Pieri-Chvalley rule, as proved in . Given a partition $`\lambda `$, we denote by $`\lambda ^{}`$ the (uniquely determined) partition obtained by removing $`m1`$ boxes from the border rim of $`\lambda `$ (recall that the border rim of a Young diagram is the set of boxes that intersect the diagram’s SE border - see also the figure below).
If $`\lambda =(\lambda _1,\mathrm{},\lambda _p)`$, note that $`\lambda ^{}`$ exists only if $`\lambda _1=mp`$ and $`\lambda _p>0`$.
###### Proposition 3.1 (Equivariant quantum Pieri-Chvalley rule - cf. , Thm. 1).
The following formula holds in $`QH_T^{}(X)`$:
$$\sigma _\lambda \sigma _{(1)}=\underset{\mu \lambda }{}\sigma _\mu +c_{\lambda ,(1)}^\lambda (t)\sigma _\lambda +q\sigma _\lambda ^{}$$
where, by the formula (2.3), $`c_{\lambda ,(1)}^\lambda (t)`$ is equal to
$$c_{\lambda ,(1)}^\lambda (t)=\underset{i=1}{\overset{p}{}}T_{mp+i\lambda _i}\underset{j=mp+1}{\overset{m}{}}T_j.$$
The last term is omitted if $`\lambda ^{}`$ does not exist.
It turns out that the equivariant quantum Pieri-Chvalley rule determines completely the equivariant quantum cohomology algebra, in the following sense:
###### Proposition 3.2 ( Corollary 7.1).
Let $`(A,)`$ be a graded, commutative, associative $`\mathrm{\Lambda }[q]`$algebra with unit such that:
1. $`A`$ has an additive $`\mathrm{\Lambda }[q]`$basis $`\{s_\lambda \}`$ (graded as usual).
2. The equivariant quantum Pieri-Chevalley rule holds, i.e.
$$s_\lambda s_{(1)}=\underset{\mu \lambda }{}s_\mu +c_{\lambda ,(1)}^\lambda (t)s_\lambda +qs_\lambda ^{}$$
where the last term is omitted if $`\lambda ^{}`$ does not exist.
Then $`A`$ is canonically isomorphic to $`QH_T^{}(Gr(p,m))`$, as $`\mathrm{\Lambda }[q]`$algebras.
This proposition will be the main tool in proving the presentation and equivariant quantum Giambelli formula from the next section.
## 4. Proof of the Theorem
The strategy for the proof is to “guess” candidates for the presentation and for the polynomial representatives, using the insight provided by the similar results in quantum cohomology (see e.g. ) and some related results in equivariant cohomology ( §6). Then one attempts to prove that the guessed polynomials form a $`\mathrm{\Lambda }[q]`$basis in the candidate presentation, and they multiply according to the EQ Pieri-Chevalley rule. Proposition 3.2 will ensure that the guessed algebra will be canonically isomorphic to $`QH_T^{}(X)`$ and that the polynomials considered will represent the equivariant quantum Schubert classes.
It turns out that each of the quantum presentations from (the usual one, involving the $`h`$ variables, and the “dual” one, involving the variables $`e`$) implies an equivariant quantum presentation (see respectively Theorems 4.3 and 4.2 below). The equivariant generalizations are obtained by taking the factorial versions, via the factorial Jacobi-Trudi formula (§2, property D) of all the expressions involved in the original quantum presentations.
Before stating the first result, we recall the notation from the introduction: $`h_1,\mathrm{},h_{mp}`$ and $`e_1,\mathrm{},e_p`$ denote two sets of indeterminates; the definitions of $`\tau ^sh_i`$, $`\tau ^se_j`$ and of $`H_k`$ ($`mp<km`$) respectively $`E_k`$ ($`p<km`$) are those given in the equations (1.1),(1.2) and (1.5),(1.6) above. For $`\lambda `$ in the $`p\times (mp)`$ rectangle recall that:
(4.1)
$$s_\lambda =det(\tau ^{1j}h_{\lambda _i+ji})_{1i,jp}$$
respectively
(4.2)
$$\stackrel{~}{s}_\lambda =det(\tau ^{j1}e_{\lambda _i^{}+ji})_{1i,jmp},$$
(cf. (1.3) and (1.4)) with the usual conventions that $`h_k=0`$ for $`k<0`$ and $`k>mp`$ respectively $`e_i=0`$ if $`i<0`$ or $`i>p`$. Before proving the theorem, we need a Nakayama-type result, which will be used several times in the paper:
###### Lemma 4.1 (cf. , Exerc. 4.6).
Let $`M`$ be an $`R`$algebra graded by nonnegative integers. Assume that $`R`$ is also graded (by nonnegative integers) and let $`I`$ be a homogeneous ideal in $`R`$ consisting of elements of positive degree. Let $`m_1,\mathrm{},m_k`$ be homogeneous elements whose images generate $`M/IM`$ as an $`R/I`$-module. Then $`m_1,\mathrm{},m_k`$ generate $`M`$ as an $`R`$module.
###### Proof.
Let $`m`$ be a nonzero homogeneous element of $`M`$. We use induction on its degree. Assume $`\mathrm{deg}m=0`$. The hypothesis implies that
(4.3)
$$m=r_1m_1+\mathrm{}+r_km_kmodIM$$
where $`r_i`$ are elements in $`R`$. Since $`I`$ contains only elements of positive degree, it follows that the equality holds in $`M`$ as well. Let now $`\mathrm{deg}m>0`$. Writing $`m`$ as in (4.3), implies that
$$m\underset{i}{}r_im_i=\underset{j}{}a_jm_j^{}$$
for some (finitely many) $`a_jI`$ and $`m_j^{}M`$. Again, since $`I`$ contains only elements of positive degree, $`\mathrm{deg}m_j^{}<\mathrm{deg}m`$ for each $`j`$. The induction hypothesis implies that each $`m_j^{}`$ is an $`R`$combination of $`m_i`$’s, which finishes the proof. ∎
We prove next the “dual version” statement from the main theorem.
###### Theorem 4.2.
There exists a canonical isomorphism of $`\mathrm{\Lambda }[q]`$algebras
$$\mathrm{\Lambda }[q][e_1,\mathrm{},e_p]/H_{mp+1},\mathrm{},H_m+(1)^pqQH_T^{}(X),$$
sending $`e_i`$ to $`\sigma _{(1)^i}`$ and $`\stackrel{~}{s}_\lambda `$ to the equivariant quantum Schubert class $`\sigma _\lambda `$.
###### Proof.
Note first that
$$e_j(x|t)=e_j(x)+f(t,x),$$
where $`f(t,x)`$ is a homogeneous polynomial in the variables $`x`$ and $`t`$, but of degree in the variables $`x`$ less than $`j`$. Since the usual elementary symmetric functions $`e_1(x),\mathrm{},e_p(x)`$ are algebraically independent over $``$, it follows that the elementary factorial Schur functions $`e_1(x|t),\mathrm{},e_p(x|t)`$ are algebraically independent over $`\mathrm{\Lambda }`$. Then there is a canonical isomorphism
$$\mathrm{\Lambda }[q][e_1(x|t),\mathrm{},e_p(x|t)]\mathrm{\Lambda }[q][e_1,\mathrm{},e_p],$$
sending $`s_\lambda (x|t)`$ to $`\stackrel{~}{s}_\lambda `$, and $`h_{mp+i}(x|t)`$ ($`1ip`$) to $`H_{mp+i}`$, by the factorial Jacobi-Trudi identity. This induces an isomorphism between
$$A:=\mathrm{\Lambda }[q][e_1(x|t),\mathrm{},e_p(x|t)]/h_{mp+1}(x|t),\mathrm{},h_m(x|t)+(1)^pq$$
and
$$\mathrm{\Lambda }[q][e_1,\mathrm{},e_p]/H_{mp+1},\mathrm{},H_m+(1)^pq.$$
By Prop. 3.2, it remains to show that the images of $`s_\lambda (x|t)`$ in $`A`$, as $`\lambda `$ varies over the partitions included in the $`p\times (mp)`$ rectangle, form a $`\mathrm{\Lambda }[q]`$-basis of $`A`$, satisfying the equivariant quantum Pieri-Chevalley rule.
Generating set. This follows from Lemma 4.1, applied to $`M=A`$, $`R=\mathrm{\Lambda }[q]`$ and $`I`$ the ideal generated by $`q`$ and $`T_1,\mathrm{},T_p`$ (in which case $`M/I`$ is the classical cohomology of $`X`$).
Linear independence. Assume that $`q^{d_\lambda }c_\lambda s_\lambda (x|t)=0`$ in $`A`$, for $`c_\lambda `$ in $`\mathrm{\Lambda }`$, where $`\lambda `$ is included in the $`p\times (mp)`$ rectangle. This implies that $`q^{d_\lambda }c_\lambda s_\lambda (x|t)`$ is in the ideal generated by $`h_{mp+1}(x|t),\mathrm{},h_m(x|t)+(1)^pq`$. By Cor. 2.1 (1), any element of this ideal can be written as:
$$\underset{\mu }{}q^{d_\mu ^{}}c_\mu ^{}s_\mu (x|t)+\underset{\nu }{}q^{d_\nu ^{}}c_\nu ^{\prime \prime }s_\nu (x|t)(h_m(x|t)+(1)^pq)$$
where $`\mu ,\nu `$ have at most $`p`$ parts, $`\mu `$ is outside the $`p\times (mp)`$ rectangle, and $`c_\mu ^{},c_\nu ^{\prime \prime }`$ are in $`\mathrm{\Lambda }`$. Note that $`s_\nu (x|t)h_m(x|t)`$ expands also into a sum of factorial Schur functions indexed by partitions outside the $`p\times (mp)`$ rectangle. Since the factorial Schur functions form a $`\mathrm{\Lambda }[q]`$basis for the polynomials in $`\mathrm{\Lambda }[q][x_1,\mathrm{},x_p]`$ symmetric in the $`x`$variables, it follows that all $`c_\lambda `$ (and $`c_\mu ^{},c_\nu ^{\prime \prime }`$) must be equal to zero, as desired.
Equivariant quantum Pieri-Chevalley. The factorial Pieri-Chevalley rule (§2, Property (C)) states that if $`\lambda `$ is included in the $`p\times (mp)`$ rectangle, then
$$s_\lambda (x|t)s_{(1)}(x|t)=\underset{\mu }{}s_\mu (x|t)+c_{(1),\lambda }^\lambda (t)s_\lambda (x|t)+s_{\overline{\lambda }}(x|t)$$
where $`\mu `$ runs over all partitions in the $`p\times (mp)`$ obtained from $`\lambda `$ by adding one box; the last term is omitted if $`\lambda _1`$, the first part of $`\lambda `$, is not equal to $`mp`$. If $`\lambda _1=mp`$ then $`\overline{\lambda }=(\overline{\lambda }_1,\mathrm{},\overline{\lambda }_p)`$ is given by adding a box to the first row of $`\lambda `$, i.e. $`\overline{\lambda }_1=mp+1`$ and $`\overline{\lambda }_i=\lambda _i`$ for $`i2`$. Since the images of $`s_\lambda (x|t)`$, as $`\lambda `$ varies in the $`p\times (mp)`$ rectangle, form a $`\mathrm{\Lambda }[q]`$basis for $`A`$, it is enough to show that
$$s_{\overline{\lambda }}(x|t)=qs_\lambda ^{}(x|t)modJ$$
where $`J`$ is the ideal generated by $`h_{mp+1}(x|t),\mathrm{},h_m(x|t)+(1)^pq`$. By the factorial Jacobi-Trudi formula (§2, Property (D)), it follows that
(4.4)
$$s_{\overline{\lambda }}(x|t)=det\left(\begin{array}{cccc}h_{mp+1}(x|t)& h_{mp+2}(x|\tau ^1t)& \mathrm{}& h_m(x|\tau ^{1p}t)\\ h_{\lambda _21}(x|t)& h_{\lambda _2}(x|\tau ^1t)& \mathrm{}& h_{\lambda _2+p2}(x|\tau ^{1p}t)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& h_{\lambda _p}(x|\tau ^{1p}t)\end{array}\right).$$
We analyze next the first row of this determinant.
Claim. Let $`i,j`$ be two integers such that $`2jip`$. Then
$$h_{mp+i}(x|\tau ^{1j}t)=h_{mp+i}(x|\tau ^{1(j1)}t).$$
###### Proof of the Claim.
Equation (2.6) implies that
$$h_j(x|\tau ^st)=h_j(x|\tau ^{s+1}t)+(t_{j+ps}t_{s+1})h_{j1}(x|\tau ^{s+1}t),$$
hence,
$$h_{mp+i}(x|\tau ^{1j}t)=h_{mp+i}(x|\tau ^{1(j1)}t)+(t_{m+1+ij}t_{2j})h_{mp+i1}(x|\tau ^{1(j1)}t).$$
The Claim follows then from the definition of $`(t_i)`$, since $`t_{m+1+ij}=t_{2j}=0`$. ∎
It follows that for any integer $`1ip`$,
(4.5)
$$h_{mp+i}(x|\tau ^{1i}t)=h_{mp+i}(x|t).$$
In particular,
$$h_{mp+i}(x|\tau ^{1i}t)=0modJ$$
if $`1ip1`$, and
$$h_m(x|\tau ^{1p}t)=(1)^{p+1}qmodJ.$$
Therefore, expanding the determinant in (4.4) after the first row, yields:
(4.6)
$$s_{\overline{\lambda }}(x|t)=(1)^{p+1}(1)^{p+1}qdet\left(\begin{array}{cccc}h_{\lambda _21}(x|t)& h_{\lambda _2}(x|\tau ^1t)& \mathrm{}& h_{\lambda _2+p3}(x|\tau ^{2p}t)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& h_{\lambda _p1}(x|\tau ^{2p}t)\end{array}\right)$$
in $`A`$. If $`\lambda _p=0`$, the last row of the determinant in (4.6) contains only zeroes; if $`\lambda _p>0`$, the determinant is equal to $`s_\lambda ^{}(x|t)`$, by the Jacobi-Trudi formula. Summarizing, $`s_{\overline{\lambda }}(x|t)`$ is equal to $`qs_\lambda ^{}(x|t)`$ in $`A`$, or it is equal to zero if $`\lambda ^{}`$ does not exists. This finishes the proof of the equivariant quantum Pieri-Chevalley rule, hence also the proof of the theorem.∎
We are ready to prove the first part of the main result, which involves the $`h`$ variables. We use the notation preceding Thm. 4.2 above.
###### Theorem 4.3.
There exist a canonical isomorphism of $`\mathrm{\Lambda }[q]`$algebras
$$\mathrm{\Lambda }[q][h_1,\mathrm{},h_{mp}]/E_{p+1},\mathrm{}.,E_m+(1)^{mp}qQH_T^{}(X),$$
such that $`h_j`$ is sent to $`\sigma _{(j)}`$ and $`s_\lambda `$ to the equivariant quantum Schubert class $`\sigma _\lambda `$.
###### Proof.
Consider the $`\mathrm{\Lambda }[q]`$algebra morphism
$$\mathrm{\Psi }:\mathrm{\Lambda }[q][h_1,\mathrm{},h_{mp}]\mathrm{\Lambda }[q][e_1(x|t),\mathrm{},e_p(x|t)]/h_{mp+1}(x|t),\mathrm{},h_m(x|t)+(1)^pq$$
sending $`h_k`$ to the image of $`h_k(x|t)=det(e_{1+ji}(x|\tau ^{1j}t))_{1i,jk}`$. Recall that the last quotient is denoted by $`A`$ and it is canonically isomorphic to $`QH_T^{}(X)`$, by the previous proof. We will show that the images under $`\mathrm{\Psi }`$ of $`E_{p+1},E_{p+2},\mathrm{},E_{m1},E_m+(1)^{mp}q`$ are equal to zero in $`A`$ (where $`E_i`$ is defined by equation (1.6)). First, we need the following claim:
Claim. The following formulae hold in $`A`$:
(4.7)
$$\mathrm{\Psi }(\tau ^sh_j)=h_j(x|\tau ^st),$$
for any nonnegative integers $`s`$ and $`j`$ with $`j<m`$, and
(4.8)
$$\mathrm{\Psi }(\tau ^{(m1)}h_m)=h_m(x|\tau ^{(m1)}t)+(1)^pq.$$
###### Proof of the Claim.
By definition, both $`\tau ^sh_j`$ and $`h_j(x|\tau ^st)`$ satisfy the same recurrence relations (given respectively by the equations (1.1) and (2.6)). This implies that there exist polynomials $`P_1(t),\mathrm{},P_s(t)`$ in $`\mathrm{\Lambda }`$, with $`\mathrm{deg}P_k(t)=k`$, such that
$$\tau ^sh_j=h_j+\underset{k=1}{\overset{s}{}}P_k(t)h_{jk},$$
respectively
(4.9)
$$h_j(x|\tau ^st)=h_j(x|t)+\underset{k=1}{\overset{s}{}}P_k(t)h_{jk}(x|t).$$
If $`jmp`$, then $`\mathrm{\Psi }(h_j)=h_j(x|t)`$ in $`A`$, by the definition of $`\mathrm{\Psi }`$, thus
(4.10)
$$\mathrm{\Psi }(\tau ^sh_j)=h_j(x|\tau ^st).$$
If $`mp+1j<m`$, $`h_j=0`$ by convention, whereas $`h_j(x|t)=0`$ in $`A`$, so equation (4.10) also holds in this case. If $`s=m1`$ and $`j=m`$, we have
$$\begin{array}{cc}\hfill \mathrm{\Psi }(\tau ^{(m1)}h_m)& =\mathrm{\Psi }(h_m+\underset{k=1}{\overset{m1}{}}P_k(t)h_{mk})\hfill \\ & =\mathrm{\Psi }(h_m)+\underset{k=1}{\overset{m1}{}}P_k(t)\mathrm{\Psi }(h_{mk})\hfill \\ & =h_m(x|t)+(1)^pq+\underset{k=1}{\overset{m1}{}}P_k(t)h_{mk}(x|t)\hfill \\ & =h_m(x|\tau ^{(m1)}t)+(1)^pq,\hfill \end{array}$$
where the third equality follows from the fact that $`h_m=0`$ and $`h_m(x|t)+(1)^pq=0`$ in $`A`$; the fourth equality follows from the expansion (4.9) of $`h_m(x|\tau ^{(m1)}t)`$.∎
By definition, $`\mathrm{\Psi }(E_i)`$ is equal to the image in $`A`$, through $`\mathrm{\Psi }`$, of
$$det\left(\begin{array}{cccccc}h_1& \tau ^1h_2& \mathrm{}& \tau ^{(s1)}h_s& \mathrm{}& \tau ^{(i1)}h_i\\ 1& \tau ^1h_1& \mathrm{}& \mathrm{}& \mathrm{}& \tau ^{(i1)}h_{i1}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0& 1& \tau ^{(i1)}h_1\end{array}\right).$$
If $`p+1i<m`$, this determinant contains only $`\tau ^sh_j`$ with $`j<m`$. Then, by the equation (4.7) from the claim, $`\mathrm{\Psi }(E_i)`$ is the image in $`A`$ of the determinant
$$det\left(\begin{array}{cccccc}h_1(x|t)& h_2(x|\tau ^1t)& \mathrm{}& h_s(x|\tau ^{(s1)}t)& \mathrm{}& h_i(x|\tau ^{(i1)}t)\\ 1& h_1(x|\tau ^1t)& \mathrm{}& \mathrm{}& \mathrm{}& h_{i1}(x|\tau ^{(i1)}t)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0& 1& h_1(x|\tau ^{(i1)}t)\end{array}\right)$$
which, by Proposition 2.4, is equal to zero, since $`i>p`$. To compute the image of $`\mathrm{\Psi }(E_m)`$ we use both the equations (4.7) and (4.8). Then $`\mathrm{\Psi }(E_m)`$ is equal to the image in $`A`$ of
$$\begin{array}{c}det\left(\begin{array}{cccccc}h_1(x|t)& h_2(x|\tau ^1t)& \mathrm{}& h_s(x|\tau ^{(s1)}t)& \mathrm{}& h_m(x|\tau ^{(m1)}t)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0& 1& h_1(x|\tau ^{(m1)}t)\end{array}\right)+\\ det\left(\begin{array}{ccccc}0& 0& \mathrm{}& 0& (1)^pq\\ 1& h_1(x|\tau ^1t)& \mathrm{}& \mathrm{}& h_{m1}(x|\tau ^{(m1)}t)\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& 0& 1& h_1(x|\tau ^{(m1)}t)\end{array}\right)\end{array}$$
The first determinant is equal to zero, by Prop. 2.4, and the second is $`(1)^{m+1}(1)^pq`$. It follows that $`\mathrm{\Psi }(E_m)+(1)^{mp}q`$ is equal to zero in $`A`$, as claimed. Thus $`\mathrm{\Psi }`$ induces a $`\mathrm{\Lambda }[q]`$algebra morphism
$$\mathrm{\Psi }^{}:\mathrm{\Lambda }[q][h_1,\mathrm{},h_{mp}]/E_{p+1},\mathrm{}.,E_m+(1)^{mp}qA.$$
Note that $`\tau ^sh_m`$ does not appear in the determinant defining $`s_\lambda `$, for any $`\lambda `$ in the $`p\times (mp)`$ rectangle, so $`\mathrm{\Psi }^{}`$ sends $`s_\lambda `$ to the image of $`s_\lambda (x|t)`$ in $`A`$. Applying Lemma 4.1 with $`M=\mathrm{\Lambda }[q][h_1,\mathrm{},h_{mp}]/E_{p+1},\mathrm{}.,E_m+(1)^{mp}q`$ and $`I`$ the ideal generated by $`q`$ and $`T_1,\mathrm{},T_m`$ implies that the polynomials $`s_\lambda `$ generate $`\mathrm{\Lambda }[q][h_1,\mathrm{},h_{mp}]/E_{p+1},\mathrm{}.,E_m+(1)^{mp}q`$ as a $`\mathrm{\Lambda }[q]`$-module. Since their images through $`\mathrm{\Psi }^{}`$ form a $`\mathrm{\Lambda }[q]`$basis, they must form a $`\mathrm{\Lambda }[q]`$basis, as well. Hence $`\mathrm{\Psi }^{}`$ is an isomorphism, as desired. ∎
Remark: The Theorems 4.2 and 4.3 are proved without using the corresponding results from quantum cohomology. In particular, we obtain a new proof for Bertram’s quantum Giambelli formula (see ). (Recall that the quantum cohomology ring of $`X`$ is a graded $`[q]`$algebra isomorphic to $`QH_T^{}(X)/(T_1,\mathrm{},T_m)`$, hence the quantum Giambelli formula is obtained by taking $`T_1=\mathrm{}=T_m=0`$ in the determinants from the above-mentioned theorems.)
## 5. Giambelli formulae in equivariant cohomology
The goal of this section is to state the equivariant Giambelli formulae implied by their equivariant quantum counterparts from the previous section. We will also use this opportunity to define rigourously the equivariant Schubert classes involved, and provide, without proof, a geometric interpretation for the factorial Schur functions.
Let $`T`$ be the usual torus, and let $`ETBT`$ be the universal $`T`$bundle. If $`X`$ is a topological space with a $`T`$action, there is an induced $`T`$action on $`ET\times X`$ given by $`t(e,x)=(et^1,tx)`$. The (topological) quotient space $`(ET\times X)/T`$ is denoted by $`X_T`$. By definition, the ($`T`$)equivariant cohomology of $`X`$, denoted $`H_T^{}(X)`$, is equal to the usual cohomology of $`X_T`$. The $`X`$bundle $`X_TBT`$ gives $`H_T^{}(X)`$ the structure of a $`\mathrm{\Lambda }`$algebra, where $`\mathrm{\Lambda }`$ denotes the equivariant cohomology of a point $`H_T^{}(pt)=H^{}(BT)`$.
Let now $`X`$ be the Grassmannian of subspaces of dimension $`p`$ in $`^m`$ with the $`T`$action induced from the usual $`GL(m)`$action. We define next the equivariant Schubert classes which determine the equivariant quantum classes $`\sigma _\lambda `$ used in previous sections (see also ). Let
$$F_{}:(0)F_1\mathrm{}F_m=^m$$
be the standard flag, so $`F_i=f_1,\mathrm{},f_i`$ and $`f_i=(0,\mathrm{},1,\mathrm{},0)`$ (with $`1`$ in the $`i`$th position). If $`\lambda =(\lambda _1,\mathrm{},\lambda _p)`$ is a partition included in the $`p\times (mp)`$ rectangle, define the Schubert variety $`\mathrm{\Omega }_\lambda (F_{})`$ by
(5.1)
$$\mathrm{\Omega }_\lambda (F_{})=\{VX:dimVF_{mp+i\lambda _i}i\}.$$
Since the flag $`F_{}`$ is $`T`$-invariant, the Schubert variety $`\mathrm{\Omega }_\lambda (F_{})`$ will be $`T`$invariant as well, so it determines a Schubert class $`\sigma _\lambda ^T`$ in $`H_T^{2|\lambda |}(X)`$. The following result is a consequence of the Theorems 4.2 and 4.3 (the notation is from the previous section):
###### Corollary 5.1.
(a) There exists a canonical isomorphism of $`\mathrm{\Lambda }`$algebras
$$\mathrm{\Lambda }[h_1,\mathrm{},h_{mp}]/E_{p+1},\mathrm{}.,E_mH_T^{}(X),$$
sending $`h_j`$ to $`\sigma _{(j)}^T`$ and $`s_\lambda `$ is the equivariant Schubert class $`\sigma _\lambda ^T`$.
(b) There exists a canonical isomorphism of $`\mathrm{\Lambda }`$algebras
$$\mathrm{\Lambda }[e_1,\mathrm{},e_p]/H_{mp+1},\mathrm{},H_mH_T^{}(X),$$
sending $`e_i`$ to $`\sigma _{(1)^i}^T`$ and $`\stackrel{~}{s}_\lambda `$ to the equivariant Schubert class $`\sigma _\lambda ^T`$.
###### Proof.
It is known (see e.g. ) that there is a canonical isomorphism of $`\mathrm{\Lambda }`$algebras
$$QH_T^{}(X)/qH_T^{}(X)$$
sending the equivariant quantum Schubert class $`\sigma _\lambda `$ from the previous section to $`\sigma _\lambda ^T`$. Then the Corollary follows from the Theorems 4.2 and 4.3. ∎
Remarks: 1. The proof of the Corollary can be given without using the equivariant quantum cohomology. There is an analogue of Prop. 3.2, stating that the Pieri-Chevalley rule determines the equivariant cohomology algebra. Then a “strictly equivariant” proof of Cor. 5.1 can be obtained by taking $`q=0`$ in all the assertions from the previous section.
2. The fact that the factorial Schur functions represent the equivariant Schubert classes can be also be deduced, indirectly, by combining the fact that the double Schubert polynomials represent the equivariant Schubert classes in the complete flag variety (cf. and ) and that, when indexed by a Grassmannian permutation, these polynomials are actually factorial Schur functions. The latter holds because the vanishing property characterizing the factorial Schur functions (§2, property (B)), is also satisfied by the double Schubert polynomials in question (see , pag. 33). However, the details of this connection are missing from the literature.
3. It is well known that the equivariant Schubert classes are determined by their restriction to the torus fixed points in $`X`$. Formulae for such restrictions have been obtained by A. Knutson - T. Tao in and, recently, by V. Lakshmibai - K.N. Raghavan - P. Sankaran in .
### 5.1. A geometric interpretation of the factorial Schur functions
Consider the tautological short exact sequence on $`X`$:
(5.2)
$$0SVQ0,$$
which is clearly $`T`$equivariant. Let $`x_1,\mathrm{},x_p`$ be the equivariant Chern roots of the bundle $`S`$. There is a weight space decomposition of the trivial (but not equivariantly trivial) vector bundle $`V`$ into a sum of $`T`$equivariant line bundles:
$$V=L_1\mathrm{}L_m.$$
Let $`T_i`$ be the equivariant first Chern class of $`L_i`$.<sup>3</sup><sup>3</sup>3All the minus signs are for positivity reasons. It turns out, for example, that $`c_1^T(L_i)`$ is the Chern class of $`𝒪_{^{\mathrm{}}}(1)`$ (see e.g. §7). Define the sequence $`(t_i)`$ as usual, using the formula from §1.1.
###### Proposition 5.2.
In $`H_T^{}(X)`$, the equivariant Schubert class $`\sigma _\lambda ^T`$ is equal to the factorial Schur polynomial $`s_\lambda (x|t)`$.
###### Idea of proof.
The equivariant Schubert class $`\sigma _\lambda ^T`$ is a cohomology class on the infinite dimensional space $`X_T`$. The first step of the proof is to approximate this class by a class $`(\sigma _\lambda )_{T,n}`$ on a finite-dimensional “approximation” $`X_{T,n}`$ ($`n0`$) of $`X_T`$. This is standard (see e.g. or ) and uses the $`T`$bundle $`(^{n+1}0)^m(^n)^m`$ which approximates the universal $`T`$bundle $`ETBT`$. Then $`X_{T,n}:=(ET_n\times X)/T`$, which, in fact, is equal to the Grassmann bundle $`𝔾(p,𝒪_{(1)}(1)\mathrm{}𝒪_{(m)}(1))`$, where $`𝒪_{(i)}(1)`$ denotes the tautological line bundle over the $`i`$th component of $`(^n)^m`$. Using this procedure one obtains the class $`(\sigma _\lambda )_{T,n}`$ as the cohomology class determined by the subvariety $`(\mathrm{\Omega }_\lambda )_{T,n}`$ of $`X_{T,n}`$. The second step is to use the definition of $`(\mathrm{\Omega }_\lambda )_{T,n}`$ to realize it as the degeneracy locus from , Thm. 14.3, whose cohomology class is given as a certain determinantal formula in the Chern classes of the vector bundles $`S_{T,n},V_{T,n}`$ and $`Q_{T,n}`$ on $`X_{T,n}`$ induced by the tautological sequence on $`X`$. Finally, one proves that the determinant in question is equal to the claimed factorial Schur polynomial, which ends the proof.∎
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# 1 Introduction
## 1 Introduction
Iwasawa decomposition of a simple complex Lie group $`D`$ has many useful applications in group theory. For instance, it gives rise to Poisson-Lie structure on the compact real form $`G`$ of $`D`$, such that the dual Poisson-Lie group is identified with the subgroup $`AN`$ of $`D=G^{}`$. In fact, $`D`$ is said to be the Lu-Weinstein Drinfeld double of $`G`$ and $`AN`$. One of our motivation for superizing the Iwasawa decomposition is to define super Lu-Weinstein Drinfeld double. The reader may be surprised that the superization of the Iwasawa decomposition has not yet been considered in the literature. The reason is simple: it was reported in that real forms of complex simple Lie superalgebra are never compact. As the compacity is an important ingredient of the standard Iwasawa decomposition, this result seemed to imply that there is no super-version of the decomposition. Recently, we have shown in that the definition of the real form given in for the supercase was too restrictive. In fact, we argued that there is a notion of graded real form which is more flexible than one. In particular, we show that working with this new concept, each supergroup of the series $`OSp(2r,s,)`$ has precisely one compact graded real form. In this article we extend the results of in two directions. First, we show the existence of compact graded real form also for $`SL(n,m,)`$ supergroup. Secondly, we argue that our new concept is very natural and powerful since it does allow to construct the super-analogue of the Iwasawa decomposition.
In section $`2`$, we give the definition of the notion of normal and graded real form. Next, we define the supergroup $`SL(n,m,)`$. Finally, we construct the compact graded real form $`SU(n,m)`$ of the supergroup $`SL(n,m,)`$ following the ideas of our paper .
In section $`3`$, we define the three ingredients used in the Iwasawa decomposition i.e. the superalgebra of complex ”functions” on respectively the real supergroup $`SL(n,m,)`$, $`SU(n,m)`$ and $`s(AN)`$. This last supergroup is the superization of the real group $`AN`$. Next, we show the existence of the Iwasawa decomposition $`SL(n,m,)=SU(n,m)s(AN)`$. Our method uses the superization of the Gram-Schmidt orthonormalisation of a family of vectors.
## 2 Compact graded real form of $`SL(n,m,)`$
We first give the definitions of complex matrix Lie supergroup and of its real form. We illustrate these notions on the supergroup $`SL(n,m,)`$.
###### Definition 2.1
A complex matrix Lie supergroup $``$ is a complex superbialgebra generated by finite set of odd and even generators subject to polynomials relations. Those relations are supposed to generate a superideal of superbialgebra such that the quotient can be given the structure of a Hopf superalgebra (i.e. the antipode can be defined).
Now we turn to the definition of normal and graded real form:
###### Definition 2.2
A normal real form of a complex Lie supergroup $``$ is a pair $`(,\sigma )`$ where $`\sigma `$ is an even map from $``$ to $``$ such that:
$$(\sigma \sigma )\mathrm{\Delta }(x)=\mathrm{\Delta }(\sigma (x)),$$
$`(1)`$
$$ϵ(\sigma (x))=\overline{ϵ(x)},$$
$`(2)`$
$$\sigma (\lambda x+\mu y)=\overline{\lambda }\sigma (x)+\overline{\mu }\sigma (y),$$
$`(3)`$
$$\sigma (xy)=\sigma (x)\sigma (y),$$
$`(4)`$
$$S\sigma S\sigma (x)=x,$$
$`(5a)`$
$$\sigma (\sigma (x))=x,$$
$`(6a)`$
with $`x,y`$ and $`\lambda ,\mu `$. If the two last properties are replaced by the following:
$$S\sigma S\sigma (x)=(1)^{|x|}x,$$
$`(5b)`$
$$\sigma (\sigma (x))=(1)^{|x|}x,$$
$`(6b)`$
then we have a graded real form (cf. ).
###### Remark 2.1
The map $`\sigma `$ is the generalisation to the supergroup framework of the concept of star structure. The latter is well-known in the Hopf algebra literature, where the real form of a complex Hopf algebra is by definition the star structure. Thus, we have adapted to the supergroup context the notion of real form, as it is defined in the Hopf algebra setting.
Now, we turn to the definition of $`SL(n,m,)`$. First, we recall the definition of the complex superbialgebra of formal power series $`[x_{ij}],i,j=1,\mathrm{},n+m`$. The coproduct and counit are defined on the generators by :
$$\mathrm{\Delta }(x_{ij})=1x_{ij}+\underset{k=1}{\overset{n+m}{}}x_{ij}1+x_{ik}x_{kj},$$
$`(7)`$
$$ϵ(x_{ij})=0.$$
$`(8)`$
Moreover, it is also enlightening to evaluate the coproduct of the elements $`y_{ij}=\delta _{ij}+x_{ij}`$. We have:
$$\mathrm{\Delta }(y_{ij})=\underset{k=1}{\overset{n+m}{}}y_{ik}y_{kj}.$$
$`(8bis)`$
These maps are defined on all elements of $`[x_{ij}]`$ by the morphism property of $`\mathrm{\Delta },ϵ`$. The gradation of the generators $`x_{ij}`$ is $`|x_{ij}|=|i|+|j|`$ where $`|i|=0,|j|=1`$ for respectively $`i=1\mathrm{}n,j=n+1\mathrm{}n+m`$. We have the following standard Grassmann rules $`x_{ij}x_{mn}=(1)^{(|i|+|j|)(|m|+|n|)}x_{mn}x_{ij}`$ for the product in $`[x_{ij}]`$.
###### Definition 2.3
The complex Lie supergroup $`SL(n,m,)`$ or better
$`Hol(SL(n,m,))`$ is the quotient of the superbialgebra $`[x_{ij}]`$ by the ideal generated by the polynomial <sup>1</sup><sup>1</sup>1 The superdeterminant of a supermatrix $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)`$ is given by $`sdet\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)=\frac{det(ABD^1C)}{det(D)}`$ where $`det`$ is the usual determinant of matrices (see ). $`sdet(\delta +x)1=0`$. The antipode is defined on the quotient by the following superalgebra-antimorphism i.e $`S(x_{ij}x_{mn})=(1)^{(|i|+|j|)(|m|+|n|)}S(x_{mn})S(x_{ij})`$:
$$S(x_{ij})=\delta _{ij}+(\delta +x)_{ij}^1,$$
$`(9)`$
with $`i,j=1\mathrm{}n+m`$.
Here $`(\delta +x)^1`$ means the inverse of the supermatrices which have for elements at the row $`i`$ and the column $`j`$: $`(\delta _{ij}+x_{ij})`$. The definition of the inverse of a supermatrix $`\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)`$ reads (see ):
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)^1=\left(\begin{array}{cc}(ABD^1C)^1& A^1B(DCA^1B)^1\\ D^1C(ABD^1C)^1& (DCA^1B)^1\end{array}\right),$$
$`(10)`$
when $`A,D`$ are even invertible matrices and nothing is assumed for the odd matrices $`B,C`$.
###### Remark 2.2
We shall use the notation $`SL(n,m,)`$ and $`Hol(SL(n,m,))`$ interchangeably. In particular, we shall adopt the latter notation when we want to stress that in the super setting we deal with the holomorphic ”functions” on the supergroup. We frenquently write ”functions” in inverted com mas, the reason is that, in fact, we are working with formal series on the Lie supergroup. But, we think that the reader will understand this abuse of notations in the sense that a lot of ideas of this article are more natural in thinking about it as if we are working on functions on some space.
Now we follow our paper and we equip $`SL(n,m,)`$ with a graded real form as follows:
###### Theorem 2.1
The even antilinear superalgebra-morphism:
$$\sigma (x_{ij})=(1)^{(|i|+|j|)|j|}S(x_{ji})$$
$`(11)`$
introduce the structure of graded real form on $`SL(n,m,)`$ in the sense of Definition 2.2.
Proof:
In fact, it is enough to prove the properties (1) to (6) just for the generators $`x_{ij}`$ because of the property of superalgebra-morphism of $`\sigma `$. Thus, the property (4) is clearly fullfilled. The antilinearity (property (3)) stems also from the definition of $`\sigma `$.
We develop respectively the expressions $`\sigma \sigma (\mathrm{\Delta }(x_{ij}))`$ and $`\mathrm{\Delta }(\sigma (x_{ij}))`$ by using the equations (7) and (11). Thereby, we obtain respectively:
$`\sigma \sigma (\mathrm{\Delta }(x_{ij}))`$ $`=`$ $`1S(x_{ji})(1)^{(|i|+|j|)|j|}+(1)^{(|i|+|j|)|j|}S(x_{ji})1`$
$`+(1)^{(|i|+|k|)|k|+(|j|+|k|)|j|}S(x_{ki})S(x_{jk}),`$
and
$`\mathrm{\Delta }(\sigma (x_{ij}))`$ $`=`$ $`(1)^{(|i|+|j|)|j|}\mathrm{\Delta }(S(x_{ji}))`$
$`=`$ $`(1)^{(|i|+|j|)|j|}\tau (SS(1x_{ji}+x_{ji}1+x_{jk}x_{ki}))`$
$`=`$ $`(1)^{(|i|+|j|)|j|}1S(x_{ji})+(1)^{((|i|+|j|)|j|)}S(x_{ji})1+`$
$`(1)^{|j|(|j|+|k|)+(|i|+|k|)|k|}S(x_{ki})S(x_{jk}).`$
Here $`\tau `$ is the flip with the following property $`\tau (fg)=(1)^{|f||g|}gf`$. Hence, $`\sigma \sigma \mathrm{\Delta }(x_{ij})=\mathrm{\Delta }(\sigma (x_{ij}))`$, which corresponds to the property (1).
Now we turn to the property (2):
$$ϵ(\sigma (x_{ij}))=ϵ((1)^{(|i|+|j|)|j|}S(x_{ji}))=(1)^{(|i|+|j|)|j|}ϵ(x_{ij})=0=\overline{ϵ(x_{ij})}.$$
Moreover, we have the property (5b) because:
$`S(\sigma (S(\sigma (x_{ij}))))`$ $`=`$ $`S(\sigma (S(S(x_{ji}))))(1)^{(|i|+|j|)|j|}`$
$`=`$ $`S(\sigma (x_{ji}))(1)^{(|i|+|j|)|j|}`$
$`=`$ $`S(S(x_{ij}))(1)^{(|i|+|j|)|j|+(|i|+|j|)|i|}`$
$`=`$ $`(1)^{|i|+|j|}x_{ij}.`$
Here we have use the property $`S(S(x))=x,xSL(n,m,)`$.
Finally, it remains to prove $`\sigma (\sigma (x_{ij}))=(1)^{|i|+|j|}x_{ij}`$. We need to evaluate $`\sigma (S(x_{ji}))`$. For this we first develop the following identity:
$$\delta _{ij}=\sigma ((\delta _{ik}+x_{ik})(\delta _{kj}+S(x_{kj}))),$$
thus from (11) and the antilinearity of $`\sigma `$ we deduce:
$$\delta _{ij}(1)^{|j|}=(\delta _{kj}+\sigma (S(x_{kj})))(\delta _{ik}+S(x_{ki}))(1)^{|k||j|}$$
Next, we multiply each menbers of the last equation by the right by $`(\delta _{ip}+x_{ip})`$ and we obtain:
$$(\delta _{kj}+\sigma (S(x_{kj})))\delta _{kp}(1)^{|k||j|}=(1)^{|j|}(\delta _{jp}+x_{jp}).$$
We deduce:
$$\sigma (S(x_{pj}))=(1)^{(|p|+|j|)|j|}x_{jp}.$$
Thus, from this last identity we have:
$$\sigma (\sigma (x_{ij}))=\sigma (S(x_{ji}))(1)^{(|i|+|j|)|j|}=(1)^{(|i|+|j|)|j|+(|i|+|j|)|i|}x_{ij}=(1)^{|i|+|j|}x_{ij},$$
which ends up the proof. $`\mathrm{}`$
###### Remark 2.3
If $`m=0`$ (i.e there are no odd genererators), our definition reduces to the standard compact real form $`SU(n,)`$ of $`SL(n,)`$. For this reason, we call the pair $`(Hol(SL(n,m,),\sigma ))`$ the compact graded real form of $`SL(n,m,)`$. With a slight ambiguity of notations, we note the compact graded real form of $`SL(n,m,)`$ by $`SU(n,m)`$.
## 3 Iwasawa decomposition of $`SL(n,m,)`$
Here, we turn to the definition of the main actors which occur in the Iwasawa decomposition of $`SL(n,m,)`$. We first recall the Iwasawa decomposition (see ) in the non-supercase. When $`(m=0)`$, the Iwasawa decomposition is the statement that the group $`SL(n,)`$ viewed as a real group can be decomposed as $`SL(n,)=GAN`$ <sup>2</sup><sup>2</sup>2More precisely, it exists two unique maps $`g,b`$ from $`SL(n,)`$ in (respectively) $`SU(n,),AN`$ such that for all $`dSL(n,)`$ we have $`d=g(d)b(d)`$. where $`G=SU(n,)`$ is the compact real form of $`SL(n,)`$ and $`AN`$ is the real subgroup of $`SL(n,)`$ of upper triangular matrices with real positive elements on the diagonal and determinant equal to one. In the supercase, we need to work with dual objects i.e. ”functions” on the supergroup. So we need: 1) the Hopf superalgebra $`Fun(SL(n,m,))`$ of ”functions” on $`SL(n,m,)`$ viewed as a real supergroup, 2) the Hopf superalgebra $`Fun(SU(n,m))`$ of ”functions” on $`SU(n,m)`$ the compact graded real form of $`SL(n,m,)`$ and 3) the Hopf superalgebra $`Fun(s(AN))`$ of ”functions” on $`s(AN)`$ which is the superization of the previous Lie group $`AN`$. We discuss the three ingredients separately.
1) For finding $`Fun(SL(n,m,))`$ we borrow some inspiration from the paper where the (non-super) q-analogue of the Iwasawa decomposition was considered. Thus the space of complex ”functions” on the real supergroup $`SL(n,m,)`$ is $`Hol(SL(n,m,))_{}\overline{Hol}(SL(n,m,))`$. $`\overline{Hol}(SL(n,m,))`$ is a ”copy” of $`Hol(SL(n,m,))`$ where the generators $`x_{ij}`$ are named $`x_{ij}^{}`$. The first copy of the tensor product corresponds to ”holomorphic functions” on $`SL(n,m,)`$ while the second copy to ”antiholomorphic functions” on $`SL(n,m,)`$. Note that we use the notation $`Hol(SL(n,m,))`$ for $`SL(n,m,)`$ viewed as the complex supergroup and the notation $`Fun(SL`$
$`(n,m,))`$ for $`SL(n,m,)`$ viewed as the real supergroup. So, we have the following definition:
###### Definition 3.1
The space of complex ”functions” $`Fun(SL(n,m,))`$ on the real supergroup $`SL(n,m,)`$ is the Hopf superalgebra:
$$Fun(SL(n,m,))=Hol(SL(n,m,))_{}\overline{Hol}(SL(n,m,)).$$
$`(12)`$
2) The construction of the Hopf superalgebra of ”functions” on the compact graded real form of the supergroup $`OSp(2r,s,)`$ was performed in detail in . Here we adapt this construction to the compact graded real form of $`SL(n,m,)`$, thud the Hopf superalgebra $`Fun(SU(n,m))`$ is the quotient: $`Hol(SL(n,m,))\overline{Hol}(SL(n,m,))/I`$. Here $`I`$ is the superideal generated by the polynomial equations $`\sigma (x_{ij})x_{ij}^{}=0`$ ($`\sigma `$ is the map of the theorem 2.1). The fact that $`I`$ is a Hopf superideal is a consequence of the relations (1-6) and the properties of the antipode for $`Fun(SL(n,m,))`$. Thus we have the definition:
###### Definition 3.2
The space of complex ”functions” on the compact graded real form $`SU(n,m)`$ of $`SL(n,m,)`$ is the following Hopf superalgebra:
$$Fun(SU(n,m))=Fun(SL(n,m,))/I.$$
$`(13)`$
Here $`I`$ is the Hopf superideal generated by the polynomial equations:
$$\sigma (x_{ij})x_{ij}^{}=0,$$
$`(14)`$
with $`\sigma (x_{ij})=(1)^{(|i|+|j|)|j|}S(x_{ji}).`$
3) Finaly, the definition of $`Fun(s(AN))`$ reads:
###### Definition 3.3
The superalgebra $`Fun(s(AN))`$ of complex ”functions” on the real supergroup $`s(AN)`$ is the following Hopf superalgebra:
$$Fun(s(AN))=Fun(SL(n,m,))/J.$$
Here, $`J`$ is the Hopf superideal generated by the following polynomials relations:
$$x_{ij}=x_{ij}^{}=0\mathrm{for}i>j,$$
$`(15)`$
$$x_{ii}x_{ii}^{}=0.$$
$`(16)`$
###### Remark 3.1
For $`m=0`$, the Hopf superalgebra $`Fun(s(AN))`$ reduces to the Hopf algebra $`Fun(AN)`$ of ”functions” on the Lie group $`AN`$. If $`m0`$, the fact that $`Fun(s(AN))`$ is a good definition follows from the fact that the super Iwasawa decomposition can be formulated with it.
###### Remark 3.2
As $`Fun(s(AN))`$ and $`Fun(SU(n,m))`$ are both quotients of $`Fun(SL(n,m,))`$, we have the canonical projections $`i`$ and $`j`$ which are the superalgebra-morphisms:
$$i:Fun(SL(n,m,))Fun(SU(n,m)),$$
$`(18)`$
$$j:Fun(SL(n,m,))Fun(s(AN)).$$
$`(19)`$
Both $`i`$ and $`j`$ map an element $`fFun(SL(n,m,))`$ into its respective cosets.
Now we turn to the main theorem of this article, the Iwasawa decomposition of the real supergroup $`Fun(SL(n,m,)`$:
###### Theorem 3.1
(Iwasawa decomposition)
There exist a unique pair ($`\varphi `$,$`\psi `$) of superalgebra-morphisms:
$$\varphi :Fun(SU(n,m))Fun(SL(n,m,))$$
$$\psi :Fun(s(AN))Fun(SL(n,m,))$$
such that:
$$\varphi (i(f_{(1)})).\psi (j(f_{(2)}))=f,fSL(n,m,).$$
$`(17)`$
Here $`\mathrm{\Delta }(f)=f_{(1)}f_{(2)}`$ and . is the standard commutative multiplication in the superalgebra $`Fun(SL(n,m,))`$.
###### Remark 3.3
This theorem in the non-super case (m=0) is the dualisation of the Iwasawa decomposition of the Lie group $`SL(n,)`$ i.e. it gives the Iwasawa decomposition on the space of complex ”functions” on $`SL(n,)`$. When $`m=0`$, note that the maps $`\varphi ,\psi `$ are, respectively, the pullbacks of the maps $`g,b`$ defined in the footnote 2 for the Lie group $`SL(n,)`$.
Before giving the proof of the theorem, we have to introduce some notations. We said that a column supervector $`X`$:
$$X=\left(\begin{array}{c}X_1\\ \mathrm{}\\ X_n\\ \chi _1\\ \mathrm{}\\ \chi _m\end{array}\right),$$
$`(20)`$
is even (odd) when $`X_i`$ are even (odd) and $`\chi _m`$ are odd (even). We define the supertranspose of the supervector $`X`$ by:
$$X^{st}=\left(\begin{array}{cccccc}(1)^{|X|}X_1& \mathrm{}& (1)^{|X|}X_n& \chi _1& \mathrm{}& \chi _m\end{array}\right).$$
$`(21)`$
We use also the notation:
$$X^{}=\left(\begin{array}{c}X_1^{}\\ \mathrm{}\\ X_n^{}\\ \chi _1^{}\\ \mathrm{}\\ \chi _m^{}\end{array}\right).$$
$`(22)`$
We give the definition of the scalar product of two supervectors $`X`$, $`Y`$:
$$(X,Y)=X^{^{st}}Y.$$
$`(23)`$
This scalar product have the following properties:
$$(X,Y)^{}=(1)^{(|X|+|Y|)|Y|}(Y,X),$$
$`(24)`$
$$(X\lambda ,Y)=(1)^{(|X|+1)|\lambda |}\lambda ^{}(X,Y),$$
$`(25)`$
where $`\lambda `$ are possible odd or even polynoms of the generators $`x_{ij},x_{ij}^{}`$. Moreover, the norm of a supervector $`X`$ is noted and defined by $`X=\sqrt{(X,X)}`$.
We define a $`(m+n)\times (m+n)`$ supermatrix $`P`$ by specifying either its entries or its column supervectors. In the first case, we note $`p_{ij}`$ the entries of the supermatrix at the ith row and jth column. In the second case, $`P`$ reads:
$$P=(P_1,\mathrm{},P_{n+m}),$$
$`(26)`$
where
$$P_i=\left(\begin{array}{c}p_{1i}\\ \mathrm{}\\ p_{ni}\\ p_{n+1i}\\ \mathrm{}\\ p_{m+ni}\end{array}\right).$$
$`(27)`$
Finaly, we end this sequence of notations with two definitions:
###### Definition 3.4
We say that a supermatrix $`P`$ with a unit superdeterminant is a $`SU(n,m)`$-supermatrix if its diagonal elements are normalized <sup>3</sup><sup>3</sup>3 A normalized formal serie is a formal serie where monomial of degree zero is equal to one. formal series and $`P`$ fullfills <sup>4</sup><sup>4</sup>4The supertranspose of a supermatrix $`N=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)`$ is $`N^{st}=\left(\begin{array}{cc}A^t& C^t\\ B^t& D^t\end{array}\right)`$, where $`A^t,B^t,C^t,D^t`$ is the usual transposition of the matrice $`A,B,C,D`$. Moreover, the entries of $`N^{}`$ at the ith row and jth column are $`n_{ij}^{}`$. $`P^{^{st}}P=1`$.
###### Definition 3.5
We say also that a supermatrix $`Q`$ with a unit superdeterminant is a $`s(AN)`$-supermatrix if its diagonal elements are real <sup>5</sup><sup>5</sup>5For the conjugation $``$, an element is real when it is equal to its conjugate. normalized formal series and $`Q`$ is an upper triangular supermatrix.
Proof of the theorem 3.1:
The proof of the theorem gets organized in two parts. In the first part, we explicitely describe the two superalgebra-morphisms $`\varphi `$ and $`\psi `$. Next, we show that they fullfill the property (17). In the second part, in turn, we show the unicity of these maps.
Let $`\stackrel{~}{\varphi },\stackrel{~}{\psi }`$ be two superalgebra-automorphisms of $`Fun(SL(n,m,))`$. Consider supermatrices $`M,\mathrm{\Phi },\mathrm{\Psi }`$ with elements $`\delta _{ij}+x_{ij}`$, $`\stackrel{~}{\varphi }(\delta _{ij}+x_{ij})`$ and $`\stackrel{~}{\psi }(\delta _{ij}+x_{ij})`$, respectively. We know from definition 3.1 that $`sdet(M)=1`$. Denote $`M_l,\mathrm{\Phi }_l,\mathrm{\Psi }_l`$ the columns of $`M,\mathrm{\Phi },\mathrm{\Psi }`$ , respectively. Set
$$\mathrm{\Phi }_l=\frac{V_l}{V_l},$$
$`(28)`$
where the supervectors $`V_l`$ are defined recursively by:
$$V_1=M_1,V_l=M_l\underset{k=1}{\overset{l1}{}}V_k\frac{(V_k,M_l)}{(V_k,V_k)},l=2,\mathrm{},n+m.$$
$`(29)`$
It is easy to see that (28) imply that $`\mathrm{\Phi }`$ is a $`SU(n,m)`$-supermatrix (cf. Definition 3.4), hence $`\mathrm{\Phi }`$ is invertible. Then, we also set:
$$\mathrm{\Psi }=\mathrm{\Phi }^1M.$$
$`(30)`$
We deduce from (30) that $`\mathrm{\Psi }`$ is a $`s(AN)`$-supermatrix (cf. Definition 3.5).
The fact that $`\mathrm{\Phi }`$ is a $`SU(n,m)`$-supermatrix imply that $`\stackrel{~}{\varphi }(I)=0`$ ($`I`$ was defined in definition 3.2). Hence, $`\stackrel{~}{\varphi }`$ gives rise to a superalgebra-morphism $`\varphi `$ from $`Fun(SU(n,m))`$ to $`Fun(SL(n,m,))`$ by:
$$\varphi (i(\delta _{ij}+x_{ij}))=\mathrm{\Phi }_{ij},\varphi (i(\delta _{ij}+x_{ij}^{}))=\mathrm{\Phi }_{ij}^{}.$$
$`(31)`$
The fact that $`\mathrm{\Psi }`$ is a $`s(AN)`$-supermatrix imply that $`\stackrel{~}{\psi }(J)=0`$ ($`J`$ was defined in definition 3.3). Hence, $`\stackrel{~}{\psi }`$ gives rise to a superalgebra-morphism $`\psi `$ from $`Fun(s(AN))`$ to $`Fun(SL(n,m,))`$ by:
$$\psi (i(\delta _{ij}+x_{ij}))=\mathrm{\Psi }_{ij},\psi (i(\delta _{ij}+x_{ij}^{}))=\mathrm{\Psi }_{ij}^{}.$$
$`(32)`$
Furthermore, the fact that $`\varphi ,\psi ,i,j`$ are superalgebra-morphisms makes sufficient to prove (17) just for the generators $`x_{ij}`$ and $`x_{ij}^{}`$. Now from eq.(30) we deduce $`M=\mathrm{\Phi }\mathrm{\Psi }`$. Hence, we have also $`M^{}=\mathrm{\Phi }^{}\mathrm{\Psi }^{}`$. These two last equalities give directly the validity of (17) for the morphisms $`(\varphi ,\psi )`$ defined by (31), (32). Thus the existence of $`(\varphi ,\psi )`$ is proved.
The reader may wish to understand better the origin of the formulas (28), (29). In fact, it is a superanalogue of the Gram-Schimdt procedure. We start with the family of column supervectors $`M_1,\mathrm{},M_{n+m}`$. In $`(n+m)`$-steps we construct a familly $`V_1,\mathrm{},V_{n+m}`$ of orthogonal supervectors. More precisely, in the $`k`$-th step of the recursion we modify the $`k`$-th column in such a way that it becomes orthogonal to the $`k1`$ previous columns. Finaly, we obtain (30) by the normalisation of the orthogonal family $`V_1,\mathrm{},V_{n+m}`$.
Now, we turn to the unicity of the maps $`\stackrel{~}{\varphi },\stackrel{~}{\psi }`$. We assume that there exist two distincts pairs of superalgebra-automorphisms of $`Fun(SL(n,m,))`$ $`(\stackrel{~}{\varphi }_k,\stackrel{~}{\psi }_k),k=1,2`$ verifying $`\stackrel{~}{\varphi }_k(I)=0,\stackrel{~}{\psi }_k(J)=0`$ and fullfilling (17). They give rise to two pairs of supermatrices $`(\mathrm{\Phi }_k,\mathrm{\Psi }_k)`$ defined for $`k=1,2`$ as follows:
$$\stackrel{~}{\varphi }_k(\delta _{ij}+x_{ij})=(\mathrm{\Phi }_k)_{ij},\stackrel{~}{\varphi }_k(\delta _{ij}+x_{ij}^{})=(\mathrm{\Phi }_k)_{ij}^{},$$
$`(33)`$
$$\stackrel{~}{\psi }_k(\delta _{ij}+x_{ij})=(\mathrm{\Psi }_k)_{ij},\stackrel{~}{\psi }_k(\delta _{ij}+x_{ij}^{})=(\mathrm{\Psi }_k)_{ij}^{}.$$
$`(34)`$
Firstly, from the fact that $`\stackrel{~}{\varphi }_k(I)=0`$ (resp. $`\stackrel{~}{\psi }_k(J)=0`$) we deduce that the supermatrix $`\mathrm{\Phi }_k`$ (resp. $`\mathrm{\Psi }_k`$) are $`SU(n,m)`$-supermatrix (resp. $`s(AN)`$-supermatrix). Secondly, we have the following equality between supermatrices:
$$\mathrm{\Phi }_1\mathrm{\Psi }_1=\mathrm{\Phi }_2\mathrm{\Psi }_2,$$
$`(35)`$
because $`\stackrel{~}{\varphi }_k,\stackrel{~}{\psi }_k`$ fullfill the relation (17). So we obtain:
$$\mathrm{\Phi }_2^1\mathrm{\Phi }_1=\mathrm{\Psi }_2\mathrm{\Psi }_1^1.$$
$`(36)`$
Finally, we remark that $`\mathrm{\Phi }_2^1\mathrm{\Phi }_1`$ is a $`SU(n,m)`$-supermatrix as a matricial product of $`SU(n,m)`$-supermatrices, whereas $`\mathrm{\Psi }_2\mathrm{\Psi }_1^1`$ is a $`s(AN)`$-supermatrix as a matricial product of $`s(AN)`$-supermatrices. The unique supermatrix which is both a $`SU(n,m)`$-supermatrix and $`s(AN)`$-supermatrix is the unit supermatrix. Hence, we deduce that:
$$\mathrm{\Phi }_1=\mathrm{\Phi }_2,\mathrm{\Psi }_1=\mathrm{\Psi }_2.$$
$`(37)`$
The unicity is therefore proved. $`\mathrm{}`$
## 4 References
F.A.Berezin, Introduction to superanalysis, edited by A.A.Kirillov, MPAM, Reidel Publishing Company, Holland (1984).
Anthony.W. Knapp, Representation theory of semisimple groups an overview based on examples, Princeton University Press, Princeton mathematical series (1986).
F.Pellegrini, Grassmann real form of $`OSp(2r,s,)`$, math.RA/0311240.
P. Podleś and S.L. Woronowicz, Quantum Deformation of Lorentz Group. Commun. Math. Phys. 130, 381-431 (1990).
V.V.Serganova, Classification of real simple Lie superalgebras and symmetric superspaces, Functional Analysis 17 n3 (July-September 1983) 46–54.
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# Dynamics of Dipolar Spinor Condensates
## I Introduction
One of the current experimental advances of considerable importance in the context of ultracold atoms is the demonstration of multicomponent Bose-Einstein condensate (BEC). In particular several groups have created spinor condensates of <sup>23</sup>Na atoms 1 and <sup>87</sup>Rb atoms 2 by transferring spin polarized condensates into far-off-resonance optical dipole traps, where the spin degree of freedom becomes active and the spinor nature of the condensates is manifested. The spin-mixing dynamics and its dependence on the magnetic field have been investigated in detail for both spin-1 and spin-2 condensates in a number of theoretical 3 and experimental works 1 ; 2 . Since the spin degree of freedom becomes accessible in an optical trap, the magnetic dipole-dipole interactions (MDDI) which arise from intrinsic or induced field of magnetic dipole moment 4 ; 5 should be taken into account in addition to the isotropic $`s`$-wave contact interaction. Due to their long-range and vectorial characters, the MDDI may largely enrich the variety of phenomena observed in condensates. Therefore this opens up new directions of research, such as the ground-state structure of spin-1 dipolar condensates in a single trap 6 as well as the ground-state magnetic properties of spinor BEC confined in deep optical lattices 7 ; 8 ; 9 , where only the dipolar interactions between condensates at different sites are considered. Recently, a beautiful experiment in Stuttgart 10 has demonstrated BEC in a gas of chromium atoms $`{}_{}{}^{52}Cr`$, the MDDI of which is a factor of 36 higher than that for alkali atoms. This achievement makes possible the studies of the anisotropic long-range interactions in degenerate quantum gases.
The present paper exploits semiclassical dynamics of spin-1 dipolar condensates in a single trap and shows the time evolutions of the population imbalance and the relative phase among different spin components 6 ; 11 . This study may help us gain some insights into the properties of this dipolar spinor condensate, such as quantum phase diffusion, spontaneous magnetization 12 and Macroscopic Quantum Self Trapping (MQST). Especially, MQST known as a novel nonlinear effect has already been predicted theoretically 13 and very recently observed experimentally in double-well trap 14 and periodic optical lattice 15 . The self trapping effect sustains a self-maintained interwell population imbalance during the nonlinear tunneling process. Here we will show that for a dipolar condensate the interplay between spin-exchange and MDDI gives rise to the intercomponent MQST naturally. Moreover, due to the adjustability of the MDDI, the experimental observation of the MQST effect in dipolar spinor condensates can be expected.
The paper is organized as follows. In Sec. II we describe the dipolar spinor BEC model and derive the equations of motion of semiclassical dynamics in terms of symmetry. Sec. III is devoted to the properties of the equations of motion such as spontaneous magnetization and spin-mixing dynamics. The time evolutions of the population imbalance and the relative phase among different spin components are investigated and quantum tunneling and self trapping among different spin components are examined carefully. Finally, a brief summary is given in Sec. IV.
## II The Model and the Equations of Motion
We consider a spin $`F=1`$ dipolar condensate with $`N`$ bosons. In second quantized form, the Hamiltonian $`\widehat{H}_{tot}`$ subject to both spin-exchanging collisions $`\widehat{H}_{sp}`$ and MDDI $`\widehat{H}_{dd}`$ reads 6
$$\widehat{H}_{tot}=\widehat{H}_{sp}+\widehat{H}_{dd}$$
with
$`\widehat{H}_{sp}`$ $`={\displaystyle 𝑑𝐫\widehat{\psi }_\alpha ^{}\left(𝐫\right)\left[\frac{^2}{2M}+V_T\left(𝐫\right)\right]\widehat{\psi }_\alpha \left(𝐫\right)}`$
$`+{\displaystyle \frac{c_0}{2}}{\displaystyle 𝑑𝐫\widehat{\psi }_\alpha ^{}\left(𝐫\right)\widehat{\psi }_\beta ^{}\left(𝐫\right)\widehat{\psi }_\alpha \left(𝐫\right)\widehat{\psi }_\beta \left(𝐫\right)}`$
$`+{\displaystyle \frac{c_2}{2}}{\displaystyle 𝑑𝐫\widehat{\psi }_\alpha ^{}\left(𝐫\right)\widehat{\psi }_\beta ^{}\left(𝐫\right)𝐅_{\alpha \alpha ^{}}𝐅_{\beta \beta ^{}}\widehat{\psi }_\beta ^{}\left(𝐫\right)\widehat{\psi }_\alpha ^{}\left(𝐫\right)}`$ (1)
and
$`\widehat{H}_{dd}`$ $`={\displaystyle \frac{c_d}{2}}{\displaystyle 𝑑𝐫𝑑𝐫^{}\frac{1}{\left|𝐫𝐫^{}\right|^3}}`$
$`\times [\widehat{\psi }_\alpha ^{}\left(𝐫\right)\widehat{\psi }_\beta ^{}\left(𝐫^{}\right)𝐅_{\alpha \alpha ^{}}𝐅_{\beta \beta ^{}}\widehat{\psi }_\beta ^{}\left(𝐫\right)\widehat{\psi }_\alpha ^{}\left(𝐫^{}\right)`$
$`3\widehat{\psi }_\alpha ^{}\left(𝐫\right)\widehat{\psi }_\beta ^{}\left(𝐫^{}\right)(𝐅_{\alpha \alpha ^{}}𝐞)(𝐅_{\beta \beta ^{}}𝐞)\widehat{\psi }_\beta ^{}\left(𝐫\right)\widehat{\psi }_\alpha ^{}\left(𝐫^{}\right)],`$ (2)
where $`\widehat{\psi }_\alpha \left(𝐫\right)`$ $`\left(\alpha =0,\pm 1\right)`$ are the field annihilation operators for an atom in the hyperfine state $`|F=1,m_F=\alpha .`$ The hyperfine spin $`𝐅`$ of one single atom is expressed in the spin-1 matrix representation, and $`M`$ is the mass of the atom. A summation over repeated indices is assumed in Eqs. (1) and (2). The external trapping potential $`V_T\left(𝐫\right)`$ is spin independent for a far off-resonant optical dipole trap which makes atomic spin degree of freedom completely accessible. The two coefficients $`c_0=4\pi \mathrm{}^2\left(a_0+2a_2\right)/3M`$, $`c_2=4\pi \mathrm{}^2\left(a_2a_0\right)/3M`$ characterize the density-density and spin-spin collisional interactions, respectively, with $`a_f`$ $`\left(f=0,2\right)`$ being the $`s`$-wave scattering length for two spin-1 atoms in the combined symmetric channel of total spin $`f.`$ And the dipolar interaction parameter is $`c_d=\mu _0\mu _B^2g_F^2/4\pi `$ with $`g_F`$ being the Landé $`g`$-factor, $`\mu _B`$ the Bohr magneton. $`𝐞=\left(𝐫𝐫^{}\right)/\left|𝐫𝐫^{}\right|`$ is a unit vector. For the two experimentally created spinor condensates (<sup>23</sup>Na and <sup>87</sup>Rb), we have $`\left|c_2\right|c_0.`$
For a condensate with its symmetry axis chosen to be along the quantization axis $`z`$ (which happens to be the most experimentally relevant cases), $`\widehat{H}_{dd}`$ takes a very simple form as shown in Ref. 6 . Under the single mode approximation (SMA), namely, $`\widehat{\psi }_\alpha \left(𝐫\right)\varphi \left(𝐫\right)\widehat{a}_\alpha ,`$ where $`\varphi \left(𝐫\right)`$ is the spin-independent condensate spatial wave function, $`\widehat{a}_\alpha `$ is the annihilation operator for $`m_F=\alpha `$ component, the total Hamiltonian apart from a trivial term is given as
$`\widehat{H}_{tot}`$ $`=\widehat{H}_{sp}+\widehat{H}_{dd}=\epsilon {\displaystyle \underset{\alpha }{}}\widehat{a}_\alpha ^{}\widehat{a}_\alpha +U_0{\displaystyle \underset{\alpha ,\beta }{}}\widehat{a}_\alpha ^{}\widehat{a}_\beta ^{}\widehat{a}_\beta \widehat{a}_\alpha `$
$`+U_2(\widehat{a}_1^{}\widehat{a}_1^{}\widehat{a}_1\widehat{a}_1+\widehat{a}_1^{}\widehat{a}_1^{}\widehat{a}_1\widehat{a}_12\widehat{a}_1^{}\widehat{a}_1^{}\widehat{a}_1\widehat{a}_1`$
$`+2\widehat{a}_1^{}\widehat{a}_0^{}\widehat{a}_1\widehat{a}_0+2\widehat{a}_1^{}\widehat{a}_0^{}\widehat{a}_1\widehat{a}_0+2\widehat{a}_0^{}\widehat{a}_0^{}\widehat{a}_1\widehat{a}_1+2\widehat{a}_1^{}\widehat{a}_1^{}\widehat{a}_0\widehat{a}_0)`$
$`+U_d(2\widehat{a}_1^{}\widehat{a}_1^{}\widehat{a}_1\widehat{a}_1+2\widehat{a}_1^{}\widehat{a}_1^{}\widehat{a}_1\widehat{a}_14\widehat{a}_1^{}\widehat{a}_1^{}\widehat{a}_1\widehat{a}_1`$
$`2\widehat{a}_1^{}\widehat{a}_0^{}\widehat{a}_1\widehat{a}_02\widehat{a}_1^{}\widehat{a}_0^{}\widehat{a}_1\widehat{a}_02\widehat{a}_0^{}\widehat{a}_0^{}\widehat{a}_1\widehat{a}_1`$
$`2\widehat{a}_1^{}\widehat{a}_1^{}\widehat{a}_0\widehat{a}_0+\widehat{a}_1^{}\widehat{a}_1+\widehat{a}_1^{}\widehat{a}_12\widehat{a}_0^{}\widehat{a}_0),`$ (3)
where $`\epsilon =d𝐫\varphi ^{}\left(𝐫\right)[^2/2M+V_T\left(𝐫\right)]\varphi \left(𝐫\right)`$ is assumed to be of the same value for all components. $`U_{0,2}=\left(c_{0,2}/2\right)\left|\varphi \left(𝐫\right)\right|^4𝑑𝐫`$ characterizes the spin-changing collisions and $`U_d=`$ $`\left(c_d/4\right)𝑑𝐫𝑑𝐫^{}\left|\varphi \left(𝐫\right)\right|^2\left|\varphi \left(𝐫^{}\right)\right|^2\left(13\mathrm{cos}^2\theta _e\right)/\left|𝐫𝐫^{}\right|^3`$ denotes the MDDI with $`\theta _e`$ being the polar angle of $`\left(𝐫𝐫^{}\right)`$. Before we proceed to examine the semiclassical dynamics of $`\widehat{H}_{tot},`$ we would like to emphasize that $`\widehat{H}_{sp}`$ is invariant under spin rotation, i.e., it possesses a $`SO(3)`$ symmetry in spin space 16 . The presence of the MDDI breaks this symmetry into $`SO(2)`$, which means $`\widehat{H}_{tot}`$ has only axial symmetry in spin space as we chosen the quantization axis along $`z`$. In this sense the MDDI plays a similar role to that of an external magnetic field.
In the mean-field theory the condensate is usually considered to be in a coherent state 11
$$|\stackrel{}{z}=\mathrm{exp}\left(\frac{1}{2}\underset{\alpha }{}\left|z_\alpha \right|^2\right)\mathrm{exp}\left(\underset{\alpha }{}z_\alpha \widehat{a}_\alpha ^{}\right)|0,$$
(4)
where $`|0`$ is the vacuum state. The complex numbers $`z_\alpha `$ are nothing but the macroscopic wave functions for the atoms in the hyperfine level $`|F=1,m_F=\alpha `$ with population $`N_\alpha `$ and phase $`\theta _\alpha ,`$ i.e.,
$$z_\alpha =\sqrt{N_\alpha }e^{i\theta _\alpha }.$$
(5)
The time-dependent variational principle of the system
$$\delta S=\delta i\mathrm{}\stackrel{}{z}|\stackrel{.}{\stackrel{}{z}}\stackrel{}{z}|H|\stackrel{}{z}dt=0$$
(6)
gives rise to the Hamiltonian equations of motion in complex coordinates
$$i\mathrm{}\stackrel{}{z}_\alpha =\frac{H_0}{z_\alpha ^{}},\text{ }i\mathrm{}\stackrel{}{z}_\alpha ^{}=\frac{H_0}{z_\alpha }$$
(7)
with
$`H_0(N_\alpha ,\theta _\alpha )`$ $`=\left(\epsilon +U_d\right)N3U_dN_0+U_0N^2+\left(U_2+2U_d\right)\left(N_1N_1\right)^2`$
$`+2\left(U_2U_d\right)N_0[N_1+N_1+2\sqrt{N_1N_1}\mathrm{cos}\left(2\theta _0\theta _1\theta _1\right)].`$ (8)
Note that the dynamics governed by Hamiltonian (3) conserves the total atom numbers $`N=_\alpha N_\alpha `$ and the total hyperfine spin of the condensate in the direction of the quantization axis $`I=\left(N_1N_1\right)\mathrm{}.`$ In terms of the following canonical variables
$$\begin{array}{c}\phi _1=\frac{\left(\theta _1+\theta _0+\theta _1\right)}{3}\hfill \\ \phi _2=\theta _0\frac{\left(\theta _1+\theta _1\right)}{2}\hfill \\ \phi _3=\theta _1\theta _1\hfill \end{array}\text{ and }\begin{array}{c}\mathrm{\Omega }_1=N_1+N_0+N_1\hfill \\ \mathrm{\Omega }_2=\frac{2}{3}N_0\frac{1}{3}\left(N_1+N_1\right)\hfill \\ \mathrm{\Omega }_3=\frac{\left(N_1N_1\right)}{2}\hfill \end{array}$$
(9)
the Hamiltonian (8) becomes cyclic in the coordinates $`\phi _1`$ and $`\phi _3`$
$`H_0(\mathrm{\Omega }_1,\mathrm{\Omega }_2,\mathrm{\Omega }_3,\phi _2)=\epsilon \mathrm{\Omega }_1+U_0\mathrm{\Omega }_1^23U_d\mathrm{\Omega }_2`$
$`+4\left(U_2+2U_d\right)\mathrm{\Omega }_3^2+2\left(U_2U_d\right)`$
$`\times [({\displaystyle \frac{1}{3}}\mathrm{\Omega }_1+\mathrm{\Omega }_2)({\displaystyle \frac{2}{3}}\mathrm{\Omega }_1\mathrm{\Omega }_2)`$
$`+\sqrt{\left({\displaystyle \frac{2}{3}}\mathrm{\Omega }_1\mathrm{\Omega }_2\right)^24\mathrm{\Omega }_3^2}({\displaystyle \frac{1}{3}}\mathrm{\Omega }_1+\mathrm{\Omega }_2)\mathrm{cos}2\phi _2].`$ (10)
Two important consequences follow as a result of the cyclic coordinates $`\phi _1`$ and $`\phi _3.`$ Firstly the mean value of the total number of atoms $`N=\mathrm{\Omega }_1`$ and that of the total hyperfine spin of the condensate $`I=2\mathrm{}\mathrm{\Omega }_3`$ are constants of motion. The dynamics, on the other hand, involves only one pair of variables $`\{\phi _2,\mathrm{\Omega }_2\}`$ as
$$\dot{\phi }_2=\frac{H_0}{\mathrm{\Omega }_2},\text{ }\dot{\mathrm{\Omega }}_2=\frac{H_0}{\phi _2},$$
Define $`\xi _i=\mathrm{\Omega }_i/N`$ $`\left(i=1,2,3\right)`$ and $`\eta =\sqrt{(\frac{2}{3}\xi _2)^24\xi _3^2}`$, we obtain
$`\mathrm{}\dot{\xi }_2`$ $`=4N\left(U_2U_d\right)\eta \left({\displaystyle \frac{1}{3}}+\xi _2\right)\mathrm{sin}2\phi _2`$
$`\mathrm{}\dot{\phi }_2`$ $`=3U_d+2N\left(U_2U_d\right)\left(2\xi _2{\displaystyle \frac{1}{3}}\right)`$
$`+2N\left(U_2U_d\right){\displaystyle \frac{\left(\frac{2}{3}\xi _2\right)\left(2\xi _2\frac{1}{3}\right)+4\xi _3^2}{\eta }}\mathrm{cos}2\phi _2.`$ (11)
That means not only $`\xi _2`$ but also $`\xi _3`$ are involved in the dynamics.
The variables $`\xi _2`$ and $`\phi _2`$ are canonically conjugate in a classical Hamiltonian
$`H`$ $`={\displaystyle \frac{\left(U_2U_d\right)N}{2}}\left(2\xi _2{\displaystyle \frac{1}{3}}\right)^2`$
$`+2\left(U_2U_d\right)N\eta \left({\displaystyle \frac{1}{3}}+\xi _2\right)\mathrm{cos}2\phi _23U_d\xi _2.`$ (12)
In a simple mechanical analogy, $`H`$ describes a nonrigid pendulum. The Bose Josephson junction (BJJ) tunneling current between different spin components is given by
$$I_{BJJ}=\frac{\dot{\xi }_2N}{2}=I_0\eta \left(\frac{1}{3}+\xi _2\right)\mathrm{sin}2\phi _2,$$
(13)
where $`I_0=2N^2\left(U_2U_d\right)`$. It differs from BJJ tunneling current of two weakly linked BEC in a double-well potential in its further nonlinearity in $`\xi _2`$ 13 .
## III Properties of the Equations of Motion
### III.1 Spontaneous magnetization
We firstly study the equilibrium configurations of the system. They are determined by the classical equations of motion (11) after setting the time derivative terms to zero
$`0`$ $`=\eta \left({\displaystyle \frac{1}{3}}+\xi _2\right)\mathrm{sin}2\phi _2`$ (14)
$`0`$ $`=\mathrm{\Lambda }\pm [2\xi _2{\displaystyle \frac{1}{3}}`$
$`+{\displaystyle \frac{\left(\frac{2}{3}\xi _2\right)\left(2\xi _2\frac{1}{3}\right)+4\xi _3^2}{\eta }}\mathrm{cos}2\phi _2].`$ (15)
We have defined a dimensionless parameter $`\mathrm{\Lambda }=\frac{3U_d}{2\left|U_2U_d\right|N}`$ characterizing the relative interaction strength of the MDDI. In Eq. (15) plus sign corresponds to the case of $`\left(U_2U_d\right)>0,`$ while minus sign corresponds to $`\left(U_2U_d\right)<0`$ (In the following we take the former case as an example). From the definition of $`\eta `$ we know that the constant of motion $`\xi _3`$ takes values in the interval $`\frac{1}{2}<\xi _3<\frac{1}{2}`$ and the dynamic variable $`\xi _2`$ in the interval $`\frac{1}{3}<\xi _2<\frac{2}{3}2|\xi _3|.`$ In the following discussion about the solutions of the equilibrium equations we separately consider three cases in order to see more clearly how these solutions depend on the relative phase $`\phi _2`$.
(1) Equilibrium configuration with $`\mathrm{cos}2\phi _2=1`$ (or $`2\phi _2=2k\pi `$).
The equilibrium value of $`\xi _2`$ is thus given by equation (15) with $`\mathrm{cos}2\phi _2=1,`$ which has only one solution in the interval $`\mathrm{}<\mathrm{\Lambda }<1+\sqrt{1\left(2\xi _3\right)^2}.`$ When $`\mathrm{\Lambda }\mathrm{}`$ the equilibrium value of $`\xi _2`$ approaches its upper boundary $`\frac{2}{3}2\left|\xi _3\right|`$. In this case the fractions of atoms $`n_\alpha =N_\alpha /N`$ ($`\alpha =0,\pm 1`$) occupying three hyperfine states are $`n_1=\left|\xi _3\right|+\xi _3,`$ $`n_0=12\left|\xi _3\right|,`$ $`n_1=\left|\xi _3\right|\xi _3`$ . On the other hand when $`\mathrm{\Lambda }=1+\sqrt{1\left(2\xi _3\right)^2},`$ $`\xi _2`$ is at the lower boundary $`\frac{1}{3},`$ where the fractions of atoms occupying the hyperfine states are $`n_1=\frac{1}{2}\left(1+2\xi _3\right),`$ $`n_0=0,`$ $`n_1=\frac{1}{2}\left(12\xi _3\right).`$ Besides, when $`\xi _3=0`$ the equilibrium value of $`\xi _2`$ is $`\xi _2=\frac{1}{2}\left(\frac{1}{3}\mathrm{\Lambda }\right)`$ in the interval $`2<\mathrm{\Lambda }<2`$ and the occupation fractions of three hyperfine states are respectively $`n_0=\frac{1}{2}\left(1\mathrm{\Lambda }\right),`$ $`n_1=n_1=\frac{1}{4}\left(1+\mathrm{\Lambda }\right).`$
(2) Equilibrium configuration with $`\mathrm{cos}2\phi _2=1`$ (or $`2\phi _2=\left(2k+1\right)\pi `$).
Now the equilibrium value of $`\xi _2`$ is given by equation (15) with $`\mathrm{cos}2\phi _2=1.`$ Similarly it has again only one solution in the interval $`\mathrm{\Lambda }>1\sqrt{1\left(2\xi _3\right)^2}.`$ When $`\mathrm{\Lambda }=`$ $`1\sqrt{1\left(2\xi _3\right)^2},`$ $`\xi _2`$ is at the lower boundary $`\frac{1}{3}.`$ In particular, when $`\xi _3=0`$ the solution only exists for vanished dipole-dipole interaction $`U_d`$.
(3) Equilibrium configuration with $`\mathrm{cos}2\phi _2=0`$ (or Equilibrium configuration which does not depend on the phase $`\phi _2`$).
In this case there exist two solutions which are independent of the value of $`\mathrm{\Lambda }.`$ One solution appears at the lower boundary $`\xi _2=\frac{1}{3}`$ and the occupation fractions are $`n_1=n_1=\frac{1}{2},n_0=0.`$ The other one appears at the upper boundary $`\xi _2=\frac{2}{3}2\left|\xi _3\right|`$ and the occupation fractions are $`n_1=\left|\xi _3\right|+\xi _3,`$ $`n_0=12\left|\xi _3\right|,`$ $`n_1=\left|\xi _3\right|+\xi _3.`$ Interestingly when $`\xi _3=0`$ the occupation fractions become $`n_1=n_1=0`$ and $`n_0=1,`$ corresponding to the “polar” state of the condensate.
The mean field theory predicted a polar ground state of spinor BEC for <sup>23</sup>Na in which the atoms interact antiferromagnetically with each other in the absence of the MDDI 16 . The intrinsic magnetic moment of particles may, however, contribute together to form a Bose-Einstein Ferromagnetism (BEF) 12 . Spinor bosons carry magnetic moments $`\mu `$ and $`M=\mu \left|\psi _1\right|^2`$ represents the magnetization. Spontaneous magnetization means that $`M`$ remains finite even if the external magnetic field $`B`$ vanishes, namely, the ground state of spinor BEC is always ferromagnetic state. Our calculations prefer the latter scenario. Specifically, above results suggest that for most given values of $`\mathrm{\Lambda }`$ and initial relative phase among the three components, when $`B=0`$ and $`U_2>0,`$ we have $`M`$ $`0`$ , i.e., the system is magnetized spontaneously once the dipolar spinor BEC is realized. As mentioned in Ref. 12 , one can provide a direct way of confirming spontaneous magnetization experimentally. Considering that both the sign and the magnitude of the MDDI coefficient $`c_d`$ depend on the geometric shape of the condensate 4 ; 10 ; 17 ; 18 , we can conveniently adjust the trap aspect ratio to manipulate the magnetic property of its ground state.
### III.2 Spin-Mixing Dynamics
A key feature of dipolar spinor BEC is that besides the usual two-body repulsive hard-core interactions, there also exist spin-exchange interaction and MDDI which lead to spin mixing within the condensates. Population can be transferred from one spin state to another under internal nonlinear interactions without the presence of external fields. Insight into the complex dynamics of our system can be gained by employing the method of action-angle variables 19 . In the subsequent section we present some numerical results and show the time evolutions of the population imbalance and the relative phase among different spin components for two cases a) $`\xi _3=0`$ and b) $`\xi _30,`$ i.e., whether total hyperfine spin of the condensate is zero or not. Time has been rescaled in units of $`2\left|U_2U_d\right|N/\mathrm{}`$ in Figures 1 and 2.
Figure 1 shows solutions of the population imbalance and relative phase of Eqs. (11) with initial conditions $`\xi _2\left(0\right)=0.12,`$ $`\xi _3\left(0\right)=0.25,`$ $`\phi _2\left(0\right)=\frac{\pi }{2}`$ for the relative interaction parameters $`\mathrm{\Lambda }=\frac{3U_d}{2\left|U_2U_d\right|N}`$ $`=0.00,0.40,0.60`$ and $`0.75`$, respectively. The left column exhibits the time evolutions of the population imbalance between $`0`$ and $`\pm 1`$ components. Josephson sinusoidal oscillations are observed as $`\mathrm{\Lambda }`$ increases in Figs. 1a, 1b. In Fig. 1c, there is a critical transition point for $`\mathrm{\Lambda }=\mathrm{\Lambda }_c=0.60`$. Increasing $`\mathrm{\Lambda }`$ further for example to $`\mathrm{\Lambda }=0.75`$, the population of each equivalent component oscillates around a nonzero time averaged value, which gives a net population imbalance $`\xi _2\left(t\right)>0`$. This is closely related to the MQST phenomenon. Simultaneously, the right column shows the time evolutions of the relative phase between $`0`$ and $`\pm 1`$ components and indicates that as $`\mathrm{\Lambda }`$ increases $`\phi _2\left(t\right)`$ varies from a monotonically increasing function of time to a periodically oscillating function of time. In the nonrigid pendulum analogy, this corresponds that the motion of our system is turned from “running-phase modes” (Figure 1a) into “$`\pi `$-phase modes” (Figures 1b-1d) or from a rotation into a vibration. Here the definitions of running-phase modes and $`\pi `$-phase modes are the same as those for two weakly coupled BEC 14 .
This critical behavior depends on $`\mathrm{\Lambda }_c=\mathrm{\Lambda }_c\left[\xi _2\left(0\right),\xi _3\left(0\right),\phi _2\left(0\right)=\frac{\pi }{2}\right]`$, as can be easily found from the energy conservation constraint and the boundness of the tunneling energy in Eq. (12). In fact, the value $`\xi _2\left(t\right)<0`$ is inaccessible at any time if
$`\mathrm{\Lambda }`$ $`>\mathrm{\Lambda }_c={\displaystyle \frac{1}{3}}\xi _2\left(0\right)+{\displaystyle \frac{2\sqrt{19\xi _3^2\left(0\right)}}{9\xi _2\left(0\right)}}`$
$`{\displaystyle \frac{\left(1+3\xi _2\left(0\right)\right)\sqrt{\left(23\xi _2\left(0\right)\right)^236\xi _3^2\left(0\right)}}{9\xi _2\left(0\right)}}.`$ (16)
When $`\phi _2\left(0\right)=0`$ the critical parameter is
$`\mathrm{\Lambda }_c`$ $`={\displaystyle \frac{1}{3}}\xi _2\left(0\right){\displaystyle \frac{2\sqrt{19\xi _3^2\left(0\right)}}{9\xi _2\left(0\right)}}`$
$`+{\displaystyle \frac{\left(1+3\xi _2\left(0\right)\right)\sqrt{\left(23\xi _2\left(0\right)\right)^236\xi _3^2\left(0\right)}}{9\xi _2\left(0\right)}},`$ (17)
and $`\mathrm{\Lambda }<\mathrm{\Lambda }_c`$ marks the regime of MQST. Specifically for $`\xi _2\left(0\right)=0.12,`$ $`\xi _3\left(0\right)=0.25,`$ the critical parameter is $`\mathrm{\Lambda }_c=0.18.`$ We take $`\mathrm{\Lambda }=0.20`$ and see MQST does set in.
In Figure 2 we show solutions of Eqs. (11) with initial conditions $`\xi _2\left(0\right)=0.12,`$ $`\xi _3\left(0\right)=0,`$ $`\phi _2\left(0\right)=0`$ for parameters $`\mathrm{\Lambda }`$ $`=0.75,0.60,0.43`$ and $`0.00,`$ respectively. The left column shows again the time evolutions of the population imbalance between $`0`$ and $`\pm 1`$ components and indicates that $`\xi _2\left(t\right)`$ is always a periodic function of time as $`\mathrm{\Lambda }`$ decreases. For $`\mathrm{\Lambda }=\mathrm{\Lambda }_c=0.43`$ the population difference is self locked to the initial value, which serves as another sign of the MQST phenomenon. The right column shows the time evolution of the relative phase between $`0`$ and $`\pm 1`$ components and indicates that as $`\mathrm{\Lambda }`$ decreases $`\phi _2\left(t\right)`$ is always a periodic function of time around its mean value $`\phi _2\left(t\right)=0`$. The dynamics corresponds to “zero-phase modes”. Moreover we observe that MQST occurs when
$$\mathrm{\Lambda }<\mathrm{\Lambda }_c=\frac{2}{3}2\xi _2\left(0\right),$$
(18)
while for $`\phi _2\left(0\right)=\pi `$ MQST never happens.
Therefore it is obvious that the spin-mixing dynamics depends on the relative interaction parameter $`\mathrm{\Lambda }`$ and is also sensitive to the initial occupations and relative phases of the three components, which can be adjusted by engineering Raman pulses 20 . In practical experiments there are usually two different ways to achieve MQST. In Figures 1 and 2, $`\xi _2\left(0\right)`$ and $`\phi _2\left(0\right)`$ are kept constants while $`\mathrm{\Lambda }`$ varies (by changing the geometry of condensates). On the other hand, one can calibrate the initial values of the population imbalance $`\xi _2\left(0\right)`$ with a fixed trap geometry (i.e., $`\mathrm{\Lambda }`$ remains constant) and $`\phi _2\left(0\right)`$ 15 .
In order to further characterize the evolution of the system we summarize the full dynamic behavior of Eq. (11) in Figure 3 that shows the $`\xi _2\left(t\right)`$-$`\phi _2\left(t\right)`$ phase portrait with constant energy contour. The distinction between the two dynamic regimes – nonlinear Josephson tunneling and MQST – becomes more apparent in Figure 3. In the regime of Josephson oscillation the dynamic variables follow a closed phase space trajectory, while in the self-trapping regime they follow an open trajectory with an unbounded phase.
We are inspired by the great expectation that tuning the contact interaction between cold atoms close to zero by Feshbach resonance 21 will make the MDDI more prominent or even the dominant interaction 10 . Based on the experimental values of the $`s`$-wave scattering lengths for <sup>23</sup>Na, $`a_0=\left(50.0\pm 1.6\right)a_B`$ and $`a_2=\left(55.0\pm 1.7\right)a_B`$ with $`a_B`$ being the Bohr radius 23 , the ratio of coefficients $`c_d`$ and $`\left|c_2\right|`$ can be shown to be $`0.007`$, while for <sup>87</sup>Rb($`a_0=101.8a_B`$ and $`a_2=100.4a_B`$) it is $`0.1`$ 6 . Evaluation with a simple variational wave function 24 gives $`U_2N6nK`$ for a sodium condensate up to $`N=5\times 10^6`$ atoms confined in an optical dipole trap with a very small trapping volume ($`10^8`$ $`cm^3`$) 1 ; 22 . On the other hand, as argued in Refs. 4 ; 6 ; 10 ; 17 ; 18 , both the sign and the magnitude of the MDDI $`U_d`$ can be greatly tuned with trapping geometry. As has been shown 6 , $`U_d/U_2`$ depends on a monotonically increasing function of the condensate aspect ratio $`\kappa `$, bounded between $`1`$ and $`2`$. This clearly shows the possibility of adjusting the MDDI very close to spin exchange interaction, i.e. $`U_d/U_21`$. At this point our parameter $`\mathrm{\Lambda }`$ may take values in a large scale and the manifestation of MQST is within reach with current technologies. We point out although for all of the alkali-metal atomic condensations the MDDI is rather weak compared to the contact potential, the experimental achievement of BEC with transition-metal chromium $`{}_{}{}^{52}Cr`$ provides us the hope because the MDDI here is 36 times stronger than that of alkali atoms. For these reasons, a degenerate quantum gas with adjustable long- and short-range interactions can be experimentally realized in near future. By altering the strengths of two kinds of interactions and the initial conditions, the nonlinear tunneling dynamics of dipolar spinor BEC consequently sustains a self-maintained population imbalance: a novel MQST effect. Such a scheme would avoid the difficulty of realizing experimentally the double-well magnetic trap.
## IV Summary
In summary we have described a semiclassical treatment of the spin-1 dipolar spinor condensates. As a result of the conservation of atom numbers and total hyperfine spin of the condensate, the classical equations of motion are derived and discussed in a similar way as in double-well BJJ. It is demonstrated that spontaneous magnetization and spin-mixing dynamics depend on both the spin-exchange interaction $`U_2`$ and the MDDI $`U_d`$ through the ratio $`\frac{3U_d}{2\left|U_2U_d\right|N}.`$ The initial population imbalance, the relative phase among the three components of the condensate as well as the total hyperfine spin of the system all play important roles in the semiclassical dynamics. Finally we have indicated the possibility of using dipolar spinor condensate as a platform for practical manipulation of MQST and Josephson oscillation.
## V Acknowledgment
This work was supported by National Natural Science Foundation of China under Grant Nos.10475053, 10175039 and 90203007.
* corresponding author: [email protected]
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# Technical Report IDSIA-13-05 Asymptotics of Discrete MDL for Online Prediction A shorter version of this paper [PH04a] appeared in COLT 2004.
## 1 Introduction
“Always prefer the *simplest* explanation for your observation,” says Occam’s razor. In Learning and Information Theory, simplicity is often quantified in terms of description length, giving immediate rise to the Minimum Description Length (MDL) principle \[WB68, Ris78, Grü98\]. Thus MDL can be seen as a strategy against overfitting. An alternative way to think of MDL is Bayesian. The explanations for the observations (the *models*) are endowed with a prior. Then the model having maximum a posteriori (MAP) probability is also a two-part MDL estimate, where the correspondence between probabilities and description lengths is simply by a negative logarithm.
How does two-part MDL perform for prediction? Some very accurate answers to this question have been already given. If the data is generated by an independently identically distributed (i.i.d.) process, then the MDL estimates are consistent \[BC91\]. In this case, an important quantity to consider is the *index of resolvability*, which depends on the complexity of the data generating process. This quantity is a tight bound on the regret in terms of coding (i.e. the excess code length). Moreover, the index of resolvability also bounds the predictive regret, namely the rate of convergence of the predictive distribution to the true one. These results apply to both discrete and continuously parameterized model classes, where in the latter case the MDL estimator must be discretized with an appropriate precision.
Under the relaxed assumption that the data generating process obeys a central limit theorem and some additional conditions, Rissanen \[Ris96, BRY98\] proves an asymptotic bound on the regret of MDL codes. Here, he also removes the coding redundancy arising if two-part codes are defined in the straightforward way. The resulting bound is very similar to that in \[CB90\] for Bayes mixture codes and i.i.d. processes, where the i.i.d. assumption may also be relaxed \[Hut03b\]. Other similar and related results can be found in \[GV01, GV04\].
In this work, we develop new methods in order to arrive at very general consistency theorems for MDL on *countable model classes*. Our setup is *online sequence prediction*, that is, the symbols $`x_1,x_2,\mathrm{}`$ of an infinite sequence are revealed successively by the environment, where our task is to predict the next symbol in each time step. Consistency is established by proving *finite cumulative bounds* on the differences of the predictive to the true distribution. Differences will be measured in terms of the relative entropy, the quadratic distance, and the Hellinger distance. Most of our results are based on the only assumption that the data generating process is *contained in the models class*. (The discussion of how strong this assumption is, will be postponed to the last section.) Our results imply regret bounds with *arbitrary* loss functions. Moreover, they can be directly applied to important general setups such as pattern classification, regression, and universal induction.
As many scientific models (e.g. in physics or biology) are smooth, much statistical work is focussed on continuous model classes. On the other hand, the largest relevant classes from a computational point of view are at most countable. In particular, the field of Algorithmic Information Theory (also known as Kolmogorov Complexity, e.g. \[ZL70, LV97, Cal02, Hut04\]) studies the class of *all lower-semicomputable semimeasures*. Then there is a one-to-one correspondence of models and programs on a fixed universal Turing Machine. (Since programs need not halt on each input, models are semimeasures instead of measures, see e.g. \[LV97\] for details). This model class can be considered the largest one which can be in the limit processed under standard computational restrictions. We will develop all our results for semimeasures, so that they can be applied in this context, which we refer to as *universal* sequence prediction.
In the universal setup, the Bayes mixture is also termed Solomonoff-Levin prior and has been intensely studied first by Solomonoff \[Sol64, Sol78\]. Its predictive properties are excellent \[Hut01, Hut04\]. Precisely one can bound the cumulative loss by the complexity of the data generating process. This is the reference performance we compare MDL to. It turns out that the predictive properties of MDL can be exponentially worse, even in the case that the model class contains only Bernoulli distributions. Another related quantity in the universal setup is *one-part MDL*, which has been studied in \[Hut03c\]. We will briefly encounter it in Section 8.4.
The paper is structured as follows. Section 2 establishes basic definitions. In Section 3, we introduce the MDL estimator and show how it can be used for sequence prediction in at least three ways. Sections 4 and 5 are devoted to convergence theorems. In Sections 6 and 7, we study the stabilization properties of the MDL estimator. Section 8 presents more general loss bounds as well as three important applications: pattern classification, regression, and universal induction. Finally, Section 9 contains the conclusions.
## 2 Prerequisites and Notation
We build on the notation of \[LV97\] and \[Hut04\]. Let the alphabet $`𝒳`$ be a finite set of symbols. We consider the spaces $`𝒳^{}`$ and $`𝒳^{\mathrm{}}`$ of finite strings and infinite sequences over $`𝒳`$. The initial part of a sequence up to a time $`t`$ or $`t1`$ is denoted by $`x_{1:t}`$ or $`x_{<t}`$, respectively. The empty string is denoted by $`ϵ`$.
A semimeasure is a function $`\nu :𝒳^{}[0,1]`$ such that
$$\nu (ϵ)1\text{ and }\nu (x)\underset{a𝒳}{}\nu (xa)\text{ for all }x𝒳^{}$$
(1)
holds. If equality holds in both inequalities of (1), then we have a measure. Intuitively, the quantity $`\nu (x)`$ can be understood as the probability that a data generating process yields a string starting with $`x`$. Then, for a measure, the probabilities of all joint continuations of $`x`$ add up to $`\nu (x)`$, while for a semimeasure, there may be a “probability leak” (1). Recall that we are interested in semimeasures (and not only in measures) because of their correspondence to programs on a fixed universal Turing machine in the universal setup and our inability to decide the halting problem.
Let $`𝒞`$ be a countable class of (semi)measures, i.e. $`𝒞=\{\nu _i:iI\}`$ with finite or infinite index set $`I`$. A (semi)measure $`\stackrel{~}{\nu }`$ dominates the class $`𝒞`$ iff for every $`\nu _i𝒞`$ there is a constant $`c_i>0`$ such that $`\nu (x)c_i\nu _i(x)`$ holds for all $`x𝒳^{}`$. A dominant semimeasure $`\stackrel{~}{\nu }`$ need not be contained in $`𝒞`$.
Each (semi)measure $`\nu 𝒞`$ is associated with a weight $`w_\nu >0`$, and we require $`_\nu w_\nu 1`$. We may interpret the weights as a prior on $`𝒞`$. Then it is obvious that the Bayes mixture
$$\xi (x)\xi _{[𝒞]}(x):=\underset{\nu 𝒞}{}w_\nu \nu (x)(\text{for }x𝒳^{})$$
(2)
dominates $`𝒞`$. Assume that there is some measure $`\mu 𝒞`$, the true distribution, generating sequences $`x_<\mathrm{}𝒳^{\mathrm{}}`$. Typically $`\mu `$ is unknown. (Note that we require $`\mu `$ to be a measure: The data stream always continues, there are no “probability leaks”.) If some initial part $`x_{<t}`$ of a sequence is given, the probability of observing $`x_t𝒳`$ as a next symbol is
$$\mu (x_t|x_{<t})=\frac{\mu (x_{<t}x_t)}{\mu (x_{<t})}\text{ if }\mu (x_{<t})>0\text{ and }\mu (x_t|x_{<t})=0\text{ if }\mu (x_{<t})=0.$$
(3)
and, for well-definedness, $`\mu (x_t|x_{<t})=0`$ if $`\mu (x_{<t})=0`$ (this case has probability zero). Note that $`\mu (x_t|x_{<t})`$ can depend on the complete history $`x_{<t}`$. We may generally define the quantity (3) for any function $`\phi :𝒳^{}[0,1]`$; we call $`\phi (x_t|x_{<t}):=\frac{\phi (x_{1:t})}{\phi (x_{<t})}`$ the $`\phi `$-prediction. Clearly, this is not necessarily a probability on $`𝒳`$ for general $`\phi `$. For a semimeasure $`\nu `$ in particular, the $`\nu `$-prediction $`\nu (|x_{<t})`$ is a semimeasure on $`𝒳`$.
We define the expectation with respect to the true probability $`\mu `$: Let $`n0`$ and $`f:𝒳^n`$ be a function, then
$$𝐄f=𝐄f(x_{1:n})=\underset{x_{1:n}𝒳^n}{}\mu (x_{1:n})f(x_{1:n}).$$
(4)
More general, the expectation may be defined as an integral over infinite sequences. But since we won’t need it, we can keep things simple. The following is a central result about prediction with Bayes mixtures in a form independent of Algorithmic Information Theory.
###### Theorem 1
For any class of (semi)measures $`𝒞`$ containing the true distribution $`\mu `$, which is a measure, we have
$$\underset{t=1}{\overset{\mathrm{}}{}}𝐄\underset{a𝒳}{}\left(\mu (a|x_{<t})\xi (a|x_{<t})\right)^2\mathrm{ln}w_\mu ^1.$$
(5)
This was found by Solomonoff (\[Sol78\]) for universal sequence prediction. A proof is also given in \[LV97\] (only for binary alphabet) or \[Hut04\] (arbitrary alphabet). It is surprisingly simple once Lemma 6 is known. A few lines analogous to (14) and (15) exploiting the dominance of $`\xi `$ are sufficient.
One should be aware that the condition $`\mu 𝒞`$ is essential in general, for both Bayes and MDL predictions \[GL04\]. On the other hand, one can show that if there is an element in $`𝒞`$ which is sufficiently close to $`\mu `$ in an appropriate sense, then still good predictive properties hold \[Hut03b\].
Note that although $`w_\nu `$ can be interpreted as a prior on the model class, we do not assume any probabilistic structure for $`𝒞`$ (e.g. a sampling mechanism). The theorem rather states that the cumulative loss is bounded by a quantity depending on the complexity $`\mathrm{ln}w_\mu ^1`$ of the true distribution. The same kind of assertion will be proven for MDL predictions later.
The bound (5) implies that the $`\xi `$-predictions converge to the $`\mu `$-predictions almost surely (i.e. with $`\mu `$-probability one). This is not hard to see, since with the abbreviation $`s_t=_a\left(\mu (a|x_{<t})\xi (a|x_{<t})\right)^2`$ and for each $`\epsilon >0`$, we have
$`𝐏(tn:s_t\epsilon )`$ $`=`$ $`𝐏\left({\displaystyle \underset{tn}{}}\left\{s_t\epsilon \right\}\right)`$ (6)
$``$ $`{\displaystyle \underset{tn}{}}𝐏\left(s_t\epsilon \right){\displaystyle \frac{1}{\epsilon }}{\displaystyle \underset{t=n}{\overset{\mathrm{}}{}}}𝐄s_t\stackrel{n\mathrm{}}{}0.`$
Actually, (5) yields an even stronger assertion, since it characterizes the speed of convergence by the quantity on the right hand side. Precisely, the expected number of times $`t`$ in which $`\xi (a|x_{<t})`$ deviates by more than $`\epsilon `$ from $`\mu (a|x_{<t})`$ is finite and bounded by $`\mathrm{ln}w_\mu ^1/\epsilon ^2`$, and the probability that the number of $`\epsilon `$-deviations exceeds $`\frac{\mathrm{ln}w_\mu ^1}{\epsilon ^2\delta }`$ is smaller than $`\delta `$. (However, we *cannot* conclude a convergence rate of $`s_t=o(\frac{1}{t})`$ from (5), since the quadratic differences generally do not decrease monotonically.)
Since we will encounter this type of convergence (5) frequently in the following, we call it convergence in mean sum (i.m.s):
$$\phi \stackrel{i.m.s.}{}\mu C>0:\underset{t=1}{\overset{\mathrm{}}{}}𝐄\underset{a𝒳}{}\left(\mu (a|x_{<t})\phi (a|x_{<t})\right)^2<\mathrm{}.$$
(7)
Then Theorem 1 states that the $`\xi `$ predictions converge to the $`\mu `$ predictions i.m.s., or “$`\xi `$ converges to $`\mu `$ i.m.s.” for short. By (6), convergence i.m.s. implies almost sure convergence (with respect to the true distribution $`\mu `$). Note that in contrast, convergence in the mean, i.e. $`𝐄[_a(\mu (a|x_{<t})\phi (a|x_{<t}))^2]\stackrel{t\mathrm{}}{}0`$, only implies convergence in probability.
Probabilities vs. Description Lengths. By the Kraft inequality, each (semi)measure can be associated with a code length or complexity by means of the negative logarithm, where all (binary) codewords form a prefix-free set. The converse holds as well. We introduce the abbreviation
$$K\mathrm{}=\mathrm{log}_2\mathrm{},\text{ e.g. }K\nu (x)=\mathrm{log}_2\nu (x)$$
(8)
for a semimeasure $`\nu `$ and $`x𝒳^{}`$ and $`K\xi (x)=\mathrm{log}_2\xi (x)`$ for the Bayes mixture $`\xi `$. It is common to ignore the somewhat irrelevant restriction that code lengths must be integer. In particular, given a class of semimeasures $`𝒞`$ together with weights, each weight $`w_\nu `$ corresponds to a description length or complexity
$$Kw(\nu )=\mathrm{log}_2w_\nu .$$
(9)
It is often only a matter of notational convenience if description lengths or probabilities are used, but description lengths are generally preferred in Algorithmic Information Theory. Keeping the equivalence in mind, we will develop the general theory in terms of probabilities, but formulate parts of the results in universal sequence prediction rather in terms of complexities.
## 3 MDL Estimator and Predictions
Assume that $`𝒞`$ is a countable class of semimeasures together with weights $`(w_\nu )_{\nu 𝒞}`$, and $`x𝒳^{}`$ is some string. Then the maximizing element $`\nu ^x`$, often called MAP (maximum a posteriori) estimator, is defined as
$$\nu ^x=\nu _{[𝒞]}^x=\mathrm{arg}\underset{\nu 𝒞}{\mathrm{max}}\{w_\nu \nu (x)\}.$$
(10)
In case of a tie, we need not specify the further choice at this point, just pick any of the maximizing elements. But for concreteness, you may think that ties are broken in favor of larger prior weights. The maximum is always attained in (10) since for each $`\epsilon >0`$ at most a finite number of elements fulfil $`w_\nu \nu (x)>\epsilon `$. Observe immediately the correspondence in terms of description lengths rather than probabilities:
$$\nu ^x=\mathrm{arg}\underset{\nu 𝒞}{\mathrm{min}}\left\{Kw(\nu )+K\nu (x)\right\}.$$
Then the minimum description length principle is obvious: $`\nu ^x`$ minimizes the joint description length of the model plus the data given the model<sup>1</sup><sup>1</sup>1The term MAP estimator is more precise. For two reasons, our definition might not be considered as MDL in the strict sense. First, MDL is often associated with a specific prior, while we admit arbitrary priors (compare the discussion section at the end of this paper). Second, when coding some data $`x`$, one can exploit the fact that once the distribution $`\nu ^x`$ is specified, only data which leads to this $`\nu ^x`$ needs to be considered. This allows for a description shorter than $`Kw(\nu ^x)`$. Nevertheless, the *construction principle* is commonly termed MDL, compare e.g. the “ideal MDL” in \[VL00\]. (see (8) and (9)). As explained before, we stick to the product notation.
For notational simplicity, let $`\nu ^{}(x)=\nu ^x(x)`$. The two-part MDL estimator is defined by
$$\varrho (x)=\varrho _{[𝒞]}(x)=w_{\nu ^x}\nu ^x(x)=\underset{\nu 𝒞}{\mathrm{max}}\{w_\nu \nu (x)\}.$$
So $`\varrho `$ chooses the maximizing element with respect to its argument. We may also use the version $`\varrho ^y(x):=w_{\nu ^y}\nu ^y(x)`$ for which the choice depends on the superscript instead of the argument. Note that the use of the term “estimator” is non-standard, since $`\varrho `$ is a product of the estimator $`\nu ^{}`$ (this use is standard) and its prior weight. There will be no confusion between these two meanings of “estimator” in the following.
For each $`x,y𝒳^{}`$,
$$\xi (x)\varrho (x)\varrho ^y(x)$$
(11)
is immediate. If $`𝒞`$ contains only measures, we have $`_a\varrho (xa)_a\varrho ^x(xa)=\varrho ^x(x)=\varrho (x)`$ for all $`x𝒳^{}`$, so $`\varrho `$ has some “anti-semimeasure” property. If $`𝒞`$ contains semimeasures, no semimeasure or anti-semimeasure property can be established for $`\varrho `$.
We can define MDL predictors according to (3). There are basically two possible ways to use MDL for prediction.
###### Definition 2
The dynamic MDL predictor is defined as
$$\varrho (a|x)=\frac{\varrho (xa)}{\varrho (x)}=\frac{\varrho ^{xa}(xa)}{\varrho ^x(x)}.$$
That is, we look for a short description of $`xa`$ and relate it to a short description of $`x=x_{<t}`$. We call this dynamic since for each possible $`a`$ we have to find a new MDL estimator. This is the closest correspondence to the Bayes mixture $`\xi `$-predictor.
###### Definition 3
The static MDL predictor is given by
$$\varrho ^{\mathrm{static}}(a|x)=\varrho ^x(a|x)=\frac{\varrho ^x(xa)}{\varrho (x)}=\frac{\varrho ^x(xa)}{\varrho ^x(x)}=\frac{\nu ^x(xa)}{\nu ^x(x)}.$$
Here obviously only one MDL estimator $`\varrho ^x`$ has to be identified. This is usually more efficient in practice.
We will define another MDL predictor, the *hybrid* one, in Section 6. It can be paraphrased as “do dynamic MDL but drop weights”. We will see that its predictive performance is weaker.
The range of the static MDL predictor is obviously contained in $`[0,1]`$. For the dynamic MDL predictor, this holds by
$$\varrho ^x(x)\varrho ^{xa}(x)\varrho ^{xa}(xa).$$
(12)
Static MDL is omnipresent in machine learning and applications, see also Section 8. In fact, many common prediction algorithms can be abstractly understood as static MDL, or rather as approximations. Namely, if a prediction task is accomplished by building a model such as a neural network with a suitable regularization<sup>2</sup><sup>2</sup>2There are however regularization methods which cannot be interpreted in this way but build on a different theoretical foundation, such as structural risk minimization. to prevent “overfitting”, this is just searching an MDL estimator within a certain class of distributions. After that, only this model is used for prediction. Dynamic MDL is applied more rarely due to its larger computational effort. For example, the similarity metric proposed in \[LCL<sup>+</sup>03\] can be interpreted as (a deterministic variant of) dynamic MDL.
We will need to convert our MDL predictors to measures by means of normalization. If $`\phi :𝒳^{}[0,1]`$ is any function, then
$$\phi _{\mathrm{norm}}(a|x):=\frac{\phi (a|x)}{_{b𝒳}\phi (b|x)}=\frac{\phi (xa)}{_{b𝒳}\phi (xb)}$$
is a measure (assume that the denominator is different from zero, which is always true with probability 1 (w.p.1) if $`\phi `$ is an MDL predictor). This procedure is known as Solomonoff normalization \[Sol78, LV97\] and results in
$$\phi _{\mathrm{norm}}(x_{1:n})=\frac{\phi (x_{1:n})}{\phi (ϵ)}\underset{t=1}{\overset{n}{}}\frac{\phi (x_{<t})}{_{a𝒳}\phi (x_{<t}a)}=\frac{\phi (x_{1:n})}{\phi (ϵ)N_\phi (x_{<n})},$$
where
$$N_\phi (x)=\underset{t=1}{\overset{\mathrm{}(x)+1}{}}\frac{_{a𝒳}\phi (x_{<t}a)}{\phi (x_{<t})}$$
(13)
is the normalizer.
We conclude this section with a simple example.
Bernoulli and i.i.d. classes. Let $`n`$, $`𝒳=\{1,\mathrm{},n\}`$, and
$$𝒞=\{\nu _\vartheta (x_{1:t})=\vartheta _{x_1}\mathrm{}\vartheta _{x_t}:\vartheta \mathrm{\Theta }\}\text{ with }\mathrm{\Theta }=\{\vartheta ([0,1])^n:\underset{i=1}{\overset{n}{}}\vartheta _i=1\}$$
be the set of all rational probability vectors with any prior $`(w_\vartheta )_{\vartheta \mathrm{\Theta }}`$. Each $`\vartheta \mathrm{\Theta }`$ generates sequences $`x_<\mathrm{}`$ of independently identically distributed (i.i.d.) random variables such that $`𝐏(x_t=i)=\vartheta _i`$ for all $`t1`$ and $`1in`$. If $`x_{1:t}`$ is the initial part of a sequence and $`\alpha \mathrm{\Theta }`$ is defined by $`\alpha _i=\frac{1}{t}|\{st:x_s=i\}|`$, then it is easy to see that
$$\nu ^{x_{1:t}}=\mathrm{arg}\underset{\vartheta \mathrm{\Theta }}{\mathrm{min}}\left\{Kw(\vartheta )\mathrm{ln}2+tD(\alpha \vartheta )\right\},$$
where $`D(\alpha \vartheta ):=_{i=1}^n\alpha _i\mathrm{ln}\frac{\alpha _i}{\vartheta _i}`$ is the Kullback-Leibler divergence. If $`|𝒳|=2`$, then $`\mathrm{\Theta }`$ is also called a Bernoulli class, and one usually takes the binary alphabet $`𝒳=\{0,1\}`$ in this case.
## 4 Dynamic MDL
We may now develop convergence results, beginning with the dynamic MDL predictor from Definition 2. The following simple lemma is crucial for all subsequent proofs.
###### Lemma 4
For an arbitrary class of (semi)measures $`𝒞`$, we have
$`(i)`$ $`\varrho (x){\displaystyle \underset{a𝒳}{}}\varrho (xa)\xi (x){\displaystyle \underset{a𝒳}{}}\xi (xa)\mathrm{and}`$
$`(ii)`$ $`\varrho ^x(x){\displaystyle \underset{a𝒳}{}}\varrho ^x(xa)\xi (x){\displaystyle \underset{a𝒳}{}}\xi (xa)`$
for all $`x𝒳^{}`$. In particular, $`\xi \varrho `$ is a semimeasure.
Proof. For all $`x𝒳^{}`$, with $`f:=\xi \varrho `$ we have
$`{\displaystyle \underset{a𝒳}{}}f(xa)`$ $`=`$ $`{\displaystyle \underset{a𝒳}{}}\left(\xi (xa)\varrho (xa)\right){\displaystyle \underset{a𝒳}{}}\left(\xi (xa)\varrho ^x(xa)\right)`$
$`=`$ $`{\displaystyle \underset{\nu \{\nu ^x\}}{}}{\displaystyle \underset{a𝒳}{}}w_\nu \nu (xa){\displaystyle \underset{\nu \{\nu ^x\}}{}}w_\nu \nu (x)=\xi (x)\varrho (x)=f(x).`$
The first inequality follows from $`\varrho ^x(xa)\varrho (xa)`$, and the second one holds since all $`\nu `$ are semimeasures. Finally, $`f(x)=\xi (x)\varrho (x)=_{\nu \{\nu ^x\}}w_\nu \nu (x)0`$ and $`f(ϵ)=\xi (ϵ)\varrho (ϵ)1`$. Hence $`f`$ is a semimeasure. $`\mathrm{}`$
The following proposition demonstrates how simple it can be to obtain a convergence result, however a weak one. Various similar results have been already obtained in the past, e.g. in \[BD62, Bar85\].
###### Proposition 5
For any class of (semi)measures $`𝒞`$ containing the true distribution $`\mu `$, we have
$$\frac{\varrho (x_t|x_{<t})}{\mu (x_t|x_{<t})}1w.\mu .p.1$$
Proof. Since $`\xi \varrho `$ is a positive semimeasure by Lemma 4, $`\frac{\xi \varrho }{\mu }`$ is a positive super-martingale. By Doob’s martingale convergence theorem (see e.g. \[Doo53\] or \[CT88\] or any textbook on advanced probability theory), it therefore converges on a set of $`\mu `$-measure one. Moreover, $`\frac{\xi }{\mu }`$ converges on a set of measure one, being a positive super-martingale as well \[LV97, Thm.5.2.2\]. Thus $`\frac{\varrho }{\mu }`$ must converge on a set of measure one. We denote this limit by $`f`$ and observe that $`fw_\mu `$ since $`\frac{\varrho }{\mu }w_\mu `$ everywhere. On this set of measure one, the denominator $`\varrho (x_{<t})/\mu (x_{<t})`$ of
$$\frac{\varrho (x_{1:t})/\mu (x_{1:t})}{\varrho (x_{<t})/\mu (x_{<t})}=\frac{\varrho (x_t|x_{<t})}{\mu (x_t|x_{<t})}$$
converges to $`f>0`$, and so does the numerator. The whole fraction thus converges to one, which was to be shown. $`\mathrm{}`$
Proposition 5 gives only a statement about “on-sequence” ($`\varrho (x_t|x_{<t})`$) convergence of the $`\varrho `$-predictions. Indeed, no conclusion about “off-sequence” convergence, i.e. $`\varrho (a|x_{<t})`$ for arbitrary $`a𝒳`$, can be drawn from the proposition, not even in the deterministic case. There, the true measure $`\mu `$ is concentrated on the particular sequence $`x_<\mathrm{}`$. So for $`ax_t`$, we have $`\mu (x_{<t}a)=0`$, and thus no assertion for $`\varrho (a|x_{<t})`$ can be made. On the other hand, an off-sequence result is essential for prediction: Even if on the *correct* next symbol the predictive probability is very close to the true value, we must be sure that this is so also for all *alternatives*. This is particularly important if we base some decision on the prediction; compare Section 8.1.
The following theorem closes this gap. In addition, it provides a statement about the speed of convergence. In order to prove it, we need a lemma establishing a relation between the square distance and the Kullback-Leibler distance, which is proven for instance in \[Hut04, Sec.3.9.2\].
###### Lemma 6
Let $`\mu `$ and $`\rho `$ be measures on $`𝒳`$, then
$$\underset{a𝒳}{}\left(\mu (a)\rho (a)\right)^2\underset{a𝒳}{}\mu (a)\mathrm{ln}\frac{\mu (a)}{\rho (a)}.$$
###### Theorem 7
For any class of (semi)measures $`𝒞`$ containing the true distribution $`\mu `$ (which is a measure), we have
$$\underset{t=1}{\overset{\mathrm{}}{}}𝐄\underset{a𝒳}{}\left(\mu (a|x_{<t})\varrho _{\mathrm{norm}}(a|x_{<t})\right)^2w_\mu ^1+\mathrm{ln}w_\mu ^1.$$
That is, $`\varrho _{\mathrm{norm}}(a|x_{<t})\stackrel{i.m.s.}{}\mu (a|x_{<t})`$ (see (7)), which implies $`\varrho _{\mathrm{norm}}(a|x_{<t})\mu (a|x_{<t})`$ with $`\mu `$-probability one.
Proof. Let $`n`$. From Lemma 6, we know
$`{\displaystyle \underset{t=1}{\overset{n}{}}}𝐄{\displaystyle \underset{a𝒳}{}}\left(\mu (a|x_{<t})\varrho _{\mathrm{norm}}(a|x_{<t})\right)^2{\displaystyle \underset{t=1}{\overset{n}{}}}𝐄{\displaystyle \underset{a𝒳}{}}\mu (a|x_{<t})\mathrm{ln}{\displaystyle \frac{\mu (a|x_{<t})}{\varrho _{\mathrm{norm}}(a|x_{<t})}}`$ (14)
$`=`$ $`{\displaystyle \underset{t=1}{\overset{n}{}}}𝐄\mathrm{ln}{\displaystyle \frac{\mu (x_t|x_{<t})}{\varrho _{\mathrm{norm}}(x_t|x_{<t})}}={\displaystyle \underset{t=1}{\overset{n}{}}}𝐄\left[\mathrm{ln}{\displaystyle \frac{\mu (x_t|x_{<t})}{\varrho (x_t|x_{<t})}}+\mathrm{ln}{\displaystyle \frac{_{a𝒳}\varrho (x_{<t}a)}{\varrho (x_{<t})}}\right].`$
Then we can estimate
$$\underset{t=1}{\overset{n}{}}𝐄\mathrm{ln}\frac{\mu (x_t|x_{<t})}{\varrho (x_t|x_{<t})}=𝐄\mathrm{ln}\underset{t=1}{\overset{n}{}}\frac{\mu (x_t|x_{<t})}{\varrho (x_t|x_{<t})}=𝐄\mathrm{ln}\frac{\mu (x_{1:n})}{\varrho (x_{1:n})}\mathrm{ln}w_\mu ^1,$$
(15)
since always $`\frac{\mu }{\varrho }w_\mu ^1`$. Moreover, by setting $`x=x_{<t}`$, using $`\mathrm{ln}uu1`$, adding an always positive max-term, and finally using $`\frac{\mu }{\varrho }w_\mu ^1`$ again, we obtain
$`𝐄\mathrm{ln}{\displaystyle \frac{_a\varrho (x_{<t}a)}{\varrho (x_{<t})}}𝐄\left[{\displaystyle \frac{_a\varrho (xa)}{\varrho (x)}}1\right]={\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}{\displaystyle \frac{\mu (x)\left[\left(_a\varrho (xa)\right)\varrho (x)\right]}{\varrho (x)}}`$ (16)
$``$ $`{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}{\displaystyle \frac{\mu (x)\left[\left(_{a𝒳}\varrho (xa)\right)\varrho (x)+\mathrm{max}\{0,\varrho (x)_{a𝒳}\varrho (xa)\}\right]}{\varrho (x)}}`$
$``$ $`w_\mu ^1{\displaystyle \underset{\mathrm{}(x)=t1}{}}\left[\left({\displaystyle \underset{a𝒳}{}}\varrho (xa)\right)\varrho (x)+\mathrm{max}\{0,\varrho (x){\displaystyle \underset{a𝒳}{}}\varrho (xa)\}\right].`$
If $`𝒞`$ contains only measures, the max-term is not necessary, since $`\varrho `$ is an “anti-semimeasure” in this case. We proceed by observing
$$\underset{t=1}{\overset{n}{}}\underset{\mathrm{}(x)=t1}{}\left[\left(\underset{a𝒳}{}\varrho (xa)\right)\varrho (x)\right]=\underset{t=1}{\overset{n}{}}\left[\underset{\mathrm{}\left(x\right)=t}{}\varrho (x)\underset{\mathrm{}\left(x\right)=t1}{}\varrho (x)\right]=\left[\underset{\mathrm{}\left(x\right)=n}{}\varrho (x)\right]\varrho (ϵ)$$
(17)
which is true since for successive $`t`$ the positive and negative terms cancel. From Lemma 4 we know $`\varrho (x)_{a𝒳}\varrho (xa)\xi (x)_{a𝒳}\xi (xa)`$ and therefore
$`{\displaystyle \underset{t=1}{\overset{n}{}}}{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}\mathrm{max}\{0,\varrho (x){\displaystyle \underset{a𝒳}{}}\varrho (xa)\}`$ $``$ $`{\displaystyle \underset{t=1}{\overset{n}{}}}{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}\mathrm{max}\{0,\xi (x){\displaystyle \underset{a𝒳}{}}\xi (xa)\}`$
$`={\displaystyle \underset{t=1}{\overset{n}{}}}{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}\left[\xi (x){\displaystyle \underset{a𝒳}{}}\xi (xa)\right]`$ $`=`$ $`\xi (ϵ){\displaystyle \underset{\mathrm{}(x)=n}{}}\xi (x).`$ (18)
Here we have again used the fact that positive and negative terms cancel for successive $`t`$, and moreover the fact that $`\xi `$ is a semimeasure. Combining (16), (17) and (18), and observing $`\varrho \xi 1`$, we obtain
$$\underset{t=1}{\overset{n}{}}𝐄\mathrm{ln}\frac{_a\varrho (x_{<t}a)}{\varrho (x_{<t})}w_\mu ^1\left[\xi (ϵ)\varrho (ϵ)+\underset{\mathrm{}\left(x\right)=n}{}\left(\varrho (x)\xi (x)\right)\right]w_\mu ^1\xi (ϵ)w_\mu ^1.$$
(19)
Therefore, (14), (15) and (19) finally prove the assertion. $`\mathrm{}`$
We point out again that the proof gets a bit simpler if $`𝒞`$ contains only measures, since then (18) becomes irrelevant. However, this case doesn’t give a tighter bound.
This is the first convergence result “in mean sum”, see (7). It implies both on-sequence and off-sequence convergence. Moreover, it asserts the convergence is “fast” in the sense that the sum of the total expected deviations is bounded by $`w_\mu ^1+\mathrm{ln}w_\mu ^1`$. Of course, $`w_\mu ^1`$ can be very large, namely $`w_\mu ^1=2^{Kw(\mu )}`$. The following example will show that this bound is sharp (save for a constant factor). Observe that in the corresponding result for mixtures, Theorem 1, the bound is much smaller, namely $`\mathrm{ln}w_\mu ^1=Kw(\mu )\mathrm{ln}2`$.
###### Example 8
Let $`𝒳=\{0,1\}`$, $`N1`$ and $`𝒞=\{\nu _1,\mathrm{},\nu _{N1},\mu \}`$. Each $`\nu _i`$ is a deterministic measure concentrated on the sequence $`z_<\mathrm{}^{(i)}=1^{i1}0^{\mathrm{}}`$, while the true distribution $`\mu `$ is deterministic and concentrated on $`x_<\mathrm{}=1^{\mathrm{}}`$. Let $`w_{\nu _i}=w_\mu =\frac{1}{N}`$ for all $`i`$. Then $`\mu `$ generates $`x_<\mathrm{}`$, and for each $`tN1`$ we have $`\varrho _{\mathrm{norm}}(0|x_{<t})=\varrho _{\mathrm{norm}}(1|x_{<t})=\frac{1}{2}`$. Hence, $`_t𝐄_a(\mu (a|x_{<t})\varrho _{\mathrm{norm}}(a|x_{<t}))^2=\frac{1}{2}(N1)\stackrel{\times }{=}w_\mu ^1`$. In Example 15 we will even see a case where the model class contains only Bernoulli distributions and nevertheless the exponential bound is sharp.
The next result implies that convergence holds also for the un-normalized dynamic MDL predictor.
###### Theorem 9
For any class of (semi)measures $`𝒞`$ containing the true distribution $`\mu `$, we have
$`(i)`$ $`{\displaystyle \underset{t=1}{\overset{\mathrm{}}{}}}𝐄\left|\mathrm{ln}{\displaystyle \underset{a𝒳}{}}\varrho (a|x_{<t})\right|2w_\mu ^1\mathrm{and}`$
$`(ii)`$ $`{\displaystyle \underset{t=1}{\overset{\mathrm{}}{}}}𝐄{\displaystyle \underset{a𝒳}{}}|\varrho _{\mathrm{norm}}(a|x_{<t})\varrho (a|x_{<t})|={\displaystyle \underset{t=1}{\overset{\mathrm{}}{}}}𝐄|1{\displaystyle \underset{a𝒳}{}}\varrho (a|x_{<t})|2w_\mu ^1.`$
Proof. $`(i)`$ Define $`u^+=\mathrm{max}\{0,u\}`$ for $`u`$, then for $`x:=x_{<t}𝒳^{t1}`$ we have
$`𝐄\left|\mathrm{ln}{\displaystyle \underset{a𝒳}{}}\varrho (a|x)\right|=𝐄\left|\mathrm{ln}{\displaystyle \frac{\underset{a}{}\varrho (xa)}{\varrho (x)}}\right|=𝐄[\left(\mathrm{ln}{\displaystyle \frac{\underset{a}{}\varrho (xa)}{\varrho (x)}}\right)^++\left(\mathrm{ln}{\displaystyle \frac{\varrho (x)}{_a\varrho (xa)}}\right)^+]`$
$``$ $`𝐄{\displaystyle \frac{\left(\underset{a}{}\varrho (xa)\varrho (x)\right)^+}{\varrho (x)}}+𝐄{\displaystyle \frac{\left(\varrho (x)\underset{a}{}\varrho (xa)\right)^+}{_a\varrho (xa)}}`$
$`=`$ $`{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}{\displaystyle \frac{\mu (x)\left(\underset{a}{}\varrho (xa)\varrho (x)\right)^+}{\varrho (x)}}+{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}{\displaystyle \frac{\mu (x)\left(\varrho (x)\underset{a}{}\varrho (xa)\right)^+}{_a\varrho (xa)}}`$
$``$ $`w_\mu ^1{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}\left(_a\varrho (xa)\varrho (x)\right)^++w_\mu ^1{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}\left(\varrho (x)_a\varrho (xa)\right)^+`$
$`=`$ $`w_\mu ^1{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}|\varrho (x)_a\varrho (xa)|=w_\mu ^1{\displaystyle \underset{\mathrm{}\left(x\right)=t1}{}}\left[_a\varrho (xa)\varrho (x)+2\left(\varrho (x)_a\varrho (xa)\right)^+\right]`$
Here, $`|u|=u^++(u)^+=u+2u^+`$, $`\mathrm{ln}uu1`$, and $`\varrho w_\mu \mu `$ have been used, the latter implies also $`_a\varrho (xa)w_\mu _a\mu (xa)=w_\mu \mu (x)`$. The last expression in this (in)equality chain, when summed over $`t=1\mathrm{}\mathrm{}`$ is bounded by $`2w_\mu ^1`$ by essentially the same arguments (16) - (19) as in the proof of Theorem 7.
$`(ii)`$ Let again $`x:=x_{<t}`$ and use $`\varrho _{\mathrm{norm}}(a|x)=\varrho (a|x)/_b\varrho (b|x)`$ to obtain
$`{\displaystyle \underset{a}{}}|\varrho _{\mathrm{norm}}(a|x)\varrho (a|x)|`$ $`=`$ $`{\displaystyle \underset{a}{}}{\displaystyle \frac{\varrho (a|x)}{_b\varrho (b|x)}}|1{\displaystyle \underset{b}{}}\varrho (b|x)|=|1{\displaystyle \underset{b}{}}\varrho (b|x)|`$
$`=`$ $`{\displaystyle \frac{\left(\underset{a}{}\varrho (xa)\varrho (x)\right)^+}{\varrho (x)}}+{\displaystyle \frac{\left(\varrho (x)\underset{a}{}\varrho (xa)\right)^+}{\varrho (x)}}.`$
Then take the expectation $`𝐄`$ and the sum $`_{t=1}^{\mathrm{}}`$ and proceed as in $`(i)`$. $`\mathrm{}`$
###### Corollary 10
For any class of (semi)measures $`𝒞`$ containing the true distribution $`\mu `$, we have
$$\underset{t=1}{\overset{\mathrm{}}{}}𝐄\underset{a𝒳}{}\left(\mu (a|x_{<t})\varrho (a|x_{<t})\right)^28w_\mu ^1.$$
That is, $`\varrho (a|x_{<t})\stackrel{i.m.s.}{}\mu (a|x_{<t})`$ (see (7)).
Proof. For two functions $`\phi _1,\phi _2:𝒳^{}[0,1]`$, let
$$\mathrm{\Delta }(\phi _1,\phi _2)=\left(\underset{t=1}{\overset{\mathrm{}}{}}𝐄\underset{a𝒳}{}\left(\phi _1(a|x_{<t})\phi _2(a|x_{<t})\right)^2\right)^{\frac{1}{2}}.$$
(21)
Then the triangle inequality holds for $`\mathrm{\Delta }(,)`$, since $`\mathrm{\Delta }`$ is (proportional to) an Euclidian distance (2-norm). Moreover, $`\mathrm{\Delta }(\mu ,\varrho _{\mathrm{norm}})\sqrt{2w_\mu ^1}`$ by Theorem 7 and $`\mathrm{ln}w_\mu ^1w_\mu ^11w_\mu ^1`$. We also have $`\mathrm{\Delta }(\varrho _{\mathrm{norm}},\varrho )\sqrt{2w_\mu ^1}`$ by multiplying $`|\varrho _{\mathrm{norm}}\varrho |`$ in Theorem 9$`(ii)`$ with another $`|\varrho _{\mathrm{norm}}\varrho |`$. Note $`|\varrho _{\mathrm{norm}}\varrho |1`$, since both $`\varrho (a|x),\varrho _{\mathrm{norm}}(a|x)[0,1]`$, for $`\varrho `$ this holds by (12). This implies $`\mathrm{\Delta }(\mu ,\varrho )\mathrm{\Delta }(\mu ,\varrho _{\mathrm{norm}})+\mathrm{\Delta }(\varrho _{\mathrm{norm}},\varrho )2\sqrt{2w_\mu ^1}`$. $`\mathrm{}`$
###### Corollary 11
For almost all $`x_<\mathrm{}𝒳^{\mathrm{}}`$, the normalizer $`N_\varrho `$ defined in (13) converges to a number which is finite and greater than zero, i.e. $`0<N_\varrho (x_<\mathrm{})<\mathrm{}`$. Moreover, the sum of the MDL posterior estimates converges to one almost surely,
$$\underset{a𝒳}{}\varrho (a|x_{<t})=\frac{_{a𝒳}\varrho (x_{<t}a)}{\varrho (x_{<t})}1\mathrm{as}t\mathrm{}w.\mu .p.1.$$
(22)
Proof. Theorem 9 implies that with probability one, the sum $`_1^n\left|\mathrm{ln}\frac{_a\varrho (x_{<t}a)}{\varrho (x_{<t})}\right|`$ is bounded in $`n`$, hence converges absolutely, hence also the limit
$$\mathrm{ln}N_\varrho (x_<\mathrm{})=\underset{t=1}{\overset{\mathrm{}}{}}\mathrm{ln}\frac{_{a𝒳}\varrho (x_{<t}a)}{\varrho (x_{<t})}$$
exists and is finite. For these sequences, $`0<N_\varrho (x_<\mathrm{})<\mathrm{}`$ and (22) follows. $`\mathrm{}`$
## 5 Static MDL
Static MDL as introduced in Definition 3 is usually more efficient and thus preferred in practice, since only one MDL estimator has to be computed. The following technical result will allow to conclude that the static MDL predictions converge in mean sum like the dynamic ones.
###### Theorem 12
For any class of (semi)measures $`𝒞`$ containing the true distribution $`\mu `$, we have
$$\underset{t=1}{\overset{\mathrm{}}{}}𝐄\underset{a𝒳}{}|\varrho _{\mathrm{norm}}^{x_{<t}}(a|x_{<t})\varrho ^{x_{<t}}(a|x_{<t})|=\underset{t=1}{\overset{\mathrm{}}{}}𝐄|1\underset{a𝒳}{}\varrho ^{x_{<t}}(a|x_{<t})|w_\mu ^1.$$
Proof. We proceed in a similar way as in the proof of Theorem 7, (16) - (18). From Lemma 4, we know $`\varrho (x)_a\varrho ^x(xa)\xi (x)_a\xi (xa)`$. Then
$`{\displaystyle \underset{t=1}{\overset{n}{}}}𝐄|1{\displaystyle \underset{a𝒳}{}}\varrho ^{x_{<t}}(a|x_{<t})|`$ $`=`$ $`{\displaystyle \underset{t=1}{\overset{n}{}}}𝐄{\displaystyle \frac{\varrho (x_{<t})_{a𝒳}\varrho ^{x_{<t}}(x_{<t}a)}{\varrho (x_{<t})}}`$
$`=`$ $`{\displaystyle \underset{t=1}{\overset{n}{}}}{\displaystyle \underset{\mathrm{}(x)=t1}{}}\mu (x){\displaystyle \frac{\varrho (x)_{a𝒳}\varrho ^x(xa)}{\varrho (x)}}`$
$``$ $`w_\mu ^1{\displaystyle \underset{t=1}{\overset{n}{}}}{\displaystyle \underset{\mathrm{}(x)=t1}{}}\left[\varrho (x){\displaystyle \underset{a𝒳}{}}\varrho ^x(xa)\right]`$
$``$ $`w_\mu ^1{\displaystyle \underset{t=1}{\overset{n}{}}}{\displaystyle \underset{\mathrm{}(x)=t1}{}}\left[\xi (x){\displaystyle \underset{a𝒳}{}}\xi (xa)\right]`$
$``$ $`w_\mu ^1\left[\xi (ϵ){\displaystyle \underset{\mathrm{}(x)=n}{}}\xi (x)\right]w_\mu ^1`$
for all $`n`$. This implies the assertion. Again we have used $`\frac{\mu }{\varrho }w_\mu ^1`$ and the fact that positive and negative terms cancel for successive $`t`$. $`\mathrm{}`$
###### Corollary 13
For any class of (semi)measures $`𝒞`$ containing the true distribution $`\mu `$, we have
$`{\displaystyle \underset{t=1}{\overset{\mathrm{}}{}}}𝐄{\displaystyle \underset{a𝒳}{}}\left(\mu (a|x_{<t})\varrho ^{x_{<t}}(a|x_{<t})\right)^2`$ $``$ $`21w_\mu ^1\text{ and }`$
$`{\displaystyle \underset{t=1}{\overset{\mathrm{}}{}}}𝐄{\displaystyle \underset{a𝒳}{}}\left(\mu (a|x_{<t})\varrho _{\mathrm{norm}}^{x_{<t}}(a|x_{<t})\right)^2`$ $``$ $`32w_\mu ^1.`$
That is, $`\varrho ^{x_{<t}}(a|x_{<t})\stackrel{i.m.s.}{}\mu (a|x_{<t})`$ and $`\varrho _{\mathrm{norm}}^{x_{<t}}(a|x_{<t})\stackrel{i.m.s.}{}\mu (a|x_{<t})`$.
Proof. Using $`\varrho (xa)\varrho ^x(xa)`$ and the triangle inequality, we see
$$\underset{a}{}|\varrho (a|x)\varrho ^x(a|x)|=|\underset{a}{}\varrho (a|x)\underset{a}{}\varrho ^x(a|x)||\underset{a}{}\varrho (a|x)1|+|1\underset{a}{}\varrho ^x(a|x)|$$
With $`\mathrm{\Delta }(,)`$ as in (21), using $`|\varrho \varrho ^x|1`$ we therefore have
$$\mathrm{\Delta }^2(\varrho ,\varrho ^{\mathrm{static}})\underset{t=1}{\overset{\mathrm{}}{}}𝐄\underset{a}{}|\varrho (a|x)\varrho ^x(a|x)|3w_\mu ^1$$
according to Theorem 9 $`(ii)`$ and Theorem 12. Since $`\mathrm{\Delta }(\mu ,\varrho )2\sqrt{2w_\mu ^1}`$ holds by Corollary 10, we obtain $`\mathrm{\Delta }(\mu ,\varrho ^{\mathrm{static}})\mathrm{\Delta }(\mu ,\varrho )+\mathrm{\Delta }(\varrho ,\varrho ^{\mathrm{static}})\sqrt{21w_\mu ^1}`$. Theorem 12 also asserts $`\mathrm{\Delta }(\varrho ^{\mathrm{static}},\varrho _{\mathrm{norm}}^{\mathrm{static}})\sqrt{w_\mu ^1}`$, hence $`\mathrm{\Delta }(\mu ,\varrho _{\mathrm{norm}}^{\mathrm{static}})\sqrt{32w_\mu ^1}`$ follows. $`\mathrm{}`$
Distance measures. The total expected square error is not the only possible choice for measuring distance of distributions and speed of convergence. In fact, looking at the proof of Theorem 7, the expected Kullback-Leibler distance may seem more natural at a first glance. However this quantity behaves well only under dynamic MDL, not static MDL. To see this, let $`𝒞\{0,\frac{1}{2}\}`$ contain two Bernoulli distributions, both with prior weight $`\frac{1}{2}`$, and let $`\mu \frac{1}{2}`$ be the uniform measure. If the first symbol happens to be 0, which occurs with probability $`\frac{1}{2}`$, then the static MDL estimate is $`\nu ^00`$. Then $`D(\mu \nu ^0)=\mathrm{}`$, hence the expectation is $`\mathrm{}`$, too. The quadratic distance behaves locally like the Kullback-Leibler distance (Lemma 6), but otherwise is bounded and thus more convenient.
Another possible choice is the *Hellinger distance*
$`h_t(\mu ,\phi )|_{x_{<t}}`$ $`=`$ $`{\displaystyle \underset{a𝒳}{}}\left(\sqrt{\mu (a|x_{<t})}\sqrt{\phi (a|x_{<t})}\right)^2\text{ and }`$ (23)
$`H_{1:n}(\mu ,\phi )`$ $`=`$ $`{\displaystyle \underset{t=1}{\overset{n}{}}}𝐄h_t(\mu ,\phi ).`$ (24)
Like the square distance, the Hellinger distance is bounded by both the relative entropy and the absolute distance:
$`h_t(\mu ,\phi )`$ $``$ $`{\displaystyle \underset{a𝒳}{}}\mu (a|x_{<t})\mathrm{ln}{\displaystyle \frac{\mu (a|x_{<t})}{\phi (a|x_{<t})}}\text{ and }`$ (25)
$`h_t(\mu ,\phi )`$ $``$ $`{\displaystyle \underset{a𝒳}{}}|\mu (a|x_{<t})\phi (a|x_{<t})|.`$ (26)
The former is e.g. shown in \[Hut04, Lem.3.11, p.114\], the latter follows from $`(\sqrt{u}\sqrt{v})^2|uv|`$ for any $`u,v`$. Therefore, the same bounds we have proven for the square distance also hold for the Hellinger distance; they are subsumed in Corollary 14 below. Although for simplicity of notation we have preferred the square distance over the Hellinger distance in the presentation so far, in Sections 8.1 and 8.3 we will meet situations where the quadratic distance is not sufficient. Then the Hellinger distance will be useful.
The following corollary recapitulates our results and states convergence i.m.s (and therefore also w.$`\mu `$-p.1) for all combinations of un-normalized/normalized and dynamic/static MDL predictions.
###### Corollary 14
Let $`𝒞`$ contain the true distribution $`\mu `$, then
$`S_<\mathrm{}(\mu ,\varrho _{\mathrm{norm}})2w_\mu ^1,`$ $`H_<\mathrm{}(\mu ,\varrho _{\mathrm{norm}})2w_\mu ^1,`$
$`S_<\mathrm{}(\mu ,\varrho )8w_\mu ^1,`$ $`H_<\mathrm{}(\mu ,\varrho )8w_\mu ^1,`$
$`S_<\mathrm{}(\mu ,\varrho ^{\mathrm{static}})21w_\mu ^1,`$ $`H_<\mathrm{}(\mu ,\varrho ^{\mathrm{static}})21w_\mu ^1,`$
$`S_<\mathrm{}(\mu ,\varrho _{\mathrm{norm}}^{\mathrm{static}})32w_\mu ^1,`$ $`H_<\mathrm{}(\mu ,\varrho _{\mathrm{norm}}^{\mathrm{static}})32w_\mu ^1,`$
where $`S_<\mathrm{}(\mu ,\phi )=_t𝐄_a\left(\mu (a|x_{<t})\phi (a|x_{<t})\right)^2`$ and $`H_<\mathrm{}`$ is as in (24).
The following example shows that the exponential bound is sharp (except for a multiplicative constant), even if the model class contains only Bernoulli distributions. It is stated in terms of static MDL, however it equally holds for dynamic MDL.
###### Example 15
Let $`N1`$ and $`𝒞\mathrm{\Theta }=\{\frac{1}{2}\}\{\frac{1}{2}+2^{k1}:1kN\}`$ be a Bernoulli class as discussed at the end of Section 3. Let $`\mu `$ be Bernoulli with parameter $`\frac{1}{2}`$, i.e. the distribution generating fair coin flips. Assume that all weights are equally $`\frac{1}{N+1}`$. Then it is shown in \[PH04b, Prop. 5\] that
$$\underset{t=1}{\overset{\mathrm{}}{}}𝐄\left(\frac{1}{2}\varrho ^{x_{<t}}(1|x_{<t})\right)^2\frac{1}{84}\left(N4\right).$$
So the bound equals $`w_\mu ^1`$ within a multiplicative constant.
This shows that in general there is no hope to improve the bounds, even for very simple model classes. But the situation is not as bad as it might seem. First, the bounds may be exponentially smaller under certain regularity conditions on the class and the weights, as \[Ris96\] and the positive assertions in \[PH04b\] show. It is open to define such conditions for more general model classes. Second, the example just given behaves differently than Example 8. There, the error remains at a significant level for $`O(w_\mu ^1)`$ time steps, which must be regarded critical. Here in contrast, the error drops to zero as $`\frac{1}{n}`$ for a very long time, namely $`O(2^{w_\mu ^1})`$ steps, and decreases more rapidly only afterwards. This behavior is tolerable in practice. Recently, \[Li99, Zha04\] have proven that this favorable case always occurs for i.i.d., if the weights satisfy the *light tails* condition $`w_\nu ^\alpha 1`$ for some $`\alpha <1`$ \[BC91\]. Precisely, they give a rapidly decaying bound on the instantaneous error. It is open if similar results also hold in more general setups than i.i.d. Example 8 shows that at least some additional assumption is necessary.
## 6 Hybrid MDL
So far, we have not cared about what happens if two or more (semi)measures obtain the same value $`w_\nu \nu (x)`$ for some string $`x`$. In fact, for the previous results, the tie-breaking strategy can be completely arbitrary. This need not be so for all thinkable prediction methods, as we will see with the hybrid MDL predictor in the subsequent example.
###### Definition 16
The hybrid MDL predictor is given by
$$\varrho ^{\mathrm{hyb}}(a|x)=\frac{\nu ^{}(xa)}{\nu ^{}(x)}$$
(compare (10)). This can be paraphrased as “do dynamic MDL and drop the weights”. It is somewhat in-between static and dynamic MDL.
###### Example 17
Let $`𝒳=\{0,1\}`$ and $`𝒞`$ contain only two measures, the uniform measure $`\lambda `$ which is defined by $`\lambda (x)=2^{\mathrm{}(x)}`$, and another measure $`\nu `$ having $`\nu (1x)=2^{\mathrm{}(x)}`$ and $`\nu (0x)=0`$. The respective weights are $`w_\lambda =\frac{2}{3}`$ and $`w_\nu =\frac{1}{3}`$. Then, for each $`x`$ starting with $`1`$, we have $`w_\nu \nu (x)=w_\lambda \lambda (x)=\frac{1}{3}2^{\mathrm{}(x)+1}`$. Therefore, for all $`x_<\mathrm{}`$ starting with $`1`$ (a set which has uniform measure $`\frac{1}{2}`$), we have a tie. If the maximizing element $`\nu ^{}`$ is chosen to be $`\lambda `$ for even $`t`$ and $`\nu `$ for odd $`t`$, then both static and dynamic MDL predict probabilities of constantly
$$\frac{1}{2}=\lambda (a|x_{<t})=\nu (a|x_{<t})=\frac{w_\lambda \lambda (x_{<t}a)}{w_\nu \nu (x_{<t})}=\frac{w_\nu \nu (x_{<t}a)}{w_\lambda \lambda (x_{<t})}$$
for all $`a\{0,1\}`$. However, the hybrid MDL predictor values $`\frac{\nu ^{}(x_{<t}a)}{\nu ^{}(x_{<t})}`$ oscillate between $`\frac{1}{4}`$ and $`1`$.
If the ambiguity in the tie-breaking process is removed, e.g. in favor of larger weights, then the hybrid MDL predictor does converge for this example. We replace (10) by this rule:
$$\nu ^x=\mathrm{arg}\mathrm{max}\left\{w_\nu :\nu \{\nu =\mathrm{arg}\underset{\nu 𝒞}{\mathrm{max}}w_\nu \nu (x)\}\right\}.$$
Then, do the hybrid MDL predictions always converge? This is equivalent to asking if the process of selecting a maximizing element eventually stabilizes. If stabilization does not occur, then hybrid MDL will necessarily fail as soon as the weights are not equal. A possible counterexample could consist of two measures the fraction of which oscillates perpetually around a certain value. We show that this can indeed happen, even for different reasons.
###### Example 18
Let $`𝒳`$ be binary, $`\mu (x)=_{i=1}^{\mathrm{}(x)}\mu _i(x_i)`$ and $`\nu (x)=_{i=1}^{\mathrm{}(x)}\nu _i(x_i)`$ with
$$\mu _i(1)=12^{2\frac{i}{2}}\mathrm{and}\nu _i(1)=12^{2\frac{i+1}{2}+1}.$$
Then one can easily see that $`\mu (111\mathrm{})=_1^{\mathrm{}}\mu _i(1)>0`$, $`\nu (111\mathrm{})=_1^{\mathrm{}}\nu _i(1)>0`$, and $`\frac{\nu (111\mathrm{})}{\mu (111\mathrm{})}`$ converges and oscillates. In fact, each sequence having positive measure under $`\mu `$ and $`\nu `$ contains eventually only ones, and the quotient oscillates.
###### Example 19
This example is a little more complex. We assume the uniform distribution $`\lambda `$ to be the true distribution. We now construct a positive martingale $`f()`$ that converges to $`\frac{3}{4}`$ with high probability and thereby oscillates infinitely often.
The martingale is defined on strings $`x`$ of successively increasing length. Of course, $`f(ϵ):=1`$. If $`f(x)`$ is already defined for strings of length $`n1`$, we extend the definition on strings of length $`n`$ in the following way: If $`f(x)>\frac{3}{4}`$, we set
$`f(x0)`$ $`:=`$ $`{\displaystyle \frac{3}{4}}2^{n2}\mathrm{and}`$
$`f(x1)`$ $`:=`$ $`2f(x)\left({\displaystyle \frac{3}{4}}2^{n2}\right).`$
This guarantees the martingale property $`f(x)=\frac{1}{2}\left(f(x0)+f(x1)\right)`$. If $`f(x)\frac{3}{4}`$ and $`f(x)\frac{3}{8}+2^{n3}`$, then we can similarly define
$`f(x0)`$ $`:=`$ $`2f(x)\left({\displaystyle \frac{3}{4}}+2^{n2}\right)\mathrm{and}`$
$`f(x1)`$ $`:=`$ $`{\displaystyle \frac{3}{4}}+2^{n2}.`$
However, if $`f(x)<\frac{3}{8}+2^{n3}`$, we cannot proceed in this way, since $`f`$ must be positive. Therefore, we set $`f(x0):=f(x1):=f(x)`$ in this case and call those $`x`$ “dead” strings. Strings that are not dead will be called “alive”. A few steps of the construction are shown in Figure 1. For example, it can be observed that the string 000 is dead, all other strings in the figure are alive.
It is obvious from the construction that $`f(x_{1:t})`$ is a martingale, it oscillates and converges to $`\frac{3}{4}`$ as $`t\mathrm{}`$ for all sequences $`x_<\mathrm{}`$ that always stay alive. The only thing we must show is that many sequences in fact stay alive.
###### Claim 20
We have $`\lambda \left(\{x_<\mathrm{}:t\text{ such that }x_{1:t}\mathrm{is}\mathrm{dead}\}\right)\frac{1}{4}`$.
Proof. After the $`n`$th step, i.e. when $`f`$ has been defined for strings of length $`n`$, $`f(x)`$ assumes the value
$$a_0^n=\frac{3}{4}2^{n2}$$
on a set of measure at most $`\frac{1}{2}`$. In the next step $`n+1`$, $`f`$ is defined to
$$a_1^n=\frac{3}{4}2^{n1}\left(1+\frac{1}{4}\right)$$
on half of the extended strings. Generally, in the $`k`$th next step, $`f`$ is defined to
$$a_k^n=\frac{3}{4}2^{n+k2}\left(\underset{j=0}{\overset{k}{}}2^{2j}\right)$$
on a $`2^k`$ fraction of the extended strings.
The extended strings stay alive as long as $`a_k^n\frac{3}{8}+2^{nk3}`$ holds. Some elementary calculations show that this is equivalent to $`kn`$. So precisely after $`n+1`$ additional steps, a fraction of $`2^{n1}`$ of the extended strings die.
We already noted that for $`A_n=\{x:\mathrm{}(x)=nf(x)=a_0^n\}`$, we have $`\lambda (A_n)\frac{1}{2}`$. Thus,
$$\lambda \left(\{x_<\mathrm{}:x_{1:n}A_n\mathrm{and}x_{1:2n+1}\mathrm{is}\mathrm{dead}\}\right)2^{n2}.$$
Hence, one can conclude
$$\lambda \left(\{x_<\mathrm{}:t\text{ such that }x_{1:t}\mathrm{is}\mathrm{dead}\}\right)\underset{n=1}{\overset{\mathrm{}}{}}2^{n2}=\frac{1}{4},$$
which proves the claim. $`\mathrm{}`$
We now define the measure $`\nu `$ by
$$\nu (x)=f(x)\lambda (x)=f(x)2^{\mathrm{}(x)},$$
and set the weights to $`w_\lambda =\frac{3}{7}`$ and $`w_\nu =\frac{4}{7}`$. Then this provides an example where the maximizing element never stops oscillating with probability at least $`\frac{3}{4}`$.
Both examples point out different possible reasons for failure of stabilizing. Example 18 works since the measure $`\mu `$ and $`\nu `$ are asymptotically very similar and close to deterministic. In contrast, in Example 19 stabilizing fails because of lack of independence: The quantity $`\nu (a|x)`$ strongly depends on $`x`$. In particular, one can note that even Markovian dependence may spoil the stabilization, since $`\nu (a|x)`$ only depends on the last symbol of $`x`$.
## 7 Stabilization
In the light of the previous section, it is therefore natural to ask when the maximizing element stabilizes (almost surely). Barron \[Bar85, BRY98\] has shown that this happens if all distributions in $`𝒞`$ are *asymptotically mutually singular*. Under this condition, the true distribution is even eventually identified almost surely.<sup>3</sup><sup>3</sup>3In general, stabilization does not imply that the true distribution is identified. Consider for instance a model class containing two measures: the true measure is concentrated on $`0^{\mathrm{}}`$ and has prior weight $`\frac{1}{8}`$, the other one assigns probability $`\nu (x_t=1)=2^t`$ independently of the past $`x_{<t}`$. Then the maximizing element will remain the incorrect distribution $`\nu `$, however with predictions rapidly converging to the truth.
The condition of asymptotic mutual singularity holds in many important cases, e.g. if the distributions are i.i.d. However, one cannot always build on it.<sup>4</sup><sup>4</sup>4Here is a simple example: let the true measure be Bernoulli($`\frac{1}{2}`$) and another measure be a product of Bernoullis with parameter rapidly converging to $`\frac{1}{2}`$. These distributions are not asymptotically mutually singular, nevertheless a.s. stabilization holds, as we will see. Therefore, in this section we give a different approach: In order to prevent stabilization, it is necessary that the ratio of two predictive distributions oscillates around the inverse ratio of the respective weights. Therefore, stabilization must occur almost surely if the ratio of two predictive distributions converges almost surely but is not *concentrated* in the limit. This is satisfied under appropriate conditions, as we will prove. We start with a general theorem which allows to conclude almost sure stabilization in a countable model class, if for any *pair* of models we have almost sure stabilization.
###### Theorem 21
Let $`𝒞`$ be a countable class of (semi)measures containing the true measure $`\mu `$. Assume that for each two $`\nu _1,\nu _2𝒞`$ the maximizing element chosen from $`\{\nu _1,\nu _2\}`$ eventually stabilizes almost surely. Then also the maximizing element chosen from all of $`𝒞`$ stabilizes almost surely.
Proof. It is immediate that the maximizing element chosen from any finite subset of $`𝒞`$ stabilizes almost surely. Now, for all $`\nu 𝒞`$ and $`c>0`$, define the set $`A_c^\nu `$ by
$$A_c^\nu =\{x_<\mathrm{}:t1\text{ such that }\frac{\nu (x_{1:t})}{\mu (x_{1:t})}c\}.$$
Then we have
$`\mu (A_c^\nu )`$ $`=`$ $`\mu \left({\displaystyle \{\mathrm{\Gamma }_x:\frac{\nu (x)}{\mu (x)}c\frac{\nu (x_{1:s})}{\mu (x_{1:s})}<cs<\mathrm{}(x)\}}\right)`$
$`=`$ $`{\displaystyle \{\mu (x):\frac{\nu (x)}{\mu (x)}c\frac{\nu (x_{1:s})}{\mu (x_{1:s})}<cs<\mathrm{}(x)\}}`$
$``$ $`{\displaystyle \{\frac{\nu (x)}{c}:\frac{\nu (x)}{\mu (x)}c\frac{\nu (x_{1:s})}{\mu (x_{1:s})}<cs<\mathrm{}(x)\}}`$
$`=`$ $`{\displaystyle \frac{1}{c}}{\displaystyle \{\nu (x):\mathrm{}\}}{\displaystyle \frac{1}{c}},`$
since $`\nu `$ is a (semi)measure and the set $`\{x𝒳^{}:\frac{\nu (x)}{\mu (x)}c\frac{\nu (x_{1:s})}{\mu (x_{1:s})}<cs<\mathrm{}(x)\}`$ is prefix-free. Let
$$B^\nu =\{x_<\mathrm{}:t1\text{ such that }\frac{w_\nu \nu (x_{1:t})}{w_\mu \mu (x_{1:t})}1\}=A_{(w_\mu /w_\nu )}^\nu ,$$
then $`\mu (B^\nu )\frac{w_\nu }{w_\mu }`$ holds. We arrange the (semi)measures $`\nu 𝒞`$ in an order $`\nu _1,\nu _2,\mathrm{}`$ such that the weights $`w_{\nu _1},w_{\nu _2},\mathrm{}`$ are descending. For each $`c1`$, we can now find an index $`k`$ and a set
$$𝒩_c=\{\nu _i:ik\}\text{ such that }\underset{\nu 𝒩_c}{}w_\nu \frac{w_\mu }{c}.$$
Defining $`B_c=_{\nu 𝒩^c}B^\nu `$, we get
$$\mu (B_c)\underset{\nu 𝒩^c}{}\frac{w_\nu }{w_\mu }\frac{1}{c}.$$
For all $`x_<\mathrm{}B^\nu `$, $`\nu `$ can never be the maximizing element. Therefore, for all $`x_<\mathrm{}B_c`$, there are only finitely many $`\nu 𝒩_c`$ having the chance of becoming the maximizing element at any time. By assumption, the maximizing element chosen from the finite set $`𝒞𝒩_c`$ stabilizes a.s. Thus, we conclude almost sure stabilization on the sequences in $`𝒳^{\mathrm{}}B_c`$. Since this holds for all $`B_c`$ and $`\mu (𝒳^{\mathrm{}}B_c)1`$ as $`c\mathrm{}`$, the maximizing element stabilizes with $`\mu `$-probability one. $`\mathrm{}`$
For the rest of this section, we assume that the model class $`𝒞`$ contains only proper measures. A measure $`\mu `$ is called factorizable if there are measures $`\mu _i`$ on $`𝒳`$ such that
$$\mu (x)=\underset{i=1}{\overset{\mathrm{}(x)}{}}\mu _i(x_i)$$
for all $`x𝒳^{}`$. That is, the symbols of sequences $`x_<\mathrm{}`$ generated by $`\mu `$ are independent. A factorizable measure $`\mu =\mu _i`$ is called uniformly stochastic, if there is some $`\delta >0`$ such that at each time $`i`$ the probability of all symbols $`a𝒳`$ is either 0 or at least $`\delta `$. That is
$$\mu _i(a)>0\mu _i(a)\delta \text{ for all }a𝒳\text{ and }i1.$$
(27)
In particular, all deterministic measures and all i.i.d. distributions are uniformly stochastic. Another simple example of a uniformly stochastic measure is a probability distribution which generates alternately random bits by fair coin flips and the digits of the binary representation of $`\pi =3.1415\mathrm{}`$
###### Lemma 22
Let $`\mu `$, $`\nu `$, and $`\stackrel{~}{\nu }`$ be factorizable and uniformly stochastic measures, where $`\mu `$ is the true distribution.
$`(i)`$ The maximizing element chosen from $`\mu `$ and $`\nu `$ stabilizes almost surely.
$`(ii)`$ If $`\mu `$ is not eventually always preferred over $`\nu `$ or $`\stackrel{~}{\nu }`$ (in which case we the maximizing element stabilizes a.s. by $`(i)`$), then the maximizing element chosen from $`\nu `$ and $`\stackrel{~}{\nu }`$ stabilizes almost surely.
Proof. We will show only $`(ii)`$, as the proof of $`(i)`$ is similar but simpler. So we assume that both $`\nu `$ and $`\stackrel{~}{\nu }`$ remain competitive in the process of choosing the maximizing element, and show that then maximizing element chosen from $`\nu `$ and $`\stackrel{~}{\nu }`$ stabilizes almost surely.
Let $`\nu =_i\nu _i`$, $`\stackrel{~}{\nu }=_i\stackrel{~}{\nu }_i`$, and $`X_i=\frac{\stackrel{~}{\nu }_i(x_i)}{\nu _i(x_i)}`$. The $`X_i`$ are independent random variables depending on the event $`x_<\mathrm{}`$. Moreover, both fractions $`\frac{\nu (x_{1:t})}{\mu (x_{1:t})}`$ and $`\frac{\stackrel{~}{\nu }(x_{1:t})}{\mu (x_{1:t})}`$ are martingales (with respect to $`\mu `$) and thus converge almost surely for $`t\mathrm{}`$. We are interested only in the events in
$$A_\nu =\{x_<\mathrm{}𝒳^{\mathrm{}}:\frac{\nu (x_{1:t})}{\mu (x_{1:t})}\mathrm{converges}\mathrm{to}\mathrm{a}\mathrm{value}>0\},$$
since otherwise $`\nu `$ eventually is no longer competitive. So we assume that $`\mu (A_\nu )>0`$, which implies $`\mu (A_\nu )=1`$ by the Kolmogorov zero-one-law (see e.g. \[CT88\]). Similarly, $`\mu (A_{\stackrel{~}{\nu }})=1`$ for the analogously defined set $`A_{\stackrel{~}{\nu }}`$. That is,
$$\underset{i=1}{\overset{t}{}}X_i=\frac{\stackrel{~}{\nu }(x_{1:t})}{\nu (x_{1:t})}=\frac{\stackrel{~}{\nu }(x_{1:t})}{\mu (x_{1:t})}/\frac{\nu (x_{1:t})}{\mu (x_{1:t})}$$
converges to a value $`>0`$ almost surely, and in particular $`0<X_i<\mathrm{}`$ a.s.
Now we will use the concentration function of a real valued random variable $`U`$,
$$Q(U,\eta )=\underset{u}{sup}\mu (uUu+\eta ),\eta 0.$$
(28)
This quantity was introduced by Lèvy, see e.g. \[Pet95\]. The concentration function is non-decreasing in $`\eta `$. Moreover, when two independent random variables $`U`$ and $`V`$ are added, we have \[Pet95, Lemma 1.11\]
$$Q(U+V,\eta )\mathrm{min}\{Q(U,\eta ),Q(V,\eta )\}\eta 0.$$
(29)
We first assume that the following set is unbounded:
$`B=\{{\displaystyle \underset{i=1}{\overset{n}{}}}\left(1Q(X_i,\eta )\right):n,\eta >0\}^+,`$ that is (30)
$`sup(B)=+\mathrm{},`$ (31)
We show that then $`\frac{\stackrel{~}{\nu }(x_{1:t})}{\nu (x_{1:t})}`$ (which converges a.s.) is not concentrated in the limit. That is, it converges to some given $`c>0`$, in particular to $`c=\frac{w_\nu }{w_{\stackrel{~}{\nu }}}`$, with $`\mu `$-probability zero. This shows that almost surely it does not oscillate around $`\frac{w_\nu }{w_{\stackrel{~}{\nu }}}`$.
Define independent random variables $`Y_i=\mathrm{ln}(X_i)`$. Let $`S_n:=_1^nY_i`$ and denote its almost everywhere existing limit by $`S=_1^{\mathrm{}}Y_i`$. The assertion is verified under condition (31), if we can show that the distribution of $`S`$ is not concentrated to any point since then also $`_1^{\mathrm{}}X_i=\mathrm{exp}(S)`$ is not concentrated to any point. In terms of the concentration function defined in (28), this reads $`Q(S,0)=0`$. According to (31), for each $`R>0`$, we find $`\eta >0`$ and $`n`$ such that $`_{i=1}^n\left(1Q(X_n,\eta )\right)>R`$. Then, because of $`X_i<\mathrm{}`$ (ignoring the measure-zero set where this may fail),
$$W=\underset{1in}{\mathrm{max}}X_i=\mathrm{max}\left\{\frac{\stackrel{~}{\nu }_i(x_i)}{\nu _i(x_i)}:1in\mathrm{and}\mu (x_i)>0\right\}$$
is finite. The mapping
$$(0,W]w\mathrm{ln}(w)(\mathrm{},\mathrm{ln}W]$$
is bijective and has derivative at least $`W^1`$. Let $`\stackrel{~}{\eta }=\frac{\eta }{W}`$. Then by definition of $`Y_i`$, we have $`Q(Y_i,\stackrel{~}{\eta })Q(X_i,\eta )`$ for $`1in`$ and consequently
$$\underset{i=1}{\overset{n}{}}\left(1Q(Y_i,\stackrel{~}{\eta })\right)>R.$$
By the Kolmogorov-Rogozin inequality (see \[Pet95, Theorem 2.15\]), there is a constant $`C`$ such that
$$Q(S_n,\stackrel{~}{\eta })C\left(\underset{i=1}{\overset{n}{}}\left(1Q(Y_i,\stackrel{~}{\eta })\right)\right)^{\frac{1}{2}}.$$
Thus, for each $`\epsilon >0`$, we can choose $`R`$ sufficiently large to guarantee $`CR^{\frac{1}{2}}<\epsilon `$. Then $`Q(S_n,\stackrel{~}{\eta })<\epsilon `$ for $`n`$ and $`\stackrel{~}{\eta }`$ as before. By (29) we conclude
$$Q(S,\stackrel{~}{\eta })=Q(S_n+\left(\underset{i=n+1}{\overset{\mathrm{}}{}}Y_i\right),\stackrel{~}{\eta })Q(S_n,\stackrel{~}{\eta })<\epsilon $$
and consequently $`Q(S,0)=0`$ since $`Q`$ is non-decreasing. This proves the assertion under assumption (31).
Now assume that $`B`$ is bounded, i.e. (31) does not hold. Then there is $`R>0`$ such that $`_1^n\left(1Q(X_i,\eta )\right)R`$ for all $`\eta >0`$ and $`n`$. Since the distribution of $`X_i`$ is a finite convex combination of point measures, for each $`i`$ there is an $`\eta >0`$ such that $`Q(X_i,\eta )=Q(X_i,0)`$ and thus $`_{i=1}^n\left(1Q(X_i,0)\right)R`$ for all $`n`$. Therefore, also $`_1^{\mathrm{}}\left(1Q(X_i,0)\right)R`$ holds. Since $`\stackrel{~}{\nu }_i(x_i)=c_i\nu _i(x_i)`$ is equivalent to $`X_i=c_i`$, this implies that there are constants $`c_i0`$ such that
$$\underset{i=1}{\overset{\mathrm{}}{}}\mu _i\{a:\stackrel{~}{\nu }_i(a)c_i\nu _i(a)\}R.$$
(32)
Next we argue that if $`c_i1`$ for infinitely many $`i`$, then either $`\nu `$ or $`\stackrel{~}{\nu }`$ is eventually not competitive. To verify this claim, let $`N_i=\{a:\stackrel{~}{\nu }_i(a)c_i\nu _i(a)\}`$ and $`M_i=𝒳N_i`$ and observe that $`\mu _i(N_i)<\delta `$ holds for sufficiently large $`i`$, since the sum (32) is bounded. On the other hand $`\mu `$ is uniformly stochastic, so there are no events of probability $`\mu _i(a)(0,\delta )`$, hence $`\mu _i(N_i)=0`$ and $`\mu _i(M_i)=1`$ for sufficiently large $`i`$. Now for these $`i`$, $`c_i>1`$ together with $`\nu _i(M_i)=1`$ implies the contradiction $`\stackrel{~}{\nu }_i(M_i)=c_i>1`$. So $`c_i>1`$ necessarily requires $`\nu _i(M_i)<1`$, hence $`\nu _i(M_i)1\delta `$, since $`\nu `$ is uniformly stochastic. If this happens infinitely often, then $`\nu `$ is eventually not competitive. A symmetric argument with $`\stackrel{~}{\nu }`$ holds for $`c_i<1`$.
The last paragraph shows that, if both $`\nu `$ and $`\stackrel{~}{\nu }`$ stay competitive, eventually $`\stackrel{~}{\nu }_i=\nu _i`$ holds a.s. In this case, $`\frac{\stackrel{~}{\nu }(x_{1:t})}{\nu (x_{1:t})}`$ is eventually constant, which completes the proof. $`\mathrm{}`$
###### Corollary 23
Let $`𝒞`$ be a countable class of factorizable and uniformly stochastic measures, then the maximizing element stabilizes almost surely.
Proof. This follows from Theorem 21 and Lemma 22. $`\mathrm{}`$
Lemma 22 and Corollary 23 are certainly not the only or the strongest assertions obtainable for stabilization. They rather give a flavor how a proof can look like, even if the distributions are not asymptotically mutually singular. On the other hand, the given result is optimal at least in some sense, as shown by the previous Examples 18 and 19. In the former example, $`\mu `$ is not uniformly stochastic but both $`\mu `$ and $`\nu `$ are factorizable, while in the latter one, $`\mu `$ is uniformly stochastic but $`\nu `$ is not factorizable.
The proof of Lemma 22 crucially relies on the independence assumption, which is necessary in order to use the Kolmogorov-Rogozin inequality. It is possible to relax this and require independent sampling only “every so often”. It is however not clear how to remove this condition completely.
## 8 Applications
In the following, we present some applications of the theory developed so far. We begin by stating general loss bounds. After that, three very general applications are discussed.
### 8.1 Loss bounds
So far we have only considered special loss functions, like the square loss, the Hellinger loss, or the relative entropy. We now show how these results, in particular the bounds for the Hellinger loss, imply regret bounds for *arbitrary loss functions*. (As we will see, square distance is not sufficient.) This parallels the bounds in \[Hut03a, Hut03b\]. The proofs are simplified, in particular Lemma 24 facilitates the analysis considerably. The reader should compare the results to the bounds for “prediction with expert advice”, e.g. \[CB97, HP05\].
In order to keep things simple, we restrict to binary alphabet $`𝒳=\{0,1\}`$ in this section. Our results extend to general alphabet by the techniques used in \[Hut03a\]. Consider a binary predictor having access to a belief probability $`\phi `$ depending on the current history, e.g. $`\phi (x_t=1|x_{<t})=\frac{1}{3}`$. Which actual prediction should he output, 0 or 1? We can answer this question if we know the *loss function*, according to which losses are assigned to the (wrong) predictions. Consider for example the 0/1 loss (also known as classification error loss), i.e. a wrong prediction gives loss of 1 and a right prediction gives no loss. Then we should predict 1 if our belief is $`\phi >\frac{1}{2}`$. This may be different under other loss functions. In general, we should predict in a *Bayes optimal* way: We should output the symbol with the least expected loss,
$$x^\phi :=\underset{\stackrel{~}{x}\{0,1\}}{\mathrm{arg}\mathrm{min}}\{(1\phi )\mathrm{}(0,\stackrel{~}{x})+\phi \mathrm{}(1,\stackrel{~}{x})\},$$
where $`\mathrm{}(x,\stackrel{~}{x})`$ is the loss incurred by prediction $`\stackrel{~}{x}`$ if the true symbol is $`x`$. In the following, we will restrict to *bounded* loss functions $`\mathrm{}(x,\stackrel{~}{x})[0,1]`$. Breaking ties in the above expression in an arbitrary deterministic way, the resulting prediction is *deterministic* for given $`\phi `$ and loss function $`\mathrm{}`$. If $`\mu `$ is the true distribution as usual, then let $`l_t^\phi :=_a\mu (a|x_{<t})\mathrm{}(a,x_t^\phi )`$ be the $`\mu `$-expected loss of the $`\phi `$-predictor. Then, by
$$L_{1:n}^\phi =𝐄[l_1^\phi +\mathrm{}+l_n^\phi ]=\underset{t=1}{\overset{n}{}}\mu (x_{<t})l_t^\phi (x_{<t})$$
we denote the cumulative $`\mu `$-expected loss of the $`\phi `$-predictor. With $`\phi `$ being the variants of the MDL predictor, we will bound the quantity $`\mathrm{\Delta }_{1:n}=L_{1:n}^\phi L_{1:n}^\mu `$, i.e. the cumulative *regret*, by an expression depending on $`L_{1:n}^\mu `$ and $`w_\mu ^1`$.
We admit arbitrary non-stationary loss functions $`\mathrm{}_{x_{<t}}`$ which may depend on the history. Our analysis considers the worst possible choice of loss functions and consists of three steps. First the cumulative regret bound is reduced to an instantaneous regret bound (Lemma 24). Then the instantaneous bound is reduced to a bound in terms of special functions of $`\mu `$ and $`\phi `$ (Lemma 25). Finally, the bound for the special functions is given (Lemma 26).
###### Lemma 24
Assume that some $`\phi `$-predictor satisfies the instantaneous regret bound $`\delta _t=l_t^\phi l_t^\mu 2h_t+2\sqrt{2h_tl_t^\mu }`$, where $`h_t=h_t(\mu ,\phi )`$ is the Hellinger distance of the instantaneous predictive probabilities (23). Then the cumulative $`\phi `$-regret is bounded in the same way:
$$\mathrm{\Delta }_{1:n}=L_{1:n}^\phi L_{1:n}^\mu 2H_{1:n}(\mu ,\phi )+2\sqrt{2H_{1:n}(\mu ,\phi )L_{1:n}^\mu }.$$
This and the following lemma hold with arbitrary constants, the choices $`2`$ and $`2\sqrt{2}`$ are the smallest ones for which Lemma 26 is true. Note that if the Hellinger distance is replaced by the relative entropy, then $`2\sqrt{2}`$ may be replaced by $`2`$. Thus, normalized dynamic MDL and Bayes mixture admit smaller bounds, compare \[Hut03a\]. However, this is not true for the other MDL variants, as we have no relative entropy bound there.
Proof. The key property is the *super-additivity* of the bound. A function $`f:[0,\mathrm{})^2[0,\mathrm{})`$ is said to be super-additive if
$$f(x_1+x_2,y_1+y_2)f(x_1,y_1)+f(x_2,y_2).$$
The function $`(H,L)\sqrt{HL}`$ satisfies this condition. We now use an inductive argument. Assume $`\mathrm{\Delta }_{2:n}^02H_{2:n}^0+2\sqrt{2H_{2:n}^0L_{2:n}^{\mu ,0}}`$, where the summation starts at $`t=2`$ and the superscript $`0`$ indicates that the first symbol of the sequence was $`0`$. Let the same hold for the first symbol $`1`$. Writing $`\mu _1=\mu (1|ϵ)`$ and using $`\delta _12h_1+2\sqrt{2h_1l_1^\mu }`$, we obtain
$`\mathrm{\Delta }_{1:n}=\delta _1+(1\mu _1)\mathrm{\Delta }_{2:n}^0+\mu _1\mathrm{\Delta }_{2:n}^1`$
$``$ $`2\left[h_1+\sqrt{2h_1\mathrm{}_1}+(1\mu _1)\left(H_{2:n}^0+\sqrt{2H_{2:n}^0L_{2:n}^{\mu ,0}}\right)+\mu _1\left(H_{2:n}^1+\sqrt{2H_{2:n}^1L_{2:n}^{\mu ,1}}\right)\right]`$
$``$ $`2\left[H_{1:n}+\sqrt{2h_1\mathrm{}_1}+\sqrt{2\left((1\mu _1)H_{2:n}^0+\mu _1H_{2:n}^1\right)\left((1\mu _1)L_{2:n}^{\mu ,0}+\mu _1L_{2:n}^{\mu ,1}\right)}\right]`$
$``$ $`2H_{1:n}+2\sqrt{2H_{1:n}L_{1:n}^\mu }.`$
Here, the first inequality is the induction hypothesis together with the instantaneous bound, the second bound is Cauchy-Schwarz’s inequality, and the last estimate is the super-additivity. $`\mathrm{}`$
###### Lemma 25
Assume that some $`\phi `$-predictor satisfies $`\stackrel{~}{\delta }2h+2\sqrt{2h\stackrel{~}{\mathrm{}}}`$ for all $`\mu ,\phi [0,1]`$, with the Hellinger distance $`h=h(\mu ,\phi )`$ and the special functions $`\stackrel{~}{\delta }(\mu ,\phi )`$ and $`\stackrel{~}{\mathrm{}}(\mu ,\phi )`$ defined in the following way, where we slightly abuse notation and abbreviate $`\mu =\mu (1|\mathrm{})`$ and $`\phi =\phi (1|\mathrm{})`$:
$$\stackrel{~}{\delta }=\frac{|\phi \mu |}{\mathrm{max}\{\phi ,1\phi \}}\text{ and }\stackrel{~}{\mathrm{}}=\{\begin{array}{cc}\mu \text{if}\hfill & \mu \phi \frac{1}{2},\hfill \\ \mu (1\phi )/\phi \text{if}\hfill & \mu \phi \frac{1}{2}\phi ,\hfill \\ 1\mu \text{if}\hfill & \frac{1}{2}\phi \mu ,\hfill \\ (1\mu )\phi /(1\phi )\text{if}\hfill & \phi \mu \phi \frac{1}{2}.\hfill \end{array}$$
Then for arbitrary bounded loss function $`\mathrm{}:\{0,1\}^2[0,1]`$, we have
$$\delta 2h+2\sqrt{2hl^\mu }.$$
(33)
Proof. First we show that we may assume $`\mathrm{}(0,0)=\mathrm{}(1,1)=0`$, i.e. we do not incur loss for correct predictions. To this end, consider the modified loss function $`\mathrm{}^{}(x,\stackrel{~}{x})=\mathrm{}(x,\stackrel{~}{x})\mathrm{}(x,x)`$ and assume w.l.o.g $`\mathrm{}^{}(x,\stackrel{~}{x})[0,1]`$. Then it is not hard to see that the regrets under the original and the modified loss functions coincide, while the expected loss of the $`\mu `$-predictor clearly decreases with the modified loss function. Thus, (33) holds for $`\mathrm{}`$ if it holds for $`\mathrm{}^{}`$. Hence we may assume $`\mathrm{}(0,0)=\mathrm{}(1,1)=0`$. For each possible outcome $`x\{0,1\}`$, we abbreviate $`\mathrm{}^x=\mathrm{}(x,1x)`$.
Now assume w.l.o.g. $`\mu \phi `$. In order to show the assertion, we need to consider the cases in the definition of $`\stackrel{~}{\mathrm{}}`$ separately. We show this only for the first case, i.e. $`\mu \phi \frac{1}{2}`$. Then $`l^\mu =\mu \mathrm{}^1`$, $`l^\phi =(1\mu )\mathrm{}^0`$. We assume that the $`\mu `$-predictor outputs 0 and the $`\phi `$-predictor 1, otherwise they give the same prediction and the $`\phi `$-predictor has no regret at all. This condition is equivalent to $`\mathrm{}^0=\mathrm{}^1\frac{u}{1u}`$ for some $`u[\mu ,\phi ]`$. We consider the worst case by maximizing $`l^\phi `$, i.e. choosing $`u`$ as large as possible. For this $`u=\phi `$, we obtain $`\mathrm{}^0=\mathrm{}^1\frac{\phi }{1\phi }`$ and
$$\delta =\mathrm{}^1[\frac{(1\mu )\phi }{1\phi }\mu ]=\mathrm{}^1\stackrel{~}{\delta }\mathrm{}^1[2h+2\sqrt{2h\stackrel{~}{\mathrm{}}}]2h+2\sqrt{2h\mathrm{}^1\stackrel{~}{\mathrm{}}}2h+2\sqrt{2hl^\mu },$$
showing (33) provided that $`\mu \phi \frac{1}{2}`$. The other cases are shown similarly. $`\mathrm{}`$
###### Lemma 26
The bound $`\stackrel{~}{\delta }2h+2\sqrt{2h\stackrel{~}{\mathrm{}}}`$ holds for all $`\mu ,\phi [0,1]`$, with the functions $`\stackrel{~}{\delta },\stackrel{~}{\mathrm{}}:[0,1]^1[0,1]`$ as defined in Lemma 25.
The technical and not very interesting proof of this lemma is omitted. The careful reader may check the assertion numerically or graphically, as it is just the boundedness of some function on the unit square. We remark that the bound does *not* hold if the Hellinger distance is replaced by the quadratic distance, not even with larger constants.
###### Theorem 27
For arbitrary non-stationary loss function which is bounded in $`[0,1]`$ and known to the MDL predictors, their respective losses are bounded by
$$L_{1:n}^{\varrho _{\mathrm{norm}}},L_{1:n}^\varrho ,L_{1:n}^{\varrho ^{\mathrm{static}}},L_{1:n}^{\varrho _{\mathrm{norm}}^{\mathrm{static}}}L_{1:n}^{\mu _{\mathrm{norm}}}+2\sqrt{2cL_{1:n}^{\mu _{\mathrm{norm}}}w_\mu ^1}+2cw_\mu ^1,$$
where the constant $`c=2,8,21,\text{ or }32`$, according to which MDL predictor is used (compare Corollary 14).
Proof. This follows from the above three lemmas and from $`H_{1:n}cw_\mu ^1`$ (Corollary 14). $`\mathrm{}`$
This shows in particular that, regardless of the loss function, the average expected per-round regret tends to zero. Again, the direct practical relevance of the bounds is limited because of the potentially huge $`w_\mu ^1`$.
### 8.2 Classification
Transferring our results to pattern classification is very easy. All we have to do is to add *inputs* to our models. That is, we consider an arbitrary input space $`𝒰`$ and (as before) a finite observation or output space $`𝒳`$. A model is now a *measure*
$$\nu (x|u)[0,1],x𝒳,u𝒰,\text{ where }\underset{x𝒳}{}\nu (x|u)=1\text{ for all }u𝒰.$$
That is, we have a distribution which is conditionalized to the input. We restrict our discussion to measures, since there is no motivation to consider semimeasures for classification. The definition of a model does not include history dependence. There is no loss of generality: We may include the history in the arbitrary input space.
Transferring the proofs in the previous sections to the present setup is straightforward. We therefore obtain immediately the following corollaries.
###### Corollary 28
Let $`𝒞`$ be a countable set of classification models containing the true distribution $`\mu `$. Then for any sequence of inputs $`u_<\mathrm{}𝒰`$, we have
$$\begin{array}{ccc}_t𝐄_a\left(\mu (a|u_t)\varrho _{\mathrm{norm}}(a|u_t,u_{<t},x_{<t})\right)^2\hfill & & 2w_\mu ^1,\hfill \\ _t𝐄_a\left(\mu (a|u_t)\varrho (a|u_t,u_{<t},x_{<t})\right)^2\hfill & & 8w_\mu ^1,\hfill \\ _t𝐄_a\left(\mu (a|u_t)\varrho ^{\mathrm{static}}(a|u_t,u_{<t},x_{<t})\right)^2\hfill & & 21w_\mu ^1.\hfill \end{array}$$
(Note that although each single model formally does not depend on the history, the MDL estimators necessarily do.)
We need not consider the normalized static variant here, since all models are measures anyway. If there is a distribution over $`𝒰`$, the result therefore also holds in expectation over the inputs. An analogue of Corollary 23 is obtained as easily. If the inputs are i.i.d., which is usually assumed for classification, then the two conditions of factorizability and uniform stochasticity are trivially satisfied. Therefore, the true distribution $`\mu `$ is eventually discovered by MDL almost surely. Note that in this case, the distributions are also asymptotically mutually singular, so that the assertion also follows from Barron’s \[Bar85\] earlier result.
Note that again, the assumption $`\mu 𝒞`$ is essential. In practical applications, if this is not clear, it may be therefore favorable to choose a different method having guarantees without this condition, compare \[GL04\].
### 8.3 Regression
We may also apply our results in the regression setup, that is for predicting continuous densities. Our use of the term regression is a bit non-standard here, since it normally refers to just estimating the mean of some prediction, where the distribution is often assumed to be Gaussian. Again the assumption $`\mu 𝒞`$ is essential, so that in practice some other method not relying on it might be preferred.
Continuous densities cause some additional difficulties. The observation space is now $``$. This implies in particular that, like for the loss bounds, the square distance is no longer appropriate for our purpose<sup>5</sup><sup>5</sup>5To see this, define a distribution $`f`$ by its density $`f_n=\frac{n}{3}\chi _{[\frac{1}{n},0]}+\frac{2n}{3}\chi _{(0,\frac{1}{n}]}`$, where $`\chi `$ is the characteristic function of an interval. Let $`\stackrel{~}{f}(x)=f(x)`$, then the quadratic distance is $`(f\stackrel{~}{f})^2𝑑x=\frac{2n}{9}\stackrel{n\mathrm{}}{}\mathrm{}`$, whereas the relative entropy $`f\mathrm{ln}(f/\stackrel{~}{f})𝑑x=\frac{\mathrm{ln}2}{3}`$ is constant. (note that our use of the squared error is completely different from the standard use in regression). So we will use the Hellinger distance instead, defined similarly to (23) by
$$h(f,\stackrel{~}{f})=\left(\sqrt{f(x)}\sqrt{\stackrel{~}{f}(x)}\right)^2𝑑x\text{ for integrable }f,\stackrel{~}{f}:[0,\mathrm{}).$$
(34)
Accordingly, $`H_{1:n}(\mu ,\phi )=_t𝐄h(\mu (|u_t),\phi (|u_t,u_{<t},x_{<t}))`$ is the cumulative Hellinger distance of two predictive distributions $`\mu `$ and $`\phi `$. Similarly as in (25) and (26), the Hellinger distance is bounded by the (continuous) relative entropy and absolute distance. This shows in particular that the integral (34) exists.
We now consider a countable class $`𝒞`$ of models that are functions $`\nu `$ from $`𝒰`$ to *uniformly bounded probability densities* on $`𝒳=`$. That is, there is some $`C>0`$ such that
$$0\nu _i(x|u)C\text{ and }_{\mathrm{}}^{\mathrm{}}\nu _i(x|u)𝑑x=1\text{ for all }i1,u𝒰,\text{ and }x.$$
(35)
for all $`i1`$, $`u𝒰`$, and $`x`$. The MDL estimator is then defined as the element which maximizes the *density*, $`\nu ^{}=\mathrm{arg}\mathrm{max}_{\nu 𝒞}\{w_\nu \nu (x_{1:n}|u_{1:n})\}`$. The uniform boundedness condition asserts that the MDL estimator exists. It may be relaxed, provided that the MDL estimator remains well-defined, such as for a family of Gaussian densities which tend to the point measure.
With these definitions, the proofs of the theorems for static and dynamic MDL can be adapted. Since the triangle inequality holds for the Hellinger distance $`\sqrt{H^2}`$, we obtain the following.
###### Corollary 29
Let $`𝒞`$ be a countable model class according to (35), containing the true distribution $`\mu `$. Then for any sequence of inputs $`u_<\mathrm{}𝒰`$, we have $`H_{1:n}(\mu ,\varrho _{\mathrm{norm}})2w_\mu ^1`$, $`H_{1:n}(\mu ,\varrho )8w_\mu ^1`$, and $`H_{1:n}(\mu ,\varrho ^{\mathrm{static}})21w_\mu ^1`$.
We may apply this for example to model classes with Gaussian noise, concluding that the mean and the variances converge to the true values, see \[PH05\] for an example. It is not immediately clear how to obtain an analogue of Corollary 23 for continuous densities.
### 8.4 Universal Induction
Since the assertions on static and dynamic MDL have been proven generally for *semimeasures*, we may apply them to the universal setup. Here $`𝒞=`$ is the countable set of all lower semicomputable (= enumerable) semimeasures on $`𝒳^{}`$. So $``$ contains stochastic models in general, and in particular all models for computable deterministic sequences. There is a one-to-one correspondence of $``$ to the class of all programs on some fixed universal monotone Turing machine $`U`$, see e.g. \[LV97\]. We will assume programs to be *binary*, in contrast to outputs, which are strings $`x𝒳^{}`$. This relation defines in particular the complexities and weights of each $`\nu `$ by
$$Kw(\nu )=\text{length of the program for }\nu \text{ on }U,\text{ and }w_\nu =2^{Kw(\nu )}.$$
(36)
We call these weights the *canonical weights*. They satisfy $`w_\nu >0`$ for all $`\nu `$ and $`_\nu w_\nu 1`$.
An enumerable semimeasure which dominates all other enumerable semimeasures is called universal. The Bayes mixture $`\xi `$ defined in (2) has this property. One can show that $`\xi `$ is equal within a multiplicative constant to Solomonoff’s prior \[Sol64, Eq. (7)\], which is the a priori probability that (some extension of) a string $`x`$ is generated provided that the input of $`U`$ consists of fair coin flips. That is
$$\xi (x)\stackrel{\times }{=}M(x)=\underset{pminimal:U\left(p\right)=x}{}2^{\mathrm{}(p)}\text{ for all }x𝒳^{}.$$
Here, we use the notations
$`f\stackrel{+}{}g:fg+O(1),`$ $`f\stackrel{+}{=}g:f\stackrel{+}{}gg\stackrel{+}{}f,`$
$`f\stackrel{\times }{}g:fgO(1),`$ $`f\stackrel{\times }{=}g:f\stackrel{\times }{}gg\stackrel{\times }{}f.`$
The MDL definitions in Section 3 directly transfer to this setup. All bounds on the cumulative square loss (subsumed in Corollary 14) therefore apply to $`\varrho =\varrho _{[]}`$. The necessary assumption now reads that $`\mu `$ must be a recursive (= computable) measure. Also, Theorem 1 implies Solomonoff’s important universal induction theorem.
In addition to $``$, we also consider the set of all recursive measures $`\stackrel{~}{}`$ together with the same canonical weights (36). We define $`\stackrel{~}{\xi }=\xi _{[\stackrel{~}{}]}`$ and $`\stackrel{~}{\varrho }=\varrho _{[\stackrel{~}{}]}`$. Then $`\stackrel{~}{\varrho }(x)\stackrel{~}{\xi }(x)\xi (x)`$ and $`\varrho (x)\xi (x)`$ for all $`x𝒳^{}`$ is immediate. It is straightforward that $`\xi (x)\stackrel{\times }{}\varrho (x)`$ since $`\xi `$. Moreover, for any string $`x𝒳^{}`$, define the *monotone complexity* $`Km(x)=\mathrm{min}\{\mathrm{}(p):U(p)=x\}`$ as the length of the shortest program such that $`U`$’s output starts with $`x`$. The following assertion holds.
###### Proposition 30
We have $`K\stackrel{~}{\varrho }\stackrel{+}{}Km`$.
Proof. We must show that given a string $`x𝒳^{}`$ and a recursive measure $`\nu `$ (which in particular may be the MDL descriptor $`\nu ^{}(x)`$) it is possible to specify a program $`p`$ of length at most $`Kw(\nu )+K\nu (x)+c`$ that outputs a string starting with $`x`$, where constant $`c`$ is independent of $`x`$ and $`\nu `$.
Consider all strings $`y_i𝒳^n`$ ($`1i|𝒳|^n`$) of length $`n=\mathrm{}(x)`$ arranged in lexicographical order. Each $`y_i`$ has measure $`P_i=\mu (y_i)`$. Let $`S_i`$ be the cumulated measures: $`S_0=0`$ and $`S_i=_{k=1}^iP_k`$. Let $`j`$ be the index of $`x`$, i.e. $`x=y_j`$. Then, the interval $`[S_{j1},S_j)[0,1)`$ has measure $`P_j`$ and therefore contains exactly one $`\mathrm{log}_2P_j`$-bit number $`z[S_{j1},S_j)`$. We describe $`x`$ with the number $`z`$, this is known as arithmetic encoding (see e.g. \[CT91\]). The coding is injective since $`[S_{i1},S_i)`$ and $`[S_{k1},S_k)`$ are disjoint for $`ik`$.
In order to decode $`z`$, we may descend the $`|𝒳|`$-ary tree of all possible strings $`y`$, first considering strings of length one, then of length two, etc. For each possible string $`y`$, we can determine its binary code by approximating $`\nu (x)`$ sufficiently accurately. Eventually we will find $`z`$, then we print the current $`y`$. At this stage, $`y`$ might be only a prefix of $`x`$, since an extension of $`y`$ might have a measure very close to $`y`$ and thus map to the same code $`z`$. Therefore we continue the procedure until all codes starting with $`z`$ are proper extensions of $`z`$ (which may never be the case, then the algorithm runs forever). In each step, the appropriate additional symbol is written on the output tape. The resulting output will be $`x`$ or some extension of $`x`$.
This algorithm can be specified in a constant $`c^{}`$ number of bits. The description of $`\nu `$ needs another $`Kw(\nu )`$ bits. Finally, $`z`$ has length $`\mathrm{log}_2P_j\mathrm{log}_2\nu (x)+1`$. Thus, the overall description has length $`Kw(\nu )+K\nu (x)+c`$ as required. $`\mathrm{}`$
It is also possible to prove the proposition indirectly using \[LV97, Thm.4.5.4\]. This implies that $`Km(x)\stackrel{+}{}Kw(\nu )+K\nu (x)`$ for all $`x𝒳^{}`$ and all recursive measures $`\nu \stackrel{~}{}`$. Then, also $`Km(x)\stackrel{+}{}\mathrm{min}\{Kw(\nu )+K\nu (x)\}=K\stackrel{~}{\varrho }(x)`$ holds.
So together with the above observations, we have
$$Km(x)\stackrel{+}{}K\stackrel{~}{\varrho }(x)\stackrel{+}{}K\stackrel{~}{\xi }(x)\stackrel{+}{}K\varrho (x)\stackrel{+}{=}KM(x).$$
(37)
On the other hand, there is a deep result in Algorithmic information theory which states that an exact coding theorem does *not* hold on continuous sample space, $`Km(x)\stackrel{+}{}KM(x)`$ \[Gác83\]. Therefore, at least one of the above $`\stackrel{+}{}`$ must be proper.
###### Problem 31
Which of the two inequalities $`K\stackrel{~}{\varrho }(x)\stackrel{\times }{}K\stackrel{~}{\xi }(x)`$ and $`K\stackrel{~}{\xi }(x)\stackrel{\times }{}K\varrho (x)`$ is proper (or are both)?
The proof in \[Gác83\] is very subtle, and the phenomenon is still not completely understood. There is some hope that by answering Problem 31, one arrives at a better understanding of the continuous coding theorem and even at a simpler proof for its failure.
## 9 Discussion
In this last section, we recapitulate the main achievements of this work and discuss their philosophical and practical consequences. In the first place, we have shown that if two-part MDL is used for predicting a stochastic sequence, then the predictive probabilities converge to the true ones in mean sum, provided that the distribution generating the sequence is contained in the model class. The two most important implications are almost sure convergence and loss bounds for arbitrary loss functions.
The guaranteed convergence is slow in general: All bounds depend linearly on $`w_\mu ^1`$, the inverse of the prior weight of the true distribution. For large model classes, this number must be regarded too huge to be relevant for practical applications. Examples show that this bound is sharp. This is in contrast to the exponentially smaller corresponding bound for the Bayes mixture. The latter predictor however is often computationally more expensive to approximate in practice. We believe that this principally indicates that with MDL, some care has to be taken when choosing the model class and the prior. Conditions which are sufficient for fast convergence have been given for instance in \[Ris96, BRY98, PH04b\]. It remains a major challenge to generalize these results in order to obtain fast convergence under assumptions that are as weak as possible. In particular for universal induction, this question is interesting and possibly difficult. Even when considering only computable Bernoulli distributions endowed with a universal prior, fast convergence possibly holds for many environments, but maybe not for all \[PH04b\]. We also need to distinguish how the large error cumulates. Either the instantaneous error remains significant for a long time, which is critical, or the instantaneous error drops just too slowly to be summable, e.g. as $`O(\frac{1}{n})`$, which is tolerable. We have seen instances for both cases; compare the discussion after Example 15. In this light, the cumulative error might not be the right quantity to assess convergence speed.
The main results have been shown under the only assumption that the data generating process is contained in the model class. This condition is essential in general, as \[GL04\] shows that in its absence MDL can fail dramatically. In the universal setup, the assumption merely requires that the data is generated in some (probabilistically) computable way. This is a very weak condition. Laplace, Zuse \[Zus67\] and successors argue that nature operates in a computable way, and consequently *all* thinkable data satisfies the assumption. On the other hand, predicting with a universal model is computationally very expensive. In particular it is provably infeasible if the thesis of computable nature holds. Despite these practical problems, the theory of universal prediction is valuable since it explores the limits of computational induction.
Acknowledgements. Thanks to Peter Grünwald and an anonymous reviewer for their very valuable comments and suggestions. This work was supported by SNF grant 2100-67712.02.
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# CERN-PH-TH/2007-173 GZK Photons as Ultra High Energy Cosmic Rays
## I I. Introduction
The cosmic rays with energies beyond the Greisen-Zatsepin-Kuzmin (GZK) cutoff gzk at $`4\times 10^{19}`$ eV present a challenging outstanding puzzle in astroparticle physics and cosmology agasa ; hires . Nucleons cannot be confined to our galaxy for energies above the “ankle”, i.e. above 10<sup>18.5</sup> eV. This and the absence of a correlation of arrival directions with the galactic plane indicate that, if nucleons are the primary particles of the ultra high energy cosmic rays (UHECR), these nucleons should be of extragalactic origin. However, nucleons with energies above $`5\times 10^{19}`$ eV could not reach Earth from a distance beyond 50 to 100 Mpc 50Mpc because they scatter off the cosmic microwave background (CMB) photons with a resonant photoproduction of pions: $`p\gamma \mathrm{\Delta }^{}N\pi `$, where the pion carries away $`20\%`$ of the original nucleon energy. The mean free path for this reaction is only $`6`$ Mpc. Photons with comparable energy pair-produce electrons and positrons on the radio background and, likewise, cannot reach Earth from beyond 10 to 40 Mpc 40Mpc (although the photon energy-attenuation length is uncertain, due to the uncertainties in the spectrum of the absorbing radio background). There only few known astrophysical sources within those distances that could produce such energetic particles, but they are not located along the arrival directions of observed cosmic rays.
Intervening sheets of large scale intense extra galactic magnetic fields (EGMF), with intensities $`B0.11\times 10^6`$ G, could provide sufficient angular deflection for protons to explain the lack of observed sources in the directions of arrival of UHECR. However, recent realistic simulations of the expected large scale EGMF, show that strong deflections could only occur when particles cross galaxy clusters. Except in the regions close to the Virgo, Perseus and Coma clusters the obtained magnetic fields are not larger than 3$`\times 10^{11}`$ G dolag2004 and the deflections expected are not important (however see Ref. Sigl:2004yk ).
Whether particles can be emitted with the necessary energies by astrophysical accelerators, such as active galactic nuclei, jets or extended lobes of radio galaxies, or even extended object such as colliding galaxies and clusters of galaxies, is still an open question. The size and possible magnetic and electric fields of these astrophysical sites make it plausible for them to produce UHECR at most up to energies of $`10^{21}`$ eV. Larger emission energies would require a reconsideration of possible acceleration models or sites.
Heavy nuclei are an interesting possibility for UHECR primaries, since they could be produced at the sources with larger maximum energies (proportional to their charges) and would more easily be deflected by intervening magnetic fields. On the other hand, both AGASA and HiRes data favor a dominance of light hadrons, consisting with being all protons, in the composition of UHECR above 10<sup>19</sup> eV. However, we should keep in mind that the inferred composition is sensitive to the interaction models used. Assuming a proton plus iron composition, HiRes Stereo data show a constant or slowly changing composition of 80% protons and 20% iron nuclei between 10<sup>18.0</sup> eV and 10<sup>19.4</sup> eV. This is consistent with the change in composition from heavy to light in the 10<sup>17</sup> eV to 10<sup>18</sup> eV range found by HiRes Prototype hires\_composition . HiRes monocular data show a 90% proton composition between 10<sup>17.6</sup> eV and 10<sup>20</sup> eV hires\_composition\_fit . Similar results were found by AGASA, which produced bounds on the iron fraction (again assuming an iron plus proton composition) of 14 $`(+16,14)`$% and 30 $`(+7,6)`$% above 10<sup>19.0</sup> eV and 10<sup>19.25</sup> eV respectively, and 1 $`\sigma `$ upper bound of 66% above 10<sup>19.5</sup> eV agasa\_composition\_2 .
In fact, a galactic component of the UHECR flux, which could be important up to energies 10<sup>19</sup> eV, should consists of heavy nuclei, given the lack of correlation with the galactic plane of events at this energy (outside the galactic plane, galactic protons would be deflected by a maximum of 15-20<sup>o</sup> at this energies galactic\_magn\_field ). For nuclei the dominant energy loss process is photodissociation through scattering with the infra-red background below 10<sup>20</sup> eV puget and with the CMB above $`10^{20}`$ eV, and pair creation on the CMB in a small energy interval around 10<sup>20</sup> eV (at energies for which the typical CMB photon energy in the rest frame of the nucleus is above threshold, i.e. above 1 MeV, but below the peak of the giant resonance, 10-20 MeV) epele . The typical attenuation length in the energy range 4$`\times 10^{19}`$ to $`1\times 10^{20}`$ eV changes from several 10<sup>3</sup> Mpc for iron and silicon to be comparable to that of nucleons for helium epele ; bertone . At energies above $`1\times 10^{20}`$ eV, the attenuation length of heavy nuclei decreases and becomes less than 10 Mpc at about $`3\times 10^{20}`$ eV for iron, $`2\times 10^{20}`$ eV for silicon and $`1\times 10^{20}`$ eV for carbon (see for example Fig. 1 of Ref. bertone ). In the realistically low EGMF of Ref. dolag2004 , most of the heavy nuclei with $`E>10^{20}`$ eV reaching us from more than 10 Mpc away with energies above those mentioned would disintegrate into protons with energy ($`1/A`$) of the original nucleus energy, where $`A`$ is the atomic number (this is 1/56 of the original energy for iron nuclei). Note also that the same photodissociation processes can destroy heavy nuclei near their sources, if the intensity of the infrared background near the sources is large enough. One should not forget that all UHECR above $`10^{18}`$ eV could be due to extragalactic protons berezinsky2002 .
The GZK cutoff at $`4\times 10^{19}`$ eV seems not to be present in the data of the AGASA ground array agasa but it appears in the data of the HiRes air fluorescence detector hires . In any case, there are events above the GZK cutoff, even in the HiRes data set, and these remain unexplained since the local Universe ($`100`$ Mpc) is devoid of strong candidate sources in the direction the events point to, and also of the large magnetic fields which could deflect the incoming particles significantly. Due to the limited statistics and different systematic errors of both experiments the discrepancy between them is not very significant. However, the presence or absence of the GZK cutoff remains an open question. This controversy will be solved conclusively by the Pierre Auger Observatory Auger , a hybrid combination of charged particles detectors and fluorescence telescopes, perhaps within the next one or two years.
The analysis of the muon content in air showers has been used by AGASA to reject photon dominance in UHECR above 10<sup>19</sup> eV agasa\_composition\_1 ; agasa\_composition\_2 . Assuming a composition of protons plus photons, AGASA quotes upper limits for the photon ratio of 34%, 59% and 63% at 10<sup>19</sup> eV, 10<sup>19.25</sup> eV and 10<sup>19.5</sup> eV respectively at the 95% confidence level agasa\_composition\_2 , and even above 10<sup>20</sup> eV they find no indication that the events they observe are mostly photons agasa\_composition\_1 . Also a reanalysis of horizontal showers at Haverah Park concluded that photons cannot constitute more that 50% of the UHECR above 4$`\times 10^{19}`$ eV haverah .
The GZK process produces pions. From the decay of $`\pi ^\pm `$ one obtains neutrinos. These “GZK neutrinos” have been extensively studied, from 1969 bere onward (see for example reviewGZKneutrinos ; reviewGZKneutrinos2 and references therein), and constitute one of the main high energy signals expected in neutrino telescopes, such as ICECUBE ICECUBE ANITA ANITA and SALSA SALSA or space based observatories such as EUSO EUSO and OWL OWL . From the decay of $`\pi ^0`$ we obtain photons, “GZK photons”, with about 0.1 of the original proton energy, which have been known to be a subdominant component of the UHECR since the work of Wdowczyk et al. in the early 1970’s wdowczyk . In 1990 it was suggested that if the extragalactic radio background and magnetic fields are small ($`B<3\times 10^{11}`$ G) GZK photons could dominate over protons and explain the super-GZK events Aharonian1990 . The dependence of the GZK photon flux on extragalactic magnetic fields was later studied in Ref. SiglOlinto95 . The argument of Ref. Aharonian1990 and its dependence on extragalactic magnetic fields was again discussed astro\_photons in connection with the possible correlation of UHECR arrival directions with BL Lacertae objects Tinyakov:2001nr . However, to our knowledge, no complete study of the expected fluxes of GZK photons was done so far, including their dependence on the initial proton fluxes, distribution of proton sources and UHECR spectrum, besides intervening backgrounds.
With the advent of the Pierre Auger Observatory, we expect to have in the near future the high statistic data that may allow to study a subdominant component of UHECR consisting of photons. The GZK photons provide a complementary handle to GZK neutrinos and other signatures to try to determine the spectrum and composition of the UHECR. The flux of GZK photons is necessarily correlated with the flux of GZK neutrinos, although the former is affected by the radio background and EGMF values which do not affect the latter.
In this paper we show that if the UHECR are mostly protons, depending on the UHECR spectrum assumed, the slope of the proton flux, distribution of sources and intervening backgrounds, between $`10^4`$ and $`10^2`$ of the UHECR above $`10^{19}`$ eV and between $`10^5`$ and $`0.6`$ of the UHECR above $`10^{20}`$ eV are GZK photons, the range being much higher for the AGASA spectrum than for the HiRes spectrum (see Fig. 17). Detection of these photons would open the way for UHECR photon astronomy.
Detection of a larger photon flux than expected for GZK photons would imply the emission of photons at the source or new physics. New physics is involved in Top-Down models, produced as an alternative to acceleration models to explain the origin of the highest energy cosmic rays. All of the Top-Down models predict photon dominance at the highest energies. Here, we estimate the minimum photon fraction Top-Down models predict, not only assuming the AGASA spectrum which these models were originally proposed to explain, but also assuming the HiRes spectrum. We show that at high energy, close to 10<sup>20</sup> eV, the maximum expected flux of GZK photons is comparable to (for the AGASA spectrum) or much smaller than (for the HiRes spectrum) the minimum flux of photons predicted by Top-Down models which fit the AGASA or the HiRes data (see Fig. 17). We try to minimize the photon ratio predicted by Top-Down models by assuming that these models explain only the highest energy UHECR (if they do not explain even those events, the models are irrelevant for UHECR). We show that the photon ratio at energies close to $`10^{20}`$ eV is a crucial test for Top-Down models, since it is always higher than about 0.5, independently of the UHECR spectrum assumed.
We also show that, surprisingly, in a limited energy range above $`10^{20}`$ eV, GZK photons could become the dominant component of the UHECR (assuming that protons could be accelerated at the source to energies as large as 10<sup>22</sup> eV). This result allows us to fit the AGASA data with an original flux of only nucleons. This seems to contradict previous estimates of the GZK photon flux in which this flux is always subdominant, however one needs to take into account the assumed initial spectrum and intervening radio background and magnetic fields (for example in Ref. reviewGZKneutrinos an average EGMF of 10<sup>-9</sup> G is assumed, much larger than the fields found later in Ref. dolag2004 ).
In section II, we explain our calculations and show the dependence of the GZK photon flux on the assumed initial proton flux and intervening background parameters. In section II we only normalize the fluxes we show to one point of the AGASA or HiRes spectrum, but we do not fit these spectra (which we do in the following section). In section III, we estimate the maximum and minimum GZK photon fractions expected either with the AGASA spectrum or with the HiRes spectrum. In section IV we estimate the minimum photon fractions predicted by several by Top-Down models and compare them with the maximum GZK photon fraction we find in section III. We also include a comparison with experimental upper bounds on photon fractions.
## II II. The GZK photon flux
We use a numerical code developed in Ref. kks1999 to compute the flux of GZK photons produced by an homogeneous distribution of sources emitting originally only protons. It calculates the propagation of protons and photons using the standard dominant processes, explained for example in Ref. reviews1 ). For protons, it takes into account single and multiple pion production, and $`e^\pm `$ pair creation. For photons, it includes $`e^\pm `$ pair production, inverse Compton scattering and double $`e^\pm `$ pair production processes. For electrons and positrons, it takes into account Compton scattering, triple pair production and synchrotron energy loss on extra galactic magnetic fields (EGMF). The propagation of protons and photons is calculated self-consistently. Namely, secondary (and higher generation) particles arising in all reactions are propagated alongside with the primaries. UHE protons and photons lose their energy in interactions with the electro-magnetic background, which consist of CMB, radio, infra-red and optical components, as well as EGMF. Protons are sensitive essentially to the CMB only, while for photons all components of the electro-magnetic background are important. Notice that the radio background is not yet well known and that our conclusions depend strongly on the background assumed. We include three models for the radio background: the background based on estimates by Clark et al. clark and the two models of Protheroe and Biermann PB , both predicting larger background than the first. To calculate the infra-red/optical background we used the same approach as in Ref. Primack:2000xp . In any event, the infra-red/optical background is not important for the production and absorption of GZK photons at high energies. This background is important to transport the energy of secondary photons in the cascade process from the 0.1 - 100 TeV energy range to the 0.1-100 GeV energy range observed by EGRET. The resulting flux in the EGRET energy range is not sensitive to details of the infra-red/optical background models.
For the EGMF only the upper bound is established observationally, $`B<10^9(\mathrm{Mpc}/l_c)^{1/2}`$ G FR (where $`l_c`$ is the reversal scale of the magnetic field in comoving coordinates). It is believed that the magnetic fields in clusters can be generated from a primordial “seed” if the later has comoving magnitude $`B10^{12}`$ G Dolag:2002 ; dolag2004 . The evolution of EGMF together with the large scale structure of the Universe has been simulated recently by two groups using independent numerical procedures Sigl:2004yk ; dolag2004 . Magnetic field strengths significantly larger than 10<sup>-10</sup> G were found only within large clusters of galaxies. In our simulations we vary the magnetic field strength in the range $`B=10^{12}10^9`$ G, assuming an unstructured field along the propagation path.
Notice that we assume that protons are produced at the source but the results at high energies would be identical if we had taken neutrons instead. The interactions of neutrons and protons with the intervening backgrounds are identical and when a neutron decays practically all of its energy goes to the final proton (while the electron and neutrino are produced with energies 10<sup>17</sup> eV or lower).
The resulting GZK photon flux depends on several astrophysical parameters. These parametrize the initial proton flux, the distribution of sources, the radio background and the EGMF. In this section, to explore the flux dependence on a given parameter, we fix all the other unknown parameters to the following values. For the radio background we take the lower estimate of Protheroe and Biermann PB , which is intermediate between the other two we consider. For the EGMF we take $`B=10^{11}`$G which is the average value found in Ref. dolag2004 . For the source distribution, we take a uniform continuous distribution of sources with zero minimum distance to us (i.e. a minimum distance comparable to the interaction length). For the maximum energy of the injected protons we use $`E_{\mathrm{max}}=10^{22}`$ eV, which is considered already a generous upper limit for acceleration in astrophysical models hillas .
With respect to cosmological parameters, we take the Hubble constant $`H=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, a dark energy density (in units of the critical density) $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ and a dark matter density $`\mathrm{\Omega }_\mathrm{m}=0.3`$. We assume the sources extend to a maximum redshift $`z_{\mathrm{max}}=2`$ (although any $`z_{\mathrm{max}}>1`$ gives the same results at the high energies we consider) and disregard a possible evolution of the sources with redshift.
### II.1 A. Dependence of the GZK photon flux on the initial proton spectrum
We parametrize the initial proton flux for any source with the following power law function,
$$F(E)=f\frac{1}{E^\alpha }\theta (E_{\mathrm{max}}E).$$
(1)
The power law index $`\alpha `$ and maximum energy $`E_{\mathrm{max}}`$ are considered free parameters. The amplitude $`f`$ is fixed by normalizing the final proton flux from all sources to the observed flux of UHECR, which we take to be either the AGASA flux or the HiRes flux.
We are implicitly assuming that the sources are astrophysical, since these are the only ones which could produce solely protons (or neutrons) as UHECR primaries. Astrophysical acceleration mechanisms often result in $`\alpha >2`$ AS2 , however, harder spectra, $`\alpha <1.5`$ are also possible, see e.g. Ref. AS1.5 . The resulting spectrum may differ from a power-law, it may even have a peak at high energies peaks . AGN cores could accelerate protons with induced electric fields, similarly to what happens in a linear accelerator. This mechanism would produce an almost monoenergetic proton flux, with energies as high as $`10^{20}`$ eV or higher mono . Here, we will consider the power law index to be in the range $`1\alpha 2.7`$.
Fig. 1 shows the GZK photon flux for three values of the power law index in Eq. (1), $`\alpha =`$1.5, 2 and 2.7. Dotted (solid) lines correspond to the resulting flux of protons (GZK photons) from all sources. A proton spectrum $`1/E^{2.7}`$ does not require an extra contribution to fit the UHECR data, except at very low energies $`E<10^{18}`$ eV outside the range we study Berezinsky:2002vt . For $`\alpha 2`$ an extra low energy component (LEC) is required to fit the UHECR data at $`E<10^{19}`$ eV. The LEC may be a galactic contribution (for example of iron nuclei, to explain the lack of correlation of arrival directions with the galactic plane), which can be parametrized as power law with an exponential cutoff as in Eq.(2) below. In this case, the “ankle” is the energy where the extragalactic protons start to dominate over the LEC. The LEC could also be due to a population of extragalactic lower energy proton sources. This latter contribution can be parametrized again as in Eq. (1), but with parameters different than those of the extragalactic proton population which dominates above the GZK energy.
Notice that in this section we just normalize the total flux to a point of the AGASA or HiRes spectrum, but we do not fit these spectra, so we do not add a LEC, even if it would be needed. We do fit the UHECR spectrum in the next section.
As seen in Fig. 1, the flux of super-GZK protons and, consequently, the flux of the GZK photons they generate, depend strongly on the power law index of the initial proton flux: they are lower for large values of $`\alpha `$. In the most conservative case of a proton flux $`1/E^{2.7}`$ the GZK photon flux at $`E=10^{19}`$ eV is as small as $`0.03\%`$ and it increases to a few $`\%`$ at $`E=2\times 10^{20}`$ eV. This means that even with the final statistics of Auger it might be difficult to detect the GZK photons in this case. On the other hand, in the optimistic case of an injection spectrum $`1/E^{1.5}`$, the GZK photons can contribute as much as 1-3% at $`E=10^{19}`$ eV and 50% or more at $`E=10^{20}`$ eV.
Let us note here, that most of the energy produced in the form of GZK photons cascades down in energy to below the pair production threshold for photons on the CMB. For $`\alpha <2`$ the diffuse extragalactic gamma-ray flux measured by EGRET EGRET at GeV energies imposes a constraint on the GZK photon flux at high energies, which we have taken into account.
The dependence of the GZK photon flux on the maximum energy $`E_{\mathrm{max}}`$ of the initial proton flux (see Eq. (1)) is shown in Fig. 2, for $`E_{\mathrm{max}}=`$ $`10^{21}`$ eV, $`10^{22}`$ eV and $`10^{23}`$ eV. We do not show here the case of $`\alpha =2.7`$ because for such a steeply falling proton flux the GZK photon flux practically does not depend on $`E_{\mathrm{max}}`$. Fig. 2a shows the case of $`\alpha =2`$ and Fig. 2b that of $`\alpha =1.5`$. These figures clearly show that the dependence on $`E_{\mathrm{max}}`$ is more significant for smaller values of the power law index $`\alpha `$. Note that not only the photon flux, but also the final UHECR proton flux above the GZK cutoff depends strongly on $`E_{\mathrm{max}}`$.
For relatively small values of the maximal energy, such as $`E_{\mathrm{max}}=10^{21}`$ eV, the GZK photon flux is very small for any power law index $`\alpha `$ (see the lowest curves in Fig. 2a and Fig. 2b). For larger values of the maximal energy, such as $`E_{\mathrm{max}}=10^{22}`$ eV and $`E_{\mathrm{max}}=10^{23}`$ eV, the GZK photon flux increases considerably for $`\alpha 2`$.
### II.2 B. Dependence of the GZK photon flux on the minimal distance to the sources
Quite often in the literature the minimal distance to the sources is taken to be negligible (i.e. comparable to the interaction length). This is one of the cases we consider as well. However we take also 50 Mpc, as inferred from the small-scale clustering of events seen in the AGASA data AGASA\_clusters , and 100 Mpc, to show how the fluxes diminish with this assumption (what proves that most photons come from smaller distances). Contrary to AGASA, HiRes does not see a clustering component in its own data HiRes\_clusters . The combined dataset shows that clustering still exists, but it is not as significant as in the data of AGASA alone agasa\_hires . Note, that the non-observation of clustering in the HiRes stereo data does not contradict the result of AGASA, because of the small number of events in the sample agasa\_hires\_ok .
Assuming proton primaries and a small EGMF (following Ref. dolag2004 ), it is possible to infer the density of the sources sources ; agasa\_hires\_ok from the clustering component of UHECR. AGASA data alone suggest a source density of $`2\times 10^5`$ Mpc<sup>-3</sup>, which makes plausible the existence of one source within 50 Mpc of us. However, the HiRes negative result on clustering requires a larger density of sources and, as a result, a smaller distance to the nearest one of them. Larger values of the EGMF (as found in Ref. Sigl:2004yk ) and/or some fraction of iron in the UHECR have the effect of reducing the required number of sources and, consequently, increasing the expected distance to the nearest one.
Fig. 3 shows the dependence of the UHECR proton and GZK photon fluxes on the assumed minimal distance to sources for an initial proton flux $`1/E^2`$ in Fig. 3a and $`1/E^{1.5}`$ in Fig. 3b. The highest, intermediate and lowest fluxes correspond to a minimal distance of 0 (labeled cont. for continuous), 50 and 100 Mpc, respectively. Notice that in all the examples presented in Fig. 3 the protons dominate the flux (i.e. the total flux is practically the proton flux). Only the highest proton fluxes shown in Fig. 3 (with negligible minimal distance) fit well the HiRes data. The intermediate and lowest proton fluxes have a sharp cutoff and do not fit the HiRes data any longer. We clearly see in the figures that most of the GZK photons with energies $`E>10^{19}`$ eV should come from nearby sources within 100 Mpc (see the impressive reduction in flux if we only take sources more than 100 Mpc away).
### II.3 C. Dependence of the GZK photon flux on the radio background
The main source of energy loss of photons with $`E>10^{19}`$ eV is pair production on the radio background (while at lower energies pair production on the CMB is more important). Fig. 4 shows GZK photon fluxes for the three different estimates of the radio background we consider: the minimal background, of Clark et al. clark , and the two estimates of Protheroe and Biermann PB , both larger than the first one. In Fig. 4a the injected proton spectrum $`1/E^2`$ and in Fig. 4b it is $`1/E^{1.5}`$. These figures show that (for the EGMF assumed, $`B=10^{11}`$ G as mentioned above) the GZK photon flux depends only mildly on the radio background at energies below $`E<10^{20}`$ eV, where we find a factor 2-3 of difference between the highest flux (with the lowest radio background from Ref. clark ) and the lowest flux (with the highest background of Ref. PB ). However, at energies above $`E>10^{20}`$ eV, the differences increase, reaching one order of magnitude or more. This behavior is due to the different shapes of the assumed radio spectra. As we see next, larger EGMF, $`B>10^{10}`$ G, increase the GZK photon absorption considerably at $`E<10^{20}`$ eV, but not close to $`E10^{20}`$ eV and above.
### II.4 D. Dependence of the GZK photon flux on EGMF
The spacial structure, amplitude and correlation length of the EGMF outside clusters of galaxies are unknown. The existing models of the EGMF attempt to evolve these fields together with the large scale structure of the Universe, starting from certain (primordial) seed values. In these models, the EGMF in the voids are close to the comoving value of the primordial field, while the EGMF in clusters of galaxies and filaments are amplified. Constrained simulations of the “local” Universe (within 100 Mpc from Earth) dolag2004 , in which the magnetic field is normalized to the values observed within clusters, yield an average $`B_{\mathrm{EGMF}}=(10^{11}10^{12})`$ G in voids. Fig. 5 shows that for $`B_{\mathrm{EGMF}}<10^{10}`$ G, the resulting GZK photon flux changes very little with $`B`$, but it decreases considerably at low energies for $`B_{\mathrm{EGMF}}>10^9`$ G. In Fig. 5 an initial proton flux $`1/E^2`$ was assumed and sources were integrated from zero distance. Assuming a minimum distance of 50 Mpc to the nearest sources (case not shown in the figures), the GZK photon fluxes differ at most by a factor of 3 when the EGMF magnitude is varied in the range $`B<10^{10}`$ G.
Fig. 5 is the only place in this paper where we used $`B_{\mathrm{EGMF}}=10^8`$ G, and this is just to show how the photon flux is affected by large $`B`$ fields. For EGMF $`10^8`$ G or larger, the photon energy is lost into synchrotron radiation as soon as the UHE photon pair produces, even for energies $`E<10^{19}`$ eV. Thus the shape of the spectrum follows the energy dependence of the photon pair production interaction length (which is dominated by the interaction with the CMB below $`10^{19}`$ eV and with the radio background above this energy). For smaller magnetic field strengths, the length of synchrotron energy loss increases and, at low energies, several steps of pair production and inverse Compton decay happen. For large enough energies, the synchrotron radiation length is smaller than the interaction length for all the EGMF values considered (i.e. even as small as $`B10^{12}`$ G) , so the the photon energy is lost into synchrotron radiation as soon the photon pair produces. Thus, only the photons which do not interact with the radio background can reach us and the spectra for all values of the EGMF converge.
Our results in Fig. 5 for $`B_{\mathrm{EGMF}}10^9`$ G are similar to those in Fig. 3 of Ref. SiglOlinto95 . In particular, both figures show that the GZK flux does not depend strongly on the magnetic field for $`B_{\mathrm{EGMF}}<10^{10}`$ G, and that for larger fields there is a suppression of the photon flux at energies $`E<10^{19}`$ eV (due to pair production on the CMB followed by synchrotron energy loss).
### II.5 E. Summary of the GZK photon flux dependence on different parameters
Figs. 4 and 5 show that given a particular UHECR proton flux the uncertainty in the resulting GZK photon flux due to our ignorance of the intervening backgrounds (minimum to maximum estimates of the radio background and EGMF from 10<sup>-11</sup> G, which is equivalent to zero, to $`10^9`$ G) is within about one order of magnitude.
Figs. 1 to 3 show much larger changes in the GZK photon flux when the parameters defining the UHECR proton flux, i.e. the power law index $`\alpha `$, maximum energy $`E_{\mathrm{max}}`$, and minimal distance to the sources, are varied. However, once the particular UHECR spectrum is fixed, these uncertainties due to the extragalactic proton model decrease and become comparable with those due to our ignorance of the intervening background. In the next section, Figs. 8 and 9 show that a particular proton dominated observed flux, the HiRes spectrum in this case, can be fitted with very different extragalactic proton fluxes, whose corresponding GZK photon fluxes differ by about one order of magnitude, for a given fixed background. In fact, the difference between the two photon lines in Fig. 8 shows the uncertainty in the GZK photon flux due to intervening background (about one order of magnitude), given a particular extragalactic proton flux, while the difference between the lower photon line of Fig. 8 and the lower photon line of Fig. 9 (both computed with the same background, i.e. maximum radio background and EGMF $`B=10^9`$ G) shows the uncertainly due to the UHECR proton flux (which is one order of magnitude too).
This means that placing an upper limit on the GZK photon flux, or measuring it, provides complementary information to that contained in the UHECR proton flux itself. However, extracting information on the extragalactic nucleon flux from the GZK photons would require to have independent information on the extragalactic magnetic fields and radio background, vice versa.
## III III. Results: possible scenarios with GZK photons
We show in Sect. II that if the UHECR above 10<sup>19</sup> eV are mostly protons (or neutrons), depending on the slope of the proton flux, the distribution of sources and the intervening backgrounds, between $`10^5`$ and $`10^2`$ of the UHECR above $`10^{19}`$ eV are photons. Much larger photon factions are predicted at $`10^{20}`$ eV in some cases.
The largest GZK photon fractions in UHECR happen for small values of $`\alpha `$, large values of $`E_{\mathrm{max}}`$, small minimal distance to the sources (which is compatible with a small frequency of clustering of the events) and small intervening backgrounds. In the most favorable cases for a large photon flux, GZK photons could dominate the UHECR flux in an energy range above 10<sup>20</sup> eV. As we show below, this allows us to fit the AGASA data, at the expense of assuming that the initial protons could have a hard spectrum $`1/E`$ and be accelerated to energies as high as $`10^{22}`$ eV. In this extreme case, the AGASA data (as shown in subsection III-A below) can be explained without any new physics, except in what the mechanism of acceleration of the initial protons is concerned. We also fit the HiRes spectrum (in III-B below). With the HiRes spectrum the GZK photons are always subdominant and can be neglected for the fit. In both cases, AGASA or HiRes data, we evaluate the minimum and maximum GZK photon fractions expected with each spectrum of UHECR.
We make a one-parameter $`\chi ^2`$ fit to the assumed total spectrum, obtained by summing up the contributions of protons, GZK photons and a low energy component (LEC) when needed.
In this section we parametrize the LEC with
$$F_{\mathrm{LEC}}E^\beta \mathrm{exp}(E/E_{\mathrm{cut}}).$$
(2)
and we fit the amplitude to the lowest energy bin in the figures. We choose the parameter $`\beta =2.72.8`$ to fit the low energy spectral points, and the parameter $`E_{\mathrm{cut}}`$ so that the minimum $`\chi ^2`$ value per degree of freedom of the fit is smaller than one.
We use the 18 highest energy data bins of AGASA and the 16 highest energy data bins of HiRes-1 monocular data. We also separately check the $`\chi ^2`$ for the AGASA events above the GZK cutoff, i.e. for the 3 highest energy AGASA data bins, with $`E>10^{20}`$ eV. We do this to exclude models which do not fit well the highest energy events but whose minimum $`\chi ^2`$ considering all the 18 bins could be good due to the LEC assumed. Additionally, we check that the number of events predicted above the end point of the AGASA spectrum (the energy above which AGASA has observed no events), i.e. at $`E>2.5\times 10^{20}`$ eV, is not larger than 4 (predicting 4 events and observing none has a very small Poisson probability of 1.8%). The number of events we predict above the end point of the HiRes spectrum, at $`E>3.2\times 10^{20}`$ eV, is always much smaller than 4.
### III.1 A. GZK photons with the AGASA spectrum
In this subsection, we will discuss fits to the AGASA data with extragalactic protons, their secondary GZK photons and a LEC as in Eq. (2) when needed. Unless we mention otherwise, here we take a zero (i.e. comparable with the interaction length) minimum distance to the sources.
The fits to the AGASA spectrum at high energy with a proton dominated flux are very poor. As shown in Fig. 1, for $`\alpha <2.7`$ a low energy component (LEC) which we parametrize as in Eq. (2), possibly consisting of galactic or extragalactic Fe and protons, is necessary to fit the data. It is well known that with extragalactic protons plus a LEC one can fit the AGASA data below the GZK cutoff, at energies $`3\times 10^{18}\mathrm{eV}<E<10^{20}\mathrm{eV}`$. In fact, we tried power law indexes $`\alpha =2.7,2,1.5,1`$ and we obtained fits with minimum $`\chi ^2=36,17.7,14,14`$ for 15 degrees of freedom, respectively. The first fit (with $`\alpha =2.7`$, which does not require a LEC) is bad, but the others (which do require a LEC) are good. Even the first fit could be improved to a minimum $`\chi ^2=18`$ by changing the power index slightly to $`\alpha =2.6`$ and increasing the number of sources in the early universe as $`(1+z)^3`$. However, the same proton fluxes fit the AGASA data at $`E>10^{20}`$ eV very poorly. We found minimum $`\chi ^2=12,12,9.8,7.8`$ for 3 degrees of freedom, respectively. The reason for these bad fits is that for $`\alpha 2`$ the proton flux at super-GZK energies is very small, and even for $`\alpha <2`$ it is still not enough.
These fits can be improved by adding a large component of GZK photons. We try to maximize the GZK photon flux by reducing the radio background and EGMF, and increasing the maximum proton energy in Eq. (1) up to $`E_{\mathrm{max}}=10^{22}`$ eV.
In Figs. 6 and 7 we show (a) the differential spectra, of each component (i.e. extragalactic $`p`$, LEC and GZK $`\gamma `$) and total (upper panels), and (b) the integrated flux fractions of different components in percentage of the total predicted flux above the energy $`E`$ (lower panels). The extragalactic protons have here an initial spectrum $`1/E`$, with maximum energy $`E_{\mathrm{max}}=10^{22}`$ eV (see Eq. (1)). The particular LEC shown has parameters $`\beta =2.7`$ and cutoff energy $`E_{\mathrm{cut}}=10^{19}`$ eV (see Eq. (2)). In both Figs. 6 and 7 the EGMF is B=10<sup>-11</sup> G. The only difference between both figures is the radio background: we took the lowest one for Fig. 6 and the intermediate one for Fig. 7. This is the only change we can impose between the maximum and the minimum GZK photon flux while not reducing the goodness of fit to the AGASA data to unacceptable levels.
The fit to the super-GZK AGASA events in Fig. 6a is now perfect, due to the GZK photons: it has a minimum $`\chi ^2=2.6`$ for 3 degrees of freedom and at $`E>10^{20}`$ eV there are 11.5 events (6.8 photons and 4.5 protons) where AGASA has observed 11. The spectrum predicts 4 events (2 photons and 2 protons) at energies above 2.5$`\times 10^{20}`$ eV, where AGASA has seen none, which we take as acceptable (the probability is small, 1.8%). Larger $`E_{\mathrm{max}}`$ or lower $`\alpha `$ values would lead to predict even more events where AGASA has seen none and would therefore not fit well the AGASA spectrum any longer.
The fit to the super-GZK AGASA events in Fig. 7a, where we try to lower the GZK flux, is not as good as that in Fig. 6a: it has a minimum $`\chi ^2=5.5`$ for 3 degrees of freedom and at $`E>10^{20}`$ eV there are 7 events (2.5 photons and 4.5 protons). But, this fit is better than that is Fig. 6a above the end-point of the AGASA spectrum: it predict only 2.7 events above the highest energy AGASA point, which has a 6.7% Poisson probability.
As we see, a good fit to the AGASA data at $`E>10^{20}`$ eV with GZK photons is strongly restricted by the total number of events on one side and by the number of events above the end-point of the AGASA spectrum on the other. Thus, Figs. 6-7 provide an estimate of the maximum and minimum GZK photon flux which fit the AGASA data.
Notice in Fig. 6b that with the maximum GZK photon flux prediction, the photon ratio increases from about 7 % at $`10^{19}`$ eV to more than 50 % above $`10^{20}`$ eV, and that the total differential flux is dominated by GZK photons at energies between 1 and 7 $`\times 10^{20}`$ eV. This large GZK photon flux is possible only under the extreme conditions chosen here. A larger radio background, or a smaller maximum proton energy quickly diminish the GZK photon flux, as Fig. 7 demonstrates.
The EGRET bound on the photon energy which cascades down to the GeV energies has been taken into account. We found that the flux predicted is about one order of magnitude below the level measured by EGRET.
The 2-$`\sigma `$ AGASA upper bounds on the Fe fraction in the integrated fluxes, of 46% and 44% above 10<sup>19.0</sup> eV and 10<sup>19.25</sup> eV respectively agasa\_composition\_2 are shown in Fig. 6b and Fig. 7b. The LEC could respect these bounds (so that the LEC could consist entirely of galactic Fe), if we assumed a somewhat softer proton spectrum than we choose for Figs. 6 and 7, possibly with $`\alpha >1.5`$. With our choice, the extragalactic proton spectrum is a bit too low at energies below the GZK energy and, consequently, the LEC is too large. The lower HiRes limit on a possible Fe low energy component hires\_composition\_fit , rejects entirely a LEC consisting mostly of iron. In this case the LEC should consist mostly of extragalactic protons with a soft spectrum $`1/E^{2.7}`$ and a small maximum energy $`E_{\mathrm{max}}10^{20}`$ eV which should come from a different class of UHECR sources (than those which produce the super-GZK UHECR).
Also shown in Fig. 6b and Fig. 7b is the bound on the photon fraction obtained with AGASA data at 10<sup>20</sup> eV agasa\_photon , which is saturated by the photon flux in Fig. 6.
### III.2 B. GZK photons with the HiRes spectrum
To estimate the possible range of photon fractions compatible with the HiRes spectrum we will here present two fits to the HiRes data, one maximizing and one minimizing the GZK photon flux. These fits are presented in Figs. 8 and 9 respectively.
Figs. 8 show (a) the differential spectra of each component (i.e. extragalactic protons, LEC and GZK photons) and total and (b) the integrated flux fractions of different components with respect to the total predicted flux shown in Fig.8a. In order to maximize the flux of GZK photons we need a relatively hard proton spectrum, thus a LEC is needed to fit the data at energies $`E<10^{19}`$ eV. The particular LEC shown has parameters $`\beta =2.7`$ and cutoff energy $`E_{\mathrm{cut}}=2\times 10^{19}`$ eV (see Eq. (2). Here we assume an extragalactic proton spectrum $`1/E`$ with maximum energy $`E_{\mathrm{max}}=10^{21}`$ eV, to maximize the number of super-GZK protons, and to minimize the photon absorption by the intervening medium, we assume the minimum radio background and $`B_{\mathrm{EGMF}}=10^{11}`$ G. This results in the higher photon curve in the figures. The lower photon curve shows how much the photon flux decreases if we keep the same proton flux and change the intervening background from minimum to maximum, i.e. if we use $`B_{\mathrm{EGMF}}=10^9`$ G and maximum radio background. The change is about an order of magnitude.
The total flux shown in Fig. 8a is dominated by protons and is insensitive to the GZK photon contribution. With this flux only one event (a proton event) is predicted above 1$`\times 10^{20}`$ eV.
Also shown in Fig. 8b are the HiRes limits on a possible LEC Fe component hires\_composition\_fit and the bound on the photon fraction obtained with AGASA data at 10<sup>20</sup> eV agasa\_photon .
In Fig. 9 we fit the HiRes data with a conservative model with a soft extragalactic proton spectrum, which does not require a low energy component. Thus, the power law index of the required proton spectrum is fixed by the observed UHECR at energies below $`10^{19}`$ eV, where the spectrum is $`1/E^{2.7}`$. This model has practically no freedom in the choice of the proton flux power law index $`\alpha `$, although this could be slightly varied in the range $`\alpha =2.42.7`$ by changing the redshift dependence of the distribution of sources. For Fig. 9 we conservatively choose $`\alpha =2.7`$ and the smallest cutoff energy which provides a good fit, which is $`E_{\mathrm{max}}=3\times 10^{20}`$ eV. We assume zero minimal distance to the sources (larger values do no provide a good fit at high energies), and, to maximize the absorption of photons, the maximum radio background and $`B_{\mathrm{EGMF}}=10^9`$ G for the lower photon curve. We also give the result for $`B_{\mathrm{EGMF}}=10^{11}`$ G and intermediate radio background (higher photon curve) to show how the photon flux increases with a less absorbing intervening background. The total flux is insensitive to the GZK photon contribution.
The difference between the lower photon line of Fig. 8 and the lower photon line of Fig. 9 (both computed with the same background) shows the uncertainly due to the UHECR proton flux (which is one order of magnitude too) for models that fit the HiRes spectrum.
Also shown in Fig. 9b are the HiRes limits on a possible LEC Fe component hires\_composition\_fit and the bound on the photon fraction obtained with AGASA data at 10<sup>20</sup> eV agasa\_photon .
We see in Fig. 9b that in this case, in which we try to minimize the GZK photons, these could contribute only $`12\times 10^4`$ at $`10^{19}`$ eV, and $`12\times 10^5`$ at $`10^{20}`$ eV, of the total integrated flux. These levels of photon fraction are out of reach for the present generation of experiments. At best Auger would detect a few GZK photons in several years of observations, while HiRes would only obtain upper limits on the number of photons at all energies.
## IV IV. Discussion: Comparison of GZK photons, minimum Top-Down photon predictions and experimental bounds
In this section we discuss the present experimental bounds on and theoretical predictions for UHECR photons, and discuss the implications of a possible future photon detection or future experimental upper limits on the photon fraction.
We start by comparing the minimal amount of photons predicted by Top-Down models of UHECR with the expected range of GZK photons discussed in Sect.III. We show that, at high energies, close to 10<sup>20</sup> eV, the maximum expected flux of GZK photons is comparable to (for the AGASA spectrum) or much smaller than (for the HiRes spectrum) the minimum flux of photons predicted by Top-Down models which fit the AGASA or the HiRes data. Thus, detection of a larger photon flux than expected for GZK photons, at those energies, would point to a Top-Down model (or to the emission of a large flux of photons at the sources). The estimate of the minimum photon ratio predicted by Top-Down models is also essential when applying to these models already existing and possible future upper bounds on the fraction of photons in UHECR
Let us recall that Top-Down models were introduced as an alternative to acceleration models to explain the highest energy cosmic rays, which the latter models have difficulty explaining. The spectra of the UHECR produced in Top-Down models are determined by the elementary particle physics of Z-boson decays and of QCD fragmentation, which predict photon domination of the spectrum at high energies.
In order to minimize the photon fraction predicted by Top-Down models while fitting the UHECR spectrum, we ask Top-Down models to explain only the highest energy events, those close to 10<sup>20</sup> eV while invoking a more conventional Bottom-Up extragalactic component (which we assume consists of nucleons) to dominate the flux at energies just below. This is an unnatural possibility which would require two completely independent mechanisms to provide UHECR at comparable levels. We consider it only because it provides the minimum amount of Top-Down photons. We will present here fits to the AGASA and HiRes data following this strategy to minimize the predicted photons for three Top-Down models: Z-bursts, topological defects (necklaces) and super heavy dark matter particles (SHDM).
### IV.1 A. Z-bursts
In the Z-burst model zburst ultra-high energy (UHE) neutrinos coming from remote sources annihilate at the Z-resonance with relic background neutrinos. The Z bosons then decay, producing secondary protons, neutrinos and photons. The Z-resonance, which acts as a new cutoff, occurs when the energy of the incoming $`\nu `$ is $`E_{res}=M_Z^2/2m_\nu =4\times 10^{21}\mathrm{eV}(\mathrm{eV}/m_\nu )`$
So far Z-burst models have been studied mostly to explain the AGASA spectrum (however, see Ref. Fodor-K-R ). Many problems have been found, which are alleviated if one assumes the HiRes spectrum. One of them is that practically no photons should be produced at the source together with the UHECR neutrinos, otherwise too many low energy photons in the EGRET region are predicted. For example, with sources emitting equal power in neutrinos and photons, the EGRET bound EGRET on the diffused GeV- $`\gamma `$ ray background is violated by two orders of magnitude (see Fig. 3 of Ref. zburst\_problem ), when the AGASA spectrum is considered. Also bounds by the GLUE GLUE and FORTE FORTE experiments on the primary neutrino flux, as well as the non-observation of UHECR events at energies above 2.5 $`\times 10^{20}`$ eV by the AGASA Collaboration imply a lower bound $`0.3`$ eV on the relic neutrino mass reviewGZKneutrinos2 ; Fodor-K-R ; GVW . Since this mass exceeds the square root of mass-squared differences inferred from oscillation physics, the bound in fact applies to all three neutrino masses. Together with the upper bound provided by CMB anisotropy and large-scale structure observations, this bound leaves only a small interval for neutrino masses around 0.3 eV, if Z-bursts are to explain the existing UHECR AGASA spectrum. These problems are somewhat alleviated if Z-bursts are to explain the ultra-GZK events in the HiRes spectrum instead of the AGASA spectrum, as can be seen in Fig. 11.
The $`p`$ and $`\gamma `$ curves in Fig. 10 and Fig. 11 show the predictions of a Z-burst model computed as in Ref. reviewGZKneutrinos2 but with a relic neutrino mass $`m_\nu =0.4`$ eV. We assume a maximum redshift $`z_{\mathrm{max}}=3`$ for the UHE neutrino sources (which emit only neutrinos and have not evolved), maximum intervening radio background and $`B_{\mathrm{EGMF}}=10^9`$ G. In our calculation we do not consider the effect of local inhomogeneities, such as the Virgo cluster RWW-new . The assumed spectrum of UHE neutrinos is shown in the figures. Only the part of this spectrum close to the resonance energy is relevant. Here we try to minimize the photon fraction predicted by Z-bursts by incorporating a low energy component of extragalactic nucleons.
In Fig. 10, a low energy component (LEC curve) parametrized as a power law (as in Eq. (1)) with index $`\alpha =2.8`$, cutoff energy $`E_{\mathrm{max}}=10^{20}`$ eV and a minimum distance to the sources of 50 Mpc, has been added to the contribution of the Z-bursts to fit the AGASA data. The fit has minimum $`\chi ^2=15`$ for 15 bins with $`E<10^{20}`$ eV. At higher energies, $`E>10^{20}`$ eV, the fit is not good, it has a min. $`\chi ^2=6.4`$ for 3 degrees of freedom. The reason is that the predicted flux is too low at these energies. However, the fit to the spectrum above the end-point of the AGASA spectrum, $`E>2.5\times 10^{20}`$ eV, is good: only two (mostly photon) events are predicted (where none were seen).
If we try to increase the Z-burst flux by minimizing the absorption of photons by the background, the fit is worse at high energies. If we take the lowest radio background and a small EGMF $`B=10^{12}`$ G, the fit to the AGASA spectrum at $`E>10^{20}`$ eV is better, with min. $`\chi ^2=4`$ for 3 degrees of freedom. However, 5.8 events (mostly photons) are predicted above the AGASA end point, which we consider unacceptable.
As shown in Fig. 10b, the gamma ray flux at low energies saturates the EGRET data. Also, as shown Fig. 17a, the predicted photon fraction saturates the upper bound on the photon fraction obtained with AGASA data at 10<sup>20</sup> eV agasa\_photon .
In Fig. 11, a low energy component (LEC curve) parametrized as a power law (see Eq. (1)) with index $`\alpha =2.7`$, maximum energy $`E_{\mathrm{max}}=10^{21}`$ eV and zero minimum distance to the sources, has been added to the contribution of the Z-bursts to fit the HiRes data. The spectrum of this model fits perfectly that of HiRes. Only 1.8 events (1 proton and 0.8 photon) are predicted above the end point of HiRes, were none were seen.
Because the super-GZK nucleon flux is here lower than with the AGASA spectrum, the predicted gamma ray flux at low energies is well under the EGRET data (see Fig. 11b). As can be seen in Fig. 17a, the predicted photon fraction is just under the upper bound obtained with AGASA data at 10<sup>20</sup> eV agasa\_photon .
### IV.2 B. Topological defects (necklaces)
The curves $`p`$ and $`\gamma `$ in Figs. 12-13 correspond to secondary protons and photons in a particular top-down model, in which topological defects (TD), such as necklaces, produce GUT-scale mass particles, which in turn decay into quarks, leptons etc (for a review see for example Ref. td\_review ). The mass scale of the parent particles provides the maximum energy of the UHECR, $`E_{\mathrm{max}}=m_X`$, thus these scenarios avoid the difficulty in astrophysical objects of accelerating the UHECR to the highest energies observed. As in Z-burst models, TD scenarios predict, therefore, a new cutoff given by the parent particle mass at energies above $`10^{20}`$ eV. The parent particles typically decay into leptons and quarks. The quarks hadronize, and some leptons decay resulting in a large cascade of photons, neutrinos, light leptons and a smaller amount of nucleons.
TD models may also have difficulties with the EGRET flux EGRET ; EGRET\_NEW on the diffused GeV- $`\gamma `$ ray background. We have taken this possible bound into account.
The TD model of Figs. 12-13 assumes a parent particle mass $`m_X=2\times 10^{13}`$ GeV, an EGMF of $`10^{12}`$ G and the low radio background predicted by Protheroe and Biermann, which is the intermediate radio background among the three we consider in this paper. Even if we are trying to minimize the photon flux at high energies, the radio background and EGMF value are not the maximal we used in this paper. This is so because, as we show here, a smaller amount of ultra-high energy photons yields a worse fit to the AGASA data. The heavy particle injection rate is assumed to be $`m_Xt^3`$, where $`t`$ is the cosmic time.
The QCD spectrum used for Figs. 12 and 13 (shown in Fig.11 of Ref. reviewGZKneutrinos ) corresponds to the decay of the heavy particles into two quarks without supersymmetry QCD-spectrum . Originally, this decay model predicts a ratio of about 10 photons per nucleon in the decay products (as does Ref. BK\_2001 ), while in more recent models SHDM-other ; Barbot-Drees ; SHDM\_2004 this ratio is only 2 - 3. So, for Fig. 12 and Fig. 13 the ratio was brought to be equal to 3. Here we fit the LEC with the function in Eq. (2) with $`\beta =2.7`$ and an exponential energy cut with $`E_{\mathrm{cut}}=8\times 10^{19}`$ eV, in order to increase the contribution of the TD model to the AGASA flux, which is still too low at high energies. Again, at energies $`E<10^{20}`$ eV the fit is good, with minimum $`\chi ^2=14`$ for 15 degrees of freedom. However, the fit of the AGASA spectrum above the GZK energy is bad, with minimum $`\chi ^2=7.4`$ per 3 degrees of freedom. This is due to the strong reduction of the TD flux above the GZK energy (due to the GZK effect, because there are more protons than in Fig. LABEL:F11), which means that in order to have a good fit at energies below the GZK energy, the flux is too small at higher energies. Now, there are only 3.7 events at $`E>10^{20}`$ eV (of which 2.7 are photons), while AGASA observed 11 events. But, if we take the minimum radio background (not shown in figures) instead of the intermediate one we use for the figures, the fit to the AGASA occupied bins above the GZK energy is good (with minimum $`\chi ^2=2.2`$ per 3 degrees of freedom), but the number of events predicted above the end-point of the AGASA spectrum (where no events were observed) becomes 10, which is again unacceptable.
From Fig. 12 we conclude that the representative TD models we study are barely consistent with the AGASA data. They either predict a flux too low at super-GZK energies or too many events above the highest energy events observed by AGASA. For the TD curve in Fig. 17a the model of Fig. 12 was used. We see in Fig. 17a that the predicted photon ratio is somewhat above the upper bound on the photon fraction obtained with AGASA data at 10<sup>20</sup> eV agasa\_photon .
In Fig. 13, a low energy component (LEC curve), parametrized as a power law (see Eq. (1)) with index $`\alpha =2.7`$ and cutoff energy $`E_{\mathrm{max}}=10^{21}`$ eV and zero minimum distance to the sources, has been added to the contribution of the TD model to fit the HiRes data. The spectrum of this model (with a $`\gamma /p`$ ratio of 3) fits well the HiRes data. This model predicts 0.4 events above the end point of the HiRes spectrum. It is clear that the fit would be good too with a larger $`\gamma /p`$ ratio in the TD decay products, since one can redistribute the protons between the LEC and the TD contribution without a significant change in the fit (but the photon fraction at the highest energies would be somewhat larger).
As mentioned above the QCD model used so far in this subsection predicts a ratio of about 10 photons per nucleon in the decay products QCD-spectrum (although we brought it artificially to 3) while in more recent models SHDM-other ; Barbot-Drees ; SHDM\_2004 this ratio is considerably smaller. We include here also the results obtained with one of these more recent models. The heavy particle decay spectrum used in Fig. 14 corresponds to the decay of the heavy particles into quark and antiquark pairs with the “gaugino set of supersymmetric parameters” taken from Ref. Barbot-Drees . We choose this particular decay mode because it is one in which the initial number of photons per nucleon produced is one of the lowest (since we want to estimate the minimum GZK photon flux produced). This decay model predicts a ratio of about 2 or less photons per nucleon in the decay products. At low energies the fragmentation functions were suppressed following Fig. 2.11 of Ref. Barbot-thesis . For $`(E/E_{\mathrm{max}})<R_o`$ the suppression factor used is $`R^{log_{10}(R/W^2)}`$, where $`R=R_o/(E/E_{\mathrm{max}})`$ and $`W`$ is the width in decades at which the spectrum is suppressed by a factor 0.1 (for $`(E/E_{\mathrm{max}})>R_o`$ there is no suppression). From the figure just mentioned, one can find the values of the parameters $`R_o`$ and $`W`$. We used $`Ro=10^6`$ and $`W=3.5`$.
Fig. 14 shows the maximum and minimum photon fractions found using the method of Ref. Gelmini:2007jy for $`E_{\mathrm{max}}<10^{23}`$ eV. In Ref.. Gelmini:2007jy the maximum and minimum GZK photon fractions were found assuming a power law spectrum of protons is injected by astrophysical sources and fitting the AGASA and HiRes UHECR spectra for energies $`E>4\times 10^{19}`$ eV. It was also assumed that any possible low energy component is irrelevant at this energies. Notice that the LEC in Fig. 12 fulfills this latter condition but that in Fig. 13 does not. To produce Fig. 14 we use the same procedure but replace the injected spectrum by that produced in the heavy particle decay. We choose the value of the amplitude of the injected spectrum by maximizing the Poisson likelihood function using the UHECR data from 4 $`\times 10^{19}`$ eV up to the last published bin of each spectrum plus one extra bin with zero observed events at higher energies. This extra bin and the highest energy empty published bins, take into account the non-observation of events above the highest occupied energy bin in the data of each collaboration, the end-point energy of each spectrum (i.e. at $`E>2.3\times 10^{20}`$ eV for AGASA agasa\_spec and $`E>1.6\times 10^{20}`$ eV for HiRes hires\_mono\_spec ), although their aperture remains constant with increasing energy. We then compute using a Monte Carlo technique the goodness of the fit, or $`p`$-value, of the distribution. Only the models with goodness of fit p-value larger than 0.05 are considered, as in Ref. Gelmini:2007jy . The maximum and minimum GZK photon fluxes depend on the intervening radio background and EGMF $`B`$ and on the value of $`E_{\mathrm{max}}=m_X/2`$. The 2006 Auger-photon-06 and 2007 Auger-photon-07 Auger upper bounds on the photon fraction are also shown in Fig. 14. The models with the minimal photon fraction for the AGASA spectrum change with energy. For $`E<1.3\times 10^{20}`$ eV the minimum photon fraction results from choosing $`E_{\mathrm{max}}=8\times 10^{22}`$, intermediate radio background and $`B=10^9`$G, while for $`E>1.3\times 10^{20}`$ eV the model with minimum photon fraction has the same $`E_{\mathrm{max}}`$ but maximal radio background, and $`B=10^{11}`$G. The model with the minimal photon fraction for the HiRes spectrum has also $`E_{\mathrm{max}}=8\times 10^{22}`$ and maximal radio background but $`B=10^9`$G.
### IV.3 C. Super Heavy Dark Matter (SHDM)
In this scenario super heavy metastable particles are produced in early Universe, and they remain at present. They form part of the dark matter of the Universe and, in particular of the dark halo of our galaxy. These particles (with colorful names such as ‘cryptons’ or ‘wimpzillas’) may decay bkv97 ; birkel-sarkar ; kr97 or annihilate BDK into the observed UHECR. The spectra of the decay or annihilation products are essentially determined by the physics of QCD fragmentation and this implies photon domination of the flux at the highest energies.
The UHECR in these models are produced predominantly within the dark halo of our own galaxy. Thus these models predict an excess of UHECR events from the galactic center dt1998 . This anisotropy is in conflict with the data on arrival directions of the SUGAR experiment SUGAR , unless SHDM are responsible for the majority of UHECR events only at energies above $`6\times 10^{19}`$ eV ks2003 . Even in this case, annihilating SHDM models are disfavored at least at the 99% C.L. by the SUGAR data, while decaying SHDM models have a probability of $`10\%`$ to be consistent with the SUGAR data ks2003 .
As seen in Fig. 17a the model we present is barely consistent with the upper bound on the photon fraction obtained with AGASA data at 10<sup>20</sup> eV agasa\_photon .
The $`p`$ and $`\gamma `$ curves in Fig. 15 are the predictions of a supersymmetric SHDM model taken from a recent calculation in Ref. SHDM\_2004 , obtained by averaging over all possible decay channels, including decays into quarks, squarks, gluons and gluinos. These predictions we use here as an example, are similar to those of previous calculations SHDM-other (see Fig. 17 of Ref. SHDM\_2004 ). In particular, the ratio of SHDM produced photons over nucleons is about 2.
Here we reduced the mass of the parent particle to $`m_X=2\times 10^{12}`$ GeV because, with the 10<sup>14</sup> GeV mass used in Ref. SHDM\_2004 to fit the AGASA data, we find that too many events are predicted above the end point of the AGASA spectrum. To be more precise, the model of Fig. 15, with $`m_X=2\times 10^{12}`$ GeV, predicts 3.0 events above the end-point of the AGASA spectrum, i.e. at $`E>2.5\times 10^{20}`$ eV. The fit has a min. $`\chi ^2=2`$ for the 3 occupied bins at energies $`E>10^{20}`$ eV.
For $`m_X=10^{14}`$ GeV, as used in Ref. SHDM\_2004 , the SHDM model predicts instead 8.5 events above the AGASA end-point. With the HiRes spectrum, there would not be any problem in using the higher $`m_X`$, since only 0.16 events are predicted with $`m_X=2\times 10^{12}`$ GeV and 0.8 events are predicted with $`m_X=10^{14}`$ GeV above the HiRes end-point (i.e. at $`E>3.2\times 10^{20}`$ eV).
We can turn this argument around and set a bound on the SHDM mass by requiring that no more than, say, 3 events are predicted above the end-point of the AGASA spectrum. At the $`95\%`$ C.L. this limit is $`m_X<2\times 10^{21}`$ eV. This should be taken as an order of magnitude limit, because AGASA assigned an energy to the events assuming proton primaries and for photon primaries the energy of some of the highest energy events can be higher Teshima\_privat . A way to alleviate this bound, at the expense of reducing the goodness of the fit, is to reduce the contribution of the SHDM model to the total UHECR spectrum. For example, one could allow for $`m_X=10^{14}`$ GeV by reducing by force the SHDM contribution above the AGASA end-point to 3 events. In this case only 7 events would be predicted at $`E>10^{20}`$ eV, where AGASA observed 11. The fit has a min. $`\chi ^2=6.7`$ for the 3 occupied bins at energies $`E>10^{20}`$ eV. Thus, reducing the contribution of the SHDM flux to the AGASA flux to allow for larger $`m_X`$ values brings SHDM models close to just extragalactic protons with a hard spectrum $`1/E`$, (with min. $`\chi ^2=7.8`$, see subsection III.A) in terms of goodness of fit.
The nucleon and photon spectra produced by the SHDM model we use is too hard, thus an additional low energy component (LEC), which we assume consists of extragalactic nucleons, is needed to fit the data. In Fig. 15a, a LEC, parametrized as a power law (see Eq. (1)) with index $`\alpha =2.8`$, maximum energy $`E_{\mathrm{max}}=10^{20}`$ eV, and with a zero minimum distance to the sources, has been added to the contribution of the SHDM model to fit the AGASA data. In Fig. 15b, the LEC shown, added to fit the HiRes spectrum, has $`\alpha =2.7`$, $`E_{\mathrm{max}}=10^{21}`$ eV and an assumed zero minimum distance to the sources.
Note that the SHDM model studied so far, with the AGASA spectrum predicts a significant photon fraction, about 10-20 %, at energies $`E>10^{19}`$ eV (see Fig. 17a) which are too high for the recent Auger limits on the the photon component of the UHECR. We discuss this issue in the following section.
Using the statistical method of Ref. Gelmini:2007jy and the heavy particle decay spectrum used in Fig. 14 (taken from Refs. Barbot-Drees ; Barbot-thesis \- see the explanations in the last paragraph of the previous subsection) we fitted the UHECR spectrum above 4 $`\times 10^{19}`$ eV just with the spectrum resulting from the superheavy particle decay, with no absorption or redshift, and obtained the maximum and minimum photon fractions of the integrated flux shown in Fig. 16. We assumed that the LEC is negligible at energies 4 $`\times 10^{19}`$ eV and above. Notice that the LEC in Fig.15b, chosen above to fit the HiRes spectrum, violates this assumption (what leads to lower predicted photon levels, since the SHDM model dominates only at higher energies). In SHDM models the maximum and minimum photon fractions depend only on the value of $`E_{\mathrm{max}}=m_X/2`$ and for each energy $`E`$ the values of $`E_{\mathrm{max}}`$ giving the maximum of the minimum photon ratio are different. We considered the range $`1\times 10^{20}`$ eV $`<E_{\mathrm{max}}<`$$`1\times 10^{23}`$ eV, However the fitting procedure shows that only the ranges 3.5 $`\times 10^{20}`$ eV $`<E_{\mathrm{max}}<1.4\times 10^{21}`$ eV and 1.2$`\times 10^{20}`$ eV $`<E_{\mathrm{max}}<7.1\times 10^{20}`$ eV provide acceptable models.
Notice that when the spectrum of SHDM is assumed to dominate the UHECR spectrum only at the highest energies, i.e. close the 10<sup>20</sup> eV as is the case of the model in Fig.15b, the resulting minimum photon fractions are smaller (about 1% at 1$`\times 10^{19}`$ eV - see Fig.17b) while if SHDM are assumed to reproduce the UHECR spectrum already at 4 $`\times 10^{19}`$ eV and above, the minimum expected photon fractions are larger (above 10% at 1$`\times 10^{19}`$ eV-see Fig.18b).
### IV.4 D. Photon fractions
In Fig. 17 we compare the range of GZK photon fractions we obtained in section III with the minimal photon fractions predicted by the Top-Down models shown in Figs. 10 to 13 and 15 and existing experimental upper bounds. Fig. 17 shows the fraction of photons as percentage of the total predicted integrated UHECR flux above the energy $`E`$ in every model.
In Fig. 17a and b the AGASA spectrum and the HiRes spectrum are assumed, respectively. The ZB, TD and SHDM curves in Fig. 17 correspond to the Z-burst, topological defects and super heavy dark matter models in Figs. 10 to 13 and 15. The pink bands show the range of GZK photons between the maximum and minimum fluxes obtained in Sect. III. The upper and lower boundaries of the pink band in Fig. 17a are the photon curve in Fig. 6b and photon curve in Fig. 7b, respectively. The upper and lower boundaries of the pink band in Fig. 17b are the highest photon curve in Fig. 8b and the lowest photon curve of Fig. 9b, respectively. Notice how the GZK photon band depends on the assumed spectrum: the band for AGASA is above the band for HiRes, entirely separated from it.
In Fig. 18 we compare the range of GZK photon fractions derived in Ref. Gelmini:2007jy with nucleons injected by the sources, with the maximum and minimum photon fractions in topological defects (necklaces) and superheavy dark matter models shown in Figs. 14 and 16. These were obtained with the same method of Ref. Gelmini:2007jy and the heavy particle decay model described in the last paragraphs of the subsections IV.B and IV.C.
From Fig. 17 and Fig. 18 we conclude that at energies above $`3\times 10^{19}`$ eV the minimum photon fraction predicted by Top-Down models is either larger or at most comparable to the maximum expected GZK photon ratio and the 2007 Auger Auger-photon-07 and the Agasa-Yakutsk AgasaYakutskLimit upper bounds on the photon fraction strongly constrain Top-Down models, in particular SHDM models.
The differences between Figs. 17 and 18 are due to the different methods and models with which the photons fractions were derived. The GZK photon fractions for the AGASA spectrum are lower in Fig. 18 than in Fig. 17 because of the different fitting procedure and the different choice of $`E_{\mathrm{max}}`$ which can be only as high as 10<sup>21</sup> eV in Ref. Gelmini:2007jy , a more conservative value, instead of 10<sup>22</sup> eV, the preferred value for the AGASA spectrum in Section III.
The SHDM photon fractions are much higher in Fig. 18 than in Fig.17. The superheavy particle fragmentation functions used to produce both figures are similar and the differences in the minimum photon fraction expected are due to the range of energies at which the SHDM is assumed to provide the bulk of UHECR: in Fig. 18 it is above 4 $`\times 10^{19}`$ eV and in Fig. 17 it is instead starting at energies closer to $`10^{20}`$ eV. However, in both cases the SHDM models studied either saturate or exceed the 2007 Auger bounds, in particular that at $`1\times 10^{19}`$ eV, and the Agasa-Yakutsk bound at $`1\times 10^{20}`$ eV. Thus, the Auger bounds by themselves already exclude as the dominant mechanism to produce UHECR the SHDM models considered here except at energies very close to $`10^{20}`$ eVSemikoz:2007wj . Also the photon fractions given in Fig.2 of Ref. Aloisio:2006yi are rejected by the 2007 Auger bound at $`1\times 10^{19}`$ eV. There is another type of SHDM models Ellis:2005jc in which the photon fraction can be smaller. Those with the smallest photon fractions among tend to correspond to superheavy particles with larger mass and the constraint on the events predicted above experimental end point is important. Some of these models are still allowed but very close to the existing photon limits, within a factor of two or so Ellis-private .
The topological defects models used in Figs. 17 and 18 are different, that of Fig. 18 being in line with the more recent estimates of fragmentation functions in which the photon fraction is smaller than in older models. This is the main reason for the minimal photon ratios expected in these models to be smaller in Fig 18 than in Fig 17. These models are not ruled out by present photon fraction bounds however the photon fractions they predict are above 10% at 1$`\times 10^{20}`$ eV. The present Agasa-Yakutsk limit upper limit of $`N_\gamma /N_{\mathrm{tot}}<36`$% strongly limits these models. So, either UHECR photons at energies close to 10<sup>20</sup> eV will be detected, or better experimental limits will be obtained in the future by Auger. An upper limit close to 10% at those energies, would reject all Top-Down models as the origin of UHECR.
Thus, the photon fraction at energies above 10<sup>19</sup> eV, is a crucial test for Top-Down models. The only caveat to this conclusion resides in considering that the evaluation PB of the extragalactic radio background could be wrong by several orders of magnitude, so that this background could be larger than those of Ref. PB by a large factor of 30 to 100 as suggested in Ref. Subir , although there are no specific arguments at present to justify these large factors.
We have shown in this paper that either the detection of UHECR photons or an improvement of the existing upper limits on the photon flux, is very important, both for Top-Down as well as for Bottom-Up mechanisms to explain the UHECR. SHDM and Z-burst models seem to be strongly disfavored by the present experimental upper bounds on photon fraction. With astrophysical sources, the GZK photon flux is important to understand the initial proton or neutron spectrum emitted at the UHECR sources and the distribution of sources. UHECR photons may help us to understand the intervening extragalactic magnetic fields and radio background. We have presented fits to both the AGASA and the HiRes UHECR spectra with extragalactic nucleons, the GZK photons they produce and, when needed, an additional low energy component at energies below 10<sup>19</sup> eV (see section III). The band of expected GZK photon flux depends clearly on the UHECR spectrum and also on the assumptions and procedure used (see Figs. 17 and 18). Once the particular UHECR spectrum is fixed, the uncertainties in this flux due to the extragalactic nucleon model and due to our ignorance of the intervening background are comparable (see subsection II.E). Thus, extracting information on the extragalactic nucleon flux from the GZK photons would require to have independent information on the extragalactic magnetic fields and radio background, and vice versa.
The detection of UHECR photons would open a new window for ultra-high energy astronomy and help establish the UHECR sources.
Acknowledgments
We thank I. Tkachev for fruitful discussions and suggestions at early stages of this work. We also thank S. Troitsky for careful reading of the manuscript and for several important suggestions and corrections. This work was supported in part by NASA grants NAG5-13399 and ATP03-0000-0057. G.G was supported in part by the US DOE grant DE-FG03-91ER40662 Task C.
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# Pseudoclassical Mechanics for the spin 0 and 1 particles.
## 1 Introduction
Investigations of particle systems with arbitrary spin was initially given by Bargmann-Wigner and Rarita-Schwinger , here the Dirac representations of the spin one half particles are the basis to the construction of higher spin theories. The formalism is based on the bispinor wave function with $`2s`$ Dirac indices (for spin $`s`$) and the total symmetrical representation is used to study the maximum spin value of the model.
On the other hand, the first ideas about the studies of classical systems that include in the phase space both commuting and anticommuting variables (pseudoclassical mechanics) was put forward by Schwinger in 1953. However it was Martin who achieved these ideas in 1959. Later in the Berezin and Marinov works a model for the description of spin one half particles was proposed, here the consistent formulation of the relativistic particle dynamics implies in the addition of a new constraint, this is because the formulation of the massive case has five grassmann variables. At the same time these models were also studied by Casalbuoni who explored the internal group symmetry and the gauge invariance of the resulting action. In this way was possible the description of spinless and spin one particles using these internal symmetries. Interaction of spinning particle systems with external Yang-Mills and gravitational fields was investigated in . The quantization of similar models are performed by means of the Dirac procedure for constrained systems.
Many other papers appeared about the study of spinning particles in the framework of pseudoclassical mechanics, for example the derivation of the equation of motions for the massive and massless spinning particles are treated in the works , where the spin description is achieved by means of the inclusion of internal group symmetries. Similarly, the case of the Dirac particle is discussed in the works . A path integral representation for obtaining a Dirac propagator was also obtained in and other studies connecting the pseudoclassical mechanics with the string theory was investigated for the free case as in interacting with an external field. Also, the pseudoclassical description of massless Weyl fermions and its path integral quantization when coupled to Yang-Mills and gravitational fields was studied in . Similarly, the path integral quantization of spinning particles interacting with external electromagnetic field was analyzed in .
Besides this, the pseudoclassical approach can be applied to other different models. This is the case of the Duffin-Kemmer-Petiau (DKP) theory which describes massive spin $`0`$ and spin $`1`$ particles in a unified representation. Questions about the equivalence of the DKP theory with theories like Klein-Gordon and Maxwell are discussed in (a good historical review of the DKP theory can be found in ). The Field theory of the massless DKP has a local gauge symmetry which describes the electromagnetic field in its spin 1 sector. It is important to notice that the massless case can not be obtained through the limit $`m0`$ of the massive case. This is due to the fact that the projections of DKP field into spin 0 and 1 sectors involve the mass as a multiplicative factor so that taking the limit $`m0`$ makes the results previously obtained useless. Moreover, if we simply make mass equal to zero in the usual massive DKP Lagrangian we obtain a Lagrangian with no local gauge symmetry. Studies in the Riemann-Cartan space time was proposed in .
Recently, a super generalization of the DKP algebra was done by Okubo where the starting point is the study of all irreducible representations by means of the Lie algebra $`so(1,4)`$ , moreover, a paraDKP (PDKP) algebra is constructed intimately related to the Lie superalgebra $`osp(1,4),`$ obtaining as result the super DKP algebra that contains the boson and fermion representations.
An extended variant including Grassmann variables for the DKP theory is very interesting for many reasons, for example a pseudoclassical version allow us to make an attempt to the construction of a supersymmetry variant of the theory where the action must be expressed in terms of (super)fields. It is also no clear about the particle states that will compose the (super)multiplet in this theory.
In this work we propose a possible action for the massless DKP theory in the pseudoclassical approach. In section 2, the pseudoclassical action is given including the correct boundary terms that yields a consistent equations of motions. We carry out the constraint analysis of the system and verify his invariance under $`\tau `$-reparametrizations, internal group $`O(N)`$ and SUSY transformations. We find the generators of corresponding transformations and give the Pauli-Lubanski vector. In section 3, the quantization is performed and proved that for the special case $`N=2`$ the both sectors of spin 0 and spin 1 of the DKP theory appear. We get the scalar and vectorial field as a first result, we also obtain the topological field solutions correspondent to the both spin sectors. In Section 4, using the SUSY principles we extend the proposed action to the Superspace formalism obtaining a consistent result as in the pseudoclassical model. Finally in section 5, we give our conclusions and comments.
## 2 Pseudoclassical Mechanics
We start with the action in the first order formalism that considers an internal group symmetry
$$S=\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \left[\left(\stackrel{.}{x}i\chi \psi \right)p+\frac{e}{2}p^2+\frac{i}{2}\psi \dot{\psi }+\frac{i}{2}f\psi \psi \right]+\frac{i}{2}\psi \left(\tau _2\right)\psi \left(\tau _1\right)$$
(1)
here $`x_\mu `$ is the space time coordinate, $`p_\mu `$ the auxiliary momentum vector; $`\psi _\mu ^k\left(\tau \right)k,l,\mathrm{}=1,2,\mathrm{}N`$ are the fermion coordinates, superpartner of $`x_\mu \left(\tau \right)`$, $`(x_\mu ,\psi _\mu ^k)`$ is the multiplet of matter; $`e\left(\tau \right)`$ is the einbein, his superpartner $`\chi _k\left(\tau \right)`$ is the unidimensional gravitino; $`f_{ik}\left(\tau \right)=f_{ki}\left(\tau \right)`$ is the gauge field for internal symmetry, $`(e,\chi _k,f_{ik})`$ is the supergravitational multiplet on the world line.
The action (1) includes the correct boundary terms that guarantee the consistence of the equations of motions for the grassmann variables. This is because in the variational principle the fermionic canonical coordinates have only one condition
$$\delta \left(\psi \left(\tau _2\right)+\psi \left(\tau _1\right)\right)=0$$
(2)
for the other coordinates only the space time coordinate is restricted to the condition
$$\delta x\left(\tau _2\right)=\delta x\left(\tau _1\right)=0$$
(3)
internal group indices in the case $`N=2`$ when $`i,k=1,2`$ are contracted by means of symbol Kroeneker $`\delta _{ik}`$ (for the group $`O\left(2\right)`$ and spin $`1`$) or Levi-Civita symbol $`ϵ_{ik}`$ (for the group $`Sp\left(1\right)`$ and spin $`0`$).
The lagrangian that follows from (1) is
$$=\left(\dot{x}i\chi \psi \right)p+\frac{e}{2}p^2+\frac{i}{2}\psi \dot{\psi }+\frac{i}{2}f\psi \psi $$
(4)
It is possible to write the action (1) in a different way, for this we perform the variation of $`S`$ with respect to $`p`$, then we get the following equation
$$p=e^1\left(\dot{x}i\chi \psi \right)$$
(5)
inserting this solution into (1) we obtain the second order formalism of the action
$`S`$ $`=`$ $`{\displaystyle \underset{\tau _1}{\overset{\tau _2}{}}}𝑑\tau \left[{\displaystyle \frac{e^1}{2}}\left(\dot{x}^22i\dot{x}\chi \psi \left(\chi \psi \right)^2\right)+{\displaystyle \frac{i}{2}}\psi \dot{\psi }+{\displaystyle \frac{i}{2}}f\psi \psi \right]`$ (6)
$`+{\displaystyle \frac{i}{2}}\psi \left(\tau _2\right)\psi \left(\tau _1\right)`$
then the lagrangian that follows from (6) is
$$=\frac{e^1}{2}\left(\dot{x}^22i\dot{x}\chi \psi \left(\chi \psi \right)^2\right)+\frac{i}{2}\psi \dot{\psi }+\frac{i}{2}f\psi \psi $$
(7)
the term $`\left(\chi \psi \right)^2=\chi _i\psi _i\chi _k\psi _k`$ appears because an internal group symmetry $`O(N)`$ was introduced in the theory.
Both formulations (1) and (6) are equivalent and as we will see later the constraint analysis gives the same result.
Equations of motions that follow from the action (1) result in
$$p_\mu \psi _k^\mu =0,\psi _{\mu i}\psi _k^\mu =0,\dot{\psi }_k^\mu =p^\mu \chi _k+f_{ik}\psi _i^\mu ,\dot{p}=0$$
we can see that for a special case $`e=1,`$ $`\chi =f=0`$ we obtain the solutions
$$x_\mu \left(\tau \right)=x_\mu \left(0\right)+p_\mu \tau ,\psi _k^\mu =const.$$
### 2.1 Constraint Analysis
Now we proceed to the constraint analysis of the theory. Using the definition for the canonical momentum: $`p_a=\frac{\stackrel{}{}}{\dot{q}^a},`$ we obtain
$`p_\mu `$ $`=`$ $`{\displaystyle \frac{}{\dot{x}^\mu }}=p_\mu ;\pi _\mu ^k={\displaystyle \frac{}{\stackrel{.}{\psi }_k^\mu }}={\displaystyle \frac{i}{2}}\psi _\mu ^k`$ (8)
$`\pi `$ $`=`$ $`{\displaystyle \frac{}{\dot{e}^\mu }}=0;\pi ^k={\displaystyle \frac{}{\dot{\chi }_k}}=0;\pi ^{ik}={\displaystyle \frac{}{\dot{f}_{ik}}}=0`$
from which a set of primary constraints appears
$`\mathrm{\Omega }_\mu ^k`$ $`=`$ $`\pi _\mu ^k{\displaystyle \frac{i}{2}}\psi _\mu ^k0,\mathrm{\Omega }_\pi =\pi 0,\mathrm{\Omega }^k=\pi ^k0,\mathrm{\Omega }^{ik}=\pi ^{ik}0`$
following the standard Dirac procedure for a theory with constraints we construct the primary hamiltonian from the lagrangian (4), $`=p_a\dot{q}^a`$,
$$^{(1)}=i\chi _k\psi _k^\mu p_\mu \frac{e}{2}p^2\frac{i}{2}f_{ik}\psi _{\mu i}\psi _k^\mu +\lambda ^a\mathrm{\Omega }_a$$
(10)
where we have included the primary constraints (LABEL:p6), $`\lambda ^a=\{\lambda _\mu ^k,\lambda _\pi ,\lambda ^k,\lambda ^{ik}\}`$ are the lagrange multipliers. When we apply the stability conditions on the primary constraints
$$\dot{\mathrm{\Omega }}_a=\{\mathrm{\Omega }_a,^{(1)}\}_{PB}=0$$
(11)
we obtain a new set of secondary constraints
$$\mathrm{\Omega }_\pi ^{(2)}=\frac{1}{2}p^20,\mathrm{\Omega }_k^{(2)}=i\psi _k^\mu p_\mu 0,\mathrm{\Omega }_{ik}^{(2)}=i\psi _{\mu i}\psi _k^\mu 0$$
(12)
the conservation of these secondary constraints in time tell us that no more constraints appear in the theory. Next the constraint classification gives the following first class
$`\mathrm{\Omega }_\pi ^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}p^20`$ (13)
$`\mathrm{\Omega }_k^{(2)}`$ $`=`$ $`i\psi _k^\mu p_\mu 0`$ (14)
$`\mathrm{\Omega }_{ik}^{(2)}`$ $`=`$ $`i\psi _i^\mu \psi _k^\mu 0`$ (15)
and the second class constraints
$$\mathrm{\Omega }_\mu ^k=\pi _\mu ^k\frac{i}{2}\psi _\mu ^k0$$
(16)
with the help of the second class constraints we construct the Dirac Bracket (DB) between the canonical variables and obtain
$$\{\psi _\mu ^i,\psi _\nu ^k\}_{DB}=i\delta ^{ik}g_{\mu \nu },\{x_\mu ,p_\nu \}_{DB}=g_{\mu \nu }$$
(17)
### 2.2 Invariance
In the theory with the action (6), we have three gauge transformations that do not change their physical sense. The $`\tau `$-reparametrization
$`\delta x`$ $`=`$ $`\epsilon \dot{x},\delta \psi =\epsilon \dot{\psi }`$ (18)
$`\delta e`$ $`=`$ $`\left(\epsilon e\right)^.,\delta \chi =\left(\epsilon \chi \right)^.,\delta f=\left(\epsilon f\right)^.`$
the invariance under local internal symmetries $`O\left(N\right)`$
$`\delta x`$ $`=`$ $`0,\delta \psi =a\psi `$ (19)
$`\delta e`$ $`=`$ $`0,\delta \chi =a\chi ,\delta f=\dot{a}+affa`$
and the invariance under local $`\left(n=1\right)`$ SUSY transformations
$`\delta x`$ $`=`$ $`i\alpha \psi ,\delta \psi =e^1\alpha \left(\dot{x}i\chi \psi \right)`$ (20)
$`\delta e`$ $`=`$ $`2i\alpha \chi ,\delta \chi =\dot{\alpha }f\alpha ,\delta f=0`$
It is interesting to commute two local $`\left(n=1\right)`$ SUSY transformations. This gives
$`[\delta _\alpha ,\delta _\beta ]x`$ $`=`$ $`\delta _{\epsilon _0}\dot{x}+\delta _{a_0}x+\delta _{\alpha _0}x`$ (21)
$`[\delta _\alpha ,\delta _\beta ]\psi `$ $`=`$ $`\delta _{\epsilon _0}\dot{\psi }+\delta _{a_0}\psi +\delta _{\alpha _0}\psi `$ (22)
$`[\delta _\alpha ,\delta _\beta ]e`$ $`=`$ $`\delta _{\epsilon _0}\dot{e}+\delta _{a_0}e+\delta _{\alpha _0}e`$ (23)
$`[\delta _\alpha ,\delta _\beta ]\chi `$ $`=`$ $`\delta _{\epsilon _0}\dot{\chi }+\delta _{a_0}\chi +\delta _{\alpha _0}\chi `$ (24)
$`[\delta _\alpha ,\delta _\beta ]f`$ $`=`$ $`\delta _{\epsilon _0}\dot{f}+\delta _{a_0}f+\delta _{\alpha _0}f`$ (25)
where the new parameters are now field dependent
$$\epsilon _0=2ie^1\alpha \beta ,\alpha _0=\epsilon _0\chi ,a_0=\epsilon _0f$$
(26)
this shows that there is no simple gauge group structure, although the invariance is still enough to secure good physical properties of the action.
The invariance of the action (6) is reached if we impose the conditions at the endpoints for the parameters
$$\epsilon \left(\tau _1\right)=\epsilon \left(\tau _2\right)=0,\alpha \left(\tau _1\right)=\alpha \left(\tau _2\right)=0$$
(27)
On the other hand it is possible to find the generators of the transformations (18)-(20). We follow the work of Casalbuoni where the generators of the transformations $`F`$ are given by
$$F=p_a\delta q^a\phi ,\delta L=\frac{d\phi }{d\tau }$$
(28)
being $`\phi `$ the generating function. To verify the correctness of found generators we use
$$\delta u=\{u,ϵF\}_{DB}$$
(29)
where $`ϵ`$ is the parameter of a given transformation.
We find for the $`\tau `$-reparametrizations
$`F`$ $`=`$ $`i\chi \psi p{\displaystyle \frac{e}{2}}p^2{\displaystyle \frac{i}{2}}f\psi \psi `$ (30)
$`\{x^\mu ,\epsilon F\}_{DB}`$ $`=`$ $`\epsilon \dot{x}^\mu ,\{\psi _k^\mu ,\epsilon F\}_{DB}=\epsilon \dot{\psi }_k^\mu `$ (31)
internal $`O\left(N\right)`$ symmetries
$`F_{ik}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\psi _i^\mu \psi _{\mu k}+\chi _i\pi _k`$ (32)
$`\{x^\mu ,aF\}_{DB}`$ $`=`$ $`0,\{\psi _i^\mu ,aF\}_{DB}=a_{ik}\psi _k^\mu ,\{\chi _i,aF\}_{DB}=a_{ik}\chi _k`$ (33)
and SUSY transformations
$`F_k`$ $`=`$ $`ip_\mu \psi _k^\mu +2i\chi _k\pi `$ (34)
$`\{x^\mu ,\alpha F\}_{DB}`$ $`=`$ $`i\alpha _k\psi _k^\mu ,\{\psi ,\alpha F\}_{DB}=e^1\alpha \left(\dot{x}i\chi \psi \right)`$ (35)
$`\{e,\alpha F\}_{DB}`$ $`=`$ $`2i\alpha \chi `$ (36)
To close the invariance we remark that the proposed theory is also invariant under Poincaré transformations, i.e.
$$\delta x^\mu =\omega ^\mu {}_{\nu }{}^{}x_{}^{\nu }+ϵ^\mu ,\delta \psi _k^\mu =\omega ^\mu {}_{\nu }{}^{}\psi _{k}^{\nu },\delta e=\delta \chi =\delta f=0$$
(37)
with the generators
$$ϵ^aF_a=ϵ^\mu P_\mu +\frac{1}{2}\omega ^{\mu \nu }M_{\mu \nu }$$
(38)
where
$$M_{\mu \nu }=L_{\mu \nu }+S_{\mu \nu },L_{\mu \nu }=x_\nu p_\mu x_\mu p_\nu ,S_{\mu \nu }=i\psi _\mu ^k\psi _\nu ^k$$
(39)
in this way is constructed the Pauli-Lubanski vector
$$W_\mu =\frac{1}{2}ϵ_{\mu \nu \lambda \rho }P^\nu M^{\lambda \rho },W^2=\frac{1}{2}\left(P^2S^2+2\left(S_{\mu \nu }P^\nu \right)^2\right)$$
(40)
## 3 Quantization
The constraint analysis which was done before takes a physical sense when the quantization is performed and a coherent interpretation of the equation of motions is given.
With the quantization the canonical variables becomes operators
$$x_\mu \widehat{x}_\mu ,p_\mu \widehat{p}_\mu ,\psi _\mu ^i\widehat{\psi }_\mu ^i$$
(41)
and the DB follows the commutator or anticommutator rules
$$\left\{\widehat{}\right\}i\mathrm{}\left\{\right\}_{DB}$$
(42)
thus we have the following commutation relations
$$\{\widehat{\psi }_\mu ^i,\widehat{\psi }_\nu ^k\}=\mathrm{}\delta ^{ik}g_{\mu \nu },[\widehat{x}_\mu ,\widehat{p}_\mu ]=i\mathrm{}g_{\mu \nu }$$
(43)
We pick out a general realization for the operator $`\widehat{\psi }_\mu ^k`$ satisfying the relation (43) and the equations of motions
$$D\left(\widehat{\psi }_\mu ^k\right)=S\left(Y\right)\left(\left(\gamma _5\right)^{\left(k1\right)}\gamma _\mu \gamma _5I^{\left(Nk\right)}\right)$$
(44)
here $`S(Y)`$ is the Young symmetrization operator, $`\gamma _\mu `$ are the Dirac matrices and $`\gamma _5`$ is given by
$$\gamma _5=i\gamma _0\gamma _1\gamma _2\gamma _3,\left(\gamma _5\right)^2=1$$
(45)
The first class constraints are applied into the vector state $`|\mathrm{\Phi }|\mathrm{\Phi }_{\alpha _1\mathrm{}\alpha _N}`$. We recall that an internal group symmetry $`O(N),`$ where $`i,k,\mathrm{}=1,2,\mathrm{},N`$, is considered in the Lagrangian (1). Thus we obtain
$`p^2|\mathrm{\Phi }_{\alpha _1\mathrm{}\alpha _N}`$ $`=`$ $`0`$ (46)
$`p^\mu \gamma _\mu ^k|\mathrm{\Phi }_{\alpha _1\mathrm{},\alpha _k,\mathrm{}\alpha _N}`$ $`=`$ $`0`$ (47)
$`\gamma ^{\mu i}\gamma _\mu ^k|\mathrm{\Phi }_{\alpha _1\mathrm{},\alpha _i,\mathrm{},\alpha _k,\mathrm{}\alpha _N}`$ $`=`$ $`0`$ (48)
the first equation is the mass shell condition in the case of a massless particle. The second one is a set of linear equations for every Dirac indices where no symmetrization on the vector state $`|\mathrm{\Phi }`$ is assumed. However when the symmetrization over the vector state is taken into account, (47) becomes the Bargmann-Wigner equation for a particle with spin $`N/2`$. The total symmetrical part of $`|\mathrm{\Phi }`$ generates a representation with the higher spin value. In our case, the third equation is a projector of the representations of DKP theory, i.e., it separates out a particular spin representation of the vector state.
In the particular choose: $`i,k=1,2,`$ i.e. when the internal group symmetry is $`O(2)`$, (46)-(48) reproduce de DKP equations for massless particles with spin $`0`$ and $`1`$. In this case the realization (44) becomes
$$D\left(\widehat{\psi }_\mu ^1\right)=i\sqrt{\frac{\mathrm{}}{2}}\left(\gamma _\mu \gamma _51\right),D\left(\widehat{\psi }_\mu ^2\right)=i\sqrt{\frac{\mathrm{}}{2}}\left(\gamma _5\gamma _\mu \gamma _5\right)$$
(49)
Let’s take only two Dirac indices in the vector state $`|\mathrm{\Phi }_{\alpha _1\alpha _2}`$, then using a complete set of Dirac matrices we decompose $`|\mathrm{\Phi }_{\alpha _1\alpha _2}`$ as follows
$`|\mathrm{\Phi }_{\alpha _1\alpha _2}`$ $`=`$ $`a\left(\gamma ^5C\right)_{\alpha _1\alpha _2}\zeta _5+a_1\left(\gamma ^5\gamma ^\mu C\right)_{\alpha _1\alpha _2}\zeta _{5\mu }+a_2C_{\alpha _1\alpha _2}\zeta `$ (50)
$`+b\left(\gamma ^\mu C\right)_{\alpha _1\alpha _2}\left(\zeta _\mu \right)+b_1\left(\mathrm{\Sigma }^{\mu \nu }C\right)_{\alpha _1\alpha _2}\zeta _{\mu \nu }`$
here $`a,a_1,b,b_1`$ and $`a_2`$ must be considered as free parameters and will be adjusted to assure the correctness of the final equations. The term: $`C_{\alpha _1\alpha _2}\zeta `$, is referred to a trivial representation and we do not consider it, therefore, we set $`a_2=0.`$
We also have
$$\mathrm{\Sigma }^{\mu \nu }=\frac{1}{2}\left(\gamma ^\mu \gamma ^\nu \gamma ^\nu \gamma ^\mu \right)$$
(51)
and $`C`$ is the charge conjugation matrix
$$C^T=C.$$
(52)
Considering the properties of the matrix $`C`$ we obtain the antisymmetrical
$$|\mathrm{\Phi }_{\left[\alpha _1\alpha _2\right]}=a\left(\gamma ^5C\right)_{\alpha _1\alpha _2}\zeta _5+a_1\left(\gamma ^5\gamma ^\mu C\right)_{\alpha _1\alpha _2}\zeta _{5\mu }$$
(53)
and the symmetrical part of the vector state.
$$|\mathrm{\Phi }_{\left\{\alpha _1\alpha _2\right\}}=b\left(\gamma ^\mu C\right)_{\alpha _1\alpha _2}\zeta _\mu +b_1\left(\mathrm{\Sigma }^{\mu \nu }C\right)_{\alpha _1\alpha _2}\zeta _{\mu \nu }.$$
(54)
Thus for the particular case of $`O(2)`$ symmetry we obtain
$`p^2|\mathrm{\Phi }_{\alpha _1\alpha _2}`$ $`=`$ $`0`$ (55)
$`p^\mu \gamma _\mu ^{(1)}|\mathrm{\Phi }_{\alpha _1\alpha _2}`$ $`=`$ $`0,p^\mu \gamma _\mu ^{(2)}|\mathrm{\Phi }_{\alpha _1\alpha _2}=0`$ (56)
$`\gamma _\mu ^{(1)}\gamma ^{\mu (2)}|\mathrm{\Phi }_{\alpha _1\alpha _2}`$ $`=`$ $`0`$ (57)
these relations give the DKP equation for spin $`0`$ and spin $`1`$. The relation (57) can be shown to be a projector that separates the corresponding sector of the vector state $`|\mathrm{\Phi }_{\alpha _1\alpha _2}.`$
### 3.1 Spin 0
Let’s take the antisymmetrical part of the vector state $`|\mathrm{\Phi }_{\alpha _1\alpha _2}`$ and replace it in one of the equations (56), then we obtain
$$\left(p_\mu \gamma ^\mu \right)_{\alpha \alpha _1}|\mathrm{\Phi }_{\left[\alpha _1\alpha _2\right]}=\left(p_\mu \gamma ^\mu \right)_{\alpha \alpha _1}\left[a\left(\gamma ^5C\right)_{\alpha _1\alpha _2}\zeta _5+a_1\left(\gamma ^5\gamma ^\nu C\right)_{\alpha _1\alpha _2}\zeta _{5\nu }\right]=0$$
(58)
multiplying on the right side by $`\left(C^1\gamma ^5\right)_{\alpha _2\alpha }`$ and considering $`\gamma _5^2=1`$, we have
$$p_\mu \left[a\left(\gamma ^\mu \right)_{\alpha \alpha }\zeta _5a_1\left(\gamma ^\mu \gamma ^\nu \right)_{\alpha \alpha }\zeta _{5\nu }\right]=0$$
(59)
with the use of the trace properties the equation (59) results in
$$a_1\left(p^\mu \zeta _{5\mu }\right)=0$$
(60)
On the other hand, if we multiply the equation (58) by $`\left(C^1\gamma ^5\gamma ^\lambda \gamma ^\rho \right)_{\alpha _2\alpha }`$ and taking the trace operation we got to
$$a_1\left(p^\mu \zeta _5^\nu p^\nu \zeta _5^\mu \right)=0$$
(61)
for $`a_10`$, one solution for the last relation is given by
$$\zeta _5^\mu =p^\mu \zeta _5$$
(62)
Thus equations (60) and (61) are the equations for the spin $`0`$ particles and the equation (60) gives the massless Klein-Gordon equation for the scalar field $`\zeta _5`$.
Now if we multiply (58) on the right side by $`\left(C^1\gamma ^5\gamma ^\lambda \right)_{\alpha _2\alpha }`$ we obtain
$$p_\mu \left[a\left(\gamma ^\lambda \gamma ^\mu \right)_{\alpha \alpha }\zeta _5a_1\left(\gamma ^\lambda \gamma ^\mu \gamma ^\nu \right)_{\alpha \alpha }\zeta _{5\nu }\right]=0$$
(63)
using again the trace properties for the Dirac matrices a third relation is obtained
$$a\left(p^\mu \zeta _5\right)=0$$
(64)
this equation is compatible with the equation (60) and (61) if only if $`a=0.`$
### 3.2 Spin 1
Now we take the symmetrical part of the vector state $`|\mathrm{\Phi }_{\alpha _1\alpha _2}`$, the equation (56) becomes
$$\left(p_\mu \gamma ^\mu \right)_{\alpha \alpha _1}\left[b\left(\gamma ^\nu C\right)_{\alpha _1\alpha _2}\zeta _\nu +b_1\left(\mathrm{\Sigma }^{\nu \lambda }C\right)_{\alpha _1\alpha _2}\zeta _{\nu \lambda }\right]=0.$$
(65)
Multiplying on the right side by $`\left(C^1\gamma ^\rho \right)_{\alpha _2\alpha }`$ we get
$$p_\mu \left[\left(\gamma ^\mu \gamma ^\nu \gamma ^\rho \right)_{\alpha \alpha }\zeta _\nu +\left(\gamma ^\mu \mathrm{\Sigma }^{\nu \lambda }\gamma ^\rho \right)_{\alpha \alpha }\zeta _{\nu \lambda }\right]=0$$
(66)
using the trace properties for the $`\gamma ^\mu `$-matrices it simplifies to give
$$b_1\left(p^\lambda \zeta _{\lambda \rho }\right)=0$$
(67)
Multiplying (65) by $`\left(C^1\gamma ^\rho \gamma ^\sigma \gamma ^\tau \right)_{\alpha _2\alpha }`$ it simplifies to be
$$p_\mu \left[\left(\gamma ^\mu \gamma ^\nu \gamma ^\rho \gamma ^\sigma \gamma ^\tau \right)_{\alpha \alpha }\zeta _\nu +\left(\gamma ^\mu \mathrm{\Sigma }^{\nu \lambda }\gamma ^\rho \gamma ^\sigma \gamma ^\tau \right)_{\alpha \alpha }\zeta _{\nu \lambda }\right]=0$$
(68)
tracing the equation above and considering the antisymmetric character of the tensor field $`\zeta ^{\rho \tau }`$ we get the Bianchi relation
$$b_1\left(p^\rho \zeta ^{\tau \sigma }+p^\sigma \zeta ^{\rho \tau }+p^\tau \zeta ^{\sigma \rho }\right)=0$$
(69)
If we set $`b_10`$, one possible solution of the relation (69) can be obtained if we put
$$\zeta ^{\mu \nu }=p^\mu \zeta ^\nu p^\nu \zeta ^\mu $$
(70)
i.e. the strength tensor of the Maxwell theory and the equation (67) becomes the Maxwell equation for the electromagnetic field .
We can obtain more two equations: the first one is gotten multiplying (65) on the right side by $`\left(C^1\right)_{\alpha _2\alpha }`$ we have
$$bp_\mu \left(\gamma ^\mu \gamma ^\nu \right)_{\alpha \alpha }\zeta _\nu =0$$
(71)
with the help of the trace properties for the Dirac matrices we obtain
$$b\left(p_\mu \zeta ^\mu \right)=0,$$
(72)
to get the second one we multiply (65) by $`\left(C^1\gamma ^\rho \gamma ^\sigma \right)_{\alpha _2\alpha }`$ and next we take the trace operation over the $`\gamma ^\mu `$-matrices to obtain
$$b\left(p^\mu \zeta ^\nu p^\nu \zeta ^\mu \right)=0$$
(73)
The equations (72) and (73) are compatible with the equations (67), (69) and (70) if and only if we set $`b=0`$.
### 3.3 Topological solutions
On the other hand we can get two additional solutions if we set $`b0`$ and $`b_1=0`$. Thus the first solution is getting when we solve the equation (72) choosing
$$\zeta ^\mu =p_\nu \zeta ^{\mu \nu }$$
(74)
where $`\zeta ^{\mu \nu }`$ is an antisymmetrical tensor field satisfying the equation (73).
And the second solution is founded when set the vector field in the equation (72) being
$$\zeta ^\mu =ϵ^{\mu \nu \alpha \beta }p_\nu \zeta _{\alpha \beta }$$
(75)
The equations (74) and (75) are topological field solutions for the spin 1 and spin 0 sectors , respectively. Such topological solutions were found in the massless DKP theory by Harish-Chandra and in the context of usual Klein-Gordon and Maxwell theories studying their higher tensor representations by Deser and Witten and Townsend .
## 4 Superspace Formulation
As a natural way we extend the previous analysis of the action and give the formulation in terms of superspace.
Firstly we consider the motion of the particle in the large superspace (big SUSY) $`(X_\mu ,\mathrm{\Theta }_\alpha )`$<sup>1</sup><sup>1</sup>1When the interaction is switched on, we must to include a complex grassmann spinor field $`\overline{\mathrm{\Theta }}_{\stackrel{.}{\alpha }}.`$ This enable us to consider theories with interacting charged particles. whose trajectory is parametrized by the proper supertime $`(\tau ,\eta _1,\eta _2)`$ of dimension $`\left(1/2\right)`$, here $`\eta _1,\eta _2`$ are the grassmann real superpartners of the convencional time $`\tau `$. In this way the coordinates of the particle are scalar superfields in the little superspace (little SUSY). For this case we have<sup>2</sup><sup>2</sup>2We recall that this form is valid only for the case of two indices $`i=1,2`$. If we want to analyse theories with a bigger internal symmetry $`O\left(N\right),`$ we need to include a more terms.
$`X_\mu (\tau ,\eta _1,\eta _2)`$ $`=`$ $`x_\mu \left(\tau \right)+i\eta _i\psi _\mu ^i\left(\tau \right)+i\eta _i\eta _jF_\mu ^{ij}\left(\tau \right)`$ (76)
$`\mathrm{\Theta }_\alpha (\tau ,\eta _1,\eta _2)`$ $`=`$ $`\theta _\alpha \left(\tau \right)+\eta _i\lambda _\alpha ^i\left(\tau \right)+\eta _i\eta _j_\alpha ^{ij}\left(\tau \right)`$ (77)
where $`i,j=1,2;`$ $`\psi _\mu ^i`$ is the grassman superpartner of the common coordinate $`x_\mu ;`$ $`\lambda _\alpha ^i`$ is a commuting majorana spinor, superpartner of the grassmann variables $`\theta _\alpha `$. $`F_\mu ^{ij}=F_\mu ^{ji}`$ and $`_\alpha ^{ij}=_\alpha ^{ji}`$ are antisymmetric fields.
In order to construct an action which is invariant under general transformations in superspace we introduce the supereinbein $`E_M^A(\tau ,\eta _1,\eta _2)`$, where $`M`$ \[$`A`$\] are a curved \[tangent\] indices and $`D_A=E_A^M_M`$ is the supercovariant general derivatives, here $`E_A^M`$ is the inverse of $`E_M^A`$. If we take a special gauge
$$E_M^\alpha =\mathrm{\Lambda }\overline{E}_M^\alpha ,E_M^a=\mathrm{\Lambda }^{1/2}\overline{E}_M^a$$
(78)
where
$$\overline{E}_\mu ^\alpha =1,\overline{E}_\mu ^a=0,\overline{E}_m^\alpha =i\eta ,\overline{E}_m^a=1$$
(79)
is the flat space supereinbein, then the superscalar field $`\mathrm{\Lambda }`$ an the derivative $`D_A`$ takes the form
$`\mathrm{\Lambda }(\tau ,\eta _1,\eta _2)`$ $`=`$ $`e\left(\tau \right)+i\eta _i\chi _i\left(\tau \right)+i\eta _i\eta _jf_{ij}\left(\tau \right),`$ (80)
$`\overline{D}_a`$ $``$ $`D_i={\displaystyle \frac{}{\eta ^i}}+i\eta _i{\displaystyle \frac{}{\tau }},\overline{D}_\alpha =_\tau `$ (81)
here $`e\left(\tau \right)`$ is the graviton field and $`\chi _i\left(\tau \right)`$ the gravitino field of the two-dimensional $`n=2`$ supergravity; $`f_{ij}=f_{ji}`$ is an antisymmetric matrix field. It is no difficult to prove that $`\left(\overline{D}_a\right)^2\left(D_i\right)^2=i_\tau `$
In this way the extension to superspace of the action (6), is given by<sup>3</sup><sup>3</sup>3The presence of the superscalar field $`\mathrm{\Lambda }`$ is to guarantee the local SUSY invariance.
$$S=\frac{1}{4}𝑑\tau 𝑑\eta _1𝑑\eta _2\mathrm{\Lambda }^1ϵ_{ij}D_iX_\mu D_jX^\mu $$
(82)
here $`ϵ_{ij}`$ is the antisymmetric matrix: $`ϵ_{12}=ϵ_{21}=1,`$ $`ϵ_{11}=ϵ_{22}=0`$. Using the property $`\mathrm{\Lambda }\mathrm{\Lambda }^1=1`$ for the supereinbein field we obtain
$`\mathrm{\Lambda }^1(\tau ,\eta _1,\eta _2)`$ $`=`$ $`e^1\left(\tau \right)ie^2\left(\tau \right)\eta _i\chi _i\left(\tau \right)ie^2\left(\tau \right)\eta _i\eta _jf_{ij}\left(\tau \right)`$ (83)
$`+e^3\left(\tau \right)\eta _i\eta _j\chi _i\left(\tau \right)\chi _j\left(\tau \right)`$
After some manipulations and integrating over the grassmann variables we have
$`S`$ $`=`$ $`{\displaystyle }d\tau ({\displaystyle \frac{1}{2}}e^1\stackrel{.}{x}^2+{\displaystyle \frac{i}{2}}e^1\psi _i\stackrel{.}{\psi }_i+{\displaystyle \frac{i}{2}}e^2\chi _i\psi _i\stackrel{.}{x}+{\displaystyle \frac{i}{2}}e^2f_{ij}\psi _i\psi _j`$ (84)
$`+{\displaystyle \frac{1}{2}}e^3\chi _i\psi _i\chi _j\psi _j+e^1F^2ie^2F_{ij}\chi _i\psi _j)`$
redefining the fields
$$\chi =e^{1/2}\chi ^{},\psi =e^{1/2}\psi ^{},f=ef^{},F=eF^{}$$
(85)
we obtain
$`S`$ $`=`$ $`{\displaystyle }d\tau ({\displaystyle \frac{1}{2}}e^1\stackrel{.}{x}^2+{\displaystyle \frac{i}{2}}\psi _i\stackrel{.}{\psi }_i+{\displaystyle \frac{i}{2}}e^1\chi _i\psi _i\stackrel{.}{x}+{\displaystyle \frac{i}{2}}f_{ij}\psi _i\psi _j`$ (86)
$`+{\displaystyle \frac{1}{2}}e^1\chi _i\psi _i\chi _j\psi _j+eF^2iF_{ij}\chi _i\psi _j)`$
we see that this action is identical to the proposed in (6) when we put $`F=\chi \psi `$, i.e. when the fermion coordinate and the gravitino field are coupled.
This shows that considering the correct inclusion of internal symmetries in the superspace formulation we obtain, in the special case, the same action proposed from the pseudoclassical point of view. The internal symmetry group $`O\left(N\right)`$ is connected to the number of grassmann variables $`\eta _i.`$
## 5 Conclusions
In this work we give an action for the massless DKP theory by using Grassmann variables and the consistence of the equations of motions are assured by means of the inclusion of boundary terms. We also verified the invariance under $`\tau `$-reparametrizations, local SUSY and internal group $`O(N)`$ transformations, the generators of these transformations are also found. We carried out the constraint analysis of the theory and verified that after quantization a possible inconsistency can appear, nevertheless the further analysis allow us to solve it with the introduction of some parameters that play a role of regulators of the theory. By the way an important result in this context was obtained, i.e. an additional topological solution for the spin 0 and 1 is derived from this model. As a natural continuation of the presented action we extended the studies to superspace formalism obtaining under some conditions the same initial pseudoclassical action.
For the further development of the theory we are working to accomplish the analysis through the most powerful method for a theory with constraints, i.e. via the BFV-BRST method, which can open the possibility of calculating the propagator of the resulting theory using the path integral representation. And, for further studies the inclusion of interactions (i.e., electromagnetic, Yang-Mills and gravitational fields) in the theory will be discussed.
### Acknowledgements
RC (grant 01/12611-7) and MP thank FAPESP and CAPES for full support, respectively, BMP thanks CNPq and FAPESP (grant 02/00222-9) for partial support, JSV thanks FAPESP (grant 00/03812-6) and FAPEMIG (grant 00193/06) for partial and full support, respectively.
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# Designable electron transport features in one-dimensional arrays of metallic nanoparticles: Monte Carlo study of the relation between shape and transport
## I Model
### I.1 Model system
In Fig. 2 a pictorial representation of the model system is shown. The system is modelled by a tight-binding tunneling Hamiltonian for spinless electrons. The free part consists of the electronic spectra of the reservoirs and the nanoparticles as well as the classical Coulomb interaction. The reservoirs shall consist of non-interacting electrons with a continuous, homogenous, infinite spectrum and they shall be in thermal equilibrium at all times. Their occupation numbers accordingly obey a Fermi distribution. The one-electron spectra on the $`Z`$ nanoparticles are considered to be discrete and on dot $`i`$ we consider explicitly $`Z_i`$ levels with level spacing $`\mathrm{\Delta }\epsilon _i`$. The Coulomb interaction is incorporated by the capacitance model which is detailed in the next section. We consider mutual capacitances between nearest neighbours like e.g. $`\stackrel{~}{C}_{12}`$ as well as between distant conductors like e.g. $`\stackrel{~}{C}_{L4}`$. Potentials can be applied to the reservoirs ($`\mathrm{\Phi }_L,\mathrm{\Phi }_R`$) and a gate ($`\mathrm{\Phi }_G`$). The uncontrollable influence of the DNA creates background charges on the dots (indicated by the random potentials $`\mathrm{\Phi }_{gi}`$).
The perturbation comprises three types of transitions: *Transitions within the array* We consider phonon assisted tunneling between the one-electron levels of nearest neighboring dots ($`w^{TA}`$). The tunneling matrix element $`t^A`$ is assumed to be equal for all tunneling transitions within the array. The phonon assisted tunneling stems from a linear coupling of a bosonic bath (bosonic bath 2) and the electronic degrees of freedom in the arrayNazarov92 ; Devoret94 . Here the bosonic bath shall model phonons in the substrate.
*Transitions between array and reservoir* We include tunneling between the continuous reservoirs and the one-electron levels of the outermost nanoparticles where no phonon assistance is considered here ($`w^{TRes}`$). The tunneling matrix element $`t^{Res}`$ for these transitions is considered to be equal for both reservoirs but generally different from $`t^A`$.
*Transitions on a single dot* These transitions among the one-electron levels on a single dot ($`w^{rel}`$) can be justified microscopically by a Fröhlich-HamiltonianMahan90 (coupling to bosonic bath 1).
We assume that the tunneling matrix elements have phases which fluctuate randomly due to slight temporary changes of the array geometry. Implicit summation over these phases prohibits any first order contributions of the tunneling matrix elements to the perturbation expansion.
### I.2 Transition energies
The transition rates which are given in the next section are determined by the transition energies, i.e. the changes in energy of the array that accompany the transitions of an electron from one level to another. It is composed of a change in electrostatic and one-electron energy. The former is computed within the capacitance model: the reservoirs, the gate and the dots are treated as macroscopic conductors with potentials $`\mathrm{\Phi }_i`$ and total charges $`Q_i`$. We call the gate and the reservoirs voltage nodes (with potentials $`\underset{¯}{\mathrm{\Phi }}_V`$ and charges $`\underset{¯}{Q}_V`$) since their potential is fixed, the dots are called charge nodes (with potentials $`\underset{¯}{\mathrm{\Phi }}_C`$ and charges $`\underset{¯}{Q}_C`$) since their charge is known and their potentials have to be determinedWasshuber01 . The potentials and the charges are related by the capacitance matrix $`C`$
$$\left(\begin{array}{c}\underset{¯}{Q}_C\\ \underset{¯}{Q}_V\end{array}\right)=C\left(\begin{array}{c}\underset{¯}{\mathrm{\Phi }}_C\\ \underset{¯}{\mathrm{\Phi }}_V\end{array}\right)=\left(\begin{array}{cc}C_c& C_b^{}\\ C_b& C_a\end{array}\right)\left(\begin{array}{c}\underset{¯}{\mathrm{\Phi }}_C\\ \underset{¯}{\mathrm{\Phi }}_V\end{array}\right)$$
(2)
The relation between the mutual capacitances $`\stackrel{~}{C}_{ij}`$ indicated in Fig. 2 and the elements of the capacitance matrix $`C_{ij}`$ is given by $`C_{ij}=\delta _{ij}_{k=0}^Z\stackrel{~}{C}_{jk}\stackrel{~}{C}_{ij}`$. The electrostatic energy can be written as:
$$E_{el}=\frac{1}{2}\left(\underset{¯}{Q}_C+\underset{¯}{Q}_C^{}+\underset{¯}{Q}_C^{bg}\right)^{}C_c^1\left(\underset{¯}{Q}_C+\underset{¯}{Q}_C^{}+\underset{¯}{Q}_C^{bg}\right)$$
(3)
where $`\underset{¯}{Q}_C^{}:=C_b^{}\underset{¯}{\mathrm{\Phi }}_V`$ is the polarization charge induced by the potentials $`\underset{¯}{\mathrm{\Phi }}_V`$ on the voltage nodes and $`\underset{¯}{Q}_C^{bg}`$ are fixed background charges which account for the electrostatic influence of the DNA and unknown fabricational details. The potentials on the charge nodes are then:
$$\underset{¯}{\mathrm{\Phi }}_C=C_c^1\left(\underset{¯}{Q}_C+\underset{¯}{Q}_C^{}+\underset{¯}{Q}_C^{bg}\right)$$
(4)
Furthermore we have to take into account the electronic spectra of the nanoparticles: the one-electron energy $`\epsilon _{li}`$ of level $`l`$ on dot $`i`$ is $`\epsilon _{li}=\epsilon _{Fi}+l\mathrm{\Delta }\epsilon _i`$ with the Fermi energy $`\epsilon _{Fi}`$ on dot $`i`$. We find the following transition energies<sup>1</sup><sup>1</sup>1We follow the considerations of Peter Hadley, presented on
http://qt.tn.tudelft.nl/$``$hadley/set/electrostatics.html :
*Transitions within the array* from level $`l`$ on dot $`i`$ to level $`l^{}`$ on dot $`i\pm 1`$
$`\mathrm{\Delta }E^{\pm ,i,l^{}l}`$ $`=l^{}\mathrm{\Delta }\epsilon _{i\pm 1}l\mathrm{\Delta }\epsilon _i`$
$`e`$ $`\left(\mathrm{\Phi }_{Ci\pm 1}\mathrm{\Phi }_{Ci}\right)+{\displaystyle \frac{e^2}{2}}\left(C_{ii}^12C_{ii\pm 1}^1+C_{i\pm 1i\pm 1}^1\right)`$ (5)
Note that the dependence of the potentials $`\underset{¯}{\mathrm{\Phi }}_C`$ on the charges $`\underset{¯}{Q}_C`$ is implied and that we include the constant terms $`\epsilon _{Fi}`$ in the definition of $`\underset{¯}{\mathrm{\Phi }}_C`$, i.e. we treat them like background charges:
$$\underset{¯}{\mathrm{\Phi }}_C=C_c^1\left(\underset{¯}{Q}_C+\underset{¯}{Q}_C^{}+\underset{¯}{\overset{~}{Q}}_C^{bg}\right)\text{ with }\underset{¯}{\overset{~}{Q}}_C^{bg}=\underset{¯}{Q}_C^{bg}\frac{1}{e}C_c\underset{¯}{\epsilon }_F$$
(6)
where the vector $`\underset{¯}{\epsilon }_F`$ contains the constants $`\epsilon _{Fi}`$.
*Transitions between array and reservoir* from/to lead $`\alpha `$ into/out of level $`l`$ on the neighbouring dot
$$\mathrm{\Delta }E^{\pm ,\alpha ,l}=\pm l\mathrm{\Delta }\epsilon _\beta \pm (e)\left(\mathrm{\Phi }_{C\beta }\mathrm{\Phi }_{V\alpha }\right)+C_{\beta \beta }^1\frac{e^2}{2}$$
(7)
where $`\beta =1`$ for $`\alpha =L`$ (left lead) and $`\beta =Z`$ for $`\alpha =R`$ (right lead). (Note that the nanoparticles are numbered consecutively from left to right.)
*Transitions on a single dot* from level $`l`$ to level $`l^{}`$ on dot $`i`$ only change the one-electron energy:
$$\mathrm{\Delta }E^{i,l^{}l}=\left(l^{}l\right)\mathrm{\Delta }\epsilon _i$$
(8)
It is convenient to express the transition energies in this way since we have to calculate the potentials $`\underset{¯}{\mathrm{\Phi }}_C`$ only once at the beginning of the simulation with (6) and then update them after each transition which is computationally cheap.
Obviously, the behaviour of the system depends strongly both on the values of the mutual capacitancesWhan96 and the one-electron level spacings. Concerning the former we determine the capacitance matrix numerically with the help of “FastCap”White91 using the geometry shown in Fig. 1. Only qualitatively we estimate the magnitude of the one-electron level spacing of a single, neutral, isolated nanoparticleHalperin86 :
$$\mathrm{\Delta }\epsilon _i=\frac{3\mathrm{}^2\pi }{2m_e\left(3\pi ^2n_i\right)^{\frac{1}{3}}}\frac{1}{r_i^3}=\frac{c}{r_i^3},c0,3010^{28}eVm^3$$
(9)
where $`n_i`$ is the electron density, $`r_i`$ the nanoparticle radius and the numerical value for $`c`$ was calculated for bulk gold. DFT calculations Barnett99 and experiments Zhang04 suggest that this estimation can at least reproduce the correct order of magnitude.
## II Method
### II.1 Transport theory
We assume that the tunneling rates defined below are much smaller than $`(k_BT)/\mathrm{}`$ so that we can treat the tunneling part of the Hamiltonian as a perturbation (weak coupling regime) and consider only the lowest non-vanishing order of the perturbation expansion (sequential tunneling regime). The reservoir and bath degrees of freedom of the density matrix are traced out which results in the reduced density matrix. We neglect its non-diagonal elements, which is justified if the broadening of the levels due to tunneling is small compared to the level spacing. The diagonal elements $`P_s`$ can then be interpreted as the probabilities of the array states $`|s`$ where the array state $`s`$ is given by all occupation numbers $`n_{li}`$ of level $`l`$ on dot $`i`$: $`|s=|\left\{n_{li}\right\}_{li}`$. The rates of electron transfer are calculated in golden rule approximation, thereby assuming that transport is incoherent and no coherent eigenstates are formed which stretch over the whole array, comparable to molecular orbitals. This is justified since in reality there are certainly processes which destroy the phase coherence – e.g. slight temporary changes in the geometry of the array.
Under the given assumptions the probabilities $`P_s`$ obey a master equation in the stationary limitSchoeller97
$$\dot{P}_s=\underset{s^{}}{}\left(w_{ss^{}}P_s^{}w_{s^{}s}P_s\right)=0$$
(10)
with golden rule rates $`w_{s^{}s}`$ from array state $`s`$ to array state $`s^{}`$. Each transition rate belongs to exactly one of the following sets:
*Transitions within the array* from level $`l`$ on dot $`i`$ to level $`l^{}`$ on dot $`i\pm 1`$
$`w^{TA}w^{\pm ,i,l^{},l}`$ $`=\mathrm{\Gamma }^AP\left(\mathrm{\Delta }E^{\pm ,i,l^{},l}\right),\mathrm{\Gamma }^A={\displaystyle \frac{2\pi }{\mathrm{}}}\left|t^A\right|^2`$
$`\text{with }P(E)`$ $`={\displaystyle \frac{1}{1+e^{\beta E}}}{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{\mathrm{\Gamma }}{E^2+\mathrm{\Gamma }^2}},\beta =\left(k_BT\right)^1`$ (11)
where $`t^A`$ is the tunnneling matrix element within the array. The function $`P(E)`$ is the probability per energy for the exchange of energy $`E`$ with the bosonic bath and therefore has to be normalized and must fulfill the condition of detailed balanceNazarov92 $`_{\mathrm{}}^{\mathrm{}}dEP(E)=1,P(E)=e^{\beta E}P(E)`$. The latter property stems from the nature of the bosonic bath: while it is always possible to emit energy, there must be excited bosons in the bath if energy shall be absorbed. The $`P(E)`$ function given above is an assumption which is preferably simple and has the required properties.
*Transitions between array and reservoir* from/to lead $`\alpha `$ into/out of level $`l`$ on the neighbouring dot
$$w^{TRes}w^{\pm ,\alpha ,l}=\mathrm{\Gamma }^\alpha f^\pm \left(\pm \mathrm{\Delta }E^{\pm ,\alpha ,l}\right),\mathrm{\Gamma }^\alpha =\frac{2\pi }{\mathrm{}}\left|t^{Res}\right|^2d_\alpha $$
(12)
where $`t^{Res}`$ is the tunnneling matrix element between a reservoir and the neighbouring dot, $`d_\alpha `$ is the density of states in the reservoir $`\alpha `$ and $`f^+(E)=(e^{\beta E}+1)^1`$, $`f^{}1f^+`$.
*Transitions on a single dot* from level $`l`$ to level $`l^{}`$ on dot $`i`$
$`w^{rel}w^{i,l^{}l}=\mathrm{\Gamma }^{rel}`$ $`(\mathrm{\Theta }(\mathrm{\Delta }E^{i,l^{}l})g(\mathrm{\Delta }E^{i,l^{}l})`$
$`+\mathrm{\Theta }(\mathrm{\Delta }E^{i,l^{}l})(1+g(\mathrm{\Delta }E^{i,l^{}l})))`$ (13)
where $`\mathrm{\Gamma }^{rel}`$ is the inverse relaxation time, $`g(E)=(e^{\beta E}1)^1`$ and $`\mathrm{\Theta }(E)`$ is the Heaviside step function. Since, as in the case of tunneling within the array, these transitions are possible due to the coupling to a bosonic bath, the rates fulfill detailed balance: $`w_{}^{i,l^{}l}=e^{\beta \mathrm{\Delta }E^{i,l^{}l}}w_{}^{i,ll^{}}\text{ for }\mathrm{\Delta }E^{i,l^{}l}>0`$. If the nanoparticles were isolated, i.e. if there were no transitions among them, it follows from the detailed balance property that in the stationary limit the occupation numbers $`n_{li}`$ obey a Fermi distribution. So while the tunneling transitions drive the electron distribution out of equilibrium the transitions on single dots effectively cool the electrons.
### II.2 Monte Carlo algorithm
If we consider $`Z_i`$ levels on dot $`i`$ there are $`2^{_iZ_i}`$ possible array states. So the master equation (10) defies a direct numerical solution except for very small systems. Therefore we employ a Monte Carlo (MC) method to retrieve the quantities of interest: the current and the shot noise. The key idea of the MC methodHonerkamp94 in this context is to discretize time and to get from the transition rate $`w_{s^{}s}`$ a transition probability $`\pi _{s^{}s}`$ by multiplying the rate with a finite time step $`\mathrm{\Delta }t(s)`$ which may depend on the present system state $`s`$: $`\pi _{s^{}s}(\mathrm{\Delta }t(s))=w_{s^{}s}\mathrm{\Delta }t(s)`$ The time step has to be chosen sufficiently small so that for the total probability $`\pi _s(\mathrm{\Delta }t(s))`$ to leave the state $`s`$ it holds
$$\pi _s(\mathrm{\Delta }t(s))=\underset{s^{}}{}\pi _{s^{}s}(\mathrm{\Delta }t(s))1s$$
(14)
The transitions among the array states, which are governed by the probabilities $`\pi _{s^{}s}`$, constitute a stochastic process which can be simulated with the help of random numbers. In contrast to the MC method used in statistical physicsBinder97 which samples a (grand) canonical ensemble, the system examined here is generally out of equilibrium. In each valid MC algorithm the probability $`\pi _s(\mathrm{\Delta }t(s))`$ must be properly represented. We write it as follows:
$$\pi _{s^{}s}(\mathrm{\Delta }t(s))=w_{s^{}s}\mathrm{\Delta }t(s)=\frac{w_{s^{}s}}{w_0}\frac{1}{D(s)}D(s)w_0\mathrm{\Delta }t(s)$$
(15)
where $`D(s)`$ is the number of possible transitions out of a state $`s`$ and $`w_0`$ is an upper bound for the transition rates which is also called attempt frequency. If we choose the time step $`\mathrm{\Delta }t(s)`$ as $`\mathrm{\Delta }t(s)=1/(D(s)w_0)`$, the transition probability reduces to:
$$\pi _{s^{}s}(\mathrm{\Delta }t(s))=\frac{w_{s^{}s}}{w_0}\frac{1}{D(s)}$$
(16)
Note that with the upper choice of the time step the requirement (14) is fulfilled since $`w_{s^{}s}w_0`$ and the number of nonzero addends is $`D(s)`$. It is straightforward to imagine the corresponding algorithm if you regard each factor in (16) as an independent probability, see Fig. 3.
We found that even though our algorithm needs two random numbers per time step and an attempt transition may be rejected, see step 3 in Fig. 3, it is convenient since only one rate has to be computed per time step. For systems with many possible transitions it is faster than algorithms which need the transition rates of all possible transitions in each time stepAmman89 ; Amman89b . Furthermore our algorithm has the convenient property to be self-adapting to the system size: since $`\mathrm{\Delta }t(s)\left(D(s)\right)^1`$ the number of performed steps scales with the amount of possible transitions if the runtime of the simulation is fixed.
Despite the optimized algorithm we can take only a finite number of one-electron levels into account, of course. In most cases, however, we want to model wide bands on the dots, i.e. we do not want to observe any impact of the finiteness of the electron spectrum. Practically we increase the number of considered levels until the results become steady for the maximum bias voltage which we want to apply. In order to minimize the computing time, i.e. real time, we consider as few levels as possible, of course. Therefore we center the spectra around the highest occupied levels in the ground state (i.e. the state assumed in equilibrium at zero temperature) since for small bias voltages and low temperatures only low-lying excitations about the Fermi edge appear. To determine the ground state the total energy of the system which includes the electrostatic energy (3) and the one-electron energies has to be minimized. Due to the discreteness of the charges $`\underset{¯}{Q}_C`$ this is a non-trivial minimization problem which we do not address here.
The initial state of a simulation run is always the ground state, i.e. all spectra are half-filled. Each run starts with an equilibration period $`T_{equi}`$ in which we let the system evolve, by iterating the algorithm (Fig. 3), without sampling the assumed states. So the system can reach its stationary state. In the subsequent measurement period $`T_{meas}`$ the quantities of interest are retrieved. The current can simply be obtained by counting the electrons that are transferred e.g. between the left lead and the first nanoparticle and dividing by the length of the measurement period $`T_{meas}`$: $`I_L=Q_L(T_{meas})/T_{meas}`$ where $`Q_L(T_{meas})`$ is the charge that is transferred during the measurement time $`T_{meas}`$. Note that in the stationary state the current through all tunneling barriers is the same. For a fixed set of parameters the simulation is repeated with different seeds for the random number generator in order to get statistically independent runs. With this ensemble we can determine the statistical standard error of the mean current. The (zero-frequency) shot noise $`S_{I_L}(0)`$ and the Fano factor $`F`$ can also be estimatedKorotkov94 :
$`S_{I_L}(0)`$ $`={\displaystyle \frac{2}{T_{meas}}}\left(Q_L(T_{meas})^2Q_L(T_{meas})^2\right)`$
$`F={\displaystyle \frac{S_{I_L}(0)}{2eI_L}}`$ (17)
where $`\mathrm{}`$ denotes the ensemble average. As for the current, in the stationary limit the shot noise is the same for all tunneling barriers. Note that all given quantities are estimators which become exact in the limit $`T_{meas}\mathrm{}`$.
The validity of our method was checked by comparing our results with the solution of the master equation for a small system. With the same benchmark sensible values for the simulation parameters (equilibration and measurement time, size of the ensemble) were obtained. Due to the self-adaptive property of the algorithm these parameters are also suitable for bigger systems. For each geometry we increased the number of considered one-electron levels until the IV curves became steady. Furthermore we found that the computing time scales linearly with the number of nanoparticles $`Z`$ and quadratically with the maximum number of considered levels on a dot ($`\mathrm{max}\{\stackrel{~}{Z}_i\}_i`$). It is exponentially smaller than the computing time needed for the direct solution of the master equation which scales as $`(2^{_i\stackrel{~}{Z}_i})^3`$. However, note that the MC method is not equivalent to a solution of the master equation: though the current and shot noise may be retrieved from a MC simulation, it is practically impossible to determine all probabilities $`P_s`$ correctly.
### II.3 Charge states
To interpret the simulated results we draw a sample out of a single simulation run and look at the probabilities of charge configurations and mean rates among them. To get the probability $`P_C`$ of a certain charge configuration $`C\{Q_{Ci}\}_i`$, we sum the MC times during which this configuration is assumed and divide by the total MC time. The mean rate $`w_{C^{}C}^{\pm ,i}`$ through a tunneling junction $`i`$ from state $`C`$ to state $`C^{}`$ to the right or left respectively is determined by counting the transitions between state $`C`$ and state $`C^{}`$ by tunneling in the given direction through junction $`i`$ and dividing the sum by the MC time that is spend in the state $`C`$. The current through a tunneling junction $`i`$ can then be written as $`I_i=_{C,C^{}}_\sigma \sigma w_{C^{}C}^{\sigma ,i}P_C\text{ with }\sigma =\pm `$. In the stationary state the currents through the junctions are all equal to one another $`I_i=I_ji=j`$, so we can write the current as
$$I=\frac{1}{Z+1}\underset{i}{}I_i=\frac{1}{Z+1}\underset{C,C^{}}{}P_Cw_{C^{}C}=\underset{C,C^{}}{}I_{C^{}C}$$
(18)
with $`w_{C^{}C}=_{i,\sigma }\sigma w_{C^{}C}^{\sigma ,i}`$ and $`I_{C^{}C}=\frac{1}{Z+1}P_Cw_{C^{}C}`$. The partial current $`I_{C^{}C}`$ from charge state $`C`$ to charge state $`C^{}`$ has a positive sign if it flows from left to right and the opposite sign for the opposite direction. If there are partial currents which flow between the same states – then they have necessarily opposite directions, we keep only the difference of them so that there is only one net partial current between two charge states.
## III Results
The following conventions hold for all shown results. The general geometrical setup is the one already shown in Fig. 1. All energies are normalized to the maximum level spacing that occurs in the array ($`\mathrm{max}\left\{\mathrm{\Delta }\epsilon _i\right\}_i`$). Instead of voltages (potentials) we use potential energies $`eV`$ and charges are given in units of $`e`$. We open the bias voltage window symmetrically (i.e. $`\mathrm{\Phi }_L=\mathrm{\Phi }_R`$) because we do not want to introduce artificially an additional asymmetry. We define $`V^{}=(e\mathrm{\Phi }_R)/(\mathrm{\Delta }\epsilon _{max})`$ so that the current is positive (i.e. flows from left to right) if $`V^{}`$ is positive. We divide all calculated golden rule rates by $`\mathrm{\Gamma }^A\mathrm{max}(P(\mathrm{\Delta }E))`$, so that the maximum rate within the array is equal to $`1`$. We set the maximum rate between array and reservoir equal to $`0.1`$ since we assume that the tunnel coupling within the array is stronger. Current and charge are expressed in units of the elementary charge $`e`$. We normalize the current to the maximum rate in the bulk of the array $`I^{}=(I/e)/(\mathrm{\Gamma }^A\mathrm{max}(P(E)))`$. The current is defined to be positive if it flows from left to right. The parameter $`\mathrm{\Gamma }`$ of the function $`P(E)`$ is set to $`2`$, which is small enough to see individual one-electron levels and big enough to give a sufficiently high current. No error bars appear in the following results since the relative statistical error is always only about 0.1%, so error bars would not be visible.
### III.1 Generalizations of results for single quantum dots
#### III.1.1 Interplay of one-electron level spacing and charging energy
In a 2 nanoparticle array we investigate the case $`\mathrm{\Delta }\epsilon _i<|\frac{e^2}{2}(C^1)_{ii}|`$, i.e. the level spacing is smaller than the typical charging energy. In agreement with results for a single dotAverin91 we find that the level spacing imposes a fine structure on the Coulomb staircase which is related to the charging energy. In Fig. 4 the corresponding IV characteristics are shown. Looking at the partial currents $`I_{C^{}C}`$ defined above we find that for the first 4 steps there is only one relevant transport path in the charge configuration space: $`(0,0)(1,0)(1,+1)(0,0)`$ where the charges are given as differences to the ground state charges. So the first 4 steps must be due to the level spacing. On the 5. step a second path is relevant: $`(0,0)(0,+1)(1,+1)(0,0)`$. For higher temperatures the fine structure – i.e. the features due to the one-electron levels – is smeared out, see Fig. 4, while the typical Coulomb staircase remains: in the middle of each plateau a new transport path becomes available.
#### III.1.2 Influence of dissipation
We consider a 3 nanoparticle array at $`T=0`$ and fixed gate voltage and study the impact of a finite relaxation rate, see Fig. 5. Generalising results for a single dotAverin90 we find that without relaxation the electrons overheat and consequently the structures in the IV characteristics are smoothed. Strong relaxation, on the other hand, which effectively cools the electrons, sharpens the steps and decreases the absolute value of the current. The IV curves for intermediate relaxation rates lie between the curves belonging to the extreme cases. These tendencies can also be observed for other gate voltages and other array lengths. They can be understood by noting that a high relaxation rate keeps the mean occupation of the one-electron levels on a dot close to the equilibrium i.e. Fermi distribution. For $`T=0`$ this leads to the formation of a defined Fermi edge on the nanoparticles. That is the reason why the steps in the IV curve become distinct. Electrons above the Fermi energy relax to lower lying levels and the corresponding transition energy is dissipated in the bosonic bath. Such electrons can perform fewer transitions as without dissipation since less energy is available. That is the reason why the current generally decreases with an increasing relaxation rate.
### III.2 Results uniquely related to the array geometry, designable effects
We examine arrays of nanoparticles with uniformly growing diameters, see Fig. 6 for the case of 4 nanoparticles. We choose this special geometry for two reasons: on the one hand the array can be enlarged in a defined way. We start with two nanoparticles and let the number of nanoparticles and therefore the array length grow so that we discover certain features which evolve systematically with increased length. On the other hand it is interesting to combine small and big nanoparticles since they differ strongly both in level spacing and charging energy. In the case of 5 nanoparticles, which is the longest array that is studied here, the diameters of the nanoparticles range from $`1nm`$ to $`1.8nm`$ in steps of $`0.2nm`$. The level spacing of the smallest particle is estimated to be about $`0.24eV`$ according to eq. (9) so that a voltage $`V^{}`$ of $`1`$ in the following curves equals then $`0.24V`$. This energy equals a temperature of about $`2800K`$ so that we find level spacing related features in the current or shot noise at finite temperatures. Since the size of the nanoparticles increases from left to right, the level spacing decreases in the same direction.
#### III.2.1 Strong asymmetry of IV characteristics
In Fig. 7 the IV characteristics of an array of 4 nanoparticles exemplify the typical IV characteristics of the array geometry. The various curves correspond to different gate voltages. We observe that for all array lengths two regions with different behaviours can be distinguished. For smaller bias voltages the curves exhibit a striking asymmetry with respect to reversal of the bias voltage. For example the threshold voltage is different for positive or negative bias voltages ($`V_T^\pm `$). This asymmetry is more pronounced for longer arrays and for curves with higher threshold voltages, especially if the $`V_T^\pm `$ exceed the offset voltage $`V_{\text{off}}`$. The position, width and height of the steps in this bias voltage region and $`V_T^\pm `$ depend strongly on the gate voltage. For larger bias voltages, however, the IV curves shown in Fig. 7 approximately coincide and become symmetric (with respect to reversal of the bias voltage). For small bias voltages the number of many-particle states or the number of paths through the charge state space that take part in transport is smaller. So it is important which states or paths actually participate. This is in turn influenced by the gate voltage since it shifts the energetic position of the charge states and determines therefore which paths are available. For bigger bias voltages many states or paths participate so it should be less important whether a certain state takes part or not: what we observe is their “average” contribution. That is why the curves for different gate voltages coincide for high bias voltages. It is also the reason for the disappearance of the asymmetry in the same region. A detailed discussion of the asymmetry can be found in the appendix A.
#### III.2.2 Overall conductance and offset voltage
As already mentioned above, the IV curves shown in Fig. 7 approximately coincide and become symmetric for larger bias voltages. In that region the current increases linearly, apart from very small steps. We fit this linear segment with a straight line and define its slope as the overall conductance $`G`$. The offset voltage $`V_{\text{off}}`$, i.e. the $`V^{}`$-intercept of the fitted line, can be thought of as a kind of “mean” threshold voltage. The actual threshold voltage of a single curve obviously depends on the gate voltage, as already mentioned. $`G`$ and $`V_{\text{off}}`$ are approximately the same for positive and negative bias voltages. Our results for $`G`$ and $`V_{\text{off}}`$ with respect to the number of nanoparticles $`Z`$ are compiled in the two left columns of Table 1 for the case without relaxation ($`w^{rel}=0`$). We find that $`G`$ and $`V_{\text{off}}`$ both increase with an increasing array length.
The increase of $`V_{\text{off}}`$ with an increasing array length has also been observed elsewhere Melsen97 ; Nguyen01 so it seems to be quite generic for arrays. One factor which contributes to this tendency is the decrease of the capacitances with increasing distance between the conductors. Especially the capacitances that couple the leads with the nanoparticles $`\stackrel{~}{C}_{\alpha i}`$, see eq. (2), are important. The cpacitances extracted from the geometry with the help of “FastCap” White91 do not decrease linearly, like e.g. for a simple parallel plate capacitor, but approximately exponential. This is reasonable since the nanoparticles partially screen the electric field. For two neighbouring nanoparticles that are in the middle of the array the difference between the coupling capacitances is small. So we have to apply a high bias voltage to create a potential difference between these two particles which permits a tunneling transition, see eq. (II.1). Therefore the threshold voltages and $`V_{\text{off}}`$ increase with increasing array length. Concerning $`G`$ the change of the number of transport paths with increasing bias voltage plays an important role. This change is in turn determined by the change of the number of possible transitions. A certain transition between two neighbouring dots becomes possible when the potential gradient between the dots becomes high enough. If we assume that this gradient is roughly the bias voltage divided by the number of tunnel junctions then we should expect that $`G`$ decreases with increasing array length. On the other hand we “grow” the array by adding bigger nanoparticles with a higher density of one-electron states. Given that the number of extra electrons in the array is small a higher number of one-electron levels within the bias voltage window results in a higher current. This effect might overcompensate the reduced potential gradient which would result in the observed behaviour of $`G`$. To check this assertion we have artificially raised the level spacing on the last dot in the 3 particle array so that it is equal to the level spacing on the first dot: in this case we find that the overall conductance $`G`$ is only $`0.1`$ compared to $`0.14`$ with the original level spacing. So $`G`$ can indeed be lowered by raising the level spacing on the last dot. This supports our assertion. Obviously the offset voltage and the overall conductance can be tailored by the choice of the nanoparticle sizes.
For the case of high relaxation rates $`G`$ and $`V_{\text{off}}`$ are compiled in the two right columns of Table 1. Comparing with the case without relaxation we observe that the offset voltage is generally bigger if there is relaxation while the overall conductance is approximately the same. As argued in section III.1.2 the current generally decreases with an increasing relaxation rate. This homogeneous downward shift of the IV curve correspondingly increases the offset voltage while the overall conductance remains unaltered.
#### III.2.3 Giant step, giant Fano factor
As already stated above we observe that in the IV characteristics for several gate voltages the first current step to the left or to the right of the Coulomb blockade region is strikingly high. One example is shown in Fig. 8. Furthermore the higher step is accompanied by a peaked giant Fano factor. To determine the origin of the big step in the current and the giant peak in the Fano factor we record the dynamics of the simulation, i.e. in Fig. 9 the transferred charge is plotted against the MC time for a single simulation run. This approach has already been used in a MC study of the electron transport properties in a different modelKoch04 in which also a giant Fano factor is found. In Fig. 9 we have recorded the dynamics for the voltages marked in Fig. 8. For the top of the peak marked with $`2.`$ we observe comparatively long periods without charge transfer and intermediate tunneling avalanches. Looking at the definition of the Fano factor $`F`$, eq. (17), there are two ways to conclude that the observed dynamics result in a super poissonian noise (i.e. $`F>1`$). On the one hand, one can regard the avalanches on a much bigger time scale as single, statistically independent transfer processes in which an effective charge bigger than $`e`$ is transported. Then one uses the Poisson value for the shot noise $`2qI`$ with an effective charge $`q`$ bigger than $`e`$ which results in a Fano factor $`F>1`$. On the other hand, one can regard the avalanches as a bunching of the tunneling processes and this positive correlation leads per definition to $`F>1`$. So both interpretations deduce the super poissonian shot noise from the observed dynamics. For bigger bias voltages, $`3.`$ and $`4.`$, the length of the periods without charge transfer is reduced so the Fano factor is reduced correspondingly. The big step height can be understood analogously: during the periods without charge transfer a quasi-blocking state is assumed , i.e. the sum of rates that lead out of this state is much smaller than the sum of the rates that lead into the same state, see Fig. 10. So if the state is visited the system rests there for a comparatively long time. If no rates led out of that state it would be a real blocking state and the current would be zero since the dynamics of the system would ultimately end up in this state. The probability $`P_s`$ of a blocking state is therefore 1. In the shown case we find that a neutral many-electron state $`(0,0,0,0)`$ serves as a quasi-blocking state (respectively blocking state in the region of Coulomb blockade). If the system leaves it, current can flow. This current is however carried by other states with much higher rates among them. Therefore the high step in the IV characteristics.
#### III.2.4 Designable NDC effect
So far we have modelled infinite spectra, see section II.2. Now we intentionally consider only a few levels on each nanoparticle. This is interesting since the electronic structure of small nanoparticles differs in general strongly from bulk materialBarnett99 . We assume on each nanoparticle a finite “band” of discrete, equally spaced energy levels. The level spacing is estimated as given above, see eq. (9). Under these assumptions we find for an array of two nanoparticles a NDC which depends strongly on the ratio between transition rates within the array ($`w^{TA}`$) and the rates between array and reservoirs ($`w^{TRes}`$), see Fig. 11. This ratio can be tuned by varying the distance between array and contacting reservoirs: the smaller the distance the higher the rate of tunneling between array and reservoirs ($`w^{TRes}`$). So this is an example of an effect which strongly depends on the geometry and which is therefore designable. The origin of the NDC can be understood by looking at the transition energies, see section I.2. In Fig. 12 we sketch the energies $`\left\{l\mathrm{\Delta }\epsilon _i+e\mathrm{\Phi }_{Ci}\right\}_{l,i}`$ for the two nanoparticles. For $`T=0`$ a transition between the nanoparticles can only happen if the difference between those energies is bigger than $`\left(C_{ii}^12C_{ii\pm 1}^1+C_{i\pm 1i\pm 1}^1\right)`$, see eq. (5). For finite temperatures this restriction is softened by the coupling to the bosonic bath which results in the $`P(E)`$ function. The offset between the sketched spectra is determined by the difference between the potentials $`e\mathrm{\Phi }_{Ci}`$ which depend in turn on the charges on the nanoparticles. If the tunneling rates between array and reservoirs are sufficiently large a high charge gradient accumulates and therefore a large offset between the spectra occurs. Since the spectra are finite, the transition energies of possible transitions are big if the offset between the spectra is large. The corresponding value of the $`P(E)`$ function is small. Consequently the transition rates among the nanoparticles become smaller with increased voltage which results in the NDC. If the $`P(E)`$ function is broad, this effect is softened, see eq. (II.1). The higher the factor $`\mathrm{\Gamma }`$ the broader the $`P(E)`$ function. Therefore the NDC effect is less distinct for higher values of $`\mathrm{\Gamma }`$.
## IV Conclusion
We have shown that an improved MC algorithm can cope with the complexity of the electronic transport through nanoparticle arrays with discrete electronic spectra. Though we have considered linear arrays only we found a huge variety of transport functionalities. Of course, other geometries would also be worth studying, e.g. ring-shaped devicesShin98 can be used as charge storage elements. Since arrays can be made self-assembling, are robust against fabricational imperfections and show, as we found out, transport features that can be designed via the geometry, we consider them to be ideal building blocks for electronic architecture on the nanoscale.
###### Acknowledgements.
The authors gratefully acknowledge the financial support of the VW Foundation and the German National Merit Foundation. The authors would like to thank M. Holtschneider, S. Jakobs, K. C. Nowack, M. Noyong, F. Reininghaus and M. R. Wegewijs for useful discussions.
*
## Appendix A Asymmetry concerning reversal of bias voltage
The occurence of IV curves which are asymmetric concerning the reversal of the bias voltage would be surprising if the transport was calculated within the Landauer-Büttiker formalism Landauer88 ; Buettiker86 . Within this approach the IV curves must be symmetric due to time reversal symmetry: for each wave which travels through the system from left to right there is a time reversed one with the opposite direction of propagation. But the Landauer-Büttiker formalism can only be applied for coherent transport without interaction. Both suppositions are violated in our case: we assume that the phase information is lost at every tunneling event and we include Coulomb interaction, see sections I.1 and II.1. Therefore, in our case, asymmetry is present in general unless special symmetries impose symmetric curves. We will discuss now two symmetries that are actually relevant in our case. By explicitly showing how these symmetries are broken by the Coulomb interaction and the distribution of the density of states (inverse level spacing) the occurence of asymmetric IV curves is rendered plausible. We adopt the framework of the orthodox theory for these considerations since the basic mechanisms can be understood by looking at charge states alone.
### A.1 Particle-hole symmetry
This symmetry means that two paths through the charge state space with opposite charges are equivalent. Equivalent here means that they appear at the same absolute value but opposite sign of the bias voltage and that the rates for corresponding transitions are equal. E.g. for a 2 nanoparticle array this means that if the path $`(1,2)(2,2)(1,3)(1,2)`$ appears at $`V=x`$ then the path $`(1,2)(2,2)(1,3)(1,2)`$ appears at $`V=x`$ and the rates for corresponding transitions are equal. This symmetry is present if there is no coupling to the gate and if there are no background charges. To conclude this we have to look at eq. (6) for the potentials $`\underset{¯}{\mathrm{\Phi }}_C`$ on the nanoparticles: $`\underset{¯}{\mathrm{\Phi }}_C=C_c^1\left(\underset{¯}{Q}_C+\underset{¯}{Q}_C^{}+\underset{¯}{\overset{~}{Q}}_C^{bg}\right)`$. If there is no coupling to the gate the polarization charges on the nanoparticles are reversed if the bias voltage is reversed: $`VV\underset{¯}{Q}_C^{}\underset{¯}{Q}_C^{}`$. If we reverse the charges $`\underset{¯}{Q}_C`$ at the same time and there are no background charges $`\underset{¯}{\overset{~}{Q}}_C^{bg}`$, the potentials $`\underset{¯}{\mathrm{\Phi }}_C`$ are reversed: $`VV,\underset{¯}{Q}_C\underset{¯}{Q}_C\underset{¯}{\mathrm{\Phi }}_C\underset{¯}{\mathrm{\Phi }}_C`$. Since the potentials $`\underset{¯}{\mathrm{\Phi }}_C`$ determine the transition rates two charge states with opposite charges are equivalent in the sense explained above. Obviously finite background charges or coupling to the gate break this symmetry.
### A.2 Inversion symmetry
This symmetry means that two paths through the charge state space which are mirror images of each other are equivalent where the mirror plane shall be situated in the middle of the array and equivalent is meant in the sense explained above. E.g. for a 2 nanoparticle array this means that if the path $`(1,2)(2,2)(1,3)(1,2)`$ appears at $`V=x`$ then the path $`(2,1)(2,2)(3,1)(2,1)`$ appears at $`V=x`$ and the rates for corresponding transitions are equal. This symmetry is obviously present if the whole setup is symmetric with respect to a mirror plane in the middle of the array. Both an asymmetric capacitance matrix and an asymmetric distribution of the density of states on the nanoparticles break this symmetry.
#### A.2.1 Asymmetric capacitance matrix
We look at an array of 2 nanoparticles with different sizes, the left nanoparticle ($`A`$) shall be smaller than the other ($`B`$). The capacitative coupling (between nodes) shall be negligible so that the capacitance matrix is diagonal. Then each nanoparticle can be characterized by a single capacitance, $`C_A`$ and $`C_B`$ respectively. Since nanoparticle $`A`$ is smaller, it holds $`C_AC_B`$. The opposite holds for the charging energies $`E_{C_A}E_{C_B}`$ so more energy is needed to charge nanoparticle $`A`$. Now it is clear that the paths $`(0,0)(1,0)(0,1)(0,0)`$ and $`(0,0)(0,1)(1,0)(0,0)`$ are not equivalent. The asymmetric capacitance matrix breaks the inversion symmetry.
#### A.2.2 Asymmetric distribution of the density of states
The capacitance matrix of the 2 nanoparticle array shall now be symmetric and still diagonal, so the charging energy is the same for both nanoparticles. The density of states, however, shall be $`D_A`$ on the left and $`D_B`$ on the right nanoparticle. (This corresponds to different level spacings in our model.) The density of states in both reservoirs shall be $`D_{Res}`$. The energy change for the three tunneling transitions that appear in the path $`(0,0)(1,0)(0,1)(0,0)`$ shall be denoted by $`\mathrm{\Delta }E_1`$,$`\mathrm{\Delta }E_2`$ and $`\mathrm{\Delta }E_3`$. (The index denotes the number of the transition e.g. $`\mathrm{\Delta }E_1`$ is the energy difference occuring at the transition $`(0,0)(1,0)`$.) The corresponding transition rates are $`w_1`$, $`w_2`$ and $`w_3`$. For zero temperature the orthodox theory rates are:
$`w_1`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}\left|t^{Res}\right|^2(\mathrm{\Delta }E_1)\mathrm{\Theta }(\mathrm{\Delta }E_1)D_{Res}D_A`$ (19)
$`w_2`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}\left|t^A\right|^2(\mathrm{\Delta }E_2)\mathrm{\Theta }(\mathrm{\Delta }E_2)D_AD_B`$ (20)
$`w_3`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}\left|t^{Res}\right|^2(\mathrm{\Delta }E_3)\mathrm{\Theta }(\mathrm{\Delta }E_3)D_BD_{Res}`$ (21)
If we reverse the bias voltage and look at the path $`(0,0)(0,1)(1,0)(0,0)`$, we get different energy differences $`\mathrm{\Delta }E_1^{}`$,$`\mathrm{\Delta }E_2^{}`$ and $`\mathrm{\Delta }E_3^{}`$ and rates $`w_1^{}`$, $`w_2^{}`$ and $`w_3^{}`$. (Note that here e.g. $`\mathrm{\Delta }E_1^{}`$ is the energy difference occuring at the transition $`(0,0)(0,1)`$.) Since the capacitance matrix is symmetric, it holds $`\mathrm{\Delta }E_i^{}=\mathrm{\Delta }E_ii`$. So we get:
$`w_1^{}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}\left|t^{Res}\right|^2(\mathrm{\Delta }E_1)\mathrm{\Theta }(\mathrm{\Delta }E_1)D_{Res}D_B`$ (22)
$`w_2^{}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}\left|t^A\right|^2(\mathrm{\Delta }E_2)\mathrm{\Theta }(\mathrm{\Delta }E_2)D_AD_B`$ (23)
$`w_3^{}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}\left|t^{Res}\right|^2(\mathrm{\Delta }E_3)\mathrm{\Theta }(\mathrm{\Delta }E_3)D_AD_{Res}`$ (24)
Obviously the rates $`\left\{w_i\right\}_i`$ are different from the rates $`\left\{w_i^{}\right\}_i`$ due to the different density of states $`\left\{D_i\right\}_i`$ so the considered paths are not equivalent. An asymmetric distribution of the densities of states breaks the inversion symmetry.
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# Superfluidity of spin-1 bosons in optical lattices
## I INTRODUCTION
Originally discovered in the system of liquid helium and later in the context of superconductors, superfluidity is a hallmark property of interacting quantum fluids and encompasses a whole class of fundamental phenomena such as the absence of viscosity, persistent currents and quantized vortices. With the achievement of Bose-Einstein condensation (BEC) in alkali-metal atoms, the weakly interacting Bose gases have served as an idealized model of the superfluid and a test ground of macroscopic quantum effects at low temperatures. Recently, rapid advances of experimental techniques in optical traps open up a prospect for the study of the superfluidity of BEC trapped in periodic potentials, which has attracted fast growing interests both experimentally and theoretically . The main reason is that the optical lattices possess controllable potential depths and lattice constants by adjusting the intensity of the laser beams in realistic experiments and moreover the lattice is basically defect free. In addition, the great advantage of optical traps is that it liberates the spin degree of freedom and provides us an opportunity to test spin-dependent quantum phenomena that are absent in the scale-condensate cases. Theoretical studies have predicted a variety of novel phenomena of spinor condensates such as fragmented condensation , skyrmion excitations and propagation of spin waves . The quantum phase transition from a superfluid to a Mott-insulator (SF-MI) phase in spinor BEC has been observed in experiments . In the Mott-insulator phase (MIP), atoms are localized; the particle-number fluctuations at each lattice site are suppressed so that there is no phase coherence across the lattice. When the tunnel coupling through the interwell barriers becomes large compared to the atom–atom interactions, the system undergoes a phase transition into the superfluid phase (SFP) in which the atom number per site is random and hence wave function exhibits long-range phase coherence. That means one can go from the regime in which the interaction energy dominates (high barrier of periodic potential) to the regime where the kinetic energy is the leading part (low barrier of periodic potential) by varying the intensity of laser beams and vice versa. The SF-MI transition has been also investigated in Refs., where the Bose–Hubbard model is introduced as the starting point, and the analytic phase-transition condition and phase diagram have been obtained by using a perturbation expansion over the superfluid order parameter with on-site zero-order energy spectrum which, although gives rise to a reasonable description of the MIP, the SFP is not described explicitly. In the present paper we study the spinor BEC in a parameter region such that its ground state stands deeply in the SFP . The Bogliubov approach is used to obtain the explicit excitation energy spectra of weakly interacting spin–1 atoms in an optical lattice and the superfluidity is explained explicitly in terms of the energy spectra. The critical velocities known as the Landau criterion for the superfluid phase are evaluated for the <sup>23</sup>Na atoms and are seen to be realizable in practical experiments. It is also demonstrated that the critical velocities are spin component dependent and controllable by adjusting of the lattice parameters. Our result may throw light on the experimental observation of the persistent atom-current density for the spinor-atom matter waves.
## II BOGLIUBOV METHOD AND ENERGY SPECTRUM
Alkali-metal atoms with nuclear spin $`I=3/2,`$ such as <sup>23</sup>Na, <sup>39</sup>K, and <sup>87</sup>Rb, behave at low temperatures like simple bosons with a hyperfine spin $`f=1`$. The most general model Hamiltonian for the dilute gas of bosonic atoms with hyperfine spin $`f=1`$ trapped in the optical potential can be written, in the second-quantization notation, as
$`\widehat{H}`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle d^3X\widehat{\psi }_\alpha ^{}\left(X\right)\left(\frac{^2}{2M}+V_0\left(X\right)+V_T\left(X\right)\right)\widehat{\psi }_\alpha \left(X\right)}`$ (3)
$`+{\displaystyle \frac{C_0}{2}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle d^3X\widehat{\psi }_\alpha ^{}\left(X\right)\widehat{\psi }_\beta ^{}\left(X\right)\widehat{\psi }_\beta \left(X\right)\widehat{\psi }_\alpha \left(X\right)}`$
$`+{\displaystyle \frac{C_2}{2}}{\displaystyle \underset{\alpha ,\beta ,\alpha ^{^{}},\beta ^{^{}}}{}}{\displaystyle d^3X\widehat{\psi }_\alpha ^{}\left(X\right)\widehat{\psi }_\beta ^{}\left(X\right)F_{\alpha \alpha ^{^{}}}F_{\beta \beta ^{^{}}}\widehat{\psi }_\beta ^{^{}}\left(X\right)\widehat{\psi }_\alpha ^{^{}}\left(X\right)},`$
where $`M`$ is the mass of a single atom; $`\widehat{\psi }_\alpha \left(X\right)`$ is the atomic field annihilation operator associated with atoms in the hyperfine spin state $`|f=1,m_f=\alpha `$ and the indices $`\alpha ,\beta ,\alpha ^{^{}},\beta ^{^{}}`$ label the three spin components $`\left(\alpha ,\beta ,\alpha ^{^{}},\beta ^{^{}}=1,0,1\right)`$. $`V_0\left(X\right)=V_0\left(\mathrm{sin}^2kX_1+\mathrm{sin}^2kX_2+\mathrm{sin}^2kX_3\right)`$ is the optical lattice potential formed by laser beams, which is assumed to be the same for all three spin components, where $`k=2\pi /\lambda `$ is the wave vector of the laser light with $`\lambda `$ being the wavelength of the laser light and $`V_0`$ is the tunable depth of the potential well and, hence, the lattice constant is $`d=\lambda /2`$. $`V_T\left(X\right)`$ denotes an additional (slowly varying) external trapping potential, e.g., a magnetic trap. The $`3\times 3`$ spin matrices $`F`$ denote the conventional three-dimensional representation (corresponding to the spin $`f=\mathit{1}`$) of the angular momentum operator with $`F_{\alpha \beta }^x=\left(\delta _{\alpha ,\beta 1}+\delta _{\alpha ,\beta +1}\right)/\sqrt{2},`$ $`F_{\alpha \beta }^y=i\left(\delta _{\alpha ,\beta 1}\delta _{\alpha ,\beta +1}\right)/\sqrt{2},`$ $`F_{\alpha \beta }^z=\alpha \delta _{\alpha \beta }.`$ The coefficients $`C_0`$ and $`C_2`$ are related to scattering lengths $`a_0`$ and $`a_2`$ of two colliding bosons with total angular momenta $`0`$ and $`2`$, respectively, by $`C_0=4\pi \mathrm{}^2\left(2a_2+a_0\right)/3M`$ and $`C_2=4\pi \mathrm{}^2\left(a_2a_0\right)/3M`$. For atoms <sup>23</sup>Na, we have $`a_2>a_0`$, that is $`C_2>0`$ and the interaction is antiferromagnetic. While, for <sup>87</sup>Rb atoms the situation is just opposite that $`a_2<a_0`$ (this leads to $`C_2<0`$) and the interaction is ferromagnetic . For the periodic potential, the energy eigenstates are Bloch states. We can expand the field operator $`\widehat{\psi }_\alpha \left(X\right)`$ in the Wannier basis, which is a superposition of Bloch states such that
$$\widehat{\psi }_\alpha \left(X\right)=\underset{i}{}\widehat{a}_{\alpha i}w\left(XX_i\right),$$
(4)
where $`w\left(XX_i\right)`$ are the Wannier functions localized in the lattice site $`i`$ and $`\widehat{a}_{\alpha i}`$ corresponds to the bosonic annihilation operator on the $`i`$th lattice site. Also Eq. (4) suggests that atoms in different spin states are approximately described by the same coordinate wave function, which is seen to be the case when the spin-symmetric interaction is strong compared with the asymmetric part, i.e., $`\left|C_0\right|\left|C_2\right|`$ . This is relevant to experimental conditions for <sup>23</sup>Na and <sup>87</sup>Rb atoms. Using Eq. (4), the general Hamiltonian (1) reduces to the Bose–Hubbard Hamiltonian
$`\widehat{H}`$ $`=`$ $`J{\displaystyle \underset{i,j}{}}{\displaystyle \underset{\alpha }{}}\widehat{a}_{\alpha i}^{}\widehat{a}_{\alpha j}+{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha }{}}\epsilon _i\widehat{a}_{\alpha i}^{}\widehat{a}_{\alpha i}`$ (7)
$`+{\displaystyle \frac{1}{2}}U_0{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha ,\beta }{}}\widehat{a}_{\alpha i}^{}\widehat{a}_{\beta i}^{}\widehat{a}_{\beta i}\widehat{a}_{\alpha i}`$
$`+{\displaystyle \frac{1}{2}}U_2{\displaystyle \underset{i}{}}{\displaystyle \underset{\alpha ,\beta ,\alpha ^{^{}},\beta ^{^{}}}{}}\widehat{a}_{\alpha i}^{}\widehat{a}_{\beta i}^{}F_{\alpha \alpha ^{^{}}}F_{\beta \beta ^{^{}}}\widehat{a}_{\beta ^{^{}}i}\widehat{a}_{\alpha ^{^{}}i}.`$
Here the first term in Eq. (7) describes the strength of the spin-symmetric tunneling, which is characterized by the hopping matrix element between adjacent sites$`i`$ and $`j.`$ The tunneling constant $`J=d^3Xw^{}\left(XX_i\right)\left[^2/2M+V_0\left(X\right)\right]w\left(XX_j\right)`$ depends exponentially on the depth of potential well $`V_0`$ and can be varied experimentally by several orders of magnitude. The second term denotes energy offset of the$`i`$th lattice site due to the external confinement of the atoms, where the parameter $`\epsilon _i=d^3Xw^{}\left(XX_i\right)V_T\left(X\right)w\left(XX_i\right)`$ is assumed to be of the same value $`\epsilon `$ for all lattice sites in the present paper. The third and fourth terms characterize the repulsion interaction between two atoms on a single lattice site, which is quantified by the on-site interaction matrix element $`U_{0(2)}=C_{0(2)}\left|w\left(XX_i\right)\right|^4d^3X`$. In our case the interaction energy is very well determined by the single parameters $`U_0`$ and $`U_2`$, due to the short range of the interactions, which is smaller than the lattice spacing.
Now we use Bogliubov approach to diagonalize the Hamiltonian and obtain the excitation energy spectrum of spin–1 bosons with weak interaction in an optical lattice and hence to study the SFP property explicitly. To this end we firstly express the site space operator $`\widehat{a}_{\alpha i}`$ in terms of the wave–vector space operator $`\widehat{a}_{k,\alpha }`$ as
$$\{\begin{array}{c}\widehat{a}_{\alpha i}=\frac{1}{\sqrt{N_s}}_k\widehat{a}_{k,\alpha }e^{ikX_i},\\ \widehat{a}_{\alpha i}^{}=\frac{1}{\sqrt{N_s}}_k\widehat{a}_{k,\alpha }^{}e^{ikX_i},\end{array}$$
(8)
where $`N_s`$ is the total number of the lattice sites and $`X_i`$ is the coordinate of site $`i`$. The wave vector $`k`$ runs only over the first Brillouin zone. In the tight–binding approximation (TBA) if we limit our description to simple cubic lattice and substitute Eq. (8) into the Hamiltonian (3), we can obtain
$`\widehat{H}`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{k}{}}\epsilon \left(k\right)\widehat{a}_{k,\alpha }^{}\widehat{a}_{k,\alpha }+\widehat{H}_{int},`$ (9)
$`\widehat{H}_{int}`$ $`=`$ $`{\displaystyle \frac{U_0}{2N_s}}{\displaystyle \underset{\alpha ,\beta }{}}{\displaystyle \underset{k,p,k^{^{}},p^{^{}}}{}}\delta _{k+p,k^{^{}}+p^{^{}}}\widehat{a}_{k,\alpha }^{}\widehat{a}_{p,\beta }^{}\widehat{a}_{p^{^{}},\beta }\widehat{a}_{k^{^{}},\alpha }`$ (11)
$`+{\displaystyle \frac{U_2}{2N_s}}{\displaystyle \underset{\alpha ,\beta ,\alpha ^{^{}},\beta ^{^{}}}{}}{\displaystyle \underset{k,p,k^{^{}},p^{^{}}}{}}\delta _{k+p,k^{^{}}+p^{^{}}}\widehat{a}_{k,\alpha }^{}\widehat{a}_{p,\beta }^{}F_{\alpha \alpha ^{^{}}}F_{\beta \beta ^{^{}}}\widehat{a}_{p^{^{}},\beta ^{^{}}}\widehat{a}_{k^{^{}},\alpha ^{^{}}},`$
where $`\epsilon \left(k\right)=\epsilon Jz\mathrm{cos}\left(kd\right)`$ with $`z`$ the number of nearest neighbors of each site. Since the number of atoms condensed in the zero-momentum state is much larger than one, we have
$$\underset{\alpha }{}\widehat{a}_{\alpha 0}\widehat{a}_{\alpha 0}^{}=\underset{\alpha }{}\widehat{a}_{\alpha 0}^{}\widehat{a}_{\alpha 0}+1\underset{\alpha }{}N_{\alpha 0}=𝒩_01.$$
(12)
$`N_{\alpha 0}`$ is the number of condensed atoms of spin-$`\alpha `$ component in the zero-momentum state and $`𝒩_0`$ is the total number of condensed atoms. So, we can replace the operator $`\widehat{a}_{\alpha 0}`$ and $`\widehat{a}_{\alpha 0}^{}`$ with a $`\mathrm{`}\mathrm{`}`$ $`c`$ $`\mathrm{"}`$ number $`\sqrt{N_{\alpha 0}}`$. Thus we have
$$N_{\alpha 0}=N_\alpha \underset{k0}{}\widehat{a}_{k,\alpha }^{}\widehat{a}_{k,\alpha },$$
(13)
where $`N_\alpha `$ is the total number of atoms of spin-$`\alpha `$ component and the term of $`k0`$ is exclusive in the wave-vector sum. Moreover, in the interacting part of the Hamiltonian when $`k0`$ we regard $`\widehat{a}_{k,\alpha }^{}`$, $`\widehat{a}_{k,\alpha }`$ as the small deviation from the operators of vanishing momentum, thus all products of four boson operators are approximated as quadratic form, for example,
$$\widehat{a}_{0,\alpha }^{}\widehat{a}_{0,\alpha }^{}\widehat{a}_{0,\alpha }\widehat{a}_{0,\alpha }=N_{\alpha 0}^2N_\alpha ^22N_\alpha \underset{k0}{}\widehat{a}_{k,\alpha }^{}\widehat{a}_{k,\alpha },$$
(14)
$$\underset{\alpha ,\beta ,k0}{}\widehat{a}_{0,\alpha }^{}\widehat{a}_{k,\beta }^{}\widehat{a}_{k,\beta }\widehat{a}_{0,\alpha }=\underset{\alpha ,\beta ,k0}{}\widehat{a}_{k,\beta }^{}\widehat{a}_{k,\beta }N_{\alpha 0}=\left[\underset{\alpha }{}\left(N_\alpha \underset{k0}{}\widehat{a}_{k,\alpha }^{}\widehat{a}_{k,\alpha }\right)\right]\underset{\beta ,k0}{}\widehat{a}_{k,\beta }^{}\widehat{a}_{k,\beta }N\underset{\alpha ,k0}{}\widehat{a}_{k,\alpha }^{}\widehat{a}_{k,\alpha },$$
(15)
where $`N=_\alpha N_\alpha `$ is the total number of atoms. With the approximation Eqs. (14) and (15) the total Hamiltonian (5) can be written as
$`\widehat{H}`$ $`=`$ $`{\displaystyle \frac{U_0}{2N_s}}N^2+N(\epsilon zJ)+{\displaystyle \frac{U_2}{2N_s}}[(N_1N_1)^2+2N_0(\sqrt{N_1}+\sqrt{N_1})^2]+{\displaystyle \underset{k0}{}}[{\displaystyle \underset{\alpha }{}}\epsilon \left(k\right)\widehat{a}_{k,\alpha }^{}\widehat{a}_{k,\alpha }`$ (19)
$`+{\displaystyle \frac{U_0}{2N_s}}{\displaystyle \underset{\alpha ,\beta }{}}\sqrt{N_{\alpha 0}}\sqrt{N_{\beta 0}}(\widehat{a}_{k,\alpha }^{}\widehat{a}_{k,\beta }^{}+\widehat{a}_{k,\alpha }\widehat{a}_{k,\beta }+2\widehat{a}_{k,\alpha }^{}\widehat{a}_{k,\beta })+{\displaystyle \frac{U_2}{2N_s}}(2N_0(\widehat{a}_{k,1}^{}\widehat{a}_{k,1}^{}+\widehat{a}_{k,1}\widehat{a}_{k,1})+{\displaystyle \underset{\gamma =\pm 1}{}}N_\gamma (\widehat{a}_{k,\gamma }^{}\widehat{a}_{k,\gamma }^{}`$
$`+\widehat{a}_{k,\gamma }\widehat{a}_{k,\gamma }+2\widehat{a}_{k,\gamma }^{}\widehat{a}_{k,\gamma })+2{\displaystyle }_{\gamma =\pm 1}\sqrt{N_{\gamma 0}}\sqrt{N_{00}}(\widehat{a}_{k,\gamma }^{}\widehat{a}_{k,0}^{}+\widehat{a}_{k,\gamma }\widehat{a}_{k,0}+\widehat{a}_{k,\gamma }^{}\widehat{a}_{k,0}+\widehat{a}_{k,0}^{}\widehat{a}_{k,\gamma }+2\widehat{a}_{k,0}^{}\widehat{a}_{k,\gamma }+2\widehat{a}_{k,\gamma }^{}\widehat{a}_{k,0})`$
$`+2\sqrt{N_{10}}\sqrt{N_{10}}(\widehat{a}_{k,0}^{}\widehat{a}_{k,0}^{}\widehat{a}_{k,1}^{}\widehat{a}_{k,1}^{}+\widehat{a}_{k,0}\widehat{a}_{k,0}\widehat{a}_{k,1}\widehat{a}_{k,1}\widehat{a}_{k,1}^{}\widehat{a}_{k,1}\widehat{a}_{k,1}^{}\widehat{a}_{k,1}2\widehat{a}_{k,0}^{}\widehat{a}_{k,0}))].`$
Note that the ratio $`U_2/U_0`$ is proportional to the ratio of $`C_2/C_0`$ for all lattice geometries and hence $`U_2/U_0`$ is small enough in general. Therefore the spin-asymmetric part of the interaction is much smaller than the spin-symmetric one.
The Hamiltonian (10) is quadratic in the operators $`\widehat{a}_{k,a}`$, $`\widehat{a}_{k,a}^{}`$ and can be diagonalized by the linear transformation
$`\widehat{a}_{k,a}`$ $`=`$ $`u_{k,\alpha }\widehat{b}_{k,\alpha }\upsilon _{k,\alpha }\widehat{b}_{k,\alpha }^{},`$ (20)
$`\widehat{a}_{k,a}^+`$ $`=`$ $`u_{k,\alpha }\widehat{b}_{k,\alpha }^{}\upsilon _{k,\alpha }\widehat{b}_{k,\alpha },`$ (21)
known as the Bogoliubov transformation. This transformation introduces a new set of operators $`\widehat{b}_{k,\alpha }`$ and $`\widehat{b}_{k,\alpha }^{}`$ to which we impose the same Bose-operator commutation relations $`[\widehat{b}_{k,\alpha },\widehat{b}_{k^{^{}},\beta }^{}]=\delta _{k,k^{^{}}}\delta _{\alpha \beta }`$ . It is easy to check that the commutation relations are fulfilled if the parameters $`u_{k,\alpha }`$ and $`\upsilon _{k,\alpha }`$ satisfy the relation
$$u_{k,\alpha }^2\upsilon _{k,\alpha }^2=1,$$
(22)
where the auxiliary parameters $`u_{k,\alpha }`$ and $`\upsilon _{k,\alpha }`$ are to be chosen in order to have the vanishing coefficients of the nondiagonal terms $`\widehat{b}_{k,\alpha }^{}\widehat{b}_{k,\alpha }^{}`$ and $`\widehat{b}_{k,\alpha }\widehat{b}_{k,\alpha }`$ in the Hamiltonian (10) (see the Appendix for detail). In virtue of the Bogliubov transformation (11), we finally obtain the diagonalized Hamiltonian as
$$\widehat{H}=E_c+\underset{k0}{}E_{k,\alpha ,\alpha }\widehat{b}_{k,\alpha }^{}\widehat{b}_{k,\alpha }$$
(23)
with
$`E_c={\displaystyle \frac{1}{2}}U_0Nn^2+N\left(\epsilon zJ\right)+{\displaystyle \frac{U_2}{2N_s}}[\left(N_1N_1\right)^2+2N_0(\sqrt{N_1}+\sqrt{N_1})^2],`$
where the energy spectra $`E_{k,\alpha ,\alpha }`$ $`\left(\alpha =0,\pm 1\right)`$ of the quasiparticle are given by
$$E_{k,\gamma ,\gamma }=\sqrt{\overline{\epsilon }_k\left(\overline{\epsilon }_k+2U_0n_\gamma +2U_2n_\gamma \right)},$$
(24)
$$E_{k,0,0}=\sqrt{\overline{\epsilon }_k\left(\overline{\epsilon }_k+2U_0n_0+4U_2\sqrt{n_1}\sqrt{n_1}\right)},$$
(25)
where $`\gamma =\pm 1`$ and
$`\overline{\epsilon }_{k\left(k0\right)}=zJ\left[1\mathrm{cos}\left(kd\right)\right].`$
The symbol $`n_\alpha =N_\alpha /N_s`$ represents the average atom number of spin-$`\alpha `$ component per lattice site. In the experiments of Ref. the number of atoms per lattice site is shown to be around $`13`$.
## III CRITICAL VELOCITY OF SUPERFLUID
The energy spectra Eqs. (24) and (25) are typical for the superfluid. To this end we look at the dispersion relations of energy spectra $`E_{k,\alpha ,\alpha }`$ for the limit case $`k0`$
$$E_{k,\gamma ,\gamma }\left[zJd^2\left(U_0+U_2\right)n_\gamma \right]^{1/2}k,$$
(26)
$$E_{k,0,0}\left[zJd^2\left(U_0n_0+2U_2\sqrt{n_1}\sqrt{n_1}\right)\right]^{1/2}k.$$
(27)
The linear wave–vector dependence of the excitation spectra $`E_{k,\alpha ,\alpha }`$ is the characteristic of the superfluid which gives rise to critical velocities of superfluid found as
$$\upsilon _{s,\gamma }=\left(\frac{E_{k,\gamma ,\gamma }}{k}\right)_{k0}=\frac{1}{\mathrm{}}\left[zJd^2\left(U_0+U_2\right)n_\gamma \right]^{1/2},$$
(28)
$$\upsilon _{s,0}=\left(\frac{E_{k,0,0}}{k}\right)_{k0}=\frac{1}{\mathrm{}}\left[zJd^2\left(U_0n_0+2U_2\sqrt{n_1}\sqrt{n_1}\right)\right]^{1/2},$$
(29)
which reduce to the critical velocity of superfluid given in Ref. for the spin-zero case when $`U_2=0`$, where $`1/\mathrm{}`$ is dimension correction. The nonvanishing velocity is nothing but the Landau criterion for the superfluid phase. As seen from the above formulas (18) and (19), whether there exist critical velocities of superfluid or not depend on appropriate values of $`J`$ and $`U_{0(2)}`$ which are related to the Wannier functions determined essentially by the potential of optical lattice. Thus, $`J`$ and $`U_{0(2)}`$ can be controlled dependently by adjusting the laser parameters. Since the spin-asymmetric interaction $`U_2`$ is typically one to two orders of magnitude less than the spin-symmetric interaction $`U_0`$, it can ensure that the nonvanishing $`\upsilon _{s,\alpha }`$ ($`\alpha =0,\pm 1`$) exist, whether for the antiferromagnetic interaction $`\left(U_2>0\right)`$ or for the case of ferromagnetic interaction $`\left(U_2<0\right)`$.
Certainly, it is important to see whether or not the superfluid phase can be realized practically with the recent progress of experiments on the confinement of atoms in the light-induced trap. To see this we evaluate the values of critical velocities of superfluid adopting the typical experimental data in Ref. for a spin-1 condensate of <sup>23</sup>Na atoms in the optical lattice created by three perpendicular standing laser beams with $`\lambda =985`$ nm. The scattering lengths for <sup>23</sup>Na atoms are $`a_0=(46\pm 5)a_B`$ and $`a_2=(52\pm 5)a_B`$, where $`a_B`$ is the Bohr radius (corresponding to a ratio value $`U_2/U_0=0.04`$). The valid condition of the Bogliubov approach that $`U<<J`$ (for example, $`U/J<<0.1`$) can be fulfilled in a region of barrier-height values of the optical lattice potential $`V_0`$ from $`0`$ to two or three times $`E_R,`$ where $`E_R`$ is the recoil energyIt turns out that the magnitude of the critical velocities of the superfluid is the order of mm/s which is seen to be in the range of experimental values .
The critical velocities of superfluid $`\upsilon _{s,\alpha }`$ different for three spin components are functions of densities. Therefore, the critical velocities of superfluid can be detected experimentally by counting the atom–number populations. Moreover, the component-dependent velocities may imply component separation of spinor BEC in an optical lattice similar to the experimentally observed component separation in a binary mixture of BECs since one can control the subsequent time evolution of the mixture of a three-condensate system and detect the relative motions of the three components which tend to preserve the density profiles, respectively. In particular, we may consider a condensate of spin–polarized atoms which are all in the state of spin component $`\alpha =0`$ at initial time $`t=0`$, i.e., $`|\psi \left(0\right)=|0,N_0,0`$. In this case a pair of atoms in the $`\alpha =0`$ state can be excited into the $`\alpha =\pm 1`$ states respectively. Thus after a time $`t_c`$ the number of atoms of the $`\alpha =0`$ component becomes steady such that $`N_0\left(t_c\right)=N_0/2`$ and the number of atoms of $`\alpha =\pm 1`$ components is about half of $`N_0\left(t_c\right)`$, i.e., $`N_1\left(t_c\right)=N_1\left(t_c\right)=N_0\left(t_c\right)/2N_0/4`$. Consequently, the critical velocities of superfluid for $`\alpha =1`$ $`(\upsilon _{s,1})`$ and $`\alpha =1`$ $`\left(\upsilon _{s,1}\right)`$ are equal and different from that for the $`\alpha =0`$ component $`\left(\upsilon _{s,0}\right)`$. There exists a simple relation that $`\upsilon _{s,1}=\upsilon _{s,1}=\upsilon _{s,0}/\sqrt{2}`$ for the case considered and therefore the atoms of the $`\alpha =0`$ component can be separated from the mixture of BECs.
## IV CONCLUSION
The energy–band structure of excitation spectra derived in terms of Bogliubov transformation for spin–1 cold bosons in an optical lattice is shown to be typical for the superfluid phase from the viewpoint of the Landau criterion. Our observation is that the critical velocities of the superfluid flow are spin-component dependent and can be controlled by adjusting the laser lights that form the optical lattice. The theoretical values of critical velocities obtained are in agreement with experimental observations and possible experiments to detect the superfluid phase are also discussed.
## V Acknowledgment
This work was supported by National Natural Science Foundation of China under Grant Nos. 10475053.
## VI APPENDIX
The inverse transformation of Eq. (21) is
$$\{\begin{array}{c}\widehat{b}_{k,\alpha }=u_{k,\alpha }\widehat{a}_{k,a}+\upsilon _{k,\alpha }\widehat{a}_{k,a}^{},\\ \widehat{b}_{k,\alpha }^+=u_{k,\alpha }\widehat{a}_{k,a}^{}+\upsilon _{k,\alpha }\widehat{a}_{k,a}.\end{array}$$
(30)
The Hamiltonian Eq. (19) can be written in terms of the quasiboson operators $`\widehat{b}_{k,\alpha }`$ and $`\widehat{b}_{k,\alpha }^{}`$ as
$$\widehat{H}=E_c+\widehat{H}_1+\widehat{H}_2,$$
(30)
where
$$E_c=\frac{U_0}{2N_s}N^2+N\left(\epsilon zJ\right)+\frac{U_2}{2N_s}[\left(N_1N_1\right)^2+2N_0(\sqrt{N_1}+\sqrt{N_1})^2]$$
(31)
is a constant. $`\widehat{H}_1`$ and $`\widehat{H}_2`$ denote the diagonal and off-diagonal parts, respectively, with
$`\widehat{H}_1`$ $`=`$ $`{\displaystyle \underset{k0}{}}([(\overline{\epsilon }_k+{\displaystyle \frac{U_0}{N_s}}N_02{\displaystyle \frac{U_2}{N_s}}\sqrt{N_{10}}\sqrt{N_{10}})(u_{k,0}^2+\upsilon _{k,0}^2)({\displaystyle \frac{2U_0}{N_s}}N_0+{\displaystyle \frac{4U_2}{N_s}}\sqrt{N_{10}}\sqrt{N_{10}})u_{k,0}\upsilon _{k,0}]\widehat{b}_{k,0}^{}\widehat{b}_{k,0}`$ ()
$`+{\displaystyle \underset{\gamma }{}}[(\overline{\epsilon }_k+{\displaystyle \frac{U_0}{N_s}}N_\gamma +{\displaystyle \frac{U_2}{N_s}}N_\gamma )(u_{k,\gamma }^2+\upsilon _{k,\gamma }^2)({\displaystyle \frac{2U_0}{N_s}}+{\displaystyle \frac{2U_2}{N_s}})N_\gamma u_{k,\gamma }\upsilon _{k,\gamma }]\widehat{b}_{k,\gamma }^{}\widehat{b}_{k,\gamma }+{\displaystyle \underset{\gamma }{}}\{[({\displaystyle \frac{U_0}{N_s}}`$
$`+{\displaystyle \frac{U_2}{N_s}})\sqrt{N_{\gamma 0}}\sqrt{N_{\gamma 0}}{\displaystyle \frac{U_2}{N_s}}N_0](u_{k,\gamma }\upsilon _{k,\gamma }+u_{k,\gamma }\upsilon _{k,\gamma })+({\displaystyle \frac{U_0}{N_s}}{\displaystyle \frac{U_2}{N_s}})\sqrt{N_{\gamma 0}}\sqrt{N_{\gamma 0}}(u_{k,\gamma }u_{k,\gamma }+\upsilon _{k,\gamma }\upsilon _{k,\gamma })\}\widehat{b}_{k,\gamma }^{}\widehat{b}_{k,\gamma }`$
$`+{\displaystyle \underset{\gamma }{}}\{[({\displaystyle \frac{U_0}{N_s}}+{\displaystyle \frac{U_2}{N_s}})\sqrt{N_{\gamma 0}}\sqrt{N_{00}}+{\displaystyle \frac{2U_2}{N_s}}\sqrt{N_{\gamma 0}}\sqrt{N_{00}}](u_{k,\gamma }u_{k,0}+\upsilon _{k,0}\upsilon _{k,\gamma })({\displaystyle \frac{U_0}{N_s}}+{\displaystyle \frac{U_2}{N_s}})\sqrt{N_{\gamma 0}}\sqrt{N_{00}}(u_{k,\gamma }\upsilon _{k,0}`$
$`+u_{k,0}\upsilon _{k,\gamma })\}(\widehat{b}_{k,\gamma }^{}\widehat{b}_{k,0}+\widehat{b}_{k,0}^{}\widehat{b}_{k,\gamma }))+const,`$
$`\widehat{H}_2`$ $`=`$ $`{\displaystyle \underset{k0}{}}([({\displaystyle \frac{U_0}{2N_s}}N_0+{\displaystyle \frac{U_2}{N_s}}\sqrt{N_{10}}\sqrt{N_{10}})(u_{k,0}^2+\upsilon _{k,0}^2)(\overline{\epsilon }_k+{\displaystyle \frac{U_0}{N_s}}N_0+2{\displaystyle \frac{U_2}{N_s}}\sqrt{N_{10}}\sqrt{N_{10}})u_{k,0}\upsilon _{k,0}](\widehat{b}_{k,0}^{}\widehat{b}_{k,0}^{}+\widehat{b}_{k,0}\widehat{b}_{k,0})`$ ()
$`+\left\{\left[\left({\displaystyle \frac{U_0}{2N_s}}{\displaystyle \frac{U_2}{N_s}}\right)\sqrt{N_{10}}\sqrt{N_{10}}+{\displaystyle \frac{U_2}{N_s}}N_0\right](u_{k,1}u_{k,1}+\upsilon _{k,1}\upsilon _{k,1})\left({\displaystyle \frac{U_0}{N_s}}{\displaystyle \frac{U_2}{N_s}}\right)\sqrt{N_{10}}\sqrt{N_{10}}u_{k,1}\upsilon _{k,1}\right\}(\widehat{b}_{k,1}^{}\widehat{b}_{k,1}^{}+\widehat{b}_{k,1}\widehat{b}_{k,1})`$
$`+\left[{\displaystyle \frac{U_0}{2N_s}}\sqrt{N_{10}}\sqrt{N_{10}}(u_{k,1}u_{k,1}+\upsilon _{k,1}\upsilon _{k,1})\left({\displaystyle \frac{U_0}{N_s}}{\displaystyle \frac{U_2}{N_s}}\right)\sqrt{N_{10}}\sqrt{N_{10}}u_{k,1}\upsilon _{k,1}\right](\widehat{b}_{k,1}^{}\widehat{b}_{k,1}^{}+\widehat{b}_{k,1}\widehat{b}_{k,1})`$
$`+{\displaystyle \underset{\gamma }{}}\left\{\left({\displaystyle \frac{U_0}{2N_s}}+{\displaystyle \frac{U_2}{2N_s}}\right)N_\gamma \left(u_{k,\gamma }^2+\upsilon _{k,\gamma }^2\right)\left[\overline{\epsilon }_k+\left({\displaystyle \frac{U_0}{N_s}}+{\displaystyle \frac{U_2}{N_s}}\right)N_\gamma \right]u_{k,\gamma }\upsilon _{k,\gamma }\right\}(\widehat{b}_{k,\gamma }^{}\widehat{b}_{k,\gamma }^{}+\widehat{b}_{k,\gamma }\widehat{b}_{k,\gamma })`$
$`+{\displaystyle \underset{\gamma }{}}\{({\displaystyle \frac{U_0}{2N_s}}+{\displaystyle \frac{U_2}{N_s}})\sqrt{N_{\gamma 0}}\sqrt{N_{00}}(u_{k,\gamma }u_{k,0}+\upsilon _{k,\gamma }\upsilon _{k,0})[({\displaystyle \frac{U_0}{N_s}}+{\displaystyle \frac{U_2}{N_s}})\sqrt{N_{\gamma 0}}\sqrt{N_{00}}+2{\displaystyle \frac{U_2}{N_s}}\sqrt{N_{\gamma 0}}\sqrt{N_{00}}]u_{k,\gamma }\upsilon _{k,0}\}(\widehat{b}_{k,\gamma }^{}\widehat{b}_{k,0}^{}`$
$`+\widehat{b}_{k,\gamma }\widehat{b}_{k,0})+{\displaystyle }_\gamma \{{\displaystyle \frac{U_0}{2N_s}}\sqrt{N_{\gamma 0}}\sqrt{N_{00}}(u_{k,0}u_{k,\gamma }+\upsilon _{k,0}\upsilon _{k,\gamma })`$
$`[({\displaystyle \frac{U_0}{N_s}}+{\displaystyle \frac{U_2}{N_s}})\sqrt{N_{\gamma 0}}\sqrt{N_{00}}+2{\displaystyle \frac{U_2}{N_s}}\sqrt{N_{\gamma 0}}\sqrt{N_{00}}]u_{k,0}\upsilon _{k,\gamma }\}(\widehat{b}_{k,0}^{}\widehat{b}_{k,\gamma }^{}+\widehat{b}_{k,0}\widehat{b}_{k,\gamma })).`$
In order to eliminate the off-diagonal part $`\widehat{H}_2`$ we require that the coefficients of all terms $`\widehat{b}_{k,\alpha }^{}\widehat{b}_{k,\alpha ^{^{}}}^{}`$ and $`\widehat{b}_{k,\alpha }\widehat{b}_{k,\alpha ^{^{}}}`$ vanish. In view of condition (12), it is easy to introduce a set of parameters $`\varphi _{k,\alpha }`$ such that
$`u_{k,\alpha }`$ $`=`$ $`\mathrm{cosh}\varphi _{k,\alpha },`$ (42)
$`\upsilon _{k,\alpha }`$ $`=`$ $`\mathrm{sinh}\varphi _{k,\alpha }`$ ()
for the convenience of calculation. Conditions (12) and (A6) lead to the useful relations
$`\mathrm{tanh}2\varphi _{k,\alpha }={\displaystyle \frac{2u_{k,\alpha }\upsilon _{k,\alpha }}{u_{k,\alpha }^2+\upsilon _{k,\alpha }^2}},`$
$`\mathrm{cosh}\left(2\varphi _{k,\alpha }\right)=u_{k,\alpha }^2+\upsilon _{k,\alpha }^2={\displaystyle \frac{1}{\sqrt{1\mathrm{tanh}^22\varphi _{k,\alpha }}}},`$
with which the Hamiltonian (10) can be finally reduced to the diagonal form as given in Eq. (23).
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# The Fubini-Furlan-Rosetti sum rule and related aspects in light of covariant baryon chiral perturbation theoryWork supported in part by the DFG through funds provided to the TR 16 “Subnuclear structure of matter”.
## 1 Introduction
The Fubini-Furlan-Rosetti (FFR) sum rule was derived in the sixties utilizing the soft-pion techniques of current algebra FFR . It relates the nucleon anomalous magnetic moment to an integral over the invariant amplitude $`A_1`$ of pion photoproduction
$$\kappa ^{v,s}=\frac{8m_N^2}{e\pi g_{\pi N}}\frac{d\nu ^{}}{\nu }\mathrm{Im}A_1^{(+,0)}(\nu ^{},t=0),$$
(1)
if one utilizes the Goldberger-Treiman relation $`g_Am_N=F_\pi g_{\pi N}`$, with $`m_N`$ the nucleon mass, $`g_A`$ the axial-vector coupling constant, $`F_\pi `$ the weak pion decay constant and $`g_{\pi N}`$ the strong pion–nucleon coupling constant. Furthermore, $`\kappa ^v=\kappa _p\kappa _n`$ and $`\kappa ^s=\kappa _p+\kappa _n`$ are the nucleon isovector and isoscalar anomalous magnetic moment, respectively. The FFR sum rule is exact in the chiral limit of QCD and thus all quantities appearing in Eq. (1) are to be understood in the limit of vanishing light quark masses, $`m_q=0`$. The FFR sum rule has recently been reexamined in Ref. PDT . In that paper<sup>1</sup><sup>1</sup>1See also that paper for references to earlier work on the FFR sum rule., pion mass corrections to the sum rule were considered (in terms of a discrepancy function $`\mathrm{\Delta }_N(\nu ,t)`$) and numerical evaluations based on a) dispersion relations and b) input from heavy baryon chiral perturbation theory (HBCHPT) were presented. It was pointed out that in the strict framework of HBCHPT the nucleon pole positions are slightly moved, which leads e.g. to an incorrect curvature of the discrepancy function for energies below the threshold. A similar behavior due to the shift of pole or cut positions in the $`1/m_N`$ expansion was observed already in the discussion of the Compton cusp at the opening of the pion threshold BKMS , the spectral functions of the nucleon isovector form factors BKMspec or the scalar form factor of the nucleon BL . Note, however, that the kinematical factors leading to the corresponding poles or cuts need not be expanded in HBCHPT as it is discussed e.g. in Ref. BKMS . Clearly, in a manifestly Lorentz-invariant formulation of baryon CHPT such problems do not arise, see e.g. BL ; KM ; MZrel . The purpose of this paper is two-fold: We analyze the FFR sum rule in the framework of infrared regularization (IR) of baryon CHPT BL and demonstrate that the energy dependence of the discrepancy function is correctly given. Second, we also take a closer look at the pion mass corrections to the sum rule, which can only be systematically calculated in chiral perturbation theory, and the related threshold multipoles in pion photoproduction. Our calculation includes all terms at third order in the chiral expansion and in addition, also the fourth order polynomial terms in the photoproduction amplitudes. Our aim is to show that within IR baryon CHPT one can describe pion photoproduction above and below threshold and that the dispersive representation can indeed be used to pin down certain low–energy constants (LECs), as suggested in PDT . It is well-known that the chiral expansion converges best in the unphysical region where all momenta can be very small. Consequently, in such regions LECs can be determined to a good precision if a corresponding dispersive representation is available. A more refined treatment including all fourth order terms and fits to the existing low–energy data from MAMI will be relegated to a future publication.
The manuscript is organized as follows: In Sec. 2 we briefly recall the formalism of pion photoproduction and collect some results for the pion mass corrections to the FFR sum rule derived in PDT . We also present an alternative way of looking at these. Sec. 3 contains the results on the FFR sum rule, the discrepancy function and the related electric dipole amplitude $`E_{0+}`$ as well as the slopes of the P-wave multipoles at threshold. We demonstrate that one can indeed determine LECs from the amplitudes in the unphysical region and end with a brief outlook.
## 2 Formalism
Consider pion photoproduction off the nucleon by real photons with $`k^2=0`$
$$\gamma (k)+N(p_1)\pi ^a(q)+N(p_2),$$
(2)
where $`p_1(p_2)`$ is the four–momentum of the incoming (outgoing) nucleon ($`N`$) and $`a`$ an isospin index ($`a=+,0,`$). For the discussion of the FFR sum rule, only the isospin $`0`$ and $`+`$ channels are of relevance, i.e. the physical channels $`\gamma p\pi ^0p`$ and $`\gamma n\pi ^0n`$. The corresponding S-matrix is given in terms of four invariant functions $`A_i`$ ($`i=1,\mathrm{},4`$) that depend on two kinematical variables. Throughout, we utilize the notation of our earlier work Relphoto2 and refer to that reference for a detailed discussion of the pertinent formalism. These invariant amplitudes can be calculated in baryon chiral perturbation theory and have the generic form (we do not display isospin quantum numbers and kinematical arguments)
$$A_i=A_i^{\mathrm{Born}}+A_i^{\mathrm{loop}}+A_i^{\mathrm{ct}},$$
(3)
where the Born terms subsume the coupling to the charge and the magnetic moment of the nucleon (note often the alternative notation of “pole terms” is used for these contributions – sometimes even calculated employing the pseudoscalar pion-nucleon coupling). All further counter terms are collected in the polynomial terms $`A_i^{\mathrm{ct}}`$. The non-trivial loop contributions (after renormalization of the single nucleon properties) are collected in the $`A_i^{\mathrm{loop}}`$. In Ref. Relphoto2 , the $`A_i`$ were calculated to third order (leading loop order) in relativistic baryon CHPT. In that formulation, a violation of the power counting through the nucleon mass term is manifest. However, from the integral representations given in Relphoto2 , it is straightforward to isolate the so-called infrared singular part BL that contains the chiral long–distance physics and leads to a one-to-one correspondence between the expansion in loops and small momenta/pion masses. This can e.g. be achieved by the prescription given in BL which we also will employ. Symbolically, it reads $`_0^1𝑑x(_0^{\mathrm{}}_1^{\mathrm{}})dx=I+R`$, with $`I`$ and $`R`$ the infrared singular (irregular) and the regular part, respectively. For the study of the FFR sum rule and related aspects, we are interested in the amplitudes at small energies and momentum transfer and thus use the variables $`\nu =(su)/4m_N`$ and $`\nu _B=(s+u2m_N^2)/4m_N`$, which are odd and even under crossing $`su`$ (for further notation, see Relphoto2 ). From the crossing properties of the $`A_i`$ and their low-energy properties as detailed in Relphoto2 , one derives the following representation for the polynomial pieces (note in particular the low-energy theorem for $`A_1`$ Relphoto1 that forbids a constant term)
$`A_1^{\mathrm{ct}}`$ $`=`$ $`a_1^1\nu ^2+a_1^2\nu _B+\mathrm{},`$
$`A_2^{\mathrm{ct}}`$ $`=`$ $`a_2^0+\mathrm{},`$
$`A_3^{\mathrm{ct}}`$ $`=`$ $`a_3^1\nu +\mathrm{},`$
$`A_4^{\mathrm{ct}}`$ $`=`$ $`a_4^0+\mathrm{}.`$ (4)
At third order in the chiral expansion, only the leading term of $`A_4`$ contributes Relphoto2 , whereas all other terms written down in Eq. (2) start at $`𝒪(q^4)`$. The mapping between these subtraction constants and the low-energy constants (LECs) used in the heavy baryon calculations HBphoto1 ; HBphoto2 ; HBphoto3 is given in the appendix. We remind the reader that although the contribution of $`A_1^{\mathrm{ct}}`$ formally starts as $`M_\pi ^2`$ at threshold, in the chiral limit the expansion coefficients are singular leading to the famous contribution to the low-energy theorem for the threshold value of $`E_{0+}`$ at next-to-leading order in the pion mass expansion Relphoto1 . We note that $`A_2^{\mathrm{ct}}`$ only feeds into the P-wave slope $`\overline{P}_2`$ and into the D-wave at threshold in such a way that it cancels in the FFR sum rule at finite pion mass. Therefore, the subtraction constant $`a_2^0`$ has to be determined completely independently of the FFR sum rule, say by fitting to the slope of $`P_2`$ at threshold. In the following, we use the third order IR representation of the pertinent one-loop graphs (which contains an infinite series of $`1/m_N`$ corrections in the heavy baryon framework) but also use the fourth-order subtraction constants displayed in Eq. (2). Based on that representation, we attempt a simultaneous description of the $`\nu `$-dependence of the FFR discrepancy function $`\mathrm{\Delta }_p(\nu ,t_{\mathrm{thr}})`$ (as defined below), the energy dependence of the electric dipole amplitude $`E_{0+}`$ for neutral pion production off protons and the P–wave threshold slopes as extracted in E0pdata . Similarly precise information is not available for the neutron, therefore in the following we mostly concentrate on the proton.
In Ref. PDT , the FFR sum rule was considered for finite mass pions and the following representation in terms of a discrepancy function $`\mathrm{\Delta }_N`$ was derived:
$`\underset{N}{\overset{}{\kappa }}\tau _3+\mathrm{\Delta }_N(\nu ,t_{\mathrm{thr}})`$
$`={\displaystyle \frac{4m_N^2}{e\pi g_{\pi N}}}{\displaystyle _{\nu _{\mathrm{thr}}}^{\mathrm{}}}{\displaystyle \frac{d\nu ^{}}{\nu }}{\displaystyle \frac{\nu ^{}\mathrm{Im}A_1^{(N,\pi ^0)}(\nu ^{},t=t_{\mathrm{thr}})}{v_{}^{}{}_{}{}^{2}\nu ^2}},`$
$`\mathrm{\Delta }_N(\nu ,t=t_{\mathrm{thr}})`$
$`={\displaystyle \frac{2m_N^2}{eg_{\pi N}}}\left(A_1^{\mathrm{loop}}(\nu ,t=t_{\mathrm{thr}})+A_1^{\mathrm{ct}}(\nu ,t=t_{\mathrm{thr}})\right).`$ (5)
A few comments on this equation are in order. First, the left-hand-side of the FFR sum rule gives the anomalous magnetic moment in the chiral limit, $`\kappa _N=\underset{N}{\overset{}{\kappa }}+𝒪(m_q^{1/2})`$. Therefore, if one uses the physical value of $`\kappa _N`$ as in Ref. PDT , one must include the corresponding loop and counter term corrections into the discrepancy function. However, for studying the pion mass corrections to the FFR sum rule, it is more appropriate to work with the chiral limit values of $`\kappa _p`$ and $`\kappa _n`$, as discussed below. Second, in the chiral limit, $`t0`$. This value can, of course, not be achieved in the physical world. We follow PDT and present our results at the minimal (threshold) value of $`t`$, $`t_{\mathrm{thr}}=M_\pi ^2/(1+M_\pi /m)=0.016`$GeV<sup>-2</sup>. Third, the Goldberger-Treiman relation is no longer exact at finite pion mass, to the order we are working, it takes the form GSS ; FMS
$$\frac{g_{\pi N}}{m_N}=\frac{g_A}{F_\pi }\left(1\frac{2M_\pi ^2}{g_A}\overline{d}_{18}\right),$$
(6)
with $`\overline{d}_{18}`$ a LEC. As long as one only works at the physical pion mass, this effect is taken care of by utilizing the physical value for the pion–nucleon coupling and the nucleon mass. If one, however, also wants to study the pion mass dependence of $`\mathrm{\Delta }_N`$, as will be done here, one explicitly has to include this pion mass dependence. However, this effect only shows up in the terms cubic in the pion mass, which will not be considered in detail here (for a more detailed discussion of this topic, see e.g. EGM ). Note this pion mass dependence can only be systematically calculated in chiral perturbation theory and, eventually, in lattice QCD.
We have explicitly worked out the quark mass expansion of the discrepancy function at threshold. For the proton, it takes the form (modulo chiral logs)
$$\mathrm{\Delta }_p(\nu =\nu _{\mathrm{thr}},t=t_{\mathrm{thr}})=\alpha _pM_\pi +\beta _pM_\pi ^2+\mathrm{}$$
(7)
with
$`\alpha _p`$ $`=`$ $`{\displaystyle \frac{m_N}{16F_\pi ^2}}\left(1{\displaystyle \frac{4g_A^2}{3\pi }}\right),`$
$`\beta _p`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2F_\pi ^2}}[3{\displaystyle \frac{\pi ^2}{8}}+(7{\displaystyle \frac{1}{2}}\underset{p}{\overset{}{\kappa }}+{\displaystyle \frac{1}{2}}\underset{n}{\overset{}{\kappa }}\left)\mathrm{ln}{\displaystyle \frac{M_\pi }{m_N}}\right]`$ (8)
$``$ $`{\displaystyle \frac{g_A^2}{48\pi ^2F_\pi ^2}}[{\displaystyle \frac{85}{3}}\pi +{\displaystyle \frac{5}{3}}\underset{p}{\overset{}{\kappa }}{\displaystyle \frac{11}{3}}\underset{n}{\overset{}{\kappa }}`$
$`+(5+7\underset{p}{\overset{}{\kappa }}\underset{n}{\overset{}{\kappa }})\mathrm{ln}{\displaystyle \frac{M_\pi }{m_N}}]`$
$`+`$ $`{\displaystyle \frac{\stackrel{~}{c}_4}{6\pi ^2F_\pi ^2}}\left({\displaystyle \frac{10}{3}}+{\displaystyle \frac{3\pi ^2}{8}}+{\displaystyle \frac{5}{2}}\mathrm{ln}{\displaystyle \frac{M_\pi }{m_N}}\right)`$
$``$ $`8m_N(2e_{105}+e_{106})+{\displaystyle \frac{2m_NF_\pi }{g_A}}\left(a_1^1{\displaystyle \frac{a_1^2}{2m_N}}\right),`$
where $`\stackrel{~}{c}_4=m_Nc_4`$ and $`c_4=3.4`$GeV<sup>-1</sup> BM . The LECs $`e_{105}`$ and $`e_{106}`$ from the dimension four chiral pion-nucleon Lagrangian contribute to the proton and neutron anomalous magnetic moment at next-to-leading loop order (for a detailed discussion, see KM ). Their values are discussed below. Further, we have set the scale of dimensional regularization equal to the nucleon mass, $`\lambda =m_N`$. Of course, $`\mathrm{\Delta }_p`$ vanishes in the chiral limit. Notice the absence of chiral logs in the terms linear in the pion mass. The representation of $`\mathrm{\Delta }_p`$ given Eq. (7) is exact to fourth order in the chiral expansion since it can be reconstructed from the HBCHPT results obtained in HBphoto1 ; HBphoto2 ; HBphoto3 . In addition, to arrive at these results, we had to include the small contribution from the slope of the D-wave combination $`D=M_{2+}E_{2+}P_2E_2`$, that is $`\overline{D}=D/q^2`$ at threshold. It is obtained from the invariant functions by
$$D(s)=\frac{5}{16}𝑑x(x^42x^2+1)_4(s,x),$$
(9)
with $`_4(s,x)`$ a combination of $`A_{2,3,4}`$, see Relphoto2 . This contribution has not been calculated before. Note further that Eq. (7) includes the terms that renormalize $`\underset{p}{\overset{}{\kappa }}`$ to its physical value, $`\kappa _p`$, since we identify the left-hand-side of the FFR with the anomalous magnetic moment in the chiral limit. The pion mass expansion of $`\mathrm{\Delta }_n(\nu =\nu _{\mathrm{thr}},t=t_{\mathrm{thr}})`$ looks very similar, we find
$`\alpha _n`$ $`=`$ $`{\displaystyle \frac{m_N}{16F_\pi ^2}}\left(1{\displaystyle \frac{4g_A^2}{3\pi }}\right),`$
$`\beta _n`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2F_\pi ^2}}[{\displaystyle \frac{7}{9}}+{\displaystyle \frac{\pi ^2}{8}}+({\displaystyle \frac{13}{3}}+{\displaystyle \frac{1}{2}}\underset{p}{\overset{}{\kappa }}{\displaystyle \frac{1}{2}}\underset{n}{\overset{}{\kappa }})\mathrm{ln}{\displaystyle \frac{M_\pi }{m_N}}]`$ (10)
$``$ $`{\displaystyle \frac{g_A^2}{48\pi ^2F_\pi ^2}}[{\displaystyle \frac{76}{3}}\pi +{\displaystyle \frac{11}{3}}\underset{p}{\overset{}{\kappa }}{\displaystyle \frac{5}{3}}\underset{n}{\overset{}{\kappa }}`$
$`+(1+\underset{p}{\overset{}{\kappa }}7\underset{n}{\overset{}{\kappa }})\mathrm{ln}{\displaystyle \frac{M_\pi }{m_N}}]`$
$``$ $`{\displaystyle \frac{\stackrel{~}{c}_4}{6\pi ^2F_\pi ^2}}\left({\displaystyle \frac{10}{3}}+{\displaystyle \frac{3\pi ^2}{8}}{\displaystyle \frac{5}{2}}\mathrm{ln}{\displaystyle \frac{M_\pi }{m_N}}\right)`$
$`+`$ $`8m_N(2e_{105}e_{106})+{\displaystyle \frac{2m_NF_\pi }{g_A}}\left(a_1^{1,n}{\displaystyle \frac{a_1^{2,n}}{2m_N}}\right),`$
where the subtraction constants refer to the neutron amplitude $`\gamma n\pi ^0n`$ as denoted by the superscript $`n`$.
## 3 Results
First, we must fix our input parameters. We use $`F_\pi =92.4`$MeV, $`M_{\pi ^+}=139.57`$MeV, $`M_{\pi ^0}=134.97`$MeV, $`m_p=938.27`$MeV, $`m_n=939.57`$MeV, $`g_{\pi N}=13.4`$, $`\kappa _p=1.793`$, $`\kappa _n=1.913`$. To the order we are working, the anomalous magnetic moment of the proton and the neutron are given in terms of two LECs from the dimension two chiral Lagrangian commonly denoted $`c_6`$ and $`c_7`$ and loop corrections that start with terms of order $`M_\pi ^3`$ BKMlec . From that one can read off the chiral limit values for the proton and the neutron anomalous magnetic moments,
$$\underset{p}{\overset{}{\kappa }}=2.37,\underset{n}{\overset{}{\kappa }}=2.84,$$
(11)
in nuclear magnetons. For comparison, at third order in HBCHPT, we have $`\underset{p}{\overset{}{\kappa }}=2.85`$, $`\underset{n}{\overset{}{\kappa }}=2.98`$ BKMlec .
The overall best fit to simultaneously describe the proton discrepancy function, the electric dipole amplitude in the threshold region and the three P-wave slopes is obtained with the following polynomial contribution to the invariant functions
$`A_1^{\mathrm{ct}}(\nu ,\nu _B)`$ $`=`$ $`(191.3\nu _B+220\nu ^2)\text{GeV}^2,`$
$`A_2^{\mathrm{ct}}(\nu ,\nu _B)`$ $`=`$ $`54.9\text{GeV}^4,`$
$`A_3^{\mathrm{ct}}(\nu ,\nu _B)`$ $`=`$ $`155\nu \text{GeV}^3,`$
$`A_4^{\mathrm{ct}}(\nu ,\nu _B)`$ $`=`$ $`181.5\text{GeV}^3.`$ (12)
While these coefficients appear large at first glance, in the appendix we show that they match quite nicely the corresponding LECs determined in fits to neutral pion photoproduction differential cross sections and the photon asymmetry. Since the LECs can be understood to a good precision in terms of resonance saturation (excitation of the $`\mathrm{\Delta }(1232)`$ and of vector mesons, see e.g. HBphoto1 ), the numbers appearing in Eq. (3) are indeed of natural size. The resulting threshold values for $`E_{0+}`$ and the $`\overline{P}_i`$ ($`i=1,2,3`$) are
$`E_{0+}`$ $`=`$ $`1.19[1.23\pm 0.08\pm 0.03],`$
$`\overline{P}_1`$ $`=`$ $`9.67[9.46\pm 0.05\pm 0.28],`$
$`\overline{P}_2`$ $`=`$ $`9.6[9.5\pm 0.09\pm 0.28],`$
$`\overline{P}_3`$ $`=`$ $`11.45[11.32\pm 0.11\pm 0.34],`$ (13)
in the conventional units of $`10^3/M_{\pi ^+}`$ and $`10^3/M_{\pi ^+}^2`$, respectively. The experimental numbers in the square brackets are from E0pdata . We note that $`\overline{P}_2`$ is obtained by adjusting the subtraction constant $`a_2^0`$. The corresponding D-wave slope is
$$\overline{D}=0.6610^3/M_{\pi ^+}^3,$$
(14)
to be compared with 0.96 (0.92) from the MAID03 analysis (the dispersive analysis of Ref. HDT ).
Let us look at these results in more detail. For the pion mass correction to the proton FFR sum rule, we find
$$\underset{p}{\overset{}{\kappa }}+\mathrm{\Delta }_p(\nu =\nu _{\mathrm{thr}},t=t_{\mathrm{thr}})=2.18$$
(15)
which should be compared with the left-hand-side of the sum rule, i.e. the proton anomalous magnetic moment in the chiral limit, cf. Eq. (11). Thus, within this framework, the FFR sum rule is fulfilled within 8%, which is of the expected size since pion mass effects are proportional to the small parameter $`\mu =M_\pi /m_N1/7`$. Note again that this way of looking at the FFR differs from what was done in PDT , where the left-hand-side of the FFR was identified with the physical value of the proton anomalous magnetic moment and the discrepancy function defined there thus differs from ours by terms $`\underset{p}{\overset{}{\kappa }}\kappa _p`$.
Next, we display the resulting discrepancy function for the proton as a function of $`\nu `$ at fixed $`t=t_{\mathrm{thr}}`$ in Fig. 1 (for better comparison with Ref. PDT , this discrepancy function is taken with respect to the physical value of the proton anomalous magnetic moment).
At $`\nu =0`$, the discrepancy functions has indeed an extremum (as pointed out in PDT ) and it increases with increasing $`\nu `$. We also find the pronounced cusp effect at $`\nu =\nu _{\mathrm{thr}}`$, as it is expected. In contrast to the prediction based on the MAID model (dashed curved), we do not observe any zero crossing for the values of the subtraction constants given in Eq. (3). Naturally, within our approach we should have a band rather than a line but we relegate a detailed error analysis to a later work, when we will also fit to all data of neutral pion photoproduction in the threshold region. Interestingly, the third order relativistic result from Relphoto2 (dot-dashed line) is very close to the IR result including the fourth order polynomial pieces. We also note that the incorrect behavior at $`\nu 0`$ observed in the HBCHPT calculation is of course not present in our manifestly covariant calculation, as it was suggested already in PDT . Note furthermore that we do not show $`\mathrm{\Delta }_P(\nu ,t_{\mathrm{thr}})`$ for values of $`\nu >155`$MeV because at $`\nu 170`$MeV, a steep rise due to the $`\mathrm{\Delta }(1232)`$ resonance sets in PDT . In the lower panel of Fig. 1 we focus on the threshold (cusp) region. We see that the relativistic and IR predictions are in good agreement with the MAID model and also with the data reconstructed form the multipoles of Ref. E0pdata (for details, see again PDT ).
We briefly discuss the pion mass dependence of $`\mathrm{\Delta }_{p,n}`$; for this purpose, we switch back to the full fourth order representation obtained from earlier heavy baryon results. Consider the proton. We find for the coefficients in Eq. (7)
$`\alpha _p`$ $`=`$ $`2.19\mathrm{GeV}^1,`$
$`\beta _p`$ $`=`$ $`(5.017.3)\mathrm{GeV}^2=22.3\mathrm{GeV}^2,`$ (16)
where we have used the physical values for $`g_A`$ and $`F_\pi `$ and the second term in $`\beta _p`$ is the contribution from the dimension four operators with the LECs $`e_{105}`$ and $`e_{106}`$. These counterterms determine the slope of $`\kappa _v`$ and $`\kappa _s`$ in the soft pion limit. Their values have been determined by using the fourth order formula for the anomalous magnetic moment from KM as chiral extrapolation functions for the lattice QCD data of FFlat . We get $`e_{105}0.45,e_{106}1.4`$ at $`\lambda =m_N`$, which are of natural size. Since these are only very rough fits to the trend of the lattice data at too high pion masses, we refrain from assigning a theoretical uncertainties to these numbers. Note that in case of the coefficient $`\beta _p`$, we have large cancellations between the loop and the counterterm contributions proportional to the LECS $`c_i`$, which are $`21.2`$GeV<sup>-2</sup> and $`16.2`$GeV<sup>-2</sup>, respectively, for $`\lambda =m_N`$. We also stress that the first two terms in the quark mass expansion of $`\mathrm{\Delta }_p`$ at threshold give $`0.13`$, which is to be compared with the full calculation that gives $`0.19`$, cf. Eq. (15) (where the contribution from the term $`2e_{105}+e_{106}`$ is $`0.34`$). For the neutron, we have no determination of the photoproduction counter terms and thus we can only give $`\alpha _n=\alpha _p=2.19`$GeV<sup>-1</sup> and $`\beta _n^{\mathrm{loop}}=18.23`$GeV<sup>-2</sup>. This value is comparable to the one of the proton. The contribution from the operators $`e_{105},e_{106}`$ is somewhat smaller than for the proton since for the neutron they appear with a different relative sign.
Also well described is the energy dependence of the electric dipole amplitude $`E_{0+}(E_\gamma )`$ shown in Fig. 2 (with $`E_\gamma =\nu (tM_\pi ^2)/4m_N`$). This description of the data is as good as the complete fourth-order HBCHPT calculation (see e.g. Fig. 3 in E0pdata ). This is not surprising since the third-order IR calculation generates all fourth-order HB corrections with fixed coefficients (from the expansion of the Dirac propagator) and these are the only fourth order loop contributions since loops with one insertion proportional to the dimension two LECs $`c_i`$ do not contribute. Related to this is the observation that most of the non-trivial energy-dependence is given by the generalized cusp function HBphoto1 ; HBphoto3
$$E_{0+}(E_\gamma )=a+b\sqrt{1\left(\frac{E_\gamma }{E_\gamma ^{\mathrm{thr}}}\right)^2},$$
(17)
where the parameter $`b`$ is proportional to the charge exchange scattering length $`a(\pi ^0p\pi ^+n)`$ (see also the detailed discussion in aron1 ; aron2 ).
As noted before, similarly precise information on the neutron is not available. The predictions for the threshold multipoles are based on resonance saturation, the only experimental information is from the SAL experiment on $`\gamma d\pi ^0d`$ SAL that is consistent with the CHPT prediction $`|E_{0+,\mathrm{thr}}^{\pi ^0n}|>|E_{0+,\mathrm{thr}}^{\pi ^0p}|`$ BBLMvK . We can fit to the energy-dependence of $`A_1^n`$ as predicted by the MAID model and the threshold multipoles as predicted by CHPT, but in the absence of more information on the energy dependence of e.g. the neutron electric dipole amplitude and any experimental verification of the predictions for the P–wave slopes based on resonance saturation, the resulting numbers are highly model-dependent and we thus refrain from showing them here.
## 4 Summary and outlook
In this paper, we have studied the FFR sum rule in the framework of covariant baryon chiral perturbation theory, extending some of the work presented in PDT . We have worked out the loop corrections to third order in the chiral expansion, corresponding to the leading loop contribution, supplemented by the polynomial pieces up-to-and-including fourth order. Since the (sub-)threshold energies considered here are small, such a procedure is justified and would only lead to a readjustment of the subtraction constants defined in Eq. (2) if one includes also the fourth order loop graphs. We have shown that within this framework one can achieve a good description of the energy dependence of the discrepancy function for the proton, defined in Eq. (2), together with the energy dependence of the electric dipole amplitude in the threshold region, cf. Fig. 2, and the P-wave slopes at threshold, see Eq. (3). The corresponding subtraction constants collected in Eq. (3) can be matched to the low-energy constants determined previously in HBCHPT studies of threshold neutral pion photoproduction and their resulting values are of natural size (as detailed in the appendix). We find that the finite pion mass corrections to the FFR sum rule for the proton are small, of the order of 8 % (cf. Eq. (15)). As shown in Fig. 1, the $`\nu `$-dependence of the discrepancy function for the proton has the proper behavior and agrees with the result of the MAID model. The unphysical behavior observed at $`\nu 0`$ in the heavy baryon scheme PDT is absent in a covariant formulation as presented here. It is also interesting to note that the third order relativistic calculation free of counter terms from Relphoto2 also gives a good description of $`\mathrm{\Delta }_p(\nu ,t_{\mathrm{thr}})`$. These findings corroborate the conjecture made in Ref. PDT that one can use the dispersive representation of the invariant amplitudes in the unphysical region to pin down low-energy constants of chiral perturbation theory (see also the appendix). It will be interesting to perform a complete fourth-order calculation and fit the corresponding LECs to the existing unpolarized and polarized threshold data of the reaction $`\gamma p\pi ^0p`$. This should further sharpen the conclusions made here. Work along such lines is under way BKMfull .
## Acknowledgments
We thank Dieter Drechsel for interesting us in this problem. We are grateful to Lothar Tiator for supplying us with the results of Ref. PDT . This research is part of the EU Integrated Infrastructure Initiative Hadron Physics Project under contract number RII3-CT-2004-506078.
## Appendix A Subtraction and low-energy constants
Here, we give the mapping between the commonly employed counter terms of the heavy baryon approach and the subtraction constants defined in Eq.(2). To third order, one has only one P–wave counter term (with the LEC $`b_p`$) that feeds into the multipole $`P_3`$ HBphoto1 ,
$$a_4^0=4\pi b_P.$$
(18)
At fourth order, there are two S–wave counter terms (which in the threshold region essentially act as one constant) HBphoto1 . The corresponding LECs $`a_1`$ and $`a_2`$ are given by the following combinations of subtraction constants:
$`12\pi a_1`$ $`=`$ $`\left(a_2^0+{\displaystyle \frac{a_4^0}{m_N}}\right),`$
$`12\pi a_2`$ $`=`$ $`\left(3a_3^1+a_2^0+3a_1^1{\displaystyle \frac{3}{2}}{\displaystyle \frac{a_1^2}{m_N}}+{\displaystyle \frac{5}{2}}{\displaystyle \frac{a_4^0}{m_N}}\right).`$ (19)
Note that in the sum $`a_1+a_2`$, the contribution from $`A_2`$ cancels (as noted earlier) and that these relations are to be taken at $`\lambda =m_N`$. Similarly, at fourth order there are two independent counter terms (with the LECs $`\xi _1`$ and $`\xi _2`$) that modify the P-waves $`P_1`$ and $`P_2`$ HBphoto3
$`{\displaystyle \frac{a_1^2}{2}}m_Na_3^1a_4^0`$ $`=`$ $`{\displaystyle \frac{g_A}{16\pi ^2F_\pi ^3}}\xi _1,`$
$`{\displaystyle \frac{a_4^0}{2}}+m_Na_3^1m_Na_2^0`$ $`=`$ $`{\displaystyle \frac{g_A}{16\pi ^2F_\pi ^3}}\xi _2.`$ (20)
At first glance, one might conclude from Eqs.(A,A) that there is a mismatch in the number of subtraction constants and counter terms. Note, however, that various subtraction constants feed into the S- and the P-waves so that finally only two independent structures remain for the S-wave and two for the P-waves.
It is interesting to compare the numbers derived from Eq. (3) with the earlier determinations of these counter terms in the HBCHPT framework. Note, however, that we did not include all fourth order loop corrections here, so that the values of the subtraction constants effectively subsume some of these effects. This is not the case for the LEC $`b_P`$ since it already appears at third order. Our value for $`a_0^4`$ translates into $`b_P=14.4`$GeV<sup>-3</sup>. This compares well with the third order fits of Ref. HBphoto1 , $`b_P=(15.8\pm 0.2)`$GeV<sup>-3</sup>, and of Ref. HBphoto2 , $`b_P=13.0`$GeV<sup>-3</sup>. Note that the corresponding value in HBphoto3 comes out smaller due to additional loop effects. Consider next the LECs contributing to the electric dipole amplitude in the threshold region. We get $`4\pi m_N(a_1+a_2)=56.2`$GeV<sup>-3</sup> from the constants in Eq. (3) compared to 31.8 GeV<sup>-3</sup>, 77.8 GeV<sup>-3</sup> and 71.2 GeV<sup>-3</sup> from HBphoto1 , HBphoto2 and HBphoto3 , respectively. Individually, we have $`a_1=3.67`$GeV<sup>-3</sup> and $`a_2=8.43`$GeV<sup>-3</sup>, which is different from but comparable in size to the free and resonance fits to the various sets of Mainz and Saskatoon data (compare e.g. table 1 in HBphoto3 ). As it was already stressed in these earlier papers, the LECs $`a_1`$ and $`a_2`$ can not be well determined individually from fits to the data in the threshold region. Next, consider the P-wave $`P_1`$. The determination of $`\xi _1=16.6`$ in HBphoto3 translates into $`a_1^2/2m_Na_3^1a_4^0=175.7`$GeV<sup>-3</sup>, which is sizably larger than the value of 59.8 obtained from Eq. (3). This is expected since in HBphoto3 it was shown that there are large cancellations between fourth-order loop and counter term contributions, which we represent by the polynomial term only. This discrepancy is even more pronounced for the combination of subtraction constants that can be obtained from $`\overline{P}_2`$. Utilizing $`\xi _1=19.7`$ from HBphoto3 and the values from Eq. (3), we obtain $`a_4^0/2+m_Na_3^1m_Na_2^0=208.5`$ GeV<sup>-3</sup> and $`3.3`$ GeV<sup>-3</sup>, respectively. This shows that the cancellations between the fourth order loop and counter term contributions are even stronger in $`P_2`$ than in $`P_1`$.
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# The distribution and kinematics of early high-𝜎 peaks in present-day haloes: implications for rare objects and old stellar populations
## 1 Introduction
In a Universe where cold dark matter (CDM) dominates structure formation, the haloes of galaxies and clusters are assembled via the hierarchical merging and accretion of smaller progenitors (e.g. Lacey & Cole 1993). This process causes structures to relax violently to a new equilibrium by redistributing energy among the collisionless mass components. Early stars formed in these progenitors behave as a collisionless system just like the dark matter particles in their host haloes, and they undergo the same dynamical processes during subsequent mergers and the buildup of larger systems like massive galaxies or clusters. It is of crucial importance in galaxy formation studies to explore the efficiency of the mixing process and see if any spatial or kinematical signatures exist in material that collapses at different epochs and within peaks of the primordial Gaussian density field of different rarity.
In this paper, we use a suite of high-resolution cosmological N-body simulations to analize the distribution and kinematics within present-day galaxy haloes of dark matter particles that originally belonged to selected branches of the merger tree. These properties are particularly relevant for the baryonic tracers of early CDM structures, e.g. the old stellar halo which may have originated from the disruption at high redshift of numerous dwarf protogalaxies (Bullock et al., 2000), the old halo globular clusters, and also giant ellipticals (Gao et al., 2004). The end product of the entire merger tree is a triaxial cuspy dark matter halo (Dubinski & Carlberg 1991; Navarro et al. 1996; Moore et al. 1999; Diemand et al. 2005): a small fraction of early progenitor systems survive the merging process and end up as dark matter substructures (Ghigna et al., 1998). Since rare, early haloes are strongly biased towards overdense regions (e.g. Cole & Kaiser 1989; Sheth & Tormen 1999), i.e. towards the centers of larger scale fluctuations that have not collapsed yet, we might expect that material originating from the earliest branches of the merger tree is today much more centrally concentrated than the overall halo. Indeed, a “non-linear” peak biasing has been discussed previously by several authors (Moore et al. 1998; White & Springel 2000; Moore 2001).
Here we show that the distribution and kinematics of “old material” within present-day galaxy haloes depends primarily on the rareness of the peaks of the primordial density fluctuation field it originally belonged to. Specifically, today’s properties of objects that formed in old rare density peaks above $`\nu \sigma (M,z)`$ \[where $`\sigma (M,z)`$ is the rms fluctuation in the density field linearly extrapolated to redshift $`z`$ smoothed with a top-hat filter of mass $`M`$. A “one $`\sigma `$ peak” corresponds to the characteristic mass “$`M_{}(z)`$” and $`\sigma (M_{},z)1.69`$\], depend largely on $`\nu `$ and not on the particular values of $`z`$ and $`M`$. Such centrally concentrated components are isotropic in the inner part, just like the host galaxy haloes, but rapidly become more radially anisotropic further out. The plan of the paper is as follows. In § 2 we describe the numerical simulations and how to define and trace high-$`\sigma `$ particle subsets. In § 3 we analyse the present-day distribution of these subsets. We derive a simple empirical fitting formula for $`\rho (r,\nu )`$, the mass density profile of all progenitors above $`\nu \sigma `$, which approximates the results of our N-body cosmological simulations for $`1<\nu <4`$. § 4 discusses the implications of our findings for old stellar populations. We argue that such centrally concentrated components are predicted to be isotropic in the inner part, just like the host galaxy halo, but to rapidly become more radially anisotropic further out. This is quantitative with the radial velocity dispersion profile of Galaxy halo stars from Battaglia et al. (2005) and with the anisotropic orbits of nearby halo stars ($`\beta 0.5`$, Chiba & Beers 2000). Finally, we present our conclusions in § 5.
## 2 Method
We identify collapsed high-$`\sigma `$ peaks at different epochs within high-resolution cosmological N-body simulations, mark them, and analyse the distribution and kinematics of this material at redshift zero. Details about the simulations and halo finding method are given in the next two subsections.
### 2.1 Simulations
The simulations have been performed using PKDGRAV, a parallel N-body treecode written by Stadel and Quinn (Stadel, 2001): cosmological and numerical parameters are the same as in Diemand et al. (2004c). The present-day haloes that we analyse are labeled “D12” (cluster) and “G0” to “G3” (galaxies). We also analyse an additional smaller galaxy halo ($`M_{\mathrm{vir}}=10^{11}\mathrm{M}_{}`$), labeled “G4”.
Figures 1 and 2 show the density field of the region containing the G0-G4 galaxy haloes. Most of the material belonging to high-$`\sigma `$ peaks in the high-resolution region can be found in the four most massive galaxy haloes, while smaller structures tend not to have any progenitors which meet the selection criteria. G4 is a special, early-forming small galaxy halo (with a high concentration of $`c=18`$) that contains similar fractions of high-$`\sigma `$ material as the more massive galaxies. From Figure 1 it is clear that G4 is not a representative system for its mass range. All “G” haloes form in the one region of $``$ (10 Mpc)<sup>3</sup> that was selected from a (90 Mpc)<sup>3</sup> box and re-simulated at high resolution. The region has a mean density of 0.5 times the background density at $`z=0`$. The properties of the six haloes are given in Table 1.
The cosmological parameters are $`(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\sigma _8,h)=(0.268,0.732,0.7,0.71)`$. The value of $`\sigma _8=0.9`$ given in Diemand et al. (2004c) is not correct: we found that, due to a mistake in the normalization, our initial conditions have less power than intended. According to linear theory, mass fluctuations grow proportional to the scale factor in a flat $`\mathrm{\Omega }_m=1`$ Universe, which is a good approximation to the adopted $`\mathrm{\Lambda }`$CDM Universe at $`z>2`$. Therefore the scale factor at collapse of a given halo in our simulations would be 0.78 times smaller in a $`\sigma _8=0.9`$ model. Throughout this paper, together with the collapse redshifts given for our low-$`\sigma _8`$ simulations, we will also state the rescaled collapse redshifts in a $`\sigma _8=0.9`$ Universe using the exact growth function for a $`(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda })=(0.3,0.7)`$ cosmology (see table 2). Using a $`\sigma _8=0.9`$ resimulation of the “D” cluster we confirm that the $`z=0`$ results (in units of the scale radius) are consistent with those of the $`\sigma _8=0.7`$ simulation when we use the $`z_{0.9}`$ outputs of the $`\sigma _8=0.9`$ run to select progenitor haloes.
### 2.2 Tracing progenitor material
Since baryonic objects (i.e. early protogalaxies, first stars) will form in the inner parts of their host gravitational potential, it is important to assess how different their $`z=0`$ distribution is when only the inner part of a progenitor halo is traced. We will therefore compare the final distribution of material marked within the central regions of early haloes with the entire marked progenitor. We will show in this section that there is no difference (except of course in the traced mass fraction) if the traced objects are early ($`z\mathrm{}>10`$) protogalaxies.
Progenitor haloes are identified using the friends-of-friend algorithm (FOF) (Davis et al., 1985). For simplicity we use a fixed linking length $`b=0.164`$ at all output times. Our results are not sensitive to this choice. For example at redshift 10.5 we found 52 haloes more massive than $`4.35\times 10^9\mathrm{M}_{}`$ in the high resolution region of the galaxy simulation. When we use a linking length two or three times smaller and mark the linked material in the same 52 groups we obviously find a smaller fraction of marked matter but very similar density profile shapes and kinematics (see Fig. 3). We also selected the cores of haloes at the same redshift by marking all particles which have a local density (calculated from an SPH kernel over 32 nearest neighbours) $`10^3`$ times higher than the mean matter density and obtained similar $`z=0`$ distributions. This demonstrates that at redshift zero, particles originating from the cores of high-$`\sigma `$ haloes are distributed in the same way as all the progenitor halo material, and justifies our choice of using all halo particles as tracers for baryonic objects which would likely form in the very center of their hosts. A large number of tracers give two important advantages. First, it allows us to reliably estimate density, velocity dispersion, and anisotropy profiles. Second, it makes the results more robust against numerical effects: since two-body relaxation completely changes the orbits of many individual particles, even in high-resolution cosmological simulations (Diemand et al., 2004a), using only one most bound particle as a tracer for (say) a Population III remnant is not a safe choice.
We have shown above that the final distribution and kinematics of particles from high-redshift, low-mass progenitor haloes is insensitive to their original location within the host potential. The orbital energy associated with the merging of these small subunits into larger systems determines the present-day distribution and dominates the relatively small differences in the orbital and potential energies of particles within their hosts. This is not true at lower redshifts and for progenitor masses closer to the mass of the $`z=0`$ parent halo. For example, progenitors above $`1\sigma `$ marked at $`z=1.6`$ (i.e. with masses above $`7.4\times 10^9\mathrm{M}_{}`$) show only mild radial and velocity bias when the entire FOF groups are marked (Fig. 4): the difference from the total dark matter distribution becomes much larger, however, if we select within these haloes only particles with overdensity greater than $`10^4`$, as this material is more concentrated, colder, and on more radial orbits (just like in the low-mass 2$`\sigma `$ selection, see Fig. 5).
## 3 Distribution and kinematics of high-$`\sigma `$ peaks
We use the linear growth approximation to structure formation to calculate the masses of 4, 3.5, 3, 2 and 1$`\sigma `$ fluctuations collapsing at a given redshift. The redshifts of the simulation outputs and the halo masses corresponding to these fluctuations are given in Table 2. These values are for the cosmological model we have simulated $`(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\sigma _8,h)=(0.268,0.732,0.7,0.71)`$. In a $`\sigma _8=0.9`$ the same fluctuations would collapse earlier (at $`z_{0.9}`$, see § 2.1).
### 3.1 Radial distribution
High-$`\sigma `$ material is strongly biased toward the center of present-day haloes (see upper panels of Figs. 4 to 12), with the rarer peaks showing stronger bias. Figures 4 to 9 show different (M,z) selections corresponding to the same value of $`\nu `$. While the shapes of these profiles are similar for a given $`\nu `$, the normalisation (or traced mass fraction) generally grows with increasing mass threshold and to lower redshifts. The density profiles are averaged over the four galaxy haloes of similar mass. The scatter from halo to halo in their total density profiles and that of their subsets is relatively small. But there are substantial halo-to-halo variations in the anisotropy parameter $`\beta (r)`$, which we discuss in § 3.2.
Mass densities are fitted with a general $`\alpha \beta \gamma `$-profile that asymptotes to a central cusp $`\rho (r)r^\gamma `$:
$$\rho (r,\nu )=\frac{\rho _s}{(r/r_\nu )^\gamma [1+(r/r_\nu )^\alpha ]^{(\beta \gamma )/\alpha }}.$$
(1)
For comparison the NFW (Navarro, Frenk & White, 1996) profile has $`(\alpha ,\beta ,\gamma )=(1,3,1)`$, while the Moore et al. (1999) profile has $`(\alpha ,\beta ,\gamma )=(1.5,3,1.5)`$. We fix $`\alpha =1`$ and the inner slope to $`\gamma =1.2`$, which is the best-fit slope for the $`D12`$ cluster when resolved at very high resolution ($`m_{\mathrm{DM}}=3.0\times 10^5\mathrm{M}_{}`$) (Diemand, Moore & Stadel, 2005). We fit the entire dark halo using an outer slope of $`\beta =3`$ to determine the scale radius $`r_s=r_{vir}/c`$, where $`c`$ is the concentration. To approximate the high-$`\sigma `$ subset profiles we use a smaller scale radius $`r_\nu r_s/f_\nu `$ (corresponding to a higher concentration $`c_\nu f_\nu c`$) and also a steeper outer slope $`\beta _\nu `$. The $`f_\nu `$ and $`\beta _\nu `$ values used in the plots are calculated with simple empirical formulae which approximately parameterise the entire range of profiles, i.e. peaks above 1 to 4$`\sigma `$ and haloes ranging from a low concentration ($`c=4.5`$) cluster halo to a small, $`c=17`$ galaxy halo:
$$r_\nu r_s/f_\nu ,f_\nu =\mathrm{exp}(\nu /2),\beta _\nu =3+0.26\nu ^{1.6}.$$
(2)
The values for the 1 to 4$`\sigma `$ peaks are given in Table 3, and the profiles are plotted in the upper left panels of Figures 4 to 12 with open triangles. The fits are just approximate but they are reproduce well at least one profile for each of the $`\nu `$ values. Other parameters and functional forms could fit some of the data better.
For easier comparison with extragalactic observational data of old stellar or globular cluster populations we also fit our high-$`\sigma `$ subset profiles with a deprojected $`R^{1/4}`$ law, using the accurate numerical approximation of Márquez et al. (2001):
$`\rho (r)`$ $`=`$ $`\rho _0(r/R_e)^pe^{b(r/R_e)^{1/4}},`$ (3)
$`p`$ $`=`$ $`1.00.6097(1/4)+0.05563(1/4)^2,`$
$`b`$ $`=`$ $`1.9992\times 40.3271.`$
The best-fit effective radii $`R_e`$ are given in Table 3. Like the $`r_\nu `$ values they scale with $`r_s`$ and not with the virial radius.
### 3.2 Kinematics
The lower left panels of Figures 4 to 12 show the three dimensional velocity dispersion profiles $`\sigma _{3\mathrm{D}}^2=\sigma _r^2+\sigma _t^2`$, where $`\sigma _r`$ and $`\sigma _t`$ are the radial and the tangential velocity dispersions ($`\sigma _t^2=\sigma _\theta ^2+\sigma _\varphi ^2`$). The lower right panels show the anisotropy parameter $`\beta =10.5\sigma _t^2/\sigma _r^2`$ as a function of radius. The high-$`\sigma `$ subsets are clearly slower and on more radial orbits than the entire dark matter component at the same radius. As with the density profiles, the differences depend on $`\nu `$, i.e. the highest $`\sigma `$ peaks have the lowest $`\sigma _{3\mathrm{D}}`$ and the largest $`\beta `$ values.
Like the density profiles also the velocity dispersions are averaged over the four galaxy-size parent haloes. There is little halo-to-halo scatter in $`\rho (r)`$ and $`\sigma _{3\mathrm{D}}(r)`$, but there is substantial scatter in the anisotropy parameter $`\beta (r)`$, especially near the virial radius. We have checked that the radial bias exists in each data set and is not just on average. The variation in $`\beta `$ over the four galaxies is about 0.1 within 10% of the virial radius, for both the parent halo and the progenitor subsets. Further out the scatter becomes much larger: particles in the parent halo (subsets) have a total spread of 0.45 (0.35) near $`r_{\mathrm{vir}}`$.
#### 3.2.1 Jeans equations
The origin of the smaller velocity dispersions of the subsets can be traced to the fact that high-$`\sigma `$ peaks form closer to the main progenitor of the present-day parent halo, hence they join the system with small infall velocities and at early times when its potential well is much shallower. Since they have been part of the parent halo for a long time, they are likely in dynamical equilibrium with the host, i.e. their density and velocity profiles should be a stationary solution to the Jeans equation in the potential, $`\mathrm{\Psi }(r)`$, of the parent halo. Setting time derivatives to zero and assuming spherical symmetry, the Jeans equation can be written as (Binney & Tremaine, 1987)
$$\frac{d}{dr}\left(\rho \sigma _r^2\right)+\frac{2\beta }{r}\rho \sigma _r^2+\rho \frac{d\mathrm{\Psi }}{dr}=0.$$
(4)
Approximating the anisotropy parameter as in Diemand et al. (2004c),
$$\beta (r)=\beta _{\mathrm{vir}}(r/r_{\mathrm{vir}})^{1/3},$$
(5)
equation (4) admits the solution
$$\rho (r)\sigma _r^2(r)=Ge^{6\beta (r)}_r^{\mathrm{}}e^{6\beta (y)}y^2\rho (y)M(y)𝑑y.$$
(6)
Figure 5 depicts this solution for 2$`\sigma `$ material. We use equation (1) with $`(\alpha ,\beta _\nu ,\gamma )=(1,5.4,1.2)`$ and $`r_\nu =r_s=r_{\mathrm{vir}}/10`$ to approximate the density profile (see long-dashed lines in Fig. 5), and equation (5) with $`\beta _{\mathrm{vir}}=0.8`$ to approximate the anisotropy profile $`\beta (r)`$. The radial velocity dispersion calculated from (6) is converted to $`\sigma _{3\mathrm{D}}`$ and plotted (filled squares) in Figure 5. It is very close to the measured velocity dispersions which confirms our expectation that high-$`\sigma `$ material is a more concentrated and colder subset of particles in dynamical equilibrium within the total dark matter potential. This is just the opposite situation to the one of surviving CDM subhaloes, which are a more extended, hot subset in equilibrium with the host potential (Diemand, Moore & Stadel, 2004b). We also calculated an isotropic \[$`\beta (r)=0`$\] and a constant $`\beta (r)=0.5`$ model (see van den Bosch et al. 2004 for the corresponding solution of eq. 4) but both cases do not to reproduce the measured $`\sigma _{3\mathrm{D}}(r)`$ very well (see circles vs. long-dashed line in the lower left panel of Fig. 5).
As a consistency check we also solve eq. (6) for the entire dark matter halo, using $`(\alpha ,\beta _\nu ,\gamma )=(1,3,1.2)`$ and $`\beta _{\mathrm{vir}}=0.45`$ to approximate the density and anisotropy profiles. The resulting $`\sigma _{3\mathrm{D}}`$ is plotted (open squares) in the lower left panel of Figure 5. The agreement with the measured dispersion is very good. As expected our average system from four isolated and relaxed galaxy haloes is very close to a stationary equilibrium solution.
#### 3.2.2 A fitting function for $`\beta (r,\nu )`$
Equation (5) is not a very good approximation to the average anisotropy $`\beta (r)`$ of the four G0-G3 parent galaxies. The function used by Mamon & Lokas (2005)
$$\beta (r)=\beta _a\frac{r}{r+r_a}$$
(7)
fits the data quite well for $`\beta _a=0.45`$ and $`r_a=0.065`$. For the high-$`\sigma `$ particles the radial anisotropy is larger and sets in further inside. This behavior is approximated by simply using the scale radii of the high-$`\sigma `$ density profiles (2) instead of $`r_a=0.065`$, and an amplitude $`\beta _a`$ that grows with $`\nu `$:
$`\beta (r,\nu )`$ $`=`$ $`\beta _{a\nu }{\displaystyle \frac{r}{r+r_{a\nu }}},r_{a\nu }r_s/f_\nu ,`$
$`f_\nu `$ $`=`$ $`\mathrm{exp}(\nu /2),\beta _{a\nu }=10.4\nu ^{0.5}.`$ (8)
These anisotropy profiles are plotted with open triangles in the lower right panels of Figures 4 to 9. They also approximate the inner $`\beta (r)`$ of the haloes D12 and G4. Outside 10% of the virial radius both the cluster D12 and the small galaxy G4 deviate substantially from the average $`\beta (r)`$ found for G0-G3.
#### 3.2.3 Density slope-velocity anisotropy relation
Hansen & Moore (2004) have recently found a relation between the logarithmic slope of a density profile and the its velocity anisotropy: $`\beta =\eta _1\eta _2(d\mathrm{log}\rho /d\mathrm{log}r)`$, with $`\eta _1`$ ranging from -0.45 to 0.05 and $`\eta _2`$ from 0.1 to 0.35. This relation approximates values measured in a variety of equilibrium N-body and SPH systems, including CDM haloes. It is interesting to ask whether a similar relation also exist for our high-$`\sigma `$ subsets of present-day galaxy haloes. By combining the fitting functions for $`\rho (r,\nu )`$ (eq. 2) and $`\beta (r,\nu )`$ (eq. 3.2.2) one indeed finds the same simple relation,
$$\beta (d\mathrm{log}\rho /d\mathrm{log}r)=\eta _2(\gamma +d\mathrm{log}\rho /d\mathrm{log}r),$$
(9)
where $`\gamma =1.2`$ is the slope of the inner density cusp ($`\rho r^{1.2}`$). The coefficient $`\eta _2`$ depends only mildly on $`\nu `$. Parent CDM haloes and low-$`\sigma `$ subsets have $`\eta _20.3`$, while high-$`\sigma `$ subsets have smaller values, $`\eta _2`$=(0.29, 0.28, 0.23, 0.19) for $`\nu `$=(1, 2, 3, 4). The scatter around these relations is about 0.1 when $`\beta <0.5`$ (i.e. in the inner regions), and is larger (up to 0.4) in the outer parts. A larger, more representative sample of CDM haloes should be studied to quantify more accurately the exact slope and scatter of this important relation. The fact that even the high-$`\sigma `$ subsets follow a nearly universal relation between velocity anisotropy and density slope supports the existence of a fundamental connection between the two quantities (Hansen & Moore, 2004). This relation breaks the mass-anisotropy degeneracy present when one uses the line-of-sight velocities of extragalactic stellar halo objects (such as globular clusters) to infer total dynamical masses.
### 3.3 Shape of high-$`\sigma `$ material
The axis ratios $`a/c`$ and $`b/c`$ are calculated with TIPSY <sup>1</sup><sup>1</sup>1Available at: http://www-hpcc.astro.washington.edu/. using the technique described in Katz (1991). First the inertia tensors of the particles within a sphere of a given radius is calculated and diagonalized, then the same is done iteratively for particles within the triaxial shape found in the previous step, until the procedure converges. Note that the shape at a given radius depends on all particles within this radius, but for density profiles less steep than $`\rho (r)r^4`$ the outer particles dominate over the inner ones in the contribution to the inertia tensor.
The average axis ratios of the four galaxy haloes G0-G3 at the virial radius ($`a/c0.6,a/b0.7`$) are close to the mean values of large samples of haloes (Jing & Suto 2002; Faltenbacher et al. 2002; Bailin & Steinmetz 2005). Figure 7 shows the average shapes of 2, 3, and 3.5$`\sigma `$ material at the present epoch. Both axis ratios become smaller for higher-$`\sigma `$ material at all radii.(This is also true for each of the four galaxies individually, not just for the average.) This result is consistent with the finding of Jing & Suto (2002) that the average axis ratios decrease with redshift, although our 2$`\sigma `$ haloes formed at higher redshift and lower virial mass than the range probed by Jing & Suto (2002). The origin of the extremely prolate shape of the high-$`\sigma `$ subsets ($`a:b:c3:1:1`$ for the 3.5$`\sigma `$ material) could be a series of correlated head-on (i.e. low-angular momentum) mergers along a filament aligned with the long axis (Moore et al., 2004).
For the total dark matter halo the axis orientations change with radius, as found in Jing & Suto (2002), and a similar behavior is observed for the high-$`\sigma `$ subsets. The highest $`\sigma `$ subsets show a slightly better alignment which is simply explained with the fact that inner particles dominate the shape calculation due to the very high concentrations of these subsets. Between the different subsets the axes are generally well aligned only when they have many particles in common, i.e. in the inner parts. As the fraction of high-$`\sigma `$ material drops with increasing radius, the alignment with the shape of the whole halo becomes worse.
### 3.4 Mass fraction of high-$`\sigma `$ material
It is well known that high-$`\sigma `$ peaks at early times are biased towards overdense regions where larger haloes will form later (see Figs. 1 and 2). But how large is the mass fraction of such peaks within present-day haloes, and how does this fraction depend on the mass of the $`z=0`$ parent host? Do haloes which correspond to higher-$`\sigma `$ peaks today (like massive clusters) contain a larger fraction of early high-$`\sigma `$ material? This question is not well defined, since the fraction of particles from peaks above 2$`\sigma `$ in a present-day galaxy halo grows when the peaks are selected with a larger mass threshold at lower redshift (see upper right panel of Fig. 5). Here we use fixed mass threshold/redshift pairs to select progenitors and compare the mass fractions they contribute to parent haloes of different sizes. The selection of a fixed progenitor mass may be motivated in studies of old stellar populations if, for example: a) Population III stars form in “minihaloes” above a molecular cooling mass, i.e. massive enough to allow efficient gas cooling via roto-vibrational levels of H<sub>2</sub>, $`M>M_{\mathrm{H}_2}6\times 10^5[(1+z)/20]^{3/2}\mathrm{M}_{}`$ (virial temperatures $`>2000`$K); and b) Metal-poor halo stars (Population II) and globular clusters form in haloes above an atomic cooling mass, i.e. massive enough to allow efficient gas cooling and fragmentation via excitation of hydrogen Ly$`\alpha `$, $`M>M_\mathrm{H}10^8[(1+z)/10]^{3/2}\mathrm{M}_{}`$ (virial temperature $`>10^4`$K).
To test whether, at the present epoch, massive clusters contain a different fraction of high-$`\sigma `$ material than field galaxies, we have used a lower-resolution N-body simulation ($`300^3`$ particles in a 90 Mpc box, for a resolution of $`10^9\mathrm{M}_{}`$) in order to obtain a larger galaxy/cluster sample. We have selected ten parent hosts at $`z=0`$ with FOF masses close to $`2\times 10^{11}\mathrm{M}_{}`$, $`10^{12}\mathrm{M}_{}`$, $`10^{13}\mathrm{M}_{}`$, and $`10^{14}\mathrm{M}_{}`$, for a total of fourty haloes. We have then marked FOF groups with more than 100 particles at redshifts 7.0, 4.3, 3.1 and 0.8, and determined the fractions of the virial mass of the parent belonging to a fixed progenitor mass/redshift ($`M,z`$) pair. The values obtained are given in Table 4: the mass fractions are practically constant for all parents which host any of the selected progenitors. As we select parents of lower masses, the number of parents hosting rarer (higher-$`\sigma `$) progenitors drops from 10 to 0 within about a decade in mass.
## 4 Some applications
We have shown that the final distribution and kinematics of dark matter particles selected from early branches of the merger tree are systematically different than those of the parent halo as a whole. These properties are also relevant for old stellar populations if these form predominantly in early low-mass progenitor haloes, as stars behave essentially as collisionless systems just like the dark matter particles in our simulations. In the following we briefly discuss a number of possible applications.
### 4.1 Remnants of the first stars
Numerical simulations performed in the context of hierarchical structure formation theories suggest that the first (Population III) stars may have formed out of metal-free gas in dark matter minihaloes of mass above $`6\times 10^5\mathrm{M}_{}`$ (Abel et al. 2000; Bromm et al. 2002; Yoshida et al. 2003; Kuhlen & Madau 2005) condensing from rare high-$`\sigma `$ peaks of the primordial density fluctuation field at $`z>20`$, and were likely very massive. Barring any fine tuning of the initial mass function (IMF) of Population III stars, intermediate-mass black holes (IMBHs) – with masses above the 5–20$`\mathrm{M}_{}`$ range of known ‘stellar-mass’ holes – may be one of the inevitable endproduct of the first episodes of pregalactic star formation (Madau & Rees 2001).
Where do relic pregalactic IMBHs lurk in present-day galaxy halos? To shed some light on this question, we have populated our 3$`\sigma `$ (3.5$`\sigma `$) simulated progenitors at $`z_{0.9}=17.9`$ ($`z_{0.9}=21.2`$) with one seed IMBH for every $`6\times 10^5`$ solar masses of halo material. As discussed by Volonteri et al. (2003), these IMBHs will undergo a variety of processes during the hierarchical buildup of larger and larger haloes, like gas accretion, binary hardening, black hole mergers, triple interactions. While we neglect all of these effects here, our dark matter simulations do correctly model the bias in the formation sites, the accretion into larger haloes, and the competing effects of dynamical friction and tidal striping within larger potential wells, as these complicated dynamical processes are dominated by the dark haloes that host the black holes. Therefore, in our toy model, the distribution of 3 or 3.5$`\sigma `$ material at $`z=0`$ describes the properties of holes wandering through within today’s galaxy haloes. The predicted IMBH number density and mass density profiles are shown as circles in Figures 8 and 9: the former may be regarded as an upper limit since we have neglected black hole mergers, while the latter have been estimated assuming that these off-muclear black holes have grown by accretion to a mean mass of $`1.5\times 10^4\mathrm{M}_{}`$, which is a rough estimate obtained from Figure 14 of (Volonteri et al.2003).
Depending on the IMF of Population III objects, some first-generation low-mass stars may have survived until today. Their number density profile $`n(r)`$ within the Milky Way can again be read following the circles in Figure 8, under the assumption that $`𝒩=1`$ metal-free star survives for every $`6\times 10^5`$ solar masses of 3$`\sigma `$ progenitor halo material at $`z_{0.9}=17.9`$. It is easy then to scale up the predicted value of $`n(r)`$ if $`𝒩1`$ such stars were to survive instead. On average, we find that about 1/3 of these remnants would lie today in the bulge, i.e. in the inner 3 kpc. This fraction fluctuates between 24% and 45% in the four galaxies G0-G3. The density in the solar neighborhood is of order $`0.1𝒩`$kpc<sup>-3</sup>, three orders of magnitude lower than in the bulge. The number density of remnants would be lower and more concentrated toward the galactic center if Population III stars only formed within rarer 3.5$`\sigma `$ peaks (Fig. 9). On average, about 59% of them would now lie within the bulge (range is 38%-84%) and the local number density would be only $`0.02𝒩`$kpc<sup>-3</sup>. Even for $`𝒩=10100`$, this is an extremely small value, many orders of magnitude below the local number density of halo stars. The above results suggest that the very oldest stars and their remnants should be best searched for within the Milky Way bulge.
### 4.2 Stellar haloes
Material from $`>2.5\sigma `$ peaks has today a density profile that is very similar to the stellar halo around the Milky Way (Moore et al., 2005). It contributes a few percent of the total virial mass, and therefore contains enough baryons to build up a $`10^9\mathrm{M}_{}`$ stellar halo with a reasonable star formation efficiency. The assumption of a common pregalactic origin between such a stellar component and the surviving Local Group dwarf galaxies provides an additional constraint and allows us to determine the progenitor mass threshold/redshift pair which best fit the data. From this argument Moore et al. (2005) identified hosts above $`10^8\mathrm{M}_{}`$ at $`z=12`$ as the progenitor haloes which the bulk of halo stars originally belonged to.
To check whether this simple model reproduces the kinematics as well as the radial distribution of halo stars, we have compared the predicted radial velocity dispersion profile with recent data from Battaglia et al. (2005) (see Fig. 10). All haloes were rescaled (i.e. $`r_{\mathrm{vir}}fr_{\mathrm{vir}}`$, $`V_cfV_c`$, $`M_{\mathrm{vir}}f^3M_{\mathrm{vir}}`$) to produce a local circular velocity of 220 km/s after taking into account the increase in circular velocity due to baryonic contraction \[$`V_c(8.5\mathrm{kpc})=1.125V_{c,\mathrm{max},\mathrm{DM}}`$\] found by Macció et al. (2005) for our halo G1. The rescaled virial masses are given in Figure 10. The predicted $`\sigma _r(r)`$ are close to the observed ones, especially for halo G0 and G4. Our model predicts a radial anisotropy that grows with radius (see lower right panel of Fig. 10): in the solar neighborhood it gives $`\sigma _{3\mathrm{D}}V_c`$ and $`\beta 0.4`$, in good agreement with the observed values (Chiba & Beers, 2000). The declining velocity dispersion near the center is characteristic of haloes containing only dark matter. In the Galaxy we expect a flatter profile in the inner, baryon-dominated regions.
Both a high concentration ($`c=18`$), highly anisotropic halo like G4 and a $`c=10`$ halo with smaller anisotropy like G0 fit the data equally well. A detailed comparison requires knowledge of the tangential components of the velocity dispersion, in order to determine the concentration and virial mass of the Milky Way halo. The recently discovered stellar halo of M31 has a similar density profile than that of our Galaxy (Guhathakurta et al., 2005). The line-of-sight velocity dispersion observed in M31 is mostly due to tangential motions: preliminary results yield a constant $`\sigma _t/\sqrt{2}`$ of about 80 km/s for $`50<r<150`$ kpc (Guhathakurta et al., 2005), close to the value expected for the G0 scenario but much higher than that of the G4 model. More data for M31 will hopefully soon become available, and many radial velocities of Milky Way halo stars will be obtained by surveys like RAVE and SDSS-II. In the next decade the planned GAIA satellite will provide very detailed phase-space positions.
Another interesting feature of our model is that it predicts significant deviations from the simple radial velocity dispersion profiles expected for a smooth stellar halo in an NFW potential (like the models in Battaglia et al. 2005, Fig. 4). These deviations are evident in the upper right panel of Figure 10 (but not in the lower left panel where we show an average over four haloes). <sup>2</sup><sup>2</sup>2The sample from Battaglia et al. (2005) and our G4 outer stellar halo contain similar numbers of ‘stars’. Both have relatively large Poisson noise and are consistent with smooth $`\sigma _r(r)`$ profiles. In haloes resolved with more particles like G0-G3 and D12, the ‘stellar’ radial velocity dispersions are clearly inconsistent with smooth profiles. The fluctuations are not an artifact of our analysis: with our usual choice of 30 logarithmic spherical bins out to the virial radius \[similar to the binning of the observed $`\sigma _r(r)`$\] the fluctuations often extend over various bins with each bin containing thousands of ‘stars’. To show fluctuations of $`\sigma _r(r)`$ on smaller scales we have binned ‘stars’ in halo G0 using linear spherical bins, and have ploted the resulting $`\sigma _r(r)`$ profile in the same radial range of the observations (circles with error bars in the upper right panel of Fig. 10). The linear bins contain at least 108 stars, the error bars represent Poisson noise. We have also plotted $`\sigma _{3\mathrm{D},\mathrm{DM}}(r)`$ of halo G0 twice using different bin sizes. Note that the velocity dispersion of the dark matter halo as a whole (instead of the material belonging to the $`>2.5\sigma `$ $`z=12`$ subset) in our simulated galaxies is generally a smooth function of radius, with small deviations becoming apparent only at $`r200`$ kpc. As fluctuations in the stellar $`\sigma _r(r)`$ are not caused by mass bound to subhaloes (the mass fraction is much too small particularly within the inner halo), this leaves tidal streams as the most probable cause of such deviations. The entire dark halo is assembled from about $`10^5`$ resolved progenitors, the tidal debris of which are well mixed and smooth. In our model, the stellar halo is made up instead from only a few hundred building blocks (and a significant fraction of it comes from the most massive ones), and is therefore expected to show some granularity. A more detailed study of tidal features in our model stellar haloes is left to a future paper.
#### 4.2.1 Comparison to recent related models for the formation of the stellar halo
During the completion of this work two related models for the formation of the stellar halo (Bullock & Johnston 2005; Abadi et al. 2005) came to our attention. Both studies are more detailed than our in the sense that they try to model the effect of star formation and chemical evolution. Bullock & Johnston (2005) use a semi-analytic approach combined with N-body simulations of satellite disruption in an external, growing galaxy potential. Like in this work the number of luminous satellites and the mass of the final stellar halo is adjusted to match the observations by assuming that reionisation suppresses star formation in later-forming small progenitors. Our stellar haloes, however, seem to be more concentrated than those of Bullock & Johnston (2005). This may due to the fact that we only trace old stellar populations, while Bullock & Johnston (2005) also allow for more recent episodes of star formation in their protogalactic building blocks. They also model the disruption of satellites assuming an external spherical potential: this is probably a good approximation to recent accretion events which build up the outer stellar halo. The inner stellar halo, however, is built up through a series of early, major mergers: these cause rapid potential fluctuations (violent relaxation) which are not accounted for in the spherical potential approximation. This limitation was also pointed out by Bullock & Johnston (2005).
The cosmological SPH simulations of Abadi et al. (2005), while not including any global mechanism for suppressing star formation in later-forming small progenitors, do lead to similar stellar halo profiles and kinematics as our model. They seem to produce, however, far too many halo stars: recent SPH simulations at higher resolution and without strong feedback effects also overpredict the number of satellites by a large factor Macció et al. (2005) and tend to produce even more halo stars. The similarity of the profile shape and kinematics of stellar haloes in all these models suggests that they are a robust and generic outcome of hierarchical structure formation (see also Hansen & Moore 2004; Dekel et al. 2005) and do not depend on the detailed formation history. Note that this does not contradict our argument about the possibility to use the radial extent of collisionless tracer populations to learn about their formation epoch and sites. It just illustrates how this arguments is limited to fossil records from very high redshift ($`z>8`$, see also § 2.2). At lower redshifts the typical haloes are larger and the exact star formation sites must be taken into account. Efficient cooling (and/or numerical losses of gas angular momentum to the dark matter in under-resolved progenitors) can produce a population of stars at some intermediate redshift with a distribution that is similar to much older populations from early high-$`\sigma `$ peaks. By $`z=0`$ the two subsets would have very similar concentrations and kinematics, which might explain the similar results found here and in Abadi et al. (2005).
### 4.3 Metal-poor globular clusters
The clustering properties of metal-poor globular clusters contain clues on their formation sites and the epoch when star formation was suppressed by feedback processes (e.g. reionisation, supernova-driven winds) in low-mass haloes. In the Milky Way metal-poor globular clusters follow the same radial profile as halo stars, suggesting within the framework of our model a common origin within early 2.5$`\sigma `$ peak progenitors at $`z12`$. Observations of the distribution of globular clusters within present-day haloes of different masses could provide information on feedback effects as a function of environment (Moore et al., 2005). The suppression of globular cluster formation at some early epoch may also explain the bimodality observed in cluster metallicities (e.g. Strader et al. 2005). The widely used assumption that globular cluster formation is a fair tracer of star formation, combined with the suppression of the formation of metal-poor globulars after reionisation, imply that the amount of high-$`\sigma `$ material in a halo is proportional to the number of metal-poor globular clusters. From the results in § 3.4 it follows then that a simple universal reionisation epoch would lead to a constant abundance of metal-poor globular clusters per virial mass. Deviations from this simplest case may provide information about the local reionisation epoch, and whether regions with more (less) metal-poor globulars per virial mass were reionised later (earlier) (see Moore et al. 2005).
### 4.4 Elliptical galaxies
The projected luminosity profiles of elliptical galaxies follow the R<sup>1/4</sup> law and resemble rescaled versions of projected CDM halo density profiles (Łokas & Mamon 2001; Merrit et al. 2005), similar to the $`z=0`$ profiles of our high-$`\sigma `$ subsets. Ordinary elliptical galaxies may form by multiple mergers or by the major merger of two disk galaxies (e.g. Dekel et al. 2005). Early formation ($`z6`$) in a large number of progenitors followed by a series of gas-poor, essentially collisionless mergers could build up the giant ellipticals (cD) in the centers of galaxy clusters. These mergers might undo some of the early dissipational contraction of the total mass distribution (Gao et al., 2004).
Figure 11 shows the distribution of material belonging to different $`\sigma `$ peaks in a $`z=0`$ cluster halo. The inner cluster regions are entirely made up of material from $`>2\sigma `$ peaks, and this selection contributes about 10% to the cluster virial mass. The central brightest galaxies in SDSS clusters have an average effective radius of $`R_e20`$kpc and presumably about 10% of the clusters luminosity (Zibetti et al., 2005). Our $`2\sigma `$ subsets contain enough mass to account for the observed cD luminosities but they are more extended. If the cD stars formed early in the inner parts of all progenitor haloes above 2$`\sigma `$ then they would be found closer to the cluster center today (see Section 2.2). A simple way to select the inner parts of 2$`\sigma `$ progenitors is to mark higher-$`\sigma `$ peaks at earlier epochs. Our 3$`\sigma `$ selections have a realistic effective radius of about $`20`$ kpc. Therefore we predict the kinematics of the average brightest cluster galaxies to follow the 3$`\sigma `$ selections shown in Figure (11) if (or where) the real cluster potentials are similar to the uncontracted potentials formed in dissipationless cosmological simulations (as suggested by Gao et al. 2004).
The anisotropy of a $`z=0`$ subset depends mostly on its density profile shape (or $`\nu `$) and not on the selection redshift or the number of mergers it has undergone (Figures 4 to 9) A few mergers are enough to establish the general $`\beta (r)`$ profile (3.2.2) and the slope – anisotropy relation (Section 3.2.3). This agrees with Moore et al. (2004) who find $`\beta 0.5`$ after only one low angular momentum merger of initially isotropic CDM haloes (see also Hansen & Moore 2004). Dekel et al. (2005) also find a similar $`\beta (r)`$ profiles after only one merger, which even starts with disk galaxies, i.e. nearly circular initial orbits for the stars (see their $`\beta `$ for the old stars, Fig. 1 of Dekel et al. 2005). Their average anisotropy at the half mass radius of the old stellar component is $`\beta =0.3`$, which is lower than our average of $`\beta =0.45`$, but considering the very different initial conditions of the stars and the larger numbers of mergers in our cosmological context the results are surprisingly similar. Therefore we argue that our collisionless results are relevant both for giant elliptical galaxies which formed in a series of dissipation-less mergers (Gao et al., 2004) and also for the old stars in smaller ellipticals.
The general relation between density pofile slope and anisotropy for high sigma subsets $`\beta (d\mathrm{log}\rho /d\mathrm{log}r)0.23(1.2+d\mathrm{log}\rho /d\mathrm{log}r)`$ can be applied to elliptical galaxies to infer $`\beta (r)`$ from the observed stellar profile. This is useful to obtain dynamical mass estimates for elliptical galaxies, where one has to assume some $`\beta (r)`$ (Mamon & Lokas, 2005) and it is probably the best guess before realistic, cosmological, hydrodynamical simulations of elliptical galaxy formation become available.
## 5 Conclusions
We have used high-resolution cosmological N-body simulations to trace the spatial distribution and kinematics in present-day CDM haloes of collisionless material that originally belonged to selected branches of the merger hierarchy. Our main results can be summarized as follows:
1. Hierarchical merging does not efficiently mix haloes up. The final distribution of dark matter particles retains a memory of when and where they collapsed initially, allowing a unique test of the popular bottom-up paradigm for the formation of cosmic structure.
2. Today’s distribution and kinematics of halo substructure that formed at very high redshift depends mostly on the rarity of the primordial density fluctuations they correspond to. For example, the $`z=0`$ density, shape, anisotropy, and velocity dispersion profiles of material originating from early 3$`\sigma `$ peaks is independent of the redshift/minimum mass pair used to select it. The mass fractions within the parent host, however, grow if such particles are choosen from lower redshift progenitors at a higher mass threshold.
3. High-$`\sigma `$ material should be looked for close to the center of massive parent haloes. The concentration and outer slope of the mass density profile are larger for the rarer peaks.
4. The anisotropy increases faster with radius for material originating within rarer peaks (eq. 3.2.2). The parameter $`\beta (r)`$ steepens with increasing $`\nu `$ at the same rate as the half-mass radius shrinks, i.e. we generally find $`\beta (r_{1/2})0.45`$.
5. Velocity dispersions are lower for material from the rarer peaks, and particle orbits are more radial. While the average velocity dispersion profiles agree with stationary solutions to the Jeans equation (4), the profiles of individual haloes show significant structure and are inconsistent with smooth $`\sigma (r)`$ solutions (see Fig. 10).
6. The high-$`\sigma `$ subsets have much more elongated shapes than their host haloes. For 3$`\sigma `$ material the mean $`c/a`$ is about 0.35, while present-day galaxy-size haloes have a mean of 0.6.
7. These properties do not depend on which regions within early ($`z>8`$) progenitors are marked and traced to the present time. Particles originating within the cores of high-$`\sigma `$ progenitors are distributed at $`z=0`$ in the same way as the entire virial mass of these early structures. This is not true at later epochs, where the marking of only the densest progenitor regions results in more concentrated density profiles (e.g. the selection of only the densest 10% of the mass of progenitor halos above 1$`\sigma `$ at $`z=1.6`$ yields a present-day distribution similar to that of all particles above 2$`\sigma `$).
8. If the first stars form at early epochs in peaks above 3.5$`\sigma `$, then half of their remnants should be found in the bulge, within 3 kpc of the galactic center. In the solar neighborhood the density of such very old population is 1000 times lower than in the bulge. Also, their characteristic velocities are 2.5 times lower than those of dark matter particles, and their orbits more radial ($`\beta 0.5`$).
9. The radial profile of the stellar halo and metal-poor globular clusters of the Milky Way suggest that these components formed in rare early peaks above 2.5$`\sigma `$ at redshift above 10. Typical outer halo objects have radial orbits, and become isotropic near the galactic center.
10. Radial orbits are a general outcome for any concentrated stellar component assembled through gas-poor mergers. The anisotropy parameter $`\beta (r,\nu )`$ correlates well with $`\rho (r,\nu )`$ and is not sensitive to the detailed assembly history of such a component.
The applications discussed above should be taken as a first attempt at predicting the spatial distribution and kinematics in present-day galaxies of objects that formed within early protogalactic systems, all in the context of hierarchical structure formation theories. Some interesting directions for future work may include, e.g., combining spatial information on Population III remnants with semi-analytic prescriptions for the growth and dynamical evolution of black holes, as well as using chemical evolution models to translate the age gradient for stellar haloes found here into metal abundance gradients. The techniques presented here also allow to study the properties of stellar streams in realistic, triaxial, clumpy CDM galaxy haloes.
## Acknowledgments
It is a pleasure to thank G. Battaglia for providing us with stellar halo data in electronic form. We are grateful to G. Battaglia, M. Beasley, A. Faltenbacher, M. Kuhlen and M. Volonteri for helpful suggestions and discussions. All computations were performed on the zBox supercomputer at the University of Zurich. Support for this work was provided by NASA grants NAG5-11513 and NNG04GK85G (P.M.), NSF grant AST-0205738 (P.M.), and by the Swiss National Science Foundation (J.D.). P.M. also acknowledges support from the Alexander von Humboldt Foundation. Part of this research was carried out at the Kavli Institute for Theoretical Physics, UC Santa Barbara, under NSF Grant No. PHY99-07949.
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# Ultrarelativistic boost of spinning black rings
## 1 Introduction
Shock pp -wave geometries describe the spacetime surrounding very fast moving objects, and are thus relevant to the study of Planckian scattering . They are also of interest in string theory, since strings may be exactly solved in such backgrounds . The prototype of shock wave solutions is the Aichelburg-Sexl spacetime, which represents the gravitational field of a massless point particle. It was originally obtained by boosting the Schwarzschild black hole to the speed of light, while rescaling the mass to zero in an appropriate way . According to recent extra-dimension scenarios, the fundamental Planck scale of (higher dimensional) gravity could be as low as a few TeV. This has stimulated renewed interest in the study of gravitational effects in high energy collisions, especially in view of the possible observation of microscopic black holes at near future colliders (see, e.g., for a recent review and for further references). It has been shown that closed trapped surfaces do indeed form in the ultrarelativistic collision of Aichelburg-Sexl point particles and of finite-size beams , which can more accurately model string-size effects. Nevertheless, it is desirable to understand how other effects could influence high energy scattering. A first step in this direction is to investigate more general shock wave solutions of higher dimensional gravity, which can naturally be obtained by applying the boosting technique of to black hole spacetimes. This has been done in any $`D4`$ for static black holes with electric charge or immersed in an external magnetic field . The ultrarelativistic limit of the Myers-Perry rotating black holes has been studied in (for the case of one non-vanishing spin). However, a striking feature of General Relativity in $`D>4`$ is the non-uniqueness of the spherical black holes of . In five-dimensional vacuum gravity, there exist also asymptotically flat rotating black rings with an event horizon of topology $`S^1\times S^2`$ . In the present contribution, we aim at studying the gravitational field generated by such rings in the Aichelburg-Sexl limit. As we will see in detail, this results in shock waves generated by extended lightlike sources (with a characteristic length-scale) which are remnants of the ring singularity of the original spacetime . Our recent results on boosted non-rotating black rings will be recovered as a special subcase. In general, the presence of spin is important because it allows black rings to be in equilibrium without introducing “unphysical” membranes via conical singularities . This will be reflected also in the shock geometry resulting from the boost. From a supergravity and string theory point of view, it is remarkable that supersymmetric black rings have been also constructed . We will conclude this article with a brief comment on the boost of such solutions. In the Appendix, we compare our results with those obtained for the ultrarelativistic limit of Myers-Perry black holes in $`D=5`$.
## 2 The black ring solution
In this section we briefly summarize the basic properties of the black ring, referring to for details. In the coordinates of ,<sup>1</sup><sup>1</sup>1Up to simple constant rescalings of $`F(\zeta )`$, $`G(\zeta )`$, $`C(\lambda ,\nu )`$, $`\psi `$ and $`\varphi `$, cf. Eqs. (2) and (4) with the corresponding ones in . In addition, multiply our $`L^2`$ by $`(1\nu )/(1\lambda )`$ to obtain the parameter used in . the line element reads
$`\mathrm{d}s^2=`$ $`{\displaystyle \frac{F(y)}{F(x)}}\left(\mathrm{d}t+C(\nu ,\lambda )L{\displaystyle \frac{1+y}{F(y)}}\mathrm{d}\psi \right)^2`$ (1)
$`+{\displaystyle \frac{L^2}{(xy)^2}}F(x)\left[{\displaystyle \frac{G(y)}{F(y)}}\mathrm{d}\psi ^2{\displaystyle \frac{\mathrm{d}y^2}{G(y)}}+{\displaystyle \frac{\mathrm{d}x^2}{G(x)}}+{\displaystyle \frac{G(x)}{F(x)}}\mathrm{d}\varphi ^2\right],`$
where
$$F(\zeta )=\frac{1+\lambda \zeta }{1\lambda },G(\zeta )=(1\zeta ^2)\frac{1+\nu \zeta }{1\nu },C(\nu ,\lambda )=\sqrt{\frac{\lambda (\lambda \nu )(1+\lambda )}{(1\nu )(1\lambda )^3}}.$$
(2)
The dimensionless parameters $`\lambda `$ and $`\nu `$ satisfy $`0\nu \lambda <1`$, and for $`\lambda =0=\nu `$ the spacetime (1) is flat. The constant $`L>0`$ represents a length related to the radius of the “central circle” of the ring. For a physical interpretation of the spacetime (1) we take $`y(\mathrm{},1]`$, $`x[1,+1]`$ (see a discussion in for other possible choices) and $`\psi `$ and $`\varphi `$ as periodic angular coordinates (see below). Surfaces of constant $`y`$ have topology $`S^1\times S^2`$. The coordinate $`\psi `$ runs along the $`S^1`$ factor, whereas $`(x,\varphi )`$ parametrize $`S^2`$ (see for illustrative pictures). Within the above range, $`y`$ parametrizes “distances” from the ring circle. At $`y\mathrm{}`$ the spacetime has a inner spacelike curvature singularity, $`y=1/\nu `$ is a horizon and $`y=1/\lambda `$ an ergosurface, both with topology $`S^1\times S^2`$. The black ring solution (1) is asymptotically flat near spatial infinity $`x,y1`$, where it tends to Minkowski spacetime in the form
$$\mathrm{d}s_0^2=\mathrm{d}t^2+\frac{L^2}{(xy)^2}\left[(y^21)\mathrm{d}\psi ^2+\frac{\mathrm{d}y^2}{y^21}+\frac{\mathrm{d}x^2}{1x^2}+(1x^2)\mathrm{d}\varphi ^2\right].$$
(3)
To avoid conical singularities at the axes $`x=1`$ and $`y=1`$, the angular coordinates must have the standard periodicity
$$\mathrm{\Delta }\varphi =2\pi =\mathrm{\Delta }\psi .$$
(4)
Centrifugal repulsion and gravitational self-attraction of the ring are in balance if conical singularities are absent also at $`x=+1`$, which requires
$$\lambda =\frac{2\nu }{1+\nu ^2}.$$
(5)
When this equilibrium condition holds, the metric (1) is a vacuum solution (of $`D=5`$ General Relativity) everywhere. With different choices (e.g., in the static limit $`\nu =\lambda `$ ), the conical singularity at $`x=+1`$ describes a disk-shaped membrane inside the ring.
The mass, angular momentum and angular velocity (at the horizon) of the black ring are
$$M=\frac{3\pi L^2}{4}\frac{\lambda }{1\lambda },J=\frac{\pi L^3}{2}\sqrt{\frac{\lambda (\lambda \nu )(1+\lambda )}{(1\nu )(1\lambda )^3}},\mathrm{\Omega }=\frac{1}{L}\sqrt{\frac{(\lambda \nu )(1\lambda )}{\lambda (1+\lambda )(1\nu )}}.$$
(6)
The algebraic type of the Weyl tensor of the ring spacetime is $`I_i`$ .
## 3 General boost
For our purposes, it is convenient to decompose the line element (1) as
$$\mathrm{d}s^2=\mathrm{d}s_0^2+\mathrm{\Delta },$$
(7)
in which $`\mathrm{d}s_0^2`$ is Minkowski spacetime (3) and
$`\mathrm{\Delta }=`$ $`\lambda {\displaystyle \frac{xy}{1+\lambda x}}\mathrm{d}t^22(1\lambda )C(\lambda ,\nu )L{\displaystyle \frac{1+y}{1+\lambda x}}\mathrm{d}t\mathrm{d}\psi `$ (8)
$`+{\displaystyle \frac{\lambda \nu }{1\nu }}{\displaystyle \frac{L^2}{1+\lambda y}}\left[\lambda {\displaystyle \frac{1+\lambda }{1\lambda }}{\displaystyle \frac{(1+y)^2}{1+\lambda x}}+{\displaystyle \frac{y^21}{xy}}\right]\mathrm{d}\psi ^2`$
$`+{\displaystyle \frac{L^2}{(xy)^2}}[\nu {\displaystyle \frac{x+1}{1\nu }}(y^21)\mathrm{d}\psi ^2+{\displaystyle \frac{\lambda (1\nu )(xy)+(\lambda \nu )(1+y)}{(1\lambda )(1+\nu y)}}{\displaystyle \frac{\mathrm{d}y^2}{y^21}}`$
$`+{\displaystyle \frac{\lambda \nu }{1\lambda }}{\displaystyle \frac{\mathrm{d}x^2}{(1x)(1+\nu x)}}+\nu {\displaystyle \frac{x+1}{1\nu }}(1x^2)\mathrm{d}\varphi ^2].`$
The above splitting is such that near infinity ($`x,y1`$) one has $`\mathrm{d}s^2\mathrm{d}s_0^2`$, while $`\mathrm{\Delta }`$ becomes “negligible” (in the sense of the “background” metric $`\mathrm{d}s_0^2`$). This enables us to define a notion of Lorentz boost using the symmetries of the asymptotic Minkowskian background $`\mathrm{d}s_0^2`$. Cartesian coordinates will visualize it most naturally. These can be introduced in two steps. First, we replace the coordinates $`(y,x)`$ with new coordinates $`(\xi ,\eta )`$ via the substitution
$$y=\frac{\xi ^2+\eta ^2+L^2}{\mathrm{\Sigma }},x=\frac{\xi ^2+\eta ^2L^2}{\mathrm{\Sigma }},$$
(9)
where
$$\mathrm{\Sigma }=\sqrt{(\eta ^2+\xi ^2L^2)^2+4L^2\eta ^2}.$$
(10)
The flat term $`\mathrm{d}s_0^2`$ in Eq. (7) now takes the form $`\mathrm{d}s_0^2=\mathrm{d}t^2+\mathrm{d}\eta ^2+\eta ^2\mathrm{d}\varphi ^2+\mathrm{d}\xi ^2+\xi ^2\mathrm{d}\psi ^2`$. Then, Cartesian coordinates adapted to the Killing vectors $`_\varphi `$ and $`_\psi `$ are given by
$$x_1=\eta \mathrm{cos}\varphi ,x_2=\eta \mathrm{sin}\varphi ,y_1=\xi \mathrm{cos}\psi ,y_2=\xi \mathrm{sin}\psi ,$$
(11)
so that $`\eta =\sqrt{x_1^2+x_2^2}`$, $`\xi =\sqrt{y_1^2+y_2^2}`$, and $`\mathrm{d}s_0^2=\mathrm{d}t^2+\mathrm{d}x_1^2+\mathrm{d}x_2^2+\mathrm{d}y_1^2+\mathrm{d}y_2^2`$. This enables us to study a boost along a general direction. Since the original spacetime (1) is symmetric under (separate) rotations in the $`(x_1,x_2)`$ and $`(y_1,y_2)`$ planes, such a direction can be specified by a single parameter $`\alpha `$, namely introducing rotated axes $`z_1`$ and $`z_2`$
$$x_1=z_1\mathrm{cos}\alpha z_2\mathrm{sin}\alpha ,y_1=z_1\mathrm{sin}\alpha +z_2\mathrm{cos}\alpha .$$
(12)
Defining now suitable double null coordinates $`(u^{},v^{})`$ by
$$t=\frac{u^{}+v^{}}{\sqrt{2}},z_1=\frac{u^{}+v^{}}{\sqrt{2}},$$
(13)
a Lorentz boost along $`z_1`$ takes the simple form
$$u^{}=ϵ^1u,v^{}=ϵv.$$
(14)
The parameter $`ϵ>0`$ is related to the standard Lorentz factor via $`\gamma =(ϵ+ϵ^1)/2`$. We are interested in “ultrarelativistic” boosts to the speed of light, i.e. in taking the limit $`ϵ0`$ in the transformation (14). While $`ϵ0`$, we will rescale the mass as $`M=\gamma ^1p_M2ϵp_M`$ , which physically means that the total energy remains finite in the limit ($`p_M>0`$ is a constant). Moreover, during the ultrarelativistic limit we wish to keep the angular velocity $`\mathrm{\Omega }`$ finite (a similar condition was imposed in ), and to allow for the possibility of black rings in equilibrium \[when the condition (5) holds\]. From Eq. (6), these requirements imply the rescalings<sup>2</sup><sup>2</sup>2This appears to be physically the most interesting and simple choice. See Footnote 3 for a subtler, slightly more general comment.
$$\lambda =ϵp_\lambda ,\nu =ϵp_\nu ,$$
(15)
where $`p_\lambda =8p_M/(3\pi L^2)`$ and $`p_\nu `$ is another positive constant such that $`p_\lambda p_\nu `$. In terms of these parameters, for $`ϵ0`$ the equilibrium condition (5) becomes
$$p_\lambda =2p_\nu .$$
(16)
Values $`p_\nu p_\lambda <2p_\nu `$ correspond to black rings (1) which are “underspinning” before the boost (and therefore balanced by a membrane of negative energy density), values $`p_\lambda >2p_\nu `$ to “overspinning” black rings (with a membrane of positive energy density). Notice, however, that under the limit $`ϵ0`$ the angular momentum $`J`$ will tend to zero (as $`ϵ`$).
We can now evaluate how the black ring metric (1) \[that is, Eq. (7) with Eqs. (3) and (8)\] transforms under the boost (14). We have first to substitute Eq. (9) into Eqs. (3) and (8). Then, we apply the sequence of substitutions (11), (12), (13) into the thus obtained expressions for $`\mathrm{d}s_0^2`$ and for $`\mathrm{\Delta }`$. Finally, we perform the boost (14) with the rescalings (15), which make $`\mathrm{\Delta }=\mathrm{\Delta }_ϵ`$ dependent on $`ϵ`$. The $`\mathrm{d}s_0^2`$ is invariant under the boost and at the end it reads
$$\mathrm{d}s_0^2=2\mathrm{d}u\mathrm{d}v+\mathrm{d}x_2^2+\mathrm{d}y_2^2+\mathrm{d}z_2^2.$$
(17)
The next step is to take the ultrarelativistic limit $`\mathrm{d}s^2=\mathrm{d}s_0^2+lim_{ϵ0}\mathrm{\Delta }_ϵ`$. This is delicate because the expansion of $`\mathrm{\Delta }_ϵ`$ in $`ϵ`$ has a different structure in different regions of the spacetime (even away from the singularity $`y=\mathrm{}`$). In particular, a peculiar behaviour is obtained for $`u=0`$, because $`\mathrm{\Delta }_ϵ`$ depends on $`u`$ through the combination
$$z_ϵ=\frac{1}{\sqrt{2}}(ϵ^1u+ϵv).$$
(18)
In order to have control over the exact distributional structure of the limit, it is convenient to isolate such dependence on $`ϵ^1u`$ by performing first an expansion of $`\mathrm{\Delta }_ϵ`$ with $`z_ϵ`$ unexpanded. This leads to an expression
$$\mathrm{\Delta }_ϵ=\frac{1}{ϵ}h(z_ϵ)\mathrm{d}u^2+\left[k_1(z_ϵ)\mathrm{d}x_2+k_2(z_ϵ)\mathrm{d}y_2+k_3(z_ϵ)\mathrm{d}z_2+k_4(z_ϵ)\mathrm{d}u\right]\mathrm{d}u+\mathrm{},$$
(19)
where the dots denote terms proportional to higher powers of $`ϵ`$, which are negligible in the limit. We have emphasized here the dependence of the functions $`h`$ and $`k_i`$ ($`i=1,\mathrm{},4`$) on $`z_ϵ`$ (and thus on $`ϵ`$), because this is essential in our limit, but they depend also on $`x_2`$, $`y_2`$ and $`z_2`$. The quantities $`k_i`$ are rather involved, but it suffices to observe here that $`lim_{ϵ0}k_i(z_ϵ)=0`$. We can thus also drop all the terms of order $`ϵ^0`$ in (19).<sup>3</sup><sup>3</sup>3A remark on the “triviality” of the $`ϵ^0`$ terms is in order, since they could be non-vanishing for certain more general scalings of the original metric parameters. While with higher order (in $`ϵ`$) corrections in Eq. (15) $`lim_{ϵ0}k_i(z_ϵ)=0`$ would still hold, we could introduce a non-vanishing contribution by allowing an $`ϵ`$-dependence in the ring “radius” via $`L_ϵ=L+c_1ϵ+c_2ϵ^2+\mathrm{}`$. The convergence of the integral (23) would then require $`c_1=0`$, but the quantity $`c_2ϵ^2`$ would affect the limit of (19) via $`lim_{ϵ0}k_4(z_ϵ)=c_2`$. The resulting term $`c_2\mathrm{d}u^2`$ is, however, obviously removable with a coordinate transformation. For $`h`$, after all the steps described above, we obtain explicitly
$`h(z_ϵ)=`$ $`p_\lambda {\displaystyle \frac{L^2}{\mathrm{\Sigma }}}+p_\nu {\displaystyle \frac{L^2}{\mathrm{\Sigma }^3}}\left[(\xi ^2\eta ^2L^2){\displaystyle \frac{y_1}{\xi }}\mathrm{sin}\alpha +2\xi x_1\mathrm{cos}\alpha \right]^2`$ (20)
$`+{\displaystyle \frac{1}{2}}(2p_\nu p_\lambda )\left(1{\displaystyle \frac{\xi ^2+\eta ^2L^2}{\mathrm{\Sigma }}}\right)\left({\displaystyle \frac{y_2^2}{\xi ^2}}\mathrm{sin}^2\alpha +{\displaystyle \frac{x_2^2}{\eta ^2}}\mathrm{cos}\alpha \right)`$
$`+\sqrt{p_\lambda (p_\lambda p_\nu )}{\displaystyle \frac{Ly_2\mathrm{sin}\alpha }{\xi ^2}}\left(1+{\displaystyle \frac{\xi ^2+\eta ^2+L^2}{\mathrm{\Sigma }}}\right)+(p_\lambda p_\nu ){\displaystyle \frac{L^2y_2^2}{\xi ^2\mathrm{\Sigma }}}\mathrm{sin}^2\alpha `$
$`+{\displaystyle \frac{1}{2}}(p_\lambda p_\nu )\left(1{\displaystyle \frac{\xi ^2+\eta ^2L^2}{\mathrm{\Sigma }}}\right).`$
Recall that the dependence of $`h`$ on $`ϵ`$ is contained in $`x_1`$ and $`y_1`$ via Eqs. (12)–(14), in $`\eta `$ and $`\xi `$ via Eq. (11) and in $`\mathrm{\Sigma }`$ via Eq. (10). In taking the limit $`ϵ0`$ of Eq. (19), we apply the distributional identity
$$\underset{ϵ0}{lim}\frac{1}{ϵ}f\left(z_ϵ\right)=\sqrt{2}\delta (u)_{\mathrm{}}^+\mathrm{}f(z)dz.$$
(21)
The final metric is thus \[cf. Eqs. (17) and (19)\]
$$\mathrm{d}s^2=2\mathrm{d}u\mathrm{d}v+\mathrm{d}x_2^2+\mathrm{d}y_2^2+\mathrm{d}z_2^2+H(x_2,y_2,z_2)\delta (u)\mathrm{d}u^2,$$
(22)
with a profile function given by
$$H(x_2,y_2,z_2)=\sqrt{2}_{\mathrm{}}^+\mathrm{}h(z)dz.$$
(23)
A black ring boosted to the speed of light in a general direction $`z_1`$ is thus described by the metric (22) with Eq. (23). This is evidently a $`D=5`$ impulsive pp -wave with wave vector $`_v`$. Such a spacetime is flat everywhere except on the null hyperplane $`u=0`$, which represents the impulsive wave front. Note that the equilibrium condition (16) has not yet been enforced in the above expression for $`h`$ (in particular, in the static limit $`p_\nu =p_\lambda `$ we recover the result of ). In order to write the solutions in a completely explicit form, it only remains to perform the integration in Eq. (23), with $`h`$ given by Eq. (20) with Eqs. (10)–(14) and (18). For any $`\alpha `$, this integral is always convergent and can in principle be expressed using elliptic integrals (because $`\mathrm{\Sigma }`$ is a square root of a fourth order polynomial in $`z`$, see for related comments). Therefore, no singular coordinate transformation of the type of has to be performed. In the following, we will explicitly calculate the integral, and study the corresponding solution in the case of two different boosts of the black ring along the privileged axes $`x_1`$ ($`\alpha =0`$) and $`y_1`$ ($`\alpha =\pi /2`$), which are respectively “orthogonal” and “parallel” to the 2-plane $`(y_1,y_2)`$ \[i.e., $`(\xi ,\psi )`$\] in which the ring rotates.
## 4 Orthogonal boost: $`\alpha =0`$
For the orthogonal boost $`\alpha =0`$, from Eq. (12) one has $`z_1=x_1`$ and $`z_2=y_1`$, so that the general pp -wave (22) reduces to
$$\mathrm{d}s^2=2\mathrm{d}u\mathrm{d}v+\mathrm{d}x_2^2+\mathrm{d}y_1^2+\mathrm{d}y_2^2+H_{_{}}(x_2,y_1,y_2)\delta (u)\mathrm{d}u^2.$$
(24)
Also, it is now convenient to rewrite $`h`$ in Eq. (20) as
$`h_{_{}}(z_ϵ)=`$ $`\left[3p_\lambda L^2(p_\lambda p_\nu )\xi ^2p_\nu (x_2^2+L^2)\right]{\displaystyle \frac{1}{2\mathrm{\Sigma }}}+p_\nu {\displaystyle \frac{4L^2\xi ^2z_ϵ^2}{\mathrm{\Sigma }^3}}`$ (25)
$`+{\displaystyle \frac{1}{2}}(2p_\nu p_\lambda )\left[{\displaystyle \frac{x_2^2(L^2\xi ^2)}{(z_ϵ^2+x_2^2)\mathrm{\Sigma }}}+{\displaystyle \frac{x_2^2}{z_ϵ^2+x_2^2}}\right]+{\displaystyle \frac{1}{2}}(p_\lambda p_\nu )\left(1{\displaystyle \frac{z_ϵ^2}{\mathrm{\Sigma }}}\right),`$
and $`\mathrm{\Sigma }`$ \[from Eq. (10)\] as
$$\mathrm{\Sigma }=\sqrt{\left[z_ϵ^2+x_2^2+(\xi +L)^2\right]\left[z_ϵ^2+x_2^2+(\xi L)^2\right]}.$$
(26)
Hereafter, it is understood that $`\xi =\sqrt{y_1^2+y_2^2}`$. In the orthogonal boost there is no contribution to $`h_{_{}}`$ from the off-diagonal term $`g_{t\psi }`$ in the metric (1). Performing the integration (23) with $`h`$ given by Eqs. (25) and (26), we find
$`H_{_{}}(x_2,y_1,y_2)=\sqrt{2}{\displaystyle \frac{3p_\lambda L^2+(2p_\nu p_\lambda )\xi ^2}{\sqrt{(\xi +L)^2+x_2^2}}}K(k)+\sqrt{2}(2p_\nu p_\lambda )`$
$`\times \left[\sqrt{(\xi +L)^2+x_2^2}E(k)+{\displaystyle \frac{\xi L}{\xi +L}}{\displaystyle \frac{x_2^2}{\sqrt{(\xi +L)^2+x_2^2}}}\mathrm{\Pi }(\rho ,k)+\pi |x_2|\mathrm{\Theta }(L\xi )\right],`$ (27)
where
$$k=\sqrt{\frac{4\xi L}{(\xi +L)^2+x_2^2}},\rho =\frac{4\xi L}{(\xi +L)^2},$$
(28)
and $`\mathrm{\Theta }(L\xi )`$ denotes the step function. In the above calculation, we have used the standard elliptic integrals and their properties summarized in the Appendix of , and the additional integral \[$`\mathrm{\Sigma }`$ given by Eq. (26) with $`z_ϵ`$ replaced by $`z`$\]
$$_0^{\mathrm{}}\left(1\frac{z^2}{\mathrm{\Sigma }}\right)dz=\sqrt{(\xi +L)^2+x_2^2}E(k).$$
(29)
In order to gain physical insight, it is useful to visualize the behaviour of the gravitational field at a large spatial distance within the wave front $`u=0`$. Defining the coordinates $`(r,\theta )`$
$$x_2=r\mathrm{cos}\theta ,\xi =r\mathrm{sin}\theta ,$$
(30)
an expansion for small values of the dimensionless parameter $`L/r`$ (using the identities summarized in ) leads to
$`H_{_{}}=`$ $`{\displaystyle \frac{\pi }{\sqrt{2}}}p_\lambda L[3{\displaystyle \frac{L}{r}}({\displaystyle \frac{5}{8}}+{\displaystyle \frac{p_\nu }{4p_\lambda }})(3\mathrm{cos}^2\theta 1){\displaystyle \frac{L^3}{r^3}}`$ (31)
$`+({\displaystyle \frac{7}{64}}+{\displaystyle \frac{p_\nu }{16p_\lambda }})(35\mathrm{cos}^4\theta 30\mathrm{cos}^2\theta +3){\displaystyle \frac{L^5}{r^5}}+O\left({\displaystyle \frac{L^7}{r^7}}\right)].`$
We recognize the standard form of multipole terms. The monopole is essentially an Aichelburg-Sexl term. The dipole and the octupole are missing, due to the geometry of the source. The quadrupole and 16-pole reflect the shape of the singularity and depend on the spin of the original black ring, but they persist even in the static limit $`p_\nu =p_\lambda `$ (when, in fact, they reach their maximal strength).
It is remarkable that for the physically more interesting case of black rings in equilibrium, i.e. those satisfying $`p_\lambda =2p_\nu `$ \[see Eq. (16)\], the profile function simplifies significantly to
$$H_{_{}}^e(x_2,y_1,y_2)=\frac{3\sqrt{2}p_\lambda L^2}{\sqrt{(\xi +L)^2+x_2^2}}K(k).$$
(32)
Interestingly, this is just the Newtonian potential generated by a uniform ring of radius $`L`$ and linear density $`\mu =3\sqrt{2}p_\lambda L/4`$ located at $`x_2=0`$ in the flat three-dimensional space $`(x_2,y_1,y_2)`$. Since for a general pp -wave (24) the only component of the Ricci tensor is $`R_{uu}=\frac{1}{2}\delta (u)𝚫H_{_{}}`$, $`𝚫`$ denoting the Laplace operator over the transverse space $`(x_2,y_1,y_2)`$, it follows that the profile function (32) represents a spacetime which is vacuum everywhere except on the circle $`u=0=x_2`$, $`\xi =L`$ \[so that $`k=1`$ in Eq. (28)\]. This lies on the wave front and corresponds to a singular ring-shaped source moving with the speed of light. It is obviously a remnant of the curvature singularity ($`y=\mathrm{}`$) of the original static black ring (1). For the non-equilibrium solution (4), the discontinuous term proportional to $`\mathrm{\Theta }(L\xi )`$ is responsible for a disk memebrane supporting the ring . We have plotted typical profile functions $`H_{_{}}`$ and $`H_{_{}}^e`$ in Fig. 1.
## 5 Parallel boost: $`\alpha =\pi /2`$
For the parallel boost $`\alpha =\pi /2`$, from Eq. (12) one has $`z_1=y_1`$ and $`z_2=x_1`$, and the general pp -wave (22) reduces to
$$\mathrm{d}s^2=2\mathrm{d}u\mathrm{d}v+\mathrm{d}x_1^2+\mathrm{d}x_2^2+\mathrm{d}y_2^2+H_{_{||}}(x_1,x_2,y_2)\delta (u)\mathrm{d}u^2.$$
(33)
The function $`h`$ can be reexpressed as
$`h_{_{||}}(z_ϵ)=`$ $`\left[(3p_\lambda +p_\nu )L^2p_\nu y_2^2+2\sqrt{p_\lambda (p_\lambda p_\nu )}Ly_2(p_\lambda p_\nu )\eta ^2\right]{\displaystyle \frac{1}{2\mathrm{\Sigma }}}p_\nu {\displaystyle \frac{4L^2\eta ^2z_ϵ^2}{\mathrm{\Sigma }^3}}`$ (34)
$`+{\displaystyle \frac{1}{2}}\left[(2p_\nu p_\lambda )y_2^22\sqrt{p_\lambda (p_\lambda p_\nu )}Ly_2\right]\left[{\displaystyle \frac{L^2+\eta ^2}{(z_ϵ^2+y_2^2)\mathrm{\Sigma }}}+{\displaystyle \frac{1}{z_ϵ^2+y_2^2}}\right]`$
$`+{\displaystyle \frac{1}{2}}(p_\lambda p_\nu )\left(1{\displaystyle \frac{z_ϵ^2}{\mathrm{\Sigma }}}\right),`$
and
$$\mathrm{\Sigma }=\sqrt{z_ϵ^4+2(y_2^2+\eta ^2L^2)z_ϵ^2+a^4},$$
(35)
with
$$a=\left[(\eta ^2+y_2^2L^2)^2+4\eta ^2L^2\right]^{1/4}.$$
(36)
It is understood that $`\eta =\sqrt{x_1^2+x_2^2}`$.
Performing the integration (23) with $`h`$ given by Eqs. (34)–(36), one obtains
$`H_{_{||}}(x_1,x_2,y_2)=[2(2p_\lambda p_\nu )L^2+(2p_\nu p_\lambda )a^2(1+{\displaystyle \frac{L^2+\eta ^2}{a^2y_2^2}})`$
$`+2\sqrt{p_\lambda (p_\lambda p_\nu )}Ly_2(1{\displaystyle \frac{L^2+\eta ^2}{a^2y_2^2}})]{\displaystyle \frac{\sqrt{2}}{a}}K(k)2\sqrt{2}(2p_\nu p_\lambda )aE(k)`$
$`+{\displaystyle \frac{\sqrt{2}}{2}}\left[(2p_\nu p_\lambda )y_22\sqrt{p_\lambda (p_\lambda p_\nu )}L\right]\left[{\displaystyle \frac{\eta ^2+L^2}{ay_2}}{\displaystyle \frac{a^2+y_2^2}{a^2y_2^2}}\mathrm{\Pi }(\rho ,k)+\pi \text{sgn}(y_2)\right],`$ (37)
where
$$k=\frac{\left(a^2\eta ^2y_2^2+L^2\right)^{1/2}}{\sqrt{2}a},\rho =\frac{(a^2y_2^2)^2}{4a^2y_2^2}.$$
(38)
Again, we refer to the Appendix of , the only additional integral used here being \[$`\mathrm{\Sigma }`$ given by Eq. (35) with $`z_ϵ`$ replaced by $`z`$\]
$$_0^{\mathrm{}}\left(1\frac{z^2}{\mathrm{\Sigma }}\right)dz=2aE(k)aK(k).$$
(39)
With the coordinates
$$y_2=r\mathrm{cos}\theta ,\eta =r\mathrm{sin}\theta ,$$
(40)
the behaviour at large spatial distances is given by
$`H_{_{||}}=`$ $`{\displaystyle \frac{\pi }{\sqrt{2}}}p_\lambda L[3{\displaystyle \frac{L}{r}}+2\sqrt{{\displaystyle \frac{p_\lambda p_\nu }{p_\lambda }}}\mathrm{cos}\theta {\displaystyle \frac{L^2}{r^2}}+({\displaystyle \frac{7}{8}}{\displaystyle \frac{p_\nu }{4p_\lambda }})(3\mathrm{cos}^2\theta 1){\displaystyle \frac{L^3}{r^3}}`$ (41)
$`+{\displaystyle \frac{3}{4}}\sqrt{{\displaystyle \frac{p_\lambda p_\nu }{p_\lambda }}}(5\mathrm{cos}^3\theta 3\mathrm{cos}\theta ){\displaystyle \frac{L^4}{r^4}}`$
$`+({\displaystyle \frac{11}{64}}{\displaystyle \frac{p_\nu }{16p_\lambda }})(35\mathrm{cos}^4\theta 30\mathrm{cos}^2\theta +3){\displaystyle \frac{L^5}{r^5}}+O\left({\displaystyle \frac{L^6}{r^6}}\right)].`$
Notice that now there appear also a dipole and an octupole term, as a remnant of the angular momentum of the black ring.
We are especially interested in black rings in equilibrium (16), for which one is left with
$`H_{_{||}}^e(x_1,x_2,y_2)=`$ $`p_\lambda L\left[{\displaystyle \frac{3\sqrt{2}L}{a}}+{\displaystyle \frac{2y_2}{a}}\left(1{\displaystyle \frac{L^2+\eta ^2}{a^2y_2^2}}\right)\right]K(k)`$ (42)
$`+p_\lambda L\left[{\displaystyle \frac{\eta ^2+L^2}{ay_2}}{\displaystyle \frac{a^2+y_2^2}{a^2y_2^2}}\mathrm{\Pi }(\rho ,k)\pi \text{sgn}(y_2)\right].`$
This function is singular at the points satisfying $`u=0=\eta `$ and $`|y_2|L`$ \[$`k=1`$ in Eq. (38)\], i.e. on a rod of length $`2L`$ contained within the wave front. This is a remnant of the curvature singularity of the original static black ring (1), which has (infinitely) Lorentz-contracted because of the ultrarelativistic boost in the plane of the ring. For the same reason, and because the original ring was rotating, the rod-source corresponding to Eq. (42) is not uniform. The profile (42) corresponds to a vacuum spacetime everywhere except on the rod. Notice also that the apparent divergences of $`H_{_{||}}^e`$ at $`y_2^2=a^2`$ and $`y_2=0`$ is only a fictitious effect: the singular behaviour of the coefficient of $`\mathrm{\Pi }`$ in Eq. (42) is exactly compensated from that of $`K`$ in the first case and from the $`\text{sgn}(y_2)`$ function in the second case \[recall also the form of $`\rho `$ in Eq. (38)\]. Finally, it is interesting to observe that the antisymmetric part (in the coordinate $`y_2`$) of $`H_{_{||}}`$ and $`H_{_{||}}^e`$ comes entirely from the off-diagonal term $`g_{t\psi }`$ in the metric (1), which was responsible for rotation before the boost \[and produces the terms even in $`\mathrm{cos}\theta `$ in the expansion (41)\]. The profile function $`H_{_{||}}^e`$ is plotted in Fig. 2.
## 6 Boost of the supersymmetric black ring
To conclude, we demonstrate that the above method can also be employed to calculate the gravitational field generated by other black rings in the ultrarelativistic limit. The first supersymmetric black ring (solution of $`D=5`$ minimal supergravity) was presented in (and subsequently generalized in ). The line element reads
$$\mathrm{d}s^2=f^2(\mathrm{d}t+\omega _\psi \mathrm{d}\psi +\omega _\varphi \mathrm{d}\varphi )^2+f^1(\mathrm{d}s_0^2+\mathrm{d}t^2),$$
(43)
with $`\mathrm{d}s_0^2`$ as in Eq. (3) and
$`f^1`$ $`=`$ $`1+{\displaystyle \frac{Qq^2}{2L^2}}(xy){\displaystyle \frac{q^2}{4L^2}}(x^2y^2),`$ (44)
$`\omega _\psi `$ $`=`$ $`{\displaystyle \frac{3}{2}}q(1+y)+{\displaystyle \frac{q}{8L^2}}(1y^2)\left[3Qq^2(3+x+y)\right],`$ (45)
$`\omega _\varphi `$ $`=`$ $`{\displaystyle \frac{q}{8L^2}}(1x^2)\left[3Qq^2(3+x+y)\right].`$ (46)
The $`S^1\times S^2`$ horizon is localized at $`y\mathrm{}`$, and asymptotic infinity at $`x,y1`$. The Maxwell field $`F=\mathrm{d}A`$ is determined by
$$A=\frac{\sqrt{3}}{2}f(\mathrm{d}t+\omega _\psi \mathrm{d}\psi +\omega _\varphi \mathrm{d}\varphi )\frac{\sqrt{3}}{4}q[(1+x)\mathrm{d}\varphi +(1+y)\mathrm{d}\psi ].$$
(47)
The net electric charge and the local dipole magnetic charge are proportional to the positive parameters $`Q`$ and $`q`$, respectively, which (for a physical interpretation) are assumed to satisfy $`Qq^2`$ and $`L<(Qq^2)/(2q)`$ . The mass and angular momenta of the ring are
$$M=\frac{3\pi }{4}Q,J_\psi =\frac{\pi }{8}q(6L^2+3Qq^2),J_\varphi =\frac{\pi }{8}q(3Qq^2).$$
(48)
In the limit $`q=0`$ the black ring becomes a static charged naked singularity, solution of the pure Einstein-Maxwell theory. In order to boost the line element (43), we can follow a procedure almost identical to the one used for the vacuum ring. The standard mass rescaling of together with the inequality $`L<(Qq^2)/(2q)`$ suggests that during the boost we rescale the charges as
$$Q=ϵp_Q,q=ϵp_q(p_Q>2Lp_q).$$
(49)
Omitting straightforward intermediate steps, in the case of a boost orthogonal to the plane $`(\xi ,\psi )`$ we obtain a shock pp -wave (24) with
$$H_{_{}}^s(x_2,y_1,y_2)=\frac{3\sqrt{2}p_Q}{\sqrt{(\xi +L)^2+x_2^2}}K(k),$$
(50)
and $`k`$ given by Eq. (28). For a parallel boost, we obtain the metric (33) with
$`H_{_{||}}^s(x_1,x_2,y_2)=`$ $`3\sqrt{2}\left[p_Q{\displaystyle \frac{1}{a}}+p_q{\displaystyle \frac{y_2}{a}}\left(1{\displaystyle \frac{L^2+\eta ^2}{a^2y_2^2}}\right)\right]K(k)`$ (51)
$`+{\displaystyle \frac{3\sqrt{2}p_q}{2}}\left[{\displaystyle \frac{\eta ^2+L^2}{ay_2}}{\displaystyle \frac{a^2+y_2^2}{a^2y_2^2}}\mathrm{\Pi }(\rho ,k)\pi \text{sgn}(y_2)\right],`$
where $`k`$ and $`\rho `$ as in Eq. (38). To obtain the field of a boosted naked singularity ($`q=0`$) just set $`p_q=0`$ in Eq. (51). Notice that the dipole charge $`q`$ has an effect only in the case of a parallel boost, since $`p_q`$ does not appear in $`H_{_{}}^s`$ \[which is in fact equivalent to the expression (32) for balanced vacuum rings\]. This is related to the “asymmetry” between the angular momenta $`J_\psi `$ and $`J_\varphi `$ in Eq. (48). In both boosts, one also finds that $`F=\mathrm{d}A`$ tends to zero together with its associated energy-momentum tensor (so that the “peculiar configuration” of does not arise here). In fact, both $`H_{_{}}^s`$ and $`H_{_{||}}^s`$ correspond to vacuum pp -waves. In principle, rescalings different from Eq. (49) can be considered if one drops the requirement $`L<(Qq^2)/(2q)`$. The detailed investigation of this and other possibilities is left for possible future work.
###### Acknowledgments.
M.O. is supported by a post-doctoral fellowship from Istituto Nazionale di Fisica Nucleare (bando n.10068/03).
## Appendix A Results for the boosted $`D=5`$ Myers-Perry black hole
Ref. analyzed the ultrarelativistic boost of $`D`$-dimensional Myers-Perry black holes with a single non-vanishing angular momentum. As in the present work, the calculation was performed in the case of two particular boosts orthogonal and parallel to the plane of rotation, and for $`D=5`$ it resulted in impulsive pp -waves of the type (24) and (33), respectively. It is thus interesting to compare the results of to ours. First of all, the angular momentum of black holes in $`D=5`$ must obey a Kerr-like bound $`a^2<\mu `$ . Since, in the Aichelburg-Sexl limit, Ref. sent the mass parameter $`\mu `$ to zero while keeping the spin parameter $`a`$ fixed, for $`D=5`$ the final metrics refer to boosted naked singularities rather than black holes . On the other hand, there is no upper limit on the spin of black rings , so that in our limit the rings do remain “black” until the final pp -wave is obtained (the same applies to the solutions of in $`D6`$, when also black holes can be ultra-spinning). In the rest of this appendix we shall present the profile functions of (for the case $`D=5`$) using an explicit form adapted to our notation<sup>4</sup><sup>4</sup>4In particular, the quantity $`L`$ will replace the original spin parameter $`a`$., and we shall compare them with our functions (4) and (5).
### A.1 Orthogonal boost
For an orthogonal boost, the result of can be rearranged as
$`\stackrel{~}{H}_{_{}}(x_2,y_1,y_2)=`$ $`{\displaystyle \frac{8\sqrt{2}p_M}{3\pi }}[{\displaystyle \frac{2\sqrt{2}}{(\xi ^2+x_2^2+L^2+b^2)^{1/2}}}K(k_1)+{\displaystyle \frac{\sqrt{2}}{L^2}}(\xi ^2+x_2^2+L^2+b^2)^{1/2}E(k_1)`$ (52)
$`{\displaystyle \frac{2\sqrt{2}}{L^2}}{\displaystyle \frac{b^2}{(\xi ^2+x_2^2+L^2+b^2)^{1/2}}}\mathrm{\Pi }(\rho _1,k_1)],`$
where
$`k_1`$ $`=`$ $`\left({\displaystyle \frac{\xi ^2+x_2^2+L^2b^2}{\xi ^2+x_2^2+L^2+b^2}}\right)^{1/2},\rho _1={\displaystyle \frac{(\xi ^2+x_2^2L^2b^2)^2}{4L^2x_2^2}},`$
$`b`$ $`=`$ $`[(\xi ^2+x_2^2L^2)^2+4x_2^2L^2]^{1/4}.`$ (53)
The above elliptic functions are singular for $`k_1=1`$, that is on a circle of radius $`L`$ given by $`x_2=0`$, $`\xi =L`$. This was already remarked in and it resembles our results of Sec. 4. Other physical properties are more “hidden” in the expression (52). First of all, for $`x_20`$ one has $`\rho _10`$ if $`\xi >L`$, whereas $`\rho _1`$ diverges if $`\xi <L`$. This implies \[with identity (A5) of \] that, when $`\xi <L`$ and $`x_2`$ is small, $`H_{_{}}`$ contains a non-smooth term proportional to $`|x_2|`$, namely there is an additional membrane at $`x_2=0`$ and $`\xi <L`$ (i.e. within the ring singularity discussed above). The presence of such a disk-shaped source is related to the structure of the singularities of the Myers-Perry solutions , and it should be contrasted with the simpler profile function (32) for balanced black rings, which has only a “uniform” circle as a source. From a complementary viewpoint, we can compare an expansion of the profile function (52) at large spatial distances with the analogous result (31) for the black ring. From Eq. (52) we obtain
$$\stackrel{~}{H}_{_{}}=\frac{1}{\sqrt{2}}\frac{8p_M}{3L}\left[3\frac{L}{r}\frac{5}{8}(3\mathrm{cos}^2\theta 1)\frac{L^3}{r^3}+\frac{7}{64}(35\mathrm{cos}^4\theta 30\mathrm{cos}^2\theta +3)\frac{L^5}{r^5}+O\left(\frac{L^7}{r^7}\right)\right].$$
(54)
The monopole term coincides with the one in the corresponding expression (31) for the black ring, which we should expect since we are boosting objects with the same mass (which scales as $`M=\gamma ^1p_M`$). However, Eqs. (54) and (31) in general differ already in the quadrupole term, in particular for the physically most interesting case of balanced rings $`p_\lambda =2p_\nu `$. They coincide only in the limiting case $`p_\nu =0`$, corresponding to $`\nu =0`$, when the black ring in fact reduces to a naked singularity isometric to that of Myers and Perry (see the discussion above about the Kerr bound). In addition, in the limit of vanishing rotation $`L=0`$ of Eq. (54) only the Aichelburg-Sexl monopole survives, which corresponds to the ultrarelativistic boost of the $`D=5`$ Schwarzschild-Tangherlini black hole.<sup>5</sup><sup>5</sup>5Recall that, instead, balanced black rings can not be static, while unbalanced static rings correspond to setting $`p_\nu =p_\lambda `$ in Eq. (31) (and not $`L=0`$).
### A.2 Parallel boost
For a parallel boost, the profile function of is
$`\stackrel{~}{H}_{_{||}}(x_1,x_2,y_2)=`$ $`{\displaystyle \frac{8\sqrt{2}p_M}{3\pi }}[{\displaystyle \frac{4}{a}}(1+{\displaystyle \frac{\eta ^2+y_2^2+L^2+a^2}{2Ly_2}})K(k)+{\displaystyle \frac{2a}{L^2}}E(k)`$ (55)
$`{\displaystyle \frac{2L+y_2}{a}}{\displaystyle \frac{\eta ^2+y_2^2+L^2+a^2}{L^2y_2}}\mathrm{\Pi }(\rho _1,k)],`$
with $`a`$ as in Eq. (36), $`k`$ as in Eq. (38) and
$$\rho _1=\frac{\eta ^2+y_2^2+L^2a^2}{2a^2}.$$
(56)
Similarly as in Sec. 5, the elliptic integrals are singular at $`k=1`$, i.e. on a rod of length $`2L`$ located at $`\eta =0`$, $`|y_2|L`$ . At large spatial distances, the expression (55) behaves as
$`\stackrel{~}{H}_{_{||}}=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{8p_M}{3L}}[3{\displaystyle \frac{L}{r}}+2\mathrm{cos}\theta {\displaystyle \frac{L^2}{r^2}}+{\displaystyle \frac{7}{8}}(3\mathrm{cos}^2\theta 1){\displaystyle \frac{L^3}{r^3}}+{\displaystyle \frac{3}{4}}(5\mathrm{cos}^3\theta 3\mathrm{cos}\theta ){\displaystyle \frac{L^4}{r^4}}`$ (57)
$`+{\displaystyle \frac{11}{64}}(35\mathrm{cos}^4\theta 30\mathrm{cos}^2\theta +3){\displaystyle \frac{L^5}{r^5}}+O\left({\displaystyle \frac{L^6}{r^6}}\right)].`$
The discussion is similar as the one above for $`\stackrel{~}{H}_{_{}}`$. Again, the monopole term coincides with the one in the corresponding expression (41) for the black ring. Higher multipoles in general differ, in particular for balanced rings. Boosted black holes reduce to the Aichelburg-Sexl monopole in the static limit $`L=0`$.
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# Security proof of practical quantum key distribution schemes
## Abstract
This paper provides a security proof of the Bennett-Brassard (BB84) quantum key distribution protocol in practical implementation. To prove the security, it is not assumed that defects in the devices are absorbed into an adversary’s attack. In fact, the only assumption in the proof is that the source is characterized. The proof is performed by lower-bounding adversary’s Rényi entropy about the key before privacy amplification. The bound reveals the leading factors reducing the key generation rate.
preprint: APS/123-QED
One of the fundamental problems in cryptography is to provide a way of sharing a secret random number between two parties, Alice and Bob, in the presence of an adversary Eve. The quantum key distribution is a solution to this problembe92 ; bb84 ; indeed it allows Alice and Bob to generate a shared secret key securely against Eve with unbounded resources of computation. The security of quantum key distribution against general attacks was first proved by Mayersma01 . Later, Shor-Preskillsp00 provided a simple security proof based on the observation that quantum key distribution (BB84 protocol) is closely related to quantum error-correcting codes (CSS codes). Gottesmann et al.gllp04 showed that the Shor-Preskill proof is still valid as long as the source and detector are perfect enough so that all defects can be absorbed into Eve’s attack (see also hm03 ; wmu04 for the rate achievability of quantum codes in the security proof). In contrast to the security proof based on quantum codes, the Mayers proof has a remarkable characteristics. Namely in the Mayers proof, although the source has to be (almost) perfect, there is no restriction on the detector; in particular, it can be uncharacterized. By exchanging the role of the source and detector in the Mayers proof, Koashi-Preskillkp03 provided a security proof which applies to the case where the detector is perfect, but the source can be uncharacterized (except that the averaged states are independent of Alice’s basis). The aim of this work is to generalize these results. We provide a security proof of the BB84 protocol in which the only assumption is that the source is characterized. In the same way as Koashi-Preskillkp03 , this can be transformed into a security proof which is based on characteristics of the detector. Further we note that the security proof also applies to the B92 protocolbe92 .
Let us first recall the BB84 protocolbb84 . Let $``$ be a Hilbert space. Let $`𝒜=\{1,\mathrm{},N\}`$, and for $`𝒜`$ denote the cardinality of $``$ by $`n_{}`$. The BB84 protocol is described as follows.
BB84 protocol: (i) Alice generates two binary strings $`a^𝒜=\{a_i\}_{i𝒜}`$ and $`x^𝒜=\{x_i\}_{i𝒜}`$ according to the probability distribution $`p(a^𝒜,x^𝒜)=_ip_{a_i,x_i}`$. (ii) Bob generates a binary string $`b^𝒜=\{b_i\}_{i𝒜}`$ according to the probability distribution $`p(b^𝒜)=_ip_{b_i}`$. (iii) Alice sends the quantum state on $`^N`$, $`\rho _{a,x}^𝒜=_{i𝒜}\rho _{a_i,x_i}`$, to Bob. (iv) Bob applies the measurement on $`^N`$, $`\{E_{b,y}^𝒜\}_{y^𝒜}=\left\{_{i𝒜}E_{b_i,y_i}\right\}_{y^𝒜\{0,1,\varphi \}^N}`$, to the received quantum state, where $`E_{0,\varphi }=E_{1,\varphi }`$ is the measurement corresponding to the result that Bob cannot detect a state. (v) Alice and Bob open $`a^𝒜`$ and $`b^𝒜`$ respectively. Let $`𝒟=\{i𝒜|y_i\varphi \}`$ and $`𝒞=\{i𝒟|a_i=b_i\}`$. Alice and Bob select a random subset $`𝒯𝒞`$ (which does not necessarily satisfy $`n_𝒯/n_𝒞1/2`$). Let $`𝒦=𝒞𝒯`$. (vi) Alice and Bob compare $`x^𝒯`$ and $`y^𝒯`$, and count the number of errors, $`n_𝒯^e=|\{i𝒯|x_iy_i\}|`$. (vii) Bob estimates $`x^𝒦`$ by exchanging error-correction information with Alice. (viii) Alice and Bob generate a secret key $`s`$ by applying a compression function to $`x^𝒦`$.
To prove the security of the BB84 protocol, the previous worksgllp04 ; kp03 ; ma01 ; sp00 assume that either Alice’s source or Bob’s detector is almost perfect in the sense that all defects in the device can be absorbed into Eve’s attack. We wish to prove the security of quantum key distribution under practical implementation. Note that the previous security proofs have been based on directly bounding Eve’s mutual information about the final key, i.e. the key after privacy amplification. In this work, we first lower-bound Eve’s Rényi entropy about the key before privacy amplification, and then apply privacy amplification in the classical information theory which makes use of a compression function in a universal hash family (see bbcm95 for the classical theory of privacy amplification).
We now provide basic definitions which will be used later (see e.g. ha05 for details). The variation distance between probability distributions $`p`$ and $`q`$ is given by $`d_V(p,q)=\frac{1}{2}_\omega |p(\omega )q(\omega )|`$. The quantum analogue of the variation distance is called the trace distance. For an Hermitian operator $`X`$ with the spectrum decomposition $`X=_ix_iE_i`$, define the projection $`\{X>0\}`$ by $`\{X>0\}=_{i:x_i>0}E_i`$. Then the trace distance between quantum states $`\rho `$ and $`\sigma `$, $`d_T(\rho ,\sigma )`$, is given by $`d_T(\rho ,\sigma )=\frac{1}{2}\mathrm{Tr}|\mathrm{\Delta }|=\frac{1}{2}\mathrm{Tr}(\mathrm{\Delta }\{\mathrm{\Delta }>0\}\mathrm{\Delta }\{\mathrm{\Delta }>0\})`$ with $`\mathrm{\Delta }=\rho \sigma `$. The trace distance can be bounded by another distance called the fidelity as $`d_T(\rho ,\sigma )\sqrt{1F(\rho ,\sigma )^2}`$, where the fidelity $`F(\rho ,\sigma )`$ between $`\rho `$ and $`\sigma `$ is given by $`F(\rho ,\sigma )=\mathrm{Tr}|\sqrt{\rho }\sqrt{\sigma }|`$.
Let $`z`$ be the output of the measurement by Eve. Then, without loss of generality, the probability distribution of the random variables can be written as
$`p_a^𝒞(x,y,z)`$ $`p(x^𝒞,y^𝒞,z|a^𝒜,b^𝒜,x^𝒯,y^𝒯,𝒟,𝒯)`$
$`=p_a^𝒞(x)\mathrm{Tr}(E_{a,y}^𝒞E_z)U(\rho _{a,x}^𝒞\rho _E)U^{}.`$
Here, $`\rho _E`$ is the initial state of an ancilla system $`_E`$ introduced by Eve, $`E_z`$ is the Eve’s measurement on the ancilla system, and $`U`$ is the Eve’s unitary operation acting on the composite system. (The quantum channel is assumed to be under Eve’s control). For $`𝒞`$ and $`p_a`$ as above, let $`p_a^{}`$ denote the marginal distribution of the random variables defined on $``$.
We begin with decomposing $`\rho _\alpha `$ ($`a,x\{0,1\}`$) as
$$\rho _{a,x}=p_{a,x}^{(0)}\rho _{a,x}^{(0)}+p_{a,x}^{(1)}\rho _{a,x}^{(1)},\rho _{a,x}^{(0)},\rho _{a,x}^{(1)}𝒮(),$$
(1)
where $`p_{a,x}^{(0)}+p_{a,x}^{(1)}=1`$, $`p_{a,x}p_{a,x}^{(0)}=p^{(0)}`$ for a positive constant $`p^{(0)}\mathrm{min}_{a,x}\{p_{a,x}\}`$, and $`\rho _{a,x}^{(0)}`$ has a Schatten decomposition of the form
$$\rho _{a,x}^{(0)}=\underset{k_{a,x}}{}\lambda _{a,x}(k_{a,x})|k_{a,x}k_{a,x}|.$$
(2)
We note that $`\rho _{a,x}`$ always has a decomposition of the above form (where we allow $`\rho _{a,x}^{(0)}=\rho _{a,x}^{(1)}`$). Let $`𝒳=\{(0,0),(0,1),(1,0),(1,1)\}`$. We now construct a set of pure states, $`\{\widehat{\rho }_\alpha \}_{\alpha 𝒳}`$, such that there exists a physical transformation from $`\{\widehat{\rho }_\alpha \}_{\alpha 𝒳}`$ to $`\{\rho _\alpha ^{(0)}\}_{\alpha 𝒳}`$. Let $`\mu _{\alpha \beta }`$ ($`\alpha ,\beta 𝒳`$) be a mapping from $`\{|k_\alpha \}_{k_\alpha }`$ to $`\{|k_\beta \}_{k_\beta }`$ with $`\mu _{\alpha \alpha }`$ being the identity on $`\{|k_\alpha \}_{k_\alpha }`$, and introduce the Gram matrix $`G`$ by writing
$$[G]_{\alpha \beta }=\underset{k_\alpha }{}\sqrt{\lambda _\alpha (k_\alpha )\lambda _\beta (k_{\alpha \beta })}k_\alpha |k_{\alpha \beta }\varphi _{k_\alpha }|\varphi _{k_{\alpha \beta }},$$
where $`|k_{\alpha \beta }=\mu _{\alpha \beta }(|k_\alpha )`$ and $`|\varphi _{k_\alpha }`$ is a state on an ancilla system $`_\varphi `$. Since $`G0`$, there exists a square matrix $`C`$ such that $`G=C^{}C`$. Further, since all the diagonal elements of $`G`$ are 1, we can define a pure state $`\widehat{\rho }_\alpha `$ ($`\alpha 𝒳`$) on a 4-dimensional Hilbert space $`_4`$ by
$$\widehat{\rho }_\alpha =|C_\alpha C_\alpha |,$$
where $`C_\alpha `$ denotes the $`\alpha `$-th column of $`C`$. It follows from this construction that there exists a physical transformation from $`\{\widehat{\rho }_\alpha \}_{\alpha 𝒳}`$ to $`\{\rho _\alpha ^{(0)}\}_{\alpha 𝒳}`$ (see cjw03 ). Now we introduce an approximation of $`\{\widehat{\rho }_\alpha \}_{\alpha 𝒳}`$ which is easier to treat in the security proof. Let $`_2`$ be a 2-dimensional subspace of $`_4`$, and $`\sigma _{a,x}`$ ($`a,x\{0,1\}`$) be states on $`_2`$ such that
$$\sigma _{0,0}+\sigma _{0,1}=\sigma _{1,0}+\sigma _{1,1}=I__2,$$
where, for a Hilbert space $``$, $`I_{}`$ denotes the identity on $``$. Note that the decompositions (1) and (2) and the choises of $`\mu _{\alpha \beta }`$, $`|\varphi _{k_\alpha }`$ and $`\sigma _{a,x}`$ are not unique; they should be determined so that the distance $`d_T(\sigma _{a,x},\widehat{\rho }_{a,x})`$ will be minimized. In the case of coherent states with no phase reference, $`\rho _\alpha =_k(\mu ^k/k!)e^\mu |k;\alpha k;\alpha |`$, for instance, we can take for $`\alpha ,\beta 𝒳`$ and $`k`$, $`\rho _\alpha ^{(0)}=\widehat{\rho }_\alpha =\sigma _\alpha =|1;\alpha 1;\alpha |`$, $`p_\alpha ^{(0)}=\mu e^\mu `$, $`\mu _{\alpha \beta }(|k;\alpha )=|k;\beta `$ and $`|\varphi _{k;\alpha }=|\varphi `$.
The decomposition (1) allows us to consider that the Alice’s source generates $`\rho _{a,x}^{(0)}`$ with probability $`p_{a,x}^{(0)}`$ and $`\rho _{a,x}^{(1)}`$ with probability $`p_{a,x}^{(1)}`$. Further, we assume that Eve is informed of partial information about each state $`\rho _A`$ generated by the Alice’s source: (i) $`\rho _A=\rho _{a,x}^{(0)}`$ or $`\rho _A=\rho _{a,x}^{(1)}`$ and (ii) $`\rho _A=\rho _{0,x}^{(1)}`$ or $`\rho _A=\rho _{1,x}^{(1)}`$ when $`\rho _A=\rho _{a,x}^{(1)}`$. This assumption is advantageous to Eve, and hence does not reduce the security of the protocol. Let $`𝒦`$ be the positions where $`\rho _{a,x}^{(0)}`$ is generated, and $`=𝒦`$. We now fix $``$ and $``$, and consider the best success probability to estimate $`x^{}`$ from $`\rho _{a,x}^{(1)}`$ and $`a^{}`$. Here note that we can estimate each bit $`x_i`$ of $`x^{}`$ separately because each state $`\rho _{a_i,x_i}`$ is generated independently of the other bits $`\{x_i^{}|i^{}i,i^{}\}`$. For $`a\{0,1\}`$, let $`\{T_{a,0},T_{a,1},T_{a,\varphi }\}`$ be a POVM on $``$ which is used to discriminate $`\rho _{a,0}^{(1)}`$ and $`\rho _{a,1}^{(1)}`$, and let $`p_a^{(1)}`$ be the conditional probability defined by $`p_a^{(1)}=(p_{a,0}p_{a,0}^{(1)}+p_{a,1}p_{a,1}^{(1)})/(p_{a,0}+p_{a,1})`$. Further, define for a constant $`\delta _{}^a>0`$,
$`p_{}^a`$ $`=(p_{}^a\delta _{}^a){\displaystyle \frac{n_𝒟^a}{n_𝒦^ap_a^{(1)}}},p_{}^a={\displaystyle \frac{n_{}^a}{n_𝒜^a}},`$
$`ϵ_{}^a`$ $`=\mathrm{exp}(n_𝒜^aD(B_1(p_{}^a)||B_1(p_{}^a\delta _{}^a))),`$
where $`n_{}^d=|\{i|a_i=d\}|`$ for $`𝒜`$, $`B_1`$ denotes the Bernoulli distribution, and $`D(p||q)`$ is the relative entropy of $`p`$ and $`q`$<sup>1</sup><sup>1</sup>1In the case where Eve is allowed to collapse Bob, $`n_{}^a`$ should be replaced by $`n_𝒞^a`$.. Here let us consider the condition $`C`$ given by
$$C:\underset{x,x^{}}{}\mathrm{Tr}\widehat{p}_{a,x}^{(1)}\rho _{a,x}^{(1)}T_{a,x^{}}p_{}^a,$$
where $`\widehat{p}_{a,x}^{(d)}=p_{a,x}p_{a,x}^{(d)}/(p_{a,0}p_{a,0}^{(d)}+p_{a,1}p_{a,1}^{(d)})`$ for $`d\{0,1\}`$. Then it can be verified that $`\mathrm{Pr}_𝒜[\neg C]ϵ_{}^a`$, where the probability $`\mathrm{Pr}_𝒜`$ is taken over the randomness in choosing $`𝒟,𝒯,𝒜`$ (see e.g. ck81 ). Suppose now that the condition $`C`$ holds. Then we have
$$n_{}^an_{}^{a+}\underset{}{\mathrm{max}}\{n_{}^a|p_{}^a1\}.$$
Also, we can write the best success probability of the discrimination as
$$s_{}^a=\underset{T_{a,0},T_{a,1}:C}{sup}\left\{\frac{_x\mathrm{Tr}\widehat{p}_{a,x}^{(1)}\rho _{a,x}^{(1)}T_{a,x}}{_{x,x^{}}\mathrm{Tr}\widehat{p}_{a,x}^{(1)}\rho _{a,x}^{(1)}T_{a,x^{}}}\right\}.$$
Let $`z^{}`$ be a random variable induced by a measurement on $`\rho _{a,x}^𝒦`$. Then, by definition of $`s_{}^a`$, it follows that
$$p_a^𝒦(x|z^{})p_a^{}(x|z^{})(s_{}^0)^{n_{}^0}(s_{}^1)^{n_{}^1}.$$
(3)
Having considered the $``$ part, we next consider the $``$ part. Let us first estimate the error rate $`p_{}^e`$ at $``$ from $`p_𝒯^e=n_𝒯^e/n_𝒯`$, the error rate at $`𝒯`$. On remembering that the error probability of the discrimination at $``$ is at least $`1s_{}^a`$ for a basis $`a`$, define for a constant $`\delta _p>0`$,
$`p_{}^+`$ $`={\displaystyle \frac{n_𝒦p_𝒯^e+n_𝒞\delta _pn_{}^0(1s_{}^0)n_{}^1(1s_{}^1)}{n_{}}},`$
$`ϵ_𝒯^e`$ $`=\mathrm{exp}(n_𝒯D(B_1(p_𝒯^e)||B_1(p_𝒯^e+\delta _p)).`$
Then we have $`\mathrm{Pr}_𝒜[p_{}^e>p_{}^+]\mu _{}ϵ_0^{}+ϵ_1^{}+ϵ_𝒯^e`$, from which, it follows that
$$\underset{x,y,z:|xy|>n_{}p_{}^+}{}p_a^{}(x,y,z)\mu _{}.$$
(4)
Now, let us consider a modified protocol in which Alice sends $`\widehat{\rho }_{\overline{a},x}^{}`$ (instead of $`\widehat{\rho }_{a,x}^{}`$), where $`\overline{a}`$ denotes the bit-wise inversion of binary string $`a`$. Let $`p_{\stackrel{~}{a}}`$ be the corresponding conditional probability in the modified protocol. It then follows from the monotonicity of the trace distance that
$`d_V(p_a^𝒯(x,\stackrel{~}{y},z),p_{\stackrel{~}{a}}^𝒯(x,\stackrel{~}{y},z))d_T(\overline{\rho }_a^{},\overline{\rho }_{\overline{a}}^{}),`$ (5)
where $`\overline{\rho }_a=\frac{1}{2}_x\widehat{\rho }_{a,x}`$ for $`a\{0,1\}`$. We note that $`d_T(\overline{\rho }_a^{},\overline{\rho }_{\overline{a}}^{})`$ can be bounded as $`d_T(\overline{\rho }_a^{},\overline{\rho }_{\overline{a}}^{})\sqrt{1F(\overline{\rho }_0,\overline{\rho }_1)^{2n_{}}}`$. From inequalities (4) and (5), it follows that
$$\underset{x,y,z:|xy|>n_{}p_{}^+}{}p_{\stackrel{~}{a}}^{}(x,y,z)\mu _{}+d_T(\overline{\rho }_a^{},\overline{\rho }_{\overline{a}}^{}).$$
(6)
Let us now introduce the POVM $`\{M_{a,yz}\}_{y,z}`$ by writing
$`p_a^{}(x,y,z)=\mathrm{Tr}\widehat{p}_a^{}(x)\widehat{\rho }_{a,x}^{}M_{a,yz}`$
with $`\widehat{p}_a^{}(x)=_i\widehat{p}_{a_i,x_i}^{(0)}=2^n_{}`$, where, for simplicity, we have omitted deviding the right-hand side by $`_{y,z}\mathrm{Tr}\overline{\rho }_a^{}M_{a,yz}`$ because it will be canceled when we will consider the conditional probability $`\widehat{p}_a^{}(x|\stackrel{~}{y},z)`$. Now, let us consider the case where Bob uses the opposite basis $`\overline{a}`$ at $``$ and introduce the notation $`\stackrel{~}{y}`$ by writing
$`p_a^{}(x,\stackrel{~}{y},z)`$ $`=\mathrm{Tr}\widehat{p}_a^{}(x)\widehat{\rho }_{a,x}^{}M_{\overline{a},yz}.`$
Note that $`E_{0,0}+E_{0,1}=E_{1,0}+E_{1,1}`$, and so $`p_a^{}(x,z)=_yp_a^{}(x,y,z)=_yp_a^{}(x,\stackrel{~}{y},z)`$. That is, the probability distribution $`p_a^{}(x,z)`$ is independent of the basis used for the Bob’s measurement. Thus, in the sequel, we will consider $`p_a(x,\stackrel{~}{y},z)`$ rather than $`p_a(x,y,z)`$.
To examine the security of the protocol, it is more convenient to treat $`\sigma _{a,x}`$ than $`\widehat{\rho }_{a,x}`$. Thus, define
$`\widehat{p}_a^{}(x,\stackrel{~}{y},z)=\mathrm{Tr}\widehat{p}_a^{}(x)\sigma _{a,x}^{}M_{\overline{a},yz}.`$
The monotonicity of the trace distance gives
$$\begin{array}{cc}& d_V(p_a^{}(x,\stackrel{~}{y},z),\widehat{p}_a^{}(x,\stackrel{~}{y},z))\nu _{},\hfill \\ & \nu _{}\underset{x}{}\widehat{p}_a^{}(x)d_T(\widehat{\rho }_{a,x}^{},\sigma _{a,x}^{}).\hfill \end{array}$$
(7)
This, together with (6), yields
$$\underset{y,z}{}\mathrm{Tr}(\overline{\sigma }_{\overline{a}}^{}\overline{\sigma }_{\stackrel{~}{y}})M_{\overline{a},yz}\mu _{}+\nu _{}+d_T(\overline{\rho }_a^{},\overline{\rho }_{\overline{a}}^{}),$$
(8)
where we have defined
$$\overline{\sigma }_{\stackrel{~}{y}}=\underset{x^{}:|xy|n_{}p_{}^+}{}\widehat{p}_a^{}(x)\sigma _{\overline{a},x}^{}.$$
Inequality (8) can be seen as a restriction on Eve’s measurement. To take advantage of this restriction, we now construct a projection on $`^n_{}`$, $`P_{\stackrel{~}{y}}`$, which sufficiently preserves $`\overline{\sigma }_{\stackrel{~}{y}}`$. For this purpose, let us first consider the problem of quantum hypothesis testing, where two hypotheses are, for fixed base $`a\{0,1\}`$, $`H_0:\rho =\sigma _{a,0}_2`$ and $`H_1:\rho =\sigma _{a,1}_2`$. If $`\{P_{a,x}\}_{x\{0,1\}}`$, defined by
$$P_{a,x}=\{\sigma _{a,x}\sigma _{a,\overline{x}}>0\},$$
is used as a test for the hypothesis testing, then the success probability $`s_{}^a`$ is given by
$`s_{}^a={\displaystyle \frac{1}{2}}(1+d_T(\sigma _{a,0},\sigma _{a,1})).`$
Suppose now that we receive a product state $`\sigma _{a,x}^{}`$ from the Alice’s source, and estimate $`x^{}`$ by applying the above hypothesis testing to each individual state. Let $`k`$ be an integer such that $`0kn_{}`$. If we allow up to $`k`$ errors in the estimation of $`n_{}`$-bit string $`x^{}`$, then the error probability $`ϵ^P`$ (i.e. the probability that we make more than $`k`$ errors) can be bounded as
$$ϵ^P\left(2^n_{}\frac{2^{n_{}h(\frac{k}{n_{}})}}{2\sqrt{n_{}}}\right)(s_{}^0)^{\overline{n}_{}^0}(s_{}^1)^{\overline{n}_{}^1}\left(\frac{1s_{}^m}{s_{}^m}\right)^k,$$
where $`s_{}^m=\mathrm{min}\{s_{}^0,s_{}^1\}`$, $`\overline{n}_{}^a=n_{}n_{}^a`$, and we have used, for $`0kn`$ and $`0q1`$,
$$\frac{2^{nh(\frac{k}{n})}}{2\sqrt{n}}\underset{i=0}{\overset{k}{}}\left(\begin{array}{c}n\\ i\end{array}\right)q^i(1q)^{ni}2^{nh(\frac{k}{n})},$$
with $`h(p)=p\mathrm{log}p(1p)\mathrm{log}(1p)`$ (see e.g. ck81 ). We are now in position to construct $`P_{\stackrel{~}{y}}`$. Let $`\delta _P=\frac{k}{n_{}}`$ and $`p^{}=p_{}^++\delta _P`$. Define the projection $`P_{\stackrel{~}{y}}`$ on $`^{}`$ by
$$P_{\stackrel{~}{y}}=\underset{x^{}:|xy|n_{}p^{}}{}\underset{i}{}P_{\overline{a}_i,x_i}.$$
Then it can be verified that $`\mathrm{Tr}\overline{\sigma }_{\stackrel{~}{y}}(I_{^{}}P_{\stackrel{~}{y}})ϵ^P\mathrm{Tr}\overline{\sigma }_{\stackrel{~}{y}}`$, which shows that $`P_{\stackrel{~}{y}}`$ is a required projection (provided that $`1s_{}^m`$ is sufficiently small).
Having constructed the projection $`P_{a,x}`$, we now bound the conditional probability $`\widehat{p}_a^{}(x|\stackrel{~}{y},z)`$. Since
$$\widehat{p}_a^{}(\stackrel{~}{y},z)=\mathrm{Tr}\overline{\sigma }_a^{}M_{\stackrel{~}{y}z}=\overline{\pi }_{}2^n_{}\mathrm{Tr}M_{\stackrel{~}{y}z}$$
with $`M_{\stackrel{~}{y}z}=M_{\overline{a},yz}`$ for short, we now bound $`\widehat{p}_a^{}(x,\stackrel{~}{y},z)`$. It follows, on using $`\mathrm{Tr}P_{\stackrel{~}{y}}2^{n_{}h(p^{})}`$, that
$`\mathrm{Tr}P_{\stackrel{~}{y}}p_a^{}(x)\rho _{a,x}^{}P_{\stackrel{~}{y}}M_{\stackrel{~}{y}z}\pi _{},`$
$`\pi _{}2^{n_{}+n_{}h(p^{})+\overline{n}_{}^0\mathrm{log}q_0+\overline{n}_{}^1\mathrm{log}q_1}\mathrm{Tr}M_{\stackrel{~}{y}z},`$
where, for $`a\{0,1\}`$, $`q_a=\mathrm{max}_{x,x^{}\{0,1\}}\{\mathrm{Tr}\sigma _{a,x}P_{\overline{a},x^{}}\}.`$ Define now
$$\widehat{p}_a^{}(x,\stackrel{~}{y},z)=\mathrm{Tr}(I_{^{}}P_{\stackrel{~}{y}})\widehat{p}_a^{}(x)\sigma _{a,x}^{}(I_{^{}}P_{\stackrel{~}{y}})M_{\stackrel{~}{y}z}.$$
Since $`\overline{\sigma }_a^{}=\overline{\sigma }_{\overline{a}}^{}=(\overline{\sigma }_{\overline{a}}^{}\overline{\sigma }_{\stackrel{~}{y}})+\overline{\sigma }_{\stackrel{~}{y}}`$, $`P_{\stackrel{~}{y}}`$ and $`\overline{\sigma }_{\stackrel{~}{y}}`$ commute, and $`_y\mathrm{Tr}\overline{\sigma }_{\stackrel{~}{y}}2^{n_{}h(p_{}^+)}`$, we have
$`{\displaystyle \underset{x,y,z}{}}\widehat{p}_a^{}(x,\stackrel{~}{y},z)={\displaystyle \underset{x,y,z}{}}\widehat{p}_a^{}(x,\stackrel{~}{y},z){\displaystyle \frac{\widehat{p}_a^{}(x,\stackrel{~}{y},z)}{\widehat{p}_a^{}(x,\stackrel{~}{y},z)}}\omega _{},`$
$`\omega _{}\mu _{}+\nu _{}+d_T(\overline{\rho }_a^{},\overline{\rho }_{\overline{a}}^{})+2^{n_{}h(p_{}^+)}ϵ^P.`$
Hence Markov’s inequality for a constant $`c>0`$ yields
$$\mathrm{Pr}_{\widehat{p}_a}[\widehat{p}_a^{}(x,\stackrel{~}{y},z)c\omega _{}\widehat{p}_a^{}(x,\stackrel{~}{y},z)]1c^1,$$
where $`c`$ should be determined so that Eve’s mutual information about the final key will be minimized. Further, Schwarz’s inequality gives
$`\mathrm{Tr}P_{\stackrel{~}{y}}\widehat{p}_a^{}(x)\sigma _{a,x}^{}(I_{^{}}P_{\stackrel{~}{y}})M_{a,yz}`$ $`(\pi _{}\widehat{p}_a^{}(x,\stackrel{~}{y},z))^{\frac{1}{2}}.`$
Therefore it follows that
$`\widehat{p}_a^{}(x,\stackrel{~}{y},z)`$ $`\left((\pi _{})^{\frac{1}{2}}+(\widehat{p}_a^{}(x,\stackrel{~}{y},z))^{\frac{1}{2}}\right)^2,`$
and so
$$\mathrm{Pr}_{\widehat{p}_a}[\widehat{p}_a^{}(x|\stackrel{~}{y},z)>\mathrm{\Pi }_{}]\frac{1}{c},\mathrm{\Pi }_{}\frac{\pi _{}}{\overline{\pi }_{}\left(1(c\omega _{})^{\frac{1}{2}}\right)^2}.$$
(9)
Now, it follows from inequality (3) that the conditional Rényi entropy $`R_a^𝒦(X|\stackrel{~}{y}^{},z)`$ can be bounded as
$`R_a^𝒦(X|\stackrel{~}{y}^{},z)`$ $`\mathrm{log}{\displaystyle \underset{x^𝒦}{}}\left(p_a^𝒦(X=x|\stackrel{~}{Y}=\stackrel{~}{y}^{},Z=z)\right)^2`$
$`R_a^{}(X|\stackrel{~}{y},z)+R_a^{},`$
where $`R_a^{}=_an_{}^a\mathrm{log}s_{}^a`$, and a capital letter (say $`X`$) denotes the random variable which samples the corresponding small letter (say $`x`$). Now, using constraints (7) and (9), let us derive another constraint of the form
$$\mathrm{Pr}_{p_a}[R_a^{}(X|\stackrel{~}{y},z)>R_a^{}]ϵ_{}.$$
If $`\nu _{}=0`$, for example, we can take $`R_a^{}=\mathrm{log}\mathrm{\Pi }_{}`$ and $`ϵ_{}=c^1`$. Define
$`R_E^𝒦=\underset{:n_{}^an_{}^{a+}}{\mathrm{min}}\{R_a^{}+R_a^{}\},`$
and let $`m`$ be an integer such that $`lR_E^𝒦m>0`$. Choose a function $`g`$ at random from a universal family of hash functions from $`\{0,1\}^n`$ to $`\{0,1\}^m`$. If Alice and Bob choose $`s=g(x^𝒦)`$ as their secret key, then the Eve’s expected information about $`S`$, given $`Z`$ and $`G`$, satisfies $`I(S:Z,G)n_{}ϵ_{}+2^l/\mathrm{ln}2`$, where we consider $`\stackrel{~}{Y}`$ as an auxiliary random variable (see bbcm95 for details). Here we note that $`R_E^𝒦`$ is not explicitly dependent on the characteristics of the detector, and hence the detector can be uncharacterized. Further, as $`n_{}\mathrm{}`$, the terms $`\nu _{}`$ and $`d_T(\overline{\rho }_a^{},\overline{\rho }_{\overline{a}}^{})`$ approach to 1 unless $`\widehat{\rho }_{a,x}=\sigma _{a,x}`$ and $`\overline{\rho }_a^{}=\overline{\rho }_{\overline{a}}^{}`$; this shows that the leading factors reducing the key generation rate are the asymmetries of the source represented by these terms.
To see that our result is consistent with the previous ones, suppose that the source and detector are perfect. In this case, we can take $`\rho _{a,x}^{(0)}=\sigma _{a,x}=\rho _{a,x}`$, $`=𝒦`$, $`\mu _{}=ϵ_𝒯^e`$, $`\nu _{}=0`$, $`d_T(\overline{\rho }_a^{},\overline{\rho }_{\overline{a}}^{})=0`$, $`\mathrm{log}q_a=1`$, $`\delta _P=0`$, $`ϵ^P=0`$. Since $`\omega _{}=ϵ_𝒯^e0`$ as $`n_𝒦\mathrm{}`$ for fixed $`\delta _p`$, $`R_E^𝒦/n_𝒦`$ approaches to $`h(p_𝒯^e)`$ for sufficiently small $`c^1`$ and $`\delta _p`$. This is consistent with the results in the previous worksgllp04 ; kp03 ; ma01 ; sp00 <sup>2</sup><sup>2</sup>2The original bound given by Mayersma01 is slightly weaker (but this can be improved by a minor modification)..
We close this paper with mentioning some extensions of this work. (i) In the same way as Koashi-Preskillkp03 , we can provide a security proof of the BB84 protocol where the only assumption is that the detector and basis dependence of the averaged states are characterized. (ii) It is also of importance to give a security proof of the B92 protocolbe92 . Suppose that the source generates $`\rho _0`$ with probability $`p_0`$ and $`\rho _1`$ with probability $`p_1`$. Then we decompose $`\rho _a`$ ($`a\{0,1\}`$) as $`\rho _a=p_0^{(0)}\rho _a^{(0)}+p_1^{(1)}\rho _a^{(1)}`$ so that $`p_0p_0^{(0)}=p_1p_1^{(1)}`$. Again we define $`\widehat{\rho }_a`$ by introducing the Gram matrix as above. Note that $`\widehat{\rho }_a`$ is a pure state on a 2-dimensional Hilbert space $`_2`$. Hence, the terms $`\nu _{}`$ and $`d_T(\overline{\rho }_a^{},\overline{\rho }_{\overline{a}}^{})`$ automatically vanish in this case, which could be considered as an advantage of the B92 protocol. More detailed investigation concerning these extensions will be the subject of future work.
The author is grateful to Dr. Keiji Matsumoto for comments. This work was supported in part by MEXT, Grant-in-Aid for Encouragement of Young Scientists (B) No. 15760289.
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# Zn and Cr abundances in damped Lyman alpha systems from the CORALS survey This work is based in part on observations collected at the European Southern Observatory, Chile (ESO Nos. 69.A-0051A & 71.A-0067A) and at the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation.
## 1 Introduction
Spectroscopic studies of absorption line systems along quasar sightlines are an important source of information regarding the chemical evolution history of neutral gas in the universe. Damped Ly$`\alpha `$ systems (DLAs, defined to have N(H i) $`2\mathrm{x}10^{20}`$ atoms cm<sup>-2</sup>; Wolfe et al. wolfe86 (1986)), which make up the high column density end of the distribution of absorption line systems, are a particularly interesting subset of absorbers to study. They dominate the neutral hydrogen content of the universe available for star formation up to $`z4`$ (e.g. Storrie-Lombardi & Wolfe storrie00 (2000); Peroux et al. peroux03 (2003)) and are thought to be the progenitors of the disk galaxies we see today. DLAs have therefore been used as a tool for probing the evolution of galaxies, especially at high redshift where direct measurements are more difficult. Wolfe et al. (wolfe95 (1995)) found that the comoving mass density in neutral hydrogen in DLAs, $`\mathrm{\Omega }_{\mathrm{DLA}}`$, is similar to the mass density in stars at redshift $`z=0`$, and therefore proposed that its redshift evolution, d$`\mathrm{\Omega }_{\mathrm{DLA}}`$/d$`z`$, could be taken as a global measure of the rate at which gas in galaxies is converted into stars.
The metallicity evolution of galaxies has also been probed by chemical abundance measurements in DLAs (e.g. Pettini et al. 1997b , pettini99 (1999); Prochaska & Wolfe prochaska99 (1999), prochaska00 (2000)), which provide the most detailed information on the chemical composition of high-$`z`$ galaxies. If indeed DLAs are representative of the galaxy population at a given redshift, then one would expect that the evolution of their neutral hydrogen fraction and metallicity would track that of the general galaxy population as a whole. However, no strong evolution of either $`\mathrm{\Omega }_{\mathrm{DLA}}`$ or metallicity ($`\mathrm{Z}_{\mathrm{DLA}}`$) has been seen in studies from $`z3.5`$ down to $`z0.5`$ (e.g. Rao & Turnshek rao00 (2000); Ryan-Weber et al. ryan-weber03 (2003); Kulkarni et al. kulkarni05 (2005)).
Furthermore, it is now well established that DLAs are generally at the low end of the metallicity distribution of galaxies at redshifts $`z=23`$ (see, for example, Figure 32 of Pettini pettini04 (2004)). The typical DLA metallicity at this epoch is only $`Z_{\mathrm{DLA}}1.2`$, or $`1/15`$ of solar (Pettini et al. pettini99 (1999); Kulkarni et al. kulkarni05 (2005)), while near-solar metallicities are common for luminous galaxies detected directly in their rest-frame ultraviolet, optical, and far-infrared light (e.g. Pettini et al. pettini02 (2002); Shapley et al. shapley04 (2004); de Mello et al. de\_mello04 (2004); Swinbank et al. swinbank05 (2005)).
This difference could have a number of causes. Selecting galaxies through H i absorption may preferentially pick out chemically unevolved systems, either because galaxies with generally low rates of star formation dominate the cross-section (Mo, Mao & White mo98 (1998)), or because the H i cross-section is largest during a stage prior to the onset of star formation. However, direct imaging of DLA galaxies shows a varied population of hosts which span a range of luminosities and morphological types (Boissier, Péroux, & Pettini boissier03 (2003); Rao et al. rao03 (2003); Chen & Lanzetta chen03 (2003); Weatherley et al. weatherley05 (2005)—see also Ryan-Weber et al. ryan-weber03 (2003)). Abundance gradients may contribute to the difference in metallicity between the outer regions—which offer the larger cross-section for absorption—and the inner regions of galaxies where star formation activity is more prominent (Pettini et al. 1994a ; Chen, Kennicutt, & Rauch chen05 (2005); Christensen et al. christensen05 (2005); Ellison, Kewley, & Mallén-Ornelas 2005b ). Quantitatively, however, the magnitude of such gradients has recently been questioned in nearby spirals (Bresolin, Garnett, & Kennicutt bresolin04 (2004)), and remains unknown at high redshifts.
A third possibility is that damped Ly$`\alpha `$ systems which are both metal- and gas-rich do exist, but are systematically underrepresented in current samples drawn from magnitude-limited QSO surveys. The hypothesis is that even moderate amounts of dust associated with intervening galaxies may be sufficient to preferentially exclude reddened QSOs from optical surveys, and that the statistics of DLAs would accordingly be skewed by such bias against dusty absorbers. It is this third possibility which we address in the present paper.
The idea of dust obscuration of QSOs has a long history (e.g. Heisler & Ostriker heisler88 (1988)) and received observational support by the work of Pei, Fall, & Bechtold (pei91 (1991)) who found the spectra of QSOs with DLAs in the redshift range $`1.77z_{\mathrm{abs}}2.80`$ to have statistically steeper continuum slopes than those of a control sample. On the basis of these results, Fall & Pei (fall93 (1993)) proposed that between 10% and 70% of bright QSOs at $`z=3`$ may have been missing from optical samples. This claim is now tempered by the recent re-analysis by Murphy & Liske (murphy04 (2004)) based on the much larger compilation of QSO spectra made available by the Sloan Digital Sky Survey (SDSS—Stoughton et al. stoughton02 (2002)). From the comparison of 70 QSOs lying behind DLAs at $`2.0<z_{\mathrm{abs}}<4.0`$) with a control sample which is one order of magnitude larger, Murphy & Liske concluded that the difference in the continuum slopes $`\alpha `$ between the two sets of spectra is only $`\mathrm{\Delta }\alpha =0.04\pm 0.05`$ (using the usual definition of the spectral index, whereby the QSO continuum flux is a power law of the form $`f_\nu \nu ^\alpha `$), corresponding to a limit on the colour excess due to SMC-like dust-reddening of $`E(BV)<0.02\mathrm{mag}`$ ($`3\sigma `$). This value is significantly lower than $`\mathrm{\Delta }\alpha =0.38\pm 0.13`$ reported by Pei et al. (pei91 (1991)). Similarly, Ellison, Hall & Lira (2005a ) find $`E(BV)<0.05\mathrm{mag}`$ ($`3\sigma `$) from a study of the optical to infrared colours of a subsample of CORALS QSOs with DLAs in the range $`1.8<z_{\mathrm{abs}}<3.5`$.
However, the SDSS results still do not preclude the possibility that highly obscured QSOs may by missing, or underrepresented in optical samples (since SDSS QSOs are optically selected). Theoretically, such selection effects have been appealed to in order to reconcile the predictions of hydrodynamic simulations (Cen et al. cen03 (2003); Churches, Nelson & Edmunds churches04 (2004); Nagamine, Springel, & Hernquist nagamine04 (2004)) and galactic chemical evolution models (Prantzos & Boissier prantzos00 (2000)) with the observations. Observationally, their importance is suggested by the apparent anti-correlation between neutral hydrogen column density $`N`$(H i) and metallicity $`Z_{\mathrm{DLA}}`$ first pointed out by Boissé et al. (boisse98 (1998)).
It is to assess quantitatively the importance of dust-induced bias for the statistics of DLAs that the Complete Optical and Radio Absorption Line System (CORALS) survey was originally conceived. As the name implies, this programme aims at measuring the properties of DLAs in a complete sample of QSOs selected at radio wavelengths, where dust obscuration is not expected to be an issue. In the first stage of the project, Ellison et al. (2001b ) identified a sample of 22 DLAs from intermediate dispersion spectroscopy of all the QSOs (66) with emission redshift $`z_{\mathrm{em}}2.2`$ in the Parkes quarter-Jansky sample of flat-spectrum radio sources (Jackson et al. jackson02 (2002); Hook et al. hook03 (2003)). The optical spectra were of sufficient quality to measure $`N`$(H i) in the 22 DLAs, enabling Ellison et al. (2001b ) to determine both the number density of DLAs per unit redshift, $`n(z)`$, and the corresponding comoving mass density of neutral gas, $`\mathrm{\Omega }_{\mathrm{DLA}}`$. The values found, $`n(z)=0.31_{0.08}^{+0.09}`$ and $`\mathrm{log}\mathrm{\Omega }_{\mathrm{DLA}}=2.59_{0.24}^{+0.17}`$ at a mean redshift $`<z>=2.37`$, are higher than the corresponding quantities previously determined from optically selected QSO samples, but only marginally so. In particular, the CORALS survey did not uncover a population of high column density ($`\mathrm{N}(\mathrm{H}\text{i})>10^{21}\mathrm{cm}^2`$) DLAs in front of faint QSOs. Within the limitations imposed by the small size of their sample, Ellison et al. (2001b ) concluded that selection effects due to intervening dust may at most account for an underestimate by a factor of $`2`$ in $`\mathrm{\Omega }_{\mathrm{DLA}}`$.
The H i results alone, however, do not tell us about the metal and dust content of CORALS DLAs and whether they are higher, on average, than those of the optically selected DLAs which have been studied extensively over the last fifteen years. These are the questions which we explore in the present work.
Specifically, we have conducted a follow-up programme of high resolution spectroscopy of the CORALS DLAs aimed at measuring in particular the abundances of zinc and chromium. Meyer, Welty & York (meyer89 (1989)) and Pettini, Boksenberg & Hunstead (pettini89 (1989), pettini90 (1990)) first drew attention to the diagnostic value of these two elements. Both are iron-peak elements, whose abundances track that of Fe to within $`\pm 0.10.2\mathrm{dex}`$ in Galactic stars of metallicities from solar to about 1/100 of solar (Chen, Nissen & Zhao chen04 (2004); Cayrel et al. cayrel04 (2004) and references therein). In the interstellar medium of the Milky Way, on the other hand, Zn is one of the few elements which show little affinity for dust grains, unlike Cr which is usually highly depleted (Savage & Sembach savage96 (1996)). In combination, therefore, these two elements can be used to obtain approximate measures of the overall degree of metal enrichment, via the \[Zn/H\] ratio, and the fraction of refractory elements locked up in solid form, via the \[Cr/Zn\] ratio.<sup>1</sup><sup>1</sup>1We use the conventional notation whereby \[Zn/H\]$`=\mathrm{log}(\mathrm{Zn}/\mathrm{H})\mathrm{log}(\mathrm{Zn}/\mathrm{H})_{}`$. Both elements have absorption lines of their dominant ionisation stages in H i regions conveniently located at $`\lambda \lambda 2026,2062`$ Å (Zn ii) and $`\lambda \lambda 2056,2062,2066`$ Å (Cr ii). All of these factors account for the fact that Zn ii and Cr ii absorption lines from DLAs have been the target of many studies since the 1990s, even though echelle spectrographs on 8-10 m class telescopes now afford a more comprehensive assessment of the overall chemical composition of QSO absorbers (e.g. Prochaska et al. prochaska01 (2001)). Accordingly, in this paper we focus on the abundances of Zn and Cr in CORALS DLAs and compare them with the large body of such measurements now available for optically selected DLAs.
The paper is organised as follows. In §2 we describe the observations, data reduction process and column density measurements, while in §3 we list the abundances measured in the DLAs in our sample. We compare the CORALS Zn and Cr abundances to previous surveys in §4 and discuss our findings, together with our conclusions in §5.
## 2 Observations and Data Reduction
The CORALS sample of QSOs found to have DLAs with absorption redshifts $`1.8z_{\mathrm{abs}}z_{\mathrm{em}}`$ consists of a total of 18 QSOs (22 DLAs; Ellison et al. 2001b ). Prior to this study, spectra of sufficiently high resolution and signal-to-noise ratio (S/N) for abundance determinations had already been obtained (either by us or by others) for five of these DLAs and have appeared in the literature. We therefore describe here only observations made on the remainder of the CORALS sample, details of which are collected in Table 1.
$`B`$-band magnitudes in column (2) have been taken from Table 2 of the recent study of Ellison et al. (2005a ) except where indicated. Emission redshifts $`z_{\mathrm{em}}`$ in column (3) have been reproduced directly from Table 3 of Ellison et al. (2001b ), while the values of DLA absorption redshift, $`z_{\mathrm{abs}}`$ listed in column (4) were determined from the observations presented here and are quoted to the precision of our measurements. One QSO from the CORALS sample, B1251$``$407, is too faint ($`B=23.7`$) for high resolution spectroscopy and no abundance measurements are therefore available for the two DLAs identified by Ellison et al. in its spectrum. However, we do not expect this small gap in our survey to affect the conclusions of the present study for two reasons. First, the two DLAs in B1251$``$407 are at higher redshifts than the rest of the CORALS sample ($`z_{\mathrm{abs}}=3.533`$ and 3.752). Second, their neutral hydrogen column densities, $`N`$(H i) $`=4\times 10^{20}`$ and $`2\times 10^{20}`$, are at the lower end of the DLA column density distribution. Thus, the metal and dust content of these two DLAs, even if they turned out to be significantly different from those of the rest of the sample, would have only a minor impact on integrated quantities such as the column density weighted metallicity considered in §4. Specifically, if we assume that, unlike most DLAs, these two systems have a solar abundance of Zn, the column density weighted mean metallicity of the CORALS sample would be increased by only 0.1 dex.
### 2.1 Data Acquisition
The observations were made between April 2002 and March 2004 during several runs on a range of telescopes and instruments. The majority of the data were obtained in service mode using the Ultraviolet and Visual Echelle Spectrograph (UVES; Dekker et al. dekker00 (2000)) on the European Southern Observatory Very Large Telescope (VLT). The spectra of B0405$``$331 and B0913$`+`$003 were obtained with the Echellette Spectrograph and Imager (ESI; Sheinis et al. sheinis02 (2002)) at the cassegrain focus of the Keck II telescope; for B0537$``$286 we used the Low-Resolution Imaging Spectrometer (LRIS; Oke et al. oke95 (1995)) at the cassegrain focus of Keck I. B1230$``$101 was observed with the newly-commissioned Magellan Inamori Kyocera Echelle spectrograph (MIKE; Bernstein et al. bernstein03 (2003)) on the 6.5m Magellan Baade telescope.
Instrument settings were chosen for each QSO observation so as to provide coverage at the wavelengths of the redshifted Zn ii and Cr ii absorption lines. We achieved spectral resolutions varying between 0.14 and 0.25 Å FWHM for the UVES and MIKE data, and of $``$1.16 Å and 3.87 Å FWHM for the ESI and LRIS data respectively. These correspond to velocity resolutions of $``$7 km s<sup>-1</sup>, $``$13 km s<sup>-1</sup>, $``$48 km s<sup>-1</sup> and $``$140 km s<sup>-1</sup> for the UVES, MIKE, ESI and LRIS spectra respectively.
### 2.2 Data Reduction
Due to the differing nature of the instruments used, the spectra were not identically reduced, although the same standard steps were incorporated in each case. The reduction of the UVES data was performed with the ESO UVES data reduction pipeline, while the other spectra were reduced using standard IRAF<sup>2</sup><sup>2</sup>2IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation routines. The two-dimensional images were trimmed and the bias level subtracted using the over-scan regions. Pixel-to-pixel variations were corrected for by dividing through by a normalized flat-field of high S/N that had been produced by co-adding several flat-field exposures. Pixels which had been affected by cosmic ray hits were corrected for, where possible, and the echelle orders or single long-slit spectrum were traced. The one-dimensional spectrum was then extracted and the sky spectrum subtracted (in the case of MIKE taking into account the slit tilt, which is not perpendicular to each order and varies over the CCD).
The individual sky-subtracted spectra were then wavelength calibrated (using comparison lamp spectra typically taken immediately following and/or prior to each science frame and extracted in the same way as above, except for sky subtraction) and then corrected to a vacuum heliocentric wavelength scale. In the next step, the spectra of each object were rebinned to a common linear wavelength scale, with a bin size close to their original size and co-added, rejecting any remaining cosmic rays or bad pixels. The co-addition was performed without any weighting and the corresponding error spectra were summed in quadrature. Gaussian fits to several emission lines from the reference lamp spectra near to the redshifted Zn ii and Cr ii lines were made to estimate the values of the spectral resolution (FWHM) listed in Table 1.
### 2.3 Spectral Line Fitting
Metal absorption lines associated with the DLA systems in each QSO spectrum were initially identified based on the redshifts reported by Ellison et al. (2001b ). The vacuum-heliocentric wavelengths measured from our spectra for the strongest component in each absorption system were then used to determine a more accurate redshift for the DLA system, based on the average of the redshifts calculated from each (non-saturated) line. The spectra were then reduced to a rest-frame wavelength scale and cut into sections corresponding to $`\pm 2000\text{km\hspace{0.17em}s}\text{-1}`$ around each absorption line. These sections were then normalised to the local continuum by division by a spline fit to portions of the spectrum judged to be free of absorption.
Ion column densities were deduced from the observed absorption lines using the line fitting software package VPFIT<sup>3</sup><sup>3</sup>3VPFIT is available from http://www.ast.cam.ac.uk/$``$rfc/ vpfit.html. VPFIT fits multiple Voigt profiles (convolved with the instrument profile) to absorption line components. Initial guesses for the column density, redshift and Doppler $`b`$ parameter for each absorption component are used as inputs into the program. VPFIT then computes the $`\chi ^2`$ goodness-of-fit (taking into account the error on each pixel) and automatically varies the parameters, iterating until $`\chi ^2`$ has been minimised and the best fit has been found. Components were rejected from the fit if their Doppler parameters were lower than approximately half the spectral resolution. For the rest-frame wavelengths and $`f`$-values of the transitions, which are inputs to VPFIT, we consulted the recent compilation by Morton (morton03 (2003)); for reference, values for the Zn ii and Cr ii multiplets are reproduced in Table 2.
The first step in the line fitting procedure was to determine the multi-component velocity structure of gas in the damped system by fitting lines spanning a wide range of $`f`$-values, preferably multiplets of Fe ii and Si ii. The parameters of the fit—that is redshift and Doppler parameter—were then kept fixed and applied to the generally weaker Zn ii and Cr ii lines. In cases where only some of the Fe ii and Si ii absorption components are detected in Zn ii and Cr ii, the total column densities of Zn<sup>+</sup> and Cr<sup>+</sup> were adjusted upwards to allow for the unseen components in the same proportion as determined for Fe ii and Si ii. The implicit assumption is that the relative proportions of absorbers in each component are the same for all the ions considered. While in principle this may not be true, because of possible differences in the degree of dust depletion of different elements, in practice element ratios are found to be remarkably uniform between the multiple velocity components within a DLA (e.g. Prochaska 2003). The correction for unseen components amounted to less than $`10`$% of the total column density in all but two cases, B0432$``$440 and B0933$``$333, where the correction for unseen components is as much as $`50`$%. Indeed, this correction for undetected components was the main reason why we used VPFIT to analyse the Zn ii and Cr ii absorption. In all cases considered here, the Zn ii and Cr ii lines are unsaturated and we would have obtained the same values of $`N`$(Zn<sup>+</sup>) and $`N`$(Cr<sup>+</sup>)—apart from the ‘incompleteness’ correction—had we analysed them with the optical depth method (Hobbs hobbs74 (1974)) which some prefer to profile fitting.
Out of the 15 DLAs listed in Table 1, we report positive detections of both the Zn ii and Cr ii multiplets in four cases (see Table 3). In two additional cases we detect Cr ii but not Zn ii, while in the remaining nine DLAs we place upper limits to both $`N`$(Zn<sup>+</sup>) and $`N`$(Cr<sup>+</sup>). The upper limits are $`3\sigma `$, deduced from the measured S/N ratios (listed in column (8) of Table 1) and the velocity spread of the absorption implied by stronger lines of Fe ii and Si ii, as explained above.
Figure 1 shows examples of the spectra in the regions of the Zn ii and Cr ii multiplets, including three of the Zn ii and Cr ii detections and one case when only Cr ii is detected.<sup>4</sup><sup>4</sup>4Spectra of all the CORALS QSOs covered in the present study are available from ftp://ftp.ast.cam.ac.uk/pub/papers/CORALS/ Also shown in the figure are examples of the stronger lines used to determine the velocity structure of each DLA, as well as the fits generated by VPFIT. Table 4 lists values of the Doppler $`b`$ parameter for the detected Zn II and Cr II lines; full details of the velocity structure of other absorption lines will be presented in a future publication (Akerman 2005).
## 3 Abundance determinations
Table 3 lists the column densities of Zn ii and Cr ii, or upper limits, measured in this study together with the 1 $`\sigma `$ uncertainties in their values. Since the first ions are the dominant ionization stages of Zn and Cr in H i regions, we can readily derive the abundances of these two elements by dividing the values of $`N`$(Zn<sup>+</sup>) and $`N`$(Cr<sup>+</sup>) by the neutral hydrogen column densities $`N`$(H<sup>0</sup>) listed in column (3) of Table 3—the results are given in columns (5) and (7) respectively. The values of $`N`$(H<sup>0</sup>) are reproduced directly from the compilation by Ellison et al. (2001b ) because (a) few of our spectra include the Ly$`\alpha `$ absorption line and (b) even in cases where we cover the transition, the intermediate resolution spectra of Ellison et al. are better suited than ours to the determination of $`N`$(H<sup>0</sup>) by fitting the damping wings of the Ly$`\alpha `$ absorption line. We have adopted throughout a conservative error of $`\pm 25`$% to $`N`$(H<sup>0</sup>), based on the typical accuracy with which such measurements have been reported in the literature.
In Table 5 we list the abundances of Zn, Cr, Fe, and Si, expressed in logarithmic units relative to solar, for each of the DLAs in the complete CORALS sample; the list includes all the new measurements reported here plus the five already available from the literature—references are given in the last column of Table 5. The only DLAs in the compilation by Ellison et al. (2001b ) which are missing from our sample are the two systems in front of the very faint QSO B1251$``$407, as explained in §2. We have adopted throughout the solar abundances proposed by Lodders (lodders03 (2003)) in her comprehensive reassessment of meteoritic and photospheric abundances; those from the more recent compilation by Asplund, Grevesse, & Sauval (asplund04 (2004)) differ by no more than 0.03 dex for the elements of interest here. All the Zn and Cr measurements collected in Table 5 (and indeed those of the larger comparison sample discussed in §4) have been reduced to the same set of $`f`$-values, as listed in Table 2, except for cases where the differences in column density would have amounted to less than a few percent.<sup>5</sup><sup>5</sup>5Most recent work has used the same set of $`f`$-values for the Zn ii and Cr ii multiplets, from the laboratory measurements by Bergeson & Lawler (bergeson93 (1993)). Small differences between the values quoted by different observers result from rounding errors. It was not possible, however, to ensure the same degree of homogeneity for the Fe ii and Si ii measurements.
## 4 Metallicity and dust in the CORALS survey
The purpose of this study is to test whether element abundances in CORALS DLAs are significantly different from those of DLAs drawn from optically selected QSO samples, and therefore to determine the extent, if any, to which previous surveys have been biased against metal-rich, high column density absorbers. We consider this question from two points of view, by examining first the metallicity distributions indicated by the \[Zn/H\] determinations, and then the degree of depletion of refractory elements implied by the \[Cr/Zn\] measures.
### 4.1 Comparison of Zn Abundances
The most extensive, homogeneous, compilation of \[Zn/H\] measurements in DLAs has recently been assembled by Kulkarni et al. (2005); it includes 51 detections and 36 upper limits over the redshift interval $`z_{\mathrm{abs}}=0.09`$ to 3.90. (Three of the Zn ii detections in the Kulkarni et al. data set are in common with the CORALS survey and have thus not been included in the control sample in the following analysis.) Although larger compilations of DLA abundance measurements have been published (e.g. Prochaska et al. 2003), they bring together data for other elements, mostly Si and Fe, in addition to Zn. The improved statistics afforded by such compilations are offset, in our view, by systematic uncertainties due to differing degrees of dust depletion and possible nucleosynthetic departures from solar relative abundances of different elements. Thus, we consider it more appropriate to restrict the present analysis to the comparison of Zn abundances alone.
The CORALS \[Zn/H\] abundances are compared with those from the compilation by Kulkarni et al. (kulkarni05 (2005)) in Figure 2 (after rescaling the latter to the solar value from Lodders lodders03 (2003)). From a visual inspection of the plot we conclude that: (a) both the CORALS and the comparison sample of optically selected DLAs are generally metal-poor, with typical values of \[Zn/H\] well below solar; and (b) there is a hint that, overall, the CORALS DLAs may have marginally higher metallicities, although the considerable number of upper limits complicates the comparison. We now address these points quantitatively.
The quantity which is of interest for ‘cosmic’ chemical evolution models (Pei & Fall pei95 (1995)) is the column density-weighted metallicity
$$[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]=\mathrm{log}(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}\mathrm{log}(\mathrm{Zn}/\mathrm{H})_{}$$
(1)
where
$$(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}=\frac{\underset{i=1}{\overset{n}{}}N(\mathrm{Zn}^+)_i}{\underset{i=1}{\overset{n}{}}N(\mathrm{H}^0)_i}$$
(2)
which is a measure of the degree of metal enrichment of the DLA population as a whole. The summation in eq. (2) is over the $`n`$ DLAs considered in a given sample. The present sample is too small to consider subsets of the data within different redshift intervals, as done by Pettini et al. (pettini99 (1999)), Prochaska et al. (prochaska03 (2003)), and Kulkarni et al. (kulkarni05 (2005)). Instead we compute eq. (2) for the CORALS sample as a whole, and compare the result with the analogous quantity for the Kulkarni et al. (kulkarni05 (2005)) sample, computed over the same redshift interval spanned by the CORALS data, $`z_{\mathrm{abs}}=1.863.45`$ for the full set (57 DLAs), and $`z_{\mathrm{abs}}=1.862.81`$ for the Zn ii detections only (27 DLAs). We find $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]=0.88\pm 0.21`$ for CORALS DLAs, and $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]=1.09\pm 0.10`$ for the control sample of Kulkarni et al. (kulkarni05 (2005)). These values, which differ at only the $`1\sigma `$ level, were calculated considering only the Zn ii detections. Repeating the calculations, but now including the upper limits as if they were detections, we obtain $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]=0.87\pm 0.13`$ and $`1.17\pm 0.07`$ respectively. The errors were estimated using bootstrap techniques (Efron & Tibshirani efron93 (1993)). For each sample (CORALS and the control sample), we constructed random datasets by drawing $`n`$ times from the real dataset with replacement, where $`n`$ is the number of \[Zn/H\] measurements in each sample. The procedure was repeated a million times to build up a distribution of values of $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]`$; we quote the standard deviation of this distribution as our estimate of the error in the quantity $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]`$.
These results are shown graphically in Figure 3. Again we see that the two samples are very similar. The marginally higher metallicity of the CORALS sample may be real. Alternatively it may be an artifact of small number statistics or, in the case where we include upper limits, it may be due to the higher proportion of upper limits skewing the results—over one half of the CORALS measures of \[Zn/H\] are upper limits, compared with one third for the control sample of Kulkarni et al. (kulkarni05 (2005)).
In order to clarify this point, we conducted a statistical test based on “survival statistics” (which takes account of upper limits) with the program ASURV (LaValley, Isobe & Feigelson lavalley92 (1992)), which implements the statistical methods of Feigelson & Nelson (feigelson85 (1985)). Using a Peto-Prentice test (Latta latta81 (1981)), we tested the null hypothesis that the two samples are drawn from the same parent population and found the hypothesis to be true at the 90% confidence level. We conclude that CORALS DLAs do not exhibit significantly different metallicities from those of existing, larger, samples of DLAs assembled from optically selected QSO surveys.
### 4.2 Comparison of Cr/Zn ratios
An indication of the degree of depletion of refractory elements onto dust grains may be obtained from the ratio of the abundances of chromium to zinc, as explained in §1. In Figure 4 we plot this ratio against the metallicity \[Zn/H\] for each of the CORALS DLAs (coloured) from Table 5 together with analogous measurements (black) from the compilations by Khare et al. (khare04 (2004)) and Kulkarni et al. (kulkarni05 (2005)), after rescaling their values to the same solar abundances used here.
Figure 4 shows the trend of increasing Cr depletion with increasing metallicity which was previously noted by Pettini et al. (1997a ) and shown by Prochaska & Wolfe (2002) to be a general feature of refractory elements in DLAs. In systems with \[Zn/H\] $`\stackrel{<}{}1.5`$, \[Cr/Zn\] is approximately solar—indicating that there is little dust depletion at such low metallicities—while when \[Zn/H\] $`>1`$, up to $`90`$% of the Cr can be ‘missing’ from the gas phase and presumably be in solid form. Even so, in none of the DLAs do we see the extreme depletions of Cr, by two orders of magnitude, commonly measured in cold clouds of the Milky Way disk (Savage & Sembach 1996). At the typical DLA metallicity, $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]1`$, approximately $`1/2`$ to $`2/3`$ of the Cr is in the dust (\[Cr/Zn\] $`0.3`$ to $`0.5`$) although there is considerable dispersion in the depleted fraction $`f_{\mathrm{Cr}}`$.
All of these facets of the depletion of refractory elements in DLAs have recently been discussed by Vladilo (2004) who linked them to the (generally early) chemical evolution of the galaxies where the absorption systems originate. The dependence of $`f_{\mathrm{Cr}}`$ on \[Zn/H\] may reflect a metallicity dependence of the efficiency of dust formation in the ejecta of core-collapse supernovae and in the winds of late-type giants. On the other hand, the large scatter which accompanies the overall trend is presumably an indication of how the detailed balance between the processes of dust formation, accretion, and destruction is affected by the local physical conditions in the ISM.
As far as the present work is concerned, it is evident from Figure 4 that the depletions of Cr in the CORALS sample conform to the overall pattern described above. The CORALS DLAs are not extreme in their values of $`f_{\mathrm{Cr}}`$, nor do they exhibit lower values of \[Cr/Zn\] at a given \[Zn/H\] than DLAs drawn from optically selected QSO samples.<sup>6</sup><sup>6</sup>6Although three out of the four CORALS DLAs with $`[\mathrm{Zn}/\mathrm{H}]<1.0`$ have $`[\mathrm{Cr}/\mathrm{Zn}]<0.5`$ compared to only five out of 16 from the control sample, these two fractions are different at less than the $`1\sigma `$ significance level even when using Poisson statistics (which overestimate the significance of any differences between samples when the number of measurements is so small). If there are any DLAs where the depletions of refractory elements approach the high values typical of cold clouds in the Milky Way disk, we have not found them yet. Based on the data in Figure 4, there is no evidence that CORALS DLAs should redden the spectra of background QSOs any more than a typical DLA from an optically selected sample.
## 5 Discussion
The work described in this paper concludes a project begun six years ago to test the extent to which existing DLA samples are biased by dust reddening against gas-rich galaxies of high metallicity. The first results, reported by Ellison et al. (2001), showed that $`\mathrm{\Omega }_{\mathrm{DLA}}`$ had not been significantly underestimated. To that conclusion we now add the findings that:
(1) At redshifts $`1.9<z<3.5`$, the metallicity of CORALS DLAs, as measured by the \[Zn/H\] ratio, is only marginally higher (at a statistically insignificant level of only $`1\sigma `$) than that of DLAs drawn from optically selected QSOs—we determine $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]=0.88\pm 0.21`$ for CORALS DLAs, and $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]=1.09\pm 0.10`$ for the control sample of Kulkarni et al. (kulkarni05 (2005)) over the same redshift interval. We note that none of the CORALS DLAs lie in the dust-‘forbidden’ zone of Prantzos & Boissier (prantzos00 (2000)) where $`[\mathrm{Zn}/\mathrm{H}]+\mathrm{log}N(\mathrm{H}\text{i})>21`$ (see Figure 5).
(2) The dust-to-metals ratio, as measured by the quantity \[Cr/Zn\], exhibits no systematic difference between the two samples—we unearthed no evidence to show that radio-selected QSOs should be more reddened by intervening systems than optically selected QSOs.
These results, together with recent reports of low dust extinction in large DLA samples drawn from the Sloan Digital Sky Survey (e.g. Murphy & Liske 2004), make it increasingly difficult to appeal to dust-induced selection effects to explain the observed properties of DLAs. It seems unlikely, for example, that the true metallicity of DLAs may have been underestimated by as much as a factor of five, as recently claimed by Vladilo & Péroux (2005). Below the redshift limit of our sample, $`z<1.8`$, the situation is less clear. The dust fraction may be higher (e.g. Vladilo 2004), but the number density of Mg ii absorbers from an extension of the CORALS radio-selected sample (in the range $`0.6<z<1.7`$) is in excellent agreement with that from optically selected surveys (Ellison 2005), suggesting that dust obscuration is not a problem at these lower redshifts either.
One caveat (often mentioned when considering the results of the CORALS survey) is the limited size of the CORALS DLA sample: are the 20 DLAs in Table 5 representative of the population as a whole, or are we being thrown off course by a statistical fluctuation? Clearly, only future observations of a larger sample of radio selected (or possibly X-ray selected) QSOs will settle this issue and test, for instance, whether the higher metallicity—by 0.21 dex—reported here is a real difference, or possibly even an underestimate of the true value of $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]`$ in an unbiased sample of DLAs. For the moment, in the absence of better statistics, we can use a Monte-Carlo approach to address such concerns, particularly if we wish to test for differences at a level as high as the factor of five claimed by Vladilo & Péroux (vladilo05 (2005)).
Specifically, we have investigated the likelihood of measuring a column density weighted metallicity $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]0.88`$, as found here, by drawing 20 DLAs at random from a much larger parent sample of DLAs with the column density and metallicity distributions proposed by Vladilo & Péroux (vladilo05 (2005)). These authors approximated the true H i column density distribution with a power law of the form $`f_{N(\mathrm{H}\mathrm{I})}N`$(H i), and the true metallicity distribution with a Schechter function $`f_Z(Z/Z_{})^\alpha e^{Z/Z_{}}`$ (motivated by luminosity-metallicity relationship of galaxies in the local universe). The least extreme set of parameters among those considered by Vladilo & Péroux has $`\beta =1.6`$, $`\alpha =0.46`$ and $`\mathrm{log}(Z_{}/Z_{})=0.19`$; with these values, their column density weighted metallicity is $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]=0.44`$. Our simulations showed that, drawing 20 DLAs at random from this parent population, one would find by chance values of $`[(\mathrm{Zn}/\mathrm{H})_{\mathrm{DLA}}]`$ as low as $`0.88`$, or lower, only in five cases out of a hundred.
The absence of a detectable dust-related bias in current magnitude-limited samples of optically selected QSOs may appear surprising to some. After all, such a bias has long been advocated by theorists to improve the match of their models to the properties of DLAs (e.g. Cen et al. cen03 (2003); Churches et al. churches04 (2004); Nagamine et al. nagamine04 (2004)). It also seemed the natural explanation for the empirical lack of DLAs with high column densities of metals highlighted by Boissé et al. (boisse98 (1998)) and still present in the larger sample assembled by Kulkarni et al. (2005) and indeed in the CORALS sample considered here (see Figure 5). The reason why most optical QSO samples do not underestimate significantly quantities such as $`\mathrm{\Omega }_{\mathrm{DLA}}`$ and $`Z_{\mathrm{DLA}}`$ compared with CORALS was clarified by Ellison et al. (ellison04 (2004)) and is related to the shape of the QSO luminosity function. Optical QSO surveys will not yield significantly skewed DLA statistics provided they reach below the break in the QSO luminosity function at $`B19`$. Brighter QSO samples, on the other hand, may show a bias from either dust extinction or lensing (the two effects would of course operate in different directions).
The main conclusion of the CORALS project so far—that there are only minor differences, if any, between DLA samples drawn from QSOs surveys at radio and at optical wavelengths—can only be regarded as ‘good news’. Its corollary is that the large data samples being made available by major projects such as the Two-degree Field and the Sloan Digital Sky Survey afford us an unfettered view of the absorber population, although surveys which concentrate only on bright QSOs may not. The challenge is now to understand, perhaps with more focused theoretical efforts, the rightful place of damped Ly$`\alpha `$ systems within the diverse population of galaxies known to inhabit the high redshift universe.
###### Acknowledgements.
We are very grateful to Varsha Kulkarni for providing us with her compilation of Zn and Cr measurements, to Kurt Adelberger and Naveen Reddy for their assistance with the Keck observations, and to the anonymous referee whose suggestions improved the paper. We also acknowledge useful discussions with Giovanni Vladilo. We wish to recognize the significant cultural role and reverence that the summit of Mauna Kea has always had within the indigenous Hawaiian community. We are most fortunate to have the opportunity to conduct observations from this mountain.
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# Uniform approximation of barrier penetration in phase space
## I Introduction
Semiclassical approaches to multidimensional tunnelling lead to very interesting problems in complexified classical dynamics, often with incompletely understood solutions. For example, recent work in references Shudo ; aY00 ; Takahashi ; TYK ; Onishi ; Julia has shown that nontrivial geometrical structure such as complex homoclinic intersections have an important role to play in multidimensional barrier penetration and that even complex chaos can be relevant. Given the difficulty inherent in a systematic treatment of multidimensional tunnelling as a result of such issues, it is perhaps surprising that a relatively simple description can be given of barrier penetration at a critical energy where classically allowed transmission mechanisms turn on and where primitive semiclassical approximations must be replaced by somewhat more complicated uniform ones.
An approach which achieves this has been proposed in references sC04 and sc05 and in this paper we apply the method explicitly to a model waveguide problem. The model is chosen to be rather simple so that fully quantum calculations are easy to perform accurately for purposes of comparison. We emphasise, however, that semiclassical aspects of the calculation are as easily applied to other problems, provided the topology is similar, and provide a description, for example, of collinear atom-diatom reactions. For that reason we use the terminology of chemical reactions in this paper and equate the probability of transmission with a probability of reaction. In fact the approach we describe here provides a natural means of visualising the quantum scattering problem in phase space and as such shows an interesting connection with classical transition state theories of chemical reaction. These classical theories have recently been of interest because the periodic-orbit dividing surface (PODS) construction Pechukas ; PODS2 has been generalised to arbitrary dimensions using the construction of normally hyperbolic invariant manifolds (NHIMS) Jaffe ; Uzer ; WW ; HCN ; LW . The classical constructions emerge naturally in our semiclassical approximation and we note that even though the illustrations offered here are in two degrees of freedom there are straightforward generalisations to higher dimensional problems where the full generality of the NHIM construction comes into play.
The approximation we use can be stated very simply as an abstract operator equation but for explicit illustration we present results in phase space, using the Wigner-Weyl calculus. In particular, we define a Weyl symbol of a transmission matrix which represents, in an averaged sense, a reaction probability as a function of phase space. We find that above threshold the support of this Weyl symbol closely mimics the shape of the classically reacting region but the Weyl symbol itself also incorporates tunnelling and other quantum effects. Notation and details for this construction are set out in section II.
We implement two versions of the theory. First, a harmonic approximation derived in sC04 is applied in section III whose classical input consists simply of an instanton orbit, along with its action and monodromy matrix. This approximation works when the classically reacting region is a small neighbourhood of the initial condition for the instanton orbit and is harmonic in the sense that it uses linearised dynamics generated by an elliptic quadratic Hamiltonian on the Poincaré section. This harmonic version covers the threshold case where classical reaction switches on as a function of energy and is relatively easy to apply. It fails however when the energy is too far above threshold and the classically reacting region is too large to be adequately described by linearised dynamics.
A semiclassical approximation that is more accurate and has greater range has been derived in sc05 and this is applied in section IV. This version uses fully nonlinear dynamics to extend further from the instanton orbit. Like the harmonic version it can be stated quite simply as an abstract operator equation but its practical implementation is more difficult. Difficulty arises primarily because we must invert an operator constructed semiclassically as an evolution operator and this inversion cannot at present be achieved in closed form. In this paper we achieve that inversion using numerical methods and while this aspect of the approach needs further work to provide an appealing semiclassical method, we can verify unambiguously that the nonlinear version of the theory is capable of describing the quantum transmission problem very accurately (see Figure 8).
We conclude this section by outlining how these results relate to existing work. The basic formalism here of relating scattering to complex barrier-crossing orbits goes back to the work of Miller and coworkers MG ; MillerSmatrix on the classical $`S`$-matrix. Our intent is to describe a simple uniform extension of this approach which applies at the boundary of classical reaction where the fate of classical orbits changes discontinuously. What allows us to make progress is that we do not directly describe the $`S`$-matrix but instead consider a transmission matrix derived from it which gives probabilities rather than amplitides. The advantage of this problem is that contributing orbits at the boundary of classical reaction depend smoothly on initial conditions, despite the singular nature of real orbits there sC04 , and give simple semiclassical expressions. This uniformisation is similar to established results relating one-dimensional transmission probabilities FW ; BM ; Smil or the cumulative reaction probability Millermicro to sums over multiple barrier crossings but includes information about how the probabilities depend on the incoming state. It is different however from uniform approximations of the scattering operator such as described in ElranKay which describe explicit matrix elements. These are uniform with respect to variation of quantum numbers whereas our approach treats the scattering operator abstractly and is uniform with respect to energy and in phase space.
Direct approximation of the scattering problem by complex trajectories has recently been examined in aY00 ; Takahashi ; TYK in the context of nonintegrable systems. It has been found there that, while intuitively one might expect tunnelling processes to be dominated by short complex orbits which cross the barrier directly, the dominant complex orbits can have have a surprisingly nontrivial topology in the deep tunnelling regime. It has even been found that the dominant complex orbits may be chaotic in related treatments of quantum propagation Shudo ; Onishi ; Julia . We also find evidence of the “fringed tunnelling” characteristic of such mechanisms in our fully quantum solutions but the theory we outline is intended to cover only the immediate vicinity of the reacting region where direct tunnelling mechanisms are dominant. In the deep tunnelling regime our uniform results revert to standard primitive approximations and we should in principle be able to marry our approach with that of aY00 ; Takahashi ; TYK . It is not obvious, however, that a fully uniform calculation could easily be applied when the contributing complex orbits are more numerous and more complicated and we do not consider that problem explicitly here.
## II Representing the scattering matrix in phase space
In the following sections we will develop semiclassical approximations for representations of the scattering matrix in phase space. In this section we illustrate these representations numerically using a two-dimensional waveguide, which serves as a simplified model of a collinear atom-diatom collision. This model consists of a particle of unit mass moving in the potential
$$V(x,y)=sech^2x+\frac{1}{2}\omega ^2(x)(ya(x))^2,$$
(1)
where
$`\omega ^2(x)`$ $`=`$ $`\mathrm{\Omega }^2+\lambda sech^2x`$
$`a(x)`$ $`=`$ $`\mu sech^2x.`$
A very similar potential has been used in aY00 ; TYK . The simplicity of this model will enable us easily to obtain accurate numerical solutions, which will be useful for later comparison with semiclassical approximations where exponentially small tunnelling effects are of interest. We emphasise, however, that none of the theory that follows is dependent on this simplicity and semiclassical aspects of the discussion can just as easily be applied to any other system, as long as the Hamiltonian is an analytic function of its arguments. The essential structural features we assume are that the waveguide should have a single bottleneck separating asymptotically decoupled channels, which we label as reactant and product channels respectively, and that the energy should be sufficiently close to threshold that recrossings of the transition state do not occur.
For future reference, it will be useful to denote the asymptotically decoupled potential by the symbol
$$V(x,y)V_{\mathrm{}}(y)=\frac{1}{2}\mathrm{\Omega }^2y^2.$$
We can therefore write the asymptotic scattering states for this problem analytically, as solutions of a harmonic oscillator.
We write the scattering matrix in block form
$$S(E)=\left(\begin{array}{cc}r_{RR}& t_{RP}\\ t_{PR}& r_{PP}\end{array}\right),$$
where the subscripts $`R`$ and $`P`$ refer to reactant and product channels respectively. For example, the transmission matrix $`t_{PR}`$ maps asymptotically incoming states on the reactant side to asymptotically outgoing states on the product side. The theory we describe is for the matrix
$$\widehat{}(E)=t_{PR}^{}t_{PR}$$
rather than for the scattering matrix itself. The matrix $`\widehat{}(E)`$ has an obvious physical role determining state-specific reaction rates. The transmission probability for an incoming state labelled $`|\psi _n`$ on the reactant side can be written as a matrix element
$$p_n=\psi _n|\widehat{}(E)|\psi _n$$
(2)
of $`\widehat{}(E)`$. In going from the scattering matrix to $`\widehat{}(E)`$ we lose information about phase and about the distribution of product states but, as described in sC04 , uniform semiclassical approximations for $`\widehat{}(E)`$ are expected to be considerably simpler than those for $`S(E)`$.
Formally, we can think of $`\widehat{}(E)`$ as an operator acting on the Hilbert space $`_R^{\mathrm{in}}(E)`$ of asymptotically propagating states in the incoming reactant channel — in this context we refer to it as the reaction operator in the following. We consider systems for which there are a finite number $`M`$ of such states so $`_R^{\mathrm{in}}(E)`$ is finite-dimensional and $`\widehat{}(E)`$ can be represented by an $`M\times M`$ matrix. The Wigner-Weyl correspondence offers an alternative representation as a function on phase space and the theory we outline in the coming sections is stated in those terms. In the remainder of this section we describe how the connection is made formally between the matrix representation and the representation in Wigner-Weyl correspondence, or the Weyl symbol of $`\widehat{}(E)`$. In preparation for that discussion, let us first describe the phase space on which these representations are defined.
As described in sc05 (see also Bog ), the natural classical analog of the Hilbert space $`_R^{\mathrm{in}}(E)`$ is a Poincaré section $`\mathrm{\Sigma }_R^{\mathrm{in}}(E)`$ defined by fixing the energy $`E`$ and an asymptotically large value of the reaction coordinate $`x`$. We use $`(y,p_y)`$ as canonical coordinates on $`\mathrm{\Sigma }_R^{\mathrm{in}}(E)`$ and we may alternatively denote points in $`\mathrm{\Sigma }_R^{\mathrm{in}}(E)`$ using the vector notation
$$𝜻=\left(\begin{array}{c}y\\ p_y\end{array}\right).$$
We define the allowed region of the $`y`$-$`p_y`$ plane by the condition $`p_y^2/2+V_{\mathrm{}}(y)<E`$ and note that, as usual in this sort of correspondence, the dimension $`M`$ of $`_R^{\mathrm{in}}(E)`$ is approximated by the Liouville area of this region divided by $`2\pi \mathrm{}`$.
To calculate the Weyl symbol $`𝒲_\widehat{}(𝜻,E)`$ of $`\widehat{}(E)`$, we denote its individual matrix elements by $`R_{nm}(E)=\psi _n|\widehat{}(E)|\psi _m`$ and write
$$𝒲_\widehat{}(𝜻,E)=2\pi \mathrm{}\underset{nm}{}R_{nm}(E)𝒲_{nm}(𝜻)$$
(3)
where $`𝒲_{nm}(𝜻)`$ are the Weyl symbols of the projectors $`|\psi _n\psi _m|/2\pi \mathrm{}`$ (the factors of $`2\pi \mathrm{}`$ are to keep the notation consistent with standard practice for Wigner functions in the case $`n=m`$). For the asymptotically harmonic potential in (1) the functions $`𝒲_{nm}(𝜻)`$ are known analytically (see Ripamonti for example).
The Weyl symbol $`𝒲_\widehat{}(𝜻,E)`$ of $`\widehat{}(E)`$ provides a remarkably transparent means of visualising the quantum transmission problem and of relating the scattering matrix to the geometry of classical phase space. To illustrate, we show examples of $`𝒲_\widehat{}(𝜻,E)`$ in Figure 2 for the model potential in (1) with energies at and above threshold. These results have been obtained by first computing the scattering matrix numerically using symplectic integration combined with the log derivative method as described in dM86 ; dM95 and then using (3). In each case we see that $`𝒲_\widehat{}(𝜻,E)`$ is effectively supported in a region of $`\mathrm{\Sigma }_R^{\mathrm{in}}(E)`$, which we can identify as the quantum-mechanically reacting region. Indeed, by rewriting the reaction probability in (2) (using standard properties of the Wigner-Weyl correspondence) in the form
$$p_n=𝒲_\widehat{}(𝜻,E)𝒲_{nn}(𝜻)d𝜻,$$
(4)
where $`\mathrm{d}𝜻=\mathrm{d}y\mathrm{d}p_y`$, and interpreting the Wigner function $`𝒲_{nn}(𝜻)`$ as a phase space pseudodensity, it is natural to identify $`𝒲_\widehat{}(𝜻,E)`$ as a probabability of reaction as a function of phase space, albeit in an averaged sense. Although the uncertainty principle prevents us from defining a point-wise transmission probability in phase space, we can construct linear combinations of incoming states with a fixed total energy (as in (5) below) whose Wigner function is supported within an area of $`O(\mathrm{})`$ in $`\mathrm{\Sigma }_R^{\mathrm{in}}(E)`$ and the appropriate modification of (4) then gives the reaction probability as an average of $`𝒲_\widehat{}(𝜻,E)`$ over that support.
For energies at or below threshold, transmission is controlled by tunnelling and $`𝒲_\widehat{}(𝜻,E)`$ is supported in a phase space region of area $`O(\mathrm{})`$, centred around an initial condition that leads to an optimal tunnelling route. An explicit semiclassical expression for $`𝒲_\widehat{}(𝜻,E)`$ in this case will be given later and one can see for the threshold case in Figure 2(a) that $`𝒲_\widehat{}(𝜻,E)`$ is indeed peaked around a single point in $`\mathrm{\Sigma }_R^{\mathrm{in}}(E)`$. As the energy increases above threshold, a classically reacting region appears, initially centred on the orbit associated with optimal tunnelling. The boundaries of the classically reacting regions are indicated in Figures 2(b) to (f) by continuous closed curves (as the energy falls to the threshold case $`E=1`$, the classically reacting region shrinks to a point corresponding to the optimal tunnelling route and around which $`𝒲_\widehat{}(𝜻,E)`$ is concentrated). One can see in each case that the reacting region closely matches the support of $`𝒲_\widehat{}(𝜻,E)`$.
Before describing how $`𝒲_\widehat{}(𝜻,E)`$ is approximated semiclassically, we should outline how the classically reacting regions in Figure 2 are defined. The boundary of the classically reacting region in full phase space is the stable manifold on the reacting side of a PODS in two dimensions or more generally a NHIM if higher-dimensional problems are treated. Extended into the incoming reactant channel, this stable manifold defines a tube, the interior of which consists of classically reacting trajectories and whose annular exterior in the classically allowed region consists of trajectories which eventually return along the outgoing reactant channel. A representation of this reacting region in a Poincaré section $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ is obtained simply by taking the a section of the tube of reacting trajectories at fixed energy and a fixed, asymptotically large value of the reaction coordinate $`x`$. Since the tube continues to evolve asymptotically (according to the dynamics of a decoupled potential $`V_{\mathrm{}}(y)`$), the shape of a reacting region defined in this way will depend on the value chosen for the coordinate $`x`$. In models where the dynamics of the reacting region is nonlinear — or potentially even chaotic in problems of higher dimension — the shape of the reacting region will not have a limit and becomes ever more complicated as $`x`$ is brought to infinity. To obtain a fixed asymptotic limit we therefore renormalise the dynamics by using the decoupled evolution of the limiting potential $`V_{\mathrm{}}(y)`$ to map the asymptotic section back to one corresponding to a fixed finite value of $`x`$. In making semiclassical comparisons the value of $`x`$ used to define this final section is dictated by the conventions used for the scattering matrix.
In the present case the asymptotic states in terms of which the scattering matrix is defined are of the form,
$$\mathrm{\Psi }_{n,E}(x,y)\frac{\mathrm{e}^{ik_nx}}{\sqrt{\mathrm{}k_n}}\psi _n(y),$$
(5)
where $`\psi _n(y)`$ are the eigenfunctions of the decoupled potential $`V_{\mathrm{}}(y)`$. The phases of these scattering states are zeroed at $`x=0`$ and asymptotic incoming states can be constructed by starting with the transverse modes at $`x=0`$ and propagating them backwards into the asymptotic region of the incoming channel using dynamics defined by $`V_{\mathrm{}}(y)`$. To make a comparison with the classical picture, the renormalisation of the classical dynamics should therefore take an asymptotic Poincaré section back to one defined by $`x=0`$. This is the convention used in Figure 2 to compare $`𝒲_\widehat{}(𝜻,E)`$ with the classically reacting region. Alternative phase conventions would lead to a reaction operator $`\widehat{}(E)`$ obtained by conjugation of the one we define by a unitary matrix which is diagonal in the basis $`|\psi _n`$. This conjugation makes no difference to the diagonal matrix elements $`\psi _n|\widehat{}|\psi _n`$ but is important for the appearance of the Weyl symbol in $`\mathrm{\Sigma }_R^{\mathrm{in}}`$. Choosing values other than $`x=0`$ for the reference section $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ would, for example, lead to a deformation of the Weyl symbol by the asymptotically decoupled dynamics.
Note that this renormalisation procedure is simply means of interpreting a term in the phase function in the classical $`S`$-matrix MillerSmatrix that fixes its asymptotic value. By accounting for this term using a conjugation of the asymptotic dynamics by the mapping in decoupled dynamics back to $`x=0`$, we can incorporate everything about the classical $`S`$-matrix into a single Poincaré mapping and present results in a more compact form, as described more fully in the coming sections.
## III Harmonic Approximation
In this section we describe a semiclassical approximation for $`\widehat{}(E)`$ based on dynamics linearised around an optimal tunnelling orbit. Although less accurate and valid over a smaller range of energies than the fully nonlinear theory described in the next section, this harmonic approximation captures the essential qualitative behaviour of $`\widehat{}(E)`$ and works well in the especially interesting range of energies around threshold where classical reaction switches on. It is also considerably simpler to apply and can be expressed in closed form using easily obtained classical data. Note that the term “harmonic” here refers simply to the fact that linearised dynamics about a tunnelling orbit are used and has nothing to do with the harmonic asymptotic behaviour of the model we use for numerical illustration.
### III.1 Operator version
A full description of the harmonic approximation to $`\widehat{}(E)`$ has been given in sC04 and we refer there for details and a derivation of the approach. Here we simply summarise the important points. The operator $`\widehat{}(E)`$ is approximated by a formula
$$\widehat{}(E)\frac{\widehat{𝒯(E)}}{1+\widehat{𝒯(E)}}$$
(6)
where $`\widehat{𝒯}(E)`$ is a “tunnelling operator” constructed from classical data. The elements needed to compute $`\widehat{𝒯}(E)`$ are as follows.
* A complex periodic orbit $`\gamma _E(t)`$, the “instanton”, is found which encircles the transition state in imaginary time. This orbit can be found for energies above and below threshold and its dynamical characteristics depend smoothly on energy there.
* At the end of the previous section it was described how the boundary of the reacting region can be first extended arbitrarily far into the asymptotic region and then renormalised by mapping back to a section $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ at $`x=0`$ using decoupled dynamics, so that the shape remains fixed as dynamics are extended into the asymptotic region. An analogous renormalisation, illustrated in Figure 2, is applied to $`\gamma _E(t)`$ so that an asymptotically fixed initial condtion for it is defined in $`\mathrm{\Sigma }_R^{\mathrm{in}}`$. Above threshold this initial condition is near the centre of the reacting region.
* The imaginary action of $`\gamma _E(t)`$ is denoted $`iK_0=_{\gamma _E}𝐩d𝐪`$. We also denote $`\theta =K_0/\mathrm{}`$. Note that for energies above threshold we have $`K_0<0`$ while $`K_0>0`$ below threshold.
* Linearised dynamics around $`\gamma _E(t)`$ are characterised by a complex monodromy matrix $`W`$, which is routinely determined as part of a numerical search for the orbit $`\gamma _E(t)`$. The eigenvalues of $`W`$ come in real reciprocal pairs $`(\mathrm{\Lambda },\mathrm{\Lambda }^1)`$, which we order so that $`\mathrm{\Lambda }>1`$.
* The matrix $`W`$ can be generated by using an elliptic quadratic Hamiltonian $`h(𝜻)`$ for an imaginary time $`i\tau _0`$. Note that this Hamiltonian generates renormalised dynamics in the section $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ and is therefore not simply a truncation of the full Hamiltonian in the transition state region.
* Canonical coordinates $`(Q,P)`$ are defined on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ so that
$$h(𝜻)=\frac{\alpha }{2}\left(Q^2+P^2\right)$$
and we have $`\mathrm{\Lambda }=\mathrm{e}^{\alpha \tau _0}`$.
The quantum analog of the classical generating Hamiltonian $`h(𝜻)`$ is denoted by $`\widehat{h}`$.
We can now write the tunnelling operator $`\widehat{𝒯}(E)`$ in the form
$$\widehat{𝒯}=\mathrm{e}^{\theta \tau _0\widehat{h}/\mathrm{}},$$
(7)
which, except for the prefactor $`\mathrm{e}^\theta `$, is an imaginary-time evolution operator generated by $`\widehat{h}`$. Since $`\widehat{h}`$ is harmonic we can explicitly construct its eigenstates $`|\phi _k`$, with $`k=0,1,\mathrm{}`$ and the corresponding eigensolutions of $`\widehat{𝒯}`$ are
$$\widehat{𝒯}|\phi _k=\tau _k|\phi _k$$
where the eigenvalues
$$\tau _k=\mathrm{e}^\theta \mathrm{\Lambda }^{(k+\frac{1}{2})}$$
(8)
are deduced simply by exponentiating the eigenvalues of $`\widehat{h}`$.
### III.2 Weyl symbol
It is shown in sC04 how closed form approximations can be deduced for phase space representations of $`\widehat{}(E)`$ as a result of substituting this exponentiated form for $`\widehat{𝒯}(E)`$ in (6) and resumming the geometric series $`\widehat{}=\widehat{𝒯}\widehat{𝒯}^2+\widehat{𝒯}^3\mathrm{}`$. This leads to an integral representation
$$\widehat{}(E)=\frac{1}{2i}_C\frac{\mathrm{e}^{\rho (\theta +\tau _0\widehat{h}/\mathrm{})}}{\mathrm{sin}\pi \rho }d\rho $$
(9)
for $`\widehat{}(E)`$, in which the contour $`C`$ ascends just to the right of the imaginary axis. Standard asymptotic approaches to this integral, such as the method of steepest descent, allow explicit asymptotic approximations to be written for $`\widehat{}`$ in various representations, including for the Weyl symbol.
These expressions are especially useful to understand the detailed structure of $`\widehat{}(E)`$ in phase space, but for the purposes of computing $`\widehat{}(E)`$ for the parameter regimes we consider here, it suffices to use an an eigenexpansion
$$\widehat{}(E)\underset{k}{}r_k|\phi _k\phi _k|$$
(10)
where
$$r_k=\frac{\tau _k}{1+\tau _k}.$$
(11)
The Weyl symbol of $`\widehat{}(E)`$ can, for example, be written as
$$𝒲_\widehat{}(𝜻,E)2\pi \mathrm{}\underset{k}{}r_k(E)\stackrel{~}{𝒲}_{kk}(𝜻)$$
where $`\stackrel{~}{𝒲}_{kk}(𝜻)`$ are the Wigner functions of the states $`|\phi _k`$ (and given analytically in Ripamonti for example). The tilde distinguishes these Wigner functions from those of the basis states $`|\psi _n`$ of the scattering operator, which are different.
The canonical coordinates $`(Q,P)`$ are centred on the initial condition for $`\gamma _E(t)`$ in $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ and are such that for energies just above threshold, the classically reacting region is circular in the $`(Q,P)`$ plane. In the original coordinate system $`(y,p_y)`$ these Wigner functions are translated, squeezed and rotated so that their level curves are aligned with the approximately elliptical reacting region. The resulting approximation for $`𝒲_\widehat{}(𝜻,E)`$ therefore describes an elliptically-shaped representation of the true quantum transmission problem
The harmonically approximated Weyl symbol $`𝒲_\widehat{}(𝜻,E)`$ is illustrated in Figure 3 for three energies at and just above threshold in the model potential with $`\mathrm{\Omega }=1`$ and $`\lambda =1/2=\mu `$. For comparison, illustrations of corresponding exact calculations can be found in the top row of Figure 2. In general we find that the harmonic approximation is in good quantitative agreement with exact results at and below threshold. The threshold case in Figure 3(a), for example, is indistinguishable from the corresponding exact result in Figure 2(a) at the level of graphical resolution used. As energy increases, the agreement deteriorates so that noticable differences are visible when $`E=1.15`$ (harmonic approximation in Figure 3(c) and exact calculation in Figure 2(c)). It should be emphasised, however, that even then, the harmonic approximation captures the essential qualitative features of $`𝒲_\widehat{}(𝜻,E)`$.
In order to make a closer comparison between exact and harmonic results, we show one-dimensional sections through the Weyl symbol in Figure 4. In each case $`𝒲_\widehat{}(𝜻,E)`$ is sampled along a horizontal line through the centre of the reacting region in $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ and plotted as a function of the $`y`$ coordinate. There is good quantitative agreement in cases (a) and (b) where the enrgy is at and just above threshold. At higher energies the harmonic approximation captures the support of the quantum-mechanically reacting region well but details of the Wigner function do not match at the centre of the reacting region. It should be remarked, however, that oscillations in the Weyl symbol are sensitive to nonlocal changes in phase space and discrepencies at the centre of the reacting region may not have a strong effect on averaged reaction probabilities as expressed in (4).
Similar one-dimensional sections are shown in Figure 5 which illustrate the harmonic approximation for different parameter sets. In each of these the energy is chosen so that the cumulative reactive flux
$$N_{\mathrm{cl}}(E)=\frac{1}{2\pi \mathrm{}}_{\mathrm{PODS}}𝐩d𝐪=\frac{E1}{\mathrm{}\sqrt{\mathrm{\Omega }^2+\lambda }}$$
is fixed (at the value 2). Figures 5(a), (b) and (c) show cases where $`\mu =0`$ and the potential has a symmetry in $`y`$. Figure 5(a) has a negative value of $`\lambda =1/2`$ for which an adiabatic approximation assuming fast transverse dynamics in the barrier region would not be expected to apply. Figure 5(b) is the separable case $`\lambda =0`$ and in Figure 5(c) we have $`\lambda =1/2`$. Note that separability does not confer a particular computational advantage in this approach, nor does it lead to particularly better accuracy of the approximation. In Figure 5(d), an example is shown in which $`\mu =1/2`$ and the potential is neither separable nor symmetric in $`y`$. The approximation works less well in that case. This is not unexpected because corrections to the harmonic approximation will be quartic rather than cubic in a symmetric problem but we note that there is still good agreement.
## IV Non-Linear Approximation
Although the harmonic approximation captures the essential qualitative features of quantum transmission and works well quantitatively near threshold, we can achieve greater range of applicability and significantly improved numerical agreement if we use fully nonlinear dynamics around the orbit $`\gamma _E(t)`$. The price to be paid for this improvement is that the resulting calculation is significantly more involved. The greatest impediment is that, although the tunnelling operator defined by nonlinear evolution can be routinely approximated semiclassically, we do not know at present how to write semiclassical approximations for the operator $`(1+\widehat{𝒯})^1`$ directly in terms of classical orbits. In this paper we simply use numerical inversion of the matrix representation of $`1+\widehat{𝒯}`$. Before describing this procedure, it is helpful to describe how nonlinear calculation is incorporated in the operator $`\widehat{𝒯}`$. This is done in section IV.1 below, followed by a description of the uniform calculation in section IV.2.
### IV.1 Primitive approximation
At energies below threshold the imaginary action of the orbit $`\gamma _E(t)`$ is positive, that is $`\theta >0`$, and the exponential prefactor $`\mathrm{e}^\theta `$ in (7) makes $`\widehat{𝒯}`$ small. We may therefore approximate the reaction operator directly by the tunnelling operator, giving
$$\widehat{}(E)\widehat{𝒯}(E),$$
which we refer to as the primitive approximation. The primitive approximation is easily extended beyond the immediate neighbourhood of $`\gamma _E(t)`$. Instead of letting $`\widehat{𝒯}`$ be the evolution operator corresponding to the classically linear evolution defined by $`W`$, as we did in the previous section, we let it be the quantum version of a nonlinear map in $`\mathrm{\Sigma }_R^{\mathrm{in}}`$.
Initial conditions near $`\gamma _E(t)`$ in $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ can be followed over a sequence of time evolutions similar to those of $`\gamma _E(t)`$ itself until they return to $`\mathrm{\Sigma }_R^{\mathrm{in}}`$, defining a surface-of-section mapping which we denote by
$$:\mathrm{\Sigma }_R^{\mathrm{in}}\mathrm{\Sigma }_R^{\mathrm{in}}.$$
As with conventional return maps, $``$ defines a canonical transformation on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$, except that it is complex, in general taking real initial conditions to complex images. Despite this complexity, the evolution has a quantum analog as an evolution operator, and this is the tunnelling operator $`\widehat{𝒯}`$.
With suitable modifications to take account of the complexity of the mapping, standard semiclassical approximations that are applied to evolution operators, such as the Van Vleck formula, can be used to approximate $`\widehat{𝒯}`$. Here we focus on an approximation derived in mB89 for the Weyl symbol of an operator (see also sC99 ). For the Weyl symbol of the operator $`\widehat{𝒯}`$ we write
$$𝒲_{\widehat{𝒯}}(𝜻,E)\frac{\mathrm{e}^{𝒜(𝜻,E)/\mathrm{}}}{\sqrt{\mathrm{det}(W_{AB}+I)/2}},$$
(12)
where $`𝒜`$ and $`W_{AB}`$ are calculated from a midpoint orbit $`𝜻_A𝜻_B`$ which is defined by the conditions
$`𝜻`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(𝜻_A+𝜻_B\right)`$
$`𝜻_B`$ $`=`$ $`(𝜻_A).`$ (13)
That is, the point $`𝜻`$ at which the Weyl symbol is to be evaluated is the midpoint of $`𝜻_A`$ and $`𝜻_B`$, where $`𝜻_A`$ evolves into $`𝜻_B`$ under the return map. The exponent $`𝒜(𝜻,E)`$ is such that
$$i𝒜(𝜻,E)=_{𝜻_A}^{𝜻_B}𝐩d𝐪p_y(y_By_A)$$
and the matrix $`W_{AB}`$ is a linearisation the map $``$ around the orbit $`𝜻_A𝜻_B`$.
It can be shown sc05 that the complex conjugate of the map $``$ is its inverse, $`^{}=^1`$, and from this a number of important symmetries follow which guarantee that $`𝒲_{\widehat{𝒯}}(𝜻,E)`$ is a real-valued function on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ that is peaked around the initial condition for $`\gamma _E(t)`$, which we denote by $`𝜻_0`$ in the following. On a formal level, the real-valuedness of $`𝒲_{\widehat{𝒯}}(𝜻,E)`$ follows from the observation that $`\widehat{𝒯}`$ is Hermitian which, as discussed in sC99 , is a quantum analog of the property $`^{}=^1`$. It is instructive, however, to see how the real-valuedness of the semiclassical approximation to $`𝒲_{\widehat{𝒯}}(𝜻,E)`$ follows directly from the symmetries of the midpoint orbit.
First we note that, given a midpoint orbit $`𝜻_A𝜻_B`$ for a real-valued $`𝜻`$, then $`𝜻_B^{}𝜻_A^{}`$ is also a midpoint orbit (for the same $`𝜻`$). This can be seen by conjugating the relations in (IV.1) and using $`𝜻_B^{}=[((𝜻_A)]^{}=^{}(𝜻_A^{})=^1(𝜻_A^{})`$ to deduce that $`𝜻_A^{}=(𝜻_B^{})`$. It turns out in fact that these two midpoint orbits coincide, so
$$𝜻_B=𝜻_A^{}$$
and
$$\mathrm{Re}𝜻_B=𝜻=\mathrm{Re}𝜻_A.$$
This is easily confirmed for the linearised map (replacing $``$ by multiplication by $`W`$ and using $`W^{}=W^1`$) and therefore holds for the nonlinear map if $`𝜻`$ is close enough to $`𝜻_0`$. The condition $`𝜻_B=𝜻_A^{}`$ can only be violated if a bifurcation is encountered and a more detailed analysis shows that this corresponds to the condition $`\mathrm{det}(W_{AB}+I)=0`$, which would lead to a caustic in (12). We will assume in this paper that no such caustics are encountered in the region of $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ which dominates reaction.
An example of a full trajectory corresponding to a typical midpoint orbit is illustrated in Figure 6. Because the initial conditions are complex, the coordinates of the trajectory are generically complex over its length, even along the segments which have been obtained by deformation of the real segments of $`\gamma _E(t)`$. A consequence of the symmetry $`𝜻_B=𝜻_A^{}`$, however, is that the time contours can be chosen so that the second half of the trajectory reverses the complex conjugate of the first half. A projection onto real configuration space, for example, is self retracing. We find as a result that the action is purely imaginary and the exponent $`𝒜(𝜻,E)`$ is a positive real number. We also find that $`W_{AB}^{}=W_{AB}^1`$ and because $`\mathrm{det}W_{AB}=1`$ this means that the amplitude term $`\mathrm{det}(W_{AB}+I)`$ in (12) is real (and positive). Therefore $`𝒲_{\widehat{𝒯}}(𝜻,E)`$ is a positive real-valued function with a maximum at $`𝜻_0`$ (for which we have $`𝒜(𝜻_0,E)=K_0(E)`$).
The harmonic approximation can be recovered by expanding the exponent to second order about $`𝜻_0`$ and approximating $`W_{AB}`$ by $`W`$, giving sC04
$$𝒲_{\widehat{𝒯}}(𝜻,E)\frac{\mathrm{e}^{\theta [(2/\beta )\mathrm{tanh}\beta /2]\tau _0h(𝜻)/\mathrm{}}}{\mathrm{cosh}\beta /2}.$$
(14)
A comparison is given in Figure 7 between this approximation, the fully nonlinear approximation of (12) and exact results. Although the harmonic approximation works well near the maximum of the Weyl symbol, the fully nonlinear result works better over a larger range. We note however that there is qualitative deviation even from the nonlinear approximation in the deep tunnelling regime where the exact calculation shows significant oscillations not captured by the harmonic or nonlinear approximations. Similar oscillatory structure, or “fringed tunnelling” in the scattering matrix has been explained in aY00 ; Takahashi ; TYK on the basis of nonintegrable complex dynamics and has been shown to involve mechanisms that also show up in chaotic tunnelling. It seems likely that the oscillations in Figure 7 have a similar origin but we have not preformed a detailed analysis. It will be an interesting problem in the future to combine the inherently nonintegrable mechanism in aY00 ; Takahashi ; TYK with the uniform approximations illustrated here.
### IV.2 Uniform approximation
Although we now have an explicit closed-form semiclassical approximation for $`\widehat{𝒯}`$, no equivalent result is currently available for the uniformisation $`\widehat{𝒯}/(1+\widehat{𝒯})`$ because inversion of the operator $`1+\widehat{𝒯}`$ cannot be done simply. In this paper we simply adopt a hybrid approach which combines semiclassical approximation of $`\widehat{𝒯}`$ with numerical inversion of $`1+\widehat{𝒯}`$. Although not a fully semiclassical method, this will allow us to verify that (6) gives an accurate reproduction of quantum transmission.
We first represent $`\widehat{𝒯}`$ as a matrix in the same asymptotic basis $`|\psi _n`$ as used for the scattering matrix. We denote individual matrix elements by $`T_{nm}=\psi _n|\widehat{𝒯}|\psi _m`$ and compute them using the Wigner-Weyl calculus by writing
$$T_{nm}(E)=𝒲_{nm}(𝜻)𝒲_{\widehat{𝒯}}(𝜻,E)d𝜻$$
(15)
and approximating $`𝒲_{\widehat{𝒯}}(𝜻,E)`$ using (12). We emphasise that while $`\widehat{𝒯}`$ is almost diagonal in the basis $`|\phi _k`$ of eigenstates of the generating Hamiltonian $`\widehat{h}`$, the same is not true in the basis $`|\psi _n`$ unless the potential is separable. The integral is performed numerically and the resulting $`M\times M`$ matrix for $`\widehat{𝒯}/(1+\widehat{𝒯})`$, whose elements are denoted $`R_{nm}^{\mathrm{sc}}(E)`$, is also computed numerically. This integration is not difficult since a grid on $`\mathrm{\Sigma }_R^{\mathrm{in}}`$ is easily filled by using Newton integration to step the midpoint orbit, starting with the known solution corresponding to $`\gamma _E(t)`$ at $`𝜻_0`$ (for which $`𝜻_A=𝜻_0=𝜻_B`$). Once the elements $`R_{nm}^{\mathrm{sc}}(E)`$ are known, the Weyl symbol for $`\widehat{𝒯}/(1+\widehat{𝒯})`$ is obtained by replacing $`R_{nm}(E)`$ with $`R_{nm}^{\mathrm{sc}}(E)`$ in (3).
We find excellent agreement between the semiclassically computed Weyl symbol $`𝒲_\widehat{}(𝜻,E)`$ and the exact result. A nonlinear version of Figure 3 is indistinguishable from the exact results shown in the top row of Figure 2 and therefore not shown. Instead we compare in Figure 8 horizontal slices of the Weyl symbol through the reacting region, in the same manner as in Figure 4. The potential used is the same as in Figure 4 and the energies treated extend somewhat higher above the threshold. The exact and semiclassical results cannot be distinguished at the resolution used, so the difference scaled by a factor of $`10`$ is also shown. We find similarly good agreement for other parameter sets we ahve investigated and we note that the quality of the approximation does not require special features of the classical dynamics such as symmetry, separability or adiabatic separation of transverse from reaction degrees of freedom.
We should remark that the current hybrid implementation of the nonlinear calculation is cumbersome and is harder to apply further above the barrier where the larger region of integration demands that we extend the midpoint trajectory deeper into complex phase space. We have not, for example, reproduced the results on the second row of Fig. 2 using this method. The purpose of this calculation is to show that the nonlinear uniform result derived in sc05 provides an accurate description of $`\widehat{}(E)`$ in the model considered and that the method therefore deserves further exploration. Even though numerical inversion was used in applying the formalism, it is built entirely on a semiclassical approximation for the tunnelling operator $`\widehat{𝒯}(E)`$ and we expect that any subsequent fully semiclassical implementation will be equally accurate.
We also remark that the current hybrid method is theoretically clumsy and obscures somewhat the deeper connections between the quantum results and the underlying classical geometry. For example, it would be especially interesting to characterise the behaviour of $`𝒲_\widehat{}(𝜻,E)`$ at the boundary of the classically reacting region where trajectories approach the PODS (or NHIM in higher dimensions) along its stable manifold and where the classical reaction probability drops sharply from $`1`$ to $`0`$. Such an analytical approximation was found in sC04 for the harmonic version in which $`𝒲_\widehat{}(𝜻,E)`$ is approximated as an integral of the Airy function near the boundary of classical reaction. Investigation is currently underway into a method to derive similar results in the nonlinear case on the basis of generating $`\widehat{𝒯}`$ as in equation (7), but with an anharmonic generator $`\widehat{h}`$ computed using classical normal form theory. Ultimately it should be possible to describe explicitly how the quantum reaction probability varies across the boundary of the classically reacting region in terms of trajectories which approach the complexified PODS along its stable manifold and evolve along it before returning to the asymptotic Poincaré section.
## V Conclusions
We have successfully treated quantum transmission across a multidimensional barrier using uniform semiclassical approximation. The method applies generically around a threshold energy and does not rely on specific features of the classical dynamics such as separability or the existence of action angle variables. In its fully nonlinear incarnation the method gives an accurate description of reaction probabilities in phase space and makes a striking connection between the quantum scattering problem and the geometry of classical reaction. We expect that it will work equally well in higher-dimensional problems, even in cases where the incoming states are chaotic in the transverse dynamics.
Although we have shown that the fully nonlinear version works well, in doing so we have resorted to numerical methods which are not in the spirit of semiclassical approximation. The theory therefore needs further development in order to achieve a fully semiclassical description of the emerging reacting region. One promising approach which is currently under investigation is to use classical normal form theory to generate dynamics around the orbit $`\gamma _E(t)`$ using a nonlinear extension of the generator $`h(𝜻)`$. Many of the explicit analytical approaches used in the harmonic case might then be adapted to the nonlinear approximation. In particular this is expected to produce a detailed analytical description of the Weyl symbol at the boundary of the reacting region which calls on intrinsic geometrical features of (the stable manifold of) the NHIM.
A second aspect of the calculation which deserves further attention is the treatment of rotational degrees of freedom in fully three-dimensional models of chemical reaction. Although at one level this is simply a question of applying the results here individually to symmetry-reduced phase spaces for given angular momentum quantum numbers, there are interesting and nontrivial problems in describing the quantum-classical correspondence compactly in operator form. This is an especially interesting issue for reactions which proceed through a collinear mechanism since the collinear configurations are a singular part of the classical reduction process.
Acknowledgements
CSD is supported by an EPSRC studentship and SCC acknowledges support by the European Network MASIE, Contract No. HPRN-CT-2000-00113..
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# Type-II Quantum Algorithms
## 1 Introduction
The existence of quantum computers with thousands of error resistant entangled qubits is many years, if not many decades, in the future. However, the success of liquid state nuclear magnetic resonance spectroscopy (NMR) quantum computation means that few ($`10`$) entangled qubit quantum computers are available now . This approach exploits the fact that the spin active nuclei in the molecules of a liquid are decoupled from each other and so each nuclear spin may act as a separate quantum bit, and each molecule as a separate few-qubit quantum computer. The states of these nuclear spins may be addressed by radio-frequency pulses. The existence of Avagadro’s number of such spins enables the ensemble-averaged results of quantum computations to be obtained via thermodynamic averaging.
The resonant frequency depends on the local magnetic field, and so placing the fluid in a non-uniform field enables one to separately address different regions of fluid . Each region of fluid then acts as a separate ensemble of quantum computers, and the idea of using a large array of these computers was first proposed in . Quantum coherence is maintained within each node on the array, and the arrays are measured after a few quantum coherent operations and then classical information representing the ensemble average of $`10^{18}`$ measurements is exchanged between nodes, providing the initial state for the next coherent quantum computations on each node. This situation is illustrated in Figure 1.
Type II quantum computers were introduced in the context of physical modelling, and all subsequent work has been focussed towards implementation of some lattice based algorithm for (usually classical) physics . Algorithms have been proposed for fluid dynamics , for the diffusion equation , the Burgers equation , for the nonlinear Schroedinger and Korteweg-de Vries equations in one-dimensions , and for MHD turbulence in one dimension . Some of these algorithms have also been implemented on NMR quantum computers .
## 2 Type II Quantum Lattice Gases
The lattice structure of the array of quantum processors and the similarity of the two step quantum computation/classical communication are very similar to the collision/propagation evolution of lattice-gas and lattice-Boltzmann algorithms . Because of this similarity Type-II algorithms are described in the literature as quantum lattice-gases. One must distinguish these Type-II models from the quantum lattice gases for the Dirac and Schroedinger equation, designed to realize exponential improvements in performance on fully coherent quantum computers . Type II quantum lattice-gases perform the collision step using quantum computation, but then sample the results of those computations before performing the subsequent propagation. Type II simulations rely on ensemble averaging, and these simulations exploit the natural existence of an ensemble within the NMR quantum computer to estimate the average values of the qubits at a site. Note carefully, however, that only information about classical averages is communicated between sites. Information about two-particle and higher correlations is not preserved by this measurement and communication step, and so the results of a Type II quantum lattice gas simulation are equivalent to the results of a classical lattice Boltzmann simulation. Indeed, they most closely resemble the early lattice-Boltzmann algorithms which used ensemble averaged lattice gas collision rules at each site .
All of the models implemented so far on NMR quantum computers have two qubits per site in one dimension. The state of the model at a site is given by a pair of single particle distribution functions, $`(f_1,f_2)`$, such that the probability that bit one is one is $`f_1`$ and the probability that bit two is one is $`f_2`$. Each separate quantum processor therefore consists of a pair of coherent qubits. The state of these qubits is initialized using $`(f_1,f_2)`$ as follows:
$$|\psi (t=0)=\sqrt{(1f_1)(1f_2)}|00+\sqrt{(1f_1)f_2}|01+\sqrt{(f_1(1f_2)}|10+\sqrt{f_1f_2}|11$$
(1)
The probabilities of the four possible lattice gas states at a site, $`\{00,01,10,11\}`$ are $`𝐧=((1f_1)(1f_2),(1f_2)f_1,f_2(1f_1),f_1f_2)`$ in the Boltzmann approximation. The diagonal components of the density matrix $`|\psi (t=0)\psi (t=0)|`$ are equal to these probabilities. The collision step of a stochastic lattice gas model is given by the multiplication of the vector n by a matrix, A, satisfying normalization and semi-detailed balance (a doubly stochastic or complete matrix). In our Type-II simulation we wish to reproduce the effect of this matrix multiplication by conjugating the density matrix by the unitary operator corresponding to the quantum mechanical part of our Type-II algorithm. The effect of this conjugation on the diagonal elements of the density matrix must be the same as the effect of matrix multiplication of vector $`𝐧`$ by the matrix $`A`$.
Now consider a Type-II quantum computer with $`b`$ qubits per node. At the beginning of each simulation step the qubits are initialized with a density matrix whose entries correspond to the probabilities of the qubit state in the molecular chaos approximation. After each local unitary update, the qubits are measured and because the existence of an ensemble is presumed in all Type-II algorithms the probability that qubit $`i`$ is one is associated with the single particle distribution function $`f_i(𝐱)`$ at site $`𝐱`$ and vector $`i`$. During the collision step the unitary action $`U`$ is applied to this density matrix and it evolves to:
$$\rho _{mr}\underset{mr}{}U_{pm}\rho _{mr}\left[U^{}\right]_{rq}$$
(2)
The density matrix formulation of quantum mechanics contains within it classical states, real linear convex combinations of classical states representing classical probability distributions, pure quantum states and mixed quantum states. For this reason Type-II algorithms may be constructed which reproduce deterministic reversible lattice gas rules, stochastic lattice gas rules, and new models which have collision operators which do not correspond to the classical average over some lattice gas. In all cases, because the measurement and averaging step precedes the classical communication step, the overall algorithm is equivalent to a classical lattice-Boltzmann scheme with collision operator given by (2).
Taking our update for the density matrix and restricting it to an update of the diagonal components of the density matrix:
$$\rho _{pp}^{}=\underset{mr}{}U_{pm}\rho _{mr}\left[U^{}\right]_{rp}$$
(3)
Where there is no implicit sum on repeated indices. Separating the action on the diagonal and off diagonal components of the density matrix:
$$\rho _{pp}^{}=\underset{m}{}U_{pm}U_{pm}^{}\rho _{mm}+\underset{mr}{}U_{pm}U_{pr}^{}\rho _{mr}$$
(4)
If the second term here is zero then the action of the unitary operator $`U`$ on the diagonal components of $`\rho `$ is identical to the action of a matrix $`A`$ whose components are the modulus squared of the components of $`U`$.
$$A_{pm}=U_{pm}U_{pm}^{}$$
(5)
The unitarity of $`U`$ implies that the matrix $`A`$ will obey normalization and semi-detailed balance.
It remains to formulate a sufficient constraint that the second term above is zero. A sufficient condition for the Type-II algorithm to reproduce the action of a stochastic lattice-gas collision operator whose average is the matrix $`A`$ is that:
$$\text{Re}[U_{pm}U_{pr}^{}]=0m<r,p$$
(6)
One realization of a Type-II algorithm on an NMR quantum computer implemented a simple stochastic lattice-gas algorithm for the diffusion equation in one-dimension, where the matrix $`A`$ is:
$$A=\frac{1}{2}\left(\begin{array}{cccc}2& 0& 0& 0\\ 0& 1& 1& 0\\ 0& 1& 1& 0\\ 0& 0& 0& 2\end{array}\right)$$
(7)
In this model is implemented using a unitary matrix $`U`$:
$$U=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& \frac{1i}{2}& \frac{1+i}{2}& 0\\ 0& \frac{1+i}{2}& \frac{1i}{2}& 0\\ 0& 0& 0& 1\end{array}\right)$$
(8)
It is clear by inspection that the modulus squared of the components of $`U`$ give the components of $`A`$. The condition (6) is also satisfied.
If a Type-II algorithm does not obey constraint (6) the action of $`U`$ on the diagonal components of the density matrix may not be described by a doubly stochastic matrix. The new values of the diagonal components will be modified by terms involving the off diagonal components of the density matrix, which in turn will involve the square roots of the single particle distribution functions. The action on the diagonal components of an arbitrary density matrix will not be linear and so (6) is a necessary as well as sufficient condition. Models violating (6) fall outside the class of lattice Boltzmann models which may be considered as averages over some underlying lattice-gas model in the Boltzmann approximation, and a detailed analysis of such lattice-Boltzmann models remains an open problem. An example of such a model is the lattice-gas model for the Burgers equation , which does not satisfy semi-detailed balance. A Type-II algorithm has been proposed for this model, and implemented in NMR .
In the case where the ensemble of molecular quantum computers at each node is initialized, not in the pure state (1), but in the completely mixed state with a density matrix with diagonal components equal to $`𝐧`$ and off diagonal components which are all zero, any unitary action reproduces the collision operator of a stochastic lattice gas satisfying detailed balance and normalization. In this case the Type-II computation is reduced to simple molecular computing of a Boltzmann approximation to a lattice-gas evolution, and no novel models are possible.
## 3 Conclusions
The characterization of current Type-II simulations as equivalent to classical lattice-Boltzmann methods illustrates the limited utility of Type-II quantum computation as presently defined. These algorithms provide at best a constant speedup and require an ensemble in order to accurately sample the classical information available through measurement which is exchanged between quantum processors. While applying techniques from medical imaging may make the realization of very large ($`>2048^3`$) grids possible, such techniques require a three dimensional field gradient to be useful. While the implementation of such field gradients is non-trivial, such large grid sizes would challenge the largest lattice Boltzmann simulations on classical supercomputers to date.
One of the strengths of Type II quantum computation is that it takes advantage of the ensemble inherent in NMR implementations. Distribution functions are obtained by averaging over the $`10^{18}`$ molecules in each small region . Of course, this necessity to obtain the average values by sampling becomes an Achilles heel once one considers more scalable architectures (such as superconducting or solid state optically addressed implementations) where the necessity for sampling requires multiplication of the hardware. Although the spirit of Type II computing is that many small quantum computers are likely to be available sooner than a single large machine, the more scalable architectures are unlikely to reproduce the huge ensemble freely available with liquid state NMR.
From the above analysis we can see that unitary matrices obeying (6) implement a doubly-stochastic Markov matrix action on the diagonal components of the density matrix. The size of this Markov matrix grows exponentially with the number of qubits. It is natural to ask whether one may use these Type-II techniques to solve Markov chain problems out of the reach of classical computers. For example, if we could address and control a pseudo-pure state of an ensemble of $`100`$ qubit systems we could address Markov problems of size $`2^{100}`$, which are out of reach of any conceivable classical digital computer. However, the limit on the size of Markov process we may implement is not given by the number of qubits but by our ability to read out the diagonal components of the density matrix by sampling from the ensemble. We require at least one member of the ensemble per state of the system in order to do this, which implies that for $`100`$ qubits one requires $`10^7`$ moles of material in our ensemble. Unfortunately this is not feasible. Current liquid state NMR realizations utilize $`10^{18}`$ molecules, implying a maximum size of problem of $`2^{60}`$, implying that the maximum Markov problem size occurs for such an ensemble with $`60`$ qubits. Such problems are classically intractable above at most $`2^{30}`$, implying that there is a range of problems which can be implemented on small NMR ensemble quantum computers which can coherently address $`3060`$ spins which are classically intractable for conventional digital computers.
This observation reinforces the point that Type-II quantum computation is really an interesting mixture of quantum computing with classical molecular computing. The approach benefits as much from the presence of a large number of realizations of the computation in a molecular ensemble as it does from the exponential growth of the dimension of the Hilbert space of the individual quantum processors.
This material is based on work supported in part by the U.S. Army Research Laboratory and the U.S. Army Research Office under grant number W911NF-04-1-0334, in part by the U.S. Air Force Office of Scientific Research under grant number FA9550-04-1-0176, and in part by the DARPA QuIST Program under AFOSR grant number F49620-01-1-0566.
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# Excitons, biexcitons and trions in self-assembled (In,Ga)As/GaAs quantum dots: Recombination energies, polarization and radiative lifetimes versus dot height
## I Introduction
Single-dot spectroscopy makes it possible to probe dot-to-dot changes in the excitonic properties of self-assembled InGaAs/GaAs quantum dots. smith\_PRL\_2005 ; ware\_PhysicaE\_2005 ; skolnick\_ARMR\_2004 ; guffarth\_PhysicaE\_2004 ; rodt\_papers ; urbaszek\_PRL\_2003 ; finley\_PRB\_2001 ; findeis\_SSC\_2000 ; landin\_papers Both single-particle and many-particle aspects of these properties depend non-trivially on the quantum dots size and shape,books+review reflecting not only simple quantum-confinement physics, but also electronic structure effects such as interband, intervalley, spin-orbit and strain-induced state coupling. The description of these effects require an atomistic multi-band approach. wang\_PRB\_1999 ; wang\_APL\_2000 ; zunger\_pss.b\_2001 ; TB\_method ; sheng\_PRB\_2005 Here we adopt a method that is based on screened pseudopotentials and the configuration-interaction approach, and address the changes with height of recombination (emission) energies, polarization and radiative lifetimes of various neutral and charged excitons: the neutral exciton ($`X^0`$), negatively- ($`X^{}`$) and positively-charged ($`X^+`$) excitons, and the biexciton ($`XX^0`$) in lens-shaped, self-assembled alloyed In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs quantum dots. Our predictions compare reasonably well with available spectroscopic data. We also compare our findings in In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs dots with those in pure, non-alloyed InAs/GaAs dots.
## II Energetics of the monoexciton, charged trions and biexciton
### II.1 Electronic structure of the excitonic manifolds
We describe the basic electronic structure of the excitonic manifolds (Fig. 1) before describing recombination processes. All excitonic states are based on mixing and excitations of the single-particle states $`\{h_0`$, $`h_1`$, $`h_2`$, $`\mathrm{}\}`$ and $`\{e_0`$, $`e_1`$, $`e_2`$, $`\mathrm{}\}`$, for holes and electrons, respectively. These states are solutions to the single-particle Schrödinger equation
$$\left\{\frac{1}{2}^2+V_{ext}(𝐑)+V_{scr}(𝐑)\right\}\psi _i=_i\psi _i,$$
(1)
where both the external (pseudo) potential $`V_{ext}(𝐑)`$ due to the ion-ion and ion-electron interaction and the screening response to such external potential $`V_{scr}(𝐑)`$ are expressed as a superposition of screened atomic pseudopotentials
$$V_{ext}(𝐑)+V_{scr}(𝐑)=V_{SO}+\underset{l}{}\underset{\alpha }{}v_\alpha [𝐑𝐑_l^{(\alpha )};\mathrm{Tr}(\stackrel{~}{\epsilon })].$$
(2)
Here, $`V_{SO}`$ is a non-local spin-orbit interaction;williamson\_PRB\_2000 $`v_\alpha `$ is a screened pseudopotential for atom of type $`\alpha `$ that depends on strain, and it has been fitted to bulk properties of GaAs and InAs, including bulk band structures, experimental deformation potentials and effective masses, as well as local-density-approximation (LDA)-determined band offsets.williamson\_PRB\_2000 The single-particle Shrödinger equation \[Eq. (1)\] includes not only quantum-confinement effects (as in simple, one-band particle-in-a-box models), but also multi-band coupling (light hole, heavy hole, conduction); inter-valley ($`\mathrm{\Gamma }`$-$`X`$-$`L`$) coupling; and spin-orbit coupling. Strain effects are present through the relaxation, via a valence force field,williamson\_PRB\_2000 of the atomic positions $`\{𝐑_l^{(\alpha )}\}`$ within the simulation supercell (quantum dot+GaAs-matrix); and directly “felt” by the potential $`v_\alpha [𝐑𝐑_l^{(\alpha )};\mathrm{Tr}(\stackrel{~}{\epsilon })]`$. States $`\{h_0`$, $`h_1`$, $`h_2`$,$`\mathrm{}\}`$ and $`\{e_0`$, $`e_1`$, $`e_2`$,$`\mathrm{}\}`$ form a basis for the excitonic states. We indicate the dominant single-particle configuration of each excitonic level by $`e_i^pe_i^{}^p^{}h_j^qh_j^{}^q^{}`$, where $`p`$, $`p^{}`$, $`q`$, and $`q^{}`$ (=0, 1, 2) indicate, correspondingly, the occupation of the electron levels $`e_i`$ and $`e_i^{}`$, and hole levels $`h_j`$ and $`h_j^{}`$. Figure 1 illustrates the excitonic manifolds using a $`75`$Å-tall lens-shaped In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs quantum dot (base diameter $`b=252`$Å), as obtained from configuration-interaction calculations based on the pseudopotential single-particle description.
Monoexciton. The monoexciton $`X^0`$ has a ground-state $`|\mathrm{\Psi }^{(0)}(X^0)`$ created by occupying $`h_0`$ and $`e_0`$, denoted $`e_0^1h_0^1`$ \[Fig 1(a)\]. However, configuration-interactionfranceschetti\_PRB\_1999 (CI) also mixes into $`|\mathrm{\Psi }^{(0)}(X^0)`$ other states such as $`e_0^1h_1^1,e_0^1h_2^1,\mathrm{}`$. This is done by expanding the monoexciton states $`\{|\mathrm{\Psi }^{(\nu )}(X^0)\}`$ ($`\nu =0,1,2,\mathrm{}`$) in a basis of Slater determinants (configurations) $`\{|\mathrm{\Phi }(X^0)\}`$ constructed in the subspace of $`N_e`$ and $`M_h`$ electron and hole confined single-particle levels, respectively:
$$|\mathrm{\Psi }^{(\nu )}(X^0)=\underset{\kappa }{}C_\kappa ^{(\nu )}(X^0)|\mathrm{\Phi }_\kappa (X^0),$$
(3)
where $`C_\kappa ^{(\nu )}(X^0)`$ are the CI coefficients and $`\kappa `$ is a composite index that labels each Slater determinant. Since $`e_0`$ and $`h_0`$ are each two-fold degenerate due to spin, $`e_0^1h_0^1`$ has four-fold degeneracy at the single-particle level of Eqs. (1) and (2). We next allow Coulomb electron-electron and hole-hole interactions. Direct Coulomb is given by $`J_{ij;ji}^{(\mu \mu )}`$ and exchange by $`J_{ij;ij}^{(\mu \mu )}`$, with the Coulomb scattering matrix elements given by
$$J_{ij;kl}^{(\mu \mu ^{})}=d𝐑d𝐑^{}\frac{\left[\psi _i^{(\mu )}(𝐑)\right]^{}\left[\psi _j^{(\mu ^{})}(𝐑^{})\right]^{}\left[\psi _k^{(\mu ^{})}(𝐑^{})\right]\left[\psi _l^{(\mu )}(𝐑)\right]}{ϵ(𝐑,𝐑^{})|𝐑𝐑^{}|}.$$
(4)
Here $`\mu ,\mu ^{}=e,h`$; and $`ϵ(𝐑,𝐑^{})`$ is a microscopic, phenomenological dielectric constant.resta\_PRB\_1977 We also allow electron-hole direct Coulomb interaction $`J_{ij;ji}^{(eh)}`$ and electron-hole exchange
$$K_{ij;kl}^{(eh)}=d𝐑d𝐑^{}\frac{\left[\psi _i^{(h)}(𝐑)\right]^{}\left[\psi _j^{(e)}(𝐑^{})\right]^{}\left[\psi _k^{(e)}(𝐑)\right]\left[\psi _l^{(h)}(𝐑^{})\right]}{ϵ(𝐑,𝐑^{})|𝐑𝐑^{}|}.$$
(5)
Inclusion of these Coulomb interactions splits $`e_0^1h_0^1`$ into four distinct levels:bester\_PRB\_2003 ; bayer\_PRB\_2002 The lowest two are spin-forbidden (“dark”) in the absence of spin-orbit coupling, and the highest two are allowed (“bright”). The bright-dark splitting of the ground-state levels shown in Fig. 1(a) is $`84\mu \mathrm{eV}`$. The magnitude of this splitting increases up to $`178\mu \mathrm{eV}`$ for a $`20`$Å-tall dot. At $`T0\mathrm{K}`$ only the dark states are populated, thus, the transition to the ground state $`e_0^0h_0^0`$ is long-lived. Figure 1(a) (shaded area) shows transitions $`\omega _{20}`$ and $`\omega _{30}`$ from the bright states to $`e_0^0h_0^0`$. The low-lying excited states of the monoexciton correspond to excitations of the hole, i.e. $`e_0^1h_1^1`$, $`e_0^1h_2^1`$ and $`e_0^1h_3^1`$, due to the much smaller spacing of hole single-particle energy levels compared to that of the electrons. Thus, the spacing of the excited states fingerprint the hole energy level structure. Figure 1(a) shows the first twelve excited states that are derived from $`e_0^1h_1^1`$, $`e_0^1h_2^1`$ and $`e_0^1h_3^1`$. Each of these states are four-fold degenerate and their fine structure is shown schematically. We find that three out of the four levels that arise from $`e_0^1h_2^1`$ are optically allowed. In particular, one of the levels, indicated as $``$, emits light that is polarized along $`[001]`$, while the remaining two emit in-plane polarized light. These excited levels become optically forbidden as the dot becomes flatter.
Biexciton. In contrast to the four levels comprising the monoexciton, the biexciton $`XX^0`$ has a singly-degenerate, closed-shell ground state ($`e_0^2h_0^2`$) that is bright. Thus, even at $`T0\mathrm{K}`$ both emissions $`\omega _{03}`$ and $`\omega _{02}`$ \[Fig. 1(b); shaded area\] of the biexciton are fast ($`\mathrm{ns}`$). The biexciton has a non-trivial ladder of excited states that is determined primarily by the relative magnitude of the direct Coulomb hole-hole interaction and the single-particle energy splittings of the hole states. The ladder begins with states derived from $`e_0^2h_0^1h_1^1`$ at about $`6\mathrm{meV}`$ above the ground state and follows with $`e_0^2h_0^1h_2^1`$ and $`e_0^2h_0^1h_3^1`$ at about $`11\mathrm{meV}`$ and $`16\mathrm{meV}`$, respectively. Due to the two-fold Kramers degeneracy of the hole levels, $`e_0^2h_0^1h_1^1`$ is four-fold degenerate at the single-particle level, these states split into four distinct states due to hole-hole and electron-hole exchange \[Fig 1(b)\]. Similarly, the four states in $`e_0^2h_0^1h_2^1`$ split in two groups of two. In this case, remarkably, the splitting is about twice as big as that in $`e_0^2h_0^1h_1^1`$ and nearly five times bigger than in $`e_0^2h_0^1h_3^1`$ ($`500\mu \mathrm{eV}`$). Similarly to the monoexciton case, the lowest split-off pair in $`e_0^2h_0^1h_2^1`$ is optically active and emits light polarized along $`[001]`$ \[$``$, Fig. 1(b)\]. Further, these states become darker as the dot becomes flatter, being forbidden at $`h=20`$Å. The subsequent excited state in the ladder derives from the singly-degenerate closed-shell state $`e_0^2h_1^2`$ and is closely-spaced with $`e_0^2h_0^1h_3^1`$. Note that albeit the splitting between these states is small, the state derived from the closed-shell configuration $`e_0^2h_1^2`$ is relatively “inert” in that it mixes weakly with other states. In particular the weight of this configuration in the CI expansion is $`83\%`$. Below, we discuss the energy, polarization and lifetime of $`\omega _{03}`$ and $`\omega _{02}`$.
Trions. The negatively- ($`X^{}`$) and positively-charged ($`X^+`$) trions have ground states that are bright, two-fold degenerate and arise from occupying $`e_0^2h_0^1`$ and $`e_0^1h_0^2`$, respectively. Due to configuration-interaction mixing, the ground states of $`X^{}`$ and $`X^+`$ mix with $`e_0^2h_1^1`$ and $`e_0^1h_0^1h_1^1`$, respectively. As in the case of $`XX^0`$, even at $`T0\mathrm{K}`$ the recombination $`\omega _{00}`$ of $`X^{}`$ \[Fig. 1(c)\] into state $`e_0^1h_0^0`$ and that of $`X^+`$ \[Fig. 1(d)\] into state $`e_0^0h_0^1`$ is fast. For both trions, while there are only two dipole-allowed transitions in the absence of spin-orbit coupling, we predict four allowed transitions. The first few excited states of $`X^{}`$ correspond to occupying states derived primarily from $`e_0^2h_1^1`$, $`e_0^2h_2^1`$ and $`e_0^2h_3^1`$ \[Fig. 1(c)\]. (Naturally, these states are mixed in with other states of $`X^{}`$.) These excited states are two-fold degenerate and lie, correspondingly, about $`10\mathrm{meV}`$, $`13\mathrm{meV}`$ and $`20\mathrm{meV}`$ above the ground state, as shown in Fig. 1(c). Note the similarity of the low-lying excited states energy spacing with that in $`X^0`$. In addition, note that one of the two states derived from $`e_0^2h_2^1`$ is optically active and emits light polarized along $`[001]`$ \[$``$ in Fig. 1(c)\]. Again, these states become optically forbidden for a $`20`$Å-tall dot. The excited states of $`X^+`$ are more complex and present a high density of states \[Fig. 1(d)\]; a consequence of the closely spaced ($`5\mathrm{meV}`$) single-particle hole levels: In contrast to $`X^{}`$, all low-lying excited states show fine-structure splitting due to electron-hole exchange. The first excited states at about $`5\mathrm{meV}`$ corresponds to occupying $`e_0^1h_0^1h_1^1`$, due to the two-fold degeneracy of each of these single-particle states, there are eight excited states that split in two groups of four, due to hole-hole exchange. Occupying $`e_0^1h_0^1h_2^1`$ leads to the next eight excited states around $`10\mathrm{meV}`$, which are also split in two quartets. In this case, there is an optically active (polarization $`[001]`$) pair of excited states that belong to the lower-energy quartet \[$``$, Fig. 1(d)\]. As in the other excitons, these states become dark as the dot height decreases. The next ten excited states about $`14\mathrm{meV}`$ are a mixture of eight configurations derived from $`e_0^1h_0^1h_3^1`$ and two from $`e_0^1h_1^2`$. Because $`e_0`$ is half-filled, $`e_0^1h_1^2`$ is very “reactive” in the sense that it heavily mixes via configuration interaction. Note that this is significantly different from the case of configuration $`e_0^2h_0^2`$ in the biexciton \[Fig. 1(b)\]. Higher excited states correspond to occupying $`e_0^1h_1^1h_2^1`$ and $`e_0^1h_1^1h_3^1`$. Below, we discuss the energy, as well as polarization and lifetime of $`\omega _{00}`$ for $`X^{}`$ and $`X^+`$.
### II.2 Recombination energies
We define the recombination energy upon recombination of an electron-hole pair in $`\chi ^q`$ as the difference between the total energies of the initial state $`|\mathrm{\Psi }^{(i)}(\chi ^q)`$ and the final state $`|\mathrm{\Psi }^{(f)}(\chi ^q1)`$. Namely,
$$\begin{array}{ccc}\hfill \omega _{i0}(X^0)& =& E^{(i)}(X^0)\hfill \\ \hfill \omega _{if}(X^{})& =& E^{(i)}(X^{})_f^{(e)}\hfill \\ \hfill \omega _{if}(X^+)& =& E^{(i)}(X^+)+_f^{(h)}\hfill \\ \hfill \omega _{if}(XX^0)& =& E^{(i)}(XX^0)E^{(f)}(X^0).\hfill \end{array}$$
(6)
where $`i`$ and $`f`$ label an initial and final state, respectively; $`_f^{(e)}`$ and $`_f^{(h)}`$ are, respectively, the single-particle energy of the electron and hole in the final state; and $`E^{(\nu )}(\chi ^q)`$ is the multi-particle, configuration-interaction energy of state $`|\mathrm{\Psi }^{(\nu )}(\chi ^q)`$. \[In Eq. (6), the energy of the ground state of the system in the absence of excited electron-hole pairs is taken to be zero.\] Optical experiments like photoluminescence probe electron-hole recombination transitions that are allowed (bright).cardona\_book Figure 1(a)-(d) indicates with arrows the lowest bright recombination transitions of exciton $`\chi ^q`$, for a $`75`$Å-tall lens-shaped In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs dot with base diameter $`b=252`$Å. The final states upon electron-hole recombination are presented in the shaded area of Fig. 1. Figure 2 shows the recombination energies $`\omega _{00}(X^{})`$, $`\omega _{00}(X^+)`$ and $`\omega _{03}(XX^0)`$ calculated at the many-particle, configuration-interaction level as a function of dot height. These recombination energies are shown relative to the lowest recombination energy $`\omega _{20}(X^0)`$ of the monoexciton (the latter energy is shown in the top axis of Fig. 2). Thus, the results correspond to spectroscopic shifts that are currently measured by several groups. Prominent features:
(i) The recombination energy $`\omega _{00}(X^{})`$ blue-shifts as height increases; in contrast, $`\omega _{00}(X^+)`$ red-shifts. These trends have been explained by Bester and Zunger in Ref. bester\_PRB\_2003b, by adopting the Hartree-Fock approximation.
(ii) For the flattest dot ($`h=20`$Å) the ordering $`\omega _{00}(X^+)>\omega _{20}(X^0)>\omega _{00}(X^{})`$ of the emission energies and the relative magnitude of the spectroscopic shifts $`\omega _{00}(X^+)\omega _{20}(X^0)<\omega _{20}(X^0)\omega _{00}(X^{})`$ agree with photoluminescence (PL) data. ware\_PhysicaE\_2005 ; skolnick\_ARMR\_2004 ; guffarth\_PhysicaE\_2004 ; finley\_PRB\_2001 These relationships are explained at the Hartree-Fock level of approximation, as in Ref. bester\_PRB\_2003b, , in which one derives
$`\omega _{20}(X^0)`$ $`=`$ $`[_0^{(e)}_0^{(h)}]J_{00}^{(eh)},`$ (7)
$`\omega _{00}(X^+)`$ $`=`$ $`[_0^{(e)}_0^{(h)}]+J_{00}^{(hh)}2J_{00}^{(eh)},`$ (8)
$`\omega _{00}(X^{})`$ $`=`$ $`[_0^{(e)}_0^{(h)}]+J_{00}^{(ee)}2J_{00}^{(eh)}.`$ (9)
Then, the relationship $`J_{00}^{(hh)}>J_{00}^{(eh)}>J_{00}^{(ee)}`$ that holds at $`h=20`$Å \[Fig. 2(b)\] reveals the ordering of the emission energies and the magnitude of the spectroscopic shifts. Regarding the magnitude of the spectroscopic shifts, $`\omega _{00}(X^+)\omega _{20}(X^0)=2.5\mathrm{meV}`$ and $`\omega _{00}(X^{})\omega _{20}(X^0)=3.7\mathrm{meV}`$ agree reasonably well with (a) $`2\mathrm{meV}`$ and $`6\mathrm{meV}`$, respectively, that Ware and coworkers have recently observed in an InAs/GaAs dot (size unspecified) with monoexciton emission at $`1.268\mathrm{eV}`$ (Ref. ware\_PhysicaE\_2005, ); also with (b) $`\omega _{00}(X^{})\omega _{20}(X^0)=5\mathrm{meV}`$ measured by Smith et al. in an (In,Ga)As/GaAs dots (size unspecified) with emission energy at $`1.319\mathrm{eV}`$ (Ref. smith\_PRL\_2005, ); and with (c) the value of $`5.8\mathrm{meV}`$ for $`\omega _{00}(X^{})\omega _{20}(X^0)`$ observed by Finley and coworkers in an (In,Ga)As/GaAs dots \[$`b=(230\pm 70)`$Å, $`h=(25\pm 10)`$Å\] with emission at $`1.263\mathrm{eV}`$ (Ref. finley\_PRB\_2001, ).
(iii) As height increases, $`\omega _{03}(XX^0)`$ red-shifts at small heights, reaches a maximum shift of nearly $`2\mathrm{meV}`$ at $`h50`$Å, and then it moderately blue-shifts for taller dots. In addition, at $`h=50`$Å the emission energy of $`XX^0`$ coincides with that of $`X^+`$. As in (ii), these results are explained at the Hartree-Fock level, which predicts
$`\omega _{03}(XX^0)\omega _{20}(X^0)`$ $`=`$ $`J_{00}^{(ee)}+J_{00}^{(hh)}2J_{00}^{(eh)},`$ (10)
$`\omega _{03}(XX^0)\omega _{00}(X^+)`$ $`=`$ $`J_{00}^{(ee)}J_{00}^{(eh)}.`$ (11)
Here, we have neglected the small ($`1`$-$`6\mu \mathrm{eV}`$) splitting of the monoexciton bright states. Then, by analyzing the height dependence of the direct Coulomb interactions $`J_{00}^{(ee)}`$, $`J_{00}^{(hh)}`$ and $`J_{00}^{(eh)}`$ \[Fig. 2(b)\] we find that Eq. (10) predicts the observed height dependence of $`\omega _{03}(XX^0)`$, although the actual magnitude of the emission is not quantitatively predicted due to correlations.shumway\_PRB\_2001 ; narvaez\_PRB\_tobepublished In addition, Eq. (11) reveals the coincidence of the emission of $`XX^0`$ and that of $`X^+`$ as a result of balancing the magnitudes of $`J_{00}^{(ee)}`$ and $`J_{00}^{(eh)}`$ at $`h=50`$Å. The latter balance arises from the similar degree of localization of $`\psi _0^{(e)}`$ and $`\psi _0^{(h)}`$ within the dot.narvaez\_condmat
(iv) $`\omega _{03}(XX^0)\omega _{20}(X^0)=2.0\mathrm{meV}`$ at $`h=50`$Å agrees with the value of $`2.0\pm 0.1\mathrm{meV}`$ for $`XX^0`$ measured by Finley and coworkers in PL experiments in an (In,Ga)As/GaAs dot \[$`b=(230\pm 70)`$Å and $`h=(25\pm 10)`$Å\] with exciton ground-state emission at $`1.345\mathrm{eV}`$ (Ref. finley\_PRB\_2001, ). In addition, this result of $`2.0\mathrm{meV}`$ for the biexciton shift agrees remarkably well with the value of $`2\mathrm{meV}`$ measured by Rodt et al in InAs/GaAs dots ($`b=100200`$Å, height unspecified) with monoexciton emission energies ranging between $`1.260\mathrm{eV}`$ and $`1.295\mathrm{eV}`$ (Ref. rodt\_papers, ); and also with the measured shift of $`2.3\mathrm{meV}`$ observed by Urbaszek and coworkers in (In,Ga)As/GaAs dots (size unspecified) with monoexciton emission energy of $`1.294\mathrm{eV}`$ (Ref. urbaszek\_PRL\_2003, ). The value of about $`1.7\mathrm{meV}`$ for $`\omega _{03}(XX^0)\omega _{20}(X^0)`$ that we find at $`h=35`$Å and $`65`$Å agrees well with the value of $`1.6\mathrm{meV}`$ observed by Bayer and coworkers in (In,Ga)As/GaAs dots \[$`b=(500\pm 30)`$Å\] with exciton emission at $`1.428\mathrm{eV}`$ (Ref. bayer\_PRB\_1998, ). Our results also agree satisfactorily with the value of $`2.7\mathrm{meV}`$ measured by Findeis and coworkers in dots (size unspecified) with monoexciton emission energy at $`1.284\mathrm{eV}`$ (Ref. findeis\_SSC\_2000, ).
### II.3 Binding energies
The binding energies of excitons $`\chi ^q`$ are defined as
$$\begin{array}{ccc}\hfill \mathrm{\Delta }(X^0)& =& \left[_0^{(e)}_0^{(h)}\right]E^{(0)}(X^0)\hfill \\ \hfill \mathrm{\Delta }(X^{})& =& \left[_0^{(e)}+E^{(0)}(X^0)\right]E^{(0)}(X^{})\hfill \\ \hfill \mathrm{\Delta }(X^+)& =& \left[_0^{(h)}+E^{(0)}(X^0)\right]E^{(0)}(X^+)\hfill \\ \hfill \mathrm{\Delta }(XX^0)& =& 2E^{(0)}(X^0)E^{(0)}(XX^0).\hfill \end{array}$$
(12)
Exciton $`\chi ^q`$ is said to be bound when the binding energy $`\mathrm{\Delta }(\chi ^q)`$ is positive. Conversely, $`\mathrm{\Delta }(\chi ^q)<0`$ implies the exciton is unbound. Note that the binding energy is defined with respect to the ground-state energy of dissociated excitonic complexes. For instance, the binding energy of $`X^0`$ is defined with respect to the energy of a non-interacting electron-hole pair. In turn, the binding energy of the biexciton $`XX^0`$ is defined with respect to the total energy of two non-interacting monoexcitons \[Eq. (12)\].
Figure 3 shows $`\mathrm{\Delta }(X^0)`$, $`\mathrm{\Delta }(X^{})`$, $`\mathrm{\Delta }(X^+)`$ and $`\mathrm{\Delta }(XX^0)`$ as well as $`E^{(0)}(X^0)`$ as a function of dot height. $`\mathrm{\Delta }(X^0)`$ decreases with increasing height, and it is well approximated by $`J_{00}^{eh}`$ as correlation effects are relatively small ($`2\mathrm{meV}`$). The height dependence of the binding energy of $`X^{}`$, $`X^+`$, and $`XX^0`$ follows the height dependence of the spectroscopic shifts shown in Fig. 2(a). This is so because the bright-dark splitting for the monoexciton $`X^0`$ is small $`80`$-$`180\mu \mathrm{eV}`$. As expected from the spectroscopic shifts results \[Fig. 2(a)\], the height dependence of the binding energy is qualitatively different for each exciton $`\chi ^q`$. (i) $`\mathrm{\Delta }(X^{})`$ is bound ($`4\mathrm{meV}`$) for the flattest dot ($`h=20`$Å) and it decreases almost linearly, becoming unbound for dots taller than $`h=65`$Å. (ii) In contrast, $`\mathrm{\Delta }(X^+)`$ is unbound for the flattest dot and increases up to $`3\mathrm{meV}`$, becoming bound slightly below $`h=35`$Å. (iii) $`\mathrm{\Delta }(XX^0)`$ does not depend monotonically on the gap, reaching a maximum around $`h=50`$Å. In addition, $`XX^0`$ is bound \[$`\mathrm{\Delta }(XX^0)>0`$\] for all heights above $`20`$Å while unbound for the flattest dot ($`h=20`$Å). The latter is due to an interplay between Hartree-Fock and correlation contributions to binding, $`\mathrm{\Delta }_{\mathrm{HF}}(XX^0)`$ and $`\delta (XX^0)`$, respectively, that results in correlation being insufficient to bind $`XX^0`$. Namely,
$`\mathrm{\Delta }(XX^0)`$ $`=`$ $`\mathrm{\Delta }_{\mathrm{HF}}(XX^0)+\delta (XX^0)`$
$`=`$ $`\{2J_{00}^{(eh)}[J_{00}^{(ee)}+J_{00}^{(hh)}]\}+\delta (XX^0)`$
$`=`$ $`1.6\mathrm{meV}+1.4\mathrm{meV}`$
$`=`$ $`0.2\mathrm{meV},`$
It should be noted that Rodt and coworkersrodt\_papers observed (in photoluminescence) a bound-unbound crossover for $`XX^0`$ as the monoexciton emission energy of InGaAs/GaAs quantum dots decreased. Those authors also calculated the binding energy of $`XX^0`$, using an 8-band $`𝐤𝐩`$ model, and suggested that the reduction of correlation effects was responsible for unbinding $`XX^0`$ as the gap increased.rodt\_papers
## III Polarization anisotropy of optical transitions
When an electron-hole pair in exciton $`\chi ^q`$ recombines optically, the transition is characterized by both the transition energy $`\omega _{if}(\chi ^q)`$ \[Eq. (6)\] and the transition dipole matrix element
$$M_{if}^{(\widehat{𝐞})}(\chi ^q)=\mathrm{\Psi }^{(f)}(\chi ^q1)|\widehat{𝐞}𝐩|\mathrm{\Psi }^{(i)}(\chi ^q).$$
(14)
Here, $`𝐩`$ is the electron momentum and $`\widehat{𝐞}`$ is the polarization vector of the electromagnetic field.cardona\_book
The dipole matrix elements $`𝐌_{if}^{(\widehat{𝐞})}(\chi ^q)`$ \[Eq. (14)\] depend on the polarization vector $`\widehat{𝐞}`$; so, it is natural to quantify what is the degree of polarization anisotropy between different polarizations $`\widehat{𝐞_\mathrm{𝟏}}`$ and $`\widehat{𝐞_\mathrm{𝟐}}`$. Therefore, we introduce the recombination (emission or pholuminescence) intensity spectrum of exciton $`\chi ^q`$ for polarization $`\widehat{𝐞}`$,
$$I^{(\widehat{𝐞})}(\omega ,T;\chi ^q)=\underset{i,f}{}\left|M_{if}^{(\widehat{𝐞})}(\chi ^q)\right|^2P_i(T;\chi ^q)\delta [\omega \omega _{if}(\chi ^q)].$$
(15)
Here,
$$P_i(T;\chi ^q)=𝒩\mathrm{exp}\{[E^{(i)}(\chi ^q)E^{(0)}(\chi ^q)]/k_BT\}$$
(16)
is the occupation (Boltzmann) probability of the initial state $`|\mathrm{\Psi }^{(i)}(\chi ^q)`$ at temperature $`T`$; $`𝒩`$ is a normalization constant such that $`_iP_i(T;\chi ^q)=1`$ and $`k_B`$ is the Boltzmann constant. Then, as in-plane polarizations $`\widehat{𝐞}_1=[110]`$ and $`\widehat{𝐞}_2=[1\overline{1}0]`$ have been probed extensively, we introduce the in-plane polarization anisotropy parameter $`\lambda `$,
$$\lambda (\omega ,T;\chi ^q)=\frac{I^{([110])}(\omega ,T;\chi ^q)I^{([1\overline{1}0])}(\omega ,T;\chi ^q)}{I^{([110])}(\omega ,T;\chi ^q)+I^{([1\overline{1}0])}(\omega ,T;\chi ^q)}.$$
(17)
Thus, $`\lambda =1`$ indicates an optical transition that is fully polarized along the $`[110]`$ direction, while $`\lambda =1`$ indicates one fully polarized along $`[1\overline{1}0]`$.
### III.1 In-plane polarization anisotropy of the lowest optical transitions of $`X^0`$, $`X^{}`$, $`X^+`$, and $`XX^0`$
Figure 4 shows the in-plane polarization anisotropy $`\lambda `$ for the lowest optical transitions of (a) $`X^0`$, (b) $`XX^0`$, (c) $`X^{}`$, and (d) $`X^+`$ as a function of dot height. Two features are prominent:
(i) The bright transitions $`\omega _{20}(X^0)`$ and $`\omega _{30}(X^0)`$ are polarized and the polarization anisotropy depends on height; see Fig. 4(a). Similarly, the lowest transitions of the biexciton $`\omega _{03}(XX^0)`$ and $`\omega _{02}(XX^0)`$ are polarized and the degree of polarization also depends on height. These transitions correspond to the decay into the two bright states of the monoexciton \[Fig. 1(d)\]. Clearly, the biexciton transitions inherit the polarization of the monoexciton transitions. For both $`X^0`$ and $`XX^0`$, we see that $`\lambda `$ switches sign as a function of height. In particular, at $`h=50`$Å the transitions of $`X^0`$ present no in-plane polarization anisotropy \[see arrow in Figs. 4(a)\] because the bright doublet formed by $`|\mathrm{\Psi }^{(2)}(X^0)`$ and $`|\mathrm{\Psi }^{(3)}(X^0)`$ is degenerate, which according to our calculation results in $`I^{([110])}(\omega _{20})I^{([1\overline{1}0])}(\omega _{20})`$.
(ii) The lowest optical transitions of $`X^{}`$ and $`X^+`$ have degenerate transition (recombination) energies \[Fig. 1(b) and (c)\] and, according to our calculations, this results in $`I^{([110])}(\omega _{00})I^{([1\overline{1}0])}(\omega _{00})`$ and, thus, negligible in-plane polarization anisotropy regardless of height.
In addition to (i) and (ii), we find that (iii) the in-plane polarization anisotropy of the lowest transitions of $`X^0`$ and $`XX^0`$ depends dramatically on the dot’s alloy randomness (disorder realization); as shown in Figure 5 for seven alloy realizations in a $`35`$Å-tall dot. For a given transition in both $`X^0`$ and $`XX^0`$, $`\lambda `$ changes sign depending on alloy randomness. Further, while some alloy realizations like $`4`$ and $`5`$ result in transitions nearly fully polarized ($`\lambda 99\%`$), others such as $`2`$, $`6`$ and $`7`$ present small anisotropy ($`\lambda 20\%`$).
## IV Radiative recombination lifetimes
The characteristic radiative lifetime $`\tau _{if}(\chi ^q)`$ of a transition $`|\mathrm{\Psi }^{(i)}(\chi ^q)|\mathrm{\Psi }^{(f)}(\chi ^q1)`$ follows from both the magnitude of the dipole matrix element of the transition $`\left|𝐌_{if}^{(\widehat{𝐞})}(\chi ^q)\right|^2`$ and the recombination energy $`\omega _{if}`$.dexter\_book Namely,
$$\frac{1}{\tau _{if}(\chi ^q)}=\frac{4}{3}\left(\frac{e^2}{m_0^2c^3\mathrm{}^2}\right)n\omega _{if}(\chi ^q)\underset{\widehat{𝐞}=\widehat{x},\widehat{y},\widehat{z}}{}\left|𝐌_{if}^{(\widehat{𝐞})}(\chi ^q)\right|^2.$$
(18)
Here, $`e`$ and $`m_0`$ are the charge and mass of the electron, respectively, and $`c`$ is the velocity of light in vacuum. In addition, the refractive index $`n`$ of the dot material accounts for the material’s effects on the photon emission. The linear dependence of $`1/\tau _{if}(\chi ^q)`$ on refractive index is applicable only when considering dot and matrix materials with similar dielectric constants, as it is the case in InGaAs/GaAs dots. In a more general case, more complicate dependences have been proposed.thranhardt\_PRB\_2002 Note that the characteristic radiative lifetime \[Eq. (18)\] does not depend on temperature nor on the occupation probability of the initial state $`|\mathrm{\Psi }^{(i)}(\chi ^q)`$, as it is a characteristic property of the transition $`|\mathrm{\Psi }^{(i)}(\chi ^q)|\mathrm{\Psi }^{(f)}(\chi ^q1)`$. On the other hand, the actual radiative lifetime $`\tau (\chi ^q)`$ of exciton $`\chi ^q`$ depends both on the probability $`n_i`$ of having the initial states $`|\mathrm{\Psi }^{(i)}(\chi ^q)`$ of exciton $`\chi ^q`$ occupied and the number of final states $`|\mathrm{\Psi }^{(f)}(\chi ^q1)`$ of exciton $`\chi ^q1`$ available for recombination, as well as on the characteristic radiative lifetimes $`\tau _{if}(\chi ^q)`$. We calculate $`\tau (\chi ^q)`$ from
$$\frac{1}{\tau (\chi ^q)}=\underset{f}{}\underset{i}{}n_i\frac{1}{\tau _{if}(\chi ^q)}.$$
(19)
Here, $`_in_i=1`$ and $`n_i=n_i(T;\chi ^q)`$ where $`T`$ is the temperature of the system. In general, the calculation of $`n_i`$ involves solving a system of rate equations.dekel\_PRB\_2000 However, in the case that the intra-level relaxation of exciton $`\chi ^q`$ is much faster than the radiative lifetimes, $`n_i`$ is given by the Boltzmann weight $`P_i(T;\chi ^q)`$ \[Eq. (16)\].
### IV.1 Radiative lifetime $`\tau (\chi ^q)`$ of the lowest optical transitions of $`X^0`$, $`X^{}`$, $`X^+`$, and $`XX^0`$
Monoexciton. At $`T=0\mathrm{K}`$, and assuming fast non-radiative relaxation to the dark ground-state of $`X^0`$, the radiative lifetime of the monoexciton $`\tau (X^0)`$ equals the characteristic radiative lifetime $`\tau _{00}(X^0)`$ and it is long ($`\mathrm{ms}`$). At finite temperatures, all four levels \[Fig. 1(a)\] of the $`X^0`$ ground state are thermally populated and emit light with their own characteristic lifetime \[Eq. (18)\]. For temperatures such that the lowest four monoexciton states are occupied, while the occupation of excited states is negligible, we calculate the radiative lifetime from Eq. (19) and find
$$\tau (X^0)=4\left[\frac{\tau _{20}(X^0)\tau _{30}(X^0)}{\tau _{20}(X^0)+\tau _{30}(X^0)}\right]2\tau _{20}(X^0).$$
(20)
Here, we have used a result of our calculations that predict $`\tau _{20}(X^0)\tau _{30}(X^0)`$ regardless of dot height. In addition, we have neglected the long-lived recombinations $`\omega _{00}`$ and $`\omega _{10}`$. Figure 6(a) shows the characteristic $`\tau _{20}(X^0)`$ versus dot height. We find that this lifetime depends weakly on height. A calculation of $`\tau (X^0)`$ at $`T=10\mathrm{K}`$ predicts a moderate decrease in the monoexciton radiative lifetime as height increases \[Fig. 6(b)\]. It should be noted that $`\tau (X^0)`$ is actually bigger than the approximate value of $`2\tau _{20}(X^0)`$ \[Eq. (20)\] due to the actual occupation probability $`n_i(10\mathrm{K},X^0)`$ of the initial states. By measuring time-resolved photoluminescence in InAs/GaAs dots ($`b=120`$-$`220`$Å, height unspecified) with monoexciton emission ranging from $`1.2`$-$`1.3\mathrm{eV}`$, Karachinsky and coworkers have recently observed the opposite trend,karachinsky\_APL\_2004 i. e. $`\tau (X^0)`$ increases with dot size from $`1\mathrm{ns}`$ in a dot with emission energy of $`1.31\mathrm{eV}`$ to $`3\mathrm{ns}`$ in one with emission energy of $`1.25\mathrm{eV}`$. The values of $`\tau (X^0)`$ we predict for tall dots ($`h=65`$Å and $`75`$Å) agree satisfactorily with the value of $`1\mathrm{ns}`$ measured by Buckle et al for an InAs/GaAs dot ($`b120`$Å, $`h\mathrm{\hspace{0.33em}30}`$Å) with gap $`1.131\mathrm{eV}`$ (Ref. buckle\_JAP\_1999, ) at a temperature of $`6\mathrm{K}`$. Further, our predictions also agree well with the value of $`1.55\mathrm{ns}`$ extracted from time-resolved photoluminescence experiments at $`10\mathrm{K}`$ performed in InAs/GaAs dots ($`b=200\AA `$, $`h=20`$Å) by Bardot and coworkers.bardot\_condmat\_2005
Biexciton. In contrast to the monoexciton, the biexciton ground state is singly-degenerate and bright \[Fig. 1(b)\]. Thus, at $`T=0\mathrm{K}`$, the radiative lifetime is given from Eq. (19) as
$$\tau (XX^0)=\left[\frac{\tau _{02}(XX^0)\tau _{03}(XX^0)}{\tau _{02}(XX^0)+\tau _{03}(XX^0)}\right]\frac{1}{2}\tau _{03}(XX^0).$$
(21)
Similarly to the $`X^0`$ case, in Eq. (21) we have used the relationship $`\tau _{03}(XX^0)\tau _{02}(XX^0)`$ that our calculations predict, and we have neglected the long-lived dark recombination channels. Figure 6(a) shows $`\tau _{03}(XX^0)`$ as a function of height. We find a weak dependence with height, as in the monoexciton case. Remarkably, we find that $`\tau _{20}(X^0)\tau _{03}(XX^0)`$ regardless of height. The latter leads to the following relationship between the radiative lifetime of the monoexciton and biexciton,
$$\tau (XX^0)\frac{1}{4}\tau (X^0).$$
(22)
Figure 6(b) shows that at $`T=10\mathrm{K}`$ the calculated $`\tau (XX^0)`$ depends weakly on height, changing by about $`0.1\mathrm{ns}`$ in the entire range of heights studied. $`\tau (XX^0)`$ is in excellent agreement with the value of $`0.5\mathrm{ns}`$ measured by Ulrich et al in an (In,Ga)As/GaAs quantum dot ($`b=150`$-$`200`$Å, $`h=10`$-$`20`$Å) with an exciton gap of $`1.337\mathrm{eV}`$ (Ref. ulrich\_condmat\_2004, ).
Trions. Both $`X^{}`$ and $`X^+`$ have a two-fold degenerate ground state that is bright \[Figs. 1(c) and (d)\]. In each of these excitons, there are four lowest transitions; namely, $`\omega _{00}(X^{})`$ \[Fig. 1(c)\], $`\omega _{10}(X^{})`$, $`\omega _{01}(X^{})`$, and $`\omega _{11}(X^{})`$; and $`\omega _{00}(X^+)`$ \[Fig. 1(d)\], $`\omega _{10}(X^+)`$, $`\omega _{01}(X^+)`$, and $`\omega _{11}(X^+)`$, for $`X^{}`$ and $`X^+`$, respectively. In turn, each of these transitions have a corresponding characteristic radiative lifetime. For the latter, our calculations predict $`\tau _{00}(X^{})=\tau _{11}(X^{})`$ and $`\tau _{01}(X^{})=\tau _{10}(X^{})`$ as well as $`\tau _{00}(X^+)=\tau _{11}(X^+)`$ and $`\tau _{01}(X^+)=\tau _{10}(X^+)`$. Thus, at $`T=0\mathrm{K}`$, the radiative lifetimes of $`X^{}`$ and $`X^+`$ are given by
$`\tau (X^{})=\left\{{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{1}{\tau _{00}(X^{})}}+{\displaystyle \frac{1}{\tau _{01}(X^{})}}\right]+{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{1}{\tau _{10}(X^{})}}+{\displaystyle \frac{1}{\tau _{11}(X^{})}}\right]\right\}^1={\displaystyle \frac{\tau _{00}(X^{})\tau _{01}(X^{})}{\tau _{00}(X^{})+\tau _{01}(X^{})}},`$
(23)
$`\tau (X^+)=\left\{{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{1}{\tau _{00}(X^+)}}+{\displaystyle \frac{1}{\tau _{01}(X^+)}}\right]+{\displaystyle \frac{1}{2}}\left[{\displaystyle \frac{1}{\tau _{10}(X^+)}}+{\displaystyle \frac{1}{\tau _{11}(X^+)}}\right]\right\}^1={\displaystyle \frac{\tau _{00}(X^+)\tau _{01}(X^+)}{\tau _{00}(X^+)+\tau _{01}(X^+)}}.`$
Figure 6(a) shows the height dependence of $`\tau _{00}(X^{})`$ and $`\tau _{01}(X^{})`$, and $`\tau _{00}(X^+)`$ and $`\tau _{01}(X^+)`$. These characteristic radiative lifetimes depend strongly and non-monotonically on height. This non-monotonic dependence translates into a rather simple and monotonic $`\tau (X^{})`$ and $`\tau (X^+)`$, as shown in Fig. 6(b). For flat dots $`\tau (X^{})\tau (X^+)`$, whereas for taller dots these lifetimes become slightly different. Our predicted $`\tau (X^{})`$ are in satisfactory agreement with the value of $`0.6\mathrm{ns}`$ recently observed by Smith and co-workers in (In,Ga)As/GaAs dots (size unspecified) with exciton ground-state emission at $`1.318\mathrm{eV}`$ (Ref. smith\_PRL\_2005, ). We find that the radiative lifetimes of the charged trions satisfy the relationship $`\tau (X^0)>\tau (X^{})\tau (X^+)>\tau (XX^0)`$.
## V Comparison of $`X^0`$, $`X^{}`$, $`X^+`$, and $`XX^0`$ in lens-shaped pure InAs/GaAs with alloyed (In,Ga)As/GaAs dots
For completeness, we briefly compare the binding and recombination energies, polarization anisotropy, and radiative lifetimes in lens-shaped pure, non-alloyed InAs/GaAs dots with In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs dots. Note that Williamson, Wang, and Zunger have already compared results for $`X^0`$ for several alloy profiles.williamson\_PRB\_2000 In Ref. narvaez\_condmat, we predicted that hole localization takes place at the dot-GaAs matrix interface as the height of these dots increases above $`35`$Å. Thus, we discuss here two flat dots ($`h=20`$Å and $`35`$Å) with base $`b=252`$Å.
(i) Recombination energies are smaller in InAs/GaAs dots than in In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs dots with the same geometry. For instance, $`\omega _{20}(X^0)=1.078\mathrm{eV}`$ and $`0.987\mathrm{eV}`$ for $`h=20`$Å and $`35`$Å, respectively.
(ii) The spectroscopic shifts show the same trends with height in flat InAs/GaAs as those in flat In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs dots. So do the binding energies $`\mathrm{\Delta }(X^{})`$ and $`\mathrm{\Delta }(X^+)`$, and $`\mathrm{\Delta }(XX^0)`$. However, there are two important differences bewteen the pure, non-alloyed InAs/GaAs dots and their In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs counterparts: (a) For the $`20`$Å-tall InAs/GaAs dot, while the binding energies still satisfy $`\mathrm{\Delta }(X^{})>\mathrm{\Delta }(XX^0)>\mathrm{\Delta }(X^+)`$, we find that $`X^+`$ and $`XX^0`$ are bound, with $`\mathrm{\Delta }(X^+)=1.6\mathrm{meV}`$ and $`\mathrm{\Delta }(XX^0)=1.5\mathrm{meV}`$, respectively. This is so because at this height the InAs/GaAs dot is in the nearly “symmetric” regime: $`J_{00}^{(hh)}=25.6\mathrm{meV}J_{00}^{(ee)}=25.1\mathrm{meV}J_{00}^{(eh)}=25.3\mathrm{meV}`$, so the Hartree-Fock component of the binding energy \[see Eq. (II.3) for $`XX^0`$ case\] is much smaller than in the In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs dot with same height. Thus, correlation becomes capable of binding $`X^+`$ and $`XX^0`$ in the pure, non-alloyed dot. (b) For the $`35`$Å-tall InAs/GaAs dot, we find ordering reversal, i.e. $`\mathrm{\Delta }(X^+)>\mathrm{\Delta }(XX^0)>\mathrm{\Delta }(X^{})`$. In In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs dots, this ordering is attained at $`h=50`$Å (Fig. 3).
(iii) In contrast to the findings in In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs dots, transitions $`\omega _{20}(X^0)`$ and $`\omega _{30}(X^0)`$ are fully polarized along $`[110]`$ ($`\lambda =1`$) and $`[1\overline{1}0]`$ ($`\lambda =1`$), respectively, regardless of height. Consequently, $`\omega _{03}(XX^0)`$ and $`\omega _{02}(XX^0)`$ are fully polarized along $`[1\overline{1}0]`$ and $`[110]`$, respectively. These polarizations are expected from a dot with $`C_{2v}`$ symmetry, like a lens-shaped pure InAs/GaAs dot.
(iv) Radiative lifetimes $`\tau (X^0)`$, $`\tau (X^+)`$ and $`\tau (X^{})`$, and $`\tau (XX^0)`$ are similar to those in In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs dots with the same geometry. For instance, $`\tau (X^0)=2.8\mathrm{ns}`$ and $`\tau (XX^0)=0.7\mathrm{ns}`$ for $`h=35`$Å; and $`\tau (X^0)=2.9\mathrm{ns}`$ and $`\tau (XX^0)=0.6\mathrm{ns}`$ for $`h=20`$Å.
## VI Summary
We have addressed the height dependence of recombination energies, polarization and radiative lifetimes of the lowest optical transitions of the neutral exciton ($`X^0`$), negatively- ($`X^{}`$) and positively-charged ($`X^+`$) trions, and the biexciton ($`XX^0`$) in lens-shaped, self assembled In<sub>0.6</sub>Ga<sub>0.4</sub>As/GaAs quantum dots. We have predicted the following.
(i) The recombination energy of the lowest transition of $`X^{}`$, $`X^+`$ and $`XX^0`$, correspondingly, $`\omega _{00}(X^{})`$, $`\omega _{00}(X^+)`$ and $`\omega _{03}(XX^0)`$ shows qualitatively different behavior for each excitonic complex. Namely, $`\omega _{00}(X^{})`$ blue-shifts as height increases, whereas that of $`\omega _{00}(X^+)`$ red-shifts. On the other hand, as height increases, $`\omega _{03}(XX^0)`$ shows a red-shift at small heights, reaches a maximum shift, and then blue-shifts for taller dots. This behavior is explained by the height dependence and relative magnitude of $`J_{00}^{(ee)}`$, $`J_{00}^{(hh)}`$ and $`J_{00}^{(eh)}`$.
(ii) The binding energies $`\mathrm{\Delta }(X^{})`$, $`\mathrm{\Delta }(X^+)`$ and $`\mathrm{\Delta }(XX^0)`$ follow the height dependence of the emission spectroscopic shifts. Changes in the dot height drives a bound-to-unbound crossover for each of these complexes.
(iii) The in-plane polarization anisotropy $`\lambda `$ of the lowest transitions of $`X^0`$ ($`\omega _{20}`$) and $`XX^0`$ ($`\omega _{03}`$) strongly depends on dot height as well as on alloy randomness (disorder realization). In contrast, the lowest transitions of $`X^{}`$ and $`X^+`$ present negligible $`\lambda `$ regardless of height.
(iv) The ground state of $`X^0`$ encompasses four states that split off in a low-energy pair that is dark and a high-energy pair that is bright, with a bright-dark splitting that increases as height decreases. Thus, at $`T=0\mathrm{K}`$ the radiative lifetime $`\tau (X^0)`$ of $`X^0`$ is long. On the other hand, at $`T=10\mathrm{K}`$ both dark and bright states are populated; so, $`\tau (X^0)`$ becomes fast, moderately decreases as height increases, and its magnitude ranges from $`2`$-$`3\mathrm{ns}`$. In contrast, $`\tau (X^{})`$, $`\tau (X^+)`$ and $`\tau (XX^0)`$ are fast even at $`T=0\mathrm{K}`$, as a consequence of these excitons having ground states that are bright. These radiative lifetimes depend weakly on height. Further, $`\tau (X^{})\tau (X^+)1.1\mathrm{ns}`$, while $`\tau (XX^0)0.5\mathrm{ns}`$.
We have compared these predictions with available data and have found them in satisfactory agreement. In addition, we compared with results in pure, non-alloyed InGa/GaAs quantum dots.
###### Acknowledgements.
We thank Alberto Franceschetti (NREL) for valuable discussions. This work has been supported by U.S. DOE-SC-BES-DMS under contract No. DE-AC36-99GO10337.
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# Oriented Percolation in One–Dimensional 1/|𝑥-𝑦|² Percolation Models
## 1 Introduction
It is well known that $`1/r^2`$ gives the “critical” falloff for percolation in one-dimensional long range independent edge percolation models. Moreover, for the one dimensional Fortuin–Kasteleyn (FK) random cluster model with weighting factor $`\kappa 1`$ and edge occupation probabilities of the form $`p_{\{x,y\}}=f(\left|xy\right|)`$, with $`r^2f(r)\beta >0`$ as $`r+\mathrm{}`$, it is known that for fixed $`f(j)<1,j2`$ and varying $`p=f(1)`$, the value $`\beta ^{}=1`$ is critical in the sense that for $`\beta 1`$ percolation cannot occur unless $`p=1`$ (see \[AN\]), while for $`\beta >1`$ there is percolation provided $`p`$ is sufficiently close to one (see \[IN\] and \[M\]). Such results are important in the description of the phase transition diagram for the one-dimensional long range Ising models studied earlier by Fröhlich and Spencer in \[FS1\] and for the corresponding Potts models (\[ACCN\],\[IN\],\[M\]), as these spin systems can be constructed by a random coloring of the clusters in the FK model with $`\kappa =2`$, or $`\kappa >2`$ integer, respectively. For the particular case of independent edge percolation models ($`\kappa =1`$) earlier results were obtained in \[NS\], where it was proven that $`\beta ^{}1`$ in this case, and that oriented percolation occurs when $`\underset{x\mathrm{}}{lim}x^sf(x)>0`$ for some $`1<s<2`$ and $`p`$ is sufficiently close to 1. The question whether oriented percolation occurs in the boundary case $`s=2`$ remained unanswered. Theorem 1.1 below gives an affirmative answer; the result is stated for the particular example of edge probabilities in (1.1) below, and oriented percolation is shown when $`\beta >1`$ and $`p<1`$ is sufficiently close to one. The proofs can be easily adapted to include any $`f()`$ satisfying $`\underset{x\mathrm{}}{lim}x^2f(x)>1`$. In this sense $`\beta ^{}=1`$ remains critical also for oriented percolation.
Our main result (Theorem 1.1) deals with the independent percolation model. On the other hand, known FKG inequalities and the above mentioned representation yield at once an application to the long range Ising (and Potts) models, which we state as Corollary 1.2. This is the reason for the preliminaries on this more general context of FK measures.
Preliminaries. Consider the infinite complete graph with set of vertices $`𝒱=`$ and set of edges $`𝔼=\{\{x,y\},xy,x,y\}`$, and let $`\mathrm{\Omega }=\{0,1\}^𝔼`$. One-dimensional long-range FK random cluster models with weighting parameter $`\kappa 1`$ are probability measures on $`\sigma (\mathrm{\Omega })`$, the usual product $`\sigma `$–algebra on $`\mathrm{\Omega }`$. To define them, let us first fix $`\nu `$ the Bernoulli product measure on $`\mathrm{\Omega }`$, with $`\nu (\omega _{\{x,y\}}=1)=p_{\{x,y\}}`$ given by
$$p_{\{x,y\}}=\{\begin{array}{cc}p\hfill & \mathrm{if}\left|xy\right|=1,\hfill \\ 1\mathrm{exp}\left\{\frac{\beta }{\left|xy\right|^2}\right\}\hfill & \mathrm{otherwise},\hfill \end{array}$$
(1.1)
where $`0<p<1`$ and $`\beta >0`$ are fixed parameters.
Notation. We write $`q_{\{x,y\}}=1p_{\{x,y\}}`$; for $`e=\{x,y\}`$ we will write $`p_e`$ instead of $`p_{\{x,y\}}`$, and say that $`e`$ “is open” if $`\omega _e=1`$. The length of an edge $`e=\{x,y\}`$ is $`|xy|`$.
Finite volume FK measures. Given $`I`$, consider $`𝔼(I)=\{\{x,y\}𝔼:x,yI\}`$, $`\mathrm{\Omega }_I=\{0,1\}^{𝔼(I)}`$ and $`\overline{\mathrm{\Omega }}_I=\{0,1\}^{𝔼𝔼(I^c)}`$, where $`I^c=\backslash I`$. Assume that $`|I|<\mathrm{}`$. The corresponding finite volume free FK-measure is the probability measure $`\mu _{\kappa ,I}^f`$ on $`\mathrm{\Omega }_I`$ given by
$$\mu _{\kappa ,I}^f(A)=\frac{{\displaystyle _A}\kappa ^{𝒞_I(\omega )}\nu _I\left(d\omega \right)}{{\displaystyle _{\mathrm{\Omega }_I}}\kappa ^{𝒞_I(\omega )}\nu _I\left(d\omega \right)},A\mathrm{\Omega }_I,$$
(1.2)
where $`\nu _I`$ is the restriction of $`\nu `$ to $`\mathrm{\Omega }_I`$, and $`𝒞_I(\omega )`$ denotes the number of disjoint connected components in the graph determined by $`\omega \mathrm{\Omega }_I`$ (i.e. the graph with vertices in $`I`$ whose edges coincide with those $`e`$ such that $`\omega _e=1`$). The corresponding wired FK-measure $`\mu _{\kappa ,I}^w`$ is a probability measure on $`\overline{\mathrm{\Omega }}_I`$, defined similarly as in (1.2), replacing $`\nu _I`$ by $`\overline{\nu }_I`$, the restriction of $`\nu `$ to $`\overline{\mathrm{\Omega }}_I`$ (so that $`A\sigma (\overline{\mathrm{\Omega }}_I`$) the usual product sigma algebra), and $`𝒞_I(\omega )`$ by $`\overline{𝒞}_I(\omega )`$, the number of disjoint connected components intersecting $`I`$ in the graph with vertices in $``$ determined by $`\overline{\omega }`$, the configuration which extends $`\omega \overline{\mathrm{\Omega }}_I`$ by setting $`\overline{\omega }_e=1`$ for all $`e𝔼(I^c)`$. Thus we may see $`\mu _{\kappa ,I}^w`$ as a measure on $`\mathrm{\Omega }`$, concentrated on the configurations for which all edges in $`𝔼(I^c)`$ are open. Analogously, we may think of $`\mu _{\kappa ,I}^f`$ as a probability measure on $`\mathrm{\Omega }`$, concentrated on the configurations $`\omega `$ such that $`\omega _e=0`$ for any $`e𝔼𝔼(I)`$. Keeping this in mind we have the following well-known property.
The infinite volume limit. On $`\mathrm{\Omega }`$ we consider the usual partial order: $`\omega \omega ^{}`$ if $`\omega _e\omega _e^{}`$ for each $`e𝔼`$. By the FKG inequality (see \[F\], \[ACCN\]), one has if $`\kappa 1`$
$$\mu _{\kappa ,I}^f(g)\mu _{\kappa ,I^{}}^f(g)\mu _{\kappa ,I^{}}^w(g)\mu _{\kappa ,I}^w(g)$$
for any finite intervals $`II^{}`$, and any non-decreasing continuous function $`g:\mathrm{\Omega }`$. Thus, as $`I`$ the limit measures $`\mu _\kappa ^f`$ and $`\mu _\kappa ^w`$ exist. Moreover, $`\mu _\kappa ^f\mu _\kappa ^w`$ in FKG sense.<sup>1</sup><sup>1</sup>1That is $`\mu \mu ^{}`$ if $`\mu (g)\mu ^{}(g)`$ for any $`g`$ continuous and increasing. If $`\kappa =1`$, trivially $`\mu _\kappa ^f=\mu _\kappa ^w=\nu `$. Since $`p_{\{x,y\}}f(|xy|)`$ the measures $`\mu _\kappa ^f`$ and $`\mu _\kappa ^w`$ are translation invariant; both are ergodic.
For a more general and complete discussion on the construction of random cluster measures, including issues in the infinite volume limit for general external conditions, see e.g. \[G1\] and \[G2\] (focused mostly in short range models). This is particularly delicate when $`0<\kappa <1`$.
Fix $`\omega \mathrm{\Omega }`$. An alternating sequence of vertices and edges $`x=x_1,e_1,x_2,\mathrm{},x_{n1},e_{n1},x_n=y,n1,`$ is called a path connecting $`x`$ to $`y`$, and we say that the path is open if $`\omega _{e_i}\omega _{\{x_i,x_{i+1}\}}=1`$, $`1in1`$. We say that $`C`$ is connected if for any two distinct vertices $`x,y`$ in $`C`$ there exists an open path $`\pi `$ connecting them. A maximal connected set is called an open cluster, and $`C_x(\omega )`$ denotes the open cluster containing $`x`$ (we write $`C_x(\omega )=\{x\}`$ if $`\omega _{\{x,y\}}=0`$, for all $`y\{x\}`$). A path $`\pi =(x_1,\mathrm{},x_n)`$ connecting $`x`$ to $`y,x<y`$, is called oriented if $`x_1=x<x_2<\mathrm{}<x_{n1}<x_n=y`$, and we write $`xy`$ when there is an open oriented path connecting $`x`$ to $`y`$. Analogously we define $`C_x^+=\{y:xy\}`$, and the event
$$[x\mathrm{}]=[|C_x^+|=\mathrm{}].$$
We are ready to state our main result.
###### Theorem 1.1
For any $`\beta >1`$, there exist $`0<p_0<1`$ such that, if $`p>p_0`$, then
$$\nu (0\mathrm{})1ϵ$$
(1.3)
holds with $`ϵ=ϵ(p)0`$ as $`p1`$.
Remarks. 1. Let $`\kappa >1`$. The statements in Theorem 4.1 of \[ACCN\] imply that $`\nu \mu _\kappa ^f`$ in FKG sense, provided the probabilities $`p_{\{x,y\}}`$ in $`\mu _\kappa ^f`$ are given by (1.1) with $`\beta `$ replaced by $`\beta ^{}\kappa \beta `$ and $`p`$, writing $`p=1e^\beta `$ for $`|xy|=1`$, replaced by $`p^{}=1(1p)^{\beta ^{}/\beta }`$. Since $`\mu _\kappa ^f\mu _\kappa ^w`$, the same holds as well for $`\mu _\kappa ^w`$.
2. Theorem 1.1 should indeed extend exactly to the FK random cluster model with $`\kappa >1`$. The authors believe that using an algebraic implementation of the multiscale analysis developed in the present work, one should be able to obtain this extension. Nevertheless, for the moment we do not have a full proof (\[MSV\]).
Some related problems. The type of questions treated here has various sources of interest and we mention only a couple of them, which have to do with our own motivations. Consider the following physical problem: take the one-dimensional Ising model with pair interactions, the couplings decaying as the inverse of square of the distance between vertices, at inverse temperature $`\beta >1`$; this is the model studied by Fröhlich and Spencer (\[FS1\]), for which a phase transition was established. Take now the finite box $`[L,L]`$ and assume the Dobrushin boundary conditions, i.e. all spins in $`(\mathrm{},L]`$ will be taken as $`+1`$, and all spins in $`[L,+\mathrm{})`$ will be taken as $`1`$. What can we say about the behaviour of this model when $`L\mathrm{}`$? Is there any sort of well defined interface? This might require a direct analysis in terms of the spin system, but it leads to a more general question for the FK model, regarding the behavior of connected components of each boundary conditioned not to touch each other. (Recall that by a random coloring of the clusters, the FK model gives origin to a spin system which interpolates the independent percolation model ($`\kappa =1`$), the Ising model ($`\kappa =2`$) and the $`q`$–states Potts model ($`\kappa =q>2`$, integer) at inverse temperature $`\beta `$ and interaction $`J_{\{x,y\}}=\beta ^1\mathrm{log}(\frac{1}{1p_{\{x,y\}}})`$, the representation being possible for some (but not all) boundary conditions. For details see \[FK, ACCN\]). Though we still do not fully understand this problem which remains unsolved, our results might shed some light on it. In \[CMPR\], the authors obtain a more precise description for very low temperatures, using cluster expansion techniques.
An interesting corollary of Theorem 1.1 is as follows. Consider the Ising model (with $`\pm 1`$–valued spins) on $`_+`$, with interaction $`J_{\{x,y\}}=|xy|^2`$ if $`\left|xy\right|2`$ and $`J_{\{x,x+1\}}=J`$ at inverse temperature $`\beta `$. Let $`m_L^{0,+}(\beta )`$ denote the average spin at the origin, with “one-sided” $`(+)`$ boundary conditions in $`[L,\mathrm{})`$. By the above mentioned FK representation (see e.g. \[F, ACCN, IN\]), we have
$$m_L^{0,+}(\beta )=\mu _{2,[0,L]}^{w_r}(0+\mathrm{}),$$
where $`\mu _{2,[0,L]}^{w_r}`$ stands for the random cluster measure on $`\{0,1\}^{𝔼(_+)}`$ with $`\kappa =2`$ and all the edges $`\{x,y\}`$ with $`xL`$ and $`yL`$ being open (wired on the right). Together with Remark 1 following Theorem 1.1, this yields the following
###### Corollary 1.2
For any $`\beta >2`$, there exist $`0<p_0<1`$ such that, if $`p>p_0`$, then
$$\underset{L\mathrm{}}{lim}m_L^{0,+}(\beta )\mu _{2,_+}^f(0\mathrm{})\nu (0\mathrm{})1ϵ$$
holds for $`ϵ=ϵ(J)0`$ as $`J\mathrm{}`$. Consequently, there exists a phase transition when the thermodynamical limit on $`_+`$ is taken with $`+`$ boundary conditions on the right side.
Remarks.1. In the above corollary there is a little change of notation with respect to the previously mentioned FK measure: the measure $`\mu _{2,_+}^f`$ is considered here on $`\{0,1\}^{𝔼(_+)}`$.
2. The content of Corollary 1.2 is an immediate consequence of Theorem 3.4 in \[IN\], with the result holding even for $`\beta >1`$.
It is also interesting to compare the result on oriented percolation and the previous corollary with the somehow similar question on the multiplicity of Gibbs states for Markov chains with infinite connections, where orientation appears naturally through the time direction. Recently Johansson and Öberg \[JO\] showed that if $`g`$ is a regular specification and
$$\text{var}_k(g)=sup\{g(\sigma )g(\sigma ^{})_1:\sigma _i=\sigma _i^{},i=1,\mathrm{},k\},$$
then $`g`$ admits a unique Gibbs measure whenever the sequence $`\{\text{var}_k(g)\}_{k=1}^+\mathrm{}`$ is in $`\mathrm{}^2`$. This tells, in particular, that there are no multiple limiting measures for chains with connections decaying as $`r^2`$, as in Example 1 in \[JO\]. This contrasts with the two-sided Ising models and, as our Theorem says, with percolation models. The understanding of Markov chains with infinite connections in the non-uniqueness regime is still very poor, and it is known as a notoriously difficult problem. There is strong evidence (see \[BHS\]) that multi-scale analysis techniques analogous to those developed in this work could be turned into a robust tool to study this question.
Heuristics of the proof. The proof relies on Fröhlich-Spencer multi-scale analysis ideas (\[FS\], \[FS1\]), and we use the version developed in \[KMP\] and \[M\]. In the next few paragraphs we outline the scheme of the proof, and comment on some key ideas, avoiding most of consuming technical points. Our goal here is only to give a very schematic and approximate picture, postponing precise formulations (which tend to be quite involved) to later in the text.
The goal. We look for an event of positive probability, whose occurrence implies not only the existence of an infinite open component, but also guarantees the presence of an oriented infinite open path. Essentially, we will construct such an event, and show that it has positive probability. Our key estimate will be: if $`\beta >1`$, we can find $`\delta >0`$, $`\delta ^{}<1`$ and $`p`$ sufficiently close to $`1`$, so that
$$\nu (\text{open path }\pi =(x_1,e_1,\mathrm{},x_n):x_1L+L^\delta ^{},x_nLL^\delta ^{},0<x_ix_{i1}L^\delta ^{},i)12L^\delta $$
(1.4)
for $`L=l_k`$ as defined below and any $`k1`$, $`l_1`$ being sufficiently large, and where $`\nu `$ stands for the product measure defined before. We will have little control on how close to one $`p`$ has to be (or, equivalently, on how large we need $`l_1`$).
Scales. We choose super-exponentially fast growing scales. Given $`1<\alpha <2`$, $`l_0=1`$ and $`l_1`$ an integer sufficiently large, let
$$l_k=l_{k1}^{\alpha 1}l_{k1},k=2,3,\mathrm{},$$
(1.5)
where as usual $`z=\mathrm{max}\{n:nz\}`$. We will use the so–called dynamical blocking argument, where the size and location of blocks<sup>2</sup><sup>2</sup>2Successive blocks share an end-vertex. will be defined along the procedure and will depend on the configuration at lower scales. Still, the length of each block $`I^{(k)}`$ of the $`k`$–th level (called $`k`$–block) will be of order $`l_k`$. More precisely, we shall see that $`l_k2l_k^{\alpha ^{}/\alpha }6l_{k1}|I^{(k)}|3l_k+6l_{k1}`$, for suitable $`1<\alpha ^{}<\alpha `$. (In particular, $`|I^{(k)}|l_{k+1}^{\alpha ^{}/\alpha }l_{k+1}`$, if $`k1`$ and $`l_1`$ is large.)
Defected and good blocks. Further we will use the following recursive definition of “defected” block. Fix $`1<\alpha ^{}<\alpha `$ to be specified later.
1) We say that the $`0`$–block $`[i,i+1]`$ is defected if the corresponding nearest neighbor edge $`\{i,i+1\}`$ is closed; otherwise the $`0`$–block is said to be good and the open nearest neighbor path from $`i`$ to $`i+1`$ is called a $`0`$pedestal;
2) For $`k1`$, a $`k`$–block $`I^{(k)}=[s,s^{}]`$ is defected if either it contains two or more defected $`(k1)`$–blocks, or it contains only one defected $`(k1)`$–block $`[i,i^{}]`$ but there is no open edge $`\{a,a^{}\}`$ of length at most $`l_k^{\alpha ^{}/\alpha }`$ , with $`ai,i^{}a^{},a\mathrm{{\rm Y}},a^{}\mathrm{{\rm Y}}^{}`$, for some $`(k1)`$–pedestals $`\mathrm{{\rm Y}},\mathrm{{\rm Y}}^{}`$ contained in $`I^{(k)}`$. Otherwise $`I^{(k)}`$ is called good.
Thus, if a $`k`$–block $`[s,s^{}]`$ is good, then it contains an oriented open path going from $`s`$ to $`s^{}`$: in the case it has no defected $`(k1)`$–blocks, this path can be obtained by concatenating $`(k1)`$–pedestals of the good $`(k1)`$–blocks which constitute the given $`k`$–block; if it has a (single) defected $`(k1)`$–block, a similar concatenation yields an oriented open path going from $`s`$ to $`a`$, which is followed by an open edge $`\{a,a^{}\}`$, and then followed by another concatenation of $`(k1)`$–pedestals of good $`(k1)`$–blocks, from $`a^{}`$ to $`s^{}`$. In both cases, such path from $`s`$ to $`s^{}`$ will be called $`k`$pedestal, and denoted by $`\mathrm{{\rm Y}}`$. The part of the cluster between $`a`$ and $`a^{}`$ is again disregarded in the future construction since we have little control on oriented connectivity in this segment. The condition $`a^{}al_k^{\alpha ^{}/\alpha }`$ will be crucial to guarantee that pedestals are quite dense sets (within the corresponding good blocks), used to push the construction to higher levels. Some care is needed when treating defects close to the boundary, which we have disregarded here.
Strategy. Being “defected” doesn’t necessarily imply that there is no oriented open path connecting the endpoints of the block. Nevertheless, in order to avoid substantial technical difficulties, we will follow two rules that simplify our construction:
a) once a block is defected, we will assume the worst possible situation, namely it will be considered as if all edges within this block were closed.
b) once we have at least two defected $`(k1)`$–blocks within a $`k`$–block $`I`$, we will not try to find connections within the $`k`$–block to fix its connectivity, but rather will “push the problem to the next level”, and try to “jump over” this troubled block $`I`$ by a longer edge of length at most $`l_{k+1}^{\alpha ^{}/\alpha }`$, which starts at the pedestal of some good $`k`$–block to the left of $`I`$, and ends similarly on the right of $`I`$.
Estimates. The scale $`l_1`$ will be taken large enough, to be determined later depending on the parameter $`\beta >1`$ and the auxiliary parameters $`\delta >0`$, $`1<\alpha ^{}<\alpha <2`$, to be chosen at the end of Sec.2 (see (2.17)–(2.20)). Once $`l_1`$ is chosen, we shall take $`p`$ so that:
$$p\left(1+\frac{(\mathrm{ln}2)^5}{128}l_1^{\delta 1}\right)^1.$$
(1.6)
For $`k2`$, let $`I^{(k)}`$ be a $`k`$–block of length<sup>3</sup><sup>3</sup>3This is not exact in general, but holds approximately, cf. (2.7). $`l_k`$, which consists of $`N_k=l_k/l_{k1}=l_{k1}^{\alpha 1}`$ blocks of level $`(k1)`$, of length $`l_{k1}`$, and written as $`\{I_j^{(k1)}\}_{j=1}^{N_k}`$. Assume that we have the following estimate
$$\nu (I_j^{(k1)}\text{is defected})l_{k1}^\delta ,1jN_k.$$
Under the above assumptions, and if $`\delta `$ is chosen to satisfy (2.18), we see that
$$\nu (\mathrm{\hspace{0.33em}1}i<jN_k:I_i^{(k1)},I_j^{(k1)}\text{ are both defected})\frac{1}{2}l_k^\delta .$$
(1.7)
When the defected $`I_i^{(k1)}`$ is unique, we assume for the moment that it stays at distance larger than $`l_k^{\alpha ^{}/\alpha }`$ from the boundary of $`I^{(k)}`$. (Otherwise a sequence of local adjustments of blocks will be needed, as we shall see in Sect. 2. The left- and right-most extremal blocks in our volume are treated differently.) In this case let $`a`$ and $`a^{}`$ be the end-vertices of the unique defected block $`I_i^{(k1)}`$. By our construction, there exists an oriented path starting from the left boundary of $`I^{(k)}`$ and ending at the vertex $`a`$ and another open oriented path starting from vertex $`a^{}`$ and going to the right boundary of $`I^{(k)}`$. Both these paths are obtained by concatenating pedestals of all good $`(k1)`$–blocks on the left side of the defected block $`I_i^{(k1)}`$ and, respectively, on the right side. We denote these new left and right pedestals by $`\mathrm{{\rm Y}}`$ and $`\mathrm{{\rm Y}}^{}`$, respectively. Given that $`I^{(k)}`$ has a unique defected $`I_i^{(k1)}=[a,a^{}]`$, and given the pedestals $`\mathrm{{\rm Y}}`$ and $`\mathrm{{\rm Y}}^{}`$, one has the following upper bound for the conditional $`\nu `$–probability of not finding an open edge $`\{x,y\}`$ with $`xa,a^{}y,x\mathrm{{\rm Y}},y\mathrm{{\rm Y}}^{}`$ and $`yxl_k^{\alpha ^{}/\alpha }`$ :
$$\underset{\begin{array}{c}x,y:xa<a^{}y,\\ yx<l_k^{\alpha ^{}/\alpha }\\ x\mathrm{{\rm Y}},y\mathrm{{\rm Y}}^{}\end{array}}{}q_{\{x,y\}}=\mathrm{exp}\left\{\underset{\begin{array}{c}x,y:xa<a^{}y,\\ yx<l_k^{\alpha ^{}/\alpha }\\ x\mathrm{{\rm Y}},y\mathrm{{\rm Y}}^{}\end{array}}{}\frac{\beta }{\left|xy\right|^2}\right\}l_{k1}^{\beta (1\eta )(\alpha ^{}1)},$$
(1.8)
where $`\eta =\eta (\alpha ,\alpha ^{},l_1)>0`$, and can be taken arbitrarily small if $`l_1\mathrm{}`$.
The precise statement and proof of the above estimate will be given in Lemma 2.1. It requires some work, and in order to obtain it for suitable $`\eta =\eta (\alpha ,\alpha ^{},l_1)>0`$ which can be taken arbitrarily small if $`l_1\mathrm{}`$ we will need to use certain geometric properties of pedestals $`\mathrm{{\rm Y}}`$ and $`\mathrm{{\rm Y}}^{}`$, which propagate inductively from each level into the next one. Namely, the pedestals are relatively dense sets (see (2.8) in Sect. 2) as the construction will show. Using the above estimate, writing
$`\{I^{(k)}`$ has a unique defected $`(k1)`$–block $`[a,a^{}]`$ and remains defected $`\}=`$
$`\{I^{(k)}`$ has unique defected $`(k1)`$–block $`[a,a^{}]\}`$
$`\left\{\text{ there is no open edge }\{x,y\}\text{ with }xa,a^{}y,x\mathrm{{\rm Y}},y\mathrm{{\rm Y}}^{}\text{ and }yxl_k^{\alpha ^{}/\alpha }\text{ }\right\}\text{ },`$
and since these events depend on disjoint sets of edges, we easily get:
$$\nu (I^{(k)}\text{ has a unique defected }I_i^{(k1)}\text{ and remains defected})l_{k1}^{\alpha 1\delta }l_{k1}^{\beta (1\eta )(\alpha ^{}1)}\frac{1}{2}l_k^\delta ,$$
(1.9)
provided
$$\beta (1\eta )(\alpha ^{}1)>(\delta +1)(\alpha 1).$$
(1.10)
Since $`\beta >1`$ and $`\eta =\eta (\alpha ,\alpha ^{},l_1)`$ can be taken very small provided $`l_1`$ is large, it will suffice to suitably fix the parameters $`\alpha `$ and $`\alpha ^{}`$ ($`\alpha ^{}`$ close enough to $`\alpha `$). This is done at the end of Sect. 2.
Difficulties. To carry on this scheme we have to go through several “unpleasant” and rather involved points. The use of a dynamical blocking argument, with the blocks of a given level depending not only on the size and location of lower level blocks, but also on their “status” (defected or good), requires a rather tight bookkeeping. This is expressed through what we call “itineraries”.
Once this is achieved, all necessary estimates follow along the scheme of \[FS1\] and \[KMP\].
In the next section we define the blocks and describe the dynamic renormalization procedure, proving Theorem 1.1.
## 2 Spatial blocks (Dynamic Renormalization)
Notation. For $`L`$, assumed to be large, the construction will involve the configuration $`\omega `$ restricted to the set of edges with both end-vertices in $`[L,L]`$, where $`[a,b]=[a,b]`$ throughout.<sup>4</sup><sup>4</sup>4Except in the proof of Lemma 2.1. We write $`\mathrm{\Omega }_L`$ as a shorthand for $`\mathrm{\Omega }_{[L,L]}`$. Scales $`\{l_k\}_k`$ are defined in the following way: $`l_0=1`$, given $`\beta >1`$ we shall take auxiliary parameters $`\delta >0`$, $`1<\alpha ^{}<\alpha <2`$ chosen according to (2.17)–(2.20), $`l_1`$ will be a suitably large integer and the parameter $`p<1`$ will be taken sufficiently close to $`1`$, depending on $`l_1`$. Then we let $`l_k`$ be given by (1.5).
Further we denote $`x_j^{(k)}=jl_k,j`$.
For the proof of Theorem 1.1 we may assume that $`L=l_M`$, for some $`M`$.
Throughout the text $`𝕀_A`$ stands for the indicator function of an event $`A`$, i.e. $`𝕀_A(\omega )=1`$ or 0 according to $`\omega A`$ or not.
Decomposition of events. Level 0. We set $`I_i^{(0)}=[i,i+1]`$. They are called $`0`$–blocks, and for $`i`$ such that $`I_i^{(0)}[L,L]`$ we define the events:
$$G(I_i^{(0)})=\{\omega :\omega _{\{i,i+1\}}=1\},B(I_i^{(0)})=\{\omega :\omega _{\{i,i+1\}}=0\}.$$
$`I_i^{(0)}`$ is said to be defected when $`B(I_i^{(0)})`$ occurs; otherwise it is said to be a good $`0`$–block.
Level 1. Consider the intervals $`\stackrel{~}{I}_j^{(1)}=[jl_1,(j+1)l_1]`$ and for each $`j`$ such that $`\stackrel{~}{I}_j^{(1)}[L,L]`$ we define the following partition of $`\mathrm{\Omega }_L`$:
$`G(\stackrel{~}{I}_j^{(1)})`$ $`=`$ $`{\displaystyle \underset{i=jl_1}{\overset{(j+1)l_11}{}}}G(I_i^{(0)}),`$
$`H_i(\stackrel{~}{I}_j^{(1)})`$ $`=`$ $`B(I_i^{(0)}){\displaystyle \underset{\begin{array}{c}s=jl_1\\ si\end{array}}{\overset{(j+1)l_11}{}}}G(I_s^{(0)})\text{ for }i[jl_1,(j+1)l_11],`$
$`H(\stackrel{~}{I}_j^{(1)})`$ $`=`$ $`{\displaystyle \underset{i=jl_1}{\overset{(j+1)l_11}{}}}H_i(\stackrel{~}{I}_j^{(1)}),`$
$`B(\stackrel{~}{I}_j^{(1)})`$ $`=`$ $`\left(G(\stackrel{~}{I}_j^{(1)})H(\stackrel{~}{I}_j^{(1)})\right)^c,`$ (2.1)
where $`G`$ stands for good, $`H`$ for hopeful and $`B`$ for bad, and accordingly, $`\stackrel{~}{I}_j^{(1)}`$ is said to be good (for given $`\omega `$) if it contains no defected $`0`$–blocks, “hopeful” if it contains only one defected $`0`$–block, and is said to be “bad” otherwise. When $`H_i(\stackrel{~}{I}_j^{(1)})`$ occurs, $`I_i^{(0)}`$ is called the defected 0-block in $`\stackrel{~}{I}_j^{(1)}`$.
Adjustment. Given $`\omega `$, we first consider the set of all $`j`$’s such that $`\omega H_{i_j}(\stackrel{~}{I}_j^{(1)})H(\stackrel{~}{I}_j^{(1)})`$ and such that the index $`i_j`$ of the (unique) defected block $`I_{i_j}^{(0)}\stackrel{~}{I}_j^{(1)}`$ verifies $`jl_1i_jjl_1+l_1^{\alpha ^{}/\alpha }1`$ (resp. $`(j+1)l_1l_1^{\alpha ^{}/\alpha }i_j(j+1)l_11`$).
If this set is empty in both cases, we set $`I_j^{(1)}=\stackrel{~}{I}_j^{(1)}`$ for all $`j`$’s, and say that $`G(I_j^{(1)})`$, $`H(I_j^{(1)})`$, $`B(I_j^{(1)})`$ occurs, according to the occurrence of the corresponding $`G(\stackrel{~}{I}_j^{(1)})`$, $`H(\stackrel{~}{I}_j^{(1)})`$, $`B(\stackrel{~}{I}_j^{(1)})`$.
If this set is not empty, we take arbitrarily one of such indices $`j`$; if $`\stackrel{~}{I}_j^{(1)}`$ is not the interval which contains $`L`$ (resp. $`L`$), to be treated in case 3) below, we check if $`\stackrel{~}{I}_{j1}^{(1)}`$ (resp. $`\stackrel{~}{I}_{j+1}^{(1)}`$) has a defected $`0`$–block in the sub-interval $`[jl_12l_1^{\alpha ^{}/\alpha },jl_1]`$ (resp. $`[(j+1)l_1,(j+1)l_1+2l_1^{\alpha ^{}/\alpha }1]`$).
1) If yes, then we consider a new interval $`I_{j1}^{(1)}=\stackrel{~}{I}_{j1}^{(1)}\stackrel{~}{I}_j^{(1)}`$ (resp. $`I_j^{(1)}=\stackrel{~}{I}_j^{(1)}\stackrel{~}{I}_{j+1}^{(1)}`$) and say that the event $`B(I_{j1}^{(1)})`$ (resp. $`B(I_j^{(1)})`$) occurs. (This is motivated by the fact that for the chosen $`\omega `$ the new interval will contain at least two defected $`0`$–blocks.)
2) If not, then we consider two new intervals $`I_{j1}^{(1)}=[(j1)l_1,jl_1l_1^{\alpha ^{}/\alpha }]`$ and $`I_j^{(1)}=[jl_1l_1^{\alpha ^{}/\alpha },(j+1)l_1]`$ (resp. $`I_j^{(1)}=[jl_1,(j+1)l_1+l_1^{\alpha ^{}/\alpha }]`$ and $`I_{j+1}^{(1)}=[(j+1)l_1+l_1^{\alpha ^{}/\alpha },(j+2)l_1]`$). We say that $`H(I_j^{(1)})`$ occurs, and that $`G(I_{j1}^{(1)})`$, $`H(I_{j1}^{(1)})`$, $`B(I_{j1}^{(1)})`$ occurs according to the occurrence of the corresponding event $`G(\stackrel{~}{I}_{j1}^{(1)})`$, $`H(\stackrel{~}{I}_{j1}^{(1)})`$, $`B(\stackrel{~}{I}_{j1}^{(1)})`$ (resp. we say that $`H(I_j^{(1)})`$ occurs, and that $`G(I_{j+1}^{(1)})`$, $`H(I_{j+1}^{(1)})`$, $`B(I_{j+1}^{(1)})`$ occurs according to the occurrence of the corresponding event $`G(\stackrel{~}{I}_{j+1}^{(1)})`$, $`H(\stackrel{~}{I}_{j+1}^{(1)})`$ $`B(\stackrel{~}{I}_{j+1}^{(1)})`$). In this case the adjustment moves the boundary “away” from the unique defected block in $`I_j^{(1)}`$, but doesn’t change the number of the defected $`0`$–blocks in the adjusted intervals.
3) If the interval $`\stackrel{~}{I}_j^{(1)}`$ under consideration is the leftmost (resp. the rightmost) interval in $`[L,L]`$, and the defect stays within distance less than $`l_1^{\alpha ^{}/\alpha }`$ from $`L`$ (resp. $`L`$), we still set $`I_j^{(1)}=\stackrel{~}{I}_j^{(1)}`$ and say that $`G(I_j^{(1)})`$ occurs.
4) We set $`I_j^{(1)}=\stackrel{~}{I}_j^{(1)}`$ if $`\stackrel{~}{I}_j^{(1)}`$ was not involved in the previous adjustment, and say that $`G(I_j^{(1)})`$, $`H(I_j^{(1)})`$, $`B(I_j^{(1)})`$ occurs if, accordingly, $`G(\stackrel{~}{I}_j^{(1)})`$, $`H(\stackrel{~}{I}_j^{(1)})`$, $`B(\stackrel{~}{I}_j^{(1)})`$ occurs.
To conclude this step, we re-numerate the intervals from left to right as $`I_j^{(1)}j=1,\mathrm{}`$. If we are still left with intervals $`I_j^{(1)}`$ for which $`H(I_j^{(1)})`$ occurs and its defected $`0`$–block stays within distance $`l_1^{\alpha ^{}/\alpha }`$ from the boundary of $`I_j^{(1)}`$, we repeat the above procedure to the intervals already adjusted in the previous step. After finitely many steps of such adjustment procedure there are left no intervals $`I_j^{(1)}`$ for which the event $`H(I_j^{(1)})`$ occurs and its defected $`0`$–block stays within distance $`l_1^{\alpha ^{}/\alpha }`$ from the boundary, and the adjustment procedure is then stopped. (Of course, due to item 3, the left- or rightmost intervals can stay with a unique defect, if this is close enough to $`L`$ or $`L`$ respectively.)
Remark. Notice that the adjustment procedure is well defined, i.e. the final partition does not depend on the order in which we do adjustments and in which order we pick the intervals that still need to be adjusted (in case we have more than one). It also has a locality property, i.e. the final modification of each initial interval $`\stackrel{~}{I}_j^{(1)}`$ depends on the values of the configuration in the nearest neighbor and, at most, in the next nearest neighbor intervals only.
Once the adjustment is completed, the obtained intervals, always re-numerated from left to right as $`I_j^{(1)},j=1,\mathrm{}`$, are called $`1`$-blocks. Notice that $`l_12l_1^{\alpha ^{}/\alpha }|I_j^{(1)}|3l_1`$, and $`_jI_j^{(1)}=[L,L]`$.
In other words, the restriction of $`\omega `$ to nearest neighbor edges of $`[L,L]`$ determines through the above procedure a random “partition” $`I^{(1)}(\omega )\{I_j^{(1)}(\omega )\}_j`$ of the interval $`[L,L]`$ into $`1`$–blocks, with the property that any two adjacent blocks share an end-vertex. This is the final state of the “adjustment” procedure. Values of $`\omega `$ on the nearest neighbor edges in $`[L,L]`$ also determine where the defected 0-blocks are located within each 1-block, and we denote by $`D_j^{(1)}(\omega )`$ the set of indices of the defected $`0`$–blocks within $`I_j^{(1)}(\omega )`$, and $`D^{(1)}(\omega )\{D_j^{(1)}(\omega )\}_j`$. The random object $`J_L^{(1)}:=\{I_j^{(1)},D_j^{(1)}\}`$ is called itinerary at level 1 or 1-itinerary.
1-Pedestals. Given the $`1`$–itinerary $`J_L^{(1)}`$, we shall attribute to each random block $`I_j^{(1)}`$ a state $`G`$ or $`B`$. We first consider the case that $`I_j^{(1)}`$ is not the leftmost (i.e. $`j1`$) nor the rightmost $`1`$–block, to be treated at the end. When $`D_j^{(1)}=\mathrm{}`$, so that all nearest neighbor edges are open, we say that $`I_j^{(1)}`$ is in state $`G`$, and we define the pedestal $`\mathrm{{\rm Y}}(I_j^{(1)})=I_j^{(1)}`$. When $`|D_j^{(1)}|=1`$, the set of vertices $`xI_j^{(1)}`$ to the left (resp. right) of the defected $`0`$–block in $`I_j^{(1)}`$ will be called left 1-pedestal of $`I_j^{(1)}`$ (resp. right 1-pedestal) and denoted by $`\mathrm{{\rm Y}}_{}(I_j^{(1)})`$ (resp. $`\mathrm{{\rm Y}}_{}(I_j^{(1)})`$). The vertices in each of these 1-pedestals are connected by open nearest neighbor edges. In this situation we say that $`I_j^{(1)}`$ is in state $`G`$ when the following event occurs:
$$[\omega :x\mathrm{{\rm Y}}_{}(I_j^{(1)}(\omega )),y\mathrm{{\rm Y}}_{}(I_j^{(1)}(\omega )),\mathrm{\hspace{0.33em}1}<yxl_1^{\alpha ^{}/\alpha }:\omega _{\{x,y\}}=1],$$
(2.2)
and otherwise we say that $`I_j^{(1)}`$ is in state $`B`$. Similarly, if $`|D_j^{(1)}|>1`$ the block $`I_j^{(1)}`$ is in state $`B`$.
For the leftmost (rightmost) $`1`$–block, there is some little difference: In the case $`|D_j^{(1)}|=1`$ and if the unique defected $`0`$–block stays within distance $`l_1^{\alpha ^{}/\alpha }`$ from $`L`$ ($`L`$), the block is said to be in state $`G`$, and the pedestal is defined as the previously defined right 1-pedestal (left 1-pedestal, resp.), $`\mathrm{{\rm Y}}(I_1^{(1)})=\mathrm{{\rm Y}}_{}(I_1^{(1)})`$ ($`\mathrm{{\rm Y}}(I_j^{(1)})=\mathrm{{\rm Y}}_{}(I_j^{(1)})`$, resp.). Except for this, the definition goes as with the other blocks.
With a little abuse of notation we use again the symbols $`G(I_j^{(1)})`$ and $`B(I_j^{(1)})`$ to denote that $`I_j^{(1)}`$ is in state $`G`$ and $`B`$ respectively. We say that $`I_j^{(1)}(\omega )`$ is defected if and only if it is in state $`B`$.
In (2.2), if the pair $`(x,y)`$ such that $`x\mathrm{{\rm Y}}_{}(I_j^{(1)}),y\mathrm{{\rm Y}}_{}(I_j^{(1)}),yxl_1^{\alpha ^{}/\alpha }`$, $`\omega _{\{x,y\}}=1`$ is not unique, we choose one in arbitrary way, and, once the pair $`(x,y)`$ is chosen, the interval $`[x+1,y1]`$ will be called defected part of $`I_j^{(1)}`$, and denoted by $`𝒟(I_j^{(1)})`$. In this case we define $`\mathrm{{\rm Y}}(I_j^{(1)})=\left(\mathrm{{\rm Y}}_{}(I_j^{(1)})\mathrm{{\rm Y}}_{}(I_j^{(1)})\right)𝒟(I_j^{(1)})`$.
In particular, a $`1`$-pedestal $`\mathrm{{\rm Y}}(I_j^{(1)})`$ is given by the vertices of an open oriented path with all edges, except possibly one, being nearest neighbor, and this larger edge has length at most $`l_1^{\alpha ^{}/\alpha }`$. For each $`1`$–block, except possibly the two which contain the extremes $`L`$ or $`L`$, the pedestal connects left and right endpoints of the interval. In the leftmost (rightmost) case, it is allowed for the $`1`$–pedestal to start (end) at a vertex within distance $`l_1^{\alpha ^{}/\alpha }+1`$ of $`L`$ ($`L`$ respectively).
Level k . Let $`2kM`$. Assume to have completed the step $`(k1)`$ of the recursion. In particular, for each $`\omega \mathrm{\Omega }_L`$ and any $`r=1,\mathrm{},k1`$ the following objects are defined:
* the collection of $`r`$–blocks $`I^{(r)}(\omega )=\{I_j^{(r)}(\omega )\}_j`$, such that $`_jI_j^{(r)}(\omega )=[L,L]`$, and any two adjacent intervals share exactly an endpoint. Moreover, the uniform bound holds:
$`l_r(2l_r^{\alpha ^{}/\alpha }+6l_{r1})<|I_j^{(r)}(\omega )|3l_r+6l_{r1},`$ (2.3)
* each of the $`I_j^{(r)}(\omega )`$ can be in two possible states $`G`$ or $`B`$:
If $`I_j^{(r)}(\omega )`$ is in state $`G`$ and it is not the leftmost or the rightmost interval of the partition, then $`\omega `$ has an $`r`$–pedestal $`\mathrm{{\rm Y}}(I_j^{(r)})`$ given by vertices of an open oriented path from the left to the right boundary of $`I_j^{(r)}(\omega )`$. If $`I_j^{(r)}(\omega )`$ is the leftmost (resp. the rightmost) interval, an $`r`$–pedestal $`\mathrm{{\rm Y}}(I_j^{(r)})`$ is given by vertices of an open oriented path which starts from some vertex $`x[L,L+2l_r^{\alpha ^{}/\alpha }]`$ and ends at the right boundary of $`I_j^{(r)}(\omega )`$ (resp. starts from the left boundary of $`I_j^{(r)}(\omega )`$ and ends at some vertex $`x[L2l_r^{\alpha ^{}/\alpha },L]`$). ($`l_1`$ being large, we may assume that the length of an $`(r1)`$–block is always bounded above by $`l_r^{\alpha ^{}/\alpha }`$, according to (2.3) for $`r`$ replaced by $`r1`$.)
* the collection $`D^{(r)}(\omega )=\{D_j^{(r)}(\omega )\}_j`$, where $`D_j^{(r)}(\omega )`$ is the set of labels of the defected $`(r1)`$–blocks which are contained in $`I_j^{(r)}(\omega )`$.
For $`\omega `$ fixed, the sequence of pairs
$$J_L^{(k1)}(\omega )=\{(I^{(1)}(\omega ),D^{(1)}(\omega )),\mathrm{},(I^{(k1)}(\omega ),D^{(k1)}(\omega ))\},$$
is called $`(k1)`$itinerary, and $`(I^{(r)},D^{(r)})`$, is called the $`r`$–th step of the itinerary, for $`1rk1`$. We shall now see how to define the $`k`$–blocks and the continuation to a $`k`$–itinerary. When $`k=M`$ we will end up with only one or two intervals.
Construction of $`k`$blocks. For any $`\omega `$ and for each $`z[L,L]`$ we set $`j_z^k=\mathrm{min}\{j:zI_j^{(k1)}\}`$, $`\widehat{ȷ}_i^k=j_{x_i^{(k)}}^k`$, cf. notation at the beginning of this section, $`i=l_M/l_k,\mathrm{},l_M/l_k1`$, and define the intervals:
$$\stackrel{~}{I}_i^{(k)}=\underset{s=\widehat{ȷ}_i^k+1}{\overset{\widehat{ȷ}_{i+1}^k}{}}I_s^{(k1)}=:[a_i^{(k)},a_{i+1}^{(k)}]$$
as well as the following partition of $`\mathrm{\Omega }_L`$:
$`G(\stackrel{~}{I}_i^{(k)})`$ $`=`$ $`{\displaystyle \underset{s=\widehat{ȷ}_i^k+1}{\overset{\widehat{ȷ}_{i+1}^k}{}}}G(I_s^{(k1)}),`$
$`H_s(\stackrel{~}{I}_i^{(k)})`$ $`=`$ $`B(I_s^{(k1)}){\displaystyle \underset{u=\widehat{ȷ}_i^k+1,us}{\overset{\widehat{ȷ}_{i+1}^k}{}}}G(I_u^{(k1)}),`$
$`H(\stackrel{~}{I}_i^{(k)})`$ $`=`$ $`{\displaystyle \underset{s=\widehat{ȷ}_i^k+1}{\overset{\widehat{ȷ}_{i+1}^k}{}}}H_s(\stackrel{~}{I}_i^{(k)}),`$
$`B(\stackrel{~}{I}_i^{(k)})`$ $`=`$ $`\mathrm{\Omega }_L\left(G(\stackrel{~}{I}_i^{(k)})H(\stackrel{~}{I}_i^{(k)})\right).`$ (2.4)
Adjustment. Given $`\omega \mathrm{\Omega }_L`$, consider all $`i`$ for which $`H_s(\stackrel{~}{I}_i^{(k)})`$ occurs for $`s`$ such that the distance of the defected $`(k1)`$–block $`I_s^{(k1)}\stackrel{~}{I}_i^{(k)}`$ to the left endpoint $`a_i^{(k)}`$ (right endpoint $`a_{i+1}^{(k)}`$, resp.) is less than $`l_k^{\alpha ^{}/\alpha }`$. If this set is non-empty take arbitrarily any such $`\stackrel{~}{I}_i^{(k)}`$.
When the selected $`\stackrel{~}{I}_i^{(k)}`$ is the leftmost (resp. the rightmost) interval in $`[L,L]`$, and the defect stays at distance less than $`l_k^{\alpha ^{}/\alpha }`$ from $`L`$ (resp. $`L`$), we set $`I_i^{(k)}=\stackrel{~}{I}_i^{(k)}`$, and say that $`G(I_i^{(k)})`$ occurs (or that $`I_i^{(k)}`$ is in $`G`$ state for this $`\omega `$). Otherwise, we then check if $`\stackrel{~}{I}_{i1}^{(k)}`$ (resp. $`\stackrel{~}{I}_{i+1}^{(k)}`$) has a defected block $`I_r^{(k1)}`$ at distance at most $`3l_k^{\alpha ^{}/\alpha }`$ from $`a_i^{(k)}`$ (resp. from $`a_{i+1}^{(k)}`$), and
1) If yes, then we consider a new interval $`I_{i1}^{(k)}=\stackrel{~}{I}_{i1}^{(k)}\stackrel{~}{I}_i^{(k)}`$ (respectively $`I_i^{(k)}=\stackrel{~}{I}_i^{(k)}\stackrel{~}{I}_{i+1}^{(k)}`$) and say that $`B(I_{i1}^{(k)})`$ (resp. $`B(I_i^{(k)})`$) occurs, or that the corresponding interval is in state $`B`$;
2) If not, then we consider two new intervals:
$`I_{i1}^{(k)}={\displaystyle \underset{s=\widehat{ȷ}_{i1}^k+1}{\overset{j_{a_i^{(k)}l_k^{\alpha ^{}/\alpha }}^k1}{}}}I_s^{(k1)},`$ $`I_i^{(k)}={\displaystyle \underset{s=j_{a_i^{(k)}l_k^{\alpha ^{}/\alpha }}^k}{\overset{\widehat{ȷ}_{i+1}^k}{}}}I_s^{(k1)}`$ (2.5)
$`(\text{respectively,}I_i^{(k)}={\displaystyle \underset{s=\widehat{ȷ}_i^k+1}{\overset{j_{a_{i+1}^{(k)}+l_k^{\alpha ^{}/\alpha }}^k}{}}}I_s^{(k1)}`$ $`I_{i+1}^{(k)}={\displaystyle \underset{s=j_{a_{i+1}^{(k)}+l_k^{\alpha ^{}/\alpha }}^k+1}{\overset{\widehat{ȷ}_{i+2}^k}{}}}I_s^{(k1)}).`$ (2.6)
In the situation of (2.5) we say that $`H(I_i^{(k)})`$ occurs, and say that $`G(I_{i1}^{(k)})`$, $`H(I_{i1}^{(k)})`$, $`B(I_{i1}^{(k)})`$ occurs according to the occurrence of the corresponding $`G(\stackrel{~}{I}_{i1}^{(k)})`$, $`H(\stackrel{~}{I}_{i1}^{(k)})`$, $`B(\stackrel{~}{I}_{i1}^{(k)})`$ (resp. in the situation of (2.6) we say that $`H(I_{i1}^{(k)})`$ occurs, and say that $`G(I_{i+1}^{(k)})`$, $`H(I_{i+1}^{(k)})`$, $`B(I_{i+1}^{(k)})`$ occurs according to the occurrence of $`G(\stackrel{~}{I}_{i+1}^{(k)})`$, $`H(\stackrel{~}{I}_{i+1}^{(k)})`$, $`B(\stackrel{~}{I}_{i+1}^{(k)})`$).
Finally we set $`I_i^{(k)}=\stackrel{~}{I}_i^{(k)}`$ if $`\stackrel{~}{I}_i^{(k)}`$ was not involved in the adjustment and say $`G(I_i^{(k)})`$, $`H(I_i^{(k)})`$, $`B(I_i^{(k)})`$ occurs if the corresponding $`G(\stackrel{~}{I}_i^{(k)})`$, $`H(\stackrel{~}{I}_i^{(k)})`$, $`B(\stackrel{~}{I}_i^{(k)})`$ does occur.
To conclude this step, we re-numerate the intervals from left to right as $`I_j^{(k)}j=1,\mathrm{}`$. If after this step we are still left with intervals $`I_i^{(k)}`$ for which $`H(I_i^{(k)})`$ occurs and its defected interval $`I_s^{(k1)}`$ stays within distance $`l_k^{\alpha ^{}/\alpha }`$ from one of the endpoints of $`I_i^{(k)}`$, then we repeat the above procedure. After finitely many steps of this adjustment procedure all $`I_i^{(k)}`$ for which $`H(I_i^{(k)})`$ occurs have their defected $`(k1)`$–block at distance larger than $`l_k^{\alpha ^{}/\alpha }`$ from the boundary of $`I_i^{(k)}`$.
Once the adjustments are completed, the final intervals, always re-numerated from left to right as $`I_j^{(k)},j=1,\mathrm{}`$, are called $`k`$–blocks. We then consider the collection $`D^{(k)}=\{D_j^{(k)}\}_j`$ where $`D_j^{(k)}`$ gives the labels of the defected $`(k1)`$–blocks contained in $`I_j^{(k)}`$.
We can always write $`I_j^{(k)}={\displaystyle _{s_0(j)}^{s_1(j)}}I_s^{(k1)}`$. It is easy to check that the procedure is well defined (measurable) and the validity of the following recursive estimate:
$$l_k(2l_k^{\alpha ^{}/\alpha }+6l_{k1})<|I_j^{(k)}|3l_k+6l_{k1}.$$
(2.7)
$`k`$-Pedestals. Given the $`k`$–itinerary we shall associate to each $`k`$–block $`I_j^{(k)}(\omega )`$ a state $`G`$ or $`B`$, and the blocks in state $`G`$ will have a $`k`$–pedestal, to be defined below. When $`|D_j^{(k)}(\omega )|2`$, the block is said to be in state $`B`$, and it has no $`k`$–pedestal.
* When $`D(I_j^{(k)})=\mathrm{}`$, all its sub-blocks $`I_s^{(k1)}`$ are in state $`G`$. In this case we define $`\mathrm{{\rm Y}}(I_j^{(k)})={\displaystyle _{s_0(j)}^{s_1(j)}}\mathrm{{\rm Y}}(I_s^{(k1)})`$.
* If $`D(I_j^{(k)})=\{r\}`$ and $`I_j^{(k)}`$ is not the leftmost (resp. rightmost) interval in $`[L,L]`$, we define $`\mathrm{{\rm Y}}_{}(I_j^{(k)})={\displaystyle _{s_0(j)}^{r1}}\mathrm{{\rm Y}}(I_s^{(k1)})`$ and $`\mathrm{{\rm Y}}_{}(I_j^{(k)})={\displaystyle _{r+1}^{s_1(j)}}\mathrm{{\rm Y}}(I_s^{(k1)}),`$ called left and right pedestals<sup>5</sup><sup>5</sup>5From the occurrence of $`G(I_s^{(k1)})`$ for all other $`(k1)`$–blocks within $`I_j^{(k)}`$, we know there exists an open oriented path connecting the left boundary of $`I_j^{(k)}`$ to the right boundary of $`I_{r1}^{(k1)}`$ and an open oriented path connecting the left boundary of $`I_{r+1}^{(k1)}`$ to the right boundary of $`I_j^{(k)}`$. These paths are obtained by concatenation of the corresponding $`\mathrm{{\rm Y}}(I_s^{(k1)}),s_0(j)sr1`$ and $`r+1ss_1(j)`$, respectively. of $`I_j^{(k)}`$, and check if there exists $`x\mathrm{{\rm Y}}_{}(I_j^{(k)})`$ and $`y\mathrm{{\rm Y}}_{}(I_j^{(k)})`$ with $`yxl_k^{\alpha ^{}/\alpha }`$ such that $`\omega _{\{x,y\}}=1`$:
– If yes, we say that $`I_j^{(k)}`$ is in state $`G`$, and if the pair $`(x,y)`$ with $`x\mathrm{{\rm Y}}_{}(I_j^{(k)}),y\mathrm{{\rm Y}}_{}(I_j^{(k)}),yxl_k^{\alpha ^{}/\alpha }`$, and $`\omega _{\{x,y\}}=1`$ is not unique, we choose one in an arbitrary way, and, once $`(x,y)`$ is chosen, denote $`𝒟(I_j^{(k)})=[x+1,y1]`$, and define
$$\mathrm{{\rm Y}}(I_j^{(k)})=\left(\mathrm{{\rm Y}}_{}(I_j^{(k)})\mathrm{{\rm Y}}_{}(I_j^{(k)})\right)𝒟(I_j^{(k)}).$$
– If such an open edge $`\{x,y\}`$ does not exist we say that $`I_j^{(k)}`$ is in $`B`$ state.
* If $`|D(I_j^{(k)})|=1`$ and $`I_j^{(k)}`$ is the leftmost (resp. rightmost) interval in $`[L,L]`$ whose unique defected $`(k1)`$–block $`I_r^{(k1)}`$ stays within distance $`l_k^{\alpha ^{}/\alpha }`$ from $`L`$ (resp. $`L`$), then we say that $`I_j^{(k)}`$ is in state $`G`$ and we define its $`k`$–pedestal as $`\mathrm{{\rm Y}}(I_j^{(k)})={\displaystyle _{r+1}^{s_1(j)}}\mathrm{{\rm Y}}(I_s^{(k1)})`$ (resp. $`\mathrm{{\rm Y}}(I_j^{(k)})={\displaystyle _{s_0(j)}^{r1}}\mathrm{{\rm Y}}(I_s^{(k1)})`$).
* Finally, if $`|D(I_j^{(k)})|=1`$ and $`I_j^{(k)}`$ is the leftmost (resp. rightmost) interval in $`[L,L]`$, but its unique defected $`(k1)`$–block $`I_r^{(k1)}`$ does not stay within distance $`l_k^{\alpha ^{}/\alpha }`$ from $`L`$ (resp. $`L`$), then we use the same procedure as if $`I_j^{(k)}`$ were not an extremal $`k`$–block.
This completes the $`k`$–th step, associating with each itinerary $`J^{(k1)}`$ its continuation with a random sequence of $`k`$–blocks $`I^{(k)}=\{I_j^{(k)}\}_j`$, re-numerated from left to right. Moreover, with each $`k`$–block we associate one of the states $`G`$ or $`B`$.
Structure of pedestals. First we state a simple geometric property of pedestals, which will be used in estimating the conditional probability that a $`k`$–block $`I_j^{(k)}`$ is in state $`G`$, given that $`|D_j^{(k)}|=1`$. Our goal is to show that there exists a positive constant $`CC(\alpha ,\alpha ^{})`$ such that if a $`k`$–block, $`k1`$, $`I^{(k)}=[s,s^{}]`$ contains only one defected $`(k1)`$–block, here denoted by $`[a,a^{}]`$, with corresponding left and right pedestals $`\mathrm{{\rm Y}}_{}`$ and $`\mathrm{{\rm Y}}_{}`$, spanning from $`s`$ to $`a`$ and from $`a^{}`$ to $`s^{}`$, respectively, then
$$\left|\mathrm{{\rm Y}}_{}[al_k^{\alpha ^{}/\alpha },a]\right|Cl_k^{\alpha ^{}/\alpha }\text{and}\left|\mathrm{{\rm Y}}_{}[a^{},a^{}+l_k^{\alpha ^{}/\alpha }]\right|Cl_k^{\alpha ^{}/\alpha }.$$
(2.8)
Inequality (2.8) follows trivially from the following recursive relation: if we have a $`k`$–block $`I^{(k)}=_{s_0}^{s_1}I_s^{(k1)}`$ which is in $`G`$ state, then
$$\left|\mathrm{{\rm Y}}(I^{(k)})\right|\underset{s:[G(I_s^{(k1)})\text{occurs}]}{}\left|\mathrm{{\rm Y}}(I_s^{(k1)})\right|l_k^{\alpha ^{}/\alpha }.$$
We now give the announced basic estimate needed for the recursive step in the previous construction, cf. (1.8). Afterwards, we fix the parameters which will determine the choice of $`p`$ close to one, as in (1.6). In the lemma below, assume that $`I_j^{(k)}`$ is a $`k`$–block and $`D_j^{(k)}=\{z\}`$, i.e. the unique defected $`(k1)`$–block within $`I_j^{(k)}`$ has index $`z`$, and by construction stays at distance larger than $`l_k^{\alpha ^{}/\alpha }`$ from the boundaries of $`I_j^{(k)}`$.
###### Lemma 2.1
There exists $`\eta \eta (\alpha ,\alpha ^{},l_1)`$ with $`\eta 0`$ as $`l_1+\mathrm{}`$ and such that the following estimate for the conditional probability with respect to the product measure (defined right above (1.1))
$$\nu [x\mathrm{{\rm Y}}_{}(I_j^{(k)}),y\mathrm{{\rm Y}}_{}(I_j^{(k)}),yxl_k^{\alpha ^{}/\alpha }:\omega _{\{x,y\}}=1||D_j^{(k)}|=1]1l_{k1}^{\beta (1\eta )(\alpha ^{}1)}$$
(2.9)
holds for $`k2`$. For $`k=1`$ the r.h.s in (2.9) is replaced by $`1l_1^{\beta (1\eta )(\alpha ^{}1)/\alpha }`$ .
Proof. We show the above estimate by conditioning on $`D_j^{(k)}=\{z\}`$, uniformly in $`z`$, and we make repeated use of the following upper and lower bounds: if $`I`$ and $`I^{}`$ are two intervals, and $`3d=\text{dist}(I,I^{})`$, then
$$C^{}J(I,I^{})\underset{\begin{array}{c}xI\\ yI^{}\end{array}}{}\frac{1}{|xy|^2}C^+J(I,I^{}),$$
(2.10)
holds with $`C^\pm =\left(1\pm 2/d\right)^2`$ and
$$J(I,I^{})=_{I\times I^{}}𝑑x𝑑y\frac{1}{|xy|^2}=\mathrm{ln}\frac{(|I|+d)(|I^{}|+d)}{d(|I|+|I^{}|+d)}.$$
(2.11)
Notice that we have $`C^{}\left(\left|xy\right|2\right)^2\left|xy\right|^2C^+\left(\left|xy\right|+2\right)^2`$ for $`\left|xy\right|d`$. We shall need also the inequality
$$J(I,I^{})4\frac{\left|I^{}\right|}{\left|I^{\prime \prime }\right|}J(I,I^{\prime \prime })$$
(2.12)
which holds for every $`I`$, $`I^{}`$ and $`I^{\prime \prime }`$ such that $`I^{}I^{\prime \prime }`$ and $`d^{}=\text{dist}(I,I^{\prime \prime })\left|I^{\prime \prime }\right|`$. Indeed, setting $`f(x)=_I𝑑y|xy|^2`$, for $`xI^{\prime \prime }`$, straightforward calculations give that under the above conditions:
$$f(x^{})4f(x^{\prime \prime })\text{for each }x^{}I^{},x^{\prime \prime }I^{\prime \prime }$$
from where the inequality (2.12) follows upon integration.
If $`k2`$ and $`D_j^{(k)}=\{z\}`$, we have the left $`k`$–pedestal $`\mathrm{{\rm Y}}_{}(I_j^{(k)})`$ spanning from the left endpoint of $`I_j^{(k)}`$ to the left endpoint of $`I_z^{(k1)}`$, and the right $`k`$ -pedestal $`\mathrm{{\rm Y}}_{}(I_j^{(k)})`$, spanning from the right endpoint of $`I_z^{(k1)}`$ to the right endpoint of $`I_j^{(k)}`$. Take two segments $`S_z^{}`$ and $`S_z^{}`$, such that $`|S_z^{}|=|S_z^{}|=l_k^{\alpha ^{}/\alpha }/3`$, lying immediately to the left and, respectively, to the right of $`I_z^{(k1)}`$. Denote
$`\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})`$ $`=`$ $`\mathrm{{\rm Y}}_{}(I_j^{(k)})S_z^{},`$
$`\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})`$ $`=`$ $`\mathrm{{\rm Y}}_{}(I_j^{(k)})S_z^{}.`$
Then
$$\nu [\text{all edges }\{x,y\},x\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)}),y\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})\text{ are closed}|D_j^{(k)}=\{z\}]$$
$$\underset{\begin{array}{c}xS_z^{}\\ yS_z^{}\end{array}}{}q_{\{x,y\}}\underset{\begin{array}{c}xS_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})\\ yS_z^{}\end{array}}{}q_{\{x,y\}}^1\underset{\begin{array}{c}xS_z^{}\\ yS_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})\end{array}}{}q_{\{x,y\}}^1.$$
(2.13)
Applying (2.10) to $`S_z^{}`$ and $`S_z^{}`$ we immediately get the following bound:
$$\underset{\begin{array}{c}xS_z^{}\\ yS_z^{}\end{array}}{}q_{\{x,y\}}=\mathrm{exp}\left\{\underset{\begin{array}{c}xS_z^{}\\ yS_z^{}\end{array}}{}\frac{\beta }{|xy|^2}\right\}l_{k1}^{\beta (\alpha ^{}1)(1b)},$$
(2.14)
where $`bb(\alpha ^{},l_1)`$ and $`b0`$ when $`l_1+\mathrm{}`$. Similar computation gives that if a 1-block $`I`$ has a unique closed edge $`\{a,a+1\}`$ with both $`a,a+1`$ at distance larger than $`l_1^{\alpha ^{}/\alpha }`$ from the endpoints of $`I`$, then the probability that there is an open edge $`\{x,y\}`$ with $`x<a<y`$, $`yxl_1^{\alpha ^{}/\alpha }`$ is larger than or equal of $`1l_1^{\beta (1\eta )(\alpha ^{}1)/\alpha }`$.
On the other hand denoting by $`𝒟_n(S_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})),\mathrm{\hspace{0.33em}0}nk2`$ (resp. $`𝒟_n(S_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})))`$ the set of vertices that belong to all defected $`n`$–blocks contained in the segment $`S_z^{}`$ (resp. $`S_z^{}`$), we get
$$\underset{\begin{array}{c}xS_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})\\ yS_z^{}\end{array}}{}q_{\{x,y\}}=\underset{n=0}{\overset{k2}{}}\underset{\begin{array}{c}x𝒟_n(S_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)}))\\ yS_z^{}\end{array}}{}q_{\{x,y\}}=\mathrm{exp}\left\{\underset{n=0}{\overset{k2}{}}\underset{\begin{array}{c}x𝒟_n(S_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)}))\\ yS_z^{}\end{array}}{}\frac{\beta }{|xy|^2}\right\}.$$
Once again, applying (2.10) for each $`0nk2`$ and taking into account the structure of $`n`$–pedestals together with (2.12), we have (uniformly on all $`l_1`$ large enough) fixed positive constants $`C_i,i=1,2,3`$ so that
$`{\displaystyle \underset{n=0}{\overset{k2}{}}}{\displaystyle \underset{\begin{array}{c}x𝒟_n(S_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)}))\\ yS_z^{}\end{array}}{}}{\displaystyle \frac{\beta }{|xy|^2}}`$ $``$ $`C_1{\displaystyle \underset{n=0}{\overset{k2}{}}}{\displaystyle \underset{\nu }{}}J(I_\nu ^{},I^{})`$
$``$ $`C_2{\displaystyle \underset{n=0}{\overset{k2}{}}}{\displaystyle \frac{l_{n+1}^{\alpha ^{}/\alpha }}{l_{n+1}}}{\displaystyle \underset{\nu }{}}J(I_\nu ^{\prime \prime },I^{})`$
$``$ $`C_3l_1^{\alpha ^{}/\alpha 1}J(I^{},I^{})`$
where $`I_\nu ^{}`$ and $`I_\nu ^{\prime \prime }`$ are intervals in $``$ so that $`_\nu \left(I_\nu ^{}\right)=𝒟_n(S_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)}))`$, the sum $`_\nu `$ is taken over all indices $`\nu `$ of $`(n+1)`$–blocks $`I_\nu ^{(n+1)}=:I_\nu ^{\prime \prime }`$ where the defected $`n`$–blocks are located, and moreover, $`I^{}=_{0nk2}_\nu I_\nu ^{\prime \prime }`$ and $`I^{}`$ is the convex envelop of $`S_z^{}`$. The condition to apply (2.12) in the first inequality above follows from $`3l_{k1}+6l_{k2}l_k^{\alpha ^{}/\alpha }`$ which is true for any $`k2`$, provided $`l_1`$ has been taken large enough. From this we can easily get that
$$\underset{\begin{array}{c}xS_z^{}\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})\\ yS_z^{}\end{array}}{}q_{\{x,y\}}l_{k1}^{\beta (\alpha ^{}1)b^{}},$$
(2.15)
where $`b^{}b^{}(\alpha ,\alpha ^{},l_1)`$ and $`b^{}0`$ when $`l_1+\mathrm{}`$. Analogous lower bound holds for the third term at the r.h.s of (2.13). Finally, from the upper bound for the length of a $`(k1)`$-block, we have
$`[\omega :x\mathrm{{\rm Y}}_{}(I_j^{(k)}),y\mathrm{{\rm Y}}_{}(I_j^{(k)}),yxl_k^{\alpha ^{}/\alpha }:\omega _{\{x,y\}}=1]^c`$ (2.16)
$``$ $`[\text{all edges }\{x,y\},x\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)}),y\widehat{\mathrm{{\rm Y}}}_{}(I_j^{(k)})\text{ are closed}],`$
the statement of the Lemma follows from (2.14) and (2.15). $`\mathrm{}`$
Fixing the parameters. For fixed $`\beta >1`$, which is the first main parameter of the model we choose the pair $`\alpha ,\alpha ^{}`$ with $`1<\alpha ^{}<\alpha <2`$ such that
$$\beta (\alpha ^{}1)\frac{2(\alpha 1)^2}{2\alpha }>\alpha 1,$$
(2.17)
i.e. $`\beta (\alpha ^{}1)>\alpha \left(\alpha 1\right)/(2\alpha )`$. We also fix
$$\delta >\frac{2(\alpha 1)}{2\alpha }$$
(2.18)
such that
$$\beta (\alpha ^{}1)\delta (\alpha 1)>\alpha 1.$$
(2.19)
By Lemma 2.1 we can fix $`l_1>1`$ so large that the parameter $`\eta =\eta (\alpha ,\alpha ^{},l_1)`$ in (2.9) becomes so close to zero, that
$$\beta (1\eta )(\alpha ^{}1)\delta (\alpha 1)>\alpha 1.$$
(2.20)
Inequalities (2.9), (2.18) and (2.20) are crucial for the inductive estimates.
Cluster of the origin. From the above estimates, and recalling (2.7), the initial heuristic discussion is indeed made rigorous: for the above choice of parameters and picking $`l_1`$ large enough we (recursively) obtain that for all $`M1`$ and at all scales $`k=1,\mathrm{},M`$,
$$\nu (I_j^{(k)}\text{ is defected })l_k^\delta .$$
(2.21)
Indeed, due to (2.7), we see that the previous analysis and the above choice of the parameters turns rigorous the discussion leading to (1.7) and (1.9). Now, for $`k=M`$, we have at most two $`M`$–blocks, denoted by $`I_i^{(M)}`$, where $`1is`$ and $`s(\omega )\{1,2\}`$. In particular, from (2.21), we immediately have the basic estimate (1.4) announced in the introduction. Next we give the uniform lower bound for
$$\nu \left(0y,\text{for some}y[L2l_M^{\alpha ^{}/\alpha },L]\right).$$
Recalling that $`j_z^k=\mathrm{min}\{j:zI_j^{(k1)}\}`$, for any $`1kM`$ we define the following events:
$$\psi ^{(k)}=\underset{i=j_0^kl_k^{\alpha ^{}/\alpha }/l_{k1}}{\overset{j_0^k+l_k^{\alpha ^{}/\alpha }/l_{k1}}{}}G(I_i^{(k1)})$$
(2.22)
and consider
$$\mathrm{\Psi }_M=\underset{j=1}{\overset{M}{}}\psi ^{(j)}.$$
(2.23)
The occurrence of $`_{k=1}^n\psi ^{(k)}`$, $`1nM`$ implies that the origin $`0`$ is the right (resp. left) end–vertex of a $`(n1)`$–block $`I_{j_0^n}^{(n1)}`$ (resp. $`I_{j_0^n+1}^{(n1)}`$) for each $`n`$, since no adjustments are performed in this case, and necessarily it belongs to the pedestals $`\mathrm{{\rm Y}}(I_{j_0^n}^{(n1)})`$ and $`\mathrm{{\rm Y}}(I_{j_0^n+1}^{(n1)})`$ for any $`\mathrm{\hspace{0.33em}1}nM`$. In particular, for $`\omega \mathrm{\Psi }_M`$ we have $`s(\omega )=2`$. Moreover, in the event $`\mathrm{\Psi }_MG(I_1^{(M)})G(I_2^{(M)})`$, the origin 0 belongs to an open oriented path connecting $`[L,L+2l_M^{\alpha ^{}/\alpha }]`$ to $`[L2l_M^{\alpha ^{}/\alpha },L]`$ as described above.
Taking into account the estimate (2.21) and the definition (2.22) we have for $`k2`$:
$$\nu (\psi ^{(k)})1(2(l_{k1})^{\alpha ^{}1}+1)(l_{k1})^\delta 13(l_{k1})^{\alpha ^{}1}(l_{k1})^\delta .$$
Since
$$\delta (\alpha ^{}1)>\delta (\alpha 1)\frac{\alpha (\alpha 1)}{2\alpha }>0,$$
we define $`u=\delta (\alpha ^{}1)>0`$ and rewrite the above inequality:
$$\nu (\psi ^{(k)})13(l_{k1})^u\text{for}k2.$$
Since $`l_k`$ grow super-exponentially fast, we get immediately that the series
$$(l_1)^u+(l_2)^u+(l_3)^u+\mathrm{}=S(l_1)$$
converges and
$$S(l_1)0,\text{when}l_1\mathrm{}.$$
This immediately implies that
$$\nu \left(\left[\underset{j=2}{\overset{M}{}}\psi ^{(j)}G(I_1^M)G(I_2^M)\right]^c\right)$$
(2.24)
can be made arbitrarily small, uniformly in $`M`$.
Finally, by choosing $`l_1`$ large enough, and then $`p`$ close enough to $`1`$ we get that $`\nu (\psi ^{(1)})`$ can be made arbitrarily close to $`1`$.
The proof of Theorem 1.1 follows at once. $`\mathrm{}`$
Acknowledgments The authors thank C. M. Newman for suggesting this problem and A.-S. Sznitman for many useful discussions. We thank E. Presutti for pointing out the “interface question” and M. Cassandro and I. Merola for discussions on this topic. We thank the referees of a previous version, for their careful reading and for their suggestions. This work was supported by Faperj grant E-26/151.905/2000, Fapesp grant 03/01366-7, CNPq grant 477259/01-4 and Fundacion Andes. VS is partially supported by CNPq grant 302221/2008-5. MEV is partially supported by CNPq grant 302796/2002-9. We also thank various institutions CBPF, IMPA, IF-USP, ETH, UC-Berkeley and MSRI for financial support and hospitality during periods when this work was carried out.
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# Introduction
## Introduction
Yang-Mills theory in two dimensions is at the intersection of many different fields of theoretical physics. It is one example of non trivial completely solvable gauge theory , in which both perturbative and non perturbative effects can be studied. Its large $`N`$ expansion has been proved to describe a two dimensional string theory , namely a theory of branched coverings on a two dimensional Riemann surface. Non trivial topological sectors in the unitary gauge also seem to be related to matrix string states . Its partition function on a torus can be described in terms of a gas of free fermions , and the kernel on a cylinder by the evolution of a system of $`N`$ free fermions on a circle, namely by the Sutherland model.
Some new connections between two dimensional gauge theories and statistical mechanical systems were pointed out in where two dimensional gauge theories of the symmetric group $`S_n`$ in the large $`n`$ limit were investigated. Gauge theories of $`S_n`$ also describe $`n`$-coverings of a Riemann surface and hence they are closely related to two dimensional Yang-Mills theories; the relation being essentially provided by Frobenius formula that relates U($`N`$) characters (in a representation with $`n`$ boxes in the Young diagram) to the corresponding characters of the $`S_n`$ group.
It was shown in that the partition function of a gauge theory on a disc or a cylinder can be interpreted in terms of random walks on the gauge group, whose initial and final positions are the holonomies at the ends of the cylinder<sup>2</sup><sup>2</sup>2In a disc the starting point is the identity of the group, and on a sphere both starting and ending points are the identity. and the number of steps is the area of the surface. Although the focus in was on the discrete group $`S_n`$ the argument can be trivially extended, as it is shown in the Appendix of the present paper, to continous Lie groups. Similar results were independently obtained in . It is well known in random walks theory that after a certain number of steps the end point of the walk becomes independent of the starting point: the walker has lost any memory of the point he started from. The critical number of steps after which that happens can be exactly calculated in a number of situations, and the corresponding transition is known as cutoff transition. Given the correspondence between random walks on the group and gauge theory on a cylinder, one expects to find the cutoff transition also in gauge theories. Indeed it was found in that for an $`S_n`$ gauge theory where the holonomy on each elementary plaquette is given by a single transposition<sup>3</sup><sup>3</sup>3In terms of the string interpretation this means that in each plaquette there is a single quadratic branch point connecting two of the $`n`$ sheets of the world sheet., a cutoff transition occurs in the large $`n`$ limit when the number of plaquettes (and hence the area) is $`\frac{1}{2}n\mathrm{log}n`$, in agreement with previous results in random walks . Models with more general Boltzmann weight for the plaquettes have a richer structure of phase diagrams . The stringy interpretation of the cutoff transition in the $`S_n`$ gauge theory is the following: beyond the transition the string world sheet is completely connected in the large $`n`$ limit, while before the transition the world sheet consists of a large connected part and of a small fraction (in fact vanishing in the large $`n`$ limit) of disconnected parts.
Another well known correspondence relates random walks with random graphs , that is graphs obtained by randomly connecting $`n`$ points with $`p`$ links. These can be put in correspondence with random walks on $`S_n`$ made of $`p`$ steps, each step consisting of a simple transposition. Two types of transitions are known in the large $`n`$ large $`p`$ double scaling limit in random graphs: a percolation transition at a critical value $`\beta =\beta _c`$ when $`p=\beta n`$ and the cutoff transition at $`\alpha =1/2`$ when $`p=\alpha n\mathrm{log}n`$. Beyond the transition, namely for $`p>1/2n\mathrm{log}n`$ all the $`n`$ points are connected whereas before the transition a vanishing fraction of disconnected points survive.
It is rather natural at this point to look for a similar transition in the large $`N`$ limit of U($`N`$) gauge theories. A large $`N`$ phase transition on a sphere and on a cylinder in two dimensional Yang-Mills theories - the Douglas-Kazakov phase transition - has been known for quite some time. However this is not a cutoff transition. In a cutoff transition the partition function on a disc for instance becomes independent on the holonomy on the border of the disc, and this is not the case in the Douglas-Kazakov transition. Besides the Douglas-Kazakov transition occurs at a finite value of the ’t Hooft coupling whereas from the previous examples it appears that the cutoff transition occurs when the area scales as $`\mathrm{log}N`$ at large $`N`$. From this point of view the Douglas-Kazakov transition appears more similar to the percolation transition in random graphs, although a precise correspondence is still to be found.
The existence of the cutoff transition in the large $`N`$ limit of $`2`$D Yang-Mills on the sphere when the ’t Hooft coupling scales as $`\mathrm{log}N`$ was proved by a simple argument in . The present paper is devoted to study such transition further on both the sphere and the cylinder, in order to characterize its phases and give some physical interpretation. The paper is organized as follows: in Section 1 we review the large $`N`$ Douglas-Kazakov transition and its physical interpretation. In Section 2 we introduce the cutoff transition on the sphere and study the phases above and below the transition. Section 3 is devoted to the transition on the disc and on the cylinder and Section 4 to the calculation of the expectation value of Wilson loops.
## 1 Large $`N`$ transition
The partition function of a pure gauge theory on an arbitrary orientable two-dimensional manifold $``$ of genus $`G`$, $`p`$ boundaries and area $`\stackrel{~}{A}`$ has been known for many years :
$`Z_{}`$ $`=`$ $`{\displaystyle [𝒟A^\mu ]e^{\frac{1}{4\lambda ^2}_{}d^2x\sqrt{g}\mathrm{Tr}(F^{\mu \nu }F_{\mu \nu })}}`$ (1)
$`=`$ $`{\displaystyle \underset{r}{}}\chi _r(g_1)\mathrm{}\chi _r(g_p)d_r^{22Gp}e^{\frac{A}{2N}C_2(r)}.`$
The sum runs over all irreducible representations of the gauge group, $`\lambda `$ is the gauge coupling, $`\chi _r(g_i)`$ is the character of the holonomy $`g_i`$ in the representation $`r`$ and $`C_2(r)`$ is the quadratic Casimir operator in the representation $`r`$. $`A`$ is related to the actual area of $``$ through $`A=\lambda ^2N\stackrel{~}{A}`$.
We consider $`G=0`$ manifolds with at most 2 boundaries; i.e. spheres, discs and cylinders. Moreover, we will confine our analysis to the unitary groups $`U(N)`$ and $`SU(N)`$. A third order phase transition in the large $`N`$ limit was discovered in the case of a sphere by Douglas and Kazakov at a critical value $`A=\pi ^2`$ of the rescaled area $`A`$. This transition appears to separate a weak coupling ($`A<\pi ^2`$) from a strong coupling ($`A>\pi ^2`$) regime. These results were generalized to the case of a cylinder in where the phase transition was also interpreted as a result of instanton condensation.
The partition function on the sphere can be written as a sum over the set of integers $`n_1>n_2>\mathrm{}>n_N`$ that label the irreducible representations of SU($`N`$) and U($`N`$)<sup>4</sup><sup>4</sup>4In the case of U($`N`$) the extra condition $`n_N\frac{N1}{2}`$ must be imposed.:
$$Z=e^{\frac{A}{24}(N^21)}\underset{n_1>n_2\mathrm{}>n_N}{}\underset{i<j}{}(n_in_j)^2e^{\frac{A}{2N}_{i=1}^N(n_i)^2},$$
(2)
The existence of a phase transition at the critical value $`A=\pi ^2`$ can be easily derived from (2) by noticing that the partition function (2) is exactly the same as the one of a gaussian hermitian matrix model but with the integral over the eigenvalues replaced by the discrete sum over the integers $`n_i`$. The solution of a gaussian hermitian matrix model in the large $`N`$ limit is given by Wigner’s semicircle distribution law for the eigenvalues:
$$|\rho (\lambda )|=\frac{A}{2\pi }\sqrt{\frac{4}{A}\lambda ^2}.$$
(3)
where the continuum variables
$$x=\frac{i}{N},\lambda (x)=\frac{n_i}{N}.$$
and the corresponding density of eigenvalues
$$\rho (\lambda )=\frac{x}{\lambda }$$
(4)
have been introduced. In (2) however the would-be eigenvalues $`n_i`$ are distinct integers and as a consequence the corresponding density $`\rho (\lambda )`$ in the large $`N`$ limit is constrained by:
$$|\rho (\lambda )|1\lambda .$$
(5)
Hence the Wigner semicircle solution is acceptable only in the weak coupling phase, namely for $`A\pi ^2`$ where the condition (5) is fulfilled. In fact the maximum of $`|\rho (\lambda )|`$ occurs at $`\lambda =0`$, it increases with the area $`A`$ and becomes equal to $`1`$ at $`A=\pi ^2`$, as easily seen from (3). The solution in the strong coupling phase $`A>\pi ^2`$ was found in and is expressed in terms of elliptic integrals. In this phase a finite fraction of the eigenvalues condenses, namely the distribution $`\rho (\lambda )`$ is flat and equal to one in a symmetric interval around $`\lambda =0`$ as shown in fig. 2.
### 1.1 Configuration space
It is well known that the partition function of two dimensional Yang-Mills theories with gauge group U($`N`$) can be interpreted in terms of a gas of $`N`$ free fermions on a circle described by a Sutherland-Calogero model. In particular, if we denote by $`𝒦_2(\theta ,\varphi ;A)`$ the kernel on a cylinder of scaled area $`A`$ and with the U($`N`$) holonomies at the two ends given by the invariant angles $`\theta _i`$ and $`\varphi _i`$, it was shown that $`𝒦_2(\theta ,\varphi ;A)`$ can be interpreted as the propagator in a time $`A`$ from an initial configuration where the positions of the $`N`$ fermions on the circle are given by $`\theta _i`$ and a final configuration with positions labeled by $`\varphi _i`$. The partition function on the sphere of area $`A`$ is a particular case where the initial and final configurations are just $`\theta _i=\varphi _i=0`$ for all $`i`$’s, namely the amplitude for a process where all fermions start at the origin and come back to the origin after a time $`A`$. By a modular transformation on the kernel of the cylinder one finds that the integers $`n_i`$ labeling the irreducible representations of U($`N`$) are just the discrete momenta of the fermions on the circle. While in the momentum representation the Douglas-Kazakov phase transition can be interpreted as fermion condensation, in the configuration representation it can be seen in terms of instantons condensation . In fact, while going from the initial $`\theta _i=0`$ configuration to the final $`\varphi _i=0`$ configuration, a fermion can in principle wind an arbitrary number of times around the circle. These winding (instantons) configurations do not contribute in the large $`N`$ limit to the weak coupling phase, as shown by the following simple argument. Consider the Wigner distribution (3) of momenta in the weak coupling phase. The maximum allowed momentum is $`n_{\mathrm{max}}=\frac{2}{\sqrt{A}}`$, hence the maximum shift in position for a single fermion in the time $`A`$ is given by
$$\mathrm{\Delta }\theta _{\mathrm{max}}=A\frac{2}{\sqrt{A}}=2\sqrt{A}$$
(6)
The existence of winding trajectories requires this shift in position to be at least $`2\pi `$, namely $`A`$ to be greater of $`\pi ^2`$. Hence the critical value of $`A`$, where the Douglas-Kazakov phase transition occurs, marks the point where instantons condense, and contribute to the functional integral in the large $`N`$ limit.
A more detailed understanding of the Douglas-Kazakov phase transition can be achieved by introducing in the large $`N`$ limit the density $`\rho (\theta ,t)`$ of fermions in the position $`\theta `$ at a given time (=area) $`t`$. Due to the compact nature of the configuration space $`\rho (\theta ,t)`$ is defined in the interval $`\pi \theta \pi `$ with $`\rho (\pi ,t)=\rho (\pi ,t)`$. Matytsin proved that if the evolution equation of the fermions is given by the Calogero-Sutherland model then the density $`\rho (\theta ,t)`$ is governed by the Das-Jevicki equation which admits Wigner semicircular distribution of radius $`r(t)`$ as a solution:
$$\rho (\theta ,t)=\frac{2}{\pi r(t)^2}\sqrt{r(t)^2\theta ^2},|\theta |r(t)$$
(7)
with $`r(t)`$ satisfying the differential equation $`\frac{d^2r(t)}{dt^2}+\frac{4}{r(t)^3}=0`$. On a sphere of area $`A`$ the boundary conditions are $`r(0)=r(A)=0`$ and the differential equation has the solution:
$$r(t)=2\sqrt{\frac{t(At)}{A}}$$
(8)
The solution given by (7) and (8) is valid provided the support of the density function $`\rho (\theta ,t)`$ is in the interval $`[\pi ,\pi ]`$ at all $`t`$, namely provided $`r(t)\pi `$. The maximum value for the radius $`r(t)`$ occurs for $`t=A/2`$ and is $`r_{max}=\sqrt{A}`$. Hence the condition for the validity of the Wigner semicircular solution is $`A\pi ^2`$. Beyond the critical value $`A=\pi ^2`$ the fermions ”realize” that the space they live in is compact, instantons effects become important and (7) is not an acceptable solution any longer. On the cylinder a similar phase transition occurs in general, at a critical value of the area that depends on the holonomies at the boundaries<sup>5</sup><sup>5</sup>5However with particular conditions at the boundary the phase transition may also be absent, see for instance .. If the distribution of the invariant angles of the holonomies at the boundaries is the Wigner semicircular distribution, then the critical area can be calculated exactly in the same way as for the sphere . For a more general discussion see .
## 2 Large $`A(N)`$ phase transition
It is apparent from the discussion in the previous section and from the explicit form (2) of the partition function that as the area of the sphere increases the distribution of the ”momenta” $`n_i`$ (i.e. of the integers that label the U($`N`$) representations) becomes more and more similar to a double step function, like the one drawn with dashed lines in fig. 3. In fact for very large areas the attractive quadratic potential tends to dominate over the repulsive force produced by the Vandermonde determinant. The double step distribution corresponds to the trivial representation of U($`N`$) in which all characters are identical irrespective of their argument. If this distribution dominates the functional integral then the kernel on a disc or on a cylinder becomes independent from the holonomies at the boundaries.
As already mentioned in the introduction and discussed in the Appendix, the partition functions on the disc and on the cylinder may be interpreted in terms of random walks on the group manifold. From this point of view the very large area phase, where the sum over the irreducible representations is dominated by the trivial representation corresponds to a walk which is so long that the walker has lost any notion of the starting point. The transition where this situation sets in is known in random walk theory as ”cutoff transition”, and the same term will be used here.
The cutoff regime occurs for areas larger than a critical $`N`$ dependent value $`A_c(N)`$ which was found in . The argument is very simple: consider the trivial double step representation $`R_0`$ that minimizes the Casimir term $`_in_i^2`$
$$R_0:\{n_1,\mathrm{},n_N\}=\{\frac{N1}{2},\mathrm{},\frac{N1}{2}\}.$$
(9)
and compare its contribution to the partition function (2) with the one coming from a representation $`R_1`$ in which $`n_1`$ has been increased by $`1`$, namely in which $`n_1=\frac{N1}{2}+1`$. In other words we look for the value of $`A`$ at which $`R_0`$ ceases to be dominant. A simple calculation shows that the ratio between the two contributions is
$$\frac{Z_0(A,N,R_0)}{Z_0(A,N,R_1)}=\frac{e^{\frac{A}{2}}}{N^2}.$$
(10)
This ratio is larger than $`1`$ (hence $`R_0`$ dominates and we are in the cutoff phase) if $`A>4\mathrm{log}N`$.
In order to study this phase transition in more detail, it is convenient to parametrize the area $`A`$ by rescaling it with $`\mathrm{log}N`$:
$$A=\alpha \mathrm{log}N+\beta $$
(11)
From the previous argument we expect the cutoff transition to occur at the critical value $`\alpha _c=4`$, separating two distinct phases . So Yang Mills theory on a sphere seems to have four phases altogether: the first two at $`\alpha =0`$ separated by the Douglas-Kazakov phase transition at $`\beta =\pi ^2`$, the other two when the area is logarithmically rescaled with $`N`$ that are separated by the cutoff transition.
The aim of this section is to find the saddle point configuration and the free energy in the two phases above and below the cutoff point $`\alpha _c=4`$: the cutoff phase $`\alpha >\alpha _c`$ has been discussed above and it is rather trivial, but the phase below $`\alpha _c`$ appears as an interesting intermediate phase between the strong coupling phase in the Douglas-Kazakov transition and the cutoff phase. Hence we shall concentrate on this in the rest of the section.
Let us consider again the partition function (2) and the corresponding action
$$S=2\underset{i>j=1}{\overset{N}{}}\mathrm{log}|n_in_j|\frac{A}{2N}\underset{i=1}{\overset{N}{}}n_i^2$$
(12)
We want to find the extremum of this action in the large $`N`$ limit, when $`A`$ is parametrized as in eq. (11). Since we expect the saddle point distribution of the ”momenta” $`n_i`$ to be symmetric with respect to the origin ( $`n_in_i`$ ) we shall perform the variation only with respect to symmetric configurations, that is we set
$$MiM,\text{with}M=\frac{N1}{2}\text{and}n_i=n_i,$$
(13)
Using this symmetry one can restrict the sums to non negative values of $`i`$ ($`i=1,\mathrm{},M`$) and write the action as:
$$S=2\underset{i>j1}{\overset{M}{}}\mathrm{log}(n_in_j)^2+\underset{i>j1}{\overset{M}{}}\mathrm{log}(n_i+n_j)^2\frac{A}{2M}\underset{i=1}{\overset{M}{}}n_i^2$$
(14)
While in the cutoff phase $`n_i=i`$ for $`i=1,\mathrm{},M`$ below the cutoff transition we expect a configuration of the type described in fig. 3, namely:
$`n_i=i`$ $`i=1,\mathrm{},Ml`$
$`n_{Ml+\alpha }=Ml+\alpha +r_\alpha `$ $`\alpha =1,\mathrm{},l`$ (15)
where the value of $`l`$ and the spectrum of the integers $`r_\alpha `$ are to be determined. With these notations the action can be written as:
$`SS_0`$ $`=`$ $`4{\displaystyle \underset{\alpha =1}{\overset{l}{}}}{\displaystyle \underset{j=1}{\overset{Ml}{}}}\mathrm{log}\left[\left(1+{\displaystyle \frac{r_\alpha }{\alpha +j}}\right)\left(1+{\displaystyle \frac{r_\alpha }{Ml+\alpha +j}}\right)\right]`$ (16)
$`+`$ $`2{\displaystyle \underset{\alpha \beta }{}}\left(\mathrm{log}\left(1+{\displaystyle \frac{r_\alpha +r_\beta }{\alpha +\beta +2(Ml)}}\right)+\mathrm{log}\left(1+{\displaystyle \frac{r_\alpha r_\beta }{\alpha \beta }}\right)\right)`$
$``$ $`{\displaystyle \frac{A}{2M}}{\displaystyle \underset{\alpha =1}{\overset{l}{}}}(2(Ml+\alpha )r_\alpha +r_\alpha ^2),`$
where $`S_0`$ represents the value of the action in the trivial representation $`R_0`$. We shall assume that as $`M\mathrm{}`$ also $`l\mathrm{}`$ but at a slower rate than $`M`$, namely $`l/M0`$. We shall also assume that $`l`$ and $`r_\alpha `$ will be of the same order in the large $`M`$ limit. These assumptions will be justified a posteriori, in the sense that they will provide a stable saddle point in the large $`M`$ limit when the area is scaled like $`\mathrm{log}M`$. They are also very reasonable assumptions: $`l`$ is of order $`N`$ in the strong coupling phase following the Douglas-Kazakov transition and one expects that with the logarithmic rescaling of the area it will shrink further by some power of $`N`$.
The first step in dealing with (16) is to make the dependence from $`M`$ explicit. By using the identity
$$\underset{j=1}{\overset{Ml}{}}(1+\frac{r_\alpha }{j+\alpha })(1+\frac{r_\alpha }{Ml+j+\alpha })=\frac{(2M2l+\alpha +r_\alpha )!}{(2M2l+\alpha )!}\frac{\alpha !}{(\alpha +r_\alpha )!}$$
(17)
we can rewrite the action as
$`SS_0`$ $`=`$ $`4{\displaystyle \underset{\alpha =1}{\overset{l}{}}}\mathrm{log}{\displaystyle \frac{(2M2l+\alpha +r_\alpha )!}{(2M2l+\alpha )!}}{\displaystyle \frac{\alpha !}{(\alpha +r_\alpha )!}}`$ (18)
$`+`$ $`2{\displaystyle \underset{\alpha \beta }{}}(\mathrm{log}(1+{\displaystyle \frac{r_\alpha +r_\beta }{\alpha +\beta +2(Ml)}})+\mathrm{log}(1+{\displaystyle \frac{r_\alpha r_\beta }{\alpha \beta }}))`$
$``$ $`{\displaystyle \frac{A}{2M}}{\displaystyle \underset{\alpha =1}{\overset{l}{}}}(2(Ml+\alpha )r_\alpha +r_\alpha ^2),`$
All ratios of factorials in (18) can be reduced to the form $`\mathrm{log}\frac{(N+C)!}{N!}`$ with $`N\mathrm{}`$, $`C\mathrm{}`$ and $`{\displaystyle \frac{C}{N}}0`$. By repeated use of Stirling formula one finds, up to terms that vanish as $`{\displaystyle \frac{C}{N}}0`$:
$$\mathrm{log}\frac{(N+C)!}{N!}C\left[\mathrm{log}N+f(\frac{C}{N})\right]$$
(19)
where
$$f(\frac{C}{N})=\mathrm{log}(1+\frac{C}{N})+\frac{\mathrm{log}(1+\frac{C}{N})}{\frac{C}{N}}1=\underset{k=1}{\overset{Ml}{}}(1)^{k1}\frac{z^k}{k(k+1)}.$$
(20)
By using this asymptotic behaviour, and introducing continuum variables in the large $`N`$ limit, namely
$$x=\frac{\alpha }{l}r(x)=\frac{r_\alpha }{l}\underset{\alpha }{}=l𝑑x.$$
(21)
the action finally takes the form
$`SS_0`$ $`=`$ $`4l^2{\displaystyle _0^1}dx[r(x)(\mathrm{log}({\displaystyle \frac{M}{l}})\mathrm{log}(x+r(x))+1+\mathrm{log}2)+x(\mathrm{log}x`$ (22)
$``$ $`\mathrm{log}(x+r(x)))]+2l^2{\displaystyle }_0^1dx{\displaystyle }_0^1dy\mathrm{log}(1+{\displaystyle \frac{r(x)r(y)}{xy}})`$
$``$ $`Al^2{\displaystyle _0^1}𝑑xr(x)`$
where subleading terms (by powers of $`\frac{l}{M}`$ ) have been neglected.
Let us now parametrize the area $`A`$ according to eq. (11) and write the action as:
$$SS_0=4l^2[F_0\mathrm{log}MF_1\mathrm{log}l+F_2]$$
(23)
where $`F_0`$, $`F_1`$ and $`F_2`$ are of order $`1`$ in the large $`M`$ and $`l`$ limit and are given by:
$`F_0`$ $`=`$ $`(1{\displaystyle \frac{\alpha }{4}}){\displaystyle _0^1}𝑑xr(x)`$
$`F_1`$ $`=`$ $`{\displaystyle _0^1}𝑑xr(x)`$
$`F_2`$ $`=`$ $`{\displaystyle _0^1}𝑑x\left[r(x)(\mathrm{log}(x+r(x))+1\mathrm{log}2{\displaystyle \frac{\beta }{4}})+x(\mathrm{log}x\mathrm{log}(x+r(x)))\right]`$
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^1}𝑑x{\displaystyle _0^1}𝑑y\mathrm{log}(1+{\displaystyle \frac{r(x)r(y)}{xy}})`$
In order to find the configuration that maximizes the functional integral in the large $`M`$ limit we take the variation of (23) with respect to both $`l`$ and $`r(x)`$. The variation with respect to $`l`$ gives the equation:
$$(1\frac{\alpha }{4})\mathrm{log}M+(\frac{F_2}{F_1}\frac{1}{2})=\mathrm{log}l$$
(24)
which shows that $`l`$ grows like $`M^{1\frac{\alpha }{4}}`$. This is consistent with what we expected: for $`\alpha 0`$ it gives $`lM`$ as in the strong phase beyond the Douglas-Kazakov transition, and at the cutoff point $`\alpha =4`$ the power vanishes as expected.
The variation with respect to $`r(x)`$ gives on the other hand
$$_0^1𝑑y\frac{1}{x+r(x)yr(y)}\mathrm{log}(x+r(x))=C$$
(25)
with
$$C=\mathrm{log}l(1\frac{\alpha }{4})\mathrm{log}M+\frac{\beta }{4}+\mathrm{log}2=\frac{F_2}{F_1}\frac{1}{2}+\frac{\beta }{4}+\mathrm{log}2$$
(26)
If one introduces the new variable $`\xi =x+r(x)`$ and the density function $`\rho (\xi )=\frac{dx}{d\xi }`$ with support in the interval $`[0,a]`$ equation (25) becomes<sup>6</sup><sup>6</sup>6The lowest extreme of the interval is $`r(0)`$ which is zero by construction, the upper end $`a`$ is, according with the definition of $`\xi `$, $`a=1+r(1)`$.:
$$_0^a𝑑\eta \frac{\rho (\eta )}{\xi \eta }\mathrm{log}\xi =C.$$
(27)
This is a standard type of equation for the density of eigenvalues in the large $`N`$ limit of matrix models and can be solved by standard analytic methods (for a detailed discussion of this equation see for instance ). The resolvent function, whose discontinuity across the cut gives the density $`\rho (\xi )`$, is given by:
$$H(\xi )=\mathrm{log}\xi +C2\mathrm{log}\left(\frac{\sqrt{\xi a}+\sqrt{\xi }}{\sqrt{a}}\right)$$
(28)
with the additional condition that for large $`\xi `$
$$H(\xi )=\frac{1}{\xi }+𝒪\left(\frac{1}{\xi ^2}\right)$$
(29)
The corresponding density is given by
$$\rho (\xi )=\frac{2}{\pi }\mathrm{arccos}(\sqrt{\frac{\xi }{a}}).$$
(30)
This solution obviously describes, through the symmetry (13), both the positive and the negative region of $`n_i`$.
The asymptotic condition (29) gives the two extra equations
$`a`$ $`=`$ $`2,`$
$`C`$ $`=`$ $`\mathrm{log}2.`$ (31)
Eq. (30), together with the first of (31) define $`r(x)`$ completely, although in an implicit way. Hence all integrals involved in the definition of $`F_1`$ and $`F_2`$ can be calculated. The calculation can actually be done analitically and gives:
$$F_1=\frac{1}{4},\frac{F_2}{F_1}=\frac{1}{2}\frac{\beta }{4}$$
(32)
The second equation could have been derived independently from (26) and (31), so it constitutes a non trivial consistency check. We can now write explicitely $`l`$ and the free energy $`F`$ in terms of the area $`A(M)`$:
$`l`$ $`=`$ $`e^{\frac{\beta }{4}}M^{1\frac{\alpha }{4}}=Me^{\frac{A(M)}{4}}`$
$`F`$ $`=`$ $`SS_0={\displaystyle \frac{1}{2}}l^2={\displaystyle \frac{1}{2}}e^{\frac{\beta }{2}}M^{2\frac{\alpha }{2}}={\displaystyle \frac{1}{2}}M^2e^{\frac{A(M)}{2}}`$ (33)
Beyond the cutoff transition we have instead $`F=l=0`$, as the dominant eigenvalue distribution is given by the double step function sketched in fig. 3.
The interpretation of these results from the point of view of the free fermion description is very clear: beyond the cutoff ($`\alpha >4`$) we are effectively in a zero temperature situation where all fermions fill the Fermi sea with no holes. Below the cutoff instead ($`\alpha <4`$) some excited fermions and the corresponding holes are present in proximity of the surface of the Fermi sea both on the positive and negative momentum side. The number of fermions above the sea level is given by $`l`$ in (33). The ratio $`\frac{l}{M}`$ vanishes like $`M^{\frac{\alpha }{4}}`$ in the large $`M`$ limit. This distinguishes this phase from the strong coupling phase of the Douglas-Kazakov transition, where such ratio remains finite , namely the number of fermions above the Fermi sea level is of order $`M`$.
In spite of being described in terms of $`N`$ free fermions, the free energy is proportional (with the standard ’t Hooft scaling) to $`N^2`$, which reflects the original number of degrees of freedom in a unitary $`N\times N`$ matrix model<sup>7</sup><sup>7</sup>7The description in the terms of $`N`$ fermions follows the integration over the angular variables that reduces the matrix model to an integral over the eigenvalues. The ensueing Vandermonde determinant makes the wave function describing the eigenvalues antisymmetric. As a consequence the total momentum of $`M`$ left moving fermions is of order $`M^2`$ rather than $`M`$, that is of the same order as the original number of bosonic degrees of freedom.. So it is not surprising that the free energy becomes proportional to $`l^2`$ in presence of $`l`$ effective fermionic degrees of freedom when the area is rescaled by a $`\mathrm{log}M`$ factor. It is as if the effective size of the original matrix had shrunk to $`l\times l`$. A full understanding of the reduction of number of degres of freedom from the point of view of the original gauge degrees of freedom is still wanted, although some light on it might be thrown by the study of the kernels on the disc and the cylinder in the following sections.
It is apparent from (33) that the number of effective degrees of freedom $`M^{2\frac{\alpha }{2}}`$ is the actual order parameter for the cutoff transition. Incidentally, this unusual dependence of the number of degrees of freedom upon $`\alpha `$, together with the choice of a large $`M`$ limit, makes it rather delicate to classify such a transition according to standard terminology.
Let us finally consider the representations of U($`N`$) and/or SU($`N`$) that correspond to the ”momentum” distribution pictured in fig. 3 and given in (30). With the group SU($`N`$)<sup>8</sup><sup>8</sup>8With SU($`N`$) the term $`_in_i^2`$ in the action should be replaced by $`_in_i^2\frac{1}{N}(_in_i)^2`$ which is invariant under $`n_in_i+a`$. However with suitable choice of $`a`$ the extra term can be set equal to zero. this corresponds to a composite representation in the sense of ref. , whose Young diagram is shown in fig. 4 where the constituent rpresentations are denoted by $`R`$ and $`S`$. The rows in the Young diagram of $`R`$ and $`S`$ (that coincide in this case) have lengths $`r_\alpha `$ and their total number of boxes $`|r|`$ is given in the large $`N`$ limit by
$$|r|=\underset{\alpha }{}r_\alpha l^2𝑑xr(x)=\frac{1}{4}l^2$$
(34)
If the group is U($`N`$) an arbitrary number of columns of length $`N`$ can be added or subtracted to the Young diagram of fig. 4, and the composite representation can be seen as the direct product of the two constituent representations $`R`$ and $`S`$ with opposite U($`1`$) charges. In fact the integers labeling the representation $`S`$ are in this case negative, which corresponds to changing signs of the invariant angles $`\theta _i`$, namely to changing $`U`$ into $`U^{}`$. In the strong couplig regime of the Douglas-Kazakov transition ($`A>\pi ^2`$) the partition function is dominated in the large $`N`$ limit by a composite representation of the same type but with the Young diagram of $`R`$ and $`S`$ made of rows and columns of lengths of order $`N`$, rather than $`l`$, and a total number of boxes of order $`N^2`$ rather than $`l^2`$.
## 3 Cutoff transition on the disc and on the cylinder
In this section we are going to consider the partition function on a disc and on a cylinder, with fixed holonomies at the boundaries, in the regime where the area $`A`$ is scaled as in (11). We shall show that the cutoff transition occurs also in this case at the same critical value of $`\alpha `$, except for some particular holonomies at the boundaries. This is in analogy to what happens in the case of the Douglas-Kazakov transition which was proved to occur on a cylinder in , except possibly for some special configurations (see ) .
We shall use the standard expression for the partition function on a cylinder with holonomies $`U`$ and $`V`$ at the boundaries, which is a particular case of (1):
$$Z_{cyl}(U,V)=e^{\frac{A}{24}(N^21)}\underset{r}{}\chi _r(U)\chi _r(V)e^{\frac{A}{2N}_in_i^2}.$$
(35)
The partition function on a disc is obtained from (35) by taking for instance $`U1`$ and the partition function on the sphere is recovered by taking the double limit $`U1`$ and $`V1`$. The area $`A`$ in (35) scales as in (11). The sum over the representations $`r`$ is replaced in the large $`N`$ limit by the saddle point, namely by a representation whose ”momentum” distribution is of the type described in fig. 3. This corresponds to a composite representation, whose constituent representations $`R`$ and $`S`$ have rows and columns of order $`l`$ with $`l/N0`$ in the large $`N`$ limit. The first step in evaluating (35) is then to calculate the characters in such limit.
### 3.1 Characters in the large $`N`$, large $`l`$ limit
For simplicity, let us consider first the case where only one constituent representation is present, which would be the case if only right (or left) moving fermions were present above the Fermi sea. This is given by:
$`n_i`$ $`=`$ $`i\mathrm{for}\mathrm{i}=1,\mathrm{},\mathrm{N}\mathrm{l}`$
$`n_i`$ $`=`$ $`Nl+\alpha +r_\alpha \mathrm{for}i=Nl+\alpha \mathrm{and}\alpha =1,\mathrm{},l`$ (36)
The dimension $`\mathrm{\Delta }_r`$ of this representation can be written in the form
$`\mathrm{\Delta }_r`$ $`=`$ $`{\displaystyle \underset{i>j}{}}{\displaystyle \frac{|(n_in_j)|}{(ij)}}={\displaystyle \underset{\alpha }{}}{\displaystyle \frac{\alpha !}{(\alpha +r_\alpha )!}}{\displaystyle \underset{\alpha >\beta }{}}\left(1+{\displaystyle \frac{r_\alpha r_\beta }{\alpha \beta }}\right){\displaystyle \underset{\alpha }{}}{\displaystyle \frac{(Nl+\alpha +r_\alpha )!}{(Nl+\alpha )!}}`$ (37)
$`=`$ $`{\displaystyle \frac{d_r}{|r|!}}{\displaystyle \underset{\alpha }{}}{\displaystyle \frac{(Nl+\alpha +r_\alpha )!}{(Nl+\alpha )!}}={\displaystyle \frac{d_r}{|r|!}}N^{{\scriptscriptstyle r_\alpha }}\left[1+𝒪\left({\displaystyle \frac{l}{N}}\right)\right],`$
where $`d_r`$ is the dimension of the representation of the symmetric group $`S_{|r|}`$ associated to the Young diagram $`r`$ of rows $`r_\alpha `$. The order of the symmetric group is the total number of boxes in the Young diagram: $`|r|=_\alpha r_\alpha `$. The dependence on $`N`$ in (37) is explicit: the ratio of factorials in (37) is a polynomial in $`N`$ of degree $`|r|`$. If $`l`$, $`\alpha `$ and $`r_\alpha `$ are all of order $`l`$ the coefficient of $`N^{|r|k}`$ in this polynomial is of order $`l^k`$ in the large $`N`$, large $`l`$ limit . Hence we can write
$$\mathrm{\Delta }_r=\frac{d_r}{|r|!}N^{|r|}\left[1+\underset{k}{\overset{|r|}{}}c_k(\frac{l}{N})^k\right]$$
(38)
where the coefficients $`c_k`$ are smooth in the large $`l`$, large $`N`$ limit. If the limit is taken keeping $`l/N`$ finite, as in the strong coupling phase of the Douglas-Kazakov transition, all terms at the r.h.s. of (38) are of the same order and cannot be neglegted. On the other hand if the double scaling limit is taken with $`l/N0`$, as in the previous section, then all terms after the $`1`$ are subleading and can be neglected.
The same argument holds if instead of the dimension of the representation we consider a character of U($`N`$). In fact the celebrated Frobenius formula gives:
$$\chi _r(U)=\frac{d_r}{|r|!}\underset{\sigma S_{|r|}}{}\frac{\widehat{\chi }_r(\sigma )}{d_r}N^{k_\sigma }\underset{j}{}\left(\frac{\mathrm{Tr}U^{s_j}}{N}\right)$$
(39)
where $`\widehat{\chi }_r(\sigma )`$ is a character of the symmetric group $`S_{|r|}`$ in the representation labeled by the same Young diagram as $`\chi _r(U)`$, $`s_j`$ are the lengths of the cycles in the cycle decomposition of $`\sigma `$ ($`_js_j=|r|`$) and $`k_\sigma `$ is the number of cycles in $`\sigma `$. By taking $`U1`$ and comparing (38) and (39) one finds that the contribution of $`\frac{\widehat{\chi }_r(\sigma )}{d_r}`$ when $`\sigma `$ consists of $`k_\sigma `$ cycles is $`l^{|r|k_\sigma }`$ in the large $`l`$ limit. So if we take the double scaling limit where both $`l`$ and $`N`$ go to infinity and the ratio $`l/N`$ goes to zero, all terms in the sum over $`\sigma `$ in (39) are subleading with respect to the one where $`\sigma `$ is the identity. We obtain:
$$\chi _r(U)=\frac{d_r}{|r|!}\left(\frac{\mathrm{Tr}U}{N}\right)^{|r|}N^{|r|}[1+𝒪(l/N)]=\mathrm{\Delta }_r\left(\frac{\mathrm{Tr}U}{N}\right)^{|r|}[1+𝒪(l/N)]$$
(40)
Notice again that in the strong coupling phase of the Douglas-Kazakov transition, where in the large $`N`$ limit $`l/N`$ is kept constant, all terms in (39) coming from different permutations $`\sigma `$ are of the same order and cannot be neglected.
### 3.2 Partition function on the cylinder and special boundary conditions
As a result of previous analysis we find that in the composite representations that are dominant in the region $`0<\alpha <4`$ the characters $`\chi _r(U)`$ depend only from $`\mathrm{Tr}U`$. Hence the partition function on the cylinder will depend only on the trace of the holonomies on the boundaries. In fact it is apparent from (40) that going from the sphere ($`U`$=1) to the cylinder just amounts to a multiplicative factor $`\left(\frac{\mathrm{Tr}U}{N}\right)^{|r|}`$. In the case of interest however the composite representation contains both chiral and anti-chiral component representations (that is $`R`$ and $`S`$ of fig. 4) and not just one as in the simplified example discussed above. However in the large $`N`$ limit the two component representations are decoupled and the character of the composite representation becomes just the product of the characters of the component representations:
$$\chi _{composite}(U)=\chi _r(U)\chi _r(U^{})=\mathrm{\Delta }_r^2\left(|\frac{\mathrm{Tr}U}{N}|\right)^{2|r|}$$
(41)
By replacing (41) into (35) we find
$$Z_{cyl}(U,V)=Z_{sphere}\left(|\frac{\mathrm{Tr}U}{N}\frac{\mathrm{Tr}V}{N}|\right)^{2|r|}$$
(42)
where of course the value of $`|r|`$ is the one determined by the saddle point equations and the equality holds, in the large $`N`$ limit, only in the regime where the area $`A`$ is scaled as in (11) with $`\alpha 0`$. If we set $`u=|\frac{\mathrm{Tr}U}{N}|`$ and $`v=|\frac{\mathrm{Tr}V}{N}|`$ it is almost immediate to see that the multiplicative factor at the r.h.s. of (42) is equivalent to replace in the action (23) the constant term $`\beta `$ in the area $`A`$ with a $`\widehat{\beta }(u,v)`$ given by:
$$\beta \widehat{\beta }(u,v)=\beta \mathrm{log}u^2\mathrm{log}v^2$$
(43)
The partition function on the cylinder is then the same as the partition function on a sphere whose area is obtained from the area of the cylinder by adding the two terms $`\mathrm{log}u^2`$ and $`\mathrm{log}v^2`$. The latter can be interpreted as the areas of the two discs necessary to go, in the given momentum configuration, from $`U=1`$ (resp. $`V=1`$) to the holonomy at the boundary with $`|\frac{\mathrm{Tr}U}{N}|=u`$ (resp $`|\frac{\mathrm{Tr}V}{N}|=v`$). The areas of the two discs are of order $`1`$, so this correction does not affect the position of the cutoff transition that remains on the cylinder at $`\alpha =4`$.
The discussion above relies on the fact that the leading term in Frobenius formula (39) comes from the identical permutation. However this is not always true: if we take the large $`N`$ and $`l`$ limit keeping $`\mathrm{Tr}U=0`$ <sup>9</sup><sup>9</sup>9This can happen for instance if the holonomy at a boundary has a symmetry of some sort, for instance of the type $`\theta \theta +\pi `$, that is preserved through the limiting process. then all terms in (39) with $`\sigma `$ containing cycles of length $`1`$ would vanish. Assuming that $`\mathrm{Tr}U^20`$, the term in Frobenius formula (39) with the highest power on $`N`$ would then come from permutations $`\sigma `$ all made out of cycles of length $`2`$ and would be of order $`N^{\frac{|r|}{2}}`$ instead of $`N^{|r|}`$. Supposing that both the trace of $`U`$ and of $`V`$ vanish the coefficient of $`\mathrm{log}M`$ in the first term of (22) would be halved. Correspondingly the critical value of $`\alpha `$ at which the cutoff transition occurs would also be halved and become $`\alpha _c=2`$
Let us make this argument general and more quantitative. Let us assume that $`\mathrm{Tr}U^{k_1}0`$ with $`\mathrm{Tr}U^j=0`$ for $`j<k_1`$ and the same for $`V`$ with $`k_2`$ at the place of $`k_1`$. The leading term in (39) will now be:
$$\chi _r(U)=\frac{d_r}{|r|!}\left(\frac{\mathrm{Tr}U^{k_1}}{N}\right)^{\frac{|r|}{k_1}}\frac{\widehat{\chi }_r(\sigma )}{d_r}N^{\frac{|r|}{k_1}}$$
(44)
where $`\sigma `$ consists of $`\frac{|r|}{k_1}`$ cycles of length $`k_1`$<sup>10</sup><sup>10</sup>10We assume here for simplicity that we take $`l\mathrm{}`$ keeping the total number of boxes in the Young diagram multiple of $`k_1`$ and $`k_2`$. From the discussion following (39) we desume that in the large $`l`$ limit $`\frac{\widehat{\chi }_r(\sigma )}{d_r}l^{|r|(1\frac{1}{k_1})}`$ By using this asymptotic behaviour and eq. (44) we can write the partition function on the cylinder in the large $`N`$ and $`l`$ limit as in (23), but with $`F_0`$ and $`F_1`$ replaced by the following expressions:
$`F_0`$ $`=`$ $`({\displaystyle \frac{1}{2k_1}}+{\displaystyle \frac{1}{2k_2}}{\displaystyle \frac{\alpha }{4}}){\displaystyle _0^1}𝑑xr(x)`$
$`F_1`$ $`=`$ $`({\displaystyle \frac{1}{2k_1}}+{\displaystyle \frac{1}{2k_2}}){\displaystyle _0^1}𝑑xr(x)`$ (45)
The dependence of $`F_2`$ from $`r(x)`$ is modified, in a so far unknown way, by next to leading terms in $`\frac{\widehat{\chi }_r(\sigma )}{d_r}`$<sup>11</sup><sup>11</sup>11A lot is known on the characters of permutations with cycles all of the same length, so an explicit expression for $`F_2`$ is probably obtainable. while the constant parameter $`\beta `$ is replaced by
$$\beta \widehat{\beta }(u_{k_1},v_{k_2})=\beta \frac{1}{k_1}\mathrm{log}u_{k_1}^2\frac{1}{k_2}\mathrm{log}v_{k_2}^2$$
(46)
where, with obvious notations, $`u_{k_1}=\frac{|\mathrm{Tr}U^{k_1}|}{N}`$ and $`v_{k_2}=\frac{|\mathrm{Tr}V^{k_2}|}{N}`$. Although the equation for $`r(x)`$ cannot be derived, the variation with respect to $`l`$ gives the scaling power of $`l`$ and the cutoff transition point:
$$lM^{1\frac{\alpha }{4(\frac{1}{2k_1}+\frac{1}{2k_2})}}$$
(47)
The cutoff transition occurs then at critical point $`\alpha _c(k_1,k_2)`$ given by:
$$\alpha _c(k_1,k_2)=4(\frac{1}{2k_1}+\frac{1}{2k_2})$$
(48)
which generalizes the result on the sphere.
## 4 Wilson loops
In this section we calculate the expectation value of a Wilson loop in a Yang-Mills theory with gauge group U($`N`$) in the large $`N`$ limit. The space-time manifold has the topology of a two dimensional sphere whose area scales like $`\mathrm{log}N`$ as in eq. (11). We shall follow th approach of Daul and Kazakov and Boulatov who did the calculation for constant areas.
We may think of the sphere of area $`A`$ as two discs of areas $`A_1`$ and $`A_2`$ ($`A_1+A_2=A`$) sewed along their common boundary and with holonomy on the boundary respectively $`U`$ and $`U^{}`$. The Wilson looop is then given by
$`W(A_1,A_2)`$ $`=`$ $`{\displaystyle \frac{1}{N}}\mathrm{Tr}U`$ (49)
$`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{R_1,R_2}{}}d_1d_2{\displaystyle 𝑑U\frac{1}{N}\mathrm{Tr}U\chi _1(U)\chi _2(U^{})e^{\frac{A_1}{2N}C_1\frac{A_2}{2N}C_2}}`$
where $`d_1`$ and $`C_1`$ are the dimension and the Casimir operator in the representation $`R_1`$, referred to the disc of area $`A_1`$; likewise for $`d_2`$ and $`C_2`$. The quantity $`𝑑U\mathrm{Tr}U\chi _1(U)\chi _2(U^{})`$ may be either $`0`$ or $`1`$, namely it is $`1`$ when the Young diagram of $`R_2`$ is obtained by adding one box to the diagram of $`R_1`$ and $`0`$ otherwise. That is, if $`R_1`$ is labeled by the integers $`n_1>n_2>\mathrm{}>n_N`$, $`R_2`$ is labeled by a set on integers where one of the $`n_i`$ is increased by one. Daul and Kazakov used this property to get rid of one summation and obtained
$`W(A_1,A_2)`$ $`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{R_1}{}}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}d_1^{\mathrm{\hspace{0.17em}2}}{\displaystyle \underset{j,ji}{}}\left(1+{\displaystyle \frac{1}{n_jn_i}}\right)`$ (50)
$`e^{\frac{A_1+A_2}{2N}C_1}e^{\frac{A_2}{N}n_i},`$
where the sum over $`i`$ corresponds to all the possible ways of adding one box to the diagram. This is not however the whole result, in fact the original expression is symmetric under exchange of $`A_1`$ and $`A_2`$, so a term with $`A_1`$ and $`A_2`$ exchanged<sup>12</sup><sup>12</sup>12This term originates from the fact that $`𝑑U\mathrm{Tr}U\chi _1(U)\chi _2(U^{})`$ is different from zero also if the representation conjugate to $`R_1`$ is obtained from the representation conjugate to $`R_2`$ by adding a box in the Young diagram. must be added to (50) and gives:
$`W(A_1,A_2)`$ $`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{R_1}{}}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i}{}}d_1^{\mathrm{\hspace{0.17em}2}}{\displaystyle \underset{j,ji}{}}\left(1+{\displaystyle \frac{1}{n_jn_i}}\right)`$ (51)
$`e^{\frac{A_1+A_2}{2N}C_1}\left(e^{\frac{A_1}{N}n_i}+e^{\frac{A_2}{N}n_i}\right).`$
Moreover, eq. (50), as well as (51), is clearly not symmetric under $`n_in_i`$, so we can not restrict our considerations to $`n_i>0`$ any longer; instead we have to consider the whole interval $`\mathrm{}<n_i<\mathrm{}`$.
Let us first compute $`W(A_1,A_2)`$ in the frozen phase where the sum over $`R_1`$ is dominated by the trivial representation of dimension $`d_1=1`$ labeled by $`n_i=i\frac{N+1}{2}`$ $`i`$ and $`R_2`$ is the fundamental representation of dimension $`d_2=N`$ with $`C_2C_1=N`$. By inserting this into (51) one finds:
$$W(A_1,A_2)=\left(e^{\frac{A_1}{2}}+e^{\frac{A_2}{2}}\right)$$
(52)
which is, as expected, a typical strong coupling result. If for instance $`A_1>>A_2`$, then we have
$$W(A_1,A_2)e^{\frac{A_2}{2}}$$
(53)
Since $`A=A_1+A_2=\alpha \mathrm{log}N+\beta `$, this situation can occur in two ways:
* $`A_2=\beta _2`$, $`A_1=\alpha \mathrm{log}N+\beta _1`$: then $`We^{\frac{\beta _2}{2}}`$
* $`A_1=\alpha _1\mathrm{log}N+\beta _1`$, $`A_2=\alpha _2\mathrm{log}N+\beta _2`$ with $`\alpha _1>\alpha _2`$: then $`WN^{\frac{\alpha _2}{2}}`$.
We proceed now to evaluate $`W(A_1,A_2)`$ in the phase before the cutoff, namely for $`\alpha <4`$. The sum over representations in (51) can be replaced by the contribution of the dominant representation in the large $`N`$ limit, calculated on the sphere in section 2. The saddle point is unaffected by the presence of the extra term
$$\underset{i}{}\underset{j,ji}{}\left(1+\frac{1}{n_jn_i}\right)\left(e^{\frac{A_1}{N}n_i}+e^{\frac{A_2}{N}n_i}\right)$$
which is subleading with respect to the action. After replacing in (51) the sum with the saddle point contribution, some simplifications occur and we get
$$W(A_1,A_2)=\frac{1}{N}\underset{i}{}\underset{j,ji}{}\left(1+\frac{1}{n_jn_i}\right)\left(e^{\frac{A_1}{N}n_i}+e^{\frac{A_2}{N}n_i}\right).$$
(54)
The representation in (54) is of the type given in (15), and the sum over $`i`$ describes all possible ways of adding a box to the Young diagram. However the replacement $`n_in_i+1`$ is impossible in the region $`lil`$ as the resulting sequence of integers would not be monotonic increasing. So the sum over $`i`$ in (54) can be replaced by a sum over $`\alpha `$ with $`1\alpha l`$. As a matter of fact we must consider only positive $`\alpha `$’s, as adding a box to the adjoint representation amounts to symmetrize with respect to $`A_1`$ and $`A_2`$, and it has been taken already into account. Hence in (54) we must replace the index $`i`$ with $`\alpha `$, $`n_i`$ with $`Ml+\alpha +r_\alpha `$ while the index $`j`$ goes from $`M`$ to $`M`$, namely it goes over both the condensed and the non condensed regions. With these substitutions the expression for the Wilson loop becomes:
$`W(A_1,A_2)`$ $`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{l}{}}}{\displaystyle \frac{1}{\xi _\alpha }}\left[1{\displaystyle \frac{l}{N}}\left(1{\displaystyle \frac{\xi _\alpha }{l}}\right)\right]`$ (55)
$`e^{_{\beta =1}^l\mathrm{log}\left(1+\frac{1}{\xi _\alpha \xi _\beta }\right)\frac{A_2}{2}\left(12\frac{l\xi _\alpha }{N}\right)}+\{A_2A_1\}`$
where
$$\xi _\alpha =\alpha +r_\alpha $$
(56)
It is convenient as usual to use in the large $`N`$ limit the continuum variables $`x=\frac{\alpha }{l}`$, $`\xi (x)=\frac{\xi _\alpha }{l}`$ and the density function $`\rho (\xi )=\frac{dx}{d\xi }`$. The crucial part of the calculation is the evaluation of $`_{\beta =1}^l\mathrm{log}\left(1+\frac{1}{\xi _\alpha xi_\beta }\right)`$ which can be done following ref. . By expanding the logarithm one finds:
$$\underset{\beta =1}{\overset{l}{}}\mathrm{log}\left(1+\frac{1}{\xi _\alpha \xi _\beta }\right)=𝑑\eta \rho (\eta )\frac{1}{\xi \eta }\underset{k=2}{}\frac{1}{k}\underset{\beta }{}\frac{1}{(\beta \alpha )^k}\rho (\xi )^k$$
(57)
The first term at the r.h.s. comes from the $`k=1`$ term of the $`\mathrm{log}`$ expansion and can be evaluated using eq. (27), the other terms can be calculated as in and give $`\mathrm{log}\frac{\mathrm{sin}\pi \rho }{\pi \rho }`$. By inserting these results into (55) and using the explicit form of the solution (30) one finally obtains (neglecting $`O(\frac{l}{N})`$ terms):
$$W(A_1,A_2)=2𝑑\xi \frac{\mathrm{sin}\pi \rho }{\pi }\left(e^{\frac{A_2}{2}}+e^{\frac{A_1}{2}}\right)=e^{\frac{A_2}{2}}+e^{\frac{A_1}{2}}$$
(58)
which is exactly the same result as in the frozen phase. The result is not trivial, but it was somehow to be expected. Both phases, before and after the transition, are strong coupling phases and the expectation value of the Wilson loop should be in both of them obtained, to the leading order, by filling the loop with elementary plaquettes in the fundamental representation. The effects of the transition are expected to appear only at the next-to-leading order ( $`\frac{l}{N}`$) which is sensitive to the $`O(\frac{l}{N})`$ degrees of freedom which are not frozen below the cutoff.
## Acknowledgments
A.A. wants to thank the ”Service de Physique Theorique (SPhT)” at Saclay for the kind hospitality of these last months. He also would like to thank I. Kostov and G. Vernizzi for many useful discussions during this period.
S.A. would like to thank the Niels Bohr Institute for warm hospitality during the early stage of this work.
## Appendix A Gauge theories as random walks
The equivalence between random walks on a group $`G`$ and two dimensional gauge theories with gauge group $`G`$ was pointed out in . This equivalence states that the partition function of a gauge theory on a disc of area $`A`$ and holonomy $`g`$ on the boundary coincides with the probability that a suitably defined random walk on the group leads, in a number of steps proportional to $`A`$, from the identity in the group to the group element $`g`$. A similar relation holds for the partition function on the cylinder. In ref. the attention was focused on discrete groups, where the number of steps can be directly identified with the area measured in suitable units, and in particular on the symmetric group $`S_n`$. In this appendix we extend, in a rather straightforward way, the argument presented there to the case of gauge theories on continous Lie groups. Similar conclusions have in the meantime appeared in the literature .
Let us consider a random walk on $`G`$ with the transition probability defined as follows: if the walker is in $`g_pG`$ after $`p`$ steps, then his position after $`p+1`$ steps is obtained by left multiplying $`g_p`$ by an element $`g`$, chosen in $`G`$ with a probability $`t(g)`$. We assume $`t(g)`$ to be a class function, whose character expansion can then be written as:
$$t(g)=\underset{r}{}\mathrm{\Delta }_r\chi _r(g)\stackrel{~}{t}_r.$$
(59)
where $`\mathrm{\Delta }_r`$ is the dimension of the representation $`r`$. Suppose that the random walk starts from an element $`g_0G`$ and denote by $`K_p(g,g_0)`$ the probability for the walker to be in $`g`$ after $`p`$ steps. Of course $`K_0(g,g_0)`$ is a delta function, namely:
$$K_0(g,g_0)=\underset{r}{}\mathrm{\Delta }_r\chi _r(g^1g_0)$$
(60)
Given $`K_p(g^{},g_0)`$, the probability for the walker to be in $`g`$ after $`p+1`$ steps is given by:
$$K_{p+1}(g,g_0)=𝑑g^{}t\left(g(g^{})^1\right)K_p(g^{},g_0),$$
(61)
where $`dg^{}`$ is the Haar measure. By using (59), (60) and (61) it is easy to show by induction that $`K_p(g,g_0)`$ is a class function of $`gg_0^1`$. In fact let us assume that this is true for $`K_p(g^{},g_0)`$, namely that $`K_p(g^{},g_0)`$ admits the character expansion
$$K_p(g^{},g_0)=\underset{r}{}\mathrm{\Delta }_rk_r^{(p)}\chi _r(g^{}g_0^1)$$
(62)
Then by replacing (62) and (59) into (61) and performing the integration over $`g^{}`$ using the characters fusion rules, one finds that $`K_{p+1}(g,g_0)`$ admits a similar expansion with
$$k_r^{(p+1)}=\stackrel{~}{t}_rk_r^{(p)}$$
(63)
As a result $`K_p(g,g_0)`$ is given by:
$$K_p(g,g_0)=\underset{r}{}\mathrm{\Delta }_r\chi _r(gg_0^1)\stackrel{~}{t}_r^p$$
(64)
Since we are dealing with a random walk on a continous Lie group, we expect the walk to be a smooth path. This is obtained by letting the number of steps $`p`$ go to infinity and at the same time the length of each step to zero. Each step corresponds then to a very small move on the group manifold, that is $`t(g)`$ is close to a delta function $`\delta (g)`$, namely according to (59), $`\stackrel{~}{t}_r1`$. We implement this by choosing
$$\stackrel{~}{t}_r=e^{\epsilon h(r)},\text{with}\epsilon 0$$
(65)
If we set $`g_0=1`$ in (64),
$$\underset{\epsilon 0}{lim}\underset{p\mathrm{}}{lim}p\epsilon =A\text{and}h(r)=\frac{C_2(r)}{2N}$$
(66)
we reproduce exactly the partition function of a gauge theory on a disc of area $`A`$, thus establishing the desired connection. A random walk with an arbitrary transition function $`h(r)`$ will be related to a generalized Yang-Mills theory with the corresponding potential .
The partition function on a cylinder may be associated to the probability of a random walk from a generic point $`g_0`$ (in general different from the identity) to $`g`$. However in order to obtain the correct answer we need to consider a random walk not on the group manifold itself, but rather in the space of orbits, where all group elements belonging to the same equivalent class are identified. This amounts to replacing in (64) $`g_0h^1g_0h`$ and integrating over $`h`$. This produces a new transition probability $`\stackrel{~}{K}_p(g,g_0)`$ which is separately a class function of both $`g`$ and $`g_0`$ and is given by:
$$\stackrel{~}{K}_p(g,g_0)=\underset{r}{}\chi _r(g)\chi _r(g_0^1)\stackrel{~}{t}_r^p$$
(67)
This coincides with the kernel on the cylinder, with the parameters identified according to (66).
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# Density matrix of a finite sub-chain of the Heisenberg anti-ferromagnet
## 1. Introduction
The present paper continues the study of correlation functions for integrable spin chains launched in . In our previous works, we found an exact expression without involving integrals for the density matrix of a finite sub-chain of the infinite XXX, XXZ and XYZ chains in the ground state. More precisely, we treated inhomogeneous models in which each site carries an independent spectral parameter. The problem of finding a compact form for the answer in the physically important homogeneous case remained open. This is exactly the problem which we solve in the present paper. We shall consider the simplest case of the XXX model.
Consider the isotropic Heisenberg antiferromagnet with the Hamiltonian
$$H=\frac{1}{2}\underset{i=\mathrm{}}{\overset{\mathrm{}}{}}\left(\sigma _i^1\sigma _{i+1}^1+\sigma _i^2\sigma _{i+1}^2+\sigma _i^3\sigma _{i+1}^3\right).$$
Take a finite sub-chain consisting of sites $`i=1,\mathrm{},n`$. The density matrix $`\rho _n`$ for this sub-chain in the infinite environment is an operator acting on $`\left(^2\right)^n`$. Its matrix elements are given by the ground state average
$$\left(\rho _n\right)_{\overline{ϵ}_1,\mathrm{},\overline{ϵ}_n}^{ϵ_1,\mathrm{},ϵ_n}=\text{vac}|\left(E_{ϵ_1}^{\overline{ϵ}_1}\right)_1\mathrm{}\left(E_{ϵ_n}^{\overline{ϵ}_n}\right)_n|\text{vac}.$$
Here $`|\text{vac}`$ is the anti-ferromagnetic ground state, $`ϵ_j,\overline{ϵ}_j=+,`$, and $`\left(E_ϵ^{\overline{ϵ}}\right)_i`$ signifies the matrix unit $`\left(\delta _{aϵ}\delta _{b\overline{ϵ}}\right)_{a,b=+,}`$ acting on the $`i`$-th tensor component. It is important to consider a vector $`h_n`$ belonging to $`\left(^2\right)^{2n}`$ instead of the matrix $`\rho _n`$ acting in $`\left(^2\right)^n`$ . The vector $`h_n`$ is given by:
$$h_n^{ϵ_1,\mathrm{},ϵ_n,\overline{ϵ}_n,\mathrm{},\overline{ϵ}_1}=\underset{j=1}{\overset{n}{}}(\overline{ϵ}_j)\left(\rho _n\right)_{\overline{ϵ}_1,\mathrm{},\overline{ϵ}_n}^{ϵ_1,\mathrm{},ϵ_n}.$$
In the sequel we refer to the tensor components of $`\left(^2\right)^{2n}`$ by the indices $`1,\mathrm{},n,\overline{n},\mathrm{},\overline{1}`$, read from left to right. The main result of our previous papers can be formulated as follows:
(1.1) $`h_n=e^{\mathrm{\Omega }_n}𝐬_n.`$
Here $`𝐬_n=s_{j,\overline{j}}`$ and $`s_{j,\overline{j}}`$ signifies the $`𝔰𝔩_2`$-singlet $`\frac{1}{2}\left(v_+v_{}v_{}v_+\right)`$ in the tensor product of two spaces $`j,\overline{j}`$, $`v_+,v_{}`$ being the standard basis of $`^2`$. The operator $`\mathrm{\Omega }_n`$ will be defined later (see (2.5)). It satisfies the condition
(1.2) $`\mathrm{\Omega }_n^{\left[\frac{n}{2}\right]+1}=0.`$
So, the series for the exponential (1.1) terminates.
## 2. Inhomogeneous case
Let us introduce inhomogeneity parameters $`\lambda _1,\mathrm{},\lambda _n`$ to the corresponding sites of the lattice (see for more details). Then the operator $`\mathrm{\Omega }_n(\lambda _1,\mathrm{},\lambda _n)`$ becomes dependent on these parameters. We define this operator following our previous works, but we shall slightly change the notation.
Let $`\{S_a\}_{a=1}^3`$ be a basis of $`𝔰𝔩_2`$ satisfying $`[S_a,S_b]=2iϵ_{abc}S_c`$. Following define the $`L`$-operator which belongs to $`U(𝔰𝔩_2)\text{End}\left(^2\right)`$:
(2.1) $`L(\lambda )={\displaystyle \frac{\rho (\lambda ,d)}{\lambda +\frac{d}{2}}}L^{(0)}(\lambda ),L^{(0)}(\lambda )=\lambda +{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{a=1}{\overset{3}{}}}S_a\sigma ^a,`$
where $`d`$ is related to the Casimir operator as $`_{a=1}^3S_a^2=d^21`$, and
$$\rho (\lambda ,d)=\frac{\mathrm{\Gamma }\left(\frac{1}{2}\frac{d}{4}+\frac{\lambda }{2}\right)\mathrm{\Gamma }\left(1\frac{d}{4}\frac{\lambda }{2}\right)}{\mathrm{\Gamma }\left(\frac{1}{2}\frac{d}{4}\frac{\lambda }{2}\right)\mathrm{\Gamma }\left(1\frac{d}{4}+\frac{\lambda }{2}\right)}.$$
In this normalization we have the unitarity and crossing symmetry in the form
$$L(\lambda )L(\lambda )=1,\sigma ^2\left(L(\lambda )\right)^t\sigma ^2=L(\lambda 1).$$
We shall consider tensor products of several spaces $`^2`$. In that case the index $`i`$ in $`L_i(\lambda )`$ denotes the tensor component as usual. In what follows the function $`\rho (\lambda ,d)`$ always comes in the combination
$$\frac{\rho (\lambda ,d)}{\lambda +\frac{d}{2}}\frac{\rho (\lambda 1,d)}{\lambda +\frac{d}{2}1}=\frac{1}{\lambda ^2\frac{d^2}{4}},$$
so the $`\mathrm{\Gamma }`$-functions will never really appear.
We shall also use the ordinary $`4\times 4`$ $`R`$-matrix obtained as the image of the $`L`$-operator (2.1) in the 2-dimensional representation of $`U(𝔰𝔩_2)`$. When acting in tensor product of two spaces $`i,j`$, it will be denoted by $`R_{i,j}(\lambda )`$. We denote the corresponding factor by
$$\rho (\lambda ):=\rho (\lambda ,2).$$
Let us explain the results of the papers in the setting of the XXX model. First, introduce the operator
$`T_n^{[1]}(\lambda ;\lambda _2,\mathrm{},\lambda _n)`$ $`:=`$ $`L_{\overline{2}}(\lambda \lambda _21)\mathrm{}L_{\overline{n}}(\lambda \lambda _n1)L_n(\lambda \lambda _n)\mathrm{}L_2(\lambda \lambda _2).`$
Notice that in this product the sites $`1`$, $`\overline{1}`$ are omitted.
In the paper , we discussed in detail the linear functional on $`U(𝔰𝔩_2)`$ called $`\mathrm{Tr}_\lambda `$. Denote by $`\varpi _d`$ the irreducible $`d`$-dimensional representation of $`U(𝔰𝔩_2)`$. By definition, the map $`\mathrm{Tr}_\lambda :U(𝔰𝔩_2)[\lambda ]`$ associates with each $`AU(𝔰𝔩_2)`$ a unique polynomial $`\mathrm{Tr}_\lambda (A)`$ of $`\lambda `$, such that $`\mathrm{Tr}_d(A)=\mathrm{tr}_^d\varpi _d(A)`$ holds for any positive integer $`d`$. Some of its main properties are
$`\mathrm{Tr}_\lambda (AB)=\mathrm{Tr}_\lambda (BA),`$
$`\mathrm{Tr}_\lambda \left(\left({\displaystyle \underset{a=1}{\overset{3}{}}}S_a^2\right)A\right)=(\lambda ^21)\mathrm{Tr}_\lambda (A),`$
$`\mathrm{Tr}\left(e^{tS^3}\right)={\displaystyle \frac{\mathrm{sinh}t\lambda }{\mathrm{sinh}t}}.`$
In the present paper we shall also use the linear functional
$$\mathrm{Tr}_{\lambda _1,\mathrm{},\lambda _k}:U(𝔰𝔩_2)^k[\lambda _1,\mathrm{},\lambda _k]$$
defined by
$$\mathrm{Tr}_{\lambda _1,\mathrm{},\lambda _k}(A_1\mathrm{}A_k)=\underset{j=1}{\overset{k}{}}\mathrm{Tr}_{\lambda _j}(A_j).$$
The main ingredient of our construction is the operator
(2.2) $`X_n(\lambda _1,\mathrm{},\lambda _n)_{1,\mathrm{},n,\overline{n},\mathrm{},\overline{1}}:=(1)^{n1}\mathrm{res}_{\lambda =\lambda _2}\mathrm{Tr}_{\lambda _1\lambda }\left(T^{[1]}_n\right({\displaystyle \frac{\lambda _1+\lambda }{2}}\left)\right)P_{1,\overline{2}}𝒫_{1,\overline{1}}^{}𝒫_{2,\overline{2}}^{},`$
where $`P_{i,j}`$ is the permutation and $`𝒫_{i,j}^{}=(1P_{i,j})/2`$ is the skew-symmetriser. Define further
$`\mathrm{\Omega }_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)=4\omega (\lambda _{i,j})X_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n),`$
$`X_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)=\stackrel{}{}_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)`$
$`\times X_n(\lambda _i,\lambda _j,\lambda _1,\mathrm{},\widehat{\lambda _i},\mathrm{},\widehat{\lambda _j},\mathrm{},\lambda _n)_{i,j,1\mathrm{}\widehat{i}\mathrm{}\widehat{j}\mathrm{}n,\overline{n},\mathrm{}\widehat{\overline{j}}\mathrm{}\widehat{\overline{i}}\mathrm{}\overline{1},\overline{j},\overline{i}}\stackrel{}{}_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n),`$
where
$`\stackrel{}{}_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)`$
$`:=R_{i,i1}(\lambda _{i,i1})\mathrm{}R_{i,1}(\lambda _{i,1})R_{\overline{i1},\overline{i}}(\lambda _{i1,i})\mathrm{}R_{\overline{1},\overline{i}}(\lambda _{1,i})`$
$`\times R_{j,j1}(\lambda _{j,j1})\mathrm{}R_{j,i+1}(\lambda _{j,i+1})R_{j,i1}(\lambda _{j,i1})\mathrm{}R_{j,1}(\lambda _{j,1})`$
$`\times R_{\overline{j1},\overline{j}}(\lambda _{j1,j})\mathrm{}R_{\overline{i+1},\overline{j}}(\lambda _{i+1,j})R_{\overline{i1},\overline{j}}(\lambda _{i1,j})\mathrm{}R_{\overline{1},\overline{j}}(\lambda _{1,j}),`$
$`\stackrel{}{}_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)`$
$`:=R_{\overline{n},\overline{i}}(\lambda _{n,i})\mathrm{}R_{\overline{j+1},\overline{i}}(\lambda _{j+1,i})R_{\overline{j1},\overline{i}}(\lambda _{j1,i})\mathrm{}R_{\overline{i+1},\overline{i}}(\lambda _{i+1,i})`$
$`\times R_{i,n}(\lambda _{i,n})\mathrm{}R_{i,j+1}(\lambda _{i,j+1})R_{i,j1}(\lambda _{i,j1})\mathrm{}R_{i,i+1}(\lambda _{i,i+1})`$
$`\times R_{\overline{n},\overline{j}}(\lambda _{n,j})\mathrm{}R_{\overline{j+1},\overline{j}}(\lambda _{j+1,j})R_{j,n}(\lambda _{j,n})\mathrm{}R_{j,j+1}(\lambda _{j,j+1})`$
and
$$\omega (\lambda )=\frac{d}{d\lambda }\mathrm{log}\rho (\lambda )+\frac{1}{2(\lambda ^21)}.$$
The definition of $`\mathrm{\Omega }^{(i,j)}`$ differs from the one used in <sup>3</sup><sup>3</sup>3See eq.(12.1) in . We have also used the fact that, in the notation there, $`L_2(\lambda _{1,2}/2)L_{\overline{2}}(\lambda _{1,2}/21)𝒫_{2,\overline{2}}^{}=0`$ inside the trace and $`P_{2,\overline{2}}s_{1,\overline{2}}s_{\overline{1},2}{}_{n2}{}^{}\mathrm{\Pi }_{n}^{}=P_{1,\overline{2}}𝒫_{1,\overline{1}}^{}𝒫_{2,\overline{2}}^{}`$. by the second product of $`R`$-matrices, but the final formula (1.1) remains unaltered by this modification.
The operators $`\mathrm{\Omega }_n^{(i,j)}`$ possess a number of properties the most important among which are
(2.3) $`[\mathrm{\Omega }_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n),\mathrm{\Omega }_n^{(k,l)}(\lambda _1,\mathrm{},\lambda _n)]=0,`$
(2.4) $`\mathrm{\Omega }_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)\mathrm{\Omega }_n^{(k,l)}(\lambda _1,\mathrm{},\lambda _n)=0\text{if}\{i,j\}\{k,l\}\mathrm{}.`$
As a function of $`\lambda _1,\mathrm{},\lambda _n`$, $`\mathrm{\Omega }_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)`$ is meromorphic. All poles are simple and located at $`\lambda _i,\lambda _j=\lambda _l`$ and $`\lambda _i,\lambda _j=\lambda _l\pm 1`$ for $`li,j`$ (which are due respectively to $`X^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)`$), and to the $`R`$-matrices) and $`\lambda _{i,j}\backslash \{0\}`$ (due to $`\omega (\lambda _{i,j})`$). We have
$`h_n(\lambda _1,\mathrm{},\lambda _n)=e^{\mathrm{\Omega }_n(\lambda _1,\mathrm{},\lambda _n)}𝐬_n,`$
where
(2.5) $`\mathrm{\Omega }_n(\lambda _1,\mathrm{},\lambda _n)={\displaystyle \underset{1i<jn}{}}\mathrm{\Omega }_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n).`$
The properties (2.3), (2.4) guarantee the nilpotency (1.2).
## 3. Homogeneous case
Our goal is to obtain the homogeneous limit $`\lambda _1=\mathrm{}=\lambda _n=0`$. In the original formula (2.5), this problem is very complicated: the singularities on the diagonal $`\lambda _i=\lambda _j`$ are present in every term of (2.5). Although these poles are absent in the sum itself, it is technically difficult to explicitly carry through the cancellation and obtain the final answer.
So, we need to rewrite the formula for $`\mathrm{\Omega }_n(\lambda _1,\mathrm{},\lambda _n)`$ in such a way that taking the homogeneous limit is easier. To this end, let us write first of all another formula for $`X_n(\lambda _1,\mathrm{},\lambda _n)`$.
Denote by $`\varpi _\lambda `$ the $`\lambda `$-dimensional irreducible representation. We use only the fact that the Casimir element reduces to $`\lambda ^21`$, and hence the following computation makes sense for non-integer $`\lambda `$ as well .
Notice that
$$P_a^+(\lambda )=\frac{1}{\lambda }\varpi _\lambda (L_a^{(0)}(\lambda /2)),P_a^{}(\lambda )=\frac{1}{\lambda }\varpi _\lambda (L_a^{(0)}(\lambda /2))$$
are orthogonal projectors. Consider now $`X_n(\lambda _1,\mathrm{},\lambda _n)`$. Using the formula $`P_{1,\overline{2}}𝒫_{1,\overline{1}}^{}𝒫_{2,\overline{2}}^{}=𝒫_{1,2}^{}𝒫_{\overline{1},\overline{2}}^{}P_{1,\overline{2}}`$, the definition of the $`L`$-operator and the crossing-symmetry, one finds
$`\varpi _{\lambda _{1,2}}\left(L_2^{(0)}({\displaystyle \frac{\lambda _{1,2}}{2}})L_1^{(0)}({\displaystyle \frac{\lambda _{2,1}}{2}})\right)𝒫_{1,2}^{}`$ $`=\varpi _{\lambda _{1,2}}\left(L_2^{(0)}({\displaystyle \frac{\lambda _{1,2}}{2}})L_2^{(0)}({\displaystyle \frac{\lambda _{1,2}}{2}}1)\right)𝒫_{1,2}^{}`$
$`=(\lambda _{1,2}1)\varpi _{\lambda _{1,2}}\left(L_2^{(0)}({\displaystyle \frac{\lambda _{1,2}}{2}})\right)𝒫_{1,2}^{},`$
$`\varpi _{\lambda _{1,2}}\left(L_{\overline{1}}^{(0)}({\displaystyle \frac{\lambda _{2,1}}{2}}1)L_{\overline{2}}^{(0)}({\displaystyle \frac{\lambda _{1,2}}{2}}1)\right)𝒫_{\overline{1},\overline{2}}^{}`$ $`=\varpi _{\lambda _{1,2}}\left(L_{\overline{1}}^{(0)}({\displaystyle \frac{\lambda _{2,1}}{2}}1)L_{\overline{1}}^{(0)}({\displaystyle \frac{\lambda _{2,1}}{2}})\right)𝒫_{\overline{1},\overline{2}}^{}`$
$`=(\lambda _{1,2}+1)\varpi _{\lambda _{1,2}}\left(L_{\overline{1}}^{(0)}({\displaystyle \frac{\lambda _{2,1}}{2}})\right)𝒫_{\overline{1},\overline{2}}^{}`$
$`=(\lambda _{1,2}+1)\varpi _{\lambda _{1,2}}\left(L_{\overline{2}}^{(0)}({\displaystyle \frac{\lambda _{1,2}}{2}}1)\right)𝒫_{\overline{1},\overline{2}}^{}.`$
Now it is easy to see that
(3.1) $`X_n(\lambda _1,\mathrm{},\lambda _n)`$
$`=(1)^{n1}\mathrm{res}_{\mu _1=\lambda _1}\mathrm{res}_{\mu _2=\lambda _2}{\displaystyle \frac{\mu _{1,2}}{\mu _{1,2}^21}}\mathrm{Tr}_{\mu _{1,2}}\left(T_n\right({\displaystyle \frac{\mu _1+\mu _2}{2}};\lambda _1,\mathrm{},\lambda _n\left)\right)P_{1,\overline{2}}𝒫_{1,\overline{1}}^{}𝒫_{2,\overline{2}}^{}`$
where $`T_n(\lambda )`$ is the complete monodromy matrix:
$`T_n(\lambda ;\lambda _1,\mathrm{},\lambda _n)`$ $`=`$ $`L_{\overline{1}}(\lambda \lambda _11)\mathrm{}L_{\overline{n}}(\lambda \lambda _n1)L_n(\lambda \lambda _n)\mathrm{}L_1(\lambda \lambda _1).`$
Using the Yang-Baxter equation one finds the following formula for $`X_n^{(i,j)}`$:
$`X_n^{(i,j)}`$ $`(\lambda _1,\mathrm{},\lambda _n)=(1)^{n1}\mathrm{res}_{\mu _1=\lambda _i}\mathrm{res}_{\mu _2=\lambda _j}{\displaystyle \frac{\mu _{1,2}}{\mu _{1,2}^21}}\mathrm{Tr}_{\mu _{1,2}}\left(T_n\right({\displaystyle \frac{\mu _1+\mu _2}{2}};\lambda _1,\mathrm{},\lambda _n\left)\right)`$
$`\times \stackrel{}{}_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)P_{i,\overline{j}}𝒫_{i,\overline{i}}^{}𝒫_{j,\overline{j}}^{}\stackrel{}{}_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n).`$
By a straightforward computation one finds
$`\stackrel{}{}_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)P_{i,\overline{j}}𝒫_{i,\overline{i}}^{}𝒫_{j,\overline{j}}^{}\stackrel{}{}_n^{(i,j)}(\lambda _1,\mathrm{},\lambda _n)`$
$`={\displaystyle \frac{1}{2}}\mathrm{Tr}_{2,2}\left(T_n(\lambda _i;\lambda _1,\mathrm{},\lambda _n)T_n(\lambda _j;\lambda _1,\mathrm{},\lambda _n)𝒫^{}\right).`$
Here the skew-symmetriser $`𝒫^{}`$ acts on the auxiliary space $`^2^2`$. Notice that $`\mathrm{Tr}_{2,2}\left(T_n(\mu _1;\lambda _1,\mathrm{},\lambda _n)T_n(\mu _2;\lambda _1,\mathrm{},\lambda _n)𝒫^{}\right)`$ is actually symmetric with respect to $`\mu _1`$, $`\mu _2`$ due to the relation
$$[R(\mu ),𝒫^{}]=0.$$
Obviously, the formula for $`\mathrm{\Omega }_n(\lambda _1,\mathrm{},\lambda _n)`$ can be rewritten now as
$`\mathrm{\Omega }_n(\lambda _1,\mathrm{},\lambda _n)`$ $`=`$ $`{\displaystyle \frac{(1)^n}{2}}{\displaystyle \frac{d\mu _1}{2\pi i}\frac{d\mu _2}{2\pi i}\omega (\mu _{1,2})\mathrm{Tr}_{\mu _{1,2}}\left(T_n(\frac{\mu _1+\mu _2}{2};\lambda _1,\mathrm{},\lambda _n)\right)}`$
$`\times \mathrm{Tr}_{2,2}\left(T_n(\mu _1;\lambda _1,\mathrm{},\lambda _n)T_n(\mu _2;\lambda _1,\mathrm{},\lambda _n)B(\mu _{1,2})\right),`$
where
$$B(\mu _{1,2})=\frac{2\mu _{1,2}}{\mu _{1,2}^21}𝒫^{},$$
the contours of integration encircle the poles $`\mu _1=\lambda _j`$, $`\mu _2=\lambda _j`$ for $`j=1,\mathrm{},n`$. The great advantage of this formula is that it allows to take the homogeneous limit $`\lambda _j=0`$. In the next formula we write $`\mathrm{\Omega }_n`$ for $`\mathrm{\Omega }_n(0,\mathrm{},0)`$, etc..
$`\mathrm{\Omega }_n`$ $`=`$ $`{\displaystyle \frac{(1)^n}{2}}{\displaystyle \frac{d\mu _1}{2\pi i}\frac{d\mu _2}{2\pi i}\omega (\mu _{1,2})\mathrm{Tr}_{\mu _{1,2}}\left(T_n\left(\frac{\mu _1+\mu _2}{2}\right)\right)}`$
$`\times \mathrm{Tr}_{2,2}\left(T_n(\mu _1)T_n(\mu _2)B(\mu _{1,2})\right),`$
where the integrals are taken around $`\mu _i=0`$.
Formulas (3), (3) are the main results of the present paper. Let us discuss them briefly.
First we note that the integrand of (3) has a pole of order $`n`$ at $`\mu _j=0`$. By evaluating residues, $`\mathrm{\Omega }_n`$ becomes a linear combination of the Taylor coefficients of $`\omega (\lambda )`$ with $`jn`$, given explicitly by
$`\omega (\lambda ){\displaystyle \frac{1}{2(\lambda ^21)}}=2\left(\mathrm{log}2+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\zeta _a(2k+1)\lambda ^{2k}\right).`$
Here $`\zeta _a(s)=(12^{1s})\zeta (s)`$, $`\zeta (s)`$ denoting the Riemann zeta function. This settles the conjecture of which states that any correlation function of the XXX model can be written as a polynomial of $`\mathrm{log}2`$ and $`\zeta (3),\zeta (5),\mathrm{}`$ with rational coefficients.
Second, formula (3) may open up a way for studying the large-distance limit. Also it would be very interesting to see if it helps for the investigation of the limit to continuous field theory. We hope to return to these problems as well as the extension to the XXZ and XYZ cases in future publications.
Acknowledgments. Research of HB is supported by the RFFI grant #04-01-00352. Research of MJ is supported by the Grant-in-Aid for Scientific Research B2–16340033. Research of TM is supported by the Grant-in-Aid for Scientific Research A1–13304010. Research of FS is supported by INTAS grant #03-51-3350 and by EC networks ”EUCLID”, contract number HPRN-CT-2002-00325 and ”ENIGMA”, contract number MRTN-CT-2004-5652. Research of YT is supported by Grant-in-Aid for Young Scientists (B) No. 17740089. This work was also supported by the grant of 21st Century COE Program at RIMS, Kyoto University.
HB is grateful to F. Göhmann, A. Klümper and J. Suzuki for discussions.
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# Beyond LISA: Exploring Future Gravitational Wave Missions
## I Introduction
The launch of the Laser Interferometer Space Antenna (LISA)PrePhaseA in the next decade will usher in a new era for gravitational wave detection. LISA will not only detect the presence of gravitational waves, but also provide detailed information about many thousands of gravitational wave sources, vastly expanding the burgeoning field of gravitational wave astronomy.
The primary sources for LISA are expected to be compact galactic binariesH\_B\_W , supermassive black hole binaries, and extreme mass ratio inspirals. The gravitational wave signals detected by LISA encode information about the parameters of the binary systems, such as sky locations, masses, and distances, which can be extracted by matched filtering. Studies into the accuracy with which the source parameters can be recovered from the LISA data have been performed for compact galactic binariesCutler ; Hellings ; Seto1 ; CrowderCornish , supermassive black hole binariesCutler ; Hellings ; Hughes ; Seto2 ; Vecchio , and extreme mass ratio inspiralsBarack .
It is hoped that LISA will be the first of several efforts to explore the low frequency portion of the gravitational wave spectrum accessible to space borne interferometers. One possible follow on mission is tentatively named the Advanced Laser Interferometer Antenna (ALIA)ALIA\_white\_paper , featuring a spacecraft configuration similar to that of LISA, with smaller arm lengths and lower noise levels. The primary mission of ALIA will be detecting intermediate mass black holes (IMBHs), with masses in the range $`5050,000M_{}`$. Information on populations, locations, and event rates could greatly enhance theories of black hole formation and evolution. Another proposed follow on mission is the Big Bang ObserverBBO\_proposal (BBO). The BBO will be an extremely sensitive antenna that is designed to detect the Gravitational Wave Background (GWB) left by the Big Bang. According to the standard cosmological picture, the GWB is a relic of the early inflationary period of the Universe. Just as the COsmic Background Explorer (COBE) and the Wilkinson Microwave Anisotropy Probe (WMAP) missions provided information about the Universe around the time of last scattering, the BBO should be able to provide information about the earliest moments in the history of the Universe. The current BBO proposal calls for four LISA-like spacecraft constellations, with the orbital configuration shown in Figure 1. Two of the constellations will be centered on a $`20^{}`$ Earth-trailing orbit, rotated $`60^{}`$ with respect to each other in the plane of the constellations. These constellations will be referred to as the star constellations, as the legs of the constellation sketch out a six pointed star. The remaining two constellations are to be placed in an Earth-like orbit $`120^{}`$ ahead and behind the star constellations. These two constellations will be referred to as the outrigger constellations. The purpose of the outrigger constellations is to provide greater angular resolution for foreground sources (see Section III for details). The star constellations provide maximum cross-correlation of gravitational wave signals between the two constellations, with minimal correlated noise Cornish\_Larson . The noise in each detector is expected to be independent, so over time the overlap of the noise between the two constellations will tend to average to zero, while the overlap of the signal will grow. The plan is to deploy the BBO in stages, starting with the star constellation, then adding the outrigger constellations at a later date.
Figure 2 shows the detector sensitivities for LISA, ALIA, and BBO, a level for the expected extra-galactic background confusion noise Phinney\_Farmer , and optimally filtered amplitude plots for equal mass binaries in their last year before coalescence at redshift $`z=1`$. The squares shown on the amplitude plots denote the frequency of the signal one week before coalescence. The signal from coalescing equal mass binaries at $`z=1`$ with masses above $`10^1M_{}`$ will be detectable by both ALIA and BBO. The BBO will also be able to detect lower mass systems such as neutron star (NS) binaries. However, for coalescing binaries, a large portion of the signal strength is due to the rapid inspiral in the last week, and for $`10M_{}`$ binaries at $`z=1`$ or greater, much of this power is deposited after the signal has crossed the ALIA sensitivity curve.
Here we begin by analyzing the parameter resolution of ALIA for IMBHs. These studies indicate that by including a second identical constellation of spacecraft in a $`20^{}`$ Earth-leading orbit, providing a constellation separation of $`40^{}`$, the parameter resolution will be greatly improved. We will call this two constellation configuration the Advanced Laser Interferometer Antenna in Stereo (ALIAS). We will also address the natural question, given the increases in precision that ALIAS provides relative to ALIA, what would be gained by adding a second constellation to the LISA mission - the Laser Interferometer Space Antenna in Stereo (LISAS). Lastly, we analyze the parameter resolution of the BBO for NSs, $`10M_{}`$, and $`100M_{}`$ black holes (BHs), and for the initial deployment phase of the BBO, which we call BBO Star, containing the two star constellations of the BBO configuration, but not the outrigger constellations. We will show that parameter resolution of the BBO and BBO Star should be sufficient to detect foreground binary systems, out to and beyond a redshift of $`z=3`$.
This paper is organized as follows. We begin with a review of parameter estimation in Section II. In Section III we describe how a pair of well separated detectors can significantly improve the angular resolution for transient sources. This is followed by our study of the ALIA mission in Section IV.1, and a study of our dual-detector variant of the ALIA mission in Section IV.2. Results for unequal mass binaries for ALIA and ALIAS are in Section IV.3. The LISA and LISAS missions are compared in Section IV.4. The BBO mission is studied in Section IV.5, and results for a down-scoped version of the BBO mission are reported in Section IV.6. Concluding remarks are made in Section V.
## II Review of Parameter Estimation
In our analysis, the orbits of the binaries are treated as quasi-circular and spin effects are neglected. This leaves nine parameters that will describe the binary systems: sky location ($`\theta ,\varphi `$); inclination and polarization angles ($`\iota ,\psi `$); reduced and chirp masses ($`\mu ,`$); time to coalescence ($`t_c`$), luminosity distance ($`D_L`$), and the initial orbital phase ($`\gamma _o`$). The signals are modeled using a truncated second-Post Newtonian ($`2PN`$) approximation2pn whereby the amplitude is kept to Newtonian order while the phase is kept to second order. In other words, we only include the dominant second harmonic of the orbital frequency.
The response of a space-borne instrument to a gravitational wave source is encoded in the Michelson-like time-delay interferometry (TDI) variables tdi , $`X_i(t)`$. Here the subscript $`i`$ denotes the vertex at which the signal is read out. In the equal-armlength limit, the TDI signal, $`X_i(t)`$, can be formed from a time-delayed combination of Michelson signals, $`M_i(t)`$, by
$$X_i(t)=M_i(t)M_i(t2L).$$
(1)
This differencing cancels the laser phase noise, while preserving the gravitational wave signal. Rather than work with the correlated $`X_i(t)`$ variables directly, we use the orthogonal signal combinations PTLA02
$`A(t)`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}\left(S_{II}(t)S_{II}(tL)\right),`$
$`E(t)`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}\left(S_I(t)S_I(tL)\right),`$ (2)
where
$`S_I(t)`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(2M_1(t)M_2(t)M_3(t)\right),`$
$`S_{II}(t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{3}}}\left(M_2(t)M_3(t)\right).`$ (3)
The $`A,E`$ combinations cancel the laser phase noise that would otherwise dominate the Michelson signals. At high frequencies, $`f>f_{}c/(2\pi L)`$, where $`L`$ is the length of the detector arms, a third independent data channel $`T(t)`$ becomes available. For frequencies below $`f_{}`$ the $`T`$-channel is insensitive to gravitational waves and can be used as an instrument noise monitor. To simplify our analysis we did not use the $`T`$-channel. Including this channel would not have had much effect on the trend of our results. As an example, looking at Figure 3 one can see that for ALIA inclusion of the $`T`$-channel has little to no effect for binaries $`<10^2M_{}`$. Higher mass binaries would have marginally improved signal strength and correspondingly improved resolution of the parameters. Table 1 contains signal-to-noise ratio (SNR) information for the combined $`A`$ and $`E`$ channels, as well as the $`T`$ channel for equal mass binaries as seen by ALIA in the year before coalescence. The SNRs in Table 1 are for the final year before coalescence, separated into the contribution from the first $`51`$ weeks and the contribution from the final week.
In the low frequency limit, the variables $`S_I`$ and $`S_{II}`$ reduce to the independent data channels defined by Cutler Cutler . We use the rigid adiabatic approximation rcp to model the gravitational wave content of the data channels. The instrument noise in the $`A`$ and $`E`$ channels was modeled by PTLA02 ; Nayak
$`S_n(f)=8\mathrm{sin}^2(f/2f_{})[(2+\mathrm{cos}(f/f_{}))S_{\mathrm{pos}}`$
$`+2(3+2\mathrm{cos}(f/f_{})+\mathrm{cos}(2f/f_{})){\displaystyle \frac{S_{\mathrm{accl}}}{(2\pi f)^4}}]/(2L)^2`$
where $`S_{\mathrm{pos}}`$ and $`S_{\mathrm{accl}}`$ are the one-way position and acceleration noise contributions. Table 2 lists the instrument parameters used in our study of the LISA, ALIA and BBO missions.
The $`A`$ and $`E`$ signals are functions of the nine parameters $`\stackrel{}{\lambda }(\theta ,\varphi ,t_c,D_L,\iota ,\psi ,,\mu ,\gamma _{})`$ that describe the source. The source parameters can be recovered from the data channels using a variety of data analysis methods. We will consider maximum likelihood estimation, using the log likelihood function
$$\mathrm{log}(\stackrel{}{\lambda }^{})=(s|h(\stackrel{}{\lambda }^{}))\frac{1}{2}(h(\stackrel{}{\lambda }^{})|h(\stackrel{}{\lambda }^{})),$$
(4)
where $`s`$ denotes the data and $`h(\stackrel{}{\lambda }^{})`$ denotes the search template. We use the standard noise-weighted inner product $`(a|b)`$ summed over the independent data channels:
$$(a|b)=2\underset{A,E}{}_0^{\mathrm{}}\frac{a^{}(f)b(f)+a(f)b^{}(f)}{S_n(f)}𝑑f.$$
(5)
For our work here the limits of the integration are set by the initial and cut-off frequencies of the binary systems. To first order these are given by:
$$f_{\mathrm{initial}}=3.1\left(\frac{10^3M_{}}{(1+z)}\right)^{5/8}\mathrm{mHz},$$
(6)
$$f_{\mathrm{cut}\mathrm{off}}=4.4\left(\frac{10^3M_{}}{M_{total}(1+z)}\right)\mathrm{Hz}.$$
(7)
where $`M_{total}=^{5/2}/\mu ^{3/2}`$ is the total binary mass. Setting $`\stackrel{}{\lambda }^{}`$ equal to the true source parameters yields the SNR
$$\mathrm{SNR}^2=2\mathrm{log}(\stackrel{}{\lambda })=(h(\stackrel{}{\lambda })|h(\stackrel{}{\lambda })).$$
(8)
Here the angle brackets $``$ denote an expectation value. The maximum likelihood estimator $`\stackrel{}{\lambda }_{\mathrm{ML}}`$ is defined by:
$$\frac{\mathrm{log}(\stackrel{}{\lambda }_{\mathrm{ML}})}{\lambda ^i}=0.$$
(9)
The Fisher Information Matrix (FIM), $`\mathrm{\Gamma }`$, is given as the negative of the expectation value of the Hessian evaluated at maximum likelihood:
$$\mathrm{\Gamma }_{ij}=\frac{^2\mathrm{log}(\stackrel{}{\lambda }_{\mathrm{ML}})}{\lambda ^i\lambda ^j}=(h_{,i}|h_{,j}),$$
(10)
where $`h_{,i}\frac{h}{\lambda ^i}`$. For large SNR, the parameter estimation uncertainties, $`\mathrm{\Delta }\lambda ^i`$, will have the Gaussian probability distribution
$$p(\mathrm{\Delta }\lambda ^i)=\sqrt{\frac{\mathrm{det}\mathrm{\Gamma }}{(2\pi )^9}}\mathrm{exp}\left(\frac{1}{2}\mathrm{\Gamma }_{ij}\mathrm{\Delta }\lambda ^i\mathrm{\Delta }\lambda ^j\right),$$
(11)
and variance-covariance matrix
$$\mathrm{\Delta }\lambda ^i\mathrm{\Delta }\lambda ^j=C^{ij}=\left(\mathrm{\Gamma }^1\right)^{ij}.$$
(12)
The uncertainties in each of the parameters are given by $`\mathrm{\Delta }\lambda ^i=(C^{ii})^{1/2}`$ (no summation). One can see from equations (5) and (12) that the parameter uncertainties scale inversely with the SNR.
## III Multiple Constellation Detection
Detectors such as LISA and ALIA determine the positions of gravitational wave sources through both amplitude and frequency modulation. The angular resolution improves with time due to the accumulation of SNR and the synthesis of a long baseline as the detectors move in their orbit. In contrast, a detector array like the BBO has widely separated elements, and thus has a built-in baseline. Adding a second widely separated constellation to LISA or ALIA would increase the SNR by a factor of $`\sqrt{2}`$, but the main gain in angular resolution for transient sources would be due to the built in baseline. The advantage of a multi-element array can be understood from the following toy model. Suppose that we have a gravitational wave of known amplitude and frequency. Neglecting the effects of amplitude modulation we have
$$h=A\mathrm{cos}(2\pi f(t+R\mathrm{sin}(\theta )\mathrm{cos}(2\pi f_mt+\kappa \varphi )))$$
(13)
Here $`f`$ is the source frequency, $`R`$ is the distance from solar barycenter to the guiding center of the constellation, $`f_m=1/\mathrm{year}`$ is the modulation frequency, $`\kappa `$ is the azimuthal location (along Earth’s orbit) of the guiding center, and $`\theta `$ and $`\varphi `$ give the source location on the sky.
With this two parameter ($`\theta `$ and $`\varphi `$) signal we can analytically derive the uncertainty in the solid angle for a source observed for a time, $`T_{obs}`$, by a single constellation
$`\mathrm{\Delta }\mathrm{\Omega }_{\mathrm{single}}`$ $`=`$ $`{\displaystyle \frac{2S_n(f)}{(A\pi fR)^2\mathrm{sin}(2\theta )T_{obs}}}`$ (14)
$`\times {\displaystyle \frac{1}{\sqrt{(1\mathrm{sinc}^2(2\pi f_mT_{obs}))}}}.`$
For small observation times $`\mathrm{\Delta }\mathrm{\Omega }_{\mathrm{single}}`$ scales as $`T_{obs}^2`$, while for large observation times it scales as $`T_{obs}^1`$.
Turning to the dual detector case, we simply add together the FIMs for each individual detector. For two constellations separated by an angle $`\mathrm{\Delta }\kappa `$, this yields a solid angle uncertainty of
$`\mathrm{\Delta }\mathrm{\Omega }_{\mathrm{dual}}`$ $`=`$ $`{\displaystyle \frac{S_n(f)}{(A\pi fR)^2\mathrm{sin}(2\theta )T_{obs}}}`$ (15)
$`\times {\displaystyle \frac{1}{\sqrt{(1\mathrm{sinc}^2(2\pi f_mT_{obs})\mathrm{cos}^2(\mathrm{\Delta }\kappa ))}}}.`$
For small observation times and non-zero $`\mathrm{\Delta }\kappa `$, $`\mathrm{\Delta }\mathrm{\Omega }_{\mathrm{dual}}`$ scales as $`T_{obs}^1`$. In other words, the built in baseline leads to a much improved angular resolution for short observation times. This is very important for coalescing binaries as most of the SNR accumulates in the final days or weeks prior to merger. Note that if the two constellations are co-located, $`\mathrm{\Delta }\kappa =0`$, the uncertainty in the solid angle is reduced by a factor of two relative to the single detector case by virtue of the increased SNR. Also, note that this toy model only includes Doppler modulation, thus the symmetry between $`\mathrm{\Delta }\kappa =0`$ and $`\mathrm{\Delta }\kappa =\pi `$. Including the amplitude modulation breaks this symmetry.
## IV Results
The data shown here is for sources at a redshift of $`z=1`$ for ALIA, ALIAS, LISA, and LISAS and redshift $`z=3`$ for BBO, and BBO Star. These correspond to luminosity distances of 6.63 Gpc and 25.8 Gpc, respectively, using the best fit WMAP cosmologydavid . Each binary system is observed for the last year before coalescence. Each data point is distilled from $`10^5`$ random samples of $`\theta `$, $`\varphi `$, $`\iota `$, $`\psi `$, and $`\gamma _o`$. For $`t_c`$, $`D_L`$, $``$, and $`\mu `$ we use logarithmic derivatives so that the uncertainties listed have been scaled by the value of the parameters. The uncertainty in sky location is simply the root of the solid angle uncertainty ($`\sqrt{\mathrm{\Delta }\mathrm{\Omega }}`$). The remaining angular parameters, $`\iota `$, $`\psi `$, and $`\gamma _o`$, have not been scaled.
The angular variables were chosen by using a Monte Carlo method. The values for $`\mathrm{cos}(\theta )`$ and $`\mathrm{cos}(\iota )`$ were chosen from a random draw on $`[1,1]`$. Values for $`\varphi `$, $`\psi `$, and $`\gamma _o`$ were each chosen from random draws on $`[0,2\pi ]`$. The parameters $``$ and $`D_L`$ were set for each Monte Carlo run (though they changed between runs). Time to coalescence, $`t_c`$, was set to $`1`$ year plus a small offset so that during the year of observation the binary did not reach a relativistic regime that would not be properly modeled by the $`2\mathrm{P}\mathrm{N}`$ approximations used.
The mean SNRs quoted in this paper are calculated by taking the square root of the average of the squares of the individual SNRs. For our analysis, positive detection will be restricted to SNRs above $`5`$.
### IV.1 Results for ALIA
Table 3 summarizes the medians and means of parameter uncertainties for detections by ALIA of equal mass binaries with masses of $`10^2M_{}`$, $`10^3M_{}`$, $`10^4M_{}`$, and $`10^5M_{}`$. The SNRs for this range of masses shows that ALIA should get positive detection of $`99^+\%`$ of coalescing IMBHs located at $`z=1`$ or closer, which would provide good information on the coalescence rates. The great precision in the measurement of $``$ and $`\mu `$ will provide a clear picture of the constituent masses of the binary systems. Furthermore, the sub-degree precision in the sky location, combined with luminosity distances known to a few percent, will facilitate the construction of a three dimensional distribution of IMBHs with which to test theoretical predictions (see Ref. 3d for a related discussion concering LISA and supermassive black holes).
The uncertainties in $`t_c`$ a month from coalescence (using $`11`$ months of data) will be on the order of a few minutes. This will provide warning time for ground-based gravitational wave detectors, as well as other systems (telescopes, neutrino detectors, etc.) to gather as much and as varied data as possible about the coalescence.
The shapes of the histograms shown in Figure 4, which are from the $`10^3M_{}`$ data, are representative of the histograms for each of the detectors covered in this work. Note that for $`\iota `$, $`\psi `$, and $`\gamma _o`$, the tails of the histograms run far beyond the range of the plots shown, raising the values of the means considerably above their respective median values. For example, while the median value of $`\gamma _o`$ for equal mass binaries with masses $`10^2M_{}`$ is $`12.2^{}`$ its mean value is $`2185^{}`$, which is well beyond the $`[0,2\pi ]`$ range of $`\gamma _o`$. Uncertainty ranges that are larger than the possible range of the parameter tell us that the parameter is indeterminate. When uncertainty ranges exceed the possible range of a parameter, one may drop that parameter from the FIM analysis. For our analysis, we did not discard any parameters. As an example of how this affects the analysis, consider the $`10^2M_{}`$ study where the mean value of the uncertainties that lie in the $`[0,2\pi ]`$ range for $`\gamma _o`$ would be $`36.3^{}`$, while its median would be $`9.79^{}`$. However, for $`16.9\%`$ of the binaries, the $`\gamma _o`$ parameter would be indeterminate.
Figure 5 shows the SNR histograms for $`10M_{}`$ equal mass binaries for ALIA (and ALIAS). As was expected from Figure 2, the SNR values for ALIA are low, with nearly $`60\%`$ too low for a positive detection. Of note, though, is that the SNR scales roughly as the luminosity distance. This suggests that ALIA should be able to detect nearly all $`10M_{}`$ binaries with luminosity distances less than $`2`$ Gpc. For binaries with masses beyond $`15M_{}`$ the low end of the range of SNRs is above $`5`$, meaning that equal mass binaries with masses in the range of IMBHs should be detectable out to $`z=1`$.
### IV.2 Results for ALIAS
The main purpose of ALIA is gathering information about IMBHs; our data shows it is capable of doing this with some success. A more accurate IMBH census could be derived from a dual constellation version of ALIA that we call ALIA in Stereo or ALIAS. Each component of the ALIAS constellation would be offset from the Earth by $`20^{}`$, one in an Earth-trailing and the other in an Earth-leading orbit, giving a $`40^{}`$ separation in order to provide increased parameter resolution in the IMBH range.
Table 4 summarizes the medians and means of the parameter resolutions that could be achieved by the ALIAS mission. Results are given for equal mass binaries with masses of $`10^2M_{}`$, $`10^3M_{}`$, $`10^4M_{}`$, and $`10^5M_{}`$. As expected, the mean SNR increases by $`\sqrt{2}`$ relative to ALIA. The improvements in parameter resolution are, however, considerably larger. At the upper end of the IMBH mass range the angular resolution improves by a factor of $`90`$ and the luminosity distance resolution improves by a factor of $`23`$.
For masses below $`10^2M_{}`$, ALIAS provides roughly a factor of $`\sqrt{2}`$ increase in the median values of each of the parameter uncertainties. This can be seen in the $`10^2M_{}`$ data in Tables 3 and 4, as well as in Figure 6 and Figure 7, which plot the uncertainty in the sky location and luminosity distance, respectively, against the chirp mass of the binaries. The trend shown in these figures holds true for $`D_L`$, $`\iota `$, and $`\psi `$ in that mass range, while the increases in $`t_c`$, $`\mu `$, $`\gamma _o`$, and $``$ are slightly larger, $`1.55`$ better than the values for ALIA. However, the median SNRs are below $`5`$ for equal mass binaries with chirp masses below $`8M_{}`$ for ALIAS (below $`10M_{}`$ for ALIA).
Figure 6 also shows that the benefits of a dual constellation becomes even more significant when the final chirp of the binary occurs above the detector noise (see Figure 2). The main improvement is in the angular resolution, but the decreased covariances between the sky location and other parameters lead to improved measurements of $`t_c`$, $`D_L`$, $`\iota `$, and $`\psi `$ at the upper end of the IMBH mass range. The resolution of these parameters improves (relative to ALIA) by a factors of $`33`$, $`23`$, $`18`$, and $`21`$, respectively. The parameters $``$, $`\mu `$, and $`\gamma _o`$ only see a factor of $`2`$ improvement in resolution.
Comparing Figure 7 with Figure 6, one can see that the increased resolution in the luminosity distance corresponds with the increased precision in sky location due to their large covariance. A similar effect was seen in the work of Hughes and Holz Hughes\_Holz , where the addition of electromagnetic information to fix the sky location of a BH merger resulted in a marked increase in resolution of the luminosity distance. In our case, the degeneracy is broken by the improved angular resolution afforded by the baseline between the ALIAS detectors. These improvements in sky location and luminosity distance resolution will provide an even more detailed three dimensional distribution of the IMBHs than that provided by ALIA. Increasing the $`3D`$ resolution of the distribution of IMBHs by roughly a factor of $`6`$ on the low end of the mass range ($`50M_{}`$) up to $`200,000`$ on the high end of the mass range ($`50,000M_{}`$), would provide much more detailed information compared to ALIA, with which to test theoretical predictions about IMBHs.
Figure 5 shows the SNR histogram for $`10M_{}`$ equal mass binaries for ALIA and ALIAS. While more than two-thirds of the sources at $`z=1`$ would be positively detected, nearly one-third would be missed. However, these results suggest that ALIAS will be able to detect nearly all of the $`10M_{}`$ binaries with luminosity distances less than $`3`$ Gpc. This is roughly in keeping with the general trend of ALIAS’s capabilities compared to ALIA.
From Figure 2 we saw that for masses below $`10M_{}`$, the source’s signal lies below the sensitivity curve for the week prior to coalescence. For larger masses, more of the binary’s final chirp will be detectable. This is why the angular resolution of ALIAS becomes so much better than ALIA’s at the upper end of the IMBH mass range. For masses above $`10^5M_{}`$ the final chirp occurs before the “sweet spot” of the sensitivity curve, which diminishes ALIAS’s advantage over ALIA.
### IV.3 Results for ALIA & ALIAS for unequal mass binaries
Equal mass binaries are the easiest to study, but they can give an overly optimistic picture of the instrument capabilities as they yield the smallest parameter uncertainties. To study this bias and provide a more realistic picture of the capabilities of ALIA and ALIAS, we now consider the case where the masses of each component in the binary are drawn from logarithmic distributions in the range $`110^8M_{}`$. The results are still presented as a function of chirp mass, but the mass ratios reflect the underlying mass distribution.
Figure 8 shows the plot of sky location uncertainty against chirp mass. The uncertainties in sky location for masses above $`10^2M_{}`$ are larger than those shown in Figure 6. The increase in sky angle uncertainty relative to the equal mass case was a factor of $`2`$ for ALIA. Similarly, other parameters show a factor of $`3`$ increase in their uncertainties relative to the equal mass study.
However, ALIAS still maintains an advantage in locating binaries. In fact, for chirp masses below $`500M_{}`$ the increase in precision for ALIAS over ALIA in locating unequal binaries is nearly the same as it was for equal mass binaries, a factor of roughly $`\sqrt{2}`$. At higher chirp masses we see less of a difference in angular resolution than in the equal mass study (factors of $`9`$ compared to $`90`$). This trend holds for $`t_c`$, $`D_L`$, $`\iota `$, and $`\psi `$, which show maximum increases in precision $`4`$ for ALIAS over ALIA, while $``$, $`\mu `$, and $`\gamma _o`$ show maximum increase in precision $`2`$ for ALIAS over ALIA (as they did for equal mass binaries). Thus the increased resolution in the $`3D`$ distribution of IMBHs for ALIAS over ALIA for unequal mass binaries ranges from a factor of $`6`$ up to $`300`$.
### IV.4 Results for LISA and LISAS
Figure 9 compares the sky location uncertainties for the LISA mission to the LISAS mission. As was seen with ALIA and ALIAS, the addition of a second constellation to LISA provides a marked increase in parameter resolution. The increases in parameter resolution for LISAS over LISA are similar to those found between ALIAS and ALIA. The angular resolution showed a maximum improvement of $`25`$ just above $`10^5M_{}`$. Also, a lesser benefit than that shown in Figure 9 occurs in the $`t_c`$, $`D_L`$, $`\iota `$, and $`\psi `$ parameters (with factors of maximum increase $`12`$ just above $`10^5M_{}`$), while $``$, $`\mu `$, and $`\gamma _o`$ show only a modest maximum improvement by a factor of $`2`$. Our results are in the same range as those found by Seto naoki for sources detected by LISA that have undergone strong gravitational lensing. In that case the extended baseline was provided by the time delay in the arrival of the signals, which effectively turned a single LISA detector into a LISAS system. Seto also gave brief consideration to the performance of a dual LISA mission, and found the ratio of angular resolution between LISAS and LISA to be larger than that seen in our simulations. The discrepancy can be traced to Seto including a large galactic confusion background which limits the time that the sources are in-band, thus magnifying the advantage of the dual configuration.
The shift in the placement of the minimum uncertainties, compared to ALIA and ALIAS, can be understood from Figure 2. The signal from a $`10^3M_{}`$ equal mass binary is passing below the LISA sensitivity curve in the last week before coalescence, and before the final chirp of the binary. This is evident in Figure 9, where the angular resolution of the single and dual configurations is similar for masses below $`10^3M_{}`$. The sweet-spot of the LISA sensitivity curve is at $`5`$ mHz, which is where $`10^5M_{}`$ binaries experience their final chirp. Above $`10^5M_{}`$, less of the final chirp is occurs near the sweet spot, and the difference in the angular resolution diminishes. Figure 10 shows how this shift corresponds precisely to the shift in the maximum SNR of the LISA and LISAS missions. The plot shows how the increase in the mean SNR ($`\sqrt{2}`$) is uniform across the range of chirp masses. Also, it shows again how the resolution in $`D_L`$ improves with the addition of more information in the form of increased angular resolution.
### IV.5 Results for BBO
Table 5 summarizes the medians and means of the parameter uncertainties for detections of equal mass binaries with masses of $`1.4M_{}`$, $`10M_{}`$, and $`10^2M_{}`$ for the BBO. Histograms for this data have the same general shapes as those shown in Figure 4 and so are not shown.
As can be seen from the data, the BBO is an extremely sensitive detector. The SNRs show that if the BBO meets design specifications there will be positive detection of coalescing stellar mass binaries out to (indeed well beyond) a redshift of $`z=3`$. Similarly, the BBO will provide a precise picture of coalescence rates, including how these rates relate to luminosity distance. With a combination of sky location and luminosity distance the BBO should be able to pick out the host galaxy of the majority of coalescing binaries. Data taken by the BBO up to a month before coalescence (using $`11`$ months of data) will determine $`t_c`$ to within seconds, again providing ample warning time for other detectors to gather data on the coalescence.
### IV.6 Results for BBO Star
As can be seen in Figure 2, the final chirp of solar mass or NS coalescing binaries occurs below the noise level of the BBO. Thus the outrigger constellations do not provide an extended baseline for the last few days of NS chirp. Positive detection of NS binaries at $`z=3`$ can be accomplished with the initial deployment of the two constellations that make up the star, while still providing the cross-correlating needed to detect the GWB.
Table 6 summarizes the medians and means of the parameter uncertainties for the detection of equal mass binaries with masses of $`1.4M_{}`$, $`10M_{}`$, and $`10^2M_{}`$ for the BBO Star. Histograms for this data have the same general shapes as those shown in Figure 4 and so are not shown.
Figure 11 plots the uncertainty in sky location against chirp mass for the BBO and BBO Star. As can be seen, the effect of the outrigger constellations is significant, providing for increased parameter resolution for the full BBO throughout the range of chirp masses shown, with a maximum increase of $`3700`$ around $`200M_{}`$. Similar to the ALIA/ALIAS comparison this extra precision increases until all of the final chirp lies above the sensitivity curve (see Figure 2), and begins to decrease as less and less of the final chirp occurs in the frequency range of the sweet-spot of the sensitivity curve (which for BBO is $`1`$ mHz). The improvement in the mean SNR for BBO over BBO Star is the expected $`\sqrt{2}`$. Also, the improvement in the resolution of $`t_c`$ ($`3550`$) is comparable to that seen in the angular resolution. More modest increases occur for $`D_L`$, $`\iota `$, and $`\psi `$, which increase by factors of $`52`$, $`35`$, and $`47`$, respectively. Only slight ($`2.5`$)increases are seen for $``$, $`\mu `$, and $`\gamma _o`$.
The SNRs from BBO Star are sufficient for positive detection of binaries with constituents less than one solar mass out to, and beyond, $`z=3`$. While the full BBO offers considerable adavantages for doing precision gravitational wave astronomy, BBO Star could fulfill the main science objective of detecting the cosmic gravitational wave background while still providing useful information about binary populations.
## V Conclusion/Discussion
While our survey is by no means comprehensive, it has helped to map out the science that can be done with the ALIA and BBO missions. We have shown that ALIAS, a modest extension to the ALIA mission, would be able to return a far more accurate census of the IMBH population. In addition we have shown that a similar extension to LISA would greatly improve its ability to locate the host galaxies of coalescing binaries. On the other hand, if we are willing to give up some of the precision astronomy offered by the full BBO, we could get by with just the first phase of the BBO deployment. The BBO Star configuration could satisfy the primary goal of detecting the CGB, while still providing a detailed binary census.
###### Acknowledgements.
This work was supported by NASA though the BBO Mission Concept Study led by Sterl Phinney.
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# Fourier Resolved Spectroscopy of the XMM-Newton Observations of MCG -06-30-15
## 1 Introduction
The “standard” picture of the Active Galactic Nuclei (AGN) accretion disk arrangement consists of a geometrically thin, optically thick disk (that produces the quasithermal optical–UV feature known as the “Blue Bump”) supplemented by an overlying, hot ($`10^8`$ K) corona responsible for the observed X-ray emission. Both these components are presumed to extend to the innermost stable particle orbits associated with the accreting black hole. This latter fact led to the suggestion that studies of the spectral and timing characteristics of these components could provide a probe of the regime of strong gravitational physics.
It has been argued that the reprocessing of X-ray radiation on the much cooler accretion disk would produce a relativistically broadened, asymmetric, Fe K$`\alpha `$ fluorescence feature at 6.4 keV (Fabian et al. 1989; Stella 1990) and a spectral hardening of the spectrum at energies $`E>10`$ keV due to reflection of the coronal X-rays by neutral matter on the accretion disk surface. It was further argued that the precise shape, EW and variability properties of this feature could allow the mapping of the space–time geometry in the black hole vicinity (Reynolds et al. 1999) and/or the geometric arrangement of the disk and the X-ray emitting plasma (Nayakshin & Kazanas 2001).
AGN observations appear to corroborate these considerations: A broad feature at the correct energy was indeed identified in the ASCA spectra of many AGN (Nandra et al. 1997b). One of the broadest such asymmetric emission features was detected in the energy spectrum of the Seyfert 1 galaxy MCG -06-30-15 (Tanaka et al. 1995). The presence of this emission feature has been confirmed by repeated observations of the source by all recent X-ray missions, including Chandra (Lee et al. 2002) and XMM-Newton (Wilms et al. 2001; Fabian et al. 2002) and thought to suggest, on the basis of its width, the presence of an extreme Kerr geometry for the accreting black hole. At energies below 2 keV MCG -06-30-15 shows great spectral complexity attributed to absorption by partially ionized material and possibly also dust (Turner et al. 2003 and references therein); alternatively, these same features have also been interpreted by some authors due to relativistically broadened soft X-ray emission features (Branduardi-Raymont et al. 2001; Sako et al. 2003).
It should be born in mind though that, however convincing, spectral features by themselves generally provide information only about the properties of the emitting (or absorbing) plasma along the observer’s line of sight, typically just its column density. However, the size or the density of the plasma, parameters needed to determine its dynamics, require additional independent information that, in the absence of sufficient spatial resolution, is provided by the sources’ variability (Kazanas, Hua & Titarchuk 1997). Hence, a more comprehensive approach should involve the use of both spectral and timing information in a combined analysis. A novel approach in this direction has been taken by Revnivtsev, Gilfanov & Churazov (1999) who used the combined RXTE variability - spectral information to produce the energy spectra of Cygnus X-1 in its hard state for different Fourier frequencies. This study showed that the Cyg X-1 spectra depend strongly on the Fourier frequency, with the Fe K$`\alpha `$ line and the reflection component becoming increasingly visible with decreasing frequency. Very similar results were obtained by the same authors for the Galactic source GX 339-4 (Revnivtsev, Gilfanov & Churazov 2001), while Cyg X-1 in its soft state showed much smaller dependence of its spectrum and Fe K$`\alpha `$ line on Fourier frequency (Revnivtsev, Gilfanov & Churazov 2000).
In the present paper we apply the method used by Revnivtsev et al. (1999) to MCG -06-30-15. We chose this source because, apart from being considered as the AGN with the archetypal Fe K$`\alpha `$ profile, it shows large amplitude variations on both short and long time scales (e.g. McHardy et al. 2005), it is bright in X-rays and has been the target of numerous X-ray observations. Of particular interest amongst them are those by XMM-Newton due to the large area of its instruments, which offer high quality, low noise spectra, albeit at energies less than 10 keV. In §2 and §3 we discuss the data sets and the methodology used and we present our results in the form of energy spectra at three different Fourier frequencies. In §4 our results are discussed briefly in the context of the standard and other models presently in the literature.
## 2 Data Reduction and Analysis Method
MCG -06-30-15 has been observed with XMM-Newton twice. The first observation was performed on 2000 July 10-11 (revolution 108) and the second on 2001 July 31 - August 5 (revolutions 301-303). In this work we use data from the EPIC-PN detector only. In both cases the source was observed on-axis and the PN camera was operated in small window mode, with medium filter applied. We processed the data using SAS v6.0.0. Source data (single and double pixel events, i.e. patterns $`04`$) were extracted from circular regions of radius $`40^{\prime \prime }`$, and background events from a source free area 2 times larger than the source extraction area. The background was in general low and stable throughout the observations, with the exception of the final few ks of each revolution where the background count rate increases. We kept 83 and 320 ks of “good” exposure time data from the 2000 and 2001 observations, respectively.
Fig. 1 shows the 0.2-10 keV, background subtracted, 400-s binned, EPIC-PN light curves. The source shows large amplitude variations at all time scales. The raw energy spectrum, using data from both observations, is shown in the upper panel of Fig. 2. The energy resolution is significantly reduced for reasons that we explain in the following section. The photon count spectrum is estimated in 22 energy bands only: we consider 6 bands from 0.2 to 0.8 keV with $`\mathrm{\Delta }E=0.1`$ keV, the bands $`0.81`$, $`11.3`$, $`1.31.6`$ and $`1.62`$ keV, 9 bands from 2 to 7 keV with $`\mathrm{\Delta }E=0.5`$ keV, and finally the bands $`78`$ and $`810`$ keV.
The main features are clearly evident even in this low resolution spectrum. The solid line in the upper panel of Fig. 2 shows a power law model spectrum of $`\mathrm{\Gamma }=2`$ and Galactic absorption only ($`N_H=4.06\times 10^{20}`$cm<sup>-2</sup>; Elvis, Wilkes, & Lockman 1989). The model normalization is adjusted to produce total number of counts equal to those of the observed spectra (this is the case for the normalization of all model spectra presented in the rest of the paper). In the lower panel of Fig. 2 we show a plot of the data over the model ratio (filled circles). The significantly asymmetric, broad iron line in the energy band $`57`$ keV and strong residuals in the soft X-ray band can be clearly seen.
### 2.1 Fourier Resolved Spectroscopy
We now discuss briefly the technique of “Fourier resolved spectroscopy” introduced by Revnivtsev et al. (1999). Suppose we have light curves of a given source at different energy bands $`E_j`$, $`j=1,2,\mathrm{}M`$. Let us denote with $`x(t_i,E_j)`$ the observed count rate at time $`t_i`$ in the energy band $`E_j`$ ($`i=\mathrm{\Delta }t,2\mathrm{\Delta }t,\mathrm{},N\mathrm{\Delta }t`$, $`N`$ is the total number of points, and $`N\mathrm{\Delta }t=T`$ is the length of the light curve). The power spectrum (PSD) of this light curve can be estimated as follows:
$$P(\nu _k,E_j)=\frac{2\mathrm{\Delta }t}{N}|a_k|^2,$$
(1)
in units of (counts s$`{}_{}{}^{1})^2`$ Hz<sup>-1</sup>, where $`\nu _k=k/N\mathrm{\Delta }t`$ $`(k=1,2,\mathrm{},(N1)/2)`$ and $`a_k=_ix(t_i,E_j)e^{i2\pi \nu _kt_i}`$. The finite Fourier representation of the observed time series $`x(t_i,E_j)`$ is equal to the sum of sinusoidal terms (“variability components”) with frequencies $`\nu _k`$. Their amplitudes, $`A_k`$, are related to $`|a_k|`$ through the relation $`A_k=(2/N)|a_k|`$. Using equation (1), this relation can be rewritten as follows:
$$A(\nu _k,E_j)=\sqrt{\frac{2P(\nu _k,E_j)}{N\mathrm{\Delta }t}}\mathrm{counts}\mathrm{s}^1.$$
(2)
$`A(\nu _k,E_j)`$, viewed as a function of $`E_j`$, represents the energy spectrum of the component with frequency $`\nu _k`$.
Although MCG -06-30-15 is a bright Seyfert 1 galaxy it is not possible to study its energy dependent variability behavior using the full energy resolution offered by XMM-Newton. Instead, we have split the 0.2-10 keV energy range in the 22 bands referred to in the previous section. These bands were so chosen to: a) be sufficiently broad to yield light curves of reasonably high signal-to-noise for an accurate determination of the power spectrum, and b) their number in the traditional “soft” ($`<2`$ keV) and “hard” ($`>2`$ keV) X-ray channels be roughly equal. The two XMM-Newton observations resulted in 5 segments during which the source was observed almost continuously for $`T32120`$ ks (see Fig. 1). For each segment we constructed background-subtracted light curves in the 22 energy bands using a bin size of $`\mathrm{\Delta }t=100`$ s. In this way a total of 105 light curves were produced. We used equation (1) to estimate their power spectrum, after subtracting the associated (constant) Poisson power. In principle, using equation (2), we can now compute the amplitudes $`A_l(\nu _k,E_j),l=1,\mathrm{},5`$. The mean of the five $`A_l(\nu _k,E_j)`$ estimates at each $`\nu _k`$, viewed as a a function of $`E_j`$, can constitute the “low-resolution” energy spectrum of the variability components with frequency $`\nu _k`$.
However, the energy spectra of the individual $`\nu _k`$’s turned out to be very noisy. In fact, at high frequencies, after the subtraction of the Poisson noise component, many of the $`P(\nu _k,E_j)`$’s are negative, thus preventing the estimation of the respective $`A(\nu _k,E_j)`$. For this reason, we decided to consider three broad frequency bands : a) $`8.3\times 10^610^4`$ Hz, b) $`13\times 10^4`$ Hz, and c) $`37\times 10^4`$ Hz (hereafter “low/LF”, “medium/MF” and “high-frequency/HF” range, respectively). These frequency bands correspond to time scales of $`0.121.4`$ days, $`310`$ ksec, and $`1.53`$ ksec, respectively. There are $`N_{\mathrm{PSD}}=39`$, 81 and 161 $`P(\nu _k,E_j)`$ estimates in each frequency bin. We computed their mean, $`\overline{P(E_j)}=_{l,k}P(\nu _k,E_j)/N_{\mathrm{PSD}}`$, and used equation (2) to estimate the average amplitude, $`\overline{A(E_j)}`$, of the variability components in each frequency band. These values comprise our final estimate of the energy spectrum of the low, medium and high-frequency variability components.
## 3 Results
The raw, count energy spectra of the LF, MF and HF components (i.e. $`\overline{A(E_j)}`$ divided by $`\mathrm{\Delta }E`$) are also shown in the upper panel of Fig. 2 ( given respectively by the open circles, filled triangles, and open squares). In order to gain some insight into their broad-band shape, we divided them by the raw total energy spectrum of the source. These ratios (shown in the lower panel of Fig. 2) should factor out the broad instrumental response that distorts the spectra and, as a result, should give a better view of their true shape, albeit only in relation to that of our entire data set.
At energies $`E<3`$ keV the LF and MF normalized spectra are consistent with constant, yielding respectively $`\chi ^2=14.4`$ and 8.7/11 dof when fit as such. This indicates that the spectra of the LF and MF components are not significantly different in shape from the total energy spectrum of the source. At higher energies, the LF and HF normalized spectra suggest that the respective frequency spectra are significantly softer and harder, respectively, than the total energy spectrum. In fact, the HF normalized spectrum suggests that the spectrum of the high frequency components is systematically harder in the entire $`0.210`$ keV band. Also, a clear decrease of the $`57`$ keV band flux is observed in the LF normalized spectrum. Such a feature should be expected if there was no significant iron line emission in the LF spectrum. A similar feature may also be present in the other two normalized spectra, although not as clear.
Irrespective of their intrinsic shape, the Fourier count spectra can be used to compute the amplitude of the various variability components in any given energy band. The integral of the Fourier spectrum over a particular energy band (say from $`E_1`$ to $`E_2`$) yields the contribution of the variability component with frequency $`\nu _k`$ to the variance of the light curve in this energy band. Consequently, the ratio of this integral over the integral of the total energy spectrum (in the same band) corresponds to the fractional root mean square amplitude ($`f_{rms})`$ of this component in this energy band.
In our case, we computed the sums of $`\overline{A(E_j)}`$ over the soft and hard energy bands for the three Fourier spectra, and divided them by the integral of the total spectrum in the same energy bands. In the soft band, we find $`f_{rms}=12.4\pm 0.5`$%, $`3.05\pm 0.06`$%, and $`1.30\pm 0.05`$% for the LF, MF and HF components, respectively. The respective values in the hard band are: $`9.7\pm 0.6`$%, $`2.90\pm 0.06`$%, and $`1.74\pm 0.05`$%. These results show that on long time scales ($`0.121.4`$ days) the amplitude of the observed variations in the soft band is larger than that in the hard band. The opposite holds true on short time scales ($`1.53`$ ksec). On intermediate time scales ($`310`$ ksec), the soft and hard band variability amplitudes are comparable.
### 3.1 Model fits to the high energy band
In order to get additional insight into the broad band energy spectra of the different frequency components we produced count spectra for power law models with only Galactic absorption and slope increasing from $`\mathrm{\Gamma }=1.5`$ to $`\mathrm{\Gamma }=2.5`$ with a step of $`\mathrm{\Delta }\mathrm{\Gamma }=0.05`$. The power law normalizations were adjusted so that the total counts in the model and the LF, MF and HF energy spectra be equal (the appropriate auxiliary files for the creation of the model spectra were produced using the RMFGEN and ARFGEN tasks of SAS). The upper panel in Fig. 3 shows the LF, MF and HF energy spectra along with the model spectrum that fits each “best” at $`E>3`$ keV, i.e. gives the minimum sum of squared residuals (weighted by the data errors square) among the model spectra considered. The lower three panels in the same Figure show the ratio of the data to the best fitting model. Our model fitting results show that the spectral slope decreases from the LF to the HF band. We find that above 3 keV the LF, MF and HF spectra are well described by a power law model with $`\mathrm{\Gamma }=2.15,1.9`$, and 1.65, respectively. When we fit the data/model ratios above 3 keV with a constant of value 1 we find $`\chi ^2=11.6,15.1,`$ and 11.2/10 dof, respectively. We conclude that there is no strong evidence for the presence of an iron line in any of the three Fourier spectra. Furthermore, although the hard power law component is variable on all time scales (i.e. all three Fourier frequency bands), its slope does not remain constant. For example, there is a significant difference of $`\mathrm{\Delta }\mathrm{\Gamma }=0.5`$ in the hard band power law slopes of the LF and HF components. In other words, we find that the high frequency component has harder spectrum than that of the lower frequency one.
### 3.2 Model fits to the full band energy spectrum
At energies below $`2`$ keV, a simple power law model with Galactic absorption does not fit well the spectra of any of the three Fourier frequency bands. In all cases, the plot of the data/model ratio as a function of energy is qualitatively similar to the respective plot in the case of the power law model fit to the total energy spectrum (shown on the top of the lower panel in Fig. 2): we observe a strong, broad absorption feature at energies $`0.72`$ keV, and a broad “hump” at lower energies, whose amplitude appears to increase with increasing frequency. We modeled the soft X-ray band spectrum both for the entire data set and the three Fourier components as follows.
Since we have undersampled the energy resolution of the EPIC-PN data, it is possible to model the soft X-ray complex features considering simple, phenomenological components such as power laws and absorption edges (as opposed to constructing physical models of the absorber and the emitting source). To this end, following Turner et al. (2003), we first considered the entire energy spectrum being the sum of two power laws, one for each of the soft and hard bands, joining at an energy $`E_{\mathrm{break}}=1`$ keV and with $`\mathrm{\Gamma }_{\mathrm{soft}}>\mathrm{\Gamma }_{\mathrm{hard}}`$. Furthermore, in addition to Galactic neutral absorption, we also considered the possibility of intrinsic cold absorption as well. We produced count spectra keeping $`\mathrm{\Gamma }_{\mathrm{hard}}`$ fixed to the value 2, and $`\mathrm{\Gamma }_{\mathrm{soft}}`$ increasing from 2.05 to 3 in steps of $`\mathrm{\Delta }\mathrm{\Gamma }=0.05`$. In each step, we produced three different spectra with $`N_{\mathrm{H},\mathrm{intr}}=10^{20},10^{21},`$ and $`10^{22}`$ cm<sup>-2</sup>. We found that the model spectra with $`\mathrm{\Gamma }_{\mathrm{soft}}2.42.6`$ and $`N_{\mathrm{H},\mathrm{intr}}=10^{20}`$ cm<sup>-2</sup> agree “best” with the observed spectrum.
We then produced a new set of model spectra with $`\mathrm{\Gamma }_{\mathrm{soft}}=2.42.6`$, and $`N_{\mathrm{H},\mathrm{intr}}=10^{20}10^{21}`$ cm<sup>-2</sup> (using $`\mathrm{\Delta }\mathrm{\Gamma }=0.05`$ and $`\mathrm{\Delta }N_{\mathrm{H},\mathrm{intr}}=10^{20}`$ cm<sup>-2</sup>, respectively). Following Wilms et al. (2001), at each spectrum corresponding to a given pair of $`(\mathrm{\Gamma }_{\mathrm{soft}},N_{\mathrm{H},\mathrm{intr}})`$ values we also added absorption edges at 0.74, 1 and 1.36 keV (these threshold energies correspond with the expected absorption edges of OVII, NeIX/MgX, and NeX). As for the optical depth, we used values of $`\tau =0.51`$ for the first edge, and $`\tau =00.5`$ in the other two cases, with a step of $`\mathrm{\Delta }\tau =0.1`$. We then compared each $`(\mathrm{\Gamma }_{\mathrm{soft}},N_{\mathrm{H},\mathrm{intr}},\tau _{0.74keV},\tau _{1\mathrm{k}\mathrm{e}\mathrm{V}},\tau _{1.36\mathrm{keV}})`$ model with the observed spectrum in order to find the one with the minimum $`\chi ^2`$.
The filled circles on the upper panel of Fig. 4 show the total energy spectrum of the source, together with our “best fitting model”: $`\mathrm{\Gamma }_{\mathrm{soft}}=2.5,N_{\mathrm{H},\mathrm{intr}}=3\times 10^{20}`$ cm<sup>-2</sup>, and $`\tau _{0.74\mathrm{keV}}=0.7,\tau _{1\mathrm{k}\mathrm{e}\mathrm{V}}=0,`$ and $`\tau _{1.36\mathrm{keV}}=0.2`$. The data/model ratio is plotted on the second panel of the same Figure. The model describes rather well the overall shape of the spectrum, except in the vicinity of the Fe K$`\alpha `$ line and the $`0.30.6`$ energy band, where a strong ‘hump’ is evident. As discussed earlier, this feature has been attributed to relativistically broadened K$`\alpha `$ emission from oxygen (Branduardi-Raymont et al. 2001, Sako et al. 2003). Another possibility is that the gas responsible for the extra absorption within MCG -06-30-15 may not be cold but partially ionized. In this case, the depth of the neutral oxygen edge at $`0.54`$ will be smaller, reducing the discrepancy between model and observed spectra in this energy band (Turner et al. 2003).
Based on the results of the best fit procedure used to model the total spectrum, we examined whether a similar procedure could provide acceptable fits to the spectra of the individual Fourier components as well. The open squares in the upper panel of Fig. 4 correspond to the HF Fourier spectrum while the solid line shows the best fitting “broken power law with intrinsic absorption and two edges” model. The model parameter values are listed in Table 1. The bottom panel in the same Figure shows the data/model ratio. Apart from the point in the $`0.50.55`$ keV band, the model fits the spectrum quite well ($`\chi ^2=24.6/21`$ dof). A similar model (without the presence of an edge at 1.36 keV) fits reasonably well the LF spectrum as well ($`\chi ^2=17.4/21`$ dof when we omit the point at $`0.5`$ keV; the normalized residuals in this case are plotted in the third panel from top in Fig. 4, and the best fitting model parameters are listed in Table 1). In the case of the MF spectrum, the addition of an extra edge at 0.85 keV (the threshold energy of the O VIII edge) is necessary to provide a good fit ($`\chi _2=26.7/21`$ dof, if the $`0.5`$ point is again not taken into account; see second panel from bottom in Fig. 4 for the data/model ratio plot in this case, and Table 1 for the model fitting results).
In summary, the overall spectrum of the source is well described by two power law components (one for the low energy and one for the high energy sections of the spectrum), intrinsic absorption, as well as two absorption edges. These distinct power law components are present in the spectra of the individual Fourier components as well, indicating that they are variable. However, contrary to the hard band power law, $`\mathrm{\Gamma }_{\mathrm{soft}}2.5`$, independent of the Fourier frequency, suggesting a different physical origin between these two components. Obviously, this result is somehow model dependent. Since the broken power law is a simple phenomenological model, there is no reason to exclude an alternative parametrization of the observed spectra, which may lead to different results, regarding the variability properties of the soft excess.
Finally, absorption features are evident in the spectra of all three Fourier components. The $`0.7`$ keV absorption edge appears to be equally strong (with $`\tau 0.7`$) in all Fourier spectra. An edge at $`1.36`$ keV ($`\tau 0.20.3`$) is evident in the MF and HF spectra, and a third edge (at $`0.85`$ keV; $`\tau 0.25`$) is present in the MF spectrum only. The differences in the presence of edges in the three Fourier spectra could be due to the fact that the MF spectrum has the highest signal to noise ratio among them. For example, the addition of two edges at 0.85 and 1.36 keV (with $`\tau =0.25`$ and 0.2, respectively) in the LF best fitting model is acceptable at a confidence level of $`10`$%, while the addition of a 0.85 keV edge ($`\tau =0.25`$) in the HF best model fitting spectrum provides an acceptable fit at the $`5`$% level.
## 4 Discussion, Conclusions
We have applied the Fourier frequency resolved spectral analysis method to the XMM-Newton data of MCG -06-30-15. This work represents the first application of this method to the X-ray data of an AGN. Our main results are the following: a) The soft and hard bands of the spectrum of MCG -06-30-15 exhibit different RMS variability at different Fourier frequencies (time scales). b) Both the soft and hard band power law components are variable on all time scales. However, while the hard band power law becomes progressively steeper with decreasing Fourier frequency by $`\mathrm{\Delta }\mathrm{\Gamma }0.25`$ between each of the HF, MF and LF bands, no spectral slope variations are observed in the soft band power law component. c) An iron line, while present and clearly broad and asymmetric in the energy spectrum of the entire observation of this source, it is apparently absent in the energy spectra of the HF, MF and LF bands. d) A significant edge-like feature at $`E0.72`$ keV is present in the spectra of all frequency bands.
Past studies based on “excess variance” analysis of $`1`$ day long light curves have shown that the variability amplitude in AGN is higher in the soft than the hard band (Nandra et al. 1997a, Leighly 1999). However the variance, being the integral of the power spectrum over frequency, say from $`\nu _{\mathrm{min}}`$ to $`\nu _{\mathrm{max}}`$, does not provide information on the variability amplitudes at the different time scales. This information, on the other hand, is provided by the Fourier resolved spectra obtained in our analysis. We thus find that the hard band variations are of larger amplitude than the soft band ones on time scales shorter than $`3`$ ksec, while the opposite effect is observed on time scales longer than 10 ksec. The distinctly different timing properties of these two bands, in conjunction with their very different spectra argue that they represent emission by different physical components, likely of different physical dimension.
However, even within the hard band itself, the spectral properties do depend on the Fourier frequency, as the spectra become softer with decreasing Fourier frequency. This can be attributed to the fact that the hard band fast varying components have spectra significantly harder than those of slower ones, or alternatively, that the higher energy variations are consistently shorter than those of lower energies. This result is not consistent with the suggestion of Vaughan & Fabian (2004) that the spectral variability of MCG -06-30-15 over the entire $`0.210`$ keV band can be explained by the presence of a constant component and a power law component which varies only in normalization.
On the other hand, this conclusion is in agreement with the power spectrum analysis of the XMM-Newton data of MCG -06-30-15 by Vaughan et al. (2003) who find that the power spectrum becomes flatter with increasing energy above the “break” frequency of $`10^4`$ Hz. Similar behavior (power spectrum hardening at high frequencies with increasing energy) has been recently observed in other AGN (NGC 7469: Papadakis, Nandra & Kazanas, 2001; Mkn 766: Vaughan & Fabian, 2003; NGC 4051: McHardy et al. 2004, and 1H 0707-495, Ton S180 with the use of structure functions, Leighly 2004), as well as in the galactic black hole candidate Cyg X-1. This result by itself, ignoring the constraints imposed by the presence and variability properties of the soft band, could be interpreted as variability due to a Comptonizing corona of decreasing electron temperature with radius, along the general lines discussed in Kazanas, Hua & Titarchuk (1997).
The same interpretation, however, is not consistent with the variability of the soft component whose slope seems independent of the Fourier frequency. At this point, in the absence of lower energy data and based on the general shape of this component we are willing to speculate that it is due to emission by a (non-standard) accretion disk. As such we have in mind the ADIOS flows of Blandford & Begelman (1999). Assuming a viscosity parameter $`\alpha 0.3`$, a black hole mass $`M10^6`$ M and a disk size $`R10R_S`$, we obtain a characteristic variability scale $`t_{\mathrm{var}}(R/c)\alpha ^1(R/R_S)^{3/2}10^4`$ s, in agreement with observations. The disk photons could then serve as the seeds needed to produce the harder power law component by hot electrons in a corona interior to and/or overlying part of this disk. In this respect, one should bear in mind that the lags between the soft (0.2-0.7 keV) and the hard (2-10 keV) bands (Vaughan et al., 2003) are in general agreement with such an interpretation and the observed tight correlation between the light curves of the soft and hard spectral components.
The results of our analysis indicate that the Fe K$`\alpha `$ line shows no significant variations on time scales shorter than $`12`$ days and are in agreement with the results of the previous studies (e.g. Vaughan & Fabian 2004, and references therein). The only grounds for concern in this conclusion is the MF energy spectrum. First of all, the residuals around $`57`$ keV in the “data/best model” ratio plot (Fig. 3 and 4) are suggestive of the presence of a line feature (albeit of small amplitude). Furthermore, if the MF spectrum were consistent with a just power law of index $`\mathrm{\Gamma }=1.9`$, we would expect a decrease in the above ratio by $`40`$% in the $`57`$ keV band (the strength of the line above the $`\mathrm{\Gamma }=1.9`$ power law in the total spectrum; see top plot in the bottom panel of Fig. 2), but we observe a smaller drop ($`15`$%), in the “MF/total spectrum” ratio (Fig. 2), a value that implies a possible difference in the true MF slope from the assumed value of $`\mathrm{\Gamma }=1.9`$. We plan to investigate this issue further with the use of more data from recent Chandra, RXTE and ASCA observations of the source.
Our conclusions are similar to those of Revnivtsev et al. (1999) for Cyg X-1 in its low/hard state. They found the Equivalent Width (EW) of the iron line in the spectrum of the $`1`$ Hz and $`10`$ Hz Fourier components to be 80 and 50 eV, respectively. We find the $`3\sigma `$ upper limit in the EW of the line to be $`60`$ eV in the spectrum of the $`10^510^4`$ Hz components in MCG -06-30-15. The ratio of these frequency bands in the two sources is $`10^5`$, roughly comparable to the ratio of the black hole mass in the two systems, if we assume a $`10`$ M and $`10^6`$M mass for Cyg X-1 and MCG -06-30-15 (McHardy et al. 2005), respectively. As noted by Revnivtsev et al. (1999), the most straightforward interpretation of the absence of a line in the frequency resolved spectra of frequency higher than $`\nu _k`$ is that the X-ray reprocessing matter is located at a distance $`R_{\mathrm{rep}}>c[\nu _k/2\pi ]^1`$. For $`\nu _k10^5`$ Hz this is $`R_{\mathrm{rep}}c15000\mathrm{s}4.5\times 10^{14}`$ cm, or $`1500R_S`$ for a $`10^6`$ M black hole mass. If this is indeed the case, the line flux light curve should be much smoother than that of the continuum, and its variations should be of amplitude smaller than that of the continuum, a fact in agreement with observations (e.g. Vaughan & Fabian 2004). However, such large distances are hard to reconcile with the observed line width that requires line emission size $`R_{\mathrm{rep}}<10R_S`$. Recently, Miniutti et al. (2003) suggested that this lack of variability could be explained by a combination of light bending effects and vertical motions of the X-ray source above the accretion disk. Alternatively, the large width of the line maybe due to its downscattering in an expanding wind, as suggested recently by Titarchuk, Kazanas & Becker (2003). Inoue & Matsumoto (2003) also argue that the large width of the line is an artifact caused by the fact that the continuum suffers from absorption by various layers of warm absorbing material with variable column density and/or covering factor on time scales $`10^5`$ s. In this case, we should observe absorption features in the $`58`$ keV band of the Fourier spectra, but none are clearly evident in the residual plots (Fig. 3 and 4).
Of particular interest is the feature at $`0.7`$ keV, which is easily discernible in the energy spectra of all 3 frequency bands. It has been attributed to a combination of O VII edge and FeI absorption (Turner et al.2003). In this case one can set a limit on the distance of the absorber, assuming that the later remains unaffected by the changes in the continuum flux; this assumption is justified by the absence of a strong dependence of this feature’s depth on the Fourier frequency. This sets a limit to the value of the photoionization parameter of the absorbing medium to values $`\xi =L/nR^2<30`$ erg cm/s. On the other hand, the depth of the feature sets also a limit on the column density of the absorbing gas $`N_H=nR10^{22}`$ cm<sup>-2</sup>, values consistent with the detailed models of Turner et al. (2003); this results, then, in $`RL/\xi N_H10\mathrm{pc}(L_{43.5}/\xi _{1.5}N_{H_{22}})`$, where the subscripts denote the base 10 exponent of the value of the corresponding parameter in cgs units.
As we have already mentioned, the same feature along with the spectral “hump” in the residual plots of Fig. 4 around 0.5-0.55 keV has been interpreted as a relativistically broadened O K$`\alpha `$ line. The presence of this feature in the spectra of all Fourier components (and in particular those of highest Fourier frequency) is, taken at face value, consistent with such an interpretation. Interestingly, an absorption feature by matter with properties obtained in the previous paragraph would also conform with the same phenomenology. On the other hand, the broad Fe K$`\alpha `$ line which is modeled as emission by plasma in the same physical location as that responsible for an O K$`\alpha `$ line, has very different timing properties, as it is absent from any of the Fourier resolved spectra. This difference in the timing properties of these two components that have presumably their origin in the kinematics of the same plasma makes the presence of a soft X-ray line suspect. In addition, when one considers the effects of the highly variable hard X-ray component on the ionization structure of the line emitting plasma that would highly ionize the accretion disk plasma (Nayakshin, Kazanas & Kallman 2000; Nayakshin & Kazanas 2001), the interpretation of this feature as a relativistically broadened O K$`\alpha `$ line becomes doubtful.
In conclusion, the Fourier resolved spectroscopy technique discussed in the present work is a powerful instrument in analyzing the spectro-temporal properties of AGN and accreting compact objects in general. By attributing specific spectral properties/features to specific Fourier frequencies/time scales it allows one to further dissect the structure of these sources. It is hoped that the additional information provided by this technique will allow for a deeper understanding of this structure and the associated physics. Our work is but a first step in this direction. At this point it is not clear whether the properties implied by our analysis represent general AGN trends or idiosyncrasies of MCG -06-30-15. We hope that analysis of additional objects along the same lines will help decide this question.
DK would like to acknowledge a stimulating discussion with Tim Kallman. Part of this work was supported by the General Sectreteriat of Research and Technology of Greece.
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# Physics of Heavy Flavour at CDF
## 1 Introduction
The upgraded Collider Detector at Fermilab (CDF II) has collected around $`800\mathrm{pb}^1`$ between February 2002 and February 2005 during the Tevatron Run II at Fermilab. At $`p\overline{p}`$ colliders a large amount of $`b`$ and $`c`$ mesons and baryons are produces within a background of hadronic particles. However the presence of heavy flavour particle’s decays can be detected by the presence of a displaced secondary vertex, because of these particles have long decay length (O($`100\mu m`$)). The issue is to be able to extract this information at trigger level. For this purpose CDF uses the Silicon Vertex Trigger (SVT) that reconstruct online the tracks providing the informations needed for the trigger decision. In this way CDF is able to efficiently select events in which the heavy meson decays in either leptonic or fully hadronic modes. The collected data samples allow to perform a wide range of measurements, from the observation of rare decays to lifetime measurement, through BR and $`CP`$ asymmetry measurements. In the following we concentrate on some selected topics.
## 2 Branching ratios and CP asymmetries measurements
Fully hadronic $`b`$ meson decays are very useful to understand the $`b`$ sector of the CKM matrix. CDF is providing interesting measurement both on two body charmless and on pure penguins decays.
### 2.1 Two body charmless decays ($`Bh^+h^{}`$)
These decays are the ones in which a $`B^0`$, $`B_s^0`$ meson goes into charged Pions and Kaons. The theoretical prediction on their BR and CP asymmetries are strongly affected by uncertainties on hadronic contributes. These unknowns can be removed by combining the informations obtained in the different modes.
The mass resolution at CDF is not enough to directly observe the different signals, but their yields can be extracted via an unbinned likelihood fit that exploit both kinematic and energy loss information. The overall yield is shown in the left plot in Fig.1. The result of the fit is shown in the right plot of Fig.1. In $`180\mathrm{pb}^1`$ we observe $`509B^0K^+\pi ^{}`$, $`134B^0\pi ^+\pi ^{}`$ and $`232B_s^0K^+K^{}`$. We measured the ratio: $`\frac{f_s(B_s^0K^+K^{})}{f_d(B^0K^+\pi ^{})}=0.50\pm 0.08(stat)\pm 0.07(syst)`$. We set also the limits on the BRs of rare $`B_s^0`$ decays as: $`\frac{f_s(B_s^0K^+\pi ^{})}{f_d(B^0K^+\pi ^{})}<0.11`$ and $`\frac{(B_s^0\pi ^+\pi ^{})}{(B_s^0K^+K^{})}<0.10`$ both at $`90\%C.L.`$. We also measure the CP asymmetry in the $`B^0K^\pm \pi ^mp`$ decay and we obtain $`0.04\pm 0.08(stat)\pm 0.01(syst)`$.
### 2.2 Pure penguin decays
B meson decays involving $`bs\overline{s}s`$ transitions can provide evidences of deviation from the SM. In particular direct CP asymmetry of $`B^\pm \varphi K^\pm `$ mode is expected to be of the order of few percent within the SM. Left plot in Fig. 2 shows the mass distribution containing the signal of this decays obtained at CDF. The number of signal events has been extracted from the background using an unbinned likelihood fit on kinematic the particle’s energy losses. We measured $`(B^\pm \varphi K^\pm )=7.6\pm 1.3(stat)\pm 0.6(syst)\times 10^6`$ and $`a_{CP}(B^\pm \varphi K^\pm )=0.07\pm 0.17(stat)_{0.02}^{+0.03}(syst)`$.
The same $`b`$ transition occurs in $`B_s^0\varphi \varphi `$ decay. First evidence of this decay has been found at CDF and the right plot in Fig. 2 shows the observer signal. We measure $`(B^\pm \varphi \varphi )=1.4_{0.5}^{+0.6}(stat)\pm 0.2(syst)\pm 0.5(BR)\times 10^6`$ where the last error come from the BR of the $`B_s^0J/\psi \varphi `$ that has been used as normalization mode. Details of these analysis can be found in Ref..
## 3 Measurement of decay width difference of $`B_s^0`$ CP eigenstates
In $`B_s^0J/\psi \varphi `$ two vector particles are produced from the decay of a pseudo scalar one. It is indeed possible to distinguish the two CP eigenstates by the relative angular distribution of the decay’s products.
Actually three linear amplitudes are possible corresponding to angular momenta 0,1 and 2. The even/odd values of angular momentum occur for CP even/odd eigenstate. We simultaneously fit the lifetime and the linear amplitudes. The plot in Fig. 3 shows the projections of the fit result on the lifetime. The yellow line is the lifetime of the heavy mass eigenstate while the lifetime of the light one is reported in red. We found $`\frac{\mathrm{\Delta }\mathrm{\Gamma }_s}{\mathrm{\Gamma }_s}=65_{33}^{+25}\%`$.
## 4 Spectroscopy
In the field of spectroscopy the measurement of the $`B_c`$ mass is important to validate the theoretical models that predict this quantity . This meson has been observed at D0 and CDF in semileptonic modes. Because in this kind of decays the events are not fully reconstructed the achieved mass resolution was not enough to constrain the theoretical predictions. We look for evidence of the fully reconstructed $`B_cJ/\psi \pi `$ decay in the mass range between $`5.6`$ and $`7.2\mathrm{GeV}/\mathrm{c}^2`$ corresponding to $`2\sigma `$ window around the previously measured value.
We optimize the selections following a blind procedure. The plot in Fig. 4 show the mass distribution in the signal region after applying the optimized cut. From the fit we measured a signal of $`18.9\pm 5.7`$ events over a background of $`10.0\pm 1.4`$. The value of the mass we obtain is $`6287.0\pm 4.8(stat)\pm 1.1(syst)\mathrm{MeV}/\mathrm{c}^2`$ where the systematic error is mainly given by the parametrization of the background.
## 5 Conclusions
We have reported some examples of the wide range of heavy flavour particles that can be detected at CDF and how their characteristics can be investigated. This analysis are still statistically limited. However the systematics are well under control and the results will be easily improved by increasing the data samples.
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# Variable structure control for parabolic evolution equations
## 1 Introduction
Variable structure control methods and in particular sliding mode controls, are by now recognised as classical tools for the regulation of systems governed by ordinary differential equations in a finite dimensional setting. For an overview of the finite-dimensional theory see .
While being easy to design, they possess attractive properties of robustness and insensitivity with respect to disturbances and unmodelled dynamics. These characteristics are all the more important when dealing with infinite-dimensional systems. In many control applications such as heat transfer processes, chemical processes, flexible manipulators the state evolution is governed by a partial differential equation. The complexity of these plants results in models having significant degrees of uncertainty. Thus motivated, recent research has been devoted to the extension of sliding mode control and therefore the use of discontinuous feedback laws, to the infinite-dimensional setting. While earlier works were confined to some special classes of systems, at present both theory and application of sliding mode control have been extended to a rather general setting . In particular in the key concept of equivalent control is introduced in a general Hilbert space framework for evolution equations governed by unbounded linear operators that generate $`C_0`$-semigroups. Also it is shown that, under some stability assumptions, the ideal sliding can be uniformly approximated by “real” motions evolving in a boundary layer of the sliding manifold, thus ensuring the validity of the method for application purposes. The relationship between the equivalent control method and generalised solutions of infinite-dimensional systems with discontinuous right-hand side is presented in .
All the results in the above cited literature only take into consideration distributed control systems, i.e. they deal with bounded input operators. In this paper we make a first attempt to consider the extension of sliding modes to a class of boundary control problems in a general setting. To the author’s knowledge there exist only a few results in this direction in the linear case , where by application of integral transformations the problem is reduced to the control of a finite-dimensional differential-difference equation. Our approach goes instead in the direction of . In Section 2 we define the general abstract variational framework in which we set up our control problem. In particular, the main assumptions we make on the operator governing the evolution, are weak continuity and coerciveness, so that both linear and non-linear operators are comprised in this setting. In Section 3 we present our main result: a Faedo-Galerkin method is used to construct a sequence of finite-dimensional approximations of the given problem. On each of these the standard variable structure control theory of can be applied. We then assume that for each approximation a control law is chosen to constrain the evolution in a boundary layer of a given sliding manifold and study the limit as the dimensions diverge. We show that, under some growth assumption on the norm of these controls, a limit motion exists, which satisfies the sliding condition. Then, in Section 4 we apply the obtained results to the Neumann boundary control of a heat equation.
## 2 Abstract setting and problem statement
In this paper we are going to consider a class of parabolic partial differential equations with controllers acting on the boundary. In particular we will study the case of Neumann boundary conditions and finite dimensional control space. Also, we suppose that a manifold $`S`$ is given, on which we want to restrict the motion of our system. We then analyse the problem of the existence of an admissible control law for which this ideal sliding motion is possible.
###### Example 2.1
Before going into the details of the precise abstract setting of the problem, we show an example of application to give an idea of the family of systems we intend to study.
Let $`\mathrm{\Omega }`$ be a bounded, open subset of $`IR^n`$ with smooth boundary $`\mathrm{\Gamma }`$, $`T>0`$ and $`\mathrm{\Delta }`$ be the laplacian differential operator on $`IR^n`$. Consider the following evolution equation
$$\begin{array}{cc}\frac{Q}{t}(t,x)=\mathrm{\Delta }Q(t,x)+q(x)Q(t,x)\hfill & t(0,T),x\mathrm{\Omega }\hfill \\ \frac{Q}{\nu }(t,\sigma )=u(t)g(\sigma )\hfill & t(0,T),\sigma \mathrm{\Gamma }\hfill \\ Q(0,x)=Q_0(x)\hfill & x\mathrm{\Omega }.\hfill \end{array}$$
(1)
Here $`Q:[0,T]\times \mathrm{\Omega }IR`$ represents the evolution of the “state vector”, $`u:[0,T]IR`$ is a scalar control law, $`g:\mathrm{\Gamma }IR`$ and $`q:IR^nIR`$ is bounded. This equation represents a model of heat conduction with both diffusion and heat generation (if $`q`$ is nonnegative). Now for $`\gamma :\mathrm{\Omega }IR`$ we can define (informally) a sliding surface $`S`$ as the set of functions $`f:\mathrm{\Omega }IR`$ such that
$$_\mathrm{\Omega }f(x)\gamma (x)𝑑x=0$$
In this case a sliding motion $`Q(t,x)`$ on $`S`$ would satisfy
$$_\mathrm{\Omega }Q(t,x)\gamma (x)𝑑x=0,t>0$$
### 2.1 Variational formulation
The setting of the abstract problem follows : let $`V`$ be a separable, reflexive Banach space, $`H`$ be a Hilbert space, $`VH`$ with continuous injection. The space $`H`$ is identified with its dual, while we denote by $`V^{}`$ the dual space of $`V`$, so that we have
$$VHV^{}.$$
For $`u_1`$, $`u_2H`$ the scalar product in $`H`$ will be denoted by $`(u_1,u_2)`$ and the derived norm by $`|u_i|`$. We will denote by $``$ the norm in $`V`$ and by $`_{}`$ that in $`V^{}`$. The dual pairing between the two spaces will be written as $`,`$. Also, we will assume that on $`V`$ it is defined a semi-norm $`[]`$ such that
$$[v]+\lambda |v|\beta v,vV,\text{for some }\lambda ,\beta >0.$$
(2)
It is assumed that all the above (infinite-dimensional) spaces are real vector spaces; results can be extended to the complex case with the necessary modifications. For any $`T>0`$ we can define the following spaces of vector-valued functions:
$$L^2(0,T;V)=\{f:[0,T]V:_0^Tf(t)^2𝑑t<+\mathrm{}\}$$
$$L^{\mathrm{}}(0,T;H)=\{f:[0,T]H:\underset{t[0,T]}{sup}|f(t)|<+\mathrm{}\}.$$
The space $`L^2(0,T;V^{})`$ can be defined analogously. Also, it is possible to define on these spaces a concept of derivative, in a distributional sense (see i. e. Chapter III). The following result will be useful in the sequel.
###### Theorem 2.1
Let
$$W(0,T)=\{fL^2(0,T;V):\frac{df}{dt}L^2(0,T;V^{})\}.$$
All functions in $`W(0,T)`$ are, after eventual modification on a null measure set, continuous from $`[0,T]`$ in $`H`$, i.e. $`W(0,T)C^0(0,T;H)`$.
For $`t(0,T)`$ let $`A(t):VV^{}`$ be an operator satisfying the following assumptions:
* for all $`v,wV`$ the map
$$tA(t)v,w\text{ is measurable;}$$
(3)
* for all $`t`$ and any $`u,v,\omega V`$ the map
$$\alpha A(t)(u+\alpha v),w\text{ is continuous;}$$
(4)
* there exist constants $`c_1>0`$, $`c_20`$ such that
$$A(t)v_{}c_1v+c_2,vV;$$
(5)
* there exist constants $`\alpha >0`$ and $`\nu IR`$ such that
$$A(t)v,v\alpha [v]^2+\nu |v|^2vV.$$
(6)
* $`A()`$ is $`2`$-weakly continuous, i.e.
$$\begin{array}{c}v_kv\text{ weakly in }W(0,T)\hfill \\ \hfill A()v_k()A()v()\text{ weakly in }L^2(0,T;V^{}).\end{array}$$
(7)
Let $`UIR^m`$ be closed and convex and let $`f:[0,T]\times UV^{}`$ satisfy the following condition: there exists a constant $`C>0`$ such that for any $`u:[0,T]U`$, $`uL^2(0,T)`$
$$_0^Tf(t,u(t))_{}^2Cu_2^2,$$
(8)
where $`u_2`$ is the usual $`L_2`$-norm (it will always be understood that control laws $`u`$ take values in $`U`$, so that we will write $`uL^2(0,T)`$ instead of $`L^2(0,T;U)`$).
We are now ready to write the abstract evolution equation we are going to study. The evolution of the system will be given by a vector-valued function $`yW(0,T)`$ satisfying the following abstract Cauchy problem
$$\{\begin{array}{c}\frac{dy}{dt}+A(t)y(t)=f(t,u(t))\mathrm{q}.\mathrm{o}.t\hfill \\ y(0)=y_0,\hfill \end{array}$$
(9)
with $`uL^2(0,T)`$ and for some $`y_0H`$ (by Theorem 2.1 this makes sense). The differential equation above as to be understood as an equality in the dual space $`V^{}`$, i.e. setting
$$a(t;v,w)=A(t)v,w,t>0,v,w,V$$
(10)
and in view of Theorem 2.1, the differential problem (9) is equivalent to the following variational formulation
$$\{\begin{array}{c}\frac{d}{dt}(y(t),v)+a(t;y(t),v)=f(t,u(t)),vvV,\hfill \\ y(0)=y_0\hfill \end{array}$$
(11)
Existence and uniqueness results of the solution of such equations, under our assumptions, can be found in under monotonicity assumptions and in for the linear case.
###### Example 2.2
Let us see how Example 2.1 fits into this framework. Let $`H=L^2(\mathrm{\Omega })`$ and
$$V=H^1(\mathrm{\Omega })=\{fH:\frac{f}{x_i}Hi=1,\mathrm{},n\}.$$
On $`V`$ we set $`[v]^2=|v|^2`$ and $`v^2=[v]^2+|v|^2`$. Let $`vV`$ be arbitrary; by scalar multiplication and using Green’s formula one finds that the solution $`Q`$ of (1) has to satisfy
$`{\displaystyle \frac{d}{dt}}(Q(t,),v)`$ $`=`$ $`{\displaystyle _\mathrm{\Omega }}\mathrm{\Delta }Q(t,x)v(x)𝑑x+(qQ(t,),v)`$
$`=`$ $`{\displaystyle _\mathrm{\Omega }}_xQ(t,v)v(x)𝑑x`$
$`+{\displaystyle _\mathrm{\Gamma }}u(t)g(\sigma )v(\sigma )𝑑\sigma +(qQ(t,),v).`$
Therefore setting $`y_0=Q_0`$ and $`y(t)=Q(t,)`$ we get the (autonomous) variational formulation of our abstract setting in the form (11) with
$$a(v,w)=(v,w)(qv,w)$$
(12)
and
$$f(t,u),v=_\mathrm{\Gamma }ug(\sigma )v(\sigma )𝑑\sigma .$$
(13)
Now (4) and (5) are easily verified and (6) follows from
$$a(v,v)=[v]^2(qv,v)[v]^2(\underset{\mathrm{\Omega }}{sup}q)|v|^2.$$
Also, the operator $`A`$ defined as $`Av,w:=a(v,w)`$ is linear and bounded, therefore it is $`2`$-weakly continuous and we have (7).
Moreover, on $`V`$ the trace operator $`\tau `$ of restriction of a function to the boundary of $`\mathrm{\Omega }`$ is well defined . The range of $`\tau `$ is the Banach space $`Z=H^{1/2}(\mathrm{\Gamma })`$ and $`\gamma `$ is continuous from $`V`$ onto $`H^{1/2}(\mathrm{\Gamma })`$. Therefore $`f`$ is well defined for any $`g`$ in the dual of $`H^{1/2}(\mathrm{\Gamma })`$, hence for example for all $`gL^2(\mathrm{\Gamma })`$ and obviously satisfies (8) with $`C=g_{L^2(\mathrm{\Gamma })}\tau _{(V,Z)}`$.
## 3 Main results
In this section we introduce the concept of sliding surface for the control problem (11) and show how sliding motions can be defined in this context.
Assume we are working in the framework set up in Section 2. Thanks to separability, there exists a countable basis for $`V`$, so that it is possible to define a family $`\{V_k\}_{kIN}`$ of finite dimensional subspaces of $`V`$
$$V_k=\mathrm{span}\{v_{1,k},\mathrm{},v_{N_k,k}\}$$
such that
$$V_kV_{k+1},\underset{kIN}{}V_k=V.$$
Then it is possible to define approximate solutions of (11) by projecting on the subspaces $`V_k`$, using the standard Faedo-Galerkin method. We thus define the following family of variational problems: find $`y_k:[0,T]V_k`$ such that
$$\{\begin{array}{c}\frac{d}{dt}(y_k(t),v)+a(t;y_k(t),v)=f(t,u_k(t)),vvV_k,\hfill \\ y_k(0)=y_{0,k}\hfill \end{array}$$
(14)
with $`y_{0,k}V_k`$ for all $`k`$ and a sequence $`\{u_k\}`$ in $`L^2(0,T)`$. Note that, since $`V_k`$ has dimension $`N_k`$, the above problem can be written as an ordinary differential equation. In fact, since $`y_k(t)V_k`$, there exists a vector $`\xi _k(t)IR^{N_k}`$ such that
$$y_k(t)=\underset{i=1}{\overset{N_K}{}}(\xi _k(t))_iv_{i,k}.$$
The differential equation in (14) is satisfied for all $`vV_k`$ iff it is valid for every element of the basis of $`V_k`$. Therefore, if
$$y_{0,k}=\underset{i=1}{\overset{N_K}{}}(\xi _{0,k})_iv_{i,k},$$
$$f_k(t)=(f(t,u_k(t)),v_{i,k})_{i=1,\mathrm{},N_k},$$
and
$`A_k(t)=(a_{ij}^{(k)}(t))_{i,j=1,\mathrm{},N_k},`$ $`a_{ij}^k(t)=a(t;v_{i,k},v_{j,k}),`$
$`M_k=(m_{ij}^{(k)}(t))_{i,j=1,\mathrm{},N_k},`$ $`m_{ij}^k(t)=(v_{i,k},v_{j,k}),`$
the differential problem (14) is equivalent to the following ordinary Cauchy problem
$$\{\begin{array}{c}M_k\dot{\xi }_k(t)+A_k\xi _k(t)=f_k(t)\hfill \\ \xi _k(0)=\xi _{0,k}.\hfill \end{array}$$
(15)
We now prove a convergence result for the approximations $`y_k`$, under some conditions on the controls sequence $`\{u_k\}`$.
###### Theorem 3.1
Let the assumptions in Section 2 be satisfied and $`\{u_k\}`$ be a sequence in $`L^2(0,T)`$. Let $`y_k`$ be the solution of (14) and suppose that $`y_{0,k}y_0`$ in $`H`$ for $`k+\mathrm{}`$. Suppose moreover that the following condition on the growth of the control norms is satisfied
$$u_k_{L^2(0,t)}^2M_0^t|y_k(s)|^2𝑑s+N,tT$$
(16)
for some non-negative constants $`M`$ and $`N`$ and that $`f`$ is the following weak continuity assumption
$$\begin{array}{c}u_ku^{}\text{ weakly in }L^2(0,T)\text{ then }\hfill \\ \hfill f(,u_k())f(,u())\text{ weakly in }L^2(0,T;V^{})\end{array}$$
(17)
Then there exist a control law $`u^{}L^2(0,T)`$ and a function $`y^{}W(0,T)`$ verifying (11), such that, for some subsequence,
$`y_ky^{}\text{ weakly in }W(0,T)`$
$`y_ky^{}\text{ weakly* in }L^{\mathrm{}}(0,T;H)`$
$`u_ku^{}\text{ weakly in }L^2(0,T).`$
Proof. Writing (14) for $`v=y_k(t)`$ we get
$$(\dot{y}_k(t),y_k(t))+a(t;y_k(t),y_k(t))=f(t,u_k(t)),y_k(t).$$
As the first term on the left is in fact the time derivative of $`|y_k(t)|^2/2`$, integrating the above identity we have
$$\begin{array}{c}\frac{1}{2}|y_k(t)|^2+_0^ta(t;y_k(s),y_k(s))𝑑s=\hfill \\ \hfill \frac{1}{2}|y_k(0)|^2+_0^tf(t,u_k(s)),y_k(s)𝑑s.\end{array}$$
By (6), (8) and (2) we obtain the following inequality
$$\begin{array}{c}\frac{1}{2}|y_k(t)|^2+\alpha _0^t[y_k(s)]^2𝑑s\hfill \\ \hfill \frac{1}{2}|y_k(0)|^2\nu _0^t|y_k(s)|^2𝑑s\\ \hfill +cu_k_2\left(_0^t[y_k(s)]^2𝑑s+_0^t[y_k(s)]^2𝑑s\right)^{1/2}\end{array}$$
for some constant $`c>0`$. Consider now for $`x0`$ the function $`h(x)=(\alpha x)/2c\sqrt{x}`$. It is easy to show that it has minimum for $`x=(c/\alpha )^2`$, therefore $`c\sqrt{x}(\alpha x+c^2/\alpha )/2`$, thus
$$\begin{array}{c}\frac{1}{2}|y_k(t)|^2+\frac{\alpha }{2}_0^t[y_k(s)]^2𝑑s\hfill \\ \hfill \frac{1}{2}|y_k(0)|^2+\left(\frac{\alpha }{2}+|v|\right)_0^t|y_k(s)|^2𝑑s+\frac{c^2}{2\alpha }u_k_2^2.\end{array}$$
Now, since by hypothesis $`|y_{0,k}y_0|`$ tends to zero, the term $`|y_k(0)|^2`$ is bounded. Moreover by (16)
$$|y_k(t)|^2+\alpha _0^t[y_k(s)]^2𝑑sc_1+c_2_0^t|y_k(s)|^2𝑑s$$
(18)
for some constants $`c_1,c_2>0`$. Since $`\alpha >0`$ we get
$$|y_k(t)|^2c_1+c_2_0^t|y_k(s)|^2𝑑s$$
Therefore, by Gronwall’s lemma we obtain for some constant $`K>0`$
$$y_k_{L^{\mathrm{}}(0,T;H)}=\underset{t[0,T]}{sup}|y_k(t)|K$$
(19)
therefore from (18) we also have
$$_0^T[y_k(s)]^2𝑑s\mathrm{const}$$
and lastly, using (2) and (5)
$`y_k_{L^2(0,T;V)}=\left({\displaystyle _0^T}y_k(s)^2𝑑s\right)\mathrm{const},`$
$`{\displaystyle _0^T}A(t)y_k(t)_{}𝑑t\mathrm{const}.`$
Since spheres are weakly compact in both $`L^2(0,T;V)`$ and $`L^2(0,T;V^{})`$, weakly\* compact in $`L^{\mathrm{}}(0,T;H)`$, we can extract a subsequence of $`\{y_k\}`$ (which for simplicity we still denote by $`\{y_k\}`$) converging to some $`y^{}L^2(0,T;V)L^{\mathrm{}}(0,T;H)`$ for both the weak topology of $`L^2(0,T;V)`$ and the weak\* topology of $`L^{\mathrm{}}(0,T;H)`$ and such that $`Ay_k`$ weakly converges to some $`\eta `$ in $`L^2(0,T;V^{})`$. By (16) we also have that $`u_k_2`$ is bounded, thus eventually passing to a further subsequence, there exists $`u^{}L^2(0,T)`$ such that $`u_k`$ converges to $`u^{}`$ weakly in $`L^2(0,T)`$. Also, by (17) we can proceed as in the proof of Theorem 1.1, p. 159 of to conclude that
$$\{\begin{array}{c}\frac{d}{dt}y^{}(t)+\eta (t)=f(t,u^{}(t))\hfill \\ y(0)=y_0.\hfill \end{array}$$
Also, by a standard argument (see i.e. , Theorem 3) one can prove that $`\dot{y}_k\dot{y}^{}`$ weakly in $`L^2(0,T;V^{})`$, i.e. $`y_ky^{}`$ weakly in $`W(0,T)`$. Thus, by (7) $`\eta (t)=A(t)y^{}(t)`$ and the proof is complete. $`\mathrm{}`$
Having achieved the above convergence result, we introduce as in , a set $`D`$ which can be either $`V`$ or a sufficiently large open subset of $`H`$ and a mapping $`s:DIR^m`$ continuously Fréchet differentiable on $`D`$. The sliding surface $`S`$ we consider is defined as $`S=\{yD:s(y)=0\}`$. Proceeding as in , by slightly modifying proofs, it is possible to prove the following
###### Corollary 3.1
Let the assumptions of Theorem 3.1 hold. Let $`z_k(t)=s(y_k(t))`$ and assume that one of the following is satisfied:
* $`D=V`$, $`s`$ is affine and $`z_k0`$ uniformly in $`t`$;
* $`_H(0,K)D`$, $`V`$ is compactly embedded in $`H`$ (here $`_H`$ denotes a ball in $`H`$, while $`K`$ is defined in (19) above) and $`z_k(t)0`$ for almost every $`t[0,T]`$.
Then the limit motion $`y^{}`$ of Theorem 3.1 belongs to the sliding manifold $`S`$.
###### Remark 3.1
Note that by (15) every $`y_k`$ solves a finite-dimensional problem, thus for the approximate solutions all results of the classical theory of variable structure systems and sliding mode control of are valid. Therefore existence results for system motions satisfying the requirements in Corollary 3.1 and design methods to achieve them are available. See also the discussion of existence under relaxed hypotheses developed in .
## 4 An application
In this Section we show an application of the obtained results on the control problem introduced in Example 2.1. We have already proved (see Example 2.2) that this partial differential equation with Neumann control fits in the abstract setting of Section 2. It is also easy to prove that for $`f`$ as in (13) the condition (17) is satisfied. In fact, if $`u_ku`$ weakly in $`L^2(0,T)`$, for any $`\phi L^2(0,T;V)`$ we have
$$\begin{array}{c}_0^T[f(t,u_k(t))f(t,u(t))],\phi (t)𝑑t=\hfill \\ \hfill _0^T[u_k(t)u(t)]_\mathrm{\Gamma }g(\sigma )\phi (t)(\sigma )𝑑\sigma 𝑑t\end{array}$$
which converges to zero since by Hölder’s inequality and continuity of the trace operator on $`V`$
$$\begin{array}{c}_0^T\left(_\mathrm{\Gamma }|g(\sigma )||\phi (t)(\sigma )|𝑑\sigma \right)^2𝑑t\hfill \\ \hfill g_{L^2(\mathrm{\Gamma })}^2_0^T\phi (t)^2𝑑t<+\mathrm{}.\end{array}$$
We then set $`s:HIR`$, $`s(x)=(x,\gamma )`$ and $`S=\mathrm{ker}S`$. For convenience we suppose that the chosen bases of the subspaces $`V_k`$ are orthonormal, so that the matrix $`M_k`$ in (15) is the identity (this is not restrictive since in the general case $`M_k`$ is symmetric, positive definite and a linear change of coordinates is sufficient to reconduct this problem to the orthonormal one). Then, setting $`g_k=((g,\tau v_{i,k})_{L^2(\mathrm{\Gamma })})_{i=1,\mathrm{},N_k}`$, (15) can be rewritten as
$$\{\begin{array}{c}\dot{\xi }_k(t)+A_k\xi _k(t)=u_k(t)g_k\hfill \\ \xi _k(0)=\xi _{0,k}.\hfill \end{array}$$
Then $`z_k(t)=s(y_k(t))=(y_k(t),\gamma )=\gamma _k^T\xi _k(t)`$, with $`\gamma _k=((v_{i,k},\gamma ))_{i=1,\mathrm{},N_k}`$. Let $`V(t)=z_k^2(t)/2`$; then
$$\dot{V}(t)=z_k(t)\dot{z}_k(t)=z_k(t)[\gamma _k^T(A_k\xi _k(t)+u_k(t)g_k)].$$
By standard finite dimensional theory a sliding mode exists on $`S_k=\{xIR^{N_k}:\gamma _k^Tx=0\}`$ if $`\gamma _k^Tg_k0`$. Also, in this case, setting
$$u_k(t)=U(t)\frac{\mathrm{sign}(z_k(t))}{\gamma _k^Tg_k}$$
with $`U(t)>|\gamma _k^TA_k\xi _k(t)|`$ the sliding surface is globally attractive and reached in finite time. Moreover, if $`\delta _k>0`$ and $`|s(y_k(0))|<\delta _k`$ the control
$$u_k(t)=\frac{U(t)}{\gamma _k^Tg_k}\frac{z_k(t)}{|z_k(t)|+\delta _k}$$
constrains the motion of the system in a $`\delta _k`$-boundary layer of $`S_k`$. Let us now consider the term $`\gamma _k^TA_k\xi _k(t)`$; since we assumed that the basis of $`V_k`$ is orthonormal, we have
$$\gamma _k^TA_k\xi _k(t)=a(y_k(t),P_k\gamma ),$$
where $`P_k:VV_k`$ is the projection on $`V_k`$. Likewise we have
$$\gamma _k^Tg_k=_\mathrm{\Gamma }g(\sigma )(P_k\gamma )(\sigma )𝑑s.$$
Thus, if for example $`\gamma V`$ and
$$_\mathrm{\Gamma }g(\sigma )\gamma (\sigma )𝑑\sigma 0,$$
since $`P_k\gamma \gamma `$ in $`V`$, there exists $`K`$ such that $`\gamma _k^Tg_k0`$ for all $`kK`$. In order to apply Theorem 3.1 we also have to show that (16) holds. Recalling that
$$a(y_k(t),P_k\gamma )=(y_k(t),P_k\gamma )(qy_k(t),P_k\gamma )$$
we just have to show that, at least for suitable $`\gamma `$-s, the first term can be estimated using $`|y_k(t)|`$. Proceeding formally, by Green’s formula we have
$$\begin{array}{c}(y_k(t),P_k\gamma )=(y_k(t),\mathrm{\Delta }(P_k\gamma ))\hfill \\ \hfill +_\mathrm{\Gamma }y_k(t)(\sigma )\frac{}{\nu }(P_k\gamma )(\sigma )𝑑\sigma .\end{array}$$
Thus (16) can be satisfied if sufficiently regular decompositions $`\{V_k\}`$ of $`H^1(\mathrm{\Omega })`$ are chosen and if the function $`\gamma `$ satisfies $`\frac{}{\nu }P_k\gamma =0`$, at least on some subsequence. For example this is true if $`\gamma V_N`$ for some $`N`$ and $`\frac{}{\nu }\gamma =0`$.
###### Remark 4.1
In this paper we have chosen a variational setting for our problem, by which we can encompass also some non-linear partial differential equations. For the linear case, another common abstract setting involves semigroup theory. In the above example our operator $`A:VV^{}`$ could be, in some sense, substituted by $`𝒜:𝒟(𝒜)HH`$, with
$$𝒟(𝒜)=\{yH^2(\mathrm{\Omega }):\frac{y}{\nu }=0\},𝒜y=\mathrm{\Delta }y+qy.$$
Note that the last condition on $`\gamma `$ above is related to “$`\gamma 𝒟(𝒜^{})`$”, which is frequently encountered in the literature on output control of infinite-dimensional systems.
###### Remark 4.2
In many applications the function $`z(t)`$ of the example represents the system’s output. The modulus of the control law we have chosen depends on the whole state norm, which could be unavailable for measurement. In observers are designed to overcome this difficulty in the case of distributed control. It would be interesting to study their application to this case also.
## 5 Conclusions and future work
In this paper we have analysed the convergence behaviour of finite dimensional Faedo-Galerkin approximations of a class of variational problems, when sliding motions are taken into consideration. We have thus shown that, under some growth hypothesis on the norms of the controls, a sliding motion exists.
This is a first attempt to extend variable structure control to boundary control problems for infinite-dimensional systems and much work has still to be done in this area. Apart from the need to extend these results to different boundary control problems, it would be interesting to study how these results are related to a notion of equivalent control, which has already by introduced in the infinite-dimensional setting and to approximability of ideal sliding motions by real ones.
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# Two Nucleon-States in a Chiral Quark-Diquark Model
## 1 Introduction
Hadron spectrum should in principle be understood by Quantum Chromodynamics (QCD). However, due to the difficulties of non-perturbative nature of QCD, in practice, QCD oriented models are often used for the description of hadrons. The non-relativistic quark model is one of successful ones. Employing a confining potential which is usually taken as a harmonic oscillator, various hadrons are described as single particle states of two (for mesons) and three (for baryons) quark states. Being not always said explicitly, the SU(4)<sub>SF</sub> spin-(two) flavor symmetry is usually assumed to work well, and possible interactions such as spin-color and spin-flavor interactions are treated perturbatively.
Such residual interactions may play an essential role for hadron properties by determining symmetries of hadrons. One of efficient methods which takes into account such interactions is to consider diquark correlations. In fact, it is known that the spin-color and spin-flavor interactions lead to a strongly attractive correlation for a color, flavor and spin antisymmetric quark-quark pair (scalar diquark) . Contrary, the color symmetric states receive repulsive correlations. If such correlations are strong enough, the resulting hadron properties will be quite different from what are expected from the SU(4)<sub>SF</sub> quark models.
In this paper, we would like to consider the case in which diquarks receive sufficiently strong correlation such that they are treated as independent degrees of freedom. This is the SU(2)$`{}_{L}{}^{}\times `$SU(2)<sub>R</sub> chiral quark-diquark model which we studied previously.
It was shown that there are two low lying diquarks with color $`\overline{3}`$, one is the scalar and isoscalar ($``$ scalar) diquark, and the other is the axial-vector and isovector ($``$ axial-vector) diquark. In a quark model picture, they are both in the ground state but with different spin and isospin configuration; the former is spin-singlet and flavor-singlet and the latter is spin-triplet and flavor-triplet ones. Adding another quark to these diquarks we can construct two states having the nucleon quantum numbers $`J^P=1/2^+`$, i.e. $`[[1/2,1/2]^0,1/2]^{1/2}`$ and $`[[1/2,1/2]^1,1/2]^{1/2}`$. In the SU(4)<sub>SF</sub> quark models, such two states do not exist independently, but only their linear combination is realized (the sum of the $`\rho `$ and $`\lambda `$ symmetric states) . If the diquark correlation exists, these two states may be treated independently.
Having the two states independently in the model set up, the two nucleon states emerge naturally. It is known that the scalar diquark is lighter than the axial-vector diquark, whose mass difference may be related to the mass difference between the nucleon and delta, which is of order 300 MeV. Hence we have the two nucleon states with the mass difference of this order to start with. Introducing a mixing interaction between them, the two states repell each other. One would then expect additional, but not too large, mass splitting due to the mixing interaction. Altogether, we expect a mass difference about 4-500 MeV. A nucleon excited state with an excitation energy of this amount is naturally identified with the Roper resonance $`N(1440)`$. The Roper state appears as a spin partner of the nucleon in the sense that it has a diquark component of different spin. Hence the mass splitting is generated by residual (hyper-fine) interactions and is likely not too large. This picture is very much different from the conventional picture of the Roper resonance of the excitation of $`2\mathrm{}\omega `$ of the harmonic oscillator potential for confinement or breathing mode of the bag model or Skrymion. In this paper, we demonstrate such a realization of the Roper resonance as a partner of the ground state of the nucleon in the quark-diquark model with a path-integral hadronization method.
This paper is organized as follows. In §2 we introduce the quark-diquark model with scalar and axial-vector diquarks. Chiral symmetry of the model is briefly discussed in the nonlinear representation. In §3 the model is hadronized to obtain an effective meson-baryon Lagrangian in the path-integral method. The obtained Hamiltonian is diagonalized in the two dimensional space of the two nucleons constructed by the scalar and axial-vector channels. In §4 numerical results are presented and the higher nucleon state is identified with the Roper resonance. The final section is devoted to a summary.
## 2 Framework
We start from the SU(2)$`{}_{L}{}^{}\times `$ SU(2)<sub>R</sub> chiral quark-diquark model of two diquarks ,
$`=\overline{\chi }_c(i\text{/}m_q)\chi _c+D_c^{}(^2+M_S^2)D_c+\stackrel{}{D}_{}^{\mu }{}_{c}{}^{}\left[(^2+M_A^2)g_{\mu \nu }_\mu _\nu \right]\stackrel{}{D}_c^\nu +_{int},`$ ( 2.1)
where $`\chi _c`$, $`D_c`$ and $`\stackrel{}{D}_{\mu c}`$ are the constituent quark, scalar diquark and axial-vector diquark fields, and $`m_q`$, $`M_S`$ and $`M_A`$ are the masses of them. The indices $`c`$ represent the color. Note that the diquarks microscopically correspond to the bi-linears of two quarks; $`D_cϵ_{abc}\stackrel{~}{\chi }^b\chi ^c,\stackrel{}{D}_{\mu c}ϵ_{abc}\stackrel{~}{\chi }^b\gamma _\mu \gamma _5\stackrel{}{\tau }\chi ^c`$. Here $`\stackrel{~}{\chi }=\chi ^TC\gamma _5i\tau _2`$ and $`ϵ_{abc}`$ is the antisymmetric tensor, then both the diquarks belong to color anti-triplet. The term $`_{int}`$ is the quark-diquark interaction, which is written as
$$_{int}=G_S\overline{\chi }_cD_c^{}D_c^{}\chi _c^{}+v(\overline{\chi }_cD_c^{}\gamma ^\mu \gamma ^5\stackrel{}{\tau }\stackrel{}{D}_{\mu c^{}}\chi _c^{}+\overline{\chi }_c\gamma ^\mu \gamma ^5\stackrel{}{\tau }\stackrel{}{D}_{\mu c}^{}D_c^{}\chi _c^{})+G_A\overline{\chi }_c\gamma ^\mu \gamma ^5\stackrel{}{\tau }\stackrel{}{D}_{\mu c}^{}\stackrel{}{\tau }\stackrel{}{D}_{\nu c^{}}\gamma ^\nu \gamma ^5\chi _c^{},$$
( 2.2)
where $`G_S`$ and $`G_A`$ are the coupling constants for the quark and scalar diquark, and for the quark and axial-vector diquark. The coupling constant $`v`$ causes the mixing between the scalar and axial-vector channels in the nucleon wave-functions.
Since chiral symmetry is important for hadron physics, we briefly discuss the properties of our model Lagrangian (2.1) under chiral transformation. In our formulation we employ the non-linear representation of chiral symmetry, therefore the constituent quarks are transformed as
$$\chi _ch(x)\chi _c.$$
( 2.3)
Here $`h(x)`$ is a local transformation depending on the global SU(2)$`{}_{L}{}^{}\times `$ SU(2)<sub>R</sub> transformation and the pion field at $`x`$. Then, the baryon field written as a product of a scalar diquark and a quark, and of an axial-vector diquark and a quark are transformed as (in detail see Appendix. A)
$`D_c\chi _c`$ $``$ $`h(x)D_c\chi _c,`$ (2.4a)
$`\stackrel{}{D}_{\mu c}\stackrel{}{\tau }\gamma ^\mu \gamma ^5\chi _c`$ $``$ $`h(x)\stackrel{}{D}_{\mu c}\stackrel{}{\tau }\gamma ^\mu \gamma ^5\chi _c.`$ (2.4b)
Note that the both baryon operators are transformed in the same way as the quark is under the chiral SU(2)$`{}_{L}{}^{}\times `$ SU(2)<sub>R</sub> transformation. In the non-linear representation, the kinetic term of the quark in Eq. (2.1) contains the mesonic currents . However, we do not show them, because they are not important in the present discussions.
## 3 Hadronization and the self-energies of nucleons
To perform the hadronization procedure, we introduce the auxiliary fields for baryons
$`=\overline{\chi }(i\text{/}m_q)\chi +D^{}(^2+M_S^2)D+\stackrel{}{D}^\mu \left[(^2+M_A^2)g_{\mu \nu }_\mu _\nu \right]\stackrel{}{D}^\nu +\overline{\psi }\widehat{G}\psi \overline{B}\widehat{G}^1B.`$ ( 3.1)
From here, we omit the color indices for brevity. $`B=(B_1,B_2)^T`$ is a two component auxiliary baryon field, whose components correspond to scalar and axial-vector channels; $`B_1D\chi `$ and $`B_2\stackrel{}{\tau }\stackrel{}{D}_\mu \gamma ^\mu \gamma ^5\chi `$. In Eq. (3.1) we have introduced matrix notations as
$`\psi `$ $`=`$ $`\left(\begin{array}{c}D\chi \\ \stackrel{}{D}_\mu \stackrel{}{\tau }\gamma ^\mu \gamma ^5\chi \end{array}\right),\overline{\psi }=\left(\begin{array}{cc}\overline{\chi }D^{},& \overline{\chi }\stackrel{}{D}_\mu ^{}\stackrel{}{\tau }\gamma ^\mu \gamma ^5\end{array}\right),`$ ( 3.5)
$`\widehat{G}`$ $`=`$ $`\left(\begin{array}{cc}G_S& v\\ v& G_A\end{array}\right).`$ ( 3.8)
Through the hadronization procedure, the quark and diquark fields are eliminated and a meson-baryon Lagrangian in the $`\text{tr}\mathrm{log}`$ form is obtained as
$$=\overline{B}\widehat{G}^1B+i\text{Tr}\mathrm{ln}(1A).$$
( 3.9)
Here the matrix $`A`$ is defined by
$`A`$ $`=`$ $`\left(\begin{array}{cc}a_{11}& a_{12}\\ a_{21}& a_{22}\end{array}\right),`$ (3.10c)
$`a_{11}`$ $`=`$ $`\mathrm{\Delta }^T\overline{B}_1SB_1,`$ (3.10d)
$`a_{12}`$ $`=`$ $`\mathrm{\Delta }^T\overline{B}_2\tau ^i\gamma ^\nu \gamma ^5SB_1,`$ (3.10e)
$`a_{21}`$ $`=`$ $`(\mathrm{\Delta }_{\rho \nu }^{lj})^T\overline{B}_1S\gamma ^\mu \gamma ^5\tau ^jB_2,`$ (3.10f)
$`a_{22}`$ $`=`$ $`(\mathrm{\Delta }_{\rho \nu }^{lj})^T\overline{B}_2\gamma ^\nu \gamma ^5\tau ^iS\gamma ^\mu \gamma ^5\tau ^jB_2.`$ (3.10g)
where $`S`$, $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }_{\mu \nu }`$ are the propagators of the quark, scalar diquark and axial-vector diquark, respectively. The expansion of the tr log in Eq. (3.9) gives various terms. The first term of the expansion gives the self-energies of the nucleons as
$`=\overline{B}\left(\begin{array}{cc}\mathrm{\Sigma }_S(p)& 0\\ 0& \mathrm{\Sigma }_A(p)\end{array}\right)B{\displaystyle \frac{1}{|\widehat{G}|}}\overline{B}\left(\begin{array}{cc}G_A& v\\ v& G_S\end{array}\right)B,`$ ( 3.15)
where $`|\widehat{G}|=\text{det}\widehat{G}=G_SG_Av^2`$. The scalar and axial-vector diquark contributions to the self-energies, $`\mathrm{\Sigma }_S`$ and $`\mathrm{\Sigma }_A`$, are shown diagrammatically in Fig. 1 and are computed as
$`\mathrm{\Sigma }_S(p)`$ $`=`$ $`iN_c{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2M_S^2}\frac{/p/k+m_q}{(pk)^2m_q^2}},`$ (3.16a)
$`\mathrm{\Sigma }_A(p)`$ $`=`$ $`iN_c{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{k^\mu k^\nu /M_A^2g^{\mu \nu }}{k^2M_A^2}\delta _{ij}\gamma _\nu \gamma _5\tau _j\frac{/p/k+m_q}{(pk)^2m_q^2}\tau _i\gamma _\mu \gamma _5}.`$ (3.16b)
Here $`N_c`$ is the number of colors. Since the self-energies Eqs. (3.16) are divergent, we regularize them by the three momentum cutoff scheme.
The self-energies $`\mathrm{\Sigma }_S`$ and $`\mathrm{\Sigma }_A`$ may be decomposed into the scalar and vector parts as
$`\mathrm{\Sigma }_S(p_0){\displaystyle \frac{1}{|\widehat{G}|}}G_A`$ $`=`$ $`Z_S^1(p_0\gamma ^0a_S),`$ (3.17a)
$`\mathrm{\Sigma }_A(p_0){\displaystyle \frac{1}{|\widehat{G}|}}G_S`$ $`=`$ $`Z_A^1(p_0\gamma ^0a_A),`$ (3.17b)
where we employ the nucleon rest frame $`p_\mu =(p_0,\stackrel{}{0})`$. The bare baryon fields $`B_{1,2}`$ are now renormalized as
$`\left(\begin{array}{c}B_1\\ B_2\end{array}\right)=\left(\begin{array}{c}\sqrt{Z_S}B_1^{}\\ \sqrt{Z_A}B_2^{}\end{array}\right),`$ ( 3.22)
with which the Lagrangian (3.15) is reduced to
$$=\overline{B}^{}(p_0\gamma ^0\widehat{M})B^{},$$
( 3.23)
where the mass matrix $`\widehat{M}`$ is given as
$`\widehat{M}=\left(\begin{array}{cc}a_S& \sqrt{Z_SZ_A}\frac{v}{|\widehat{G}|}\\ \sqrt{Z_SZ_A}\frac{v}{|\widehat{G}|}& a_A\end{array}\right).`$ ( 3.26)
When there is no mixing interaction $`(v=0)`$, $`B_{1,2}^{}`$ become the physical baryon fields and $`a_{S,A}`$ the physical masses. On the contrary, in the presence of the mixing, the physical states are obtained after the diagonalization of the mass matrix by an unitary transformation:
$`B^{}`$ $`=`$ $`U^{}N,`$ ( 3.27)
$`U\widehat{M}U^{}`$ $`=`$ $`\left(\begin{array}{cc}M_1& 0\\ 0& M_2\end{array}\right).`$ ( 3.30)
One finds
$$=\overline{N_1}(p_0\gamma ^0M_1)N_1+\overline{N_2}(p_0\gamma ^0M_2)N_2,$$
( 3.31)
where the physical eigenvalues $`M_{1,2}`$ and eigenvectors $`N=(N_1,N_2)^T`$ are obtained as
$`M_{1,2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[a_S+a_A\pm \sqrt{(a_Sa_A)^2+4Z_SZ_A\left({\displaystyle \frac{v}{|\widehat{G}|}}\right)^2}\right],`$ ( 3.32)
$`N_1`$ $`=`$ $`\mathrm{cos}\varphi B_1^{}+\mathrm{sin}\varphi B_2^{}`$ ( 3.33)
$`N_2`$ $`=`$ $`\mathrm{sin}\varphi B_1^{}+\mathrm{cos}\varphi B_2^{},`$ ( 3.34)
and the mixing angle $`\varphi `$ is given by
$$\mathrm{tan}2\varphi =\frac{2\sqrt{Z_SZ_A}v}{(a_Aa_S)|\widehat{G}|}$$
( 3.35)
## 4 Numerical results
For numerical calculations, let us first discuss model parameters. The constituent mass of the $`ud`$ quarks $`m_q`$ and the three momentum cutoff $`\mathrm{\Lambda }`$ are determined in such a way that they reproduce meson properties in the NJL model. The masses of the diquarks may be also calculated in the NJL model . Here we choose $`m_q`$=390 MeV, $`M_S`$=650 MeV, $`M_A`$=1050 MeV and $`\mathrm{\Lambda }`$=600 MeV, which are within the reasonable range known from the previous study of diquarks in the NJL model. The mass difference $`M_AM_S`$ may be related to that of the nucleon and delta. In the quark-diquark models, the delta is expressed as a bound state of an axial-vector diquark and a quark, while the nucleon is a superposition of the two components:
$`|\mathrm{\Delta }`$ $``$ $`\stackrel{}{D}_\mu \chi ,`$
$`|N`$ $``$ $`\mathrm{cos}\theta D\chi +\mathrm{sin}\theta \stackrel{}{\tau }\stackrel{}{D}_\mu \gamma ^\mu \gamma _5\chi .`$
Therefore, in a simple additive picture, the $`N\mathrm{\Delta }`$ mass difference can be expressed as
$`M_\mathrm{\Delta }M_N=\mathrm{cos}^2\theta (M_AM_S).`$ ( 4.1)
In deriving this relation we have ignored possible interactions such as the binding effect of the quarks and diquarks, and the interaction between the two diquark channels. However, for a rough estimation we may use the relation (4.1) and the mixing angle $`\mathrm{cos}^2\theta \underset{}{>}1/2`$, assuming that the nucleon state is rather dominated by the scalar diquark component. This qualitatively justifies the mass difference $`M_AM_S400`$ MeV that we adopt. The masses of the two nucleon states are then studied in the following two cases.
(i) In the first case, $`G_S`$ and $`G_A`$ are fixed such that the binding energies of the two quark-diquark bound states become 50 MeV when there is no coupling $`v`$. The resulting coupling strengths are $`G_S54`$ and $`G_A5.9`$, which generates the masses $`M_1=0.99`$ GeV and $`M_2=`$1.39 GeV at $`v=`$ 0 GeV<sup>-1</sup>. In the previous works, this binding energy was chosen in order to obtain a reasonable size of the nucleon. The masses and the mixing angle $`\varphi `$ of the two nucleon states are then calculated as functions of the coupling strength $`v`$, which are shown in Figs. 2. The effect of the $`v`$-coupling appears not only in the off-diagonal elements but also in the diagonal elements of the mass matrix (3.26). The effect of the $`v`$-coupling on the diagonal elements is from the term of $`|\widehat{G}|`$ in Eqs. (3.17), which is shown by dashed lines of Fig. 2. As is shown in the figure, the $`v`$-coupling acts repulsively to both the diagonal elements $`a_{S,A}`$. The non-zero off-diagonal elements then split $`M_1`$ and $`M_2`$ as is shown by the solid lines of Fig. 2. The off-diagonal coupling decreases the mass of the nucleon and increases that of the heavier state, while both the two diagonal elements increase as $`v`$ is increased. This is why these two contributions cancel each other for the nucleon, but they are enhanced for the heavier state. We find that $`M_1=0.99`$ GeV and $`M_2=`$ 1.44 GeV when $`v9`$ GeV<sup>-1</sup>. The mass of the second nucleon is close to that of the Roper resonance $`N`$(1440). At this strength of $`v`$, as is shown in the right panel of Fig. 2, the mixing angle is rather small, $`\varphi \underset{}{<}10`$ degree. Even at the small mixing angle, the effect of the axial-vector component is significant, since the self-energy $`\mathrm{\Sigma }_A`$ is much larger than $`\mathrm{\Sigma }_S`$ .
(ii) As have been mentioned, the masses $`M_{1,2}`$ depend not only on the off-diagonal elements of the mass matrix (3.26), but also on the diagonal elements. In the second choice we determine the parameters $`G_S`$ and $`G_A`$ such that they always produce the same amount of the binding energy of 50 MeV as the coupling strength $`v`$ is varied, or $`a_S`$ and $`a_A`$ are fixed at $`a_{S,A}=m_q+M_{S,A}50`$ MeV. In this way we can study the effect of the off-diagonal coupling $`v`$ just as in a simple two level problem with fixed values of the diagonal elements, which is shown in Fig.3. In this case dependence of $`M_{1,2}`$ on the values of $`v`$ is larger than that of the case (i). We find that $`M_1=`$ 0.94 GeV, $`M_2`$=1.44 GeV and $`\varphi `$ =18 degree at $`v`$ 22 GeV<sup>-1</sup>.
In both the cases of (i) and (ii), we can obtain the reasonable mass $`M_2`$=1.4 $``$1.5 GeV. Hence we identify the second state with the Roper resonance.
The present identification of the Roper resonance is very much different from the conventional picture; in the quark model, it is described as an excited state of $`2\mathrm{}\omega `$ with $`(n,l)=(1,0)`$, where $`(n,l)`$ are the principle and angular momentum quantum numbers of the harmonic oscillator. The excitation energy of such a state is as high as $`2\mathrm{}\omega 1`$ GeV for the oscillator parameter $`\omega 0.5`$ GeV, and many mechanisms have been proposed to lower the energy . In the present picture the two nucleons are described as quark-diquark bound states, but with different diquarks with different spins, isospins and masses. In a quark model picture, the two diquarks correspond to the quark pair of the $`\rho \rho `$ and $`\lambda \lambda `$ mixed symmetry (see Appendix B). In the limit of SU(4)<sub>SF</sub> symmetry, these state can not be independent degrees of freedom due to the Pauli principle when constructing the nucleon state of $`J^P=1/2^+`$, only their symmetric combination is allowed. In the present case, because the strong correlation between the quarks violates the SU(4)<sub>SF</sub> symmetry, the two states can be independent. In the picture of the harmonic oscillator basis, the two quarks of the diquarks are in the ground state but with being correlated. The energy difference is therefore supplied not by the difference in the single particle energies, but by the residual but significant correlation between the quarks. In any event, the scale of the energy of such correlations is expected to be of order of a few to several hundreds MeV. Concerning the wave functions, the two nucleon states of the quark-diquark model may be given in terms of quarks as
$`B_1^{}`$ $`=`$ $`D\chi [[1/2,1/2]^0,1/2]^{1/2},`$
$`B_2^{}`$ $`=`$ $`D\chi [[1/2,1/2]^1,1/2]^{1/2}.`$
From these structure, we may say that the Roper resonance appears as a spin partner of the nucleon, which has different internal spin structure.
## 5 Summary
In this paper, we have studied the nucleon and Roper resonance using the chiral quark-diquark model. It was shown that in the chiral symmetric construction of the model, there appear two states as the physical states with the quantum numbers of the nucleon. It was also shown that if we identified the low-lying state with the nucleon, the mass of the higher state became about 1.4 $``$ 1.5 GeV, with model parameters chosen so as to reproduce the mass difference of the nucleon and delta. Hence we consider this explanation of the Roper resonance is natural, although our current result have some dependence on the model parameters. Although some authors have studied the Roper resonance by using the lattice QCD the nature of the Roper resonance is still puzzle. The explanation here is different from many conventional approach, it may reveal the nature of the Roper resonance. One interesting aspect is the size of the Roper resonance. In our picture, the nucleon is expressed as a superposition of the quark and scalar diquark bound state and the quark and axial-vector diquark bound state, while the Roper resonance is another superposition orthogonal to the nucleon state. Then, the difference between the wave-functions of the quark and scalar diquark bound state and that of the quark and axial-vector diquark bound state would make that of the nucleon and Roper resonance different. It is interesting to study the effect of the change in the wave-functions on various properties, such as radii, the decay width and so on.
## Acknowledgments
We thank V. Dmitrasinovic and N. Ishii for fruitful discussions and useful advice.
## Appendix A Chiral transformation of baryon operators
Here we explicitly show the transformation of the baryon operators under the group SU(2)$`{}_{L}{}^{}\times `$ SU(2)<sub>R</sub>. The constituent quark field is defined through the Weinberg rotation
$$\chi =\xi _5q,$$
( A.1)
where $`\xi _5=\mathrm{exp}(i\gamma _5\stackrel{}{\tau }\stackrel{}{\mathrm{\Phi }}/2f_\pi )`$ with the pion field $`\stackrel{}{\mathrm{\Phi }}`$. Using the right- and left-handed spinors $`q_R`$ and $`q_L`$ of linear chiral representations, Eq. (A.1) is written as
$`\chi =\left(\begin{array}{c}\xi q_R\\ \xi ^{}q_L\end{array}\right),`$ ( A.4)
where $`\xi =\mathrm{exp}(i\stackrel{}{\tau }\stackrel{}{\mathrm{\Phi }}/2f_\pi )`$. Here $`\xi `$ and $`\xi _5`$ are related to each other
$`\xi _5+\xi _5^{}`$ $`=`$ $`\xi +\xi ^{},`$ (A.5a)
$`\xi _5\xi _5^{}`$ $`=`$ $`\gamma _5(\xi \xi ^{}).`$ (A.5b)
Under the group SU(2)$`{}_{L}{}^{}\times `$ SU(2)<sub>R</sub>, $`q_R`$, $`q_L`$ and $`\xi `$ are transformed as
$`q_{R,L}`$ $``$ $`g_{R,L}q_{R,L},`$ ( A.6)
$`\xi `$ $``$ $`g_L\xi h^{}=h\xi g_R^{},`$ ( A.7)
then the transformations of the constituent quark field and its Dirac and charge conjugated fields are given as
$`\chi `$ $``$ $`h\chi ,`$ (A.8a)
$`\overline{\chi }`$ $``$ $`\overline{\chi }h^{},`$ (A.8b)
$`\stackrel{~}{\chi }`$ $``$ $`\stackrel{~}{\chi }h^{}.`$ (A.8c)
Then, the two diquark states are transformed as
$`\stackrel{~}{\chi }\chi `$ $``$ $`\stackrel{~}{\chi }\chi ,`$ (A.9a)
$`\stackrel{~}{\chi }\gamma _\mu \gamma _5\tau ^a\chi `$ $``$ $`R^{ab}(x)\stackrel{~}{\chi }\gamma _\mu \gamma _5\tau ^b\chi .`$ (A.9b)
Here $`R^{ab}(x)`$ is the three dimensional rotation matrix, which is defined by $`h^{}\tau ^ah=R^{ab}(x)\tau ^b`$. Finally, the transformations of the baryon operators Eqs.(2.4) are obtained as
$`D\chi `$ $``$ $`h(x)D\chi ,`$ (A.10a)
$`D_\mu ^a\gamma ^\mu \gamma _5\tau ^a\chi `$ $``$ $`h(x)D_\mu ^a\gamma ^\mu \gamma _5\tau ^a\chi .`$ (A.10b)
## Appendix B Non Relativistic limit
It is useful to consider the connection to the non-relativistic quark model. The interaction Eq. (2.2) may be rewritten as
$`_{int}`$ $`=`$ $`a\overline{\chi }(\mathrm{cos}\theta D^{}+{\displaystyle \frac{1}{3}}\mathrm{sin}\theta \gamma ^\mu \gamma ^5\stackrel{}{\tau }\stackrel{}{D}_\mu ^{})(\mathrm{cos}\theta D+{\displaystyle \frac{1}{3}}\mathrm{sin}\theta \stackrel{}{\tau }\stackrel{}{D}_\mu \gamma ^\mu \gamma ^5)\chi `$ ( B.1)
$`+`$ $`b\overline{\chi }(\mathrm{sin}\theta D^{}+{\displaystyle \frac{1}{3}}\mathrm{cos}\theta \gamma ^\mu \gamma ^5\stackrel{}{\tau }\stackrel{}{D}_\mu ^{})(\mathrm{sin}\theta D+{\displaystyle \frac{1}{3}}\mathrm{cos}\theta \stackrel{}{\tau }\stackrel{}{D}_\mu \gamma ^\mu \gamma ^5)\chi ,`$
with the relations between the parameters
$`a\mathrm{cos}^2\theta +b\mathrm{sin}^2\theta `$ $`=`$ $`G_S,`$ (B.2a)
$`a\mathrm{sin}^2\theta +b\mathrm{cos}^2\theta `$ $`=`$ $`9G_A,`$ (B.2b)
$`(ab)\mathrm{sin}\theta \mathrm{cos}\theta `$ $`=`$ $`3v.`$ (B.2c)
Here we note that the baryon fields, in terms of a quark and a diquark, are related to the nucleon wave-functions of the constituent quark model by
$`D\chi =2\varphi _\rho \chi _\rho ,`$ (B.3a)
$`\stackrel{}{D}_\nu \stackrel{}{\tau }\gamma ^\nu \gamma _5\chi =6\varphi _\lambda \chi _\lambda ,`$ (B.3b)
in the non-relativistic limit, where $`\varphi _\rho ,\varphi _\lambda `$ and $`\chi _\rho ,\chi _\lambda `$ are the standard three quark spin and isospin wave-functions . By the use of these expressions, the interaction Eq. (B.1) is reduced to
$`_{int}`$ $`=`$ $`2a(\mathrm{cos}\theta \varphi _\rho \chi _\rho +\mathrm{sin}\theta \varphi _\lambda \chi _\lambda )^{}(\mathrm{cos}\theta \varphi _\rho \chi _\rho +\mathrm{sin}\theta \varphi _\lambda \chi _\lambda )`$ ( B.4)
$`+`$ $`2b(\mathrm{sin}\theta \varphi _\rho \chi _\rho +\mathrm{cos}\theta \varphi _\lambda \chi _\lambda )^{}(\mathrm{sin}\theta \varphi _\rho \chi _\rho +\mathrm{cos}\theta \varphi _\lambda \chi _\lambda ).`$
If we take $`\mathrm{tan}\theta =1`$, we realize SU(4)<sub>SF</sub> symmetry for the first term in Eq. (B.4), while the second term become the SU(4)<sub>SF</sub> forbidden state. Since we employ the diquark correlations violating SU(4)<sub>SF</sub> symmetry, the SU(4)<sub>SF</sub> forbidden state can be included. Then, the two channel type interaction Eq. (2.2) gives the two bound states of a quark and a diquark.
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# The exact solution of the Riemann problem in relativistic MHD with tangential magnetic fields
## 1 Introduction
The decay of a discontinuity separating two constant initial states (Riemann problem) has played a very important role in the development of numerical codes for classical (Newtonian) hydrodynamics after the pioneering work of Godunov (1959). Nowadays, most modern high-resolution shock-capturing methods (LeVeque 1992) are based on the exact or approximate solution of Riemann problems between adjacent numerical cells and the development of efficient Riemann solvers has become a research field in numerical analysis in its own (see, e.g., Toro 1997). The success of high-resolution shock-capturing methods in many areas of computational fluid dynamics has triggered their extension to classical magnetohydrodynamics (e.g., Brio & Wu 1988; for an up-to-date discussion of the issue, see Balsara 2004).
As in other fields in physics, during the last two decades astrophysics, relativity and cosmology have become computational sciences. Modeling and understanding fluid dynamics in astrophysical scenarios is now a key part in research projects involving supernovae, relativistic jets, neutron star instabilities, or accretion onto compact objects, who share a common distinctive feature: either special or general relativity effects are relevant. With this motivation, Riemann solvers have been introduced in numerical relativistic hydrodynamics since the beginning of the nineties (Martí, Ibáñez & Miralles 1991). Presently, the use of high-resolution shock-capturing methods based on Riemann solvers is considered the best strategy to solve the equations of relativistic hydrodynamics which has caused the rapid development of Riemann solvers for both special and general relativistic hydrodynamics (see, e.g., Martí & Müller 2003, Font 2003).
The finding by Martí & Müller (1994) of the analytical solution for initial states where the flow is normal to the initial discontinuity boosted the efforts to develop exact Riemann solvers for relativistic hydrodynamics. Pons, Martí & Müller (2000) extended the domain of solutions to problems with arbitrary initial velocities. Later, Rezzolla & Zanotti (2001), for purely normal flow, and Rezzolla, Zanotti & Pons (2003), for the general case, proposed a new procedure to find the solution of the Riemann problem that uses the relativistically invariant relative velocity between the unperturbed initial states. However, up to date, no analytical or exact solution of the equations of the Riemann problem in relativistic magnetohydrodynamics has been derived.
The equations of both classical and relativistic magnetohydrodynamics form a non-strictly hyperbolic system. A consequence of the non-strict hyperbolicity of the MHD are the degeneracies in the wave speeds (that lead, e.g., to compound waves, admisible solutions of the planar MHD that involve intermediate shocks), which must be handled analytically with care and hinder the development of both exact and approximate Riemann solvers using the characteristic information for the MHD equations. In the relativistic case the difficulties in the development of such solvers are increased by the higher non-linearity of the RMHD system of equations. The characteristic structure of the equations of RMHD is analyzed in the books by Lichnerowicz (1967) and Anile (1989). Approximate Riemann solvers using the characteristic information have been developed by Romero et al. (1996), for the same particular magnetohydrodynamic configuration considered here, and by Balsara (2001), Komissarov (1999) and Koldoba, Kuznetsov & Ustyugova (2002), for the general case. A number of analytical solutions involving only shocks, rarefactions and Alfven waves (Komissarov 1999, Komissarov 2003) have also been recently derived.
In this paper we describe the solution of the Riemann problem for the particular case in which the flow speed has two non-vanishing components and the magnetic field is orthogonal to them. Besides this, we force the flow to have dependence on one spatial coordinate taken along one of the two non-vanishing velocity components. For this particular setup, the Riemann structure degenerates to only three waves making the solution attainable. The solution reveals interesting and distinct properties of RMHD and could serve as a guide in the way to the general RMHD Riemann solution.
In this paper we closely follow the structure and notation used in Pons, Martí & Müller (2000; hereafter PMM). The paper is organized as follows. Section 2 collects the relevant equations. Sections 3 and 4 describe, respectively, the flow across rarefactions and shocks setting the ingredients for the Riemann solution, which is discussed in Sect. 5. Conclusions are gathered in Sect. 6.
## 2 Equations
Let $`J^\mu `$, $`T^{\mu \nu }`$ and $`F^{\mu \nu }`$ ($`\mu ,\nu =0,1,2,3`$) be the components of the density current, the energy–momentum tensor and the Maxwell dual tensor of an ideal magneto-fluid, respectively
$$J^\mu =\rho u^\mu $$
(2.1)
$$T^{\mu \nu }=\rho \widehat{h}u^\mu u^\nu +\eta ^{\mu \nu }\widehat{p}b^\mu b^\nu $$
(2.2)
$$F^{\mu \nu }=u^\mu b^\nu u^\nu b^\mu $$
(2.3)
where $`\rho `$ is the proper rest-mass density, $`\widehat{h}=1+ϵ+p/\rho +b^2/\rho `$ is the specific enthalpy including the contribution from the magnetic field ($`b^2`$ stands for $`b^\mu b_\mu `$), $`ϵ`$ is the specific internal energy, $`p`$ the thermal pressure, $`\widehat{p}=p+b^2/2`$ the total pressure, and $`\eta ^{\mu \nu }=diag(1,1,1,1)`$ the Minkowski metric in Cartesian coordinates. Throughout the paper we use units in which the speed of light is $`c=1`$.
The four-vectors representing the fluid velocity and the magnetic field in the fluid rest frame, $`u^\mu `$ and $`b^\mu `$, satisfy the conditions $`u^\mu u_\mu =1`$ and $`u^\mu b_\mu =0`$, and there is an equation of state $`p=p(\rho ,ϵ)`$ that closes the system. All the discussion will be valid for a general equation of state but results will be shown for an ideal gas, for which $`p=(\gamma 1)\rho ϵ`$, where $`\gamma `$ is the adiabatic exponent.
The equations of ideal RMHD correspond to the conservation of rest mass and energy-momentum, and the Maxwell equations. In flat space-time and Cartesian coordinates, these equations read:
$$J_{,\mu }^\mu =0$$
(2.4)
$$T_{,\mu }^{\mu \nu }=0$$
(2.5)
$$F_{,\mu }^{\mu \nu }=0$$
(2.6)
We consider a particular case in which the flow speed has two components and the magnetic field is orthogonal to them. Besides this, we force the flow to have dependence on one spatial coordinate ($`x`$) taken along one of the two non-vanishing velocity components. Specifically we set $`u^\mu =W(1,v^x,0,v^z)`$, $`b^\mu =(0,0,b,0)`$, where $`W`$ is the flow Lorentz factor. With these restrictions, the above system can be written as a system of conservation laws
$$\frac{𝐔}{t}+\frac{𝐅}{x}=0$$
(2.7)
where
$$𝐔=(D,\widehat{S}^x,\widehat{S}^z,\widehat{\tau },B)^\mathrm{T}$$
(2.8)
is the state vector of conserved quantities and
$$𝐅=(Dv^x,\widehat{S}^xv^x+\widehat{p},\widehat{S}^zv^x,\widehat{S}^x,Bv^x)^\mathrm{T}$$
(2.9)
is the corresponding vector of fluxes, with
$$D=\rho W$$
(2.10)
$$\widehat{S}^i=\rho \widehat{h}W^2v^i,(i=x,z)$$
(2.11)
and
$$\widehat{\tau }=\rho \widehat{h}W^2\widehat{p}$$
(2.12)
being the rest-mass, momentum and total energy densities, and
$$B=bW$$
(2.13)
the $`y`$-component of the magnetic field as measured in the laboratory frame. Hence, according to these equations, the particular initial configuration chosen together with the imposed symmetry prevent the generation of new components of the velocity and magnetic field.
It is worth noting that for the particular configuration chosen, the term $`(b^\mu b^\nu )_{,\mu }`$ appearing in the equation of conservation of the stress-energy tensor, vanishes and the RMHD equations reduce to the purely hydrodynamical case with the only contributions from the magnetic field appearing in the pressure and specific enthalpy, and an additional continuity equation for the evolution of the transversal magnetic field. This fact is considered in the Appendix where we explore the possibility of including the magnetic effects of the present configuration in the definition of the equation of state.
According to the previous discussion, the magnetized flow under consideration falls in one of the two degeneracies of the RMHD system (Degeneracy I; Komissarov 1999), for which a description in terms of just three characteristic waves (namely the entropy wave and the two fast magnetosonic waves) is adequate. Turning now towards the solution of the Riemann problem in this particular case, the discontinuity in the initial states breaks down into a couple of left and right propagating rarefaction waves (self-similar continuous flows) and/or shocks and a central tangential discontinuity across which the total pressure, $`\widehat{p}`$, is constant. Both thermal and total pressure increase at fast magnetosonic shocks. Hence, we would use the comparison of the total pressure at the two limiting states to select between shocks and rarefaction waves.
## 3 Flow across rarefactions
Rarefaction waves are self-similar solutions of the equations that depend on $`x`$ and $`t`$ only through the combination $`\xi =x/t`$. By imposing such a dependence in system (2.7) above we obtain the following set of equations
$$(v^x\xi )\frac{d\rho }{d\xi }+\rho (1+v^xW^2(v^x\xi ))\frac{dv^x}{d\xi }+\rho W^2v^z(v^x\xi )\frac{dv^z}{d\xi }=0$$
(3.14)
$$W^2\rho \widehat{h}(v^x\xi )\frac{dv^x}{d\xi }+b(1v^x\xi )\frac{db}{d\xi }+(1v^x\xi )\frac{dp}{d\xi }=0$$
(3.15)
$$W^2\rho \widehat{h}(v^x\xi )\frac{dv^z}{d\xi }v^zb\xi \frac{db}{d\xi }v^z\xi \frac{dp}{d\xi }=0$$
(3.16)
$$\frac{dp}{d\xi }=hc_s^2\frac{d\rho }{d\xi }$$
(3.17)
$$b(1+v^xW^2(v^x\xi ))\frac{dv^x}{d\xi }+bW^2v^z(v^x\xi )\frac{dv^z}{d\xi }+(v^x\xi )\frac{db}{d\xi }=0,$$
(3.18)
similar to the system obtained in PMM. The quantity $`c_s`$ is the sound speed, defined by
$$c_s=\sqrt{\frac{1}{h}\frac{p}{\rho }|_s}$$
(3.19)
where $`s`$ is the specific entropy and $`h=1+\epsilon +p/\rho `$, the specific enthalpy.
Non-trivial similarity solutions exist only if the determinant of system (3.14)-(3.18) vanish. This leads to the condition
$$\xi ^\pm =\frac{v^x(1\omega ^2)\pm \omega \sqrt{(1v^2)\left[1v^2\omega ^2(v^x)^2(1\omega ^2)\right]}}{1v^2\omega ^2}$$
(3.20)
where $`\omega ^2=c_s^2+v_A^2c_s^2v_A^2`$ and $`v_A^2=b^2/\rho \widehat{h}`$ is the Alfvén velocity. The plus and minus signs correspond to rarefaction waves propagating to the left $`_{}`$ and right $`_{}`$, respectively. Note that the values of $`\xi `$ reduce to those obtained in PMM by replacing $`\omega `$ by $`c_s`$ (i.e., $`b=0`$).
From system (3.14)-(3.18), after some algebraic manipulations, we get:
$$W^2\rho \widehat{h}(v^x\xi )\frac{dv^x}{d\xi }+b(1v^x\xi )\frac{db}{d\xi }+(1v^x\xi )\frac{dp}{d\xi }=0$$
(3.21)
$$\frac{dp}{d\xi }=hc_s^2\frac{d\rho }{d\xi }$$
(3.22)
$$\widehat{h}Wv^z=\mathrm{constant}$$
(3.23)
$$\frac{b}{\rho }=\mathrm{constant}.$$
(3.24)
Now, using (3.22) and (3.24) to eliminate the differentials of $`b`$ and $`\rho `$, and defining $`\widehat{}=b/\rho `$, the ODE (3.21) can be rewritten as
$$\frac{dv^x}{dp}=\frac{\left(1+\frac{\widehat{}^2\rho }{hc_s^2}\right)}{\rho \widehat{h}W^2}\frac{(1\xi v^x)}{(\xi v^x)},$$
(3.25)
in complete analogy with Eq. (3.20) in Rezzolla, Zanotti & Pons (2003) in the case $`b=0`$. The analogy can also be extended to the solution procedure. If we define $`\widehat{𝒜}=\widehat{h}Wv^z`$ then, from Eq. (3.23),
$$(v^z)^2=\widehat{𝒜}^2\left(\frac{1(v^x)^2}{\widehat{h}^2+\widehat{𝒜}^2}\right),$$
(3.26)
which allows to eliminate the dependence on $`v^z`$ in (3.20). Now, from the definition of the Lorentz factor, it can be derived
$$W^2=\frac{\widehat{h}^2+\widehat{𝒜}^2}{\widehat{h}^2(1(v^x)^2)}$$
(3.27)
and, after some algebra,
$$\frac{1\xi v^x}{\xi v^x}=\pm \frac{\sqrt{\widehat{h}^2+\widehat{𝒜}^2(1\omega ^2)}}{\widehat{h}\omega }$$
(3.28)
(where $`\omega `$ is the positive root of $`\omega ^2`$). Finally, substituting these last two expressions in (3.25), we get
$$\frac{dv^x}{1(v^x)^2}=\pm \frac{(1+\frac{\widehat{}^2\rho }{hc_s^2})\sqrt{\widehat{h}^2+\widehat{𝒜}^2(1\omega ^2)}}{\widehat{h}^2+\widehat{𝒜}^2}\frac{dp}{\rho \omega }.$$
(3.29)
The left hand side of this expression can be integrated analytically and the right hand side involves only thermodynamical variables and constants. Considering that in a Riemann problem the state ahead of the rarefaction wave is known, the integration of (3.29) allows one to connect the states ahead ($`a`$) and behind ($`b`$) the rarefaction wave. The normal velocity behind the rarefaction wave can be directly obtained as
$$v_b^x=\mathrm{tanh}\widehat{𝒞},$$
(3.30)
where
$$\widehat{𝒞}=\frac{1}{2}\mathrm{log}\left(\frac{1+v_a^x}{1v_a^x}\right)\pm _{p_a}^{p_b}\frac{(1+\frac{\widehat{}^2\rho }{hc_s^2})\sqrt{\widehat{h}^2+\widehat{𝒜}^2(1\omega ^2)}}{\widehat{h}^2+\widehat{𝒜}^2}\frac{dp}{\rho \omega }.$$
(3.31)
The differential of $`p`$ in the last integral is taken along the adiabats of the equation of state. The isentropic character of rarefaction waves fixes the entropy to that of state $`a`$, $`s_a`$. Having this in mind, the ODE can be integrated, the solution being only a function of $`p_b`$. This can be stated in compact form as
$$v_b^x=_{}^a(p_b).$$
(3.32)
It is interesting to have an expression for the normal velocity inside the rarefaction wave in terms of the total pressure. This expression can be built taking into account that, from the definition of $`\widehat{p}`$,
$$d\widehat{p}=dp+\widehat{}^2\rho d\rho $$
(3.33)
and that along an adiabat,
$$dp=hc_s^2d\rho ,$$
(3.34)
which combined result in
$$d\widehat{p}=\left(1+\frac{\widehat{}^2\rho }{hc_s^2}\right)dp.$$
(3.35)
Substitution in Eq. (3.31) gives
$$\widehat{𝒞}=\frac{1}{2}\mathrm{log}\left(\frac{1+v_a^x}{1v_a^x}\right)\pm _{\widehat{p}_a}^{\widehat{p}_b}\frac{\sqrt{\widehat{h}^2+\widehat{𝒜}^2(1\omega ^2)}}{\widehat{h}^2+\widehat{𝒜}^2}\frac{d\widehat{p}}{\rho \omega }.$$
(3.36)
Note that the previous expression is identical to the one derived by Rezzolla, Zanotti & Pons (2003) in the case of (non-magnetized) relativistic hydrodynamics after removing the hats and substituting $`\omega `$ by the sound speed. The corresponding compact notation for the function giving $`v_b^x`$ in terms of $`\widehat{p}_b`$ will be
$$v_b^x=\widehat{}_{}^a(\widehat{p}_b).$$
(3.37)
Function $`_{}^a(p)`$ is shown in Fig. 1, for different values of the invariant $`\widehat{}`$, the various branches of the curves corresponding to rarefaction waves propagating towards or away from $`a`$. Rarefaction waves move towards (away from) $`a`$, if the pressure inside the rarefaction is smaller (larger) than $`p_a`$. The last assertion also applies for the total pressure, $`\widehat{p}`$. In a Riemann problem the state $`a`$ is ahead of the wave and only those branches corresponding to waves propagating towards $`a`$ in Fig. 1 must be considered. Moreover, one can discriminate between waves propagating towards the left and right by taking into account that the initial left (right) state can only be reached by a wave propagating towards the left (right). The addition of a transverse magnetic field in the limiting state forces the value of the normal velocity within the rarefaction wave to have larger absolute values. This effect is a consequence of the fact that the absolute value of the slope of the function $`_{}^a(p)`$ at $`p_a`$, $`|_{}^a^{}(p_a)|`$, is an increasing function of $`\widehat{}`$ (or $`b_a`$). Moreover it can be easily proved that
$$|_{}^a^{}(p_a;\widehat{}\mathrm{})|=\frac{1}{c_{sa}}|_{}^a^{}(p_a;\widehat{}=0)|,$$
(3.38)
where $`c_{sa}`$ is the sound speed at state $`a`$. Hence the sound speed at state $`a`$ limits in practice the range of values of the normal velocity within the rarefaction wave (see Fig. 1 where the curves corresponding to $`\widehat{}=3`$ and $`10`$ almost coincide). Another consequence of the previous result is that in the extreme case for which $`c_{sa}`$ tends to the light speed, the presence of a transverse magnetic field in the limiting state would have no practical effect on the rarefaction wave. The effect of the magnetic field must be combined with the one coming from the presence of tangential velocities in state $`a`$ which operates in the opposite direction (see PMM).
## 4 Jumps across shocks.
If $`\mathrm{\Sigma }`$ is a hyper-surface in Minkowski space time across which $`\rho `$, $`u^\mu `$, $`T^{\mu \nu }`$ and $`F^{\mu \nu }`$ are discontinuous, the Rankine-Hugoniot conditions are given by (Lichnerowicz 1967, Anile 1989)
$$[\rho u^\mu ]n_\mu =0,$$
(4.39)
$$[T^{\mu \nu }]n_\nu =0,$$
(4.40)
$$\left[F^{\mu \nu }\right]n_\nu =0$$
(4.41)
where $`n_\mu `$ is the unit normal to $`\mathrm{\Sigma }`$, and where we have used the notation
$$[G]=G_aG_b,$$
(4.42)
$`G_a`$ and $`G_b`$ being the boundary values of $`G`$ on the two sides of $`\mathrm{\Sigma }`$.
Considering $`\mathrm{\Sigma }`$ as the hyper-surface in four-dimensional space describing the evolution of a shock wave normal to the $`x`$-axis, the unitarity of $`n_\nu `$ allows one to write it as
$$n^\nu =W_s(V_s,1,0,0),$$
(4.43)
where $`V_s`$ is interpreted as the coordinate velocity of the surface that defines the position of the shock wave and $`W_s`$ is the Lorentz factor of the shock,
$$W_s=\frac{1}{\sqrt{1V_s^2}}.$$
(4.44)
Equations (4.39) and (4.41) allow one to introduce two invariants across the shock
$$jW_sD_a(V_sv_a^x)=W_sD_b(V_sv_b^x),$$
(4.45)
$$fW_sB_a(V_sv_a^x)=W_sB_b(V_sv_b^x),$$
(4.46)
where $`B=bW`$. Quantity $`j`$ represents the mass flux across the shock and according to our definition, $`j`$ is positive for shocks propagating to the right (the same convention as the one used in Martí & Müller 1994, and PMM). Dividing Eq. (4.45) by Eq. (4.46), we get that the quantity $`B/D`$ (or, equivalently, $`b/\rho `$) is constant across shocks as it was through rarefaction waves.
Next, the Rankine-Hugoniot conditions (4.39), (4.40) can be written in terms of the conserved quantities $`D`$, $`\widehat{S}^j`$ and $`\widehat{\tau }`$, and $`j`$ as follows
$$\left[v^x\right]=\frac{j}{W_s}\left[\frac{1}{D}\right]$$
(4.47)
$$\left[\widehat{p}\right]=\frac{j}{W_s}\left[\frac{\widehat{S}^x}{D}\right]$$
(4.48)
$$\left[\frac{\widehat{S}^z}{D}\right]=0$$
(4.49)
$$\left[v^x\widehat{p}\right]=\frac{j}{W_s}\left[\frac{\widehat{\tau }}{D}\right].$$
(4.50)
Now from Eq. (4.49) we have that the quantity $`\widehat{h}Wv^z`$ is constant across the shock, as it is through rarefactions.
We note that in deriving equations (4.47)–(4.50) we have made use of the fact that the mass flux is nonzero across a shock. The conditions across a tangential discontinuity imply continuous total pressure and normal velocity (by setting $`j=0`$ in equations (4.47), (4.48) and (4.50)), and an arbitrary jump in the tangential velocity and transverse magnetic field.
Our aim now is to write $`v_b^x`$, the normal flow speed in the post-shock state, as a function of the post-shock pressure $`p_b`$. As a first step, we write $`v_b^x`$ as a function of $`p_b^{}`$, $`j`$ and $`V_s`$ (and the preshock state, $`a`$). Given the complete analogy between the jump conditions Eqs. (4.45), (4.47)-(4.50), and the corresponding expressions in PMM, we write
$$v_b^x=\left(\widehat{h}_aW_av_a^x+\frac{W_s(\widehat{p}_b\widehat{p}_a)}{j}\right)\left(\widehat{h}_aW_a+(\widehat{p}_b\widehat{p}_a)\left(\frac{W_sv_a^x}{j}+\frac{1}{\rho _aW_a}\right)\right)^1.$$
(4.51)
The dependence on the magnetic field in the pre- and post-shock states is hidden in the definitions of $`\widehat{p}`$ and $`\widehat{h}`$.
The shock speed $`V_s`$ can be eliminated using the definition of mass flux to obtain
$$V_s^\pm =\frac{\rho _a^2W_a^2v_a^x\pm |j|\sqrt{j^2+\rho _a^2W_a^2(1v_{a}^{x}{}_{}{}^{2})}}{\rho _a^2W_a^2+j^2}$$
(4.52)
where $`V_s^+`$ ($`V_s^{}`$) corresponds to shocks propagating to the right (left).
Proceeding in the same way as in PMM (i.e., $`\left[T^{\mu \nu }\right]n_\nu \{(\widehat{h}u_\mu )_a+(\widehat{h}u_\mu )_b\}=0`$) to derive the Taub’s adiabat, we can now obtain
$$\left[\widehat{h}^2\right]=\left(\frac{\widehat{h}_b}{\rho _b}+\frac{\widehat{h}_a}{\rho _a}\right)\left[\widehat{p}\right],$$
(4.53)
i.e, the Lichnerowicz’s adiabat (Anile 1989) particularized to our special setup. Figure 2 represents the function
$$^a(\widehat{h};\widehat{p}_b)\widehat{h}^2\widehat{h}_a^2\left(\frac{\widehat{h}}{\rho (\widehat{h},\widehat{p}_b)}+\frac{\widehat{h}_a}{\rho _a}\right)(\widehat{p}_b\widehat{p}_a)$$
(4.54)
for an ideal gas equation of state although the general shape of the curve (positive asymptotic branches; negative value for $`\widehat{h}=1`$) is independent of the equation of state. The (unique) root at the right of $`\widehat{h}=1`$ defines the thermodynamical post-shock state.
Equation (4.53) together with the definitions of $`\widehat{p}`$ and $`\widehat{h}`$, the equation of state and the constancy of $`b/\rho `$ through the shock, allows to write $`\rho _b`$ as a function of $`p_b`$ and the preshock state $`a`$. Next, multiplying (4.40) by $`n_\mu `$ and using the definition of relativistic mass flux one obtains
$$j^2=\frac{\left[\widehat{p}\right]}{\left[\widehat{h}/\rho \right]}.$$
(4.55)
Using the positive (negative) root of $`j^2`$ for shock waves propagating towards the right (left), equation (4.55) allows one to obtain the desired relation between the post-shock normal velocity $`v_b^x`$ and the post-shock pressure $`p_b`$. In a compact way the relation reads
$$v_b^x=𝒮_{}^a(p_b).$$
(4.56)
Alternatively, the relation can be written as a function of $`\widehat{p}_b`$
$$v_b^x=\widehat{𝒮}_{}^a(\widehat{p}_b).$$
(4.57)
Let us note that the expressions used to build up the function $`\widehat{𝒮}_{}^a`$, namely Eqs. (4.51), (4.52), (4.53) and (4.55), are formally identical to those corresponding to the pure (i.e., non-magnetized) relativistic hydrodynamical case. The difference appears in the definition of the function $`\rho =\rho (\widehat{h},\widehat{p})`$, which leads to different roots of the function $`^a(\widehat{h};\widehat{p}_b)`$, Eq. (4.54). Function $`𝒮_{}^a(p)`$ is shown in Fig. 3, for different values of the invariant $`\widehat{}`$, the various branches of the curves corresponding to shock waves propagating towards or away from $`a`$. In order to select the relevant branch of the function $`𝒮_{}^a(p)`$ the same argumentation as in the case of rarefaction waves can be used (see §3). As in the case of rarefaction waves, the addition of a transverse magnetic field in the limiting state forces the value of the normal velocity in the pre/post shock state to have larger absolute values. Again, this effect must be combined with the one coming from the presence of tangential velocities in state $`a`$ (see PMM).
Once $`v_b^x`$ is know, $`v_b^z`$ can be obtained through
$$(v_b^z)^2=\widehat{𝒜}^2\left(\frac{1(v_b^x)^2}{\widehat{h}_b^2+\widehat{𝒜}^2}\right),$$
(4.58)
where we have defined $`\widehat{𝒜}=\widehat{h}_aW_av_a^z`$. Analogously, $`b_b=\widehat{}\rho _b`$ with $`\widehat{}=b_a/\rho _a`$.
## 5 The solution of the Riemann problem.
As discussed in § 2, for the particular case under consideration (magnetic field orthogonal to both the fluid velocity and the wave propagation direction), the time evolution of a Riemann problem with initial states $`L`$ (left) and $`R`$ (right) can be represented as:
$$LRL𝒲_{}L_{}𝒞R_{}𝒲_{}R$$
(5.59)
where $`𝒲`$ and $`𝒞`$ denote a (fast magnetosonic-) shock or rarefaction, and a contact discontinuity, respectively. The arrows ($``$ / $``$) indicate the direction (left / right) from which fluid elements enter the corresponding wave.
The solution of the Riemann problem consists in finding the intermediate states, $`L_{}`$ and $`R_{}`$, as well as the positions of the waves separating the four states (which only depend on $`L`$, $`L_{}`$, $`R_{}`$ and $`R`$). The functions $`𝒲_{}`$ and $`𝒲_{}`$ allow one to determine the functions $`v_R^x(\widehat{p})`$ and $`v_L^x(\widehat{p})`$, respectively. The pressure $`\widehat{p}_{}`$ and the flow velocity $`v_{}^x`$ in the intermediate states are then given by the condition
$$v_R^x(\widehat{p}_{})=v_L^x(\widehat{p}_{})=v_{}^x.$$
(5.60)
The functions $`v_S^x(\widehat{p})`$ are defined by
$$v_S^x(\widehat{p})=\{\begin{array}{cc}\widehat{}^S(\widehat{p})\hfill & \text{if }\widehat{p}\widehat{p}_S\hfill \\ \widehat{𝒮}^S(\widehat{p})\hfill & \text{if }\widehat{p}>\widehat{p}_S\text{ ,}\hfill \end{array}$$
(5.61)
where $`\widehat{}^S(\widehat{p})`$ ($`\widehat{𝒮}^S(\widehat{p})`$) denotes the family of all states which can be connected through a rarefaction (shock) with a given state $`S`$ ($`L,R`$) ahead of the wave. Once $`\widehat{p}_{}`$ and $`v_{}^x`$ have been obtained the remaining quantities can be computed.
Figure 4 shows the solution of a particular Riemann problem for different values of the magnetic field $`b=0,1.0,2.0,4.0`$ in the initial states. The crossing point of any two lines gives the pressure and the normal velocity in the intermediate states.
It must be noted that the resolution of the Riemann problem under consideration can be formally done in the same way as the pure (relativistic) hydrodynamical problem with some modifications (see the Appendix). First of all, $`p`$ and $`h`$ have to be replaced by $`\widehat{p}`$ and $`\widehat{h}`$. Secondly, the sound speed in the integrand of Eq. (3.31) has to be replaced by $`\omega `$. Finally, the equation of state that provides the rest-mass density as a function of the pressure and the enthalpy, $`\rho (\widehat{h},\widehat{p})`$, has to be modified to include the contributions from the magnetic field. Given the parallelism between the present particular (relativistic) magnetohydrodynamical case and the purely hydrodynamical one, the effects concerning the smooth transition from one wave pattern to another when the tangential velocities in the initial states are changed (Rezzolla & Zanotti 2002, Rezzolla, Zanotti & Pons 2003) will extend to the present case for fixed initial values of the magnetic field. Hence we concentrate in the effects on the solution induced by varying the initial magnetic fields. Figure 5 shows the solution of a Riemann problem with a) vanishing magnetic field, and b) $`b_R=0.8`$. Whereas in the purely hydrodynamical case the Riemann solution gives rise to a left propagating rarefaction wave and right propagating contact and shock waves, the case with non vanishing magnetic field leads to a couple of shock waves and a left propagating contact discontinuity. The reason for this qualitative change (rarefaction/shock to shock/shock) can be found in the increase of the total pressure in the initial right state of the magnetized case. This increases the total pressure in the intermediate states at the two sides of the contact discontinuity. When this pressure becomes larger than that at the initial left state then a left propagating shock instead a rarefaction is produced.
Also noticeable from Fig. 5 is the increase of velocity of the shock propagating towards the right in the magnetized case. The increase of the velocity of propagation of fast waves for increasing magnetic fields (approaching the light speed in the fluid rest frame for strong enough fields, much larger than equipartition) is well-known in RMHD (e.g., Anile 1989). In the particular magnetic problem under consideration it leads to the increase of velocity of propagation of rarefaction heads and tails, and of shocks and can be understood as follows. Equation (3.20) gives the propagation speed of the head/tail of a right/left propagating rarefaction wave on a given state. Taking $`v^x,v=0`$ in Eq. (3.20), one gets $`|\xi ^\pm |=\omega `$ and $`\omega 1`$ for large enough $`b`$. One should remember that for purely hydrodynamical (relativistic) flows, $`|\xi ^\pm |`$ has the sound speed as limiting value. A similar result holds for shocks. For preshock states at rest, $`|V_s^\pm |1`$ as the magnetic field in the preshock state is increased. This can be seen by remembering that the shocks that appear in our configurations are super-magnetosonic, and that the fast magnetosonic speed in the preshock state ($`\omega _a`$) tends to light speed when $`b_a\mathrm{}`$.
Finally, let us note that the drift towards solutions involving only discontinuous waves (shock waves and rarefactions of negligible width) for increasing magnetic fields as concluded in the previous paragraph, is consistent with the fact that in the limit of strong magnetization, the equations of RMHD reduce to the equations of force-free degenerate electrodynamics, whose Riemann problem only involves (linearly degenerate) discontinuous waves (Komissarov 2002).
## 6 Summary and conclusions
We have obtained an exact solution of the Riemann problem for multidimensional relativistic magnetohydrodynamics in the particular case in which the magnetic field is normal to the fluid velocity. In this particular problem, the complex 7–wave pattern of RMHD is reduced to two fast magnetosonic waves and a contact discontinuity, which allows to use the same procedure as in the non magnetic case. Alternatively, we have shown that the problem can be understood as a purely RHD situation with a modified equation of state (see Appendix A for details).
Two interesting features arise from our results. First, for fixed initial thermodynamical states, it is possible to change continuously from one wave pattern to another (shock/shock, shock/rarefaction, rarefaction/rarefaction) analogously to what happens when tangential velocities are introduced (Rezzolla & Zanotti 2002). Secondly, we recover the result for RMHD flows with general magnetic field configurations stablishing the tendency of the wave speeds to the light speed when the magnetic field dominates the thermodynamical pressure and energy. For our particular configuration of the magnetic field, this results in fast moving shock waves and rarefaction waves in which the distance between the head and the tail is progressively reduced as the magnetic field increases. The drift towards solutions involving only discontinuous waves (shock waves and rarefactions of negligible width) for increasing magnetic fields is consistent with the fact that in the limit of strong magnetization, the equations of RMHD reduce to the equations of force-free degenerate electrodynamics, whose Riemann problem only involves discontinuous waves (Komissarov 2002).
In addition to the theoretical interest of our results by their own, having obtained an exact solution of the RMHD Riemann problem is relevant for the development of numerical codes. Up to now, in order to test the various algorithms and approximate Riemann solvers developed for numerical applications, one could only increase the spatial resolution and hope that the numerical solution converged to the physical one. Having an exact solution to compare with, even if it is just a particular case, allows for a more rigorous testing and error estimation. Last but not least, from a pedagogical point of view, it is more convenient to start understanding and solving a simpler case before attempting the solution of the full problem, which is the next natural extension of this work.
###### Acknowledgements.
It is a pleasure to thank L. Antón and L. Rezzolla for useful discussions and comments. Financial support for this research has been provided by Spanish MEC grant AYA2004-08067-C03. J.A.P. is supported by a Ramón y Cajal contract.
## A A hydrodynamical approach
Equations (2.7) are identical to those for the purely hydrodynamical case by replacing
$$p\widehat{p}=p+\frac{b^2}{2}$$
(A 62)
$$h\widehat{h}=h+\frac{b^2}{\rho },$$
(A 63)
indicating that a description of the present particular RMHD problem based on a purely hydrodynamical approach with a different equation of state may be possible. In this Appendix, we explore such a possibility, first suggested by Romero et al. (1996). The key point is to eliminate the magnetic field from the equations by building up a thermodynamically consistent equation of state including the effects of the magnetic field.
It follows from eqs. (2.7) that
$$\frac{D(b/\rho )}{Dt}=0,$$
(A 64)
where $`D/Dt`$ stands for the standard convective derivative, implying that the evolution of the fluid elements is along states keeping $`b/\rho =`$ constant. Then for a particular fluid element, $`\widehat{p}`$ and $`\widehat{h}`$ can be written as
$$\widehat{p}=p+\widehat{}^2\rho ^2/2$$
(A 65)
$$\widehat{h}=1+\epsilon +\frac{\widehat{p}}{\rho }+\widehat{}^2\rho /2,$$
(A 66)
where $`\widehat{}`$ is a constant.
Consistency of the two previous expressions with the definition of $`\widehat{h}`$ is fulfilled by defining a new specific internal energy, $`\widehat{\epsilon }`$,
$$\widehat{\epsilon }=\epsilon +\widehat{}^2\rho /2.$$
(A 67)
The evolution of the fluid elements in a perfect fluid is adiabatic. Hence now we look for the adiabats of the new equation of state, $`\widehat{p}=\widehat{p}(\rho ,\widehat{\epsilon })`$. The fact that the evolution of the fluid elements keeps $`\widehat{}`$ = constant draw us to consider that $`\widehat{}`$ is constant along the adiabats of the new equation of state, $`\widehat{s}=`$ constant. We shall use this fact to look for the desired relation between the entropies of the two equations of state, $`s`$ and $`\widehat{s}`$. To do this, we differentiate eq. (A 67) along transformations keeping $`\widehat{s}=`$ constant. We get
$$d\widehat{\epsilon }=d\epsilon +\widehat{}^2d\rho /2.$$
(A 68)
On the other hand, according to the first law of thermodynamics, for an adiabatic transformation,
$$d\widehat{\epsilon }=\frac{\widehat{p}}{\rho ^2}d\rho .$$
(A 69)
Substitution of $`\widehat{p}`$ in the previous expression gives
$$d\widehat{\epsilon }=\frac{p}{\rho ^2}d\rho +\widehat{}^2d\rho /2.$$
(A 70)
Finally, comparison with eq. (A 68) leads to
$$d\epsilon =\frac{p}{\rho ^2}d\rho ,$$
(A 71)
which is formally identical to the variation of internal energy in an adiabatic transformation $`s=`$ constant. Taking into account that the differentials were taken along the adiabats of the new equation of state, the conclusion is that the entropy in the new equation of state must be a function of the entropy in the original equation of state only.
Now the sound speed of the new equation of state,
$$\widehat{c}_s=\sqrt{\frac{1}{\widehat{h}}\frac{\widehat{p}}{\rho }|_{\widehat{s}}}$$
(A 72)
can be derived. The substitution of $`\widehat{p}`$ following eq. (A 65) and the equivalence of the adiabats of the two equations of state leads to
$$\widehat{c}_s=\sqrt{\frac{h}{\widehat{h}}c_s^2+\frac{\widehat{}^2\rho }{\widehat{h}}},$$
(A 73)
where $`h`$ and $`c_s`$ stand for the enthalpy and the sound speed of the original equation of state, respectively. Finally, a bit of algebra allows us to write,
$$\widehat{c}_s=\sqrt{(1v_A^2)c_s^2+v_A^2},$$
(A 74)
where $`v_A`$ is the Alfven speed for our particular case in which only one component of the magnetic field is non-zero, $`v_A=b/\sqrt{\rho \widehat{h}}`$ $`\left(=\widehat{}\sqrt{{\displaystyle \frac{\rho }{h+\widehat{}^2\rho }}}\right)`$. Note that $`\widehat{c}_s`$ coincides with the quantity $`\omega `$ defined in § 2.
It is interesting to note that although the original equation of state could have a sound speed significantly smaller than light speed (e.g., $`1/\sqrt{3}`$, for an ultra-relativistic non-degenerate ideal gas), the sound speed of the new equation of state (that represents the true propagation speed of perturbations in our magnetized fluid) approaches the light speed for large enough values of $`\widehat{}`$ (or $`b`$). Finally, notice that the convexity of the EOS is ensured, since $`^2`$ is a positive defined quantity and therefore $`{\displaystyle \frac{^2\widehat{p}}{\rho ^2}}|_s>0`$.
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# 1 The Energy-Momentum
## 1 The Energy-Momentum
### 1.1 Introduction
This paper proves a positive energy-momentum theorem under the (well known in general relativity) dominant energy condition, for AdS-asymptotically hyperbolic manifolds. An AdS-asymptotically hyperbolic manifold is by definition a manifold $`(M,g,k)`$ such that at infinity, the Riemannian metric $`g`$ and the symmetric 2-tensor $`k`$ tend respectively to the metric and second fundamental form of a standard hyperbolic slice of Anti-de Sitter (AdS).
Chruściel and Nagy recently defined the notion of energy-momentum of an asymptotically hyperbolic manifold, which generalizes the analogous notion in the asymptotically flat case. Besides Chruściel and Herzlich recently proved a positive mass theorem for asymptotically hyperbolic spin Riemannian manifolds (with zero extrinsic curvature).
The aim of the present paper is to extend this result to the non-zero extrinsic curvature case.
### 1.2 Some Definitions and Notations
We consider a Lorentzian manifold $`N^{n+1}`$ and a Riemannian spacelike hypersurface $`M`$. Using geodesic coordinates along $`M`$, we shall write a neighbourhood of $`M`$ in $`N`$ as a subset of $`]ϵ,ϵ[\times M`$, endowed with the metric $`\gamma =\text{d}t^2+g_t`$. The Riemannian n-manifold $`M`$ has induced metric $`g_0=g`$ and second fundamental form $`k:=(\frac{1}{2}\frac{\text{d}}{\text{d}t}g_t)_{|t=0}`$. We assume that $`(M,g,k)`$ is AdS-asymptotically hyperbolic that is to say, the metric $`g`$ and the second fundamental form $`k`$ are asymptotic at infinity to the metric and the second fundamental form of a standard hyperbolic slice in AdS. More precisely we adopt the following
Definition. $`(M,g,k)`$ is said to be AdS-asymptotically hyperbolic if there exists some compact K, a positive number R and a homeomorphism $`MK^nB(0,R)`$ called a chart at infinity such that in this chart we have
$$\{\begin{array}{ccc}e:=gb=O(e^{\tau r}),\hfill & e=O(e^{\tau r}),\hfill & ^2e=O(e^{\tau r}),\hfill \\ k=O(e^{\tau r}),\hfill & k=O(e^{\tau r}),\hfill & \end{array}$$
for $`\tau >n/2`$ and where $``$ is taken with respect to the hyperbolic metric $`b=\text{d}r^2+\mathrm{sinh}^2rg_{𝕊^{n1}}`$ with $`g_{𝕊^{n1}}`$ the standard metric of $`𝕊^{n1}.`$
AdS space-time is merely denoted by $`(\stackrel{~}{N},\beta )`$. If one considers $`\stackrel{~}{N}`$ as $`^{n+1}`$ then we will write $`\beta =\text{d}t^2+b_t`$, with $`b_0=b`$ the hyperbolic metric.
The motivation for the definition of the energy-momentum comes from the study of the constraints map which by definition is
$$\begin{array}{cccc}\mathrm{\Phi }:\hfill & \times \mathrm{\Gamma }(S^2T^{}M)& \hfill & C^{\mathrm{}}(M)\times \mathrm{\Gamma }(T^{}M)\\ & (h,p)& \hfill & \left(\begin{array}{c}\text{Scal}^h+(tr_hp)^2\left|p\right|_h^2\\ 2\left(\delta _hp+\text{d}tr_hp\right)\end{array}\right),\end{array}$$
where $``$ is the set of Riemannian metrics on the manifold $`M`$. Let us denote by $`(\dot{h},\dot{p})`$ an infinitesimal deformation of $`(h,p)`$. Now if we take a couple $`(f,\alpha )C^{\mathrm{}}(M)\times \mathrm{\Gamma }(T^{}M)`$ then we compute
$`(f,\alpha ),(\mathrm{\Phi }(h+\dot{h},p+\dot{p})\mathrm{\Phi }(h,p))`$ $`=`$ $`\delta (f(\delta \dot{h}+\text{dtr}\dot{h})+i_f\dot{h}(\text{tr}\dot{h})\text{d}f+2i_\alpha \dot{p}2(\text{tr}\dot{p})\alpha )`$
$`+`$ $`\delta (<p,\dot{h}>\alpha +<h,\dot{h}>i_\alpha p2i_{i_\alpha p}\dot{h})`$
$`+`$ $`\text{d}\mathrm{\Phi }_{(h,p)}^{}(f,\alpha ),(\dot{h},\dot{p})+Q(f,\alpha ,h,p,\dot{h},\dot{k}),`$
where \<,\> is the metric extended to all tensors, $`\delta `$ is the $`h`$-divergence operator, $`\text{d}\mathrm{\Phi }_{(h,p)}^{}`$ is the formal adjoint of the linearized constraints map at the point $`(h,p)`$, traces are taken with respect to $`h`$ and $`Q(f,\alpha ,h,p,\dot{h},\dot{k})`$ is a remainder which is linear with respect to $`(f,\alpha )`$ and at least quadratic with respect to $`(\dot{h},\dot{p})`$. Now considering the constraints map along the hyperbolic space embedded in AdS, that is to say $`(h,k)=(b,0)`$ and $`(\dot{h},\dot{k})=(gb=e,k)`$ one finds
$`(f,\alpha ),(\mathrm{\Phi }(g,k)\mathrm{\Phi }(b,0))`$ $`=`$ $`\delta (f(\delta e+\text{dtr}e)+i_fe(\text{tr}e)\text{d}f+2i_\alpha k2(\text{tr}k)\alpha )`$
$`+`$ $`\text{d}\mathrm{\Phi }_{(b,0)}^{}(f,\alpha ),(e,k)+Q(f,\alpha ,b,k,e).`$
As a consequence if we assume that $`(M,g,k)`$ is AdS-asymptotically hyperbolic and if the function $`(f,\alpha ),(\mathrm{\Phi }(g,k)\mathrm{\Phi }(b,0))`$ is integrable on $`M`$ with respect to the measure $`\text{dVol}_b`$, then the energy-momentum $``$ can be defined as a linear form on Ker $`\text{d}\mathrm{\Phi }_{(b,0)}^{}`$
$$\text{}.$$
The integrand in the formula of $``$ is in index notation
$$f(e_{j,i}^ie_{i,j}^i)f^{,i}e_{ij}+(e_i^i)f_{,j}2\alpha ^ik_{ij}+2(k_i^i)\alpha _j,$$
where “,” stands for the $`b`$-derivatives and where $`h^i=b^{ij}h_j`$ for any tensor $`h`$. This integrand is the same as the one of Chruściel and Nagy since each Killing vector fields on AdS is decomposable into the sum of some normal and tangential components (with respect to a standard hyperbolic slice) which are in our case given by the couple $`(f,\alpha )`$ (see also ). More precisely, one can show, using Moncrief argument , that Ker $`\text{d}\mathrm{\Phi }_{(b,0)}^{}𝔎𝔦𝔩𝔩(\text{AdS})`$ where $`𝔎𝔦𝔩𝔩(\text{AdS})`$ denotes the Lie algebra of Killing vector fields on AdS, since it satisfies the Einstein equations with a (negative) cosmological constant. The isometry group of AdS is O(n,2), and thereby $`𝔎𝔦𝔩𝔩(\text{AdS})𝔰𝔬(n,2)N_b𝔰𝔬(n,1)N_b𝔎𝔦𝔩𝔩(^n)`$, where we have set $`N_b=\{fC^{\mathrm{}}(M)|\text{Hess}f=fb\}`$. It is well known , that the application
$$\begin{array}{ccc}^{n,1}\hfill & & N_b\hfill \\ y_k\hfill & & x_k:=y_{k|^n}\hfill \end{array},$$
(where $`(y_k)_{k=0}^n`$ are the standard coordinates) is an isometry, and the mass part of the energy-momentum is a linear form on $`N_b`$ which is causal and positively oriented as soon as $`\text{Scal}^gn(n1)=\text{Scal}^b`$. Remark that the sharpest integrability conditions in order to make $``$ well defined and invariant under asymptotic isometries have been found by Chruściel and Nagy still in . However for the sake of simplicity one can use instead of the integrability condition $`(f,\alpha ),(\mathrm{\Phi }(g,k)\mathrm{\Phi }(b,0))L^1(M,\text{dVol}_b)`$, the less general but more convenient condition $`\left|\mathrm{\Phi }(g,k)\mathrm{\Phi }(b,0)\right|e^rL^1(M,\text{dVol}_b)`$.
Remark. In the asymptotically flat situation, the energy-momentum is also a linear form on $`^{n,1}𝔰𝔬(n,1)`$ where the first component corresponds to translational isometries and the second one to rotations. This interpretation gave rise to the respective terminology of linear and angular momentum. In the AdS-asymptotically hyperbolic situation, one cannot identify some linear momentum in the decomposition $`𝔰𝔬(n,2)^{n,1}𝔰𝔬(n,1)`$, since the first component $`^{n,1}`$ in $`𝔰𝔬(n,2)`$ is not of translational nature. This whole first component of the energy is then called the mass functional and it only remains some angular momentum. Moreover physicists often call the limit of integrals $`(f,\alpha )`$ global charges and so the positive energy-momentum theorem could be consequently renamed global inequalities theorem. Some supplementary details on the physical interpretation of our result can be found in the forthcoming note by Chruściel and the author . $`\mathrm{}`$
### 1.3 Statement of the Theorems and Comments
As a matter fo fact, we know that, given a chart at infinity, $``$ can be considered as a vector of $`^{n,1}𝔰𝔬(n,1)`$ and will be denoted by $`M\mathrm{\Xi }`$. The vector $`M`$ is the mass part of $``$, and $`\mathrm{\Xi }`$ is the angular momentum. We will prove the existence of a Hermitian application
$$Q:\text{},$$
which has to be non-negative under some energy assumptions but nonetheless $`Q`$ quite difficult to explicite in general.
However, in dimension $`n=3`$, we can be more specific giving the explicite formula of $`Q`$ in terms of the components $`M𝔐\text{M}(2,)`$ (cf. section 2.4 for the definition of $`𝔐`$) and $`\mathrm{\Xi }𝔰𝔩(2,)`$ of the energy-momentum $``$. More precisely we will show that
$$Q=2\left(\begin{array}{cc}\widehat{M}& \mathrm{\Xi }\\ \mathrm{\Xi }^{}& M\end{array}\right),$$
where $`\widehat{M}`$ is the transposed comatrix of $`M`$. We will also treat the case where the slice $`M`$ has a compact inner boundary $`M`$ whose induced metric and second fundamental form are respectively denoted by $`\stackrel{˘}{g}`$ and $`\stackrel{˘}{k}`$. To this end, we have to define the vector field $`\stackrel{}{k}:=(\text{tr}\stackrel{˘}{k}+(n1))e_0+k(\nu )`$ along the boundary $`M`$. We can now state the
Positive Energy-Momentum Theorem. Let $`(M^n,g,k)`$ be an AdS-asymptotically hyperbolic spin Riemannian manifold satisfying the decay conditions stated in section 1.2 and the following conditions
(i) $`(f,\alpha ),(\mathrm{\Phi }(g,k)\mathrm{\Phi }(b,0))L^1(M,\text{dVol}_b)`$ for every $`(f,\alpha )N_b𝔎𝔦𝔩𝔩(M,b),`$
(ii) the relative version of the dominant energy condition (cf. section 2.2) holds, that is to say $`(\mathrm{\Phi }(g,k)\mathrm{\Phi }(b,0))`$ is a positively oriented causal (n+1)-vector along $`M`$,
(iii) in the case where M has a compact boundary $`M`$, we assume moreover that $`\stackrel{}{k}`$ is causal and positively oriented along $`M`$.
Then there exists a (hardly explicitable) map $`^{n,1}𝔰𝔬(n,1)\text{Herm}(C^d)`$ which sends, under the assumptions (i-iii), the energy-momentum on a non-negative Hermitian form Q.
Moreover, when n=3, we can explicite Q in terms of the components of the energy-momentum as described above.
Classical algebra results give the non negativity of each principal minors of $`Q`$ which provide a set of inequalities on the coefficients of $``$ that are explicitely written in the appendix in dimension $`n=3`$ (cf. section 5.1).
This result is new (even though many formal arguments where given by Gibbons, Hull and Warner in ) and based on the recent global charge definition of Chruściel and Nagy for AdS-asymptotically hyperbolic manifolds, which comes from the Hamiltonian description of General Relativity. Our approach is purely Riemannian, the Lorentzian connection and manifold introduced are auxiliary since everything is restricted to the Riemannian slice $`M`$ (contraryly to ). In the other hand, the positive mass theorem for Minkowski-asymptotically hyperbolic initial data sets of Chruściel, Jezierski and Łȩski is also different from ours, since their Riemannian hypersurface is supposed to be asymptotic at infinity to a standard hyperbolic slice of Minkowski space-time (in that case the extrinsic curvature does not tend to 0) . Then considering the translational Killing vector fields of Minkowski, they defined a hyperbolic 4-momentum (usually called Trautman-Bondi mass) and proved that it is timelike and future directed under the dominant energy condition (and some other technical assumptions).
Remark finally that our result extends the positive mass theorem of Chruściel and Herzlich in dimension $`n`$, since if one supposes that $`k=0`$ then we recover their result: the mass functional $`M`$ has to be time-like future directed.
As regards the rigidity part we have the
Theorem. Under the assumptions of the positive energy-momentum theorem, $`\text{tr}Q=0`$ implies that $`(M,g,k)`$ is isometrically embeddable in AdS<sup>n,1</sup>.
This result is optimal in the sense that one could not reasonably hope better than being able to embbed isometrically our triple $`(M,g,k)`$ in AdS.
Some additional but partial results will be proved (also for the Trautman-Bondi 4-momentum) in order to weaken the defining condition of rigidity.
### 1.4 Organisation of the Paper
In section 2, we give the necessary geometric background by recalling some basic facts on spinors and defining the Killing connection used in the remainder of the paper. We also prove the Bochner-Lichnerowicz-Weitzenböck-Witten formula with respect to our Killing connection and deduce an integration formula.
In section 3, we prove the positive energy-momentum theorem: we remark that the boundary contribution of the integrated Bochner-Lichnerowicz-Weitzenböck-Witten formula can be identified to the global charges $`(f,\alpha )`$, for some choices of $`(f,\alpha )`$. This can be done using the same ideas as in but in a Lorentzian situation, and extends the quite technical computations of in a non-trivial way, since the algebraic structures are different (spinors, Hermitian scalar product, gauge etc…) and since new terms (involving the extrinsic curvature) appeared and had to be identified. We then make the analysis of the Dirac operator (we also treat the case where $`M`$ has a compact boundary) which gives the non-negativity of the globlal charges $`(f,\alpha )`$ when the couple $`(f,\alpha )`$ comes from an imaginary Killing spinor of AdS<sup>n,1</sup>. Then we restrict to dimension $`n=3`$, and completely study the imaginary Killing spinors of AdS<sup>3,1</sup> in order to interpret the non-negativity of the global charges as the non-negativity of the Hermitian matrix $`Q`$ on $`^4`$.
Section 4 is devoted to the proof of the rigidity results.
The last section is an appendix which gives the non-negativity of $`Q`$ in dimension $`n=3`$ seen through its coefficients, and proves some rigidity results for the Trautman-Bondi mass .
## 2 Geometric Background
All the definitions and conventions of this section will be for any $`n3`$, where $`n`$ is the dimension of the AdS-asymptotically hyperbolic slice, except if the dimension is explicitely mentioned to be 3.
### 2.1 Connections and curvatures
$`,\overline{}`$ denote respectively the Levi-Civita connections of $`\gamma `$ and g. Let us take a spinor field $`\psi \mathrm{\Gamma }(\mathrm{\Sigma })`$ and a vector field $`X\mathrm{\Gamma }(TM)`$, then
$$\{\begin{array}{ccc}_X\psi & =\hfill & \overline{}_X\psi \frac{1}{2}k(X)e_0\psi \hfill \\ k(X),Y_\gamma & =\hfill & _XY,e_0_\gamma \hfill \end{array}.$$
In these formulae $``$ denotes the Clifford action with respect to the metric $`\gamma `$, and $`e_0=_t`$. We will use different notations when we have to make the difference between the Clifford action with respect to the metric $`\gamma `$ or $`\beta `$.
Definition. The Killing equation on a spinor field $`\tau \mathrm{\Gamma }(\mathrm{\Sigma })`$ is
$$\widehat{D}_X\tau :=D_X\tau +\frac{i}{2}X_\beta \tau =0X\mathrm{\Gamma }(TM),$$
where $`D`$ denotes the Levi-Civita connection of AdS along $`M`$. Such a $`\widehat{D}`$-parallel spinor field is called a $`\beta `$-imaginary Killing spinor and we denote $`\tau IKS(\mathrm{\Sigma })`$. In the same way, a $`\widehat{}`$-parallel spinor field (where $`\widehat{}_X:=_X+\frac{\text{i}}{2}X_\gamma `$) is called a $`\gamma `$-imaginary Killing spinor.
Notice that the equation $`\widehat{D}\tau =0`$ is neither the Killing equation in AdS nor in $`^n`$, but the Killing equation in AdS along $`^n`$ (in particular the imaginary Killing spinors considered here are not the one of , ).
Now if $`R,\widehat{R}`$ are the respective curvatures of $``$ and $`\widehat{}`$, we have the relation
$$\widehat{R}_{X,Y}=R_{X,Y}\frac{1}{4}(XYYX),$$
where we use the convention of for the curvature.
### 2.2 Bochner-Lichnerowicz-Weitzenböck-Witten Formula and the Dominant Energy Condition
From now on $`\left(e_k\right)_{k=0}^n`$ is an orthonormal basis at the point with respect to the metric $`\gamma `$. We define the Dirac-Witten operators
$$𝔇\psi =\underset{k=1}{\overset{n}{}}e_k_{e_k}\psi ,\widehat{𝔇}\psi =\underset{k=1}{\overset{n}{}}e_k\widehat{}_{e_k}\psi ,$$
where $`n`$ is the dimension of the spacelike slice.
Lemma.(Bochner-Lichnerowicz-Weitzenböck-Witten formula)
$$\widehat{𝔇}^{}\widehat{𝔇}=\widehat{}^{}\widehat{}+\widehat{},$$
where $`\widehat{}:=\frac{1}{4}\left(\text{Scal}^\gamma +n(n1)+4\text{Ric}^\gamma (e_0,e_0)+2e_0\text{Ric}^\gamma (e_0)\right).`$
Proof. The Dirac-Witten operator $`𝔇`$ is clearly formally self adjoint, and we have the classical Bochner-Lichnerowicz-Weitzenböck formula (cf. , for instance) $`𝔇^{}𝔇=𝔇^2=^{}+`$, where $`:=\frac{1}{4}\left(\text{Scal}^\gamma +4\text{Ric}^\gamma (e_0,e_0)+2e_0\text{Ric}^\gamma (e_0)\right)`$. We also know that $`\widehat{𝔇}=𝔇\text{i}\frac{n}{2}`$ and so we get
$$\widehat{𝔇}^{}\widehat{𝔇}=^{}++\frac{n^2}{4},$$
but finally remarking that $`\widehat{}^{}\widehat{}=^{}+\frac{n}{4}`$, we obtain our formula.$`\mathrm{}`$
We derive an integration formula from the Bochner-Lichnerowicz-Weitzenböck-Witten identity considering the 1-form $`\theta `$ on $`M`$ defined by $`\theta (X)=\widehat{}_X\psi +X\widehat{𝔇}\psi ,\psi _\gamma `$, where $`\psi `$ is a spinor field. Straightforward computations lead to the following $`g`$-divergence formula
$$\text{div}\theta =\widehat{𝔇}\psi ,\widehat{𝔇}\psi _\gamma \widehat{}\psi ,\psi _\gamma \widehat{}\psi ,\widehat{}\psi _\gamma .$$
Let $`S_r`$ the $`g`$-geodesic sphere of radius $`r`$ and centered in a point of M. The radius $`r`$ is supposed to be as large as necessary. We denote by $`M_r`$ the interior domain of $`S_r`$ and $`\nu _r`$ the (pointing outside) unit normal. Integrating our divergence formula over $`M_r`$ and using Stokes theorem, we get
$$_{M_r}\left|\widehat{𝔇}\psi \right|_\gamma ^2=_{M_r}\left(\left|\widehat{}\psi \right|_\gamma ^2+\widehat{}\psi ,\psi _\gamma \right)_{S_r}\widehat{}_{\nu _r}\psi +\nu _r\widehat{𝔇}\psi ,\psi _\gamma \text{dVol}_{S_r}.$$
Let us now consider the Einstein tensor $`G=\text{Ric}^\gamma \frac{1}{2}\text{Scal}^\gamma \gamma `$ with respect to the metric $`\gamma `$. The dominant energy condition says that the speed of energy flow of matter is always less than the speed of light. More precisely, for every positively oriented time-like vector field $`v`$, the energy-momentum current of density of matter $`G(v,.)^{\mathrm{}}`$ must be time-like or null, with the same orientation as $`v`$. The assumption we make in order to prove the positive energy-momentum theorem is a relative version of the dominant energy condition: $`\left(G\frac{n(n1)}{2}\gamma \right)(e_0)`$ is a positively oriented time-like or null vector along $`M`$. Some easy computations give
$`\text{Scal}^\gamma `$ $`=`$ $`2(G(e_0,e_0)\text{Ric}^\gamma (e_0,e_0))`$
$`e_0\text{Ric}^\gamma (e_0)`$ $`=`$ $`e_0G_{|TM}(e_0)\text{Ric}^\gamma (e_0,e_0),`$
where $`G_{|TM}(e_0)=_{k=1}^3G(e_0,e_k)e_k`$. Thereby
$`\widehat{}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(2G(e_0,e_0)+\left(n(n1)\right)+2e_0G_{|TM}(e_0))`$
$`=`$ $`{\displaystyle \frac{1}{2}}((G(e_0,e_0)+{\displaystyle \frac{n(n1)}{2}})e_0G_{|TM}(e_0))e_0`$
$`=`$ $`{\displaystyle \frac{1}{2}}((G(e_0,e_0){\displaystyle \frac{n(n1)}{2}}\gamma (e_0,e_0))e_0G_{|TM}(e_0))e_0`$
$`=`$ $`{\displaystyle \frac{1}{2}}(G{\displaystyle \frac{n(n1)}{2}}\gamma )(e_0)e_0.`$
Our assumption gives the non negativity of the spinorial endomorphism $`\widehat{}`$ that is to say $`\widehat{}\psi ,\psi 0`$ for every spinor field $`\psi `$.
Remark. We can express the dominant energy condition in terms of the constraints as in section 1.3 since
$$\left(G\frac{n(n1)}{2}\gamma \right)(e_0)=\frac{1}{2}\left(\mathrm{\Phi }(g,k)\mathrm{\Phi }(b,0)\right).$$
$`\mathrm{}`$
### 2.3 Spinorial Gauge
In the same way as Andersson and Dahl , but in a Lorentzian situation, we compare spinors in $`\mathrm{\Sigma }`$ (along $`M`$) with respect to the two different metrics $`\beta `$ and $`\gamma `$. This can be done according to as soon as the tubular neighbourhood of $`M`$ in $`N`$ is small enough. Consequently we suppose that both metrics are written in Gaussian coordinates $`\beta =\text{d}t^2+g_t,\gamma =\text{d}t^2+b_t`$ on $`]ϵ,+ϵ[\times M`$ for $`ϵ`$ small enough. We define the spinorial gauge $`𝒜\mathrm{\Gamma }(\text{End}(𝕋))`$ with the relations
$$\{\begin{array}{ccc}\gamma (𝒜X,𝒜Y)\hfill & =\hfill & \beta (X,Y)\hfill \\ \gamma (𝒜X,Y)\hfill & =\hfill & \gamma (X,𝒜Y)\hfill \end{array},$$
where $`𝕋`$ is $`TN`$ restricted to $`M`$. The first relation says that $`𝒜`$ sends $`\beta `$-orthonormal frames on $`\gamma `$-orthonormal frames whereas the second one means that the endomorphism $`𝒜`$ is symmetric. We notice that these relations are only satisfied along $`M=\left\{t=0\right\}`$ and can also be written in the following way
$$\{\begin{array}{ccc}\hfill 𝒜e_0& =\hfill & e_0\hfill \\ \hfill g(𝒜X,𝒜Y)& =\hfill & b(X,Y)\hfill \\ \hfill g(𝒜X,Y)& =\hfill & g(X,𝒜Y)\hfill \end{array}.$$
Consequently $`𝒜`$ is an application $`\text{P}_{\text{SO}_0(n,1)}(\beta )_{|M}\text{P}_{\text{SO}_0(n,1)}(\gamma )_{|M}`$, which can be covered by an application still denoted $`𝒜:\text{P}_{\text{Spin}_0(n,1)}(\beta )_{|M}\text{P}_{\text{Spin}_0(n,1)}(\gamma )_{|M}`$. This application carries $`\beta `$-spinors on $`\gamma `$-spinors so that we have the compatibility relation about the Clifford actions of $`\beta `$ and $`\gamma `$
$$𝒜(X_\beta \sigma )=(𝒜X)_\gamma (𝒜\sigma ),$$
for every $`X\mathrm{\Gamma }(𝕋),\sigma \mathrm{\Gamma }(\mathrm{\Sigma })`$ and where $`_\beta ,_\gamma `$ denotes the Clifford actions respectively of $`\beta `$ and $`\gamma `$. Remark that our gauge is more sophisticated that the one of since it deals with the trace of Lorentzian structures (metrics, spinors, Hermitian scalar product etc…) along the spacelike slice $`M`$.
We define a new connection $`\stackrel{~}{}X=𝒜(\overline{D}𝒜^1X)`$ along $`M`$. It is easy to check that $`\stackrel{~}{}`$ is $`g`$-metric and has torsion $`\stackrel{~}{T}(X,Y)=((\overline{D}_X𝒜)𝒜^1Y(\overline{D}_Y𝒜)𝒜^1X)`$. We extract some formulae for later use
$$2g(\stackrel{~}{}_XY\overline{}_XY,Z)=g(\stackrel{~}{T}(X,Y),Z)g(\stackrel{~}{T}(X,Z),Y)g(\stackrel{~}{T}(Y,Z),X).$$
Now we intend to compare the connexions $`\overline{}`$ and $`\stackrel{~}{}`$ on $`\mathrm{\Sigma }`$. $`(\sigma _s)_s`$ denotes the spinorial frame corresponding to the orthonormal frame $`(e_k)_{k=0}^n`$, and $`\overline{\omega },\stackrel{~}{\omega }`$ are the connection 1-forms respectively of $`\overline{}`$ and $`\stackrel{~}{}`$
$$\begin{array}{ccc}\overline{\omega }_{ij}\hfill & =\hfill & g(\overline{}e_i,e_j)\hfill \\ \stackrel{~}{\omega }_{ij}\hfill & =\hfill & g(\stackrel{~}{}e_i,e_j),\hfill \end{array}$$
and if we take a general spinor $`\phi =\phi ^s\sigma _s`$, their derivatives are given by
$`\overline{}\phi =\text{d}\phi ^s\sigma _s+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i<j}{}}\overline{\omega }_{ij}e_i_\gamma e_j_\gamma \phi `$
$`\stackrel{~}{}\phi =\text{d}\phi ^s\sigma _s+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i<j}{}}\stackrel{~}{\omega }_{ij}e_i_\gamma e_j_\gamma \phi ,`$
and as a consequence
$$(\overline{}\stackrel{~}{})\phi =\frac{1}{4}\underset{i,j=0}{\overset{n}{}}(\overline{\omega }_{ij}\stackrel{~}{\omega }_{ij})e_i_\gamma e_j_\gamma \phi .$$
### 2.4 Tangent and Spinor Bundles
In this paper, the model spaces AdS<sup>n,1</sup> and $`^n`$ are considered as symmetric spaces:
$$\text{},$$
so that every section of any natural fiber bundle above $`^n`$ can be seen as a function $`\text{Spin}_0(n,1)^d`$ which is $`\text{Spin}(n)`$–equivariant (with $`d`$ depending upon $`n`$). We can be more explicite when we take $`n=3`$ (this fact is due to the exceptional isomorphisms of Lie groups below)
$$\text{},$$
with $`\text{SU}(2)\text{Spin}(3)`$ and $`\text{SL}(2,)\text{Spin}_0(3,1)`$.
The spinor bundle of AdS is $`\mathrm{\Sigma }_{AdS}=\text{Spin}_0(3,2)\times _{\stackrel{~}{\rho }}^4`$, where $`\text{Spin}_0(3,2)`$ is the bundle of the $`\text{Spin}_0(3,1)`$-frames in AdS, and $`\stackrel{~}{\rho }`$ is the standard représentation of $`\text{SL}(2,)`$ on $`^4^2\overline{^2}^{}`$. In other words
$$\begin{array}{cccc}\stackrel{~}{\rho }:& \text{SL}(2,)& & \text{M}_4()\hfill \\ & \stackrel{~}{g}& & \left(\begin{array}{cc}\stackrel{~}{g}& 0\\ 0& (\stackrel{~}{g}^{})^1\end{array}\right)\hfill \end{array},$$
where $`A^{}=^t\overline{A},A\text{M}_2()`$. When we restrict this bundle to the hypersurface $`^3`$ we have $`\mathrm{\Sigma }=\text{SL}(2,)\times _{\stackrel{~}{\rho }_{|\text{SU}(2)}}^4`$.
Proposition. $`\mathrm{\Sigma }`$ and $`^3\times ^4`$ are isomorphic thanks to the following trivialisation:
$$\begin{array}{cccc}T:& \mathrm{\Sigma }& & ^3\times ^4\hfill \\ & \{\stackrel{~}{e},w\}& & ([\stackrel{~}{e}],\stackrel{~}{\rho }(\stackrel{~}{e})w)\hfill \end{array},$$
where $`\{\stackrel{~}{e},w\}`$ denotes the class of $`(\stackrel{~}{e},w)\text{SL}(2,)\times ^4`$ in $`\mathrm{\Sigma }`$ , and $`[\stackrel{~}{e}]`$ denotes the class of $`\stackrel{~}{e}\text{SL}(2,)`$ in $`^3=\text{SL}(2,)/\text{SU}(2)`$.
The construction of $`𝕋_{AdS}`$, the tangent bundle of AdS, is quite similar to the construction of the spinor bundle. Still noticing that the principal bundle of $`\text{SO}_0(3,1)`$-frames in AdS is isomorphic to $`\text{SO}_0(3,2)`$, we write $`𝕋_{AdS}=\text{SO}_0(3,2)\times _\rho ^4`$, where $`\rho `$ is the standard representation of $`\text{SO}_0(3,1)`$ on $`^4`$. By restriction to the hypersurface $`^3`$, we obtain $`𝕋=\text{SO}_0(3,1)\times _{\rho _{|\text{SO}(3)}}^4`$, where $`\text{SO}(3)`$ is by definition the isotropy group of $`f_0`$ if $`\left(f_k\right)_{k=0}^3`$ denotes the canonical basis of $`^4`$.
Proposition. $`𝕋`$ and $`^3\times ^4`$ are isomorphic thanks to the following trivialisation:
$$\begin{array}{cccc}T:& 𝕋& & ^3\times ^4\hfill \\ & \{e,u\}& & ([e],\rho (e)u)\hfill \end{array},$$
where $`\{e,u\}`$ denotes the class of $`(e,u)\text{SO}_0(3,1)\times ^4`$ in $`𝕋`$, and $`[e]`$ denotes the class of $`e\text{SO}_0(3,1)`$ in $`^3=\text{SO}_0(3,1)/\text{SO}(3)`$.
We are going to define the Clifford action on $`\mathrm{\Sigma }`$, in the same way as in . To this end, we denote by $`(^4,q)`$ the Minkowski space-time of signature (3,1), where $`q=\text{d}y_0^2+\text{d}y_1^2+\text{d}y_2^2+\text{d}y_3^2`$. This space is isometric to a subspace of $`\text{M}_2()`$ via
$$\begin{array}{cccc}\mathrm{\Lambda }:\hfill & (^4,q)& \hfill & 𝔐:=(\left\{A\text{M}_2()/A^{}=A\right\},\text{det})\\ & y=\left(y_i\right)_{i=0}^3& \hfill & \left(\begin{array}{cc}y_0+y_1\hfill & y_2+iy_3\hfill \\ y_2iy_3\hfill & y_0y_1\hfill \end{array}\right)\end{array}.$$
We have thus the following real vector space isomorphisms:
$$\begin{array}{ccc}\text{M}_2()\hfill & \hfill & 𝔲(2)𝔐\hfill \\ 𝔰𝔩_2()\hfill & \hfill & 𝔰𝔲(2)\left(𝔐𝔰𝔩_2()\right)\hfill \\ & \hfill & 𝔰𝔲(2)𝔊,\hfill \end{array}$$
and $`𝔊^3`$. In order to make the value of the sectional curvature of $`^3`$ equal to -1, when we consider $`^3=\text{SL}(2,)/\text{SU}(2)`$ as a symmetric space, we have to consider $`^4`$ endowed with $`4q`$ and not $`q`$, and consequently the embedding of the Clifford algebra $`\text{C}\mathrm{}_{3,1}`$ in $`\text{M}_4()`$ becomes
$$\mathrm{\Theta }:X𝔐\left(\begin{array}{cc}0& 2X\\ 2\widehat{X}& 0\end{array}\right),$$
where $`\widehat{X}`$ means the transposed comatrix of X.
It will be convenient to see $`𝕋`$ as $`\text{SL}(2,)\times _\mu 𝔐`$, where $`\mu `$ is the universal covering of $`\text{SO}_0(3,1)`$ by $`\text{SL}(2,)`$, which is given by:
$$\begin{array}{cccc}\mu :& \text{SL}(2,)& & \text{SO}_0(3,1)\\ & \stackrel{~}{g}& & \left(\stackrel{~}{g}:X𝔐\stackrel{~}{g}X\stackrel{~}{g}^{}\right)\end{array}.$$
We can now define the Clifford action. Let us take $`e\text{SO}_0(3,1)`$ and $`\stackrel{~}{e}\text{SL}(2,)`$ such that $`e=\mu (\stackrel{~}{e})`$. A vector $`X=X[e]`$ tangent at the point $`[e]=[\stackrel{~}{e}]^3`$, is a class $`\{e,u\}𝕋`$. A spinor $`\sigma =\sigma [\stackrel{~}{e}]`$ at the same point is likewise a class $`\{\stackrel{~}{e},w\}\mathrm{\Sigma }`$. The result of the Clifford action of $`X`$ on $`\sigma `$ is the spinor $`(X\sigma )[\stackrel{~}{e}]=\{e,u\}\{\stackrel{~}{e},w\}=\{\stackrel{~}{e},\mathrm{\Theta }(u)w\}`$. We define a sesquilinear inner product (not definite positive) $`(,)`$ on $`^4^2\overline{^2}^{}`$ as in $`(\xi ,\eta ):=\xi _1,\eta _2_^2+\xi _2,\eta _1_^2`$, where $`\xi =\left(\genfrac{}{}{0pt}{}{\xi _1}{\xi _2}\right),\eta =\left(\genfrac{}{}{0pt}{}{\eta _1}{\eta _2}\right)^4`$ and where $`,_^2`$ is the standard Hermitian product on $`^2`$. This induces a sesquilinear product on $`\mathrm{\Sigma }`$ by $`(\{\stackrel{~}{e},\xi \},\{\stackrel{~}{e},\eta \}):=(\xi ,\eta )`$. In the same way we define a scalar product on $`\mathrm{\Sigma }`$ setting
$$\begin{array}{ccc}\{\stackrel{~}{e},\xi \},\{\stackrel{~}{e},\eta \}\hfill & :=& (\frac{1}{2}f_0\{\stackrel{~}{e},\xi \},\{\stackrel{~}{e},\eta \})\hfill \\ & =& (\{\stackrel{~}{e},\frac{1}{2}\mathrm{\Theta }(f_0)\xi \},\{\stackrel{~}{e},\eta \})\hfill \\ & =& \xi ,\eta _^4,\hfill \end{array}$$
where $`,_^4`$ denotes the standard Hermitian product on $`^4`$.
Since $`\text{SL}(2,)`$ is the 2-sheeted covering of $`\text{SO}_0(3,1)`$, there exists a natural (left) action of $`\text{SL}(2,)`$ on $`\mathrm{\Sigma }`$ which is derived from the natural (left) action of $`\text{SO}_0(3,1)`$ on $`𝕋`$: the action of the group of the isometries of AdS preserving the slice $`^3`$ that is $`\stackrel{~}{g}\{\stackrel{~}{e},w\}=\{\stackrel{~}{g}\stackrel{~}{e},w\}`$, with $`\stackrel{~}{g}\text{SL}(2,)`$ and $`\sigma [\stackrel{~}{e}]=\{\stackrel{~}{e},w\}`$ a spinor at $`[\stackrel{~}{e}]`$. To have the action on a section $`\sigma \mathrm{\Gamma }(\mathrm{\Sigma })`$ we set as usual $`(\stackrel{~}{g}\sigma )[\stackrel{~}{e}]=\stackrel{~}{g}\sigma (\stackrel{~}{g}^1\stackrel{~}{e})`$.
## 3 Positive Energy-Momentum Theorem
In this section, the dimension will be $`n3`$ expect if $`n`$ is explicitely mentioned to be 3.
Moreover $`f`$ will denote a smooth cutoff function which is 0 on $`M`$ except on a small neighbourhood of the infinity boundary of $`M`$ where $`f1`$, and $`H(a)`$, $`H_{}(a)`$ are Hilbert spaces of spinor fields defined in the subsections below. We will prove the
Proposition. For every $`\sigma IKS(\mathrm{\Sigma })`$ there exists a unique $`\xi _0H(a)`$ (resp. $`H_{}(a)`$ if $`M`$ has a boundary) such that
$$\xi =f𝒜\sigma +\xi _0\text{Ker}\widehat{𝔇}(\text{resp.}\text{Ker}\widehat{𝔇}H_{}(a))\text{and}(V_\sigma ,\alpha _\sigma )0,$$
where $`V_\sigma =\sigma ,\sigma `$ and $`\alpha _\sigma (X)=Xe_0\sigma ,\sigma `$ .
In section 3.3, it will proved that the couple $`(V_\sigma ,\alpha _\sigma )`$ belongs to $`N_b𝔎𝔦𝔩𝔩(^n)`$ so that $`(V_\sigma ,\alpha _\sigma )`$ is actually well-defined.
The computations we will make in section 3.1 prove that if we integrate the Bochner-Lichnerowicz-Weitzenböck-Witten formula with an asymptotically imaginary Killing spinor $`f𝒜\sigma `$, then the boundary integrals tend to some global charge $`(V_\sigma ,\alpha _\sigma )`$, when $`r`$ goes to infinity. In fact it is still true if we perturb $`f𝒜\sigma `$ with a smooth compactly supported spinor field $`\xi _0`$ (that is to say if we consider $`f𝒜\sigma +\xi _0`$ instead of $`f𝒜\sigma `$). Actually we will show in section 3.2 that we can find a perturbation $`\xi _0`$ in a relevant Hilbert space such that $`\xi _0`$ has no contribution at infinity, and $`f𝒜\sigma +\xi _0`$ belongs to the kernel of $`\widehat{𝔇}`$.
This will naturally imply the non-negativity of $`(V_\sigma ,\alpha _\sigma )`$ when $`\sigma `$ is a $`\beta `$-imaginary Killing spinor. This is the reason why we focus on the study of the Killing equation in section 3.3 so as to interpret the non-negativity of the $`(V_\sigma ,\alpha _\sigma )`$.
### 3.1 Energy-Momentum and Imaginary Killing Spinors
The aim of this section is to show the
Proposition. Let $`\xi =f𝒜\sigma +\xi _0`$ , where $`\sigma IKS(\mathrm{\Sigma })`$ and $`\xi _0`$ is a compactly supported spinor field. Then we have
$`(V_\sigma ,\alpha _\sigma )`$ $`=`$ $`4\underset{r+\mathrm{}}{lim}{\displaystyle _{S_r}}\widehat{}_{𝒜\nu _r}\xi +𝒜\nu _r\widehat{𝔇}\xi ,\xi _\gamma `$
$`=`$ $`4{\displaystyle _M}\left(\left|\widehat{}\xi \right|_\gamma ^2+\widehat{}\xi ,\xi _\gamma \right)4{\displaystyle _M}\left|\widehat{𝔇}\xi \right|_\gamma ^2.`$
Remark that the only important data is the exact $`\beta `$-imaginary Killing spinor $`\sigma `$ involved in the definition of the couple $`(V_\sigma ,\alpha _\sigma )`$.
Proof. Remember that
$$_{M_r}\left|\widehat{𝔇}\psi \right|_\gamma ^2=_{M_r}\left(\left|\widehat{}\psi \right|_\gamma ^2+\widehat{}\psi ,\psi _\gamma \right)_{S_r}\widehat{}_{𝒜\nu _r}\psi +𝒜\nu _r_\gamma \widehat{𝔇}\psi ,\psi _\gamma \text{dVol}_{S_r},$$
where $`\nu _r`$ denotes the $`b`$-normal of $`S_r`$, $`e_0=_t`$ and we set $`e_1=𝒜\nu _r`$ for the remainder of the proof.
We have to work on the expression $`\widehat{}_{𝒜\nu _r}\psi +𝒜\nu _r_\gamma \widehat{𝔇}\psi ,\psi _\gamma `$ in order to identify the integrand used to compute the energy-momentum for some couple $`(f,\alpha )`$. We start with noticing that $`e_1_\gamma e_1_\gamma \widehat{}_{e_1}=\widehat{}_{e_1}=\widehat{}_{𝒜\nu _r}`$ so that
$$\widehat{}_{𝒜\nu _r}\psi +𝒜\nu _r_\gamma \widehat{𝔇}\psi =𝒜\nu _r_\gamma \left(\underset{j=2}{\overset{n}{}}e_j_\gamma \widehat{}_{e_j}\right)\psi .$$
From now on we work on
$$𝒜\nu _r_\gamma \left(\underset{j=2}{\overset{n}{}}e_j_\gamma \widehat{}_{e_j}\right)_\gamma \phi ,\phi _\gamma .$$
Let us take $`\sigma `$ a $`\beta `$-imaginary Killing spinor, that is to say a spinor field solution, by definition, of $`\widehat{D}_X\sigma =D_X\sigma +\frac{\text{i}}{2}X_\beta \sigma =0`$, for every vector field $`X\mathrm{\Gamma }(TM)`$. Consider $`f`$ a smooth cutoff function which is 0 on $`M`$ except on a compact neighbourhood of the infinity boundary of $`M`$ where $`f1`$. Then we have
$$\begin{array}{ccccc}\widehat{}_X(f𝒜\sigma )\hfill & =\hfill & \text{d}f(X)𝒜\sigma \hfill & +f\widehat{}_X(𝒜\sigma )\hfill & \\ & =\hfill & \text{d}f(X)𝒜\sigma \hfill & +f(\overline{}_X\stackrel{~}{}_X)(𝒜\sigma )\hfill & +f(\stackrel{~}{}_X+\frac{\text{i}}{2}X_\gamma \frac{1}{2}k(X)_\gamma e_0_\gamma )(𝒜\sigma ),\hfill \end{array}$$
but since $`\stackrel{~}{}_X(𝒜\sigma )=𝒜\overline{D}_X\sigma =\frac{\text{i}}{2}𝒜(X_\beta \sigma )=\frac{i}{2}(𝒜X)_\gamma (𝒜\sigma )`$, we obtain
$$\widehat{}_X(f𝒜\sigma )=\text{d}f(X)𝒜\sigma +f(\overline{}_X\stackrel{~}{}_X)(𝒜\sigma )\frac{1}{2}f\left(k(X)_\gamma e_0+\text{i}(𝒜Id)X\right)_\gamma (𝒜\sigma ),$$
that we restrict to the neighbourhood where $`f1`$
$$\widehat{}_X(𝒜\sigma )=(\overline{}_X\stackrel{~}{}_X)(𝒜\sigma )\frac{1}{2}(k(X)_\gamma e_0+\text{i}(𝒜Id)X)_\gamma (𝒜\sigma ).$$
As a consequence our boundary term becomes for $`r`$ great enough
$$\underset{j=2}{\overset{n}{}}𝒜\nu _r_\gamma e_j_\gamma ((\overline{}_{e_j}\stackrel{~}{}_{e_j})\frac{1}{2}(k(e_j)_\gamma e_0+\text{i}(𝒜Id)e_j)_\gamma )(𝒜\sigma ),𝒜\sigma _\gamma .$$
We will estimate this boundary term in several steps. From the decay assumptions stated section in 1.2, the gauge is supposed to be of the form $`𝒜=Id+B+O(|B|^2)`$, where B has the same decay to 0 as $`e=gb`$. In the following $`(ϵ_k=𝒜^1e_j)_{j=0}^n`$ is a $`\beta `$-orthonormal frame.
We begin with the easiest term
$`{\displaystyle \underset{j=2}{\overset{n}{}}}𝒜\nu _r_\gamma e_j_\gamma k(e_j)_\gamma e_0_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$ $`=`$ $`{\displaystyle \underset{j=2}{\overset{n}{}}}𝒜\nu _r_\gamma 𝒜ϵ_j_\gamma k(𝒜ϵ_j)_\gamma 𝒜ϵ_0_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$
$`=`$ $`{\displaystyle \underset{j=2}{\overset{n}{}}}\nu _r_\beta ϵ_j_\beta 𝒜^1k𝒜(ϵ_j)_\beta ϵ_0_\beta \sigma ,\sigma _\beta .`$
But we note that $`𝒜^1k𝒜=kBk+kB+O(|B|^2)`$. Now B has the same decay as $`k`$ so $`Bk+kB=O(|B|^2)`$, terms that we can neglect since the energy-momentum is computed by a limit procedure of integrals over large spheres. We conclude that $`𝒜^1k𝒜k`$, where for convenience the relation $``$ means that $`||`$ is at least a $`O(e^{2\tau r})`$ when $`r`$ goes to infinity. Moreover
$`\nu _r_\beta {\displaystyle \underset{j=2}{\overset{n}{}}}ϵ_j_\beta k(ϵ_j)`$ $`=`$ $`ϵ_1_\beta \left({\displaystyle \underset{j=1}{\overset{n}{}}}ϵ_j_\beta k(ϵ_j)ϵ_1_\beta k(ϵ_1)\right)`$
$`=`$ $`k(\nu _r)(\text{tr}_bk)\nu _r,`$
which implies
$`{\displaystyle \underset{j=2}{\overset{n}{}}}𝒜\nu _r_\gamma e_j_\gamma k(e_j)_\gamma e_0_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$ $``$ $`k(\nu _r)(\text{tr}_bk)\nu _r_\beta ϵ_0_\beta \sigma ,\sigma _\beta `$
$`=`$ $`\left(i_{\alpha _\sigma }k(\text{tr}_bk)\alpha _\sigma \right)(\nu _r),`$
where $`\alpha _\sigma (X)=X_\beta ϵ_0_\beta \sigma ,\sigma _\beta `$.
The second term we study is
$`\text{i}{\displaystyle \underset{j=2}{\overset{n}{}}}𝒜\nu _r_\gamma e_j_\gamma (𝒜Id)(e_j)_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$ $`=`$ $`\text{i}{\displaystyle \underset{j=2}{\overset{n}{}}}𝒜\nu _r_\gamma 𝒜ϵ_j_\gamma (𝒜Id)(𝒜ϵ_j)_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$
$`=`$ $`\text{i}{\displaystyle \underset{j=2}{\overset{n}{}}}\nu _r_\beta ϵ_j_\beta 𝒜^1(𝒜Id)𝒜(ϵ_j)_\beta \sigma ,\sigma _\beta `$
$``$ $`\text{i}{\displaystyle \underset{j=2}{\overset{n}{}}}\nu _r_\beta ϵ_j_\beta B(ϵ_j)_\beta \sigma ,\sigma _\beta ,`$
but thanks to the same property as above
$$\nu _r_\beta \underset{j=2}{\overset{n}{}}ϵ_j_\beta B(ϵ_j)=B(\nu _r)(\text{tr}_bB)\nu _r,$$
which induces
$`\text{i}{\displaystyle \underset{j=2}{\overset{n}{}}}𝒜\nu _r_\gamma e_j_\gamma (𝒜Id)(e_j)_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$ $``$ $`\text{i}B(\nu _r)(\text{tr}_bB)\nu _r_\beta \sigma ,\sigma _\beta `$
$`=`$ $`\left(i_{V_\sigma }B(\text{tr}_bB)\text{d}V_\sigma \right)(\nu _r),`$
where $`\text{d}V_\sigma (X)=\text{i}X_\beta \sigma ,\sigma _\beta `$.
The last term we have to study is certainly the most difficult (summation convention $`k\{2,3,\mathrm{},n\},l\{1,2,\mathrm{},n\},m\{1,2,\mathrm{},n\}`$)
$$\begin{array}{ccc}𝒜\nu _r_\gamma e_k_\gamma (\overline{}_{e_k}\stackrel{~}{}_{e_k})(𝒜\sigma ),𝒜\sigma _\gamma \hfill & =\hfill & \frac{1}{4}(\overline{\omega }_{lm}\stackrel{~}{\omega }_{lm})(e_k)𝒜\nu _r_\gamma e_k_\gamma e_l_\gamma e_m_\gamma (𝒜\sigma ),𝒜\sigma _\gamma \hfill \\ & =\hfill & \frac{1}{4}(\overline{\omega }_{lm}\stackrel{~}{\omega }_{lm})𝒜(ϵ_k)\nu _r_\beta ϵ_k_\beta ϵ_l_\beta ϵ_m_\beta \sigma ,\sigma _\beta \hfill \\ & =\hfill & \frac{1}{4}S\hfill \end{array}$$
$`S`$ $`=`$ $`{\displaystyle \underset{k,l,m=2}{\overset{n}{}}}(\overline{\omega }_{lm}\stackrel{~}{\omega }_{lm})(e_k)𝒜\nu _r_\gamma e_k_\gamma e_l_\gamma e_m_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$
$`+`$ $`2{\displaystyle \underset{k,l=2}{\overset{n}{}}}(\overline{\omega }_{1l}\stackrel{~}{\omega }_{1l})(e_k)𝒜\nu _r_\gamma e_k_\gamma e_1_\gamma e_l_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$
$`=`$ $`S_1+2S_2.`$
We will give estimates of each $`S_k`$, keeping in mind that they are real and that every term that is at least $`O(|B|^2)`$ can be neglected when $`r+\mathrm{}`$ for the computations of the global charge integrals.
Estimate of $`S_1`$
$$S_1=\underset{k,l,m=2}{\overset{n}{}}(\overline{\omega }_{lm}\stackrel{~}{\omega }_{lm})(e_k)𝒜\nu _r_\gamma e_k_\gamma e_l_\gamma e_m_\gamma (𝒜\sigma ),𝒜\sigma _\gamma .$$
We can keep only the subscripts $`lm`$ because of the skew-symmetry of $`(\omega \stackrel{~}{\omega })`$. Besides if we suppose that $`k=l`$, we have terms like $`𝒜\nu _r_\gamma e_m_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$ which belong to $`\text{i}`$. So we can sum over $`k,l,m`$ distinct subscripts without any loss of generality. On the other hand
$$(\overline{\omega }_{lm}\stackrel{~}{\omega }_{lm})(e_k)=\frac{1}{2}\left(g(\stackrel{~}{T}(e_k,e_l),e_m)+g(\stackrel{~}{T}(e_k,e_m),e_l)+g(\stackrel{~}{T}(e_l,e_m),e_k)\right)$$
where the two last terms of the right-hand side member are symmetric with respect to $`(l,k)`$, so they vanish when we sum over $`k`$ and $`l`$ distinct. Consequently
$$\begin{array}{ccc}(\overline{\omega }_{lm}\stackrel{~}{\omega }_{lm})(e_k)ϵ_k_\beta ϵ_l_\beta ϵ_m\hfill & =\hfill & \frac{1}{2}b(𝒜^1(\overline{D}_{e_k}𝒜)ϵ_l𝒜^1(\overline{D}_{e_l}𝒜)ϵ_k,ϵ_m)ϵ_k_\beta ϵ_l_\beta ϵ_m\hfill \\ & =\hfill & b(𝒜^1(\overline{D}_{e_k}𝒜)ϵ_l,ϵ_m)ϵ_k_\beta ϵ_l_\beta ϵ_m,\hfill \end{array}$$
but
$$\begin{array}{ccc}b(𝒜^1(\overline{D}_{e_k}𝒜)ϵ_l,ϵ_m)\hfill & =\hfill & b(𝒜^1(\overline{D}_{e_k}(𝒜ϵ_l)𝒜\overline{D}_{e_k}ϵ_l),ϵ_m)\hfill \\ & \hfill & b(\overline{D}_{ϵ_k}(Bϵ_l)B(\overline{D}_{ϵ_k}ϵ_l),ϵ_m)\hfill \\ & =\hfill & b((\overline{D}_{ϵ_k}B)ϵ_l,ϵ_m),\hfill \end{array}$$
expression which is symmetric with respect to $`(l,m)`$, since $`\overline{D}B`$ is a symmetric endomorphism. Consequently
$$\underset{\text{k,l,m}\text{ distinct}}{}(\overline{\omega }_{lm}\stackrel{~}{\omega }_{lm})(e_k)ϵ_k_\beta ϵ_l_\beta ϵ_m0,$$
when $`r+\mathrm{}`$.
Estimate of $`S_2`$
$`S_2`$ $`=`$ $`{\displaystyle \underset{k,l=2}{\overset{n}{}}}(\overline{\omega }_{1l}\stackrel{~}{\omega }_{1l})(e_k)𝒜\nu _r_\gamma e_k_\gamma e_1_\gamma e_l_\gamma (𝒜\sigma ),𝒜\sigma _\gamma `$
$`=`$ $`{\displaystyle \underset{k=2}{\overset{n}{}}}(\overline{\omega }_{1k}\stackrel{~}{\omega }_{1k})(e_k)\sigma ,\sigma _\beta `$
$`+{\displaystyle \underset{kl}{}}(\overline{\omega }_{1l}\stackrel{~}{\omega }_{1l})(e_k)e_k_\gamma e_l_\gamma (𝒜\sigma ),𝒜\sigma _\gamma ,`$
but the second sum is in $`\text{i}`$, so it remains
$$\mathrm{}e(S_2)=\left(\underset{k=1}{\overset{n}{}}(\overline{\omega }_{1k}\stackrel{~}{\omega }_{1k})(e_k)\right)V_\sigma .$$
We only have to compute
$`{\displaystyle \underset{k=1}{\overset{n}{}}}(\overline{\omega }_{1k}\stackrel{~}{\omega }_{1k})(e_k)`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}g((\overline{D}_{e_1}𝒜)𝒜^1e_k,e_k(\overline{D}_{e_k}𝒜)𝒜^1e_1,e_k)`$
$`=`$ $`S_2^{}S_2^{\prime \prime }`$
We focus on
$`S_2^{\prime \prime }`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}g((\overline{D}_{e_k}𝒜)𝒜^1e_1,e_k)`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}b(𝒜^1(\overline{D}_{e_k}𝒜)ϵ_1,ϵ_k)`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}b(𝒜^1\overline{D}_{e_k}(𝒜ϵ_1)\overline{D}_{e_k}ϵ_1,ϵ_k)`$
$``$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}b((\overline{D}_{e_k}B)ϵ_1,ϵ_k)`$
$``$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}b(ϵ_1,(\overline{D}_{ϵ_k}B)ϵ_k)`$
$`=`$ $`\text{div}_bB(\nu _r).`$
As regards the first term $`S_2^{}`$, we decompose the gauge endomorphism $`𝒜`$ as follows:$`𝒜ϵ_i=_{k=0}^n𝒜_i^kϵ_k`$. We remind that $`𝒜ϵ_0=ϵ_0`$, $`𝒜(TM)TM`$ and so we have $`𝒜_0^k=𝒜_k^0=0,k1`$.
$`S_2^{}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}g((\overline{D}_{e_1}𝒜)𝒜^1e_k,e_k)`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}b((𝒜^1\overline{D}_{e_1}𝒜)ϵ_k,ϵ_k)`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}b(𝒜^1\overline{D}_{e_1}(𝒜ϵ_k)\overline{D}_{e_1}ϵ_k,ϵ_k)`$
$``$ $`{\displaystyle \underset{k,l=1}{\overset{n}{}}}(e_1𝒜_k^l)\left\{b(ϵ_l,ϵ_k)b(Bϵ_l,ϵ_k)\right\}{\displaystyle \underset{k=1}{\overset{n}{}}}b(\overline{D}_{e_1}ϵ_k,ϵ_k)`$
$`+{\displaystyle \underset{k,l=1}{\overset{n}{}}}𝒜_k^l\left\{b(\overline{D}_{e_1}ϵ_l,ϵ_k)b(B\overline{D}_{e_1}ϵ_l,ϵ_k)\right\}`$
$``$ $`(ϵ_1\text{tr}_bB){\displaystyle \underset{k=1}{\overset{n}{}}}b(\overline{D}_{e_1}ϵ_k,ϵ_k)+{\displaystyle \underset{k,l=1}{\overset{n}{}}}(\delta _l^k+B_k^l)\left\{b(\overline{D}_{e_1}ϵ_l,ϵ_k)b(B\overline{D}_{e_1}ϵ_l,ϵ_k)\right\}`$
$``$ $`(ϵ_1\text{tr}_bB){\displaystyle \underset{k=1}{\overset{n}{}}}\left\{b(B\overline{D}_{e_1}ϵ_k,ϵ_k)b(\overline{D}_{e_1}ϵ_k,Bϵ_k)\right\}`$
$`=`$ $`\text{d}(\text{tr}_bB)(\nu _r),`$
that entails
$$\mathrm{}e(S_2)V_\sigma (\text{d}(\text{tr}_bB)+\text{div}_bB)(\nu _r).$$
We can conclude, taking $`B=\frac{1}{2}e`$, that the real part of our boundary integrand is nothing but
$$\frac{1}{4}\left(V_\sigma (\delta _be+\text{d}tr_be)i_{^bV_\sigma }e+(tr_be)\text{d}V_\sigma 2i_{\alpha _\sigma ^{\mathrm{}}}k+2(tr_bk)\alpha _\sigma \right)(\nu _r),$$
what achieves the proof. $`\mathrm{}`$
### 3.2 Analysis of $`\widehat{𝔇}`$
This section is devoted to the study of the analytical properties of $`\widehat{𝔇}`$. The first paragraph deals with the case where $`M`$ has no boundary whereas the second one deals with the case where $`M`$ has a compact and connected boundary denoted as usual by $`M`$.
#### 3.2.1 $`M`$ without Boundary
Proposition. For every $`\sigma IKS(\mathrm{\Sigma })`$ there exists a unique $`\xi _0H(a)`$ such that $`\xi =f𝒜\sigma +\xi _0\text{Ker}\widehat{𝔇}`$ and
$$(V_\sigma ,\alpha _\sigma )=\underset{r+\mathrm{}}{lim}_{S_r}\widehat{}_{𝒜\nu _r}\xi +𝒜\nu _r\widehat{𝔇}\xi ,\xi _\gamma =4_M\left(\left|\widehat{}\xi \right|_\gamma ^2+\widehat{}\xi ,\xi _\gamma \right)0.$$
Proof. We study in a usual way the analytical properties of $`\widehat{𝔇}`$. Let us consider $`C_0^{\mathrm{}}(\mathrm{\Sigma })=C_0^{\mathrm{}}`$ the space of smooth and compactly supported spinors. We define a sesquilinear form on $`C_0^{\mathrm{}}`$ by
$$a(\phi ,\psi )=_M\widehat{𝔇}\phi ,\widehat{𝔇}\psi _\gamma \text{d}\mu _g,$$
where $`\text{d}\mu _g`$ denotes the standard volum form of the metric g. The form $`a`$ is clearly bounded and non-negative on $`C_0^{\mathrm{}}`$. We define the usual Sobolev space
$$H^1(\mathrm{\Sigma })=\left\{\psi \mathrm{\Sigma }/_M\right|\psi |_\gamma ^2+\left|\psi \right|_\gamma ^2<\mathrm{}\}.$$
Definition. We set $`H(a):=\overline{C_0^{\mathrm{}}}^a.`$
Remark. Weighted Poincaré inequality
$$\omega L_{loc}^1(M,\text{dVol}_g)ess_Minf\omega >0uC_0^1_M\omega |u|^2\text{dVol}_g_M|\widehat{}u|^2\text{dVol}_g.$$
It is easy to see that $`\mathrm{\Gamma }`$, the symmetric part of the connection $`\widehat{}`$ is given by $`\mathrm{\Gamma }_X=\frac{1}{2}\{k(X)e_0\text{i}X\}`$, and so satisfies the conditions (cf. ) in order to have the existence of a weighted Poincaré inequality that is to say $`\mathrm{\Gamma }L_{loc}^n(M)\text{and}lim\; sup_{x0}|x\mathrm{\Gamma }_x|<\frac{n1}{2}`$. Such a weighted Poincaré inequality insures the continuity of the embedding of (cf. section 3.2.2 for a proof of this fact when $`M`$ has a compact boundary).$`\mathrm{}`$
We notice that for $`r`$ great enough and for $`\sigma IKS(\mathrm{\Sigma })`$
$`\widehat{}_X(𝒜\sigma )`$ $`=`$ $`_X(𝒜\sigma )+{\displaystyle \frac{\text{i}}{2}}X_\gamma (𝒜\sigma )`$
$`=`$ $`(\overline{}_X\stackrel{~}{}_X)(𝒜\sigma ){\displaystyle \frac{1}{2}}(k(X)_\gamma e_0+\text{i}(𝒜Id)X)_\gamma (𝒜\sigma ).`$
But the relations
$$\{\begin{array}{ccc}\hfill \stackrel{~}{T}(X,Y)& =\hfill & ((\overline{D}_X𝒜)𝒜^1Y(\overline{D}_Y𝒜)𝒜^1X)\hfill \\ \hfill 2g(\stackrel{~}{}_XY\overline{}_XY,Z)& =\hfill & g(\stackrel{~}{T}(X,Y),Z)g(\stackrel{~}{T}(X,Z),Y)g(\stackrel{~}{T}(Y,Z),X)\hfill \end{array}$$
tell us that $`|(\overline{\omega }_{ij}\stackrel{~}{\omega _{ij}})(e_k)|C|𝒜^1||\overline{D}𝒜|`$. We get an estimate
$$|\widehat{𝔇}(𝒜\sigma )|C|𝒜|(|\overline{D}𝒜|+|𝒜Id|+|k|)|\sigma |L^2(M,\text{d}\mu _g),$$
which infers that $`\widehat{𝔇}(f𝒜\sigma )L^2(M,\text{d}\mu _g)`$. We now consider the linear form l on $`H(a)`$ defined by
$$l(\psi )=_M\widehat{𝔇}(f𝒜\sigma ),\widehat{𝔇}\psi _\gamma \text{d}\mu _g.$$
Thanks to our estimate above we get $`|l(\psi )|^2\widehat{𝔇}(f𝒜\sigma )_{L^2}^2a(\psi ,\psi )`$, that gives the continuity of l in $`H(a)`$. We can claim, thanks to Lax-Milgram theorem, that there exists a unique $`\xi _0H(a)`$ such that $`l=a(\xi _0,)`$. In other words
$$_M(\widehat{𝔇})^{}\widehat{𝔇}(f𝒜\sigma +\xi _0),\psi _\gamma =0.$$
Since $`\widehat{𝔇}^{}=\widehat{𝔇}+\text{i}n`$, we have in the distributional sense $`(\widehat{𝔇}+\text{i}n)\widehat{𝔇}\xi =0`$, where we have set $`\xi =f𝒜\sigma +\xi _0`$. By an elliptic regularity argument, $`\widehat{𝔇}\xi `$ is in fact smooth and $`(\widehat{𝔇})^k\xi `$ are $`L^2`$, for every $`k`$. It follows
$`{\displaystyle _M}(\widehat{𝔇})^2\xi ,(\widehat{𝔇})^2\xi _\gamma `$ $`=`$ $`{\displaystyle _M}(\widehat{𝔇}+\text{i}n)(\widehat{𝔇})^2\xi ,\widehat{𝔇}\xi _\gamma `$
$`=`$ $`{\displaystyle _M}\widehat{𝔇}(\widehat{𝔇}+\text{i}n)\widehat{𝔇}\xi ,\widehat{𝔇}\xi _\gamma `$
$`=`$ $`0,`$
that implies $`(\widehat{𝔇})^2\xi =0`$, but we already know that $`(\widehat{𝔇}+\text{i}n)\widehat{𝔇}\xi =0`$, and thereby $`\widehat{𝔇}\xi =0`$. We now apply our integration formula to $`\xi `$
$`(V_\sigma ,\alpha _\sigma )`$ $`=`$ $`\underset{r+\mathrm{}}{lim}{\displaystyle _{S_r}}\widehat{}_{𝒜\nu _r}\xi +𝒜\nu _r\widehat{𝔇}\xi ,\xi _\gamma `$
$`=`$ $`4{\displaystyle _M}\left(\left|\widehat{}\xi \right|_\gamma ^2+\widehat{}\xi ,\xi _\gamma \right)4{\displaystyle _M}\left|\widehat{𝔇}\xi \right|_\gamma ^2`$
$`=`$ $`4{\displaystyle _M}\left(\left|\widehat{}\xi \right|_\gamma ^2+\widehat{}\xi ,\xi _\gamma \right)0,`$
and the proof is complete.$`\mathrm{}`$
#### 3.2.2 $`M`$ with Boundary
We will consider, in this section, a Riemannian slice $`M`$ that has a non empty inner boundary $`M`$. $`\stackrel{˘}{g},\stackrel{}{},\stackrel{˘}{k}`$ will denote respectively the induced metric, the connection and the second fundamental form which is defined by
$`\overline{}_XY`$ $`=`$ $`\underset{X}{\overset{}{}}Y\stackrel{˘}{k}(X,Y)\nu `$
$`\overline{}_X\psi `$ $`=`$ $`\underset{X}{\overset{}{}}\psi {\displaystyle \frac{1}{2}}\stackrel{˘}{k}(X)\nu \psi ,`$
where $`\nu `$ is the normal to $`M`$ pointing toward infinity (that is to say pointing inside), and $``$ still denotes the Clifford action with respect to the metric $`\gamma `$. Consequently our integration formula has another boundary term
$$_{M_r}\left|\widehat{𝔇}\psi \right|_\gamma ^2=_{M_r}\left(\left|\widehat{}\psi \right|_\gamma ^2+\widehat{}\psi ,\psi _\gamma \right)_{S_r}\widehat{}_{𝒜\nu _r}\psi +𝒜\nu _r\widehat{𝔇}\psi ,\psi _\gamma +_M\widehat{}_\nu \psi +\nu \widehat{𝔇}\psi ,\psi _\gamma .$$
But if $`\psi `$ is a compactly supported smooth spinor field then, making $`r\mathrm{}`$ one finds
$$_M\left|\widehat{𝔇}\psi \right|_\gamma ^2=_M\left(\left|\widehat{}\psi \right|_\gamma ^2+\widehat{}\psi ,\psi _\gamma \right)+_M\widehat{}_\nu \psi +\nu \widehat{𝔇}\psi ,\psi _\gamma .$$
We then have to estimate the boundary integrand $`\widehat{}_\nu \psi +\nu \widehat{𝔇}\psi ,\psi `$.
Lemma. If $`(\nu =e_1,e_2,\mathrm{},e_n)`$ is a local orthonormal frame of $`TM_{|M}`$ then
$$\widehat{}_\nu \psi +\nu \widehat{𝔇}\psi =\nu \underset{k=2}{\overset{n}{}}\widehat{}_{e_k}\psi .$$
Proof. Just remark that $`\widehat{}_\nu \psi =e_1e_1\widehat{}_{e_1}\psi `$ $`\mathrm{}`$
Lemma. Keeping our orthonormal frame $`(\nu =e_1,e_2,\mathrm{},e_n),`$ we have
$$\widehat{}_\nu \psi +\nu \widehat{𝔇}\psi =\underset{k=2}{\overset{n}{}}\nu e_k\underset{e_k}{\overset{}{}}\psi +\frac{1}{2}\{\text{tr}\stackrel{˘}{k}(n1)\text{i}\nu +(\text{tr}k)\nu e_0k(\nu )e_0\}\psi .$$
Proof. Using the formula above, we then express $`\widehat{}`$ in term of the $`(n1)`$-dimensional connection and second form, and the $`n`$-dimensional second form.$`\mathrm{}`$
Let us define $`FEnd(\mathrm{\Sigma }_{|M})`$ by $`F(\psi )=\text{i}\nu \psi `$. We sum up some basic properties of $`F`$ in the following
Proposition. The endomorphism F is symmetric, isometric with respect to $`,`$, commutes to the action of $`\nu `$ and anticommutes to each $`e_k,(k1).`$
Lemma. If $`F(\psi )=\psi `$ then
$$\widehat{}_\nu \psi +\nu \widehat{𝔇}\psi ,\psi _{|M}=\frac{1}{2}e_0\left((\text{tr}\stackrel{˘}{k}+(n1))e_0+k(\nu )\right)\psi ,\psi .$$
Proof. Using the proposition above we know that $`\nu e_k(k1)`$ anticommutes with $`F`$ and the formula follows since $`F`$ respects $`,`$.$`\mathrm{}`$
Assumption. Let us suppose that the 4-vector $`\stackrel{}{k}:=(\text{tr}\stackrel{˘}{k}+(n1))e_0+k(\nu )`$ is causal and positively oriented, that is to say $`\gamma (\stackrel{}{k},\stackrel{}{k})0`$ and $`\text{tr}\stackrel{˘}{k}(n1).`$
This assumption (which is exactly the same as for $`\widehat{}`$) guarantees the non-negativity of the boundary integrand term $`\widehat{}_\nu \psi +\nu \widehat{𝔇}\psi ,\psi _{|M}=\frac{1}{2}e_0\stackrel{}{k}\psi ,\psi ,`$ whenever the boundary condition $`F(\psi )=\psi `$ is satisfied. Although this assumption is vectorial, it clearly extends the one given in .
Let us define $`H_{}(a)=\left\{\psi H(a)/F(\psi )=\psi \right\}`$ where $`H(a)`$ has been defined in section 3.6.1. Still taking $`\psi `$ a compactly supported smooth spinor field in $`H_{}(a)`$, we have
$$a(\psi ,\psi )=_M\left(\left|\widehat{}\psi \right|_\gamma ^2+\widehat{}\psi ,\psi _\gamma \right)+\frac{1}{2}_Me_0\stackrel{}{k}\psi ,\psi ,$$
whose each single term is non-negative tanks to our assumption.
Lemma. $`H_{}(a)`$ continuously embeds in $`H_{loc}^1`$ and furthermore
$$\left(\text{}\right)(\text{}\text{and}\mathrm{\Omega }M|\mathrm{\Omega }|<\mathrm{}\text{}),$$
with a weighted Poincaré inequality.
Proof. Let $`\left(\psi _k\right)_k(C_0^{\mathrm{}})^{}`$ a Cauchy sequence with respect to the form $`a`$ whose elements satisfy the boundary condition $`F(\psi _k)=\psi _k`$. Then we have
$$_M\left|\widehat{𝔇}\psi _k\right|_\gamma ^2=\left(\left|\widehat{}\psi _k\right|_\gamma ^2+\widehat{}\psi _k,\psi _k_\gamma \right)+\frac{1}{2}_Me_0\stackrel{}{k}\psi _k,\psi _k,$$
and thus thanks to the weighted Poincaré inequality
$$\mathrm{\Omega }M|\mathrm{\Omega }|<\mathrm{}\text{}\text{and}\text{}.$$
Now let us take a $`\phi C_0^1`$ such that $`\text{Supp}\phi K(MM)`$ ($`K`$ compact without boundary) and then
$$\text{},$$
and therefore $`\rho =\widehat{}\psi `$ in the distributional sense.$`\mathrm{}`$
We consider the linear form $`l`$ on $`H_{}(a)`$ defined by
$$l(\psi )=_M\widehat{𝔇}(f𝒜\sigma ),\widehat{𝔇}\psi _\gamma \text{d}\mu _g.$$
It still is a continuous linear form on the Hilbert space $`H_{}(a)`$ (it is complete since the condition $`F(\psi )=\psi `$ is closed) and applying again Lax-Milgram theorem we get the existence of a unique $`\xi _0H(a)`$ such that $`l=a(\xi _0,)`$. In other words
$$\psi H_{}(a)_M\chi ,\widehat{𝔇}\psi =0,$$
where we have set $`\xi =f𝒜\sigma +\xi _0`$ and $`\chi =\widehat{𝔇}\xi `$.
Remark. We have for any compactly supported smooth spinor fields $`\phi _k,k=1,2`$, the integration by parts formula
$$_M\phi _1,\widehat{𝔇}\phi _2=_M\widehat{𝔇}^{}\phi _1,\phi _2+_M\nu \phi _1,\phi _2.$$
$`\mathrm{}`$
For any $`\psi C_0^1`$ we have
$$_M\chi ,\widehat{𝔇}\psi =0=_M\widehat{𝔇}^{}\chi ,\psi +_M\nu \chi ,\psi .$$
But remembering that $`C_0^{\mathrm{}}(MM)`$ the space of smooth spinor fields compactly supported in $`MM`$ is dense in $`L^2(M)`$ then we obtain that $`\widehat{𝔇}^{}\chi =0`$ and $`\chi H_+(a)=\left\{\psi H(a)/F(\psi )=+\psi \right\}`$. By ellipticity $`\chi `$ is smooth and $`\widehat{𝔇}^k\chi L^2(M)`$ for every $`k`$. Finally we notice that
$`{\displaystyle _M}\left|\widehat{𝔇}\chi \right|^2`$ $`=`$ $`{\displaystyle _M}\widehat{𝔇}^{}\widehat{𝔇}\chi ,\chi +{\displaystyle _M}\nu \widehat{𝔇}\chi ,\chi `$
$`=`$ $`0+{\displaystyle _M}\text{i}n\nu \chi ,\chi `$
$`=`$ $`n{\displaystyle _M}\left|\chi \right|^2,`$
and therefore $`\widehat{𝔇}\chi =0`$ which implies that $`\chi =0`$. We can conclude with the
Proposition. For every $`\sigma IKS(\mathrm{\Sigma })`$ there exists a unique $`\xi _0H_{}(a)`$ such that
$`\xi =f𝒜\sigma +\xi _0\text{Ker}\widehat{𝔇}H_{}(a)`$ and
$`(V_\sigma ,\alpha _\sigma )`$ $`=`$ $`\underset{r+\mathrm{}}{lim}{\displaystyle _{S_r}}\widehat{}_{𝒜\nu _r}\xi +𝒜\nu _r\widehat{𝔇}\xi ,\xi _\gamma `$
$`=`$ $`4{\displaystyle _M}\left(\left|\widehat{}\xi \right|_\gamma ^2+\widehat{}\xi ,\xi _\gamma \right)+2{\displaystyle _M}e_0\stackrel{}{k}\psi ,\psi 0.`$
### 3.3 Imaginary Killing Spinors
In general (that is to say whatever the dimension), the space of imaginary Killing spinors of AdS<sup>n,1</sup> is a finite dimensional complex vector space which trivializes the spinor bundle of AdS<sup>n,1</sup> (cf. for instance). Consequently there exists an integer $`d`$ depending upon the dimension $`n+1`$, such that $`IKS(\mathrm{\Sigma })^d`$ where $`\mathrm{\Sigma }`$ still is the spinor bundle of AdS<sup>n,1</sup> restricted to a standard hyperbolic slice. Thereby there exists a quadratic Hermitian application
$$𝒦:\text{}$$
which is difficult to explicite except when the dimension of the slice is $`n=3`$ (because of exceptional isomorphisms for Lie groups). In order to make out the meaning of the non-negativity of the energy-momentum, we will consider in this section the particular case of dimension $`n=3`$.
The aim of this section is to solve explicitely the Killing equation of section 2.2. As a matter of fact, representation theory provides us good candidates for the imaginary Killing spinors. Thanks to Schur’s lemma, we have an isomorphism
$$\begin{array}{ccc}^2& & \text{Hom}^{\text{SU}(2)}(^2,^2^2)\\ \left(\genfrac{}{}{0pt}{}{z_1}{z_2}\right)& & \left(\genfrac{}{}{0pt}{}{z_1\text{I}_2}{z_2\text{I}_2}\right)\end{array},$$
We are now considering two families of spinors which are derived from representation theory. To this end, we will denote $`wz^2\text{Hom}^{\text{SU}(2)}(^2,^2^2)`$ thanks to the isomorphism above.
Definition. Let $`wz^2^2`$ and set $`\sigma _{wz}^1[\stackrel{~}{g}]=\{\stackrel{~}{g},z(\stackrel{~}{g}^1w)\},\sigma _{wz}^{}[\stackrel{~}{g}]=\{\stackrel{~}{g},z(\stackrel{~}{g}^{}w)\}.`$
Let us consider a spinor field $`\tau \mathrm{\Gamma }(\mathrm{\Sigma })`$ and a vector field $`X\mathrm{\Gamma }(𝕋)`$ tangent to $`^3`$. We can write $`\tau [\stackrel{~}{g}]=\{\stackrel{~}{g},v(\stackrel{~}{g})\}`$ and $`X[g]=\{g,\zeta (g)\}`$, where $`v:^3^4`$ and $`\zeta :^3𝔊`$ are respectively $`\text{SO}(3)`$ and $`\text{SU}(2)`$-équivariant functions. We can now differentiate $`\tau `$ in the direction of $`X`$ and write down
$$(D_X\tau )[\stackrel{~}{g}]=\{\stackrel{~}{g},v_{}(X)_{[\stackrel{~}{g}]}+\stackrel{~}{\rho }_{}s^{}\theta (\zeta )_{[\stackrel{~}{g}]}v[\stackrel{~}{g}]\},$$
where $`\theta `$ is the connection 1-form of the bundle of $`\text{SL}(2,)`$-frames, restricted to $`^3`$. If one remembers that $`\theta `$ is only the projection on the first factor in the decomposition $`𝔰𝔩_2()𝔰𝔲(2)𝔊`$, we can conclude that $`\stackrel{~}{\rho }_{}s^{}\theta (\zeta )_{[\stackrel{~}{g}]}v[\stackrel{~}{g}]`$ vanishes. Besides we will apply this formula to spinors in $`\left\{\sigma _{wz}^1,\sigma _{uz}^{},w,u^2\right\}`$ so that we can only derive at the point $`\stackrel{~}{g}=1`$ unity in $`\text{SL}(2,)`$ since we have the
Proposition. The set $`\left\{\sigma _{wz}^1,\sigma _{uz}^{},w,u^2\right\}`$ is stable under the $`\text{SL}(2,)`$ action. More precisely for every $`\stackrel{~}{e}\text{SL}(2,)`$ we have $`\stackrel{~}{e}\sigma _{wz}^1=\sigma _{\stackrel{~}{e}wz}^1`$ and $`\stackrel{~}{e}\sigma _{uz}^{}=\sigma _{(\stackrel{~}{e}^{})^1uz}^{}.`$
We obtain
$$\{\begin{array}{ccc}(D_X\sigma _{wz}^1)[1]\hfill & =\hfill & \{1,z(\zeta w)\}\hfill \\ (D_X\sigma _{uz}^{})[1]\hfill & =\hfill & \{1,z(\zeta u)\}\hfill \end{array},$$
where $`\zeta =\zeta (1)`$. We also compute the Clifford action of X on $`\sigma _{wz}^1,\sigma _{uz}^{}`$ at the point 1:
$$\{\begin{array}{ccc}X\sigma _{wz}^1[1]\hfill & =\hfill & \{1,\mathrm{\Theta }(\zeta )z(w)\}\hfill \\ X\sigma _{uz}^{}[1]\hfill & =\hfill & \{1,\mathrm{\Theta }(\zeta )z(u)\}\hfill \end{array}.$$
We must precise $`\mathrm{\Theta }_{|𝔊}:\zeta \left(\begin{array}{cc}0& 2\zeta \\ 2\zeta & 0\end{array}\right)`$, and if we introduce the sections $`\sigma _{w\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^1`$ and $`\sigma _{w\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^{}`$, for any $`w^2`$, we have on one hand
$$\{\begin{array}{ccccc}\frac{i}{2}X\sigma _{w\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^1[1]\hfill & =\hfill & i\{1,i\zeta w\zeta w\}\hfill & =\hfill & \{1,\zeta wi\zeta w\}\hfill \\ \frac{i}{2}X\sigma _{u\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^{}[1]\hfill & =\hfill & i\{1,i\zeta u\zeta u\}\hfill & =\hfill & \{1,\zeta ui\zeta u\}\hfill \end{array},$$
and on the other hand
$$\{\begin{array}{ccc}\left(D_X\sigma _{w\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^1\right)[1]& =\hfill & \{1,\zeta wi\zeta w\}\hfill \\ \left(D_X\sigma _{u\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^{}\right)[1]& =\hfill & \{1,\zeta ui\zeta u\}\hfill \end{array}.$$
Since $`\left\{\sigma _{w\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^1+\sigma _{u\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^{}/w,u^2\right\}`$ is a 4-dimensional complex vector space, we obviously obtain the
Proposition. The space of imaginary Killing spinors denoted by $`IKS(\mathrm{\Sigma })`$ is generated by
$$\left\{\sigma _{w\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^1,\sigma _{u\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^{},w,u^2\right\}.$$
Let $`\sigma `$ an imaginary Killing spinor and set $`V_\sigma :=<\sigma ,\sigma >`$ which is a function on $`^3`$, and if $`e_0`$ denotes a unit normal of $`^3`$ in AdS, we set $`\alpha _\sigma (Y):=Ye_0\sigma ,\sigma `$ which is a real 1-form on $`^3`$. The goal of the two next paragraphs is to define some $`\text{SL}(2,)`$-equivariant application
$$\begin{array}{cccc}𝒦:\hfill & \hfill IKS(\mathrm{\Sigma })^2^2& \hfill & \left(𝔐𝔰𝔩_2()\right)^{}\hfill \\ & \hfill wu& \hfill & 𝒦_{wu}:=(V_{wu}\alpha _{wu}).\hfill \end{array}$$
The functions $`V_\sigma `$
We compute the functions $`V_\sigma `$ which are by definition
$$\begin{array}{ccc}V_\sigma [\stackrel{~}{g}]\hfill & =\hfill & \left|\sigma _{w\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^1[\stackrel{~}{g}]\right|_^4^2+\left|\sigma _{u\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^{}[\stackrel{~}{g}]\right|_^4^2+2\mathrm{}e\left(\sigma _{w\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^1[\stackrel{~}{g}],\sigma _{u\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^{}[\stackrel{~}{g}]_^4\right)\hfill \\ & =\hfill & 2\left|\stackrel{~}{g}^1w\right|_^2^2+2\left|\stackrel{~}{g}^{}u\right|_^2^2\hfill \end{array}$$
Remark. $`\sigma _{w\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^1`$ and $`\sigma _{u\left(\genfrac{}{}{0pt}{}{1}{i}\right)}^{}`$ are orthogonal spinors for every $`u,w^2`$.$`\mathrm{}`$
If $`\stackrel{~}{g}\text{SL}(2,)`$, the corresponding base point is $`\stackrel{~}{g}\stackrel{~}{g}^{}^3𝔐^{3,1}`$ whose coordinates are given by $`(x_k)_{k=0}^3=\mathrm{\Lambda }^1(\stackrel{~}{g}\stackrel{~}{g}^{})`$.
Proposition. $`V_\sigma `$ is a causal element of $`N_b.`$
Proof. Let $`U=\left(\genfrac{}{}{0pt}{}{u_1}{\overline{w_2}}\right)^2,V=\left(\genfrac{}{}{0pt}{}{u_2}{\overline{w_1}}\right)^2`$. We notice that
$$V_\sigma [\stackrel{~}{g}]=x_0(|U|^2+|V|^2)+x_1(|U|^2|V|^2)+2x_2\mathrm{}e(<U,V>)2x_3\mathrm{}m(<U,V>),$$
so that the norm of $`V_\sigma `$ is $`|V_\sigma [\stackrel{~}{g}]|^2=4(|<U,V>|^2|U|^2|V|^2)0`$, thanks to the Cauchy-Schwarz inequality for the standard Hermitian form on $`^2`$.$`\mathrm{}`$
More conceptually we see that $`V_\sigma [\stackrel{~}{g}]=2(w^{}\widehat{W}w+u^{}Wu)`$, where we have set $`W:=\stackrel{~}{g}\stackrel{~}{g}^{}^3𝔐`$. Thereby we can define by extension an application
$$\begin{array}{ccc}^2^2\hfill & \hfill & 𝔐^{}\hfill \\ wu\hfill & \hfill & \left(V_{wu}:W2(w^{}\widehat{W}w+u^{}Wu)\right)\hfill \end{array}.$$
The 1-forms $`\alpha _\sigma `$
The positively oriented unit normal of $`^3`$ in AdS is given by $`e_0[\stackrel{~}{g}]=\{\stackrel{~}{g},\frac{1}{2}\mu (\stackrel{~}{g})I_2\}`$ and for any $`\xi 𝔊`$ satisfying $`\text{det}\xi =1`$ we set $`X^\xi [\stackrel{~}{g}]=\{\stackrel{~}{g},\frac{1}{2}\mu (\stackrel{~}{g})\xi \}`$. Just remember that $`\alpha _\sigma (X^\xi )_{[\stackrel{~}{g}]}:=X^\xi e_0\sigma ,\sigma _{[\stackrel{~}{g}]}`$. As we suppose that $`\sigma IKS(\mathrm{\Sigma })`$, we can easily compute the first derivative of $`\alpha _\sigma `$
$$D_{X^\eta }\alpha _\sigma (X^\xi )_{[\stackrel{~}{g}]}=\frac{\text{i}}{2}(X^\eta X^\xi X^\xi X^\eta )e_0\sigma ,\sigma _{[\stackrel{~}{g}]},$$
which is a real skew symmetric 2-form and hence $`\alpha _\sigma `$ is a Killing form on $`^3`$ . From now on we set $`\alpha _\sigma =(\alpha _\sigma )_1`$ and $`D\alpha _\sigma =(D\alpha _\sigma )_1,`$ that we will write as function of $`wu`$. After some computations we find
$$\{\begin{array}{ccc}\alpha _\sigma (\xi )\hfill & =\hfill & 2(w^{}\xi u+u^{}\xi w)\hfill \\ D\alpha _\sigma (\eta ,\xi )\hfill & =\hfill & (w^{}(\xi \eta \eta \xi )uu^{}(\xi \eta \eta \xi )w)\hfill \end{array}.$$
We have to notice that $`\xi \eta \eta \xi \text{i}𝔊`$ so that $`D\alpha _\sigma `$ is naturally a linear form on $`\text{i}𝔊`$. As a consequence we define, thanks to the Killing 1-form $`\alpha _\sigma `$, the following application
$$\begin{array}{ccc}^2^2& \hfill & 𝔰𝔩_2()^{}\hfill \\ wu& \hfill & \left(\alpha _{wu}:\xi 2(w^{}\xi u+u^{}\xi ^{}w)\right),\hfill \end{array}$$
where <sup>∗R</sup> stands for the duality with respect to the reals. We then define
$$𝒦_{wu}=V_{wu}\alpha _{wu}$$
and conclude with the
Proposition. The application $`𝒦`$ is $`\text{SL}(2,)`$-equivariant. More precisely, for every $`\stackrel{~}{e}\text{SL}(2,)`$
$$𝒦_{\stackrel{~}{e}(wu)}=\left(V_{wu}\mu (\stackrel{~}{e}^1)\right)\left(\alpha _{wu}\text{Ad}(\stackrel{~}{e}^{})\right).$$
Proof. We must compute for every $`W𝔐`$ and $`\xi 𝔰𝔩_2()`$
$$\begin{array}{ccc}𝒦_{\stackrel{~}{e}(wu)}(W,\xi )\hfill & =\hfill & 𝒦_{\stackrel{~}{e}w(\stackrel{~}{e}^{})^1u}(W,\xi )\hfill \\ & =\hfill & 2(w^{}\stackrel{~}{e}^{}\widehat{W}\stackrel{~}{e}w+u^{}\stackrel{~}{e}^1W(\stackrel{~}{e}^{})^1u+w^{}\stackrel{~}{e}^{}\xi (\stackrel{~}{e}^{})^1u+u^{}\stackrel{~}{e}^1\xi ^{}\stackrel{~}{e}w)\hfill \\ & =\hfill & V_{wu}\mu (\stackrel{~}{e}^1)(W)\alpha _{wu}\text{Ad}(\stackrel{~}{e}^{})(\xi ).\hfill \end{array}$$
$`\mathrm{}`$
Remark. The norm of imaginary Killing spinors
Classical considerations on Lie algebras show that $`𝔰𝔬(3,2)`$ endowed with its Killing form, is isometric to $`(𝔐,\text{det})(𝔰𝔩_2(),\mathrm{}e(\text{det}))`$ which is a 10-dimensional real vector space of signature (6,4). The norm of $`𝒦(wu)`$ with respect to the Killing form is, up to a multiplicative and positive constant $`|𝒦(wu)|^2=|<U,V>|^2|U|^2|V|^2+\mathrm{}e(\chi ^2)`$, where we have set $`\chi =\overline{u_1}w_1+\overline{u_2}w_2`$. Besides, if $`V_{wu}`$ is isotropic in $`𝔐`$ then $`\alpha _{wu}`$ and $`𝒦(wu)`$ are also isotropic respectively in $`𝔰𝔩_2()^{}`$ and $`\left(𝔐𝔰𝔩_2()\right)^{}`$. Indeed the equality case in the Cauchy-Schwarz inequality occurs if and only if U and V satisfy $`\text{det}_^2(U,V)=\overline{\chi }=0.`$$`\mathrm{}`$
### 3.4 End of the Proof
Whatever the dimension is, we obtain a Hermitian application
$$Q:\text{},$$
which has to be non-negative in vertue of the non-negativity results of sections 3.2.1 and 3.2.2. This completes the proof of the positivity theorem stated in section 1.3.
In dimension $`n=3`$, we can be more specific giving the explicite formula of $`Q`$ in terms of the components of the energy-momentum $``$. More precisely, on one hand we have found a quadratic application
$$\begin{array}{cccc}𝒦:\hfill & \hfill IKS(\mathrm{\Sigma })^2^2& \hfill & \left(𝔐𝔰𝔩_2()\right)^{}\text{Ker d}\mathrm{\Phi }_{(b,0)}^{}\hfill \\ & \hfill wu& \hfill & (V_{wu}\alpha _{wu})\hfill \end{array},$$
which is $`\text{SL}(2,)`$-equivariant. On the other hand we know that the energy-momentum functional $``$ can be seen as a real linear form on $`\left(𝔐𝔰𝔩_2()\right)^{}`$ that is to say, as a vector $`=M\mathrm{\Xi }𝔐𝔰𝔩_2()`$. In the following, we will adopt the notations $`\mathrm{\Xi }=N\text{i}R𝔊\text{i}𝔊`$, and $`M=\mathrm{\Lambda }(m_0,m),N=\mathrm{\Lambda }(0,n),R=\mathrm{\Lambda }(0,r)`$, where $`\mathrm{\Lambda }`$ is the isomorphism defined in section 2.4. Now applying the non-negativity result of section 3.2.1 or 3.2.2, we know that (even if our AdS-asymptotically hyperbolic manifold has a compact boundary such that $`\stackrel{}{k}`$ is causal and positively oriented)
$$\sigma IKS(\mathrm{\Sigma })(V_\sigma ,\alpha _\sigma )0.$$
In other words, for each $`wu^4`$, we have $`(𝒦_{wu})0`$. But the complete study of $`IKS(\mathrm{\Sigma })`$ of section 3.3 implies that actually
$`(𝒦_{wu})`$ $`=`$ $`V_{wu}(M)+\alpha _{wu}(\mathrm{\Xi })`$
$`=`$ $`2(w^{}\widehat{M}w+u^{}Mu)+2(w^{}\mathrm{\Xi }u+u^{}\mathrm{\Xi }^{}w),`$
and consequently the application $`wu(𝒦_{wu})`$ is a Hermitian form on $`^2^2`$ whose matrix is
$$Q=2\left(\begin{array}{cc}\widehat{M}& \mathrm{\Xi }\\ \mathrm{\Xi }^{}& M\end{array}\right)=2\left(\begin{array}{cc}\mathrm{\Lambda }(m_0,m)& \mathrm{\Lambda }(0,n)+\text{i}\mathrm{\Lambda }(0,r)\\ \mathrm{\Lambda }(0,n)\text{i}\mathrm{\Lambda }(0,r)& \mathrm{\Lambda }(m_0,m)\end{array}\right).$$
It is easy to conclude since we have the identity
$$wu^4(V_{wu}\alpha _{wu})=Q(wu,wu)0,$$
which ends the proof of the
Positive Energy-Momentum Theorem. Let $`(M^n,g,k)`$ be an AdS-asymptotically hyperbolic spin Riemannian manifold satisfying the decay conditions stated in section 1.2 and the following conditions
(i) $`(f,\alpha ),(\mathrm{\Phi }(g,k)\mathrm{\Phi }(b,0))L^1(M,\text{dVol}_b)`$ for every $`(f,\alpha )N_b𝔎𝔦𝔩𝔩(M,b),`$
(ii) the relative version of the dominant energy condition (cf. section 2.2) holds, that is to say $`(\mathrm{\Phi }(g,k)\mathrm{\Phi }(b,0))`$ is a positively oriented causal (n+1)-vector along $`M`$,
(iii) in the case where M has a compact boundary $`M`$, we assume moreover that $`\stackrel{}{k}`$ is causal and positively oriented along $`M`$.
Then there exists a (hardly explicitable) map $`^{n,1}𝔰𝔬(n,1)\text{Herm}(C^d)`$ which sends, under the assumptions (i-iii), the energy-momentum on a non-negative Hermitian form Q.
Moreover, when n=3, we can explicite Q in terms of the components of the energy-momentum as described above.
The end of this section is devoted to the 3-dimensional case.
As the invariance of the non-negativity of $`Q`$ under asymptotic hyperbolic isometries, was proved in , one can be interested in the description of the orbit of the energy-momentum under the action of $`\text{SL}(2,)`$.
Proposition. If M is timelike, there exists a (non-unique) representative element of the orbit of $`=M\mathrm{\Xi }`$ under the natural action (cf. section 3.3) of $`\text{SL}(2,)`$ on $`𝔐𝔰𝔩_2()`$ which can be written
$$m_0\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)n_1\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\text{i}\left(\begin{array}{cc}r_1& r_2\\ r_2& r_1\end{array}\right),m_0,n_1,r_1,r_2.$$
The positive energy-momentum theorem then reduces to $`m_0\sqrt{(|n_1|+|r_2|)^2+r_1^2}.`$
Proof. Let us suppose that $`M𝔐`$ is timelike. Thus considering the action of $`\text{SL}(2,)`$ on $`𝔐𝔰𝔩_2()`$ (cf. section 3.3), then there exists an element in the orbit of $``$ that can be written $`m_0\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\mathrm{\Xi }^{}`$. Since the isotropy group of $`\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ is $`\text{SU}(2)`$ whose action on $`𝔊`$ is transitive, then there exists an element in the orbit of $``$ that can be written $`m_0\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)n_1\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\text{i}R^{}`$. But the isotropy group of $`\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ is the one parameter group $`\left\{\left(\begin{array}{cc}e^{\text{i}\theta }& 0\\ 0& e^{\text{i}\theta }\end{array}\right),\theta \right\}`$. Finally there exists an element (not unique since the isotropy group of $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ is isomorphic to $`_2`$) in the orbit of $``$ that can be written as announced in the proposition. The corresponding Hermitian matrix is
$$Q=2\left(\begin{array}{cccc}m_0& 0& n_1+\text{i}r_1& \text{i}r_2\\ 0& m_0& \text{i}r_2& n_1\text{i}r_1\\ n_1\text{i}r_1& \text{i}r_2& m_0& 0\\ \text{i}r_2& n_1+\text{i}r_1& 0& m_0\end{array}\right).$$
Since $`Q`$ is non negative we have
$`m_0`$ $``$ $`0`$
$`m_0(m_0^2(n_1^2+r_1^2+r_2^2))`$ $``$ $`0`$
$`(m_0^2(n_1^2+r_1^2+r_2^2))^2`$ $``$ $`4(n_1r_2)^2,`$
which can be summerized with $`m_0\sqrt{(|n_1|+|r_2|)^2+r_1^2}`$.$`\mathrm{}`$
Remark. The $`\left\{t=0\right\}`$ slices of the Kerr AdS metrics are AdS-asymptotically hyperbolic and parametrized by 2 real parameters: the mass and the angular momentum. The proposition above then shows that there exists some energy-momenta that could not be obtained by the action of $`\text{SL}(2,)`$ on a Kerr-AdS solution. As a consequence, an interesting question would be to find some (new) AdS-asymptotically hyperbolic metrics which have an energy-momentum of the form given in the proposition above with non-zero coefficients $`m_0,n_1,r_1,r_2`$, and which satisfy the dominant energy condition or the (stronger) cosmological vacuum constraints.
## 4 Rigidity Theorems
Theorem. Under the assumptions of the positive energy-momentum theorem, $`\text{tr}Q=0`$ implies that $`(M,g,k)`$ is isometrically embeddable in AdS<sup>n,1</sup>.
Proof. The vanishing of $`\text{tr}Q`$ implies from the non-negativity of $`Q`$, that $`Q=0`$. Consequently our spinor bundle $`\mathrm{\Sigma }`$ is trivialized by a basis of $`\gamma `$-imaginary Killing spinors. We denote by $`\xi `$ any $`\gamma `$-imaginary Killing spinors of this basis. We will need the following spinorial Gauss-Codazzi equation.
Proposition. For every $`X,Y\mathrm{\Gamma }(TM)`$ we have
$$R_{X,Y}^\gamma =R_{X,Y}^g\frac{1}{2}(\text{d}^\overline{}k(X,Y)e_0+\frac{1}{2}(k(X)k(Y)k(Y)k(X))),$$
where $``$ denotes the Clifford action with respect to the metric $`\gamma `$.
Proof of the proposition. It is a straightforward computation where we use vector fields $`X,Y`$ satisfying at the point $`\overline{}_XY=\overline{}_YX=0`$.
$`_X_Y`$ $`=`$ $`_X(\overline{}_Y{\displaystyle \frac{1}{2}}k(Y)e_0)`$
$`=`$ $`\overline{}_X\overline{}_Y{\displaystyle \frac{1}{2}}k(X)e_0\overline{}_Y`$
$`{\displaystyle \frac{1}{2}}(_Xk(Y)e_0+k(Y)(_Xe_0)+k(Y)e_0_X)`$
$`=`$ $`\overline{}_X\overline{}_Y{\displaystyle \frac{1}{2}}(k(X)e_0\overline{}_Y+k(Y)e_0\overline{}_Xkk(X,Y)e_0)`$
$`({\displaystyle \frac{1}{2}}\overline{}_Xk(Y)e_0{\displaystyle \frac{1}{4}}k(Y)k(X)),`$
and the curvature formula above follows.$`\mathrm{}`$
Using the fact that $`\xi `$ is a $`\gamma `$-imaginary Killing spinor one gets
$$R_{X,Y}^g\xi \frac{1}{4}\left(XYYX+k(X)k(Y)k(Y)k(X)\right)\xi ,\xi =\frac{1}{2}\text{d}^\overline{}k(X,Y)e_0\xi ,\xi ,$$
where $`R_{X,Y}^g\xi ,\xi `$ and $`\left(XYYX+k(X)k(Y)k(Y)k(X)\right)\xi ,\xi `$ are purely imaginary terms whereas $`\text{d}^\overline{}k(X,Y)e_0\xi ,\xi `$ is real. As a consequence $`\text{d}^\overline{}k(X,Y)e_0\xi ,\xi =0`$ for any $`\xi `$ of our $`\gamma `$-imaginary Killing spinor basis and so $`\text{d}^\overline{}k=0`$. This implies
$$R_{X,Y}^g=\frac{1}{4}(XYYX+k(X)k(Y)k(Y)k(X)),$$
and using the natural isomorphism between $`\text{C}\mathrm{}_0(^{3,1})`$ and $`\mathrm{\Lambda }^2(^{3,1})`$ (cf. proposition 6.2) we get that
$`R^g`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(g\wedge ⃝g+k\wedge ⃝k\right)`$
$`\text{d}^\overline{}k`$ $`=`$ $`0.`$
Let us denote by $`V`$ the function $`<\xi ,\xi >`$, $`\alpha `$ the real 1-form defined by $`\alpha (Y)=Ye_0\xi ,\xi `$. Then the couple $`(V,Y):=(V,\alpha ^{\mathrm{}})`$ is a Killing Initial Data (KID) . If we consider $`(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{k})`$ the universal Riemannian covering of $`(M,g,k)`$, then we can make the Killing development of $`(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{k})`$ with respect to the KID $`(\stackrel{~}{V},\stackrel{~}{Y})`$ which by definition is $`\times \stackrel{~}{M}`$ endowed with the Lorentzian metric $`\stackrel{~}{\gamma }=\left(\stackrel{~}{N}^2+|\stackrel{~}{Y}|^2\right)\text{d}u^2+2\stackrel{~}{Y}^{\mathrm{}}\text{d}u+\stackrel{~}{g}`$. By construction, $`\stackrel{~}{M}`$ is embedded in $`(\times \stackrel{~}{M},\stackrel{~}{\gamma })`$ with induced metric $`\stackrel{~}{g}`$ and second fundamental form $`\stackrel{~}{k}`$. Besides $`\times \stackrel{~}{M}`$ is the universal covering of $`N`$, and $`\stackrel{~}{\gamma }`$ which has sectional curvature -1, is a stationary solution of the vacuum Einstein equations with cosmological constant that is to say $`G^{\stackrel{~}{\gamma }}=\frac{n(n1)}{2}\stackrel{~}{\gamma }`$. But $`(\stackrel{~}{M},\stackrel{~}{g})`$ is complete since $`(M,g)`$ is complete and therefore $`(\times \stackrel{~}{M},\stackrel{~}{\gamma })`$ is geodesically complete. It follows that $`(\times \stackrel{~}{M},\stackrel{~}{\gamma })`$ is AdS<sup>n,1</sup> (in vertue of Proposition 23 \[p.227\] of ). It only remains to show that $`M`$ is simply connected. We know that $`\times \stackrel{~}{M}^{n+1}`$ and thereby using the following compactly supported de Rham cohomology isomorphisms $`\left\{0\right\}=H_{dR,c}^2(\times \stackrel{~}{M})=H_{dR,c}^2(^{n+1})=H_{dR,c}^1(\stackrel{~}{M})`$ (cf. Proposition 4.7 and Corollary 4.7.1 \[p. 39\] of for instance), we obtain that $`\stackrel{~}{M}`$ has only one asymptotic end. This last fact compels the universal covering map $`\stackrel{~}{M}M`$ to be trivial and as a consequence $`(M,g,k)(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{k})`$ is isometrically embedded in AdS$`{}_{}{}^{n,1}(\stackrel{~}{N},\stackrel{~}{\gamma })`$. This completes the proof of the theorem. $`\mathrm{}`$
Then end of this section is devoted to weaken the condition defining the rigidity case in dimension $`n=3`$. Namely we prove the following
Theorem. Let us suppose that $`(M^3,g,k)`$ satisfies the assumptions of the positive energy-momentum theorem and that the matrix Q is degenerate. Then there exists some $`\widehat{}`$-parallel spinor field $`\xi `$ such that $`\widehat{}\xi ,\xi =0`$ and consequently $`(M,g,k)`$ is isometrically embeddable in a stationary pp-wave space-time.
If furthermore the constant function $`(\xi ,\xi )`$ is non-zero then $`(M,g,k)`$ admits a vacuum Killing development which is a solution of the Einstein equations (with the cosmological constant -3).
Remark. A pp-wave space-time is a Lorentzian manifolds such that its stress-energy tensor satisfies $`T_{\mu \nu }=\lambda Z_\mu Z_\nu `$ where $`Z^\mu `$ is an isotropic Killing vector field and $`\lambda `$ a function on the manifold. Some results were also proved by Siklos in and by Leitner in for Lorentzian manifolds admitting a Killing spinor.$`\mathrm{}`$
Proof. The degenerate character of $`Q`$ implies the existence of a non-zero $`wu^4`$ and a unique $`\xi _0`$ such that $`\xi =f𝒜\sigma _{wu}+\xi _0`$ satisfies the conditions
$`\widehat{}\xi `$ $`=`$ $`0`$
$`\widehat{}\xi ,\xi `$ $`=`$ $`0.`$
By the same argument as above we get that $`\text{d}^\overline{}k(X,Y)e_0\xi ,\xi =0`$ (which can also be thought as $`\overline{}_Xk(Y,\alpha )=\overline{}_Yk(X,\alpha )`$). Now since $`\xi `$ is $`\widehat{}`$-parallel we get
$`\mathrm{}e{\displaystyle \underset{k=1}{\overset{3}{}}}e_kR_{X,e_k}^\gamma \xi ,Y\xi `$ $`=`$ $`{\displaystyle \frac{1}{4}}\mathrm{}e{\displaystyle \underset{k=1}{\overset{3}{}}}e_k(Xe_ke_kX)\xi ,Y\xi `$
$`=`$ $`X,YX,Y\mathrm{\Gamma }(TM).`$
On the other hand a direct computation leads to
$`{\displaystyle \underset{k=1}{\overset{3}{}}}e_kR_{X,e_k}^\gamma `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{l,m=1}{\overset{3}{}}}\text{d}^\overline{}k(X,e_l,e_m)e_le_0e_m+{\displaystyle \frac{1}{2}}{\displaystyle \underset{l,m=1}{\overset{3}{}}}R^\gamma (X,e_l,e_l,e_m)e_le_le_m`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{l,m=1}{\overset{3}{}}}\text{d}^\overline{}k(X,e_l,e_m)e_le_0e_m`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{l,m=1}{\overset{3}{}}}\{R^g(X,e_l,e_l,e_m)k(X,e_l)k(e_l,e_m)+k(X,e_m)k(e_l,e_l)\}e_m`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{l,m=1}{\overset{3}{}}}\text{d}^\overline{}k(X,e_l,e_m)e_le_0e_m{\displaystyle \frac{1}{2}}(E(X)2X),`$
where we have set $`E=\text{Ric}^g+2g+(\text{tr}k)kkk`$. It is then clear that
$$\mathrm{}e\underset{k,l=1}{\overset{3}{}}\text{d}^\overline{}k(X,e_l,e_m)e_le_0e_m\psi ,Y\psi =VE(X,Y)$$
In the following computation we will set $`Y=e_s`$. We recall that
$`\mathrm{}e{\displaystyle \underset{l,m=1}{\overset{3}{}}}\text{d}^\overline{}k(X,e_l,e_m)e_le_0e_m\psi ,e_s\psi `$ $`=`$ $`{\displaystyle \underset{l,m=1}{\overset{3}{}}}\text{d}^\overline{}k(X,e_l,e_m)\mathrm{}ee_le_0e_m\psi ,e_s\psi `$
$`=`$ $`{\displaystyle \underset{l=1}{\overset{3}{}}}\text{d}^\overline{}k(X,e_l,e_l)\mathrm{}ee_le_0e_l\psi ,e_s\psi `$
$`+{\displaystyle \underset{lm}{}}\text{d}^\overline{}k(X,e_l,e_m)\mathrm{}ee_le_0e_m\psi ,e_s\psi `$
$`=`$ $`(I+II)(X,e_s),`$
and we will treat $`I`$ and $`II`$ seperately for convenience. The easiest one is
$`I(X,e_s)`$ $`=`$ $`e_se_0\psi ,\psi {\displaystyle \underset{l=1}{\overset{3}{}}}(\overline{}_Xk(e_l,e_l)\overline{}_{e_l}k(X,e_l))`$
$`=`$ $`\left((\delta _gk+\text{dtr}_gk)\alpha \right)(X,e_s).`$
Thereby we can conclude that $`I=(\delta _gk+\text{dtr}_gk)\alpha `$. We compute now $`II(X,e_s)`$.
$`II(X,e_s)`$ $`=`$ $`{\displaystyle \underset{ls}{}}\text{d}^\overline{}k(X,e_l,e_s)\mathrm{}ee_le_0e_s\psi ,e_s\psi `$
$`+{\displaystyle \underset{ms}{}}\text{d}^\overline{}k(X,e_s,e_m)\mathrm{}ee_se_0e_m\psi ,e_s\psi `$
$`+{\displaystyle \underset{lm,ls,ms}{}}\text{d}^\overline{}k(X,e_l,e_m)\mathrm{}ee_le_0e_m\psi ,e_s\psi ,`$
but the last sum is zero since $`e_ke_0e_m\psi ,e_s\psi `$ is purely imaginary whenever $`k,m,s`$ are distinct indices. Thereby it comes out that $`II(X,e_s)=\left(\overline{}_{e_s}k(X,\alpha )\overline{}_\alpha k(X,e_s)\right)`$, so that we can conclude
$$II(X,Y)=\overline{}_Yk(X,\alpha )\overline{}_\alpha k(X,Y),$$
and consequently
$`V\left(\text{Ric}^g+2g+(\text{tr}_gk)kkk\right)(X,Y)`$ $`=`$ $`\left((\delta _gk+\text{dtr}_gk)\alpha \right)(X,Y)`$
$`+\left(\overline{}_Yk(X,\alpha )\overline{}_\alpha k(X,Y)\right).`$
Moreover the couple $`(V,\alpha )`$ satisfies the following differential equations
$`\overline{}_X\alpha (Y)`$ $`=`$ $`Vk(X,Y)+{\displaystyle \frac{𝐢}{2}}((XYYX)\xi ,\xi )`$
$`\delta _g^{}\alpha `$ $`=`$ $`Vk`$
$`\text{d}V(X)`$ $`=`$ $`k(X,\alpha )+\text{i}X\psi ,\psi `$
and
$`(\text{Hess}^gV)(X,Y)`$ $`=`$ $`\overline{}_Yk(X,\alpha )V(kk)(X,Y)+Vg(X,Y)+\overline{}_X\alpha (k(Y))+\overline{}_Y\alpha (k(X))`$
$`=`$ $`\overline{}_Yk(X,\alpha )\overline{}_\alpha k(X,Y)V(kk)(X,Y)+Vg(X,Y)+_\alpha k(Y,X)`$
$`=`$ $`V\left(\text{Ric}^g+3g+(\text{tr}_gk)k2(kk)\right)(X,Y)`$
$`+\left((\delta _gk+\text{dtr}_gk)\alpha \right)(X,Y)+_\alpha k(Y,X)`$
It is clear that the couple $`(V,W):=(V,\alpha ^{\mathrm{}})`$ satisfies the first KID equation , and using the second KID equation for defining the symmetric tensor $`\tau `$ of , namely
$`V\left(\tau {\displaystyle \frac{1}{2}}(\text{tr}_g\tau \rho )g\right)`$ $`=`$ $`V\left(\text{Ric}^g+(\text{tr}_gk)k2(kk)\right)_Wk(\text{Hess}^gV)`$
$`=`$ $`\left((\delta _gk+\text{dtr}_gk)W^{\mathrm{}}\right)3Vg,`$
where $`2\rho :=\text{Scal}^g+(\text{tr}k)^2\left|k\right|^2`$. Taking the trace of last equation one gets $`\text{tr}_g\tau \rho =12`$ and consequently
$$V\tau =\left((\delta _gk+\text{dtr}_gk)W^{\mathrm{}}\right)+3Vg.$$
Now the equation $`\widehat{}\xi ,\xi =0`$ implies $`\left(\text{Scal}^g+n(n1)+(\text{tr}_gk)^2\left|k\right|_g^2\right)=2\left|\delta _gk+\text{dtr}_gk\right|`$. We also know that $`V\left(\text{Scal}^g+6+(\text{tr}_gk)^2\left|k\right|_g^2\right)=2(\delta _gk+\text{dtr}_gk),W^{\mathrm{}}`$, and thereby it is clear that there exists some function on $`M`$ denoted by $`\vartheta `$ such that $`W=\vartheta (\delta _gk+\text{dtr}_gk)`$ and so
$$V\left(\text{Scal}^g+6+(\text{tr}_gk)^2\left|k\right|_g^2\right)=2\left|\vartheta \right|\left|\delta _gk+\text{dtr}_gk\right|^2$$
and therefore
$$2V\left|\delta _gk+\text{dtr}_gk\right|=\left(2\rho +6\right)\left|W\right|,$$
which shows that in the Killing development the Killing vector field $`(V,W)`$ will be colinear to the cosmological contraints 4-vector $`(2\rho +6,2(\delta _gk+\text{dtr}_gk))`$ which is isotropic. It follows that the Killing vector field $`(N,W)`$ is also isotropic in the Killing development. We finally obtain the relation $`V^2(\tau 3g)=\frac{1}{2}(2\rho +6)W^{\mathrm{}}W^{\mathrm{}}`$ which means that the Killing development is a stationary pp-wave space-time.
Supposing furthermore that the constant function $`(\xi ,\xi )`$ is non-zero, we get $`(\delta _gk+\text{dtr}_gk)(\xi ,\xi )=0`$ and so $`(\delta _gk+\text{dtr}_gk)=0`$ by tracing the equation $`\text{d}^\overline{}k(X,Y)e_0\xi ,\xi =0`$. It follows by the dominant energy condition that $`\text{Scal}^g+(\text{tr}k)^2\left|k\right|^2=6`$, and we obtain finally that $`\tau =3g`$. Thereby the couple $`(V,W)`$ is a cosmological vacuum KID. It is known that in that case $`(M,g,k)`$ has a cosmological vacuum Killing development denoted by $`(\overline{N},\overline{\gamma })`$ which is a stationary 4-dimensional Lorentzian manifold satisfying $`G^{\overline{\gamma }}=3\overline{\gamma }`$ and carrying a Killing vector field which is the natural extension of the KID $`(V,W)`$.$`\mathrm{}`$
Remark. It is clear that expecting $`m_0=0`$ so as to define the rigidity situation is much stronger than expecting the degenerate character of $`Q`$. A good issue would certainly be to use the geometry at infinity in the same way as in but in the AdS-asymptotically hyperbolic context, in order to prove under the degenerate character of $`Q`$ the existence of an isometric embedding of $`(M,g,k)`$ in AdS.$`\mathrm{}`$
## 5 Appendix
### 5.1 Non-negativity of $`Q`$ seen through its coefficients when $`n=3`$
Classical linear algebra results state that every principal minor of $`Q`$ must be non-negative which give rise to a set of inequalities on the coefficients of $``$.
$$\begin{array}{c}m_0+m_10\hfill \\ m_0m_10\hfill \end{array}$$
$$\begin{array}{ccc}\hfill m_0^2|m|^2& \hfill & 0\hfill \\ \hfill (m_0+m_1)^2(n_2+r_3)^2(r_2n_3)^2& \hfill & 0\hfill \\ \hfill (m_0m_1)^2(n_2r_3)^2(r_2+n_3)^2& \hfill & 0\hfill \\ \hfill m_0^2m_1^2n_1^2r_1^2& \hfill & 0\hfill \end{array}$$
$$\begin{array}{cc}\hfill (m_0+m_1)(m_0^2(|m|^2+n_1^2+r_1^2))(m_0m_1)((n_2+r_3)^2+(n_3r_2)^2)& \\ \hfill 2((n_2+r_3)(m_2n_1+m_3r_1)+(n_3+r_2)(m_2r_1m_3n_1)& 0\hfill \\ \hfill (m_0m_1)(m_0^2(|m|^2+n_1^2+r_1^2))(m_0+m_1)((n_2r_3)^2+(n_3+r_2)^2)& \\ \hfill +2((n_2r_3)(m_2n_1m_3r_1)+(n_3+r_2)(m_2r_1+m_3n_1)& 0\hfill \end{array}$$
$$\begin{array}{cc}(m_0^2(|m|^2+|n|^2+|r|^2))^24(|m|^2|n|^2+|m|^2|r|^2+|n|^2|r|^2)\hfill & \\ +4(<m,n>^2+<m,r>^2+<n,r>^2)+8m_0\text{det}_^3(m,n,r)\hfill & 0.\hfill \end{array}$$
### 5.2 Rigidity Results for the Trautman-Bondi Mass
Oppositely to the rest of the paper, we consider here the situation of the Trautman-Bondi mass , namely $`(M,g,k)`$ is assumed to be Minkowski-asymptotically hyperbolic which means that the triple $`(M,g,k)`$ is asymptotic at infinity to a standard hyperbolic slice of Minkowski space-time. It has been proved (cf. Theorem 5.4 of ) that the Trautman-Bondi four-momentum $`p_\mu `$ is timelike and future directed under the dominant energy condition (and some other technical assumptions). The aim of this section is to prove some rigidity results for the Trautman-Bondi four-momentum which are analogous to the statements of section 4. More precisely
Theorem. Under the assumptions of Theorem 5.4 of , and if the component $`p_0`$ of the Trautman-Bondi four-momentum $`p_\mu `$ vanishes, then $`(M,g,k)`$ can be isometrically embedded in Minkowski space-time.
Proof. This can be done in the same way as our rigidity theorem: since $`p_\mu `$ is timelike, the condition $`p_0=0`$ implies that $`p_\mu `$ actually vanishes. Consequently there exists a basis of $``$-parallel spinor fields on $`M`$, where $``$ is the connection on some cylinder $`]ϵ,+ϵ[\times M`$ endowed with some Lorentzian metric $`\gamma =\text{d}t^2+g_t`$ (such that $`M`$ has induced metric $`g`$ and extrinsic curvature $`k`$ satisfying the conditions of ). Now if $`\overline{}`$ denotes the Levi-Civita connection of $`g=g_0`$, we still have the relation $`_X\xi =\overline{}_X\xi \frac{1}{2}k(X)e_0\xi `$ where $``$ is the Clifford action with respect to $`\gamma `$. Our spinorial Gauss-Codazzi formula is still valid, that is
$$R_{X,Y}^\gamma =R_{X,Y}^g\frac{1}{2}(\text{d}^\overline{}k(X,Y)e_0+\frac{1}{2}(k(X)k(Y)k(Y)k(X)))=0,$$
and so
$`R^g`$ $`=`$ $`{\displaystyle \frac{1}{2}}k\wedge ⃝k`$
$`\text{d}^\overline{}k`$ $`=`$ $`0.`$
Furthermore, the couple $`(V,W):=(V,\alpha ^{\mathrm{}})`$ is a vacuum KID if one defines $`V=<\xi ,\xi >`$ and the real 1-form $`\alpha `$ by $`\alpha (Y)=Ye_0\xi ,\xi `$. We consider again the Killing development of $`(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{k})`$ with respect to $`(\stackrel{~}{V},\stackrel{~}{W})`$, and observe that it must be a geodesically complete stationary solution of vacuum Einstein equations of zero sectional curvature and thereby must be Minkowski space-time (cf. Proposition 23 \[p.227\] of ). Now the same cohomological arguments give $`(\stackrel{~}{M},\stackrel{~}{g},\stackrel{~}{k})=(M,g,k)`$ which is by construction embedded in its Killing development that is Minkowski. $`\mathrm{}`$
Theorem. Let us suppose that $`(M,g,k)`$ satisfies the assumptions of Theorem 5.4 of and that $`p_\mu `$ is null. Then there exists some $``$-parallel spinor field $`\xi `$ such that $`\xi ,\xi =0`$ and consequently $`(M,g,k)`$ is isometrically embeddable in a stationary pp-wave space-time.
If furthermore the constant function $`(\xi ,\xi )`$ is non-zero then $`(M,g,k)`$ admits a vacuum Killing development which is a stationary solution of the Einstein equations.
Proof. $`p_\mu `$ is null implies the existence of a spinor field $`\xi `$ satisfying the conditions
$`\xi `$ $`=`$ $`0`$
$`\xi ,\xi `$ $`=`$ $`0.`$
Then in the same way as in the last Theorem of section 4, but defining here the 2-tensor $`E=:\text{Ric}^g+(\text{tr}k)kkk`$ we obtain that the couple $`(V,W)`$ is a vacuum KID and the corresponding Killing development satisfies $`V^2\tau =\rho (W^{\mathrm{}}W^{\mathrm{}})`$ which means that it is a stationary pp-wave space-time.
Still using the same computations as in the last Theorem of section 4 and assuming that the constant function $`(\xi ,\xi )`$ is non-zero we find that the constraints equations are satisfied (because of the dominant energy condition) and that $`\tau =0`$. Thereby $`(V,W)`$ is a vacuum KID and it is known that in this case $`(M,g,k)`$ has a stationary vacuum Killing development.
$`\mathrm{}`$
Remark. It is clear that expecting $`p_0=0`$ so as to define the rigidity situation is much stronger than expecting the null character of $`p_\mu `$. As in our situation (cf. the remark at the end of section 4), a good issue would certainly be to use the geometry at infinity in the same way as in but in the Minkowski-asymptotically hyperbolic context, in order to prove under the equality case of Theorem 5.4 of , the existence of an isometric embedding of $`(M,g,k)`$ in Minkowski. $`\mathrm{}`$
Acknowledgements. I wish to thank M. Herzlich for his helpful comments as regard the redaction of this text.
Institut de Mathématiques et de Modélisation de Montpellier (I3M)
Université Montpellier II
UMR 5149 CNRS
Place Eugène Bataillon
34095 MONTPELLIER (FRANCE)
email:[email protected]
fax: 33 (0) 4 67 14 35 58
tel: 33 (0) 4 67 14 48 47
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# The 𝐿¹ Ehrenpreis Conjecture
## 1. Introduction
Let $`\mathrm{\Sigma },\mathrm{\Sigma }^{}`$ be closed Riemann surfaces of genus greater than 1. The most succinct formulation of the Ehrenpreis conjecture (EC) uses the fact that $`\mathrm{\Sigma },\mathrm{\Sigma }^{}`$ may be regarded as riemannian manifolds with metrics of curvature $`1`$. While it is an elementary fact that the riemannian universal covers $`\stackrel{~}{\mathrm{\Sigma }},\stackrel{~}{\mathrm{\Sigma }}^{}`$ are isometric to $`^2`$, the EC asserts a similar, asymptotic phenomenon for the family of finite riemannian covers:
###### Ehrenpreis Conjecture (Hyperbolic).
For any $`ϵ>0`$, there are finite degree isometric covers $`Z\mathrm{\Sigma }`$, $`Z^{}\mathrm{\Sigma }^{}`$ whose total spaces are $`(1+ϵ)`$-quasiisometric.
In , we announced the solution of an $`L^1`$-version of this conjecture. In this paper, we provide the proof.
The traditional or conformal version of the EC can be described in terms of Teichmüller theory. Let $`\mathrm{\Sigma }`$ be a fixed compact surface of genus greater than 1, $`𝒯(\mathrm{\Sigma })`$ its Teichmüller space, $`d_{𝒯(\mathrm{\Sigma })}`$ the Teichmüller metric. Then given $`\mu ,\nu 𝒯(\mathrm{\Sigma })`$, the conformal version of the EC states
###### Ehrenpreis Conjecture (Conformal).
For any $`ϵ>0`$, there exists a surface $`Z`$ and finite covers $`\rho ,\sigma :Z\mathrm{\Sigma }`$ such that
$$d_{𝒯(Z)}(\rho ^{}\mu ,\sigma ^{}\nu )<ϵ.$$
In the above, $`\rho ^{},\sigma ^{}:𝒯(\mathrm{\Sigma })𝒯(Z)`$ are the isometric inclusions induced by $`\rho ,\sigma `$. We remark that this version of the EC makes sense in genus 1, where it is not difficult to verify .
However, it is the genus independent or solenoidal version of the EC that will be most important for us. Let $`\widehat{\mathrm{\Sigma }}`$ be the algebraic universal cover of the closed surface $`\mathrm{\Sigma }`$, by definition the inverse limit of the total spaces of finite covers $`\rho :Z\mathrm{\Sigma }`$ (one cover for each homotopy class of cover). The algebraic universal cover is a surface solenoid (a surface lamination with Cantor transversals), and as such has a Teichmüller space $`𝒯(\widehat{\mathrm{\Sigma }})`$ of marked conformal structures . The mapping class group $`𝖬(\widehat{\mathrm{\Sigma }},\mathrm{})`$ of homotopy classes of homeomorphisms of $`\widehat{\mathrm{\Sigma }}`$ fixing a base leaf $`\mathrm{}`$ may be identified with the group of homotopy classes of lifts of correspondences $`\sigma \rho ^1:\mathrm{\Sigma }\mathrm{\Sigma }`$.
###### Ehrenpreis Conjecture (Solenoidal).
$`𝖬(\widehat{\mathrm{\Sigma }},\mathrm{})`$ acts on $`𝒯(\widehat{\mathrm{\Sigma }})`$ with dense orbits.
The proof that all of these versions of the EC are equivalent can be found in .
The genus independent version of EC says that the ”universal moduli space”
$$(\widehat{\mathrm{\Sigma }})=𝒯(\widehat{\mathrm{\Sigma }})/𝖬(\widehat{\mathrm{\Sigma }},\mathrm{}),$$
although an uncountable set, has the topology of a point (the discrete topology). It has the virtue of giving a certain explanation of the “moduli-rigidity gap” that separates the theory of compact hyperbolic surfaces from that of compact hyperbolic manifolds in dimension three or greater. From a practical point of view, working with $`\widehat{\mathrm{\Sigma }}`$ allows us to regard Teichmüller theory of closed hyperbolic surfaces as concerning complex structures on a single topological type (as in the genus 1 case). In this way, we may isolate geometric properties of closed Riemann surfaces of hyperbolic type that do not depend on genera.
In this paper, we shall formulate and prove an $`L^1`$ version of the EC. On $`𝒯(\widehat{\mathrm{\Sigma }})`$ there are – in addition to the Teichmüller metric – three other metrics coming from the $`L^1`$, the $`L^{\mathrm{}}`$ and the $`L^2`$ structures on the cotangent bundle of $`𝒯(\widehat{\mathrm{\Sigma }})`$. The $`L^1`$ version of the EC is obtained by asking that $`𝖬(\widehat{\mathrm{\Sigma }},\mathrm{})`$ act densely on $`𝒯(\widehat{\mathrm{\Sigma }})`$ with regard to the $`L^1`$ geometry.
The proof of the $`L^1`$ EC is as follows. A dense set of pairs $`\widehat{\mu },\widehat{\nu }𝒯(\widehat{\mathrm{\Sigma }})`$ (dense in the Teichmüller geometry) lie along the axis $`A`$ of a pseudo-Anosov homeomorphism $`\widehat{\mathrm{\Phi }}:\widehat{\mathrm{\Sigma }}\widehat{\mathrm{\Sigma }}`$. $`A`$ is a Teichmüller geodesic and the action of $`\widehat{\mathrm{\Phi }}`$ on $`𝒯(\widehat{\mathrm{\Sigma }})`$ stabilizes $`A`$, translating points along $`A`$ a distance of $`\frac{1}{2}\mathrm{log}\lambda `$, where $`\lambda `$ is the entropy of $`\widehat{\mathrm{\Phi }}`$. Given $`n`$, by an $`L^1`$ $`n`$th root of $`\mathrm{\Phi }`$ we mean a sequence of pseudo Anosov homeomorphisms $`\{\widehat{\mathrm{\Psi }}_m\}`$ in which
* The entropies $`\lambda _m`$ of $`\widehat{\mathrm{\Psi }}_m`$ converge to $`\lambda ^{1/n}`$.
* The axes $`A_m`$ of $`\widehat{\mathrm{\Psi }}_m`$ converge to $`A`$ in the $`L^1`$ Hausdorff topology.
We shall show that $`L^1`$ $`n`$th roots exist for every (lifted) pseudo Anosov homeomorphism of $`\widehat{\mathrm{\Sigma }}`$. This will then imply the $`L^1`$ version of EC.
Acknowledgements: This work benefited greatly from conversations with Dennis Sullivan and Yair Minksy.
## 2. Topology of the Algebraic Universal Cover
Let $`\mathrm{\Sigma }`$ be a fixed compact surface of genus at least two. We describe in this section the topology of the algebraic universal cover $`\widehat{\mathrm{\Sigma }}`$. Unless otherwise noted, all proofs of statements in this section can be found in .
Let $`\pi =\pi _1\mathrm{\Sigma }`$. For every finite index normal subgroup $`H<\pi `$, choose a pointed cover $`\rho :(Z,x_Z)(\mathrm{\Sigma },x)`$ for which $`\rho _{}\pi _1Z=H`$. By adding to this collection of covers all covers $`\tau :ZZ^{}`$ between total spaces for which $`\rho ^{}\tau =\rho `$, we obtain an inverse system of surfaces. The limit of this system $`\widehat{\mathrm{\Sigma }}`$ is called the algebraic universal cover of $`\mathrm{\Sigma }`$, a compact topological space. Its topological type is independent of the choice of covers.
If we denote by $`\widehat{\pi }`$ the profinite completion of $`\pi `$, then $`\widehat{\mathrm{\Sigma }}`$ is homeomorphic to the quotient
(1)
$$\left(\stackrel{~}{\mathrm{\Sigma }}\times \widehat{\pi }\right)/\pi ,$$
where $`\pi `$ acts diagonally, and so has the structure of a surface solenoid: a surface lamination whose model transversals are Cantor sets. It also follows from (1) that $`\widehat{\mathrm{\Sigma }}`$ is connected, its path components are its leaves, and each leaf is homeomorphic to $`^2`$ and dense in $`\widehat{\mathrm{\Sigma }}`$. The point $`\widehat{x}=(x_Z)\widehat{\mathrm{\Sigma }}`$ – defined by the string of basepoints of the surfaces in the defining system – is contained in a leaf $`\mathrm{}`$ which we call the base leaf. The Haar measure on $`\widehat{\pi }`$ induces a transverse invariant measure $`\eta `$ on $`\widehat{\mathrm{\Sigma }}`$ which gives measure $`1/\mathrm{deg}Z`$ to the fibers of the natural projection $`\widehat{\mathrm{\Sigma }}Z`$.
Any pointed finite cover $`\sigma :(Y,y)(\mathrm{\Sigma },x)`$ lifts to a base leaf preserving homeomorphism $`\widehat{\sigma }:\widehat{Y}\widehat{\mathrm{\Sigma }}`$, where $`\widehat{Y}`$ is the algebraic universal cover of $`Y`$. If $`\rho :(Y,y)(\mathrm{\Sigma },x)`$ is another such cover, the correspondence $`\sigma \rho ^1`$ lifts to the homeomorphism $`\widehat{\sigma }\widehat{\rho }^1`$ of $`\widehat{\mathrm{\Sigma }}`$ preserving $`\mathrm{}`$. Let $`𝖬(\widehat{\mathrm{\Sigma }},\mathrm{})`$ denote the the group of homotopy classes of orientation preserving homeomorphisms of $`\widehat{\mathrm{\Sigma }}`$ preserving $`\mathrm{}`$. See , for a proof of the following
###### Theorem 1.
Every class $`[h]𝖬(\widehat{\mathrm{\Sigma }},\mathrm{})`$ contains an element of the form $`\widehat{\sigma }\widehat{\rho }^1`$.
## 3. Measured Laminations on $`\widehat{\mathrm{\Sigma }}`$
We begin by recalling a few facts about measured laminations on $`\mathrm{\Sigma }`$, a closed Riemann surface of genus $`g>1`$. See , , , for further discussion. A measured lamination $`𝔣`$ on $`\mathrm{\Sigma }`$ is a closed 1-dimensional lamination smoothly embedded in $`\mathrm{\Sigma }`$ and possessing a transverse invariant measure $`m_𝔣`$. Two measured laminations are equivalent if they are isotopic through an isotopy taking one measure to the other. The set of equivalence classes of measured laminations is denoted $`(\mathrm{\Sigma })`$. Let $`𝒞(\mathrm{\Sigma })`$ denote the set of isotopy classes of simple closed curves in $`\mathrm{\Sigma }`$. Given $`𝔣(\mathrm{\Sigma })`$ and $`c𝒞(\mathrm{\Sigma })`$, the intersection pairing is defined
$$\text{I}(𝔣,c)=inf_{𝔣c}𝑑m_𝔣,$$
where the infimum is taken over representatives of the classes of $`𝔣`$ and $`c`$. The intersection topology on $`(\mathrm{\Sigma })`$ is the weak topology with respect to the intersection pairing. The space of projective classes of measured laminations is denoted $`𝒫(\mathrm{\Sigma })`$ and is homeomorphic to a sphere of dimension $`6g7`$. We have $`𝒞(\mathrm{\Sigma })𝒫(\mathrm{\Sigma })`$ with dense image. The intersection pairing extends to a map $`(\mathrm{\Sigma })\times (\mathrm{\Sigma })`$ via the formula
$$\text{I}(𝔣,𝔤)=inf_{𝔣𝔤}𝑑m_𝔣𝑑m_𝔤.$$
A word is in order here regarding the allied concept of a measured foliation, a singular foliation $``$ of $`\mathrm{\Sigma }`$ equipped with a transverse invariant measure: these typically arise as trajectories of holomorphic quadratic differentials on $`\mathrm{\Sigma }`$ . Two measured foliations are equivalent if after a finite number of Whitehead moves are applied to their singular leaves, they are isotopic through an isotopy taking one measure to the other. There is a bijective correpondence between classes of measured foliations and classes of measured laminations . For example, to obtain a measured lamination starting with a measured foliation $``$, one chooses a nonsingular leaf from each minimal component of $``$, pulls each such leaf geodesic (with respect to the hyperbolic metric of $`\mathrm{\Sigma }`$) and completes the resulting space. Our default will be to work with measured laminations, and – with the exception of the proof of Theorem 4, where we revert back to measured foliations – whenever a measured foliation happens to arise, we will assume it has been converted into its associated measured lamination.
A (homotopy class of) homeomorphism $`\mathrm{\Phi }:\mathrm{\Sigma }\mathrm{\Sigma }`$ induces a homeomorphism of $`(\mathrm{\Sigma })`$ via pullback of measures, in particular inducing a homeomorphism of $`𝒫(\mathrm{\Sigma })`$. According to the classification of surface diffeomorphisms, , , , $`\mathrm{\Phi }`$ is called pseudo Anosov if its induced action on $`𝒫(\mathrm{\Sigma })`$ fixes precisely two classes $`[𝔣^u]`$ and $`[𝔣^s]`$. If $`\lambda `$ is the entropy of $`\mathrm{\Phi }`$, then $`\lambda >1`$; and if $`𝔣^u[𝔣^u]`$ ($`𝔣^s[𝔣^s]`$) is a representative in $`(\mathrm{\Sigma })`$, then there is a representative diffeomorphism in the class of $`\mathrm{\Phi }`$ (also denoted $`\mathrm{\Phi }`$) such that $`\mathrm{\Phi }(𝔣^u)=\lambda 𝔣^u`$ ($`\mathrm{\Phi }(𝔣^s)=\lambda ^1𝔣^s`$).
In this paper, we will be interested in the following class of pseudo Anosov homeomorphisms. Let $`𝒞`$, $`𝒟`$ be families of pairwise nonisotopic simple closed curves, for which elements of $`𝒞`$ intersect minimally in their isotopy classes with elements of $`𝒟`$, and for which $`\mathrm{\Sigma }(𝒞𝒟)`$ consists of a union of disks (the families are then said to be filling). For $`c𝒞`$, $`d𝒟`$, let $`F_c`$ resp. $`G_d`$ denote the right Dehn twist about $`c`$ resp. $`d`$. Then a homeomorphism of the form
$$\mathrm{\Phi }=G_{d_k}^{N_k}\mathrm{}G_{d_1}^{N_1}F_{c_j}^{M_j}\mathrm{}F_{c_1}^{M_1},$$
where the exponents $`M_1,\mathrm{},M_j,N_1,\mathrm{},N_k`$ are positive and where all curves in $`𝒞`$, $`𝒟`$ occur, is pseudo Anosov, , . We call these pseudo Anosovs of Thurston-Penner type.
We now extend the above considerations to $`\widehat{\mathrm{\Sigma }}`$. A measured lamination $`\widehat{𝔣}`$ on $`\widehat{\mathrm{\Sigma }}`$ is a collection of measured laminations $`\{𝔣_{\mathrm{}}\}`$, one on each leaf $`\mathrm{}`$ of $`\widehat{\mathrm{\Sigma }}`$, which have the same transversal model $`𝖳`$ and which vary in the following way with respect to the transversals of $`\widehat{\mathrm{\Sigma }}`$. Let $`𝒪D\times \widehat{T}`$ be a flowbox for $`\widehat{\mathrm{\Sigma }}`$, such that $`𝔣_t:=\widehat{𝔣}|_{D\times \{t\}}`$ is a flowbox for $`𝔣_{\mathrm{}}`$ if $`D\times \{t\}\mathrm{}`$. Then
1. The family of flowboxes $`𝔣_t`$ varies continuously in $`t`$. (Thus $`\widehat{𝔣}`$ gives rise to a smooth 1-dimensional sublamination of $`\widehat{\mathrm{\Sigma }}`$ with transversal models $`\widehat{T}\times 𝖳`$.)
2. Given any continuous family of test transversals $`𝖳_tD\times \{t\}`$, the function obtained by pairing with the measures of $`\widehat{𝔣}`$ is continuous in $`t`$.
The space of equivalence classes of measured laminations is denoted $`(\widehat{\mathrm{\Sigma }})`$. If $`c`$ is a simple closed curve occurring in some surface $`Z`$ in the defining system of $`\widehat{\mathrm{\Sigma }}`$, its preimage $`\widehat{c}`$ in $`\widehat{\mathrm{\Sigma }}`$ is a $`1`$-dimensional solenoid which we call a simple closed solenoid (abusively, since such a $`\widehat{c}`$ always has infinitely many connected components). The set of such is denoted $`𝒞(\widehat{\mathrm{\Sigma }})`$.
The intersection pairing $`\text{I}(\widehat{𝔣},\widehat{c})`$ between a measured lamination and a simple closed solenoid is defined using the transverse invariant measure $`\eta `$,
$$\text{I}(\widehat{𝔣},\widehat{c})=_{\widehat{𝔣}\widehat{c}}𝑑m_{\widehat{𝔣}}𝑑\eta .$$
We equip $`(\widehat{\mathrm{\Sigma }})`$ with the resulting weak topology. With its induced topology, $`𝒫(\widehat{\mathrm{\Sigma }})`$ is precompact (being essentially a space of probability measures), but owing to its infinite-dimensionality, $`𝒫(\widehat{\mathrm{\Sigma }})`$ is not compact. The simple closed solenoids $`𝒞(\widehat{\mathrm{\Sigma }})`$ are dense in $`𝒫(\widehat{\mathrm{\Sigma }})`$.
If $`Z`$ is a surface occurring in the defining limit of $`\widehat{\mathrm{\Sigma }}`$, then we may pullback measured laminations on $`Z`$ to measured laminations on $`\widehat{\mathrm{\Sigma }}`$. The result is a direct system of inclusions $`(Z)(\widehat{\mathrm{\Sigma }})`$ whose limit $`^{\mathrm{}}:=lim_{}(Z)`$ has dense image in $`(\widehat{\mathrm{\Sigma }})`$. In this paper we will work exclusively with $`^{\mathrm{}}`$.
A pseudo Anosov diffeomorphism $`\mathrm{\Phi }:ZZ`$ lifts to a diffeomorphism of $`\widehat{\mathrm{\Sigma }}`$ fixing precisely the lifted projective classes $`[\widehat{𝔣}^s]`$ and $`[\widehat{𝔣}^u]`$. In what follows, the terminology “pseudo Anosov homeomorphism of $`\widehat{\mathrm{\Sigma }}`$” will always mean such a lift.
## 4. Train Tracks
We recall first some facts about train tracks on $`\mathrm{\Sigma }`$: details may be found in . Let $`\tau \mathrm{\Sigma }`$ be a smooth 1-dimensional branched manifold: thus $`\tau `$ is a 1-dimensional CW-complex in which the interiors of edges are smooth curves, and the field of tangent lines $`\mathrm{T}_x\tau `$, $`x\tau \{\text{vertices}\}`$, extends to a continuous line field on $`\tau `$. We say that $`\tau `$ is a train track if it satisfies the following additional properties:
1. The valency of any vertex is at least 3, except for simple closed curve components, which have a single vertex of valence 2.
2. If $`D(S)`$ is the double of a component $`S\mathrm{\Sigma }\tau `$, then the Euler characteristic of $`D(S)`$ is negative.
We shall follow the custom of referring to the vertices of a train track as switches.
A bigon track is a smooth 1-dimensional branched manifold $`\tau \mathrm{\Sigma }`$ satisfying item (1) and which satisfies (2) after collapsing bigon complementary regions to curves. Bigon tracks arise naturally from a pair $`𝒞`$, $`𝒟`$ of transverse filling curves, by turning each intersection of a $`𝒞`$ curve with a $`𝒟`$ curve into a pair of 3-valent vertices as in Figure 1. Since such bigon tracks will be the only ones appearing in this article, we will assume from now on that all switches in bigon tracks have valency no more than three.
Denote by $`E`$ the set of edges of the bigon track $`\tau `$. In a small disk neighborhood of a switch $`v`$, the ends of edges incident to $`v`$ may be divided into two classes, which for convenience we refer to as ”incoming” and ”outgoing”: each class consists of ends that are asymptotic to one another, and the decision of naming one class incoming, the other outgoing, is arbitrary. We write $`e\mathrm{𝗂𝗇}(v)`$ or $`e\mathrm{𝗈𝗎𝗍}(v)`$ if $`e`$ has an end belonging to the appropriate class. See Figure 2. (Note: it can happen that $`e`$ belongs to both $`\mathrm{𝗂𝗇}(v)`$ and $`\mathrm{𝗈𝗎𝗍}(v)`$.)
A switch-additive measure on $`\tau `$ is a function $`m:E_+`$ for which
$$\underset{e\mathrm{𝗂𝗇}(v)}{}m(e)=\underset{e\mathrm{𝗈𝗎𝗍}(v)}{}m(e)$$
for all switches $`v`$. The set of all switch-additive measures forms a linear cone $`𝖢_\tau `$ in $`^E`$.
Let $`N(\tau )`$ be a tubular neighborhood of $`\tau `$ equipped with a (singular) foliation by line segments transverse to $`\tau `$. See Figure 3.
A measured lamination $`𝔣\mathrm{\Sigma }`$ is said to be carried by $`\tau `$ if it may by isotoped into $`N(\tau )`$ transverse to its foliation. We write in this case $`𝔣<\tau `$. The subspace of isotopy classes of measured laminations carried by $`\tau `$ is denoted $`_\tau (\mathrm{\Sigma })`$. There is an open surjection
$$𝖢_\tau _\tau (\mathrm{\Sigma })$$
which is a homeomorphism if $`\tau `$ is a train track.
Let $`\mathrm{\Phi }:\mathrm{\Sigma }\mathrm{\Sigma }`$ be a pseudo Anosov diffeomorphism. We say that $`\mathrm{\Phi }`$ acts on $`\tau `$ if $`\mathrm{\Phi }(\tau )`$ may be isotoped into $`N(\tau )`$ transverse to its foliation. We write then $`\mathrm{\Phi }(\tau )<\tau `$. Fix a leaf $`t_iN(\tau )`$ through each edge $`e_i`$ of $`\tau `$. The carrying matrix of $`\mathrm{\Phi }`$ is by definition $`M_\mathrm{\Phi }=(a_{ij})`$ where
$$a_{ij}=|\mathrm{\Phi }(e_i)t_j|.$$
$`M_\mathrm{\Phi }`$ induces an inclusion $`𝖢_{\mathrm{\Phi }(\tau )}𝖢_\tau `$ which when precomposed with the pushforward map $`𝖢_\tau 𝖢_{\mathrm{\Phi }(\tau )}`$ defines a linear map
$$M_\mathrm{\Phi }:𝖢_\tau 𝖢_\tau .$$
Note that the carrying matrix $`M_\mathrm{\Phi }`$ is non-negative. Such a matrix has a unique eigenvalue of greatest modulus, which is positive-real and simple . This eigenvalue is called the Perron root. A corresponding eigenvector may be taken non-negative, and is called a Perron vector. For $`M_\mathrm{\Phi }`$, the Perron root coincides with the entropy $`\lambda `$ of $`\mathrm{\Phi }`$, and the Perron vector parametrizes in track coordinates the unstable measured lamination $`𝔣^u`$ of $`\mathrm{\Phi }`$.
When $`\mathrm{\Phi }`$ is a pseudo Anosov of Thurston-Penner type, one can recover the carrying matrix and all of its Perron data from a simpler matrix which records the action of $`\mathrm{\Phi }`$ on the curves in the families $`𝒞`$ and $`𝒟`$. Indeed, let $`\tau `$ be the bigon track formed from $`𝒞𝒟`$. If $`e_i,e_jE_\tau `$ are edges contained in say $`c`$, $`d`$ resp. then given $`e_i^{}c`$ another edge, there exists a unique $`e_j^{}d`$ with $`a_{i^{}j^{}}=a_{ij}`$. Conversely, for $`e_j^{}d`$, there exists $`e_i^{}c`$ with $`a_{i^{}j^{}}=a_{ij}`$. It follows that the carrying matrix of $`\mathrm{\Phi }`$ can be subdivided into blocks indexed by pairs $`(c,d)`$, which are of the form $`a_{c,d}I`$ where $`I`$ is a square matrix in which each column and row has exactly one non zero entry = 1. The matrix $`(a_{c,d})`$ whose columns and rows are indexed by $`𝒞𝒟`$ has exactly the same Perron root as $`M_\mathrm{\Phi }`$, and its Perron vector gives that of $`M_\mathrm{\Phi }`$ in the obvious way. We shall call this matrix the curve matrix of the Thurston-Penner type pseudo Anosov $`\mathrm{\Phi }`$, and we shall denote it $`M_\mathrm{\Phi }`$ as well.
For example, if $`𝒞=\{c\}`$, $`𝒟=\{d\}`$ with $`|cd|=r`$ and $`\mathrm{\Phi }=G_d^NF_c^N`$, then
$$M_\mathrm{\Phi }=\left(\begin{array}{cc}1\hfill & rN\hfill \\ rN\hfill & (rN)^2+1\hfill \end{array}\right).$$
We note that using the quadratic formula, it is easy to see that the eigenvalue not equal to the Perron root is $`<1`$.
We now discuss tracks on the solenoid $`\widehat{\mathrm{\Sigma }}`$: in fact, we will only require tracks pulled back from surfaces appearing in its defining inverse system. Thus, if $`\tau `$ is a train track on such a surface $`Z`$, its preimage $`\widehat{\tau }`$ is a smooth 1-dimensional branched solenoid with edge set $`\widehat{E}E\times \widehat{T}_Z`$, where $`E`$ is the edge set of $`\tau `$ and $`\widehat{T}_Z`$ is the fiber over a point of $`Z`$, homeomorphic to the Cantor group $`\widehat{\pi }_1Z`$. With respect to this decomposition we define a measure on $`\widehat{E}`$ by the formula
$$\mu _{\widehat{E}}=\mu _E\times \eta _{\widehat{T}_Z}$$
where $`\mu _E(e)=1`$ for each edge of $`E`$ and $`\eta _{\widehat{T}_Z}`$ is the restriction to $`\widehat{T}_Z`$ of the transverse invariant measure of $`\widehat{\mathrm{\Sigma }}`$. In addition, if $`\tau `$ is equipped with a switch additive measure $`\upsilon `$, the pullback $`\widehat{\upsilon }`$ is a transversally continuous switch-additive measure on $`\widehat{\tau }`$; the cone of such measures on $`\widehat{\tau }`$ is denoted $`𝖢_{\widehat{\tau }}`$. The relation $`\widehat{𝔣}<\widehat{\tau }`$ has exactly the same meaning as in the case of a surface.
## 5. Intersection Formulas
Let $`\tau ,\kappa `$ be bigon tracks in $`\mathrm{\Sigma }`$ that intersect transversally and minimally with edge sets $`E_\tau `$ and $`E_\kappa `$; let $`𝔣,𝔤`$ be measured laminations carried by them, parametrized by weights $`\upsilon ,\omega `$. The intersection pairing may be calculated by the following formula :
(2)
$$\text{I}(𝔣,𝔤)=\underset{eE_\tau ,e^{}E_\kappa }{}\upsilon (e)\omega (e^{}).$$
It is useful to re-express (2) as a sum over edges in $`E_\tau `$ only. Thus if we write
$$\omega (e)=\underset{e^{}E_\kappa }{}\omega (e^{})|ee^{}|$$
then
(3)
$$\text{I}(𝔣,𝔤)=\underset{eE_\tau }{}\upsilon (e)\omega (e).$$
Suppose now that $`\widehat{𝔣}`$, $`\widehat{𝔤}`$ are measured laminations obtained as preimages of measured laminations $`𝔣Y`$ and $`𝔤Z`$, surfaces occurring in the defining system of $`\widehat{\mathrm{\Sigma }}`$. Let $`W`$ be a surface finitely covering each of $`Y,Z`$, and let $`\stackrel{~}{𝔣}`$, $`\stackrel{~}{𝔤}`$ be the preimages in $`W`$ of $`𝔣`$, $`𝔤`$. Let $`\mathrm{deg}(W)`$ be the degree of the covering $`W\mathrm{\Sigma }`$.
###### Proposition 1.
$`\text{I}(\widehat{𝔣},\widehat{𝔤})=\frac{\text{I}(\stackrel{~}{𝔣},\stackrel{~}{𝔤})}{\mathrm{deg}(W)}`$.
###### Proof.
The intersection locus of $`\widehat{𝔣}`$ and $`\widehat{𝔤}`$ is of the form
$$(\stackrel{~}{𝔣}\stackrel{~}{𝔤})\times \widehat{T}_W$$
where $`\widehat{T}_W`$ is a fiber of $`\widehat{\mathrm{\Sigma }}W`$. Since $`\widehat{T}_W`$ has $`\eta `$-measure $`1/\mathrm{deg}(W)`$, the result follows. ∎
Let $`\widehat{𝔣},\widehat{𝔤}`$ be as in the previous paragraphs. Suppose now that $`𝔣`$ is parametrized by $`\upsilon :E_+`$ a weight on a bigon track $`\tau Y`$ and $`𝔤`$ is parametrized by $`\omega :E^{}_+`$ a weight on a bigon track $`\tau ^{}Z`$. The preimages $`\stackrel{~}{𝔣},\stackrel{~}{𝔤}`$ are parametrized by the pullback weights $`\stackrel{~}{\upsilon }`$, $`\stackrel{~}{\omega }`$ on the preimages $`\stackrel{~}{\tau },\stackrel{~}{\tau }^{}`$. Rewriting $`\omega `$ as above as a function of the edge set $`E`$, we have
$$\text{I}(\stackrel{~}{𝔣},\stackrel{~}{𝔤})=\underset{\stackrel{~}{e}\stackrel{~}{E}}{}\stackrel{~}{\upsilon }(\stackrel{~}{e})\stackrel{~}{\omega }(\stackrel{~}{e}).$$
Let $`\widehat{\tau }`$ be the preimage of $`\tau `$ in $`\widehat{\mathrm{\Sigma }}`$, and let $`\widehat{\upsilon }`$, $`\widehat{\omega }`$ be the pullbacks along the projection $`\widehat{E}\stackrel{~}{E}`$.
###### Proposition 2.
$`\text{I}(\widehat{𝔣},\widehat{𝔤})=_{\widehat{E}}\widehat{\upsilon }\widehat{\omega }𝑑\mu _{\widehat{E}}`$ where $`\mu _{\widehat{E}}`$ is the edge measure on $`\widehat{E}`$.
###### Proof.
A calculation:
$$\text{I}(\widehat{𝔣},\widehat{𝔤})=\left(\underset{\stackrel{~}{e}\stackrel{~}{E}}{}\stackrel{~}{\upsilon }(\stackrel{~}{e})\stackrel{~}{\omega }(\stackrel{~}{e})\right)\frac{1}{\mathrm{deg}(W)}=\left(\underset{\stackrel{~}{e}\stackrel{~}{E}}{}\stackrel{~}{\upsilon }(\stackrel{~}{e})\stackrel{~}{\omega }(\stackrel{~}{e})\right)|\widehat{T}_W|=_{\widehat{E}}\widehat{\upsilon }\widehat{\omega }𝑑\mu _{\widehat{E}}.$$
## 6. Teichmüller Theory of $`\widehat{\mathrm{\Sigma }}`$
References for material in this section are , , .
The definition of the Teichmüller space $`𝒯(\widehat{\mathrm{\Sigma }})`$ and of its metric $`d_{𝒯(\widehat{\mathrm{\Sigma }})}`$ copies that of a surface. In particular,
1. A conformal structure on $`\widehat{\mathrm{\Sigma }}`$ is determined by a conformal structure on each leaf. These structures are required to vary continuously in the transverse direction.
2. Elements of $`𝒯(\widehat{\mathrm{\Sigma }})`$ are represented by marked solenoids i.e. by homeomorphisms $`\mu :\widehat{\mathrm{\Sigma }}\widehat{\mathrm{\Sigma }}_\mu `$, where $`\widehat{\mathrm{\Sigma }}_\mu `$ is presumed to have a conformal structure.
3. The marked solenoid $`\mu ^{}:\widehat{\mathrm{\Sigma }}\widehat{\mathrm{\Sigma }}_\mu ^{}`$ is equivalent to $`\mu `$ if there exists an isomorphism $`\sigma :\widehat{\mathrm{\Sigma }}_\mu \widehat{\mathrm{\Sigma }}_\mu ^{}`$ such that $`\sigma \mu \mu ^{}`$.
$`𝒯(\widehat{\mathrm{\Sigma }})`$ has the structure of a separable Banach manifold. The canonical projection $`\widehat{p}:\widehat{\mathrm{\Sigma }}Z`$ onto any surface $`Z`$ in the defining inverse system induces a direct system of isometric inclusions $`\widehat{p}^{}:𝒯(Z)𝒯(\widehat{\mathrm{\Sigma }})`$.
###### Theorem 2 ().
The induced inclusion
$$i:\underset{}{lim}𝒯(Z)𝒯(\widehat{\mathrm{\Sigma }})$$
is isometric with dense image.
For most of our purposes, it will be sufficient to work with the dense subspace
$$𝒯^{\mathrm{}}:=i(\underset{}{lim}𝒯(Z)),$$
which is an incomplete metric space with respect to the direct limit of the Teichmüller metrics and a pre-Banach manifold. Unless otherwise said, all structures $`\widehat{\mu }`$ considered below will be assumed to be in $`𝒯^{\mathrm{}}`$.
Let $`\widehat{\mathrm{\Sigma }}_{\widehat{\mu }}`$ be as above. By a holomorphic quadratic differential $`\widehat{q}`$ on $`\widehat{\mathrm{\Sigma }}_{\widehat{\mu }}`$, we shall always mean the pull-back of a holomorphic quadratic differential $`q`$ occurring on some surface $`Z_\mu `$, where $`\widehat{\mu }`$ is the pull-back of $`\mu `$. Thus, $`\widehat{q}`$ is a choice of holomorphic quadratic differential on each leaf, constant along the fiber transversals $`\widehat{T}_Z`$ over $`Z`$. The tangent space to $`𝒯^{\mathrm{}}`$ at $`\widehat{\mu }`$ may be identified with the direct limit
$$Q_{\widehat{\mu }}^{\mathrm{}}=\underset{}{lim}Q_\mu (Z).$$
The tangent bundle of $`𝒯^{\mathrm{}}`$ is then identified with $`𝒬^{\mathrm{}}=lim_{}𝒬(Z)`$ where $`𝒬(Z)`$ is the space of holomorphic quadratic differentials on $`Z`$ (with respect to all possible complex structures).
The $`L^1`$ norm on $`Q_{\widehat{\mu }}^{\mathrm{}}`$ is defined
$$\widehat{q}=_{\widehat{\mathrm{\Sigma }}_\mu }|\widehat{q}|𝑑\eta .$$
From the pre-Finsler norm $``$ we induce a path metric $`d_{L^1}`$ on $`𝒯^{\mathrm{}}`$ that defines 1) the $`L^1`$ topology on $`𝒯^{\mathrm{}}`$ and 2) along with $``$, the $`L^1`$ topology on $`𝒬^{\mathrm{}}`$.
To any quadratic diferential $`\widehat{q}`$ one associates two transverse, measured laminations $`\widehat{𝔣}^h`$ and $`\widehat{𝔣}^v`$ (those that correspond to the horizontal and vertical trajectories of $`\widehat{q}`$). We have the following generalization to $`\widehat{\mathrm{\Sigma }}`$ of a well-known formula for surfaces:
###### Lemma 1.
$`\text{I}(\widehat{𝔣}^h,\widehat{𝔣}^v)=\widehat{q}`$ .
###### Proof.
$`\widehat{q}`$ is the lift of a holomorphic quadratic differential on some surface $`Z_\mu `$, where $`\widehat{\mu }`$ is the lift of $`\mu `$. The result now follows from the classical formula and the fact that the lift of an area measure on $`Z_\mu `$ scales by $`1/(\mathrm{deg}Z)`$ in $`\widehat{\mathrm{\Sigma }}`$. ∎
In the same way, we may also avail ourselves of a direct limit version of the theorem of Hubbard and Masur :
###### Theorem 3.
Any pair of measured laminations $`\widehat{𝔣},\widehat{𝔤}^{\mathrm{}}`$ determines a unique quadratic differential $`\widehat{q}`$.
###### Theorem 4.
Let $`\widehat{q},\widehat{q}_i𝒬^{\mathrm{}}`$, $`i=1,2,`$, be quadratic differentials. If $`\widehat{𝔣}_i^h\widehat{𝔣}^h`$ and $`\widehat{𝔣}_i^v\widehat{𝔣}^v`$ in the intersection topology, then $`\widehat{q}_i\widehat{q}`$ in the $`L^1`$ topology.
###### Proof.
We assume first that there exists $`\widehat{\mu }𝒯^{\mathrm{}}`$ with $`\widehat{q}_i,\widehat{q}Q_{\widehat{\mu }}^{\mathrm{}}`$. Let $`\widehat{𝖿}^h`$, $`\widehat{𝖿}^v`$ and $`\widehat{𝖿}_i^h`$, $`\widehat{𝖿}_i^v`$ be the pairs of measured foliations which are the horizontal and vertical line fields of $`\widehat{q}`$ resp. $`\widehat{q}_i`$. By the comments of §3, the hypothesis on the convergence of measured laminations is equivalent to the corresponding statement for measured foliations. In particular we have,
(4)
$$lim\text{I}(\widehat{𝖿}^h,\widehat{𝖿}_i^h)=\mathrm{\hspace{0.33em}\hspace{0.33em}0}\text{and}lim\text{I}(\widehat{𝖿}^v,\widehat{𝖿}_i^v)=\mathrm{\hspace{0.33em}\hspace{0.33em}0}.$$
Consider smooth measured foliations $`\widehat{𝗀}_i^h`$, $`\widehat{𝗀}_i^v`$ equivalent to $`\widehat{𝖿}_i^h`$, $`\widehat{𝖿}_i^v`$ whose heights with respect to $`\widehat{𝖿}^h`$, $`\widehat{𝖿}^v`$ nearly give the intersections $`\text{I}(\widehat{𝖿}^h,\widehat{𝖿}_i^h)`$, $`\text{I}(\widehat{𝖿}^v,\widehat{𝖿}_i^v)`$. More precisely, for $`ϵ_i0`$,
(5)
$$_{\widehat{𝗀}_i^h}|\mathrm{𝖨𝗆}\sqrt{\widehat{q}}|𝑑\eta \text{I}(\widehat{𝖿}^h,\widehat{𝖿}_i^h)<ϵ_i\text{ and }_{\widehat{𝗀}_i^v}|\mathrm{𝖱𝖾}\sqrt{\widehat{q}}|d\eta \text{I}(\widehat{𝖿}^v,\widehat{𝖿}_i^v)<ϵ_i.$$
Let $`\widehat{q}_i^{}`$ denote the smooth quatratic differential whose horizontal and vertical foliations are $`\widehat{𝗀}_i^h`$, $`\widehat{𝗀}_i^v`$. Given $`\delta >0`$, let $`\widehat{A}_i`$ be the set of points for which $`|\widehat{q}\widehat{q}_i^{}|`$ is uniformly $`\delta `$ small. Let $`\widehat{B}_i=\widehat{\mathrm{\Sigma }}\widehat{A}_i`$.
We begin by showing that
(6)
$$_{\widehat{B}_i}|\widehat{q}\widehat{q}_i^{}|𝑑\eta 0$$
as $`i0`$. By (4) and (5), $`_{\widehat{B}_i}|\widehat{q}|𝑑\eta 0`$. If for $`i`$ large, there is a $`m>0`$ with
$$0<m_{\widehat{B}_i}|\widehat{q}\widehat{q}_i^{}|d\eta $$
we must also have
(7)
$$0<m_0_{\widehat{B}_i}|\widehat{q}_i^{}|d\eta $$
for some $`m_0>0`$. The fact that $`_{\widehat{B}_i}|\widehat{q}|𝑑\eta 0`$ whereas $`_{\widehat{B}_i}|\widehat{q}_i^{}|𝑑\eta `$ does not would violate (4) and (5) as well. Thus $`_{\widehat{B}_i}|\widehat{q}_i|𝑑\eta 0`$, proving (6). In particular, we have
$$lim|\widehat{q}\widehat{q}_i^{}|𝑑\eta =lim_{\widehat{A}_i}|\widehat{q}\widehat{q}_i^{}|.$$
Now $`\widehat{q}_i^{}`$ is measure equivalent to $`\widehat{q}_i`$, hence $`\widehat{q}\widehat{q}_i`$ is measure equivalent to $`\widehat{q}\widehat{q}_i^{}`$. By the second minimal norm property , it follows that
$$\widehat{q}\widehat{q}_i\widehat{q}\widehat{q}_i^{}.$$
Letting $`ϵ0`$, we obtain $`\widehat{q}\widehat{q}_i0`$. This proves the theorem in the special case where $`\widehat{\mu }=\widehat{\mu }_i`$ for $`i`$ large.
Now we suppose that $`\widehat{q}Q_{\widehat{\mu }}^{\mathrm{}}`$, $`\widehat{q}_iQ_{\widehat{\mu }_i}^{\mathrm{}}`$ with $`\widehat{\mu }_i\widehat{\mu }`$. Let $`\widehat{C}_i`$ be the set of points where the foliations $`\widehat{𝖿}_i^h`$, $`\widehat{𝖿}_i^v`$ are uniformly $`ϵ`$-close to $`\widehat{𝖿}^h`$, $`\widehat{𝖿}^v`$ in the $`\widehat{q}`$-metric. Then for $`i`$ large, in $`\widehat{C}_i`$ the complex structures defined by $`\widehat{q}_i,\widehat{q}`$ are uniformally nearly conformal. On the other hand, in $`\widehat{D}_i=\widehat{\mathrm{\Sigma }}\widehat{C}_i`$ they are not, but the $`|\widehat{q}|`$-volume of this set limits to zero. Therefore, if $`\widehat{p}_iQ_{\widehat{\mu }}^{\mathrm{}}`$ generates the Teichmüller geodesic connecting $`\widehat{\mu }`$ to $`\widehat{\mu }_i`$ in time 1, it follows that $`\widehat{p}_i0`$ so that $`\widehat{\mu }_i`$ converges to $`\widehat{\mu }`$ in the $`L^1`$ path metric. Moreover, the induced flow of quadratic differentials takes $`\widehat{q}`$ $`L^1`$ close to $`\widehat{q}_i`$ so that $`\widehat{q}_i`$ converges to $`\widehat{q}`$ in the $`L^1`$ topology on $`𝒬^{\mathrm{}}`$. ∎
## 7. The $`L^1`$ Ehrenpreis Conjecture
The $`L^1`$ EC is the following statement:
###### $`L^1`$ Ehrenpreis Conjecture.
The mapping class group $`𝖬(\widehat{\mathrm{\Sigma }},\mathrm{})`$ acts with $`L^1`$ dense orbits on $`𝒯(\widehat{\mathrm{\Sigma }})`$.
Since $`𝒯^{\mathrm{}}`$ is dense in $`𝒯(\widehat{\mathrm{\Sigma }})`$ and $`𝖬(\widehat{\mathrm{\Sigma }},\mathrm{})`$ stabilizes $`𝒯^{\mathrm{}}`$ , it will be enough to demonstrate that $`𝖬(\widehat{\mathrm{\Sigma }},\mathrm{})`$ acts with $`L^1`$ dense orbits on $`𝒯^{\mathrm{}}`$.
In it is shown that for a closed surface $`Z`$, a Teichmüller dense subset of pairs $`\mu ,\nu 𝒯(Z)`$ lie on the axes of pseudo Anosov homeomorphisms. By its definition as an isometric direct limit, $`𝒯^{\mathrm{}}`$ enjoys the same property. Fix a pair $`\widehat{\mu },\widehat{\nu }𝒯^{\mathrm{}}`$; without loss of generality, we may then assume that $`\widehat{\mu },\widehat{\nu }`$ lie on the axis $`A`$ of a pseudo Anosov diffeomorphism $`\widehat{\mathrm{\Phi }}`$ which is the lift of a pseudo Anosov $`\mathrm{\Phi }:ZZ`$, for some surface $`Z`$ occurring in the defining system of $`\widehat{\mathrm{\Sigma }}`$.
By an $`L^1`$ $`\mathrm{n}`$th root of $`\widehat{\mathrm{\Phi }}`$ we mean a sequence $`\{\widehat{\mathrm{\Psi }}_m\}`$ of pseudo Anosov homeomorphisms for which
1. If $`\lambda _m`$ is the entropy of $`\widehat{\mathrm{\Psi }}_m`$ then $`lim\lambda _m=\lambda ^{1/n}`$.
2. If $`A_m`$ is the axis of $`\widehat{\mathrm{\Psi }}_m`$ then $`A_mA`$ converges in the Hausdorff topology induced from the $`L^1`$ metric.
###### Theorem 5.
If for every pseudo Anosov $`\widehat{\mathrm{\Phi }}`$ and every $`n>0`$, $`\widehat{\mathrm{\Phi }}`$ has an $`L^1`$ $`n`$th root, then the $`L^1`$ EC is true.
###### Proof.
Suppose that $`\widehat{\mu },\widehat{\nu }`$ lie on the axis $`A`$ of $`\widehat{\mathrm{\Phi }}`$ and let $`\{\widehat{\mathrm{\Psi }}_m\}`$ be an $`L^1`$ $`n`$th root, $`n`$ large. Since the $`A_mA`$ in the $`L^1`$ Hausdorff topology, there exists $`\widehat{\mu }^{},\widehat{\nu }^{}`$ lying on some axis $`A_m`$ with $`d_{L^1}(\widehat{\mu },\widehat{\mu }^{})<ϵ`$, $`d_{L^1}(\widehat{\nu },\widehat{\nu }^{})<ϵ`$. On the axis $`A_m`$, we may move via a power of $`\widehat{\mathrm{\Psi }}_m`$ $`\widehat{\mu }^{}`$ close to $`\widehat{\nu ^{}}`$, which implies that $`\widehat{\mu }`$ is moved close to $`\widehat{\nu }`$ by the same power of $`\widehat{\mathrm{\Psi }}_m`$ as well. ∎
###### Note 1.
The existence of Teichmüller roots for all $`\widehat{\mathrm{\Phi }}`$ and all $`n`$ implies the classical $`\mathrm{𝖤𝖢}`$.
## 8. Directional Density
A family $`𝒫`$ of pseudo Anosov homeomorphisms is said to be directionally dense (in $`𝒬^{\mathrm{}}`$) if the set of quadratic differentials tangent to axes of elements of $`𝒫`$ is Teichmüller dense in $`𝒬^{\mathrm{}}`$. By , the family of all pseudo Anosov maps is directionally dense. In fact, it follows easily from the arguments in that the family of lifts of pseudo Anosovs $`\mathrm{\Phi }`$ of the type $`\mathrm{\Phi }=G^{2N}F^{2N}`$, where $`F`$, $`G`$ are right Dehn twists about simple closed curves $`c,d`$ that fill $`Z`$, where $`Z`$ ranges over all surfaces in the defining system of $`\widehat{\mathrm{\Sigma }}`$, is directionally dense. By Corollary 2.6 in , the subfamily obtained by demanding that $`c`$ is nonseparating is also directionally dense.
Now given a nonseparating simple closed curve $`\gamma Z`$, let $`\rho _\gamma :Z_\gamma Z`$ be the degree 2 cover obtained by cutting two copies of $`Z`$ along $`\gamma `$ and gluing ends. We say that a pair $`(c,d)`$ of filling, simple closed curves is interlacing if there exists a pair of nonseparating simple closed curves $`\alpha ,\beta `$ such that $`\rho _\alpha ^1(c)`$ is connected whereas $`\rho _\alpha ^1(d)`$ is not and $`\rho _\beta ^1(d)`$ is connected whereas $`\rho _\beta ^1(c)`$ is not. See Figure 4.
Let $`𝒫`$ be the family of pseudo Anosov homeomorphisms of $`\widehat{\mathrm{\Sigma }}`$ which are lifts of pseudo Anosovs of the form
$$\mathrm{\Phi }=G^{2N}F^{2N}:ZZ,$$
where
1. $`F`$, $`G`$ are right Dehn twists about $`c,dZ`$.
2. $`c`$ is nonseparating and $`(c,d)`$ is an interlacing pair.
3. $`Z`$ ranges over all surfaces in the defining system of $`\widehat{\mathrm{\Sigma }}`$.
###### Lemma 2.
$`𝒫`$ is directionally dense.
###### Proof.
It is enough to show that for a fixed surface $`Z`$, the family of maps satisfying (1) and (2) is directionally dense in $`𝒬(Z)`$. Assume that $`c`$ and $`d`$ are filling, generating the pseudo Anosov homeomorphism $`\mathrm{\Phi }=G^{2N}F^{2N}`$. If $`(c,d)`$ is not an interlacing pair, there exists a simple closed curve $`\delta `$ for which the pair $`(c,\delta )`$ is interlacing, though not necessarily filling. Indeed, one may assume after a homeomorphism that $`c`$ is the curve appearing in Figure 5; then taking $`\delta ,\alpha ,\beta `$ as indicated there, $`(c,\delta )`$ is interlacing with respect to the pair $`(\alpha `$, $`\beta )`$. Now for $`j`$ large, $`\delta _j=G^j(\delta )`$ is close to $`d`$, hence $`(c,\delta _j)`$ is eventually filling. If $`j`$ is in addition even, $`G^j`$ lifts to the total space of any degree 2 cover of $`\mathrm{\Sigma }`$, thus the pair $`(c,\delta _j)`$ is interlacable with respect to the same curves interlacing $`(c,\delta )`$. For $`j`$ large, the pseudo Anosov $`\mathrm{\Phi }_j=G_{\delta _j}^{2N}F^{2N}`$ has axis close to that of $`\mathrm{\Phi }`$, and since the maps of the form $`\mathrm{\Phi }`$ are already directionally dense, we are done. ∎
## 9. Necklace Roots
Let $`n`$ and let $`\widehat{\mathrm{\Phi }}𝒫`$, the family appearing in Lemma 2, so that in particular $`\widehat{\mathrm{\Phi }}`$ is the lift of a pseudo Anosov of the form $`\mathrm{\Phi }=G^{2N}F^{2N}:ZZ`$. Let $`\rho _{mn}:Z_{mn}Z`$ be the cover obtained by cutting $`2mn`$ copies of $`Z`$ along a pair $`\alpha `$ and $`\beta `$ interlacing $`c,d`$ and gluing in a circular fashion. We call $`\rho _{mn}`$ the necklace cover associated to $`(c,d)`$. In Figure 6, we illustrate the construction of the necklace $`Z_{mn}`$ and the formation of the lifts of the curve $`c`$. In Figure 7 we display the finished necklace.
There are $`mn`$ lifts $`c_1,\mathrm{},c_{mn}`$ and $`d_1,\mathrm{},d_{mn}`$ of each of $`c`$ and $`d`$, each mapping with degree two onto their ancestor. On $`Z_{mn}`$, $`\mathrm{\Phi }`$ lifts to
$$\stackrel{~}{\mathrm{\Phi }}=G_{mn}^N\mathrm{}G_1^NF_{mn}^N\mathrm{}F_1^N$$
where $`F_i`$, $`G_i`$ is the right Dehn twist about $`c_i`$, $`d_i`$. Let $`\chi `$ denote the clockwise rotation of $`Z_{mn}`$ by an angle of $`2\pi /n`$, so that the pair $`c_i`$, $`d_i`$ is taken to $`c_{j+m}`$, $`d_{j+m}`$ (indices taken mod $`mn`$). We define the necklace $`\mathrm{n}`$th root to be the sequence of lifts of pseudo Anosovs $`\{\widehat{\sqrt[n]{\mathrm{\Phi }}_m}\}`$ to $`\widehat{\mathrm{\Sigma }}`$ where
$$\sqrt[n]{\mathrm{\Phi }}_m=\chi G_m^N\mathrm{}G_1^NF_m^N\mathrm{}F_1^N,$$
$`m=2,3,\mathrm{}`$. The necklace $`n`$th root is the basic contruction used in the formation of $`L^1`$ roots. The construction of $`\sqrt[n]{\mathrm{\Phi }}_m`$ is a generalization of one that first appeared in , where branched covers were used.
###### Lemma 3.
$`\sqrt[n]{\mathrm{\Phi }}_m`$ is pseudo Anosov for all $`m`$.
###### Proof.
For $`i=1,\mathrm{},n`$, let
$$T_i=G_{im}^N\mathrm{}G_{(i1)m+1}^NF_{im}^N\mathrm{}F_{(i1)m+1}^N.$$
Then it is easy to see that
$$(\sqrt[n]{\mathrm{\Phi }}_m)^n=T_2\mathrm{}T_nT_1,$$
which is of Thurston-Penner type, hence $`(\sqrt[n]{\mathrm{\Phi }}_m)^n`$ is pseudo Anosov, implying $`\sqrt[n]{\mathrm{\Phi }}_m`$ is pseudo Anosov as well. ∎
## 10. Existence of $`L^1`$ Roots
Denote by $`𝒞_Z(\widehat{\mathrm{\Sigma }})`$ the family of simple closed solenoids which are lifts of simple closed curves on $`Z`$. We begin by constructing a family $`\{\widehat{\mathrm{\Psi }}_m\}`$ whose stable and unstable laminations intersection converge to those of $`\widehat{\mathrm{\Phi }}`$ with respect to test solenoids in $`𝒞_Z(\widehat{\mathrm{\Sigma }})`$.
For each $`m=2,3,\mathrm{}`$ let $`\sqrt[mn]{\mathrm{\Phi }}_m`$ denote the $`mth`$ element in the sequence of pseudo Anosovs whose lifts define the $`mn`$th necklace root of $`\widehat{\mathrm{\Phi }}`$. Define the sequence $`\{\widehat{\mathrm{\Psi }}_m\}`$ as the lifts to $`\widehat{\mathrm{\Sigma }}`$ of the pseudo Anosov homeomorphisms
$$\mathrm{\Psi }_m=\left(\sqrt[nm]{\mathrm{\Phi }}_m\right)^m.$$
Observe that the stable and unstable foliations of $`\mathrm{\Psi }_m`$ and $`\sqrt[nm]{\mathrm{\Phi }}_m`$ are equal.
###### Note 2.
$`\widehat{\mathrm{\Psi }}_m`$ is not the same as $`\widehat{\sqrt[n]{\mathrm{\Phi }}_m}`$. In fact, if we lift $`\sqrt[n]{\mathrm{\Phi }}_m`$ to $`Z_{m^2n}`$ – where $`\mathrm{\Psi }_m`$ is defined – we see that this lift twists along $`m`$ disjoint “blocks” of curves, each block consisting of a succession of $`m`$ lifts of $`c`$ and $`d`$. On the other hand, $`\mathrm{\Psi }_m`$ consists of twists along one block consisting of a succession of $`m^2`$ lifts of $`c`$ and $`d`$. As we shall see, the stable and unstable laminations of the family $`\{\widehat{\mathrm{\Psi }}_m\}`$ have better convergence properties than those of $`\{\widehat{\sqrt[n]{\mathrm{\Phi }}_m}\}`$.
Denote by $`\widehat{𝔣}_m^u,\widehat{𝔣}_m^s`$ and by $`\widehat{𝔣}^u,\widehat{𝔣}^s`$ the unstable and stable laminations of $`\widehat{\mathrm{\Psi }}_m`$ and $`\widehat{\mathrm{\Phi }}`$.
###### Lemma 4.
For all $`\widehat{c}𝒞_Z(\widehat{\mathrm{\Sigma }})`$,
$$\text{I}(\widehat{𝔣}_m^u,\widehat{c})\text{I}(\widehat{𝔣}^u,\widehat{c})\text{ and }\text{I}(\widehat{𝔣}_m^s,\widehat{c})\text{I}(\widehat{𝔣}^s,\widehat{c}).$$
###### Proof.
The proof will be through examination of curve matrices. We begin with $`\mathrm{\Phi }`$. Let $`r=\text{I}(c,d)`$. The action of $`\mathrm{\Phi }`$ along the curves $`c,d`$ is given by the matrix
$$M_\mathrm{\Phi }=\left(\begin{array}{cc}1& 2rN\\ 2rN& (2rN)^2+1\end{array}\right).$$
Let $`Z_{m^2n}`$ be the surface where $`\mathrm{\Psi }_m`$ is defined. The curve families $`𝒞=\{c_1,\mathrm{},c_{m^2n}\}`$, $`𝒟=\{d_1,\mathrm{}d_{m^2n}\}`$ are filling, and the action of $`\stackrel{~}{\mathrm{\Phi }}`$ on $`𝒞𝒟`$ is prescribed schematically by the matrix in Figure 8, where all entries not contained in the boxed vectors are zero, and where the “broken” vectors indicate that only that portion of the corresponding vector is used. For example, in the upper right hand corner we have the entry $`rN`$, which is the bottom half of the $`B`$-vector; in the lower right hand corner, we have the vector entry $`\left((rN)^2\mathrm{\hspace{0.33em}\hspace{0.33em}2}(rN)^2+1\right)^𝖳`$, which is the top two thirds of the vector $`D`$, and so on. The curve matrix of $`(\mathrm{\Psi }_m)^n=(\sqrt[mn]{\mathrm{\Phi }}_m)^{mn}`$ is displayed in Figure 9. The black vectors indicate regions where $`M_{(\mathrm{\Psi }_m)^n}=(M_{\mathrm{\Psi }_m})^n=(M_{\sqrt[mn]{\mathrm{\Phi }}_m})^{mn}`$ differs from $`\stackrel{~}{M}_\mathrm{\Phi }`$. Figure 10 contains the curve matrix of $`\sqrt[mn]{\mathrm{\Phi }}_m`$.
Denote by $`\lambda _m`$ the Perron root of $`M_{\mathrm{\Psi }_m}`$ and by
$$\upsilon _m=(a_1\mathrm{},a_{m^2n},b_1,\mathrm{},b_{m^2n})$$
the corresponding Perron vector, normalized to have $`L^1`$ norm 1 i.e. so that $`\upsilon _m`$ is a probability vector. Let $`\upsilon _m^{\mathrm{avg}}`$ be the vector
$$\upsilon _m^{\mathrm{avg}}=(a_1+\mathrm{}+a_{m^2n},b_1+\mathrm{}+b_{m^2n}),$$
which is also a probability vector. In the case of $`\mathrm{\Phi }`$ and $`\stackrel{~}{\mathrm{\Phi }}`$, the Perron roots are identical and will be denoted $`\lambda `$; if $`\upsilon =(a,b)^𝖳`$ is the $`L^1`$ norm 1 Perron vector of $`M_\mathrm{\Phi }`$, $`\stackrel{~}{\upsilon }=(1/2m^2n)(a,\mathrm{},a,b,\mathrm{},b)^𝖳`$ is the $`L^1`$ norm 1 Perron vector for $`M_{\stackrel{~}{\mathrm{\Phi }}}`$. We recall that this spectral data has the following interpretation:
1. The Perron roots $`\lambda _m`$, $`\lambda `$ are equal to the entropies of $`\mathrm{\Psi }_m`$, $`\mathrm{\Phi }`$.
2. Let $`\tau _m`$ be the bigon track formed from the curve families $`𝒞`$, $`𝒟`$. The measures $`\mu _m`$, $`\stackrel{~}{\mu }`$ formed from the Perron vectors $`\upsilon _m`$, $`\stackrel{~}{\upsilon }`$, parametrize the unstable laminations $`𝔣_m^u`$, $`\stackrel{~}{𝔣}^u`$ in the cone $`𝖢(\tau _m)`$.
Note that the column sums of $`M_{(\mathrm{\Psi }_m)^n}`$ have uniform upper and lower bounds $`B`$ and $`b>1`$. We thus obtain the bound
$$1<b<(\lambda _m)^n<B.$$
We may then assume, after passing to a subsequence if necessary, that the $`\lambda _m`$ converge to some value $`\lambda ^{}>1`$. We shall need to control the following entries of $`\upsilon _m`$:
###### Claim 1.
$`a_m,a_{m+1},b_m,b_{m+1}0`$ as $`m\mathrm{}`$.
Proof of claim. Let $`\xi _m`$ be the Perron root of $`M_{\sqrt[nm]{\mathrm{\Phi }}_m}`$: thus $`(\xi _m)^m=\lambda _m`$. If one of the four entries listed in the statement does not converge to 0, then consideration of the matrix $`M_{(\mathrm{\Psi }_m)^n}`$ shows that none of them do. Thus, let us suppose that $`a_m\to ̸0`$. It follows that eventually $`a_{2m}\delta `$ for some positive $`\delta `$. Examination of the action of $`M_{\sqrt[nm]{\mathrm{\Phi }}_m}`$ on $`\upsilon _m`$ shows that $`a_{m+im}=\xi _m^{i1}a_{2m}`$ for $`i=2,\mathrm{}mn`$. However $`\xi _m^{i1}>1`$ for all $`i`$, and since $`\upsilon _m`$ is a probability vector, this would imply that $`mn\delta <mna_{2m}<1`$, impossible since $`m\mathrm{}`$. This proves the claim.
Let us shorten notation, writing $`\widehat{𝔣}=\widehat{𝔣}^u`$ and $`\widehat{𝔣}_m=\widehat{𝔣}_m^u`$ for the unstable foliations of $`\widehat{\mathrm{\Phi }}`$ and $`\widehat{\mathrm{\Psi }}_m`$. Let $`\widehat{c}`$ be any simple closed solenoid in $`𝒞_Z(\widehat{\mathrm{\Sigma }})`$.
Then
$$\text{I}((\lambda _m)^n\widehat{𝔣}_m,\widehat{c})=\text{I}((\mathrm{\Psi }_m)^n\widehat{𝔣}_m,\widehat{c})=_{\widehat{E}}\widehat{M_{(\mathrm{\Psi }_m)^n}\upsilon _m}|\widehat{\tau }\widehat{c}|,$$
where $`\widehat{M_{(\mathrm{\Psi }_m)^n}\upsilon _m}`$ is the lift of the vector $`M_{(\mathrm{\Psi }_m)^n}\upsilon _m`$ to $`\widehat{E}`$. Now since $`\widehat{c}`$ is the lift of a simple closed curve $`c`$ in $`Z`$, we have that
$$\text{I}(\widehat{𝔣}_m,\widehat{c})=\frac{1}{\mathrm{deg}(Z)}\text{I}(𝔣_m^{\mathrm{avg}},c),$$
where $`𝔣_m^{\mathrm{avg}}`$ is the measured lamination in $`Z`$ corresponding to the weight $`\upsilon _m^{\mathrm{avg}}`$. However an examination of the matrices $`M_{\stackrel{~}{\mathrm{\Phi }}}`$ and $`M_{(\mathrm{\Psi }_m)^n}`$ yields
$$(\lambda _m)^n\upsilon _m^{\mathrm{avg}}=M_\mathrm{\Phi }\upsilon _m^{\mathrm{avg}}+ϵ_m$$
where
$$ϵ_m=((rN)^2(a_m+a_{m+1}),(rN)^2b_{m+1}+(rN)^3a_m+((rN)^2(rN)+1)a_{m+1}+(rN)^4b_m)^T.$$
By Claim 1, it follows that $`ϵ_m(0,0)^T`$ or $`\upsilon _m^{\mathrm{avg}}`$ converges to an eigenvector of $`M_\mathrm{\Phi }`$ of eigenvalue $`\lambda _{}^n`$. But since $`\lambda _{}^n>1`$ and the second eigenvalue of $`M_\mathrm{\Phi }`$ is strictly less than 1, we must have that $`\upsilon _m^{\mathrm{avg}}\upsilon `$ and $`\lambda _{}=\lambda ^{1/n}`$. In particular,
$$\underset{m\mathrm{}}{lim}\text{I}(𝔣_m^{\mathrm{avg}},c)=\text{I}(𝔣,c).$$
This takes care of the unstable part of the theorem; the stable part is handled by repeating the above argument for $`\mathrm{\Phi }^1`$ and $`\mathrm{\Psi }_m^1`$. ∎
###### Theorem 6.
Every pseudo Anosov $`\widehat{\mathrm{\Phi }}`$ has an $`L^1`$ $`n`$th root for all $`n`$.
###### Proof.
Since $`𝒫`$ is directionally dense, there exists a sequence $`\{\widehat{\mathrm{\Phi }}^{(g)}\}𝒫`$, where $`\widehat{\mathrm{\Phi }}^{(g)}`$ is the lift of a pseudo Anosov $`\mathrm{\Phi }^{(g)}:X_gX_g`$ in which $`X_g`$ is a surface of genus $`g\mathrm{}`$, and the axes $`A_gA`$ = axis of $`\widehat{\mathrm{\Phi }}`$. For each $`g`$, let $`\{\widehat{\mathrm{\Psi }}_m^{(g)}\}`$ be the sequence of pseudo Anosovs constructed above. Then by Lemma 4 we may obtain an $`L^1`$ $`n`$th root $`\{\widehat{\mathrm{\Psi }}_m\}`$ of $`\widehat{\mathrm{\Phi }}`$ by extracting a suitable diagonal subsequence of $`\{\widehat{\mathrm{\Psi }}_m^{(g)}\}`$. Indeed, a suitable diagonal subsequence $`\{\widehat{\mathrm{\Psi }}_m\}`$ yields a sequence of pseudo Anosov homeomorphisms whose stable and unstable laminations intersection converge to those of $`\widehat{\mathrm{\Phi }}`$. By Theorems 3 and 4, this gives rise to a sequence of quadratic differentials $`\widehat{q}_i`$ along the associated axes $`A_i`$ which $`L^1`$-converge to the quadratic differential $`\widehat{q}`$ determined by $`\widehat{𝔣}^u`$ and $`\widehat{𝔣}^s`$. ∎
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# Storage and retrieval of light pulses in atomic media with “slow” and “fast” light
## Abstract
We present experimental evidence that light storage, i.e. the controlled release of a light pulse by an atomic sample dependent on the past presence of a writing pulse, is not restricted to small group velocity media but can also occur in a negative group velocity medium. A simple physical picture applicable to both cases and previous light storage experiments is discussed.
All-optical information processing requires the use of photons as fast and reliable carriers of information. Photons are quantum objects and the search for media where the quantum state of photons can be preserved and processed is of great significance. Recently, broad attention has been focussed on the possibility of “light storage” (LS) which is the preservation of the information carried by a light pulse for controllable later release. Such a possibility was suggested in theory Fleischhauer and Lukin (2000) and subsequent experimental results were presented in support of this suggestion Phillips et al. (2001); Liu et al. (2001); Zibrov et al. (2002); Mair et al. (2002). In all of these experiments, a weak light pulse was “written” into an atomic medium driven by a stronger field and, after a dark interval, retrieved from the medium by turning on the strong (drive) field.
All observations of LS were achieved under conditions of electromagnetically induced transparency (EIT) since, according to Fleischhauer and Lukin (2000), EIT and “slow light” (small group velocity associated with EIT) play a key role. The storage effect is seen as a consequence of the slowing and compression of the light pulse in the atomic medium, the propagation of a mixed light-matter excitation (dark-state polariton), the transformation of the dark-state polariton in the absence of light into a pure atomic spin excitation and finally the release of a light pulse once the drive field is turned on Lukin (2003).
The purpose of this letter is to present experimental results and theoretical considerations which broaden the scope of the subject by demonstrating that a LS effect, analogous to that previously reported Phillips et al. (2001); Liu et al. (2001); Zibrov et al. (2002); Mair et al. (2002), can take place in media where EIT does not occur and where the probe pulse group velocity is negative (“fast light”) as a result of large anomalous dispersion.
Steep anomalous dispersion exists in driven atomic media in connection with electromagnetically induced absorption (EIA) Akulshin et al. (1999). EIA occurs when resonant light interacts with a two-level atomic transition in which the Zeeman degeneracy of the excited level is higher than that of the lower level; namely $`F_e>F_g>0`$ where $`F_e`$ and $`F_g`$ are the total angular momenta of the excited level and ground level, respectively. A resonant increase in the probe absorption occurs under the condition of a Raman resonance with ground state Zeeman sublevels Akulshin et al. (1998). In particular, at zero magnetic field the EIA resonance condition is achieved for the two orthogonal polarization components of a single monochromatic optical field. Superluminar pulse propagation in an atomic vapor under the conditions of EIA has been demonstrated in Akulshin et al. (2003).
The experimental scheme used is very similar to the one presented in Ref. Phillips et al. (2001). LS was studied in a 5 cm long vapor cell containing a natural isotopic mixture of rubidium. The cell was placed at the center of a cylindrical $`\mu `$-metal shield. A solenoid inside the magnetic shield allows tuning of the longitudinal magnetic field $`B`$. The cell was heated ($`70^{}`$C) to produce almost $`100\%`$ linear absorption and $`5080\%`$ absorption at maximum light intensity. Two extended cavity diode lasers were used for the $`D_1`$ and $`D_2`$ transitions. Fast switching on and off of the laser light was achieved with an acousto-optic modulator (AOM). After the AOM, a polarizer fixes the polarization of the drive field. A Pockels cell along the light beam was used as an electro-optic modulator (EOM) to generate a probe pulse with orthogonal polarization relative to the drive field. The probe pulse was 5 times weaker than the drive field. Care was taken to use the EOM in the $`\lambda /2`$ configuration in order to produce in-phase probe pulses relative to the drive field. The light beam was expanded to a $`1`$ cm diameter before the cell. The maximum laser power available at the cell was $`0.6`$ mW and $`2`$ mW for the $`D_1`$ and $`D_2`$ lines, respectively. After traversing the atomic vapor, the drive and probe polarization components of the light were separated by a polarizing beam splitter and collected by fast photodiodes. The electronic control and detection response times were shorter than 1 $`\mu s`$.
Our setup allows the use of linear perpendicular or opposite circular drive and probe field polarization combinations. Qualitatively similar LS effects were observed for both choices. We describe in the following the signals obtained with perpendicular linear drive and probe fields polarizations. LS under conditions of EIT was observed on both the transitions $`5S_{1/2}\left(F=1\right)5P_{3/2}`$ ($`D_2`$ line) and $`5S_{1/2}\left(F=2\right)5P_{1/2}\left(F^{}=1\right)`$ ($`D_1`$ line) of <sup>87</sup>Rb. We present results obtained for the latter transition which was also used in previous experiments Phillips et al. (2001); Zibrov et al. (2002); Mair et al. (2002). The width of the EIT resonance, measured at maximum light power by varying the magnetic field, was 65 kHz. Trace a of Fig. 1, obtained with a Gaussian shape probe pulse, reproduces the essential features of the previous LS experiments Phillips et al. (2001); Liu et al. (2001); Zibrov et al. (2002); Mair et al. (2002). Since the inverse of the pulse duration is small compared to the EIT resonance width, the Gaussian pulse shape is well preserved during propagation. A pulse delay (relative to vacuum propagation) of approximately $`1`$ $`\mu `$s is observed corresponding to a group velocity $`v_g5\times 10^4`$ m/s. We subsequently used square probe pulses to simplify further comparison with numerical modelling. Trace b was obtained for the same transition as for a with a square probe pulse. The transmitted probe pulse presents a characteristic distortion, indicative of a slow light medium Akulshin et al. (2003). In both traces the retrieved pulse has an exponential decay with approximately the same decay time ($`2\mu `$s). An exponential decay time of $`8\pm 2\mu `$s was measured for the retrieved pulse amplitude as a function of the dark interval, in agreement with the transverse atomic time-of-flight estimate.
We now turn our attention to the $`D_2`$ line where EIA leads to fast light propagation. Trace c of Fig. 1 was obtained with the laser tuned to the $`5S_{1/2}\left(F=2\right)5P_{3/2}\left(F^{}=3\right)`$ transition of the $`D_2`$ line. We notice a characteristic probe pulse distortion with a leading edge less absorbed than the rest of the pulse. Such distortion is indicative of the pulse advancement expected in a fast light medium Akulshin et al. (2003). We observe that in the fast light medium the retrieved pulse has similar characteristics to the pulse retrieved from the slow light medium. In either case, in the absence of a magnetic field, the retrieved pulse has an exponentially decaying slope with a time constant dependent on the drive field intensity. In the presence of a nonzero magnetic field similar damped oscillations are seen on the retrieved pulse envelope for the two types of transitions. The same dependence of the retrieved pulse amplitude on the dark interval duration is observed in the two cases.
To model the observed results, we have solved the Bloch equations for an homogeneous ensemble of atoms with two energy levels with Zeeman degeneracy Valente et al. (2002). The excited state decays spontaneously to the ground state at a rate $`\mathrm{\Gamma }`$ and, in order to mimic time-of-flight relaxation, all states in the system decay at rate $`\gamma `$ while “fresh atoms” are isotropically injected in the ground levels at the same rate. The atoms interact with an optical field whose polarization is decomposed into two orthogonal (linear or circular) components with arbitrary amplitude and phase. The light propagation in the atomic sample is considered to the lowest order in the sample optical length Walser and Zoller (1994). Consequently, it is assumed that all atoms in the sample see the same (undepleted) incident field. Under such an assumption, the transmitted field is $`\stackrel{}{E}_T\stackrel{}{E}_0+i\alpha \stackrel{}{P}`$ where $`\stackrel{}{E}_0`$ is the total incident field, $`\stackrel{}{P}=Tr\left(\rho \stackrel{}{D}\right)`$is the atomic optical polarization per atom, $`\stackrel{}{D}`$ is the electric dipole operator and $`\alpha `$ is a real constant proportional to the atomic density and the sample length that is adjusted to the observed absorption. The transmitted signal can be computed for either polarization component. The pulse sequence used for the calculation is shown in the upper traces of Fig. 2. Trace a corresponds to the calculated transmission (at the probe field polarization) for parameters corresponding to the transition $`5S_{1/2}\left(F=2\right)5P_{1/2}\left(F^{}=1\right)`$ of <sup>87</sup>Rb. Trace b was calculated for the transition $`5S_{1/2}\left(F=2\right)5P_{3/2}\left(F^{}=3\right)`$ of <sup>87</sup>Rb. The pulse distortions observed in the experiment are well reproduced and a retrieved pulse is obtained for both transitions after the dark interval. The amplitude of the retrieved pulse depends on the probe pulse amplitude and duration and decays exponentially with a decay time depending on the drive field intensity alone. Similar results to those presented in Fig. 2 are obtained if orthogonal circular polarization are considered for both fields.
The above results can be explained on the basis of a unique physical picture. Consider the basic light-atom interaction process presented in Fig. 3a where an optical field couples the atomic ground-state $`C`$ (coupled state) to the excited state $`e`$ while the second ground-state $`D`$ (dark state) is unaffected by the field. After excitation in state $`e`$ the atom can decay spontaneously to both lower states. This is an optical pumping process Happer (1972). If the atoms interact with the light for a sufficiently long time and if states $`C`$ and $`D`$ are stationary, the atoms will eventually end in state $`D`$ and the medium will become transparent. During this process, transitions from state $`e`$ to state $`D`$ occur spontaneously but may also take place through stimulated Raman emission. The condition for stimulated Raman emission is the existence of nonzero coherence between states $`C`$ and $`D`$, i.e. $`\rho _{DC}0`$ where $`\rho _{DC}`$ is the element of the density matrix between states $`D`$ and $`C`$. Under such a condition the interaction with the applied field results in optical coherence between states $`e`$ and $`D`$ ($`\rho _{eD}0`$) which in turn radiates a field for this transition which is coherent with the initial field. As a general rule, *a light field interacting with a coherently prepared atom in the ground state will induce coherent emission of light in all allowed optical transitions as long as the (ground state) atomic coherence is preserved* (Fig. 3b).
Direct generalization of the situation represented in Fig 3a to the more complex case of a two-level system with arbitrary Zeeman degeneracy interacting with an optical field is possible. In general, given the field polarization, after a suitable basis transformation, one can identify several ($`C`$) states in the ground level, each coupled to a corresponding excited state sublevel, and two, one or zero dark ($`D`$) states (Fig. 3c). The last case occurs when the Zeeman degeneracy of the excited level is higher than that of the ground level Smirnov et al. (1989); Taichenachev et al. (2004). In all cases the general rule stated above applies.
Two types of systems have commonly been considered for the experimental study of coherent light-atom interaction dynamics. In the first class one has the Hanle type experiments Kastler (1973); Hannaford (1997); Dancheva et al. (2000); Renzoni et al. (2001); Failache et al. (2003) where a single optical field with well defined polarization interacts with two atomic levels with Zeeman sublevels whose energy is tuned with a static magnetic field. If the light polarization is linear it is straightforward that a scheme similar to that of Fig. 3c results by choosing the quantization axis along the direction of the light polarization. Even in the general case of arbitrary elliptical light polarization the scheme of Fig. 3c can be applied after a suitable transformation of the state basis that depends on the incident light polarization Taichenachev et al. (2004). An immediate consequence of this description is that in the presence of ground state coherence the system will respond by the coherent emission of light with a polarization orthogonal to that of the incident field. Such a process was recently discussed in terms of the polarization change that the light experiences while propagating in an atomic medium with light induced anisotropy and an alternative explanation of the LS effect in such terms was proposed Kozlov et al. (2004).
A second class of experiments for which the picture presented in Fig. 3a-b can be conveniently applied, concerns a system of an excited level and two ground levels with energy separation $`\mathrm{\Delta }`$ driven by two optical fields of frequency $`\omega _1`$ and $`\omega _2`$. This system (hereafter designated as the $`\mathrm{\Lambda }`$ system), has been extensively studied theoretically and experimentally in connection with EIT. Since in a $`\mathrm{\Lambda }`$ system each field interacts with a different transition, it is possible after a time dependent unitary transformation to describe the atomic dynamics by a time-independent Hamiltonian and to identify a dark state $`D`$ and a coupled state $`C`$ Scully and Zubairy (1997). This “coupled-uncoupled” description exactly corresponds to the situation of Fig. 3a-b. Both levels $`C`$ and $`D`$ are linear combinations of the two lower levels that explicitly depend on the amplitudes and phases of the applied fields. Here again, if coherence is present between levels $`C`$ and $`D`$, the atomic medium will react to the optical excitation by the coherent emission of fields at frequencies $`\omega _1`$ and $`\omega _2`$. These fields are “orthogonal” to the applied fields in the sense that they *together* only couple to state $`D`$ (while the incident fields only couple to $`C`$). As an example we consider the situation corresponding to the retrieval process in the LS experiment reported in Liu et al. (2001). After the dark interval, the drive field ($`\omega _1`$) alone excites the $`\mathrm{\Lambda }`$ system. Since (previously created) coherence is present between the ground levels (that coincide with $`C`$ and $`D`$ in this case), a field of frequency $`\omega _2=\omega _1\mathrm{\Delta }`$ is emitted by the medium. In the general case, in the presence of coherence between $`C`$ and $`D`$, two fields will be emitted with frequencies $`\omega _1`$ and $`\omega _2`$ that will interfere with the incident fields.
The previous discussion leads to a very general qualitative understanding of the transient behavior of coherently driven atomic systems. For given excitation conditions, after a long enough time, the system will reach a steady state. If one or more dark states exists the steady state of the system will be the dark state(s) $`D`$. If no dark state exists the steady state will generally be a statistical mixture (diagonal density matrix) of states $`C`$. A rapid modification of the excitation conditions, for instance a light polarization change in a Hanle experiment or a change in the amplitude or phase of the fields in a $`\mathrm{\Lambda }`$ scheme, will determine a new set ($`C^{}`$, $`D^{}`$) of coupled and dark states. The change in the excitation conditions needs to be rapid with respect to the optical pumping time but may otherwise be slow with respect to other characteristic times. Since the previous state of the system is generally a linear combination of the new $`C^{}`$, $`D^{}`$ states, coherence among the latter states exists and results in the transient emission of an “orthogonal” field. Such emission will last as long as the coherence between $`C^{}`$ and $`D^{}`$ survives. The decay of the transient emission will be purely exponential with a time constant given by the optical pumping time between states $`C^{}`$ and $`D^{}`$ if those states are stationary. If they are not stationary (nonzero magnetic field for a Hanle experiment or nonzero Raman detuning in a $`\mathrm{\Lambda }`$ system), the decay will show damped oscillations Valente et al. (2002). If the two different excitation conditions are separated in time by a dark interval then the transient emission of an “orthogonal” field will occur after the dark interval provided that this interval is not long compared to the ground state coherence lifetime. This mechanism applies to the LS experiments previously reported Phillips et al. (2001); Liu et al. (2001); Mair et al. (2002); Zibrov et al. (2002).
The experimental results and the discussion above demonstrate that EIT and propagation in a “slow light” medium is not an essential requirement for the storage and retrieval of a light pulse in an atomic medium. One point in the simple description presented here, which departs from previous theoretical treatments Fleischhauer and Lukin (2000), is the simple consideration of the field propagation in the atomic medium where the spatial variation of the incident field along the sample is essentially neglected. This crude approximation is nevertheless appropriate under conditions where the duration of the probe pulse is comparable or longer than the light propagation time. Such a situation occurred in most LS experiments reported so far (including ours) with the exception of Liu et al. (2001). In the picture presented here, the light retrieval transient appears as a consequence of the irreversible relaxation of the system towards a new steady state. In this context, only exponential decaying pulses can be obtained preventing the retrieval of information about the state of the initial probe pulse other than its presence (one bit information). The storage and recovery of information on the state of the incoming probe pulse are beyond the scope of this simple picture and possibly of most experimental conditions achieved to date.
In conclusion, we have achieved storage and retrieval of light pulses in slow-light and fast-light media and given a unified description of both cases. The results presented in this letter suggest that further theoretical and experimental work is needed for the understanding and practical realization of information preserving storage in an atomic medium. A necessary first step in this direction would be the realization of an atomic medium in which more than one probe pulse (more than one bit of information) can be contained in the atomic sample and propagate without significant distortion Boyd et al. (2005); Matsko et al. (2005).
This work was supported by a Swinburne University RDGS grant and Fondo Clemente Estable (Uruguay).
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# Gravitational Waves from Compact Sources 11footnote 1Proceedings of the 5th International Workshop “New Worlds in Astroparticle Physics”
## 1 Introduction
The new generation of gravitational wave (GW) detectors is already collecting data, improving the achieved sensitivity by at least one order of magnitude compared to operating resonant bar detectors. Broadband GW detectors are sensitive to frequencies between 50 and a few hundred Hz. In their advanced stage, the current GW detectors will have a broader bandwidth but will still not be sensitive enough to frequencies above 500 to 600Hz. Nevertheless, improved sensitivity can be achieved at high frequencies through narrow-band operation . In addition, there are proposals for constructing wide band resonant detectors in the kHz band .
In this short review we will discuss some of the sources that are in the high frequency band ($`500600`$Hz), where the currently operating interferometers are sensitive enough only if they are narrowbanded. Since there exist a variety of GW sources with very interesting physics to be explored, this high-frequency window deserves special attention. If either resonant or narrow-band interferometers achieve the required sensitivity, a plethora of unique information can be deduced from detected signals.
## 2 Gravitational collapse
Core collapse. One of the most spectacular astrophysical events is the core collapse of massive stars, leading to the formation of a neutron star (NS) or a black hole (BH). The outcome of core collapse depends sensitively on several factors: mass, angular momentum and metallicity of progenitor, existence of a binary companion, high-density equation of state, neutrino emission, magnetic fields, etc. Partial understanding of each of the above factors is emerging, but a complete and consistent theory for core collapse is still years away.
Roughly speaking, isolated stars more massive than $`810M_{}`$ end in core collapse and $`90\%`$ of them are stars with masses $`820M_{}`$. After core bounce, most of the material is ejected and if the progenitor star has a mass $`M20M_{}`$ a neutron star is left behind. On the other hand, if $`M20M_{}`$ fall-back accretion increases the mass of the formed proton-neutron star (PNS), pushing it above the maximum mass limit, which results in the formation of a black hole. Furthermore, if the progenitor star has a mass of roughly $`M45M_{}`$, no supernova explosion is launched and the star collapses directly to a BH.
The above picture is, of course, greatly simplified. In reality, the metallicity of the progenitor, the angular momentum of the pre-collapse core and the presence of a binary companion will decisively influence the outcome of core collapse. Very massive stars lose mass through strong stellar winds. The mass loss rate is sensitive to the metallicity of the star and can be very high, allowing a 60 $`M_{}`$ star to leave behind a typical 1.4 $`M_{}`$ neutron star, instead of a much more massive black hole. Since mass-loss is a complex phenomenon, current observations are still insufficient to constrain the different possible outcomes of core collapse for stars with $`M45M_{}`$. Through mass transfer or common envelope episodes, a binary companion can cause a massive star to enter the rapid mass-losing (Wolf-Rayet) phase earlier, making wind mass loss more effective. Roughly half of all stars are in suffiently close binaries that binary interactions must be taken into account when trying to predict to outcome of core collapse. Rotation influences the collapse by changing dramatically the properties of the convective region above the proto-neutron star core. Centrifugal forces slow down infalling material in the equatorial region compared to materiall falling in along the polar axis, yielding a weaker bounce. This asymmetry between equator and poles also strongly influences the neutrino emmission and the revival of the stalled shock by neutrinos .
The supernova event rate is 1-2 per century per galaxy and about 5-40% of them produce BHs in delayed collapse (through fall-back accretion), or direct collapse.
Initial rotation rates. Of considerable importance is the initial rotation rate of proto-neutron stars, since (as will be detailed in the next sections) most mechanisms for emission of detectable gravitational waves from compact objects require very rapid rotation at birth (rotational periods of the order of a few milliseconds or less). Since most massive stars have non-negligible rotation rates (some even rotate near their break-up limit), simple conservation of angular momentum would suggest a proto-neutron star to be strongly differentially rotating with very high rotation rates and this picture is supported by numerical simulations of rotating core collapse.
On the other hand, observationally, we know of pulsars (that have been associated with a supernova remnant) rotating only as fast as with a 16ms period, which suggest a period of at least several milliseconds at birth. A possible explanation for this discrepancy is that the pre-supernova core has been slowed down by e.g. magnetic torques. The most recent evolutionary models of rotating core-collapse progenitors (including the effect of magnetic torques) suggest that progenitors of neutron stars with typical masses of $`1.4M_{}`$ are indeed slowly rotating, producing remnants with periods at birth of the order of 10 ms. Nevertheless, the same study finds that very massive progenitors evolve so rapidly that the angular momentum transfer out of the core, by magnetic torques, is diminished, yielding heavy proto-neutron stars with rotational periods of the order of a few milliseconds (provided mass-loss takes place to prevent black hole formation). Binary interactions can also accelerate some phases of progenitor evolution, allowing for fast initial spins. If the magnetic torques operate efficiently, it is clear that rapid rotation at birth will have an event rate much smaller than the usual galactic supernova rate. Still, strong emission mechanisms (e.g. bar-mode instabilities) could yield detectable signals at acceptable event rates.
Other ways to form a rapidly rotating proto-neutron star would be through fall-back accretion, through the accretion-induced collapse of a white dwarf or through the merger of binary white dwarfs in globular clusters. It is also relevant to take into account current gamma-ray-burst models. The collapsar model requires high rotation rates of a proto-black hole. In addition, a possible formation scenario for magnetars involves a rapidly rotating protoneutron star formed through the collape of a very massive progenitor and some observational evidence is already emerging.
Gravitational wave emission. Gravitational waves from core collapse have a rich spectrum, reflecting the various stages of this event. The initial signal is emitted due to the changing axisymmetric quadrupole moment during collapse. In the case of neutron star formation, the quadrupole moment typically becomes larger, as the core spins up during contraction. In contrast, when a rapidly rotating neutron star collapses to form a Kerr black hole, the axisymmetric quadrupole moment first increases but is finally reduced by a large factor when the black hole is formed.
A second part of the gravitational wave signal is produced when gravitational collapse is halted by the stiffening of the equation of state above nuclear densities and the core bounces, driving an outwards moving shock. The dense fluid undergoes motions with relativistic speeds ($`v/c0.20.4`$) and a rapidly rotating proto-neutron star thus oscillates in several of its axisymmetric normal modes of oscillation. This quasi-periodic part of the signal could last for hundreds of oscillation periods, before being effectively damped. If, instead, a black hole is directly formed, then black hole quasinormal modes are excited, lasting for only a few oscillation periods. A combination of neutron star and black hole oscillations will appear if the proton-neutron star is not stable but collapses to a black hole.
In a rotating proto-neutron star, nonaxisymmetric processes can yield additional types of gravitational wave signals. Such processes are dynamical instabilities, secular gravitational-wave driven instabilities or convection inside the proto-neutron star and in its surrounding hot envelope. Anisotropic neutrino emission is accompanied by a gravitational wave signal. Nonaxisymmetries could already be present in the pre-collapse core and become amplified during collapse. Furthermore, if there is persistent fall-back accretion onto a proto-neutron star or black hole, these can be brought into ringing.
Below, we discuss in more detail those processes which result in high frequency gravitational radiation.
### 2.1 Neutron star formation
Core collapse as a potential source of GWs has been studied for more than three decades (some of the most recent simulations can be found in ). The main differences between the various studies are the progenitor models (slowly or rapidly rotating), equation of state (polytropic or realistic), gravity (Newtonian or relativistic) and neutrino emission (simple, sophisticated or no treatment). In general, the gravitational wave signal from neutron star formation is divided into a core bounce signal, a signal due to convective motions and a signal due to anisotropic neutrino emission.
Core bounce signal. The core bounce signal is produced due to rotational flattening and excitation of normal modes of oscillations, the main contributions coming from the axisymmetric quadrupole ($`l=2`$) and quasi-radial ($`l=0`$) modes (the latter radiating through its rotationally acquired $`l=2`$ piece). If detected, such signals will be a unique probe for the high-density EOS of neutron stars. The strength of this signal is sensitive to the available angular momentum in the progenitor core. If the progenitor core is rapidly rotating, then core bounce signals from Galactic supernovae ($`d10`$kpc) are detectable even with the initial LIGO/Virgo sensitivity at frequencies $``$1kHz. In the best-case scenario, advanced LIGO could detect signals from distances of 1Mpc, but not from the Virgo cluster ($``$15Mpc), where the event rate would be high. The typical GW amplitude from 2D numerical simulations for an observer located in the equatorial plane of the source is
$$h9\times 10^{21}\epsilon \left(\frac{10\mathrm{k}\mathrm{p}\mathrm{c}}{d}\right)$$
(1)
where $`\epsilon 1`$ is the normalized GW amplitude. For such rapidly rotating initial models, the total energy radiated in GWs during the collapse is $`10^610^8M_{}c^2`$. If, on the other hand, progenitor cores are slowly rotating (due to e.g. magnetic torques), then the signal strength is significantly reduced, but, in the best case, is still within reach of advanced LIGO for galactic sources.
Normal mode oscillations, if excited in an equilibrium star at a small to moderate amplitude, would last for hundreds to thousands of oscillation periods, being damped only slowly by gravitational wave emission or viscosity. However, the proto-neutron star immediately after core bounce has a very different structure than a cold equilibrium star. It has a high internal temperature and is surrounded but an extended, hot envelope. Nonlinear oscillations excited in the core after bounce can penetrate into the hot envelope. Through this damping mechanism, the normal mode oscillations are damped on a much shorter timescale (on the order of ten oscillation periods), which is typically seen in the core collapse simulations mentioned above.
Convection signal. The post-shock region surrounding a proto-neutron star is convectively unstable to both low-mode and high-mode convection. Neutrino emission also drives convection in this region. The most realistic 2D simulations of core collapse to date have shown that the gravitational wave signal from convection significantly exceeds the core bounce signal for slowly rotating progenitors, being detectable with advanced LIGO for galactic sources, and is detectable even for nonrotating collapse. For slowly rotating collapse, there is a detectable part of the signal in the high-frequency range of 700Hz-1kHz, originating from convective motions that dominate around 200ms after core bounce. Thus, if both core a bounce signal and a convection signal would be detected in the same frequency range, these would be well separated in time.
Neutrino signal. In many simulations the gravitational wave signature of anisotropic neutrino emission has also been considered. This type of signal can be detectable by advanced LIGO for galactic sources, but the main contribution is at low frequencies for a slowly rotating progenitor. For rapidly rotating progenitors, stronger contributions at high frequencies could be present, but would probably be burried within the high-frequency convection signal.
Numerical simulations of neutron star formation have gone a long way, but a fully consistent 3D simulation including relativistic gravity, neutrino emission and magnetic fields is still missing. The combined treatment of these effects might not change the above estimations by orders of magnitude but it will provide more conclusive answers. There are also issues that need to be understood such as pulsar kicks (velocities exceeding 1000 km/s) which suggest that in a fraction of newly-born NSs (and probably BHs) the formation process may be strongly asymmetric. Better treatment of the microphysics and construction of accurate progenitor models for the angular momentum distributions are needed. All these issues are under investigation by many groups.
### 2.2 Neutron star ringing through fall-back accretion
A possible mechanism for the excitation of oscillations in a proto-neutron star after core bounce is the fall-back accretion of material that has not been expelled by the revived supernova shock. The isotropy of this material is expected to be broken due to e.g. rotation or nonaxisymmetric convective motions, thus a large number of oscillation modes will be excited as this material falls back onto the neutron star. This process is, of course, complex and the detectability of gravitational waves from these oscillations will depend on several factors, such as the fall-back accretion rate, the degree of asymmetry of the fall-back material the structure of the proto-neutron star envelope, the presence of magnetic fields etc.
Recently, the ringing of a neutron star through fall-back accretion has been modeled through relativistic 2D nonlinear hydrodynamical simulations. Quadrupolar shells of matter were accreted on a static neutron star (in the approximation that the background spacetime remains unchanged). Gravitational waves were then extracted through the Zerilli-Moncrief formalism. The gravitational wave signal from such a process comprises a narrow peak at the $`l=2`$ normal mode frequency of the neutron star and a very broad peak, featuring interference fringes, centered at a much higher frequency. Since the frequency of the broad peak is still too low to be identified with a $`w`$-mode of the star, the interpretation for this part of the signal is that it is related to the motion of the fluid shell and the reflection of the gravitational-wave pulse from this motion in the external Zerilli potential, which also creates the interference fringes. The accretion of a quadrupolar shell containing 1% of the mass of the star releases gravitational waves with a total energy similar to the energy emitted immediately after core bounce. It is thus interesting to consider this mechanism in more detail, since the excitation of the normal modes in the neutron star happens when the star has already cooled somewhat (compared to the proto-neutron star immediately after core bounce) which simplifies the identification of observed oscillations with normal modes of cold neutron star models. The formation of a dense torus as as a result of stellar gravitational collapse, binary neutron star merger or disruption. Such a system either becomes unstable to the runaway instability or exhibit a regular oscillatory behavior, resulting in a quasi-periodic variation of the accretion rate as well as of the mass quadrupole suggesting a new sources of potentially detectable gravitational waves.
### 2.3 Black hole formation
The gravitational-wave emission from the formation of a Kerr BH is a sum of two signals: the collapse signal and the BH ringing. The collapse signal is produced due to the changing multipole moments of the spacetime during the transition from a rotating iron core or proto-neutron star to a Kerr BH. A uniformly rotating neutron star has an axisymmetric quadrupole moment given by
$$Q=a\frac{J^2}{M}$$
(2)
where $`a`$ depends on the equation of state and is in the range of $`28`$ for 1.4$`M_{}`$ models. This is several times larger in magnitude than the corresponding qudrupole moment of a Kerr black hole ($`a=1`$). Thus, the reduction of the axisymmetric quadrupole moment is the main source of the collapse signal. Once the BH is formed, it continues to oscillate in its axisymmetric $`l=2`$ QNM, until all oscillation energy is radiated away and the stationary Kerr limit is approached.
The numerical study of rotating collapse to BHs was pioneered by Nakamura but first waveforms and gravitational-wave estimates were obtained by Stark and Piran . These simulations we performed in 2D, using approximate initial data (essentially a spherical star to which angular momentum was artificially added). A new 3D computation of the gravitational wave emission from the collapse of unstable uniformly rotating relativistic polytropes to Kerr BHs finds that the energy emitted is
$$\mathrm{\Delta }E1.5\times 10^6(M/M_{}),$$
(3)
significantly less than the result of Stark and Piran. Still, the collapse of an unstable 2$`M_{}`$ rapidly rotating neutron star leads to a characteristic gravitational-wave amplitude $`h_c3\times 10^{21}`$, at a frequency of $`5.5`$kHz, for an event at 10kpc. Emission is mainly through the ”+” polarization, with the ”$`\times `$” polarization being an order of magnitude weaker.
Whether a BH forms promptly after collapse or a delayed collapse takes place depends sensitively on a number of factors, such as the progenitor mass and angular momentum and the high-density EOS. The most detailed investigation of the influence of these factors on the outcome of collapse has been presented recently in, where it was found that shock formation increases the threshold for black hole formation by $`2040\%`$, while rotation results in an increase of at most 25%.
### 2.4 Black hole ringing through fall-back or hyper-accretion
Single events. A black hole can form after core collapse, if fall-back accretion increases the mass of the proto-neutron star above the maximum mass allowed by axisymmetric stability. Material falling back after the black hole is formed excites the black hole quasi-normal modes of oscillation. If, on the other hand, the black hole is formed directly through core collapse (without a core bounce taking place) then most of the material of the progenitor star is accreted at very high rates ($`12M_{}`$/s) into the hole. In such hyper-accretion the black hole’s quasi-normal modes (QNM) can be excited for as long as the process lasts and until the black hole becomes stationary. Typical frequencies of the emitted GWs are in the range 1-3kHz for $`310M_{}`$ BHs.
The frequency and the damping time of the oscillations for the $`l=m=2`$ mode can be estimated via the relations
$`\sigma `$ $``$ $`3.2\mathrm{kHz}M_{10}^1\left[10.63(1a/M)^{3/10}\right]`$ (4)
$`Q`$ $`=`$ $`\pi \sigma \tau 2\left(1a\right)^{9/20}`$ (5)
These relations together with similar ones either for the 2nd QNM or the $`l=2`$, $`m=0`$ mode can uniquely determine the mass $`M`$ and angular momentum parameter $`a`$ of the BH if the frequency and the damping time of the signal have been accurately extracted . The amplitude of the ring-down waves depends on the BH’s initial distortion, i.e. on the nonaxisymmetry of the blobs or shells of matter falling into the BH. If matter of mass $`\mu `$ falls into a BH of mass $`M`$, then the gravitational wave energy is roughly
$$\mathrm{\Delta }Eϵ\mu c^2(\mu /M)$$
(6)
where $`ϵ`$ is related to the degree of asymmetry and could be $`ϵ0.01`$ . This leads to an effective GW amplitude
$$h_{\mathrm{eff}}2\times 10^{21}\left(\frac{ϵ}{0.01}\right)\left(\frac{10\mathrm{M}\mathrm{p}\mathrm{c}}{d}\right)\left(\frac{\mu }{M_{}}\right)$$
(7)
Resonant driving. If hyper-accretion proceeds through an accretion disk around a rapidly spinning Kerr BH, then the matter near the marginally bound orbit radius can become unstable to the magnetorotational (MRI) instability, leading to the formation of large-scale asymmetries. Under certain conditions, resonant driving of the BH QNMs could take place. Such a continuous signal could be integrated, yielding a much larger signal to noise ratio than a single event. For a 15$`M_{}`$ nearly maximal Kerr BH created at 27Mpc the integrated signal becomes detectable by LIGO II at a frequency of $`1600`$Hz, especially if narrow-banding is used.
## 3 Rotational instabilities
If proto-neutron stars rotate rapidly, nonaxisymmetric dynamical instabilities can develop. These arise from non-axisymmetric perturbations having angular dependence $`e^{im\varphi }`$ and are of two different types: the classical bar-mode instability and the more recently discovered low-$`T/|W|`$ bar-mode and one-armed spiral instabilities, which appear to be associated to the presence of corotation points. Another class of nonaxisymmetric instabilities are secular instabilities, driven by dissipative effects, such as fluid viscosity or gravitational radiation.
### 3.1 Dynamical instabilities
Classical bar-mode instability. The classical $`m=2`$ bar-mode instability is excited in Newtonian stars when the ratio $`\beta =T/|W|`$ of the rotational kinetic energy $`T`$ to the gravitational binding energy $`|W|`$ is larger than $`\beta _{\mathrm{dyn}}=0.27`$. The instability grows on a dynamical time scale (the time that a sound wave needs to travel across the star) which is about one rotational period and may last from 1 to 100 rotations depending on the degree of differential rotation in the PNS.
The bar-mode instability can be excited in a hot PNS, a few milliseconds after core bounce, or, alternatively, it could also be excited a few tenths of seconds later, when the PNS cools due to neutrino emission and contracts further, with $`\beta `$ becoming larger than the threshold $`\beta _{dyn}`$ ( $`\beta `$ increases roughly as $`1/R`$ during contraction). The amplitude of the emitted gravitational waves can be estimated as $`hMR^2\mathrm{\Omega }^2/d`$, where $`M`$ is the mass of the body, $`R`$ its size, $`\mathrm{\Omega }`$ the rotation rate and $`d`$ the distance of the source. This leads to an estimation of the GW amplitude
$$h9\times 10^{23}\left(\frac{ϵ}{0.2}\right)\left(\frac{f}{3\mathrm{k}\mathrm{H}\mathrm{z}}\right)^2\left(\frac{15\mathrm{M}\mathrm{p}\mathrm{c}}{d}\right)M_{1.4}R_{10}^2.$$
(8)
where $`ϵ`$ measures the ellipticity of the bar, $`M`$ is measured in units of $`1.4M_{}`$ and $`R`$ is measured in units of 10km. Notice that, in uniformly rotation Maclaurin spheroids, the GW frequency $`f`$ is twice the rotational frequency $`\mathrm{\Omega }`$. Such a signal is detectable only from sources in our galaxy or the nearby ones (our Local Group). If the sensitivity of the detectors is improved in the kHz region, signals from the Virgo cluster could be detectable. If the bar persists for many ($``$ 10-100) rotation periods, then even signals from distances considerably larger than the Virgo cluster will be detectable. Due to the requirement of rapid rotation, the event rate of the classical dynamical instability is considerably lower than the SN event rate.
The above estimates rely on Newtonian calculations; GR enhances the onset of the instability, $`\beta _{\mathrm{dyn}}0.24`$ and somewhat lower than that for large compactness (large $`M/R`$). Fully relativistic dynamical simulations of this instability have been obtained, including detailed waveforms of the associated gravitational wave emission. A detailed investigation of the required initial conditions of the progenitor core, which can lead to the onset of the dynamical bar-mode instability in the formed PNS, was presented in. The amplitude of gravitational waves was due to the bar-mode instability was found to be larger by an order of magnitude, compared to the axisymmetric core collapse signal.
Low-$`T/|W|`$ instabilities. The bar-mode instability may be excited for significantly smaller $`\beta `$, if centrifugal forces produce a peak in the density off the source’s rotational center. Rotating stars with a high degree of differential rotation are also dynamically unstable for significantly lower $`\beta _{\mathrm{dyn}}0.01`$ . According to this scenario the unstable neutron star settles down to a non-axisymmetric quasi-stationary state which is a strong emitter of quasi-periodic gravitational waves
$$h_{\mathrm{eff}}3\times 10^{22}\left(\frac{R_{\mathrm{eq}}}{30\mathrm{k}\mathrm{m}}\right)\left(\frac{f}{800\mathrm{H}\mathrm{z}}\right)^{1/2}\left(\frac{100\mathrm{M}\mathrm{p}\mathrm{c}}{d}\right)M_{1.4}^{1/2}.$$
(9)
The bar-mode instability of differentially rotating neutron stars is an excellent source of gravitational waves, provided the high degree of differential rotation that is required can be realized. One should also consider the effects of viscosity and magnetic fields. If magnetic fields enforce uniform rotation on a short timescale, this could have strong consequences regarding the appearance and duration of the dynamical nonaxisymmetric instabilities.
An $`m=1`$ one-armed spiral instability has also been shown to become unstable in proto-neutron stars, provided that the differential rotation is sufficiently strong . Although it is dominated by a “dipole” mode, the instability has a spiral character, conserving the center of mass. The onset of the instability appears to be linked to the presence of corotation points (a similar link to corotation points has been proposed for the low-$`T/|W|`$ bar mode instability ) and requires a very high degree of differential rotation (with matter on the axis rotating at least 10 times faster than matter on the equator). The $`m=1`$ spiral instability was recently observed in simulations of rotating core collapse, which started with the core of an evolved 20$`M_{}`$ progenitor star to which differential rotation was added. Growing from noise level ($`10^6`$) on a timescale of 5ms, the $`m=1`$ mode reached its maximum amplitude after $`100`$ms. Gravitational waves were emitted through the excitation of an $`m=2`$ nonlinear harmonic at a frequency of $`800`$Hz with an amplitude comparable to the core-bounce axisymmetric signal.
### 3.2 Secular gravitational-wave-driven instabilities
In a nonrotating star, the forward and backward moving modes of same $`(l,|m|)`$ (corresponding to $`(l,+m)`$ and $`(l,m)`$) have eigenfrequencies $`\pm |\sigma |`$. Rotation splits this degeneracy by an amount $`\delta \sigma m\mathrm{\Omega }`$ and both the prograde and retrograde modes are dragged forward by the stellar rotation. If the star spins sufficiently rapidly, a mode which is retrograde (in the frame rotating with the star) will appear as prograde in the inertial frame (a nonrotating observer at infinity). Thus, an inertial observer sees GWs with positive angular momentum emitted by the retrograde mode, but since the perturbed fluid rotates slower than it would in the absence of the perturbation, the angular momentum of the mode in the rotating frame is negative. The emission of GWs consequently makes the angular momentum of the mode increasingly negative, leading to the instability. A mode is unstable when $`\sigma (\sigma m\mathrm{\Omega })<0`$. This class of frame-dragging instabilities is usually referred to as Chandrasekhar-Friedman-Schutz (CFS) instabilities.
$`f`$-mode instability. In the Newtonian limit, the $`l=m=2`$ $`f`$-mode (which has the shortest growth time of all polar fluid modes) becomes unstable when $`T/|W|>0.14`$, which is near or even above the mass-shedding limit for typical polytropic EOSs used to model uniformly rotating neutron stars. Dissipative effects (e.g. shear and bulk viscosity or mutual friction in superfluids) leave only a small instability window near mass-shedding, at temperatures of $`10^9`$K. However, relativistic effects strengthen the instability considerably, lowering the required $`\beta `$ to $`0.060.08`$ for most realistic EOSs and masses of $`1.4M_{}`$ (for higher masses, such as hypermassive stars created in a binary NS merger, the required rotation rates are even lower).
Since PNSs rotate differentially, the above limits derived under the assumption of uniform rotation are too strict. Unless uniform rotation is enforced on a short timescale, due to e.g. magnetic braking, the $`f`$-mode instability will develop in a differentially rotating background, in which the required $`T/|W|`$ is only somewhat larger than the corresponding value for uniform rotation, but the mass-shedding limit is dramatically relaxed. Thus, in a differentially rotating PNS, the $`f`$-mode instability window is huge, compared to the case of uniform rotation and the instability can develop provided there is sufficient $`T/|W|`$ to begin with.
The $`f`$-mode instability is an excellent source of GWs. Simulations of its nonlinear development in the ellipsoidal approximation have shown that the mode can grow to a large nonlinear amplitude, modifying the background star from an axisymmetric shape to a differentially rotating ellipsoid. In this modified background the $`f`$-mode amplitude saturates and the ellipsoid becomes a strong emitter of gravitational waves, radiating away angular momentum until the star is slowed-down towards a stationary state. In the case of uniform density ellipsoids, this stationary state is the Dedekind ellipsoid, i.e. a nonaxisymmetric ellipsoid with internal flows but with a stationary (nonradiating) shape in the inertial frame. In the ellipsoidal approximation, the nonaxisymmetric pattern radiates gravitational waves sweeping through the LIGO II sensitivity window (from 1kHz down to about 100Hz) which could become detectable out to a distance of more than 100Mpc.
Two recent hydrodynamical simulations (in the Newtonian limit and using a post-Newtonian radiation-reaction potential) essentially confirm this picture. In a differentially rotating, $`N=1`$ polytropic model with a large $`T/|W|0.20.26`$ is chosen as the initial equilibrium state. The main difference of this simulation compared to the ellipsoidal approximation comes from the choice of EOS. For $`N=1`$ Newtonian polytropes it is argued that the secular evolution cannot lead to a a stationary Dedekind-like state does not exist. Instead, the $`f`$-mode instability will continue to be active until all nonaxisymmetries are radiated away and an axisymmetric shape is reached. This conclusion should be checked when relativistic effects are taken into account, since, contrary to the Newtonian case, relativistic $`N=1`$ uniformly rotating polytropes are unstable to the $`l=m=2`$ $`f`$-mode – however it has not become possible, to date, to construct relativistic analogs of Dedekind ellipsoids.
In the other recent simulation , the initial state was chosen to be a uniformly rotating, $`N=0.5`$ polytropic model with $`T/|W|0.18`$. Again, the main conclusions reached in are confirmed, however, the assumption of uniform initial rotation limits the available angular momentum that can be radiated away, leading to a detectable signal only out to about $`40`$Mpc. The star appears to be driven towards a Dedekind-like state, but after about 10 dynamical periods, the shape is disrupted by growing short-wavelength motions, which are suggested to arise because of a shearing type instability, such as the elliptic flow instability.
$`r`$-mode instability. Rotation does not only shift the spectra of polar modes; it also lifts the degeneracy of axial modes, give rise to a new family of inertial modes, of which the $`l=m=2`$ $`r`$-mode is a special member. The restoring force, for these oscillations is the Coriolis force. Inertial modes are primarily velocity perturbations. The frequency of the $`r`$-mode in the rotating frame of reference is $`\sigma =2\mathrm{\Omega }/3`$. According to the criterion for the onset of the CFS instability, the $`r`$-mode is unstable for any rotation rate of the star. For temperatures between $`10^710^9`$K and rotation rates larger than 5-10% of the Kepler limit, the growth time of the unstable mode is smaller than the damping times of the bulk and shear viscosity. The existence of a solid crust or of hyperons in the core and magnetic fields , can also significantly affect the onset of he instability (for extended reviews see ). The suppression of the $`r`$-mode instability by the presence of hyperons in the core is not expected to operate efficiently in rapidly rotating stars, since the central density is probably too low to allow for hyperon formation. Moreover, a recent calculation finds the contribution of hyperons to the bulk viscosity to be two orders of magnitude smaller than previously estimated. If accreting neutron stars in Low Mass X-Ray Binaries (LMXB, considered to be the progenitors of millisecond pulsars) are shown to reach high masses of $`1.8M_{}`$, then the EOS could be too stiff to allow for hyperons in the core (for recent observations that support a high mass for some millisecond pulsars see ).
The unstable $`r`$-mode grows exponentially until it saturates due to nonlinear effects at some maximum amplitude $`\alpha _{max}`$. The first computation of nonlinear mode couplings using second-order perturbation theory suggested that the $`r`$-mode is limited to very small amplitudes (of order $`10^310^4`$) due to transfer of energy to a large number of other inertial modes, in the form of a cascade, leading to an equilibrium distribution of mode amplitudes. The small saturation values for the amplitude are supported by recent nonlinear estimations based on the drift, induced by the r-modes, causing differential rotation. On the other hand, hydrodynamical simulations of limited resolution showed that an initially large-amplitude $`r`$-mode does not decay appreciably over several dynamical timescales , but on a somewhat longer timescale a catastrophic decay was observed indicating a transfer of energy to other modes, due to nonlinear mode couplings and suggesting that a hydrodynamical instability may be operating. A specific resonant 3-mode coupling was identified as the cause of the instability and a perturbative analysis of the decay rate suggests a maximum saturation amplitude $`\alpha _{max}<10^2`$. A new computation using second-order perturbation theory finds that the catastrophic decay seen in the hydrodynamical simulations can indeed be explained by a parametric instability operating in 3-mode couplings between the $`r`$-mode and two other inertial modes. Whether the maximum saturation amplitude is set by a network of 3-mode couplings or a cascade is reached, is, however, still unclear.
A neutron star spinning down due to the $`r`$-mode instability will emit gravitational waves of amplitude
$$h(t)10^{21}\alpha \left(\frac{\mathrm{\Omega }}{1\mathrm{k}\mathrm{H}\mathrm{z}}\right)\left(\frac{100\mathrm{k}\mathrm{p}\mathrm{c}}{d}\right)$$
(10)
Since $`\alpha `$ is small, even with LIGO II the signal is undetectable at large distances (VIRGO cluster) where the SN event rate is appreciable, but could be detectable after long-time integration from a galactic event. However, if the compact object is a strange star, then the instability may not reach high amplitudes ($`\alpha 10^310^4`$) but it will persist for a few hundred years (due to the different temperature dependence of viscosity in strange quark matter) and in this case there might be up to ten unstable stars in our galaxy at any time . Integrating data for a few weeks could lead to an effective amplitude $`h_{\mathrm{eff}}10^{21}`$ for galactic signals at frequencies $`7001000`$Hz. The frequency of the signal changes only slightly on a timescale of a few months, so that the radiation is practically monochromatic.
Other unstable modes. The CFS instability can also operate for core g-mode oscillations but also for w-mode oscillations, which are basically spacetime modes. In addition, the CFS instability can operate through other dissipative effects. Instead of the gravitational radiation, any radiative mechanism (such as electromagnetic radiation) can in principle lead to an instability.
### 3.3 Secular viscosity-driven instability
A different type of nonaxisymmetric instability in rotating stars is the instability driven by viscosity, which breaks the circulation of the fluid . The instability is suppressed by gravitational radiation, so it cannot act in the temperature window in which the CFS-instability is active. The instability sets in when the frequency of a prograde $`l=m`$ mode goes through zero in the rotating frame. In contrast to the CFS-instability, the viscosity-driven instability is not generic in rotating stars. The $`m=2`$ mode becomes unstable at a high rotation rate for very stiff stars and higher $`m`$-modes become unstable at larger rotation rates.
In Newtonian polytropes, the instability occurs only for stiff polytropes of index $`N<0.808`$ . For relativistic models, the situation for the instability becomes worse, since relativistic effects tend to suppress the viscosity-driven instability (while the CFS-instability becomes stronger). For the most relativistic stars, the viscosity-driven bar mode can become unstable only if $`N<0.55`$. For $`1.4M_{}`$ stars, the instability is present for $`N<0.67`$.
An investigation of the viscosity-driven bar mode instability, using incompressible, uniformly rotating triaxial ellipsoids in the post-Newtonian approximation finds that the relativistic effects increase the critical $`T/|W|`$ ratio for the onset of the instability significantly. More recently, new post-Newtonian and fully relativistic calculations for uniform-density stars show that the viscosity-driven instability is not as strongly suppressed by relativistic effects as suggested in . The most promising case for the onset of the viscosity-driven instability (in terms of the critical rotation rate) would be rapidly rotating strange stars , but the instability can only appear if its growth rate is larger than the damping rate due to the emission of gravitational radiation - a corresponding detailed comparison is still missing.
## 4 Accreting neutron stars in LMXBs
Spinning neutron stars with even tiny deformations are interesting sources of gravitational waves. The deformations might results from various factors but it seems that the most interesting cases are the ones in which the deformations are caused by accreting material. A class of objects called Low-Mass X-Ray Binaries (LMXB) consist of a fast rotating neutron star (spin $`270650`$Hz) torqued by accreting material from a companion star which has filled up its Roche lobe. The material adds both mass and angular momentum to the star, which, on timescales of the order of tenths of Megayears could, in principle, spin up the neutron star to its break up limit. One viable scenario suggests that the accreted material (mainly hydrogen and helium) after an initial phase of thermonuclear burning undergoes a non-uniform crystallization, forming a crust at densities $`10^810^9`$g/cm<sup>3</sup>. The quadrupole moment of the deformed crust is the source of the emitted gravitational radiation which slows-down the star, or halts the spin-up by accretion.
An alternative scenario has been proposed by Wagoner as a follow up of an earlier idea by Papaloizou-Pringle. The suggestion was that the spin-up due to accretion might excite the $`f`$-mode instability, before the rotation reaches the breakup spin. The emission of gravitational waves will torque down the star’s spin at the same rate as the accretion will torque it up, however, it is questionable whether the $`f`$-mode instability will ever be excited for old, accreting neutron stars. Following the discovery that the $`r`$-modes are unstable at any rotation rate, this scenario has been revived independently by Bildsten and Andersson, Kokkotas and Stergioulas . The amplitude of the emitted gravitational waves from such a process is quite small, even for high accretion rates, but the sources are persistent and in our galactic neighborhood the expected amplitude is
$$h10^{27}\left(\frac{1.6\text{ms}}{P}\right)^5\frac{1.5\text{kpc}}{D}.$$
(11)
This signal is within reach of advanced LIGO with signal recycling tuned at the appropriate frequency and integrating for a few months. This picture is in practice more complicated, since the growth rate of the $`r`$-modes (and consequently the rate of gravitational wave emission) is a function of the core temperature of the star. This leads to a thermal runaway due to the heat released as viscous damping mechanisms counteract the r-mode growth . Thus, the system executes a limit cycle, spinning up for several million years and spinning down in a much shorter period. The duration of the unstable part of the cycle depends critically on the saturation amplitude $`\alpha _{max}`$ of the $`r`$-modes . Since current computations suggest an $`\alpha _{max}10^310^4`$, this leads to a quite long duration for the unstable part of the cycle of the order of $`1Myear`$.
The instability window depends critically on the effect of the shear and bulk viscosity and various alternative scenarios might be considered. The existence of hyperons in the core of neutron stars induces much stronger bulk viscosity which suggests a much narrower instability window for the $`r`$-modes and the bulk viscosity prevails over the instability even in temperatures as low as $`10^8`$K . A similar picture can be drawn if the star is composed of “deconfined” $`u`$, $`d`$ and $`s`$ quarks - a strange star . In this case, there is a possibility that the strange stars in LMXBs evolve into a quasi-steady state with nearly constant rotation rate, temperature and mode amplitude emitting gravitational waves for as long as the accretion lasts. This result has also been found later for stars with hyperon cores . It is interesting that the stalling of the spin up in millisecond pulsars (MSPs) due to $`r`$-modes is in good agreement with the minimum observed period and the clustering of the frequencies of MSPs .
## 5 Binary mergers
Depending on the high-density EOS and their initial masses, the outcome of the merger of two neutron stars may not always be a black hole, but a hypermassive, differentially rotating compact star (even if it is only temporarily supported against collapse by differential rotation). A recent detailed simulation in full GR has shown that the hypermassive object created in a binary NS merger is nonaxisymmetric. The nonaxisymmetry lasts for a large number of rotational periods, leading to the emission of gravitational waves with a frequency of $`3`$kHz and an effective amplitude of $`67\times 10^{21}`$ at a large distance of 50Mpc. Such large effective amplitude may be detectable even by LIGO II at this high frequency.
The tidal disruption of a NS by a BH or the merging of two NSs may give valuable information for the radius and the EoS if we can recover the signal at frequencies higher than 1 kHz.
## 6 Gravitational-wave asteroseismology
If various types of oscillation modes are excited during the formation of a compact star and become detectable by gravitational wave emission, one could try to identify observed frequencies with frequencies obtained by mode-calculations for a wide parameter range of masses, angular momenta and EOSs. . Thus, gravitational wave asteroseismology could enable us to estimate the mass, radius and rotation rate of compact stars, leading to the determination of the ”best-candidate” high-density EoS, which is still very uncertain. For this to happen, accurate frequencies for different mode-sequences of rapidly rotating compact objects have to be computed.
For slowly rotating stars, the frequencies of $`f`$, $`p`$ and $`w`$ modes are still unaffected by rotation, and one can construct approximate formulae in order to relate observed frequencies and damping times of the various stellar modes to stellar parameters. For example, for the fundamental oscillation ($`l=2`$) mode ($`f`$-mode) of non-rotating stars one obtains
$`\sigma (\mathrm{kHz})`$ $``$ $`0.8+1.6M_{1.4}^{1/2}R_{10}^{3/2}+\delta _1m\overline{\mathrm{\Omega }}`$ (12)
$`\tau ^1(\mathrm{secs}^1)`$ $``$ $`M_{1.4}^3R_{10}^4\left(22.914.7M_{1.4}R_{10}^1\right)+\delta _2m\overline{\mathrm{\Omega }}`$ (13)
where $`\overline{\mathrm{\Omega }}`$ is the normalized rotation frequency of the star, and $`\delta _1`$ and $`\delta _2`$ are constants estimated by sampling data from various EOSs. The typical frequencies of NS oscillation modes are larger than 1kHz. Since each type of mode is sensitive to the physical conditions where the amplitude of the mode is largest, the more oscillations modes can be identified through gravitation waves, the better we will understand the detailed internal structure of compact objects, such as the existence of a possible superfluid state of matter.
If, on the other hand, some compact stars are born rapidly rotating with moderate differential rotation, then their central densities will be much smaller than the central density of a nonrotating star or same baryonic mass. Correspondingly, the typical axisymmetric oscillation frequencies will be smaller than 1kHz, which is more favorable for the sensitivity window of current interferometric detectors. Indeed, axisymmetric simulations of rotating core-collapse have shown that if a rapidly rotating NS is created, then the dominant frequency of the core-bounce signal (originating from the fundamental $`l=2`$ mode or the $`l=2`$ piece of the fundamental quasi-radial mode) is in the range 600Hz-1kHz.
If different type of signals are observed after core collapse, such as both an axisymmetric core-bounce signal and a nonaxisymmetric one-armed instability signal, with a time separation of the order of 100ms, this would yield invaluable information about the angular momentum distribution in the proto-neutron stars.
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# Computing with matrix invariants
## 1. Introduction to invariant theory
All considerations in this paper are over an arbitrary field $`K`$ of characteristic 0. If not explicitly stated, all vector spaces and algebras are over $`K`$, and the algebras are unitary and commutative.
To get some idea about invariant theory, we start with the following well known result from the undergraduate algebra course.
Let $`n2`$ be an integer and let
$$A=K[X]=K[x_1,\mathrm{},x_n]$$
be the algebra of polynomials in $`n`$ variables. The algebra of symmetric polynomials $`A^{S_n}=K[X]^{S_n}`$ consists of all polynomials $`f(x_1,\mathrm{},x_n)A`$ such that
$$f(x_{\sigma (1)},\mathrm{},x_{\sigma (n)})=f(x_1,\mathrm{},x_n)$$
for all permutations $`\sigma `$ in the symmetric group $`S_n`$.
###### Theorem 1.1.
(i) The algebra $`A^{S_n}`$ is generated by the elementary symmetric functions
$$e_1=x_1+\mathrm{}+x_n=x_i,$$
$$e_2=x_1x_2+x_1x_3+\mathrm{}+x_{n1}x_n=\underset{i<j}{}x_ix_j,$$
$$\mathrm{}$$
$$e_n=x_1\mathrm{}x_n.$$
(ii) The presentation of the symmetric polynomials as polynomials of $`e_1,\mathrm{},e_n`$ is unique, i.e.
$$A^{S_n}K[y_1,\mathrm{},y_n].$$
As a measure how many are the symmetric polynomials in $`n`$ variables we may use the dimensions of the vector spaces $`A_k^{S_n}`$ of the homogeneous symmetric polynomials of degree $`k`$. For example, for $`n=3`$, the vector space $`A_k^{S_3}`$ has a basis consisting of all products $`e_1^{k_1}e_2^{k_2}e_3^{k_3}`$, where $`k_1+2k_2+3k_3=k`$. Instead, we consider the generating function of the sequence $`\text{dim}(A_k^{S_n})`$, namely the formal power series
$$H(A^{S_n},t)=\underset{k0}{}\text{dim}(A_k^{S_n})t^k,$$
which is called the Hilbert series of $`A^{S_n}`$. Since the set of products
$$\{e_1^{k_1}e_2^{k_2}\mathrm{}e_n^{k_n}k_1,\mathrm{},k_n0\}$$
is a basis of $`A^{S_n}`$, it is easy to see that
$$H(A^{S_n},t)=\underset{k_j0}{}t^{k_1}t^{2k_2}\mathrm{}t^{nk_n}$$
$$=(1+t+t^2+\mathrm{})(1+t^2+t^4+\mathrm{})\mathrm{}(1+t^n+t^{2n}+\mathrm{})=\underset{i=1}{\overset{n}{}}\frac{1}{1t^i}.$$
Let us consider another example.
###### Example 1.2.
Let $`n=3`$, $`A=K[x_1,x_2,x_3]`$ and let $`G`$ be the subgroup of $`S_3`$ generated by the cycle $`\rho =(123)`$. We denote by $`A^G`$ the algebra of “cyclicly symmetric polynomials”,
$$A^G=\{f(x_1,x_2,x_3)Af(x_2,x_3,x_1)=f(x_1,x_2,x_3)\}.$$
One can show that $`A^G`$ is generated by the elementary symmetric functions $`e_1,e_2,e_3`$ and one more polynomial
$$f_4=x_1^2x_2+x_2^2x_3+x_3^2x_1,$$
which is not symmetric. The sum
$$a=(x_1^2x_2+x_2^2x_3+x_3^2x_1)+(x_1x_2^2+x_2x_3^2+x_3x_1^2)$$
of $`f_4`$ and its “other half” $`x_1x_2^2+x_2x_3^2+x_3x_1^2`$, is symmetric. The polynomial $`b=f_4(af_4)`$ is also symmetric. Hence both $`a,b`$ belong to $`K[e_1,e_2,e_3]`$ and $`f_4`$ satisfies the relation $`f_4^2af_4+b=0`$ with coefficients
$$a=\alpha (e_1,e_2,e_3),b=\beta (e_1,e_2,e_3)K[e_1,e_2,e_3].$$
There exists a natural homomorphism $`\pi :K[y_1,y_2,y_3,y_4]A^G`$ defined by
$$y_1e_1,y_2e_2,y_3e_3,y_4f_4,$$
and
$$r(y_1,y_2,y_3,y_4)=y_4^2\alpha (y_1,y_2,y_3)y_4+\beta (y_1,y_2,y_3)$$
belongs to the kernel of $`\pi `$. One can show that the ideal $`\text{Ker}(\pi )`$ is generated by $`r(y_1,y_2,y_3,y_4)`$. Hence $`A^G`$ has a basis
$$\{e_1^{k_1}e_2^{k_2}e_3^{k_3},e_1^{k_1}e_2^{k_2}e_3^{k_3}f_4k_1,k_2,k_30\}$$
and, if $`A_k^G`$ is the vector space of the homogeneous cyclicly symmetric polynomials of degree $`k`$, then the Hilbert series of $`A^G`$ is
$$H(A^G,t)=\underset{k0}{}\text{dim}(A_k^G)t^k$$
$$=\underset{k_i0}{}t^{k_1}(t^2)^{k_2}(t^3)^{k_3}+\left(\underset{k_i0}{}t^{k_1}(t^2)^{k_2}(t^3)^{k_3}\right)t^3=\frac{1+t^3}{(1t)(1t^2)(1t^3)}.$$
Let $`V_n`$ be the vector space with basis $`X=\{x_1,\mathrm{},x_n\}`$ and let the general linear group $`GL_n=GL_n(K)=GL(V_n)`$ act canonically on $`V_n`$. If we identify $`gGL_n`$ with the invertible matrix $`g=(g_{ij})`$, $`g_{ij}K`$, then the action of $`g`$ on $`x_j`$ is defined by
$$g(x_j)=g_{1j}x_1+\mathrm{}+g_{nj}x_n.$$
This action is extended diagonally on $`A=K[X]=K[x_1,\mathrm{},x_n]`$:
$$g(f(x_1,\mathrm{},x_n))=f(g(x_1),\mathrm{},g(x_n)),gGL_n,f(x_1,\mathrm{},x_n)A.$$
###### Definition 1.3.
Let $`G`$ be any subgroup of $`GL_n`$. (One says that the group $`G`$ is linear.) The polynomial $`f(x_1,\mathrm{},x_n)`$ is called a $`G`$-invariant if
$$g(f(x_1,\mathrm{},x_n))=f(x_1,\mathrm{},x_n)$$
for all $`gG`$. We denote the algebra of $`G`$-invariants by $`A^G`$.
For example, if $`G`$ is the subgroup of all permutational matrices
$$g_\sigma =\underset{j=1}{\overset{n}{}}e_{\sigma (j),j},\sigma S_n,$$
where the $`e_{ij}`$’s are the matrix units, then $`GS_n`$ and the algebra of $`G`$-invariants consists of all symmetric polynomials in $`n`$ variables.
For $`n=3`$ and $`G`$ being the cyclic group of order 3 generated by the matrix
$$g=\left(\begin{array}{ccc}0& 0& 1\\ 1& 0& 0\\ 0& 1& 0\end{array}\right),$$
the algebra of invariants $`K[x_1,x_2,x_3]^G`$ consists of all cyclicly symmetric polynomials.
It is easy to see that the algebra of invariants is graded with respect to the usual degree of polynomials. This means that if the polynomial $`f=f_0+f_1+\mathrm{}+f_p`$ is a $`G`$-invariant and $`f_k`$ is its homogeneous component of degree $`k`$, then $`f_k`$ is also $`G`$-invariant. If $`K[X]_k^G`$ is the homogeneous component of degree $`k`$ of $`K[X]^G`$, then one may define the Hilbert series of $`K[X]^G`$ as the formal power series
$$H(K[X]^G,t)=\underset{k0}{}\text{dim}(K[X]_k^G)t^k.$$
We state of list of problems which are among the main problems in invariant theory. For more details, see some book on invariant theory, e.g. Dieudonné and Carrell \[DC\] or Springer \[Sp\]; for computational aspects of invariant theory see Sturmfels \[St\].
Recall that if $`R`$ is a finitely generated commutative algebra, then $`R`$ is a homomorphic image of the polynomial algebra with the same number of generators. Hence $`RK[y_1,\mathrm{},y_m]/I`$ for some $`m`$ and for some ideal $`I`$. The equations $`f(y_1,\mathrm{},y_m)=0`$, $`fI`$, are called relations of the algebra $`R`$ with respect to the given set of generators. The elements of any generating set of the ideal $`I`$ are called defined relations of $`R`$. A finitely generated algebra with a finite system of defining relations is called finitely presented.
###### Problems 1.4.
(i) For a given linear group $`G`$, find a set of generators of the algebra of $`G`$-invariants.
(ii) Is it true that the algebra of invariants $`K[X]^G`$ is finitely generated for any subgroup $`G`$ of $`GL_n`$?
(iii) If $`K[X]^G`$ is finitely generated, find its defining relations. Is $`K[X]^G`$ always finitely presented?
(iv) When $`K[X]^G`$ is isomorphic to a polynomial algebra (i.e. has no defining relations with respect to a suitable system of generators)?
(v) Calculate the Hilbert series of $`K[X]^G`$. Is it a rational function? What kind of properties of $`K[X]^G`$ can be recovered from its Hilbert series?
We give short comments on the solutions of the above problems in the case of finite groups. Problem 1.4 (i) is solved by the following theorem of Emmy Noether.
###### Theorem 1.5.
(Endlichkeitssatz, \[No\]) The algebra of invariants $`K[X]^G`$ is finitely generated for any finite subgroup $`G`$ of $`GL_n`$, by a set of invariants of degree $`|G|`$, the order of the group $`G`$.
It is well known that the linear operator of $`K[X]`$ (called the Reynolds operator)
$$\rho :f(x_1,\mathrm{},x_n)\frac{1}{|G|}\underset{gG}{}g(f)$$
is a projection of $`K[X]`$ onto $`K[X]^G`$, i.e. $`\rho ^2=\rho `$ and $`\rho (K[X])=K[X]^G`$. If we take $`f=x_1^{k_1}\mathrm{}x_n^{k_n}`$ and consider all monomials of degree $`|G|`$, then we obtain a system of generators of $`K[X]^G`$.
Problem 1.4 (ii) was the main motivation for the Hilbert 14-th Problem from the famous list of 23 open problems presented by Hilbert at the International Congress of Mathematicians in Paris, 1900, \[H2\]. Only in the late 1950’s Nagata \[Na2\] found the first counterexamples to the 14-th Hilbert Problem. We recommend the recent survey by Freudenburg \[Fr\] for the state-of-the-art of the problem and for different approaches to construct other counterexamples.
###### Theorem 1.6.
(Hilbert Basissatz, or Hilbert Basis Theorem, \[H1\]) Every ideal of $`K[y_1,\mathrm{},y_m]`$ is finitely generated.
As an immediate consequence of the Hilbert Basissatz and the theorem of Emmy Noether we obtain that for finite groups $`G`$ the algebra $`K[X]^G`$ is finitely presented, i.e. the positive solution of Problem 1.4 (iii). Since the degree of the generators of $`K[X]^G`$ is bounded, in principle, one can answer also the question which are the defining relations of $`K[X]^G`$.
Problem 1.4 (iv) is solved by the following theorem of Chevalley-Shephard-Todd \[Ch, ST\].
###### Theorem 1.7.
For a finite group $`G`$, the algebra of invariants $`K[X]^G`$ is isomorphic to a polynomial algebra (in $`n=|X|`$ variables) if and only if $`GGL_n`$ is generated by pseudo-reflections.
Recall that the matrix $`gGL_n`$ is called a pseudo-reflection, if it is similar to a matrix of the form
$$\text{diag}(\xi ,1,\mathrm{},1)=\left(\begin{array}{cccc}\xi & 0& \mathrm{}& 0\\ 0& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 1\end{array}\right),$$
it is a reflection if $`\xi =1`$. The symmetric group $`S_n`$ is generated by the transpositions $`\tau _j=(1j)`$, $`j=2,\mathrm{},n`$, and the images $`g_{\tau _j}=1e_{11}e_{jj}+e_{j1}+e_{1j}`$ of $`\tau _j`$ under the embedding of $`S_n`$ into $`GL_n`$ are reflections. Hence $`K[X]^{S_n}`$ is isomorphic to a polynomial algebra. On the other hand, the matrix $`g`$ from Example 1.2 (as well as its square $`g^2`$) is similar to
$$g=\left(\begin{array}{ccc}1& 0& 0\\ 0& \epsilon & 0\\ 0& 0& \epsilon ^2\end{array}\right),$$
where $`\epsilon `$ is a primitive third root of 1 in some extension of $`K`$. This implies that $`G=\{1,g,g^2\}GL_3`$ cannot be generated by pseudo-reflections and $`K[x_1,x_2,x_3]^G`$ is not isomorphic to a polynomial algebra.
Finally, the following Molien formula \[Mo\] answers Problem 1.4 (v).
###### Theorem 1.8.
For any finite group $`GGL_n`$
$$H(K[X]^G,t)=\frac{1}{|G|}\underset{gG}{}\frac{1}{\text{det}(1gt)}.$$
Here $`\text{det}(1gt)`$ is the determinant of the $`n\times n`$ matrix $`1gt`$. In particular, the Hilbert series $`H(K[X]^G,t)`$ is a rational function. We refer to the book of Stanley \[St2\] for the relationship between the algebraic properties of graded algebras and their Hilbert series.
The natural generalization of the case of finite linear groups is the case of reductive groups. Recall that, by the Maschke theorem, the finite subgroups $`G`$ of $`GL_n`$ are completely reducible. This means that if, for a suitable basis of $`V_n`$, the group $`G`$ consists of matrices of the block triangular form
$$\left(\begin{array}{cc}& \\ 0& \end{array}\right),$$
then, changing the basis we may assume that the matrices are of the block diagonal form
$$\left(\begin{array}{cc}& 0\\ 0& \end{array}\right).$$
A similar property holds for the class of reductive subgroups of $`GL_n`$. In particular, this holds for all classical groups ($`GL_k`$, $`SL_k`$, $`O_k`$, $`SO_k`$, $`Sp_k`$, $`U_k`$).
###### Theorem 1.9.
(i) (Hilbert-Serre, see e.g. \[DC\]) For any reductive subgroup $`G`$ of $`GL_n`$ the algebra of invariants $`K[X]^G`$ is finitely generated.
(ii) (Molien-Weyl formula, see \[W1\]) If $`GGL_n()`$ is compact, then one can define Haar measure on $`G`$, replace in the Molien formula the sum with an integral and obtain the formula for the Hilbert series of the algebra of invariants $`[X]^G`$.
## 2. Matrix invariants
In the sequel, we do not use the symbol $`X`$ for the set of variables $`\{x_1,\mathrm{},x_n\}`$ and we denote by $`X,X_i,Y,Y_i`$, etc. $`n\times n`$ matrices. We start with a special case of the main object of this paper. Let
$$X=(x_{ij})=\left(\begin{array}{ccc}x_{11}& \mathrm{}& x_{1n}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ x_{n1}& \mathrm{}& x_{nn}\end{array}\right)$$
be an $`n\times n`$ matrix with $`n^2`$ algebraically independent entries $`x_{ij}`$. Such a matrix is called a generic $`n\times n`$ matrix because every $`n\times n`$ matrix $`A=(a_{ij})`$ with entries $`\alpha _{ij}`$ from a commutative $`K`$-algebra $`S`$ can be obtained from $`X`$ by specialization of the variables $`x_{ij}`$ (i.e. a homomorphism $`K[x_{ij}]=K[x_{ij}i,j=1,\mathrm{},n]S`$ defined by $`x_{ij}\alpha _{ij}`$, $`i,j=1,\mathrm{},n`$). The group $`GL_n`$ acts on $`X`$ by conjugation:
$$gXg^1=Y=(y_{ij}),gGL_n,$$
where the $`y_{ij}`$’s are linear combinations of the $`x_{ij}`$’s. We define an action of $`GL_n`$ on the algebra $`K[x_{ij}]`$ by $`g:x_{ij}y_{ij}`$, $`i,j=1,\mathrm{},n`$, and are interested in the algebra of invariants $`K[x_{ij}]^{GL_n}`$. The obvious invariants come from the Cayley-Hamilton theorem. If
$$f_X(t)=\text{det}(1tX)=(1)^n(t^n+a_1t^{n1}+\mathrm{}+a_{n1}t+a_n)$$
is the characteristic polynomial of $`X`$, then the coefficients $`a_1,\mathrm{},a_n`$ depend on the entries of $`X`$. In particular, $`a_1=\text{tr}(X)`$ and $`a_n=(1)^n\text{det}(X)`$. The Cayley-Hamilton theorem states that $`f_X(X)=0`$. Since $`gf_X(X)g^1=f_X(gXg^1)`$ for any $`gGL_n`$, we obtain that the coefficients $`a_1,\mathrm{},a_n`$ are $`GL_n`$-invariants. The following theorem is a partial case of Theorem 2.5 below, but probably was known much before it. It seems impossible to trace its origin.
###### Theorem 2.1.
With respect to the above action of $`GL_n`$, the algebra of invariants $`K[x_{ij}]^{GL_n}`$ is generated by the coefficients of the characteristic polynomial of $`X`$.
Up to a sign, the coefficients of the characteristic polynomial are equal to the elementary symmetric polynomials evaluated on the eigenvalues of the matrix. Let
$$p_k=x_1^k+\mathrm{}+x_n^k$$
be the $`k`$-th power symmetric function. Since the characteristic of $`K`$ is equal to 0, the Newton formulas
$$p_ke_1p_{k1}+\mathrm{}+(1)^{k1}e_{k1}p_1+(1)^kke_k=0,k=1,\mathrm{},n,$$
allow to express the elementary symmetric polynomials $`e_1,\mathrm{},e_n`$ in terms of the power symmetric polynomials. Hence $`K[x_1,\mathrm{},x_n]^{S_n}`$ is generated by $`p_1,\mathrm{},p_n`$. If $`\xi _1,\mathrm{},\xi _n`$ are the eigenvalues of the matrix $`X`$ in the algebraic closure of the field of rational functions $`K(x_{ij})`$, then $`X`$ is similar to the upper triangular matrix
$$\left(\begin{array}{cccc}\xi _1& & \mathrm{}& \\ 0& \xi _2& \mathrm{}& \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& \xi _n\end{array}\right).$$
(It is known that $`X`$ is similar even to the diagonal matrix $`\text{diag}(\xi _1,\mathrm{},\xi _n)`$ with the same diagonal entries as the above matrix.) The coefficients of the characteristic polynomial $`f_X(X)`$ are
$$a_1=e_1(\xi _1,\mathrm{},\xi _n),a_2=e_2(\xi _1,\mathrm{},\xi _n),\mathrm{},a_n=(1)^ne_n(\xi _1,\mathrm{},\xi _n).$$
Similarly,
$$p_1(\xi _1,\mathrm{},\xi _n)=\text{tr}(X),p_2(\xi _1,\mathrm{},\xi _n)=\text{tr}(X^2),\mathrm{},p_n(\xi _1,\mathrm{},\xi _n)=\text{tr}(X^n).$$
In this way, as an immediately consequence of Theorem 2.1, we obtain:
###### Corollary 2.2.
For the above action of $`GL_n`$ on $`K[x_{ij}]`$, the algebra $`K[x_{ij}]^{GL_n}`$ is generated by the traces $`\text{tr}(X),\text{tr}(X^2),\mathrm{},\text{tr}(X^n)`$ of the powers of $`X`$.
Now we consider the more general situation, when the group $`GL_n`$ acts on several generic matrices. We fix an integer $`d2`$. Let
$$X_i=\left(x_{pq}^{(i)}\right),p,q=1,\mathrm{},n,i=1,\mathrm{},d,$$
be $`d`$ generic $`n\times n`$ matrices and let
$$\mathrm{\Omega }=\mathrm{\Omega }_{nd}=K[x_{pq}^{(i)}p,q=1,\mathrm{},n,i=1,\mathrm{},d]$$
be the polynomial algebra in $`dn^2`$ variables. The action of $`GL_n`$ on $`\mathrm{\Omega }`$ is as in the case of one generic matrix. If
$$gX_ig^1=Y_i=(y_{pq}^{(i)}),gGL_n,$$
then we define
$$g:x_{pq}^{(i)}y_{pq}^{9i)},p,q=1,\mathrm{},n,i=1,\mathrm{},d.$$
This action of $`GL_n`$ on $`\mathrm{\Omega }`$ is called the action of $`GL_n`$ by simultaneous conjugation on $`d`$ generic $`n\times n`$ matrices. Here are the main problems concerning the algebra of matrix invariants $`\mathrm{\Omega }_{nd}^{GL_n}`$.
###### Problem 2.3.
Describe the algebraic properties of $`\mathrm{\Omega }_{nd}^{GL_n}`$. In particular, find concrete sets of generators and defining relations, calculate the Hilbert series, etc.
We refer e.g. to the books by Procesi \[P1\], Formanek \[F4\], and Drensky and Formanek \[DF\] as a further reading on invariant theory of matrices.
Since the action of $`GL_n`$ on $`\mathrm{\Omega }_{nd}`$ is reductive, Classical Invariant Theory gives:
###### Theorem 2.4.
The algebra $`\mathrm{\Omega }_{nd}^{GL_n}`$ is finitely generated.
Traditionally, a result giving the explicit generators of the algebra of invariants of a linear group $`G`$ is called a first fundamental theorem of the invariant theory of $`G`$ and a result describing the relations between the generators is a second fundamental theorem. As in the case of a single generic matrix, the coefficients of the Cayley-Hamilton polynomial of any matrix expressed in terms of $`X_1,\mathrm{},X_d`$ is $`GL_n`$-invariant. In particular, all traces $`\text{tr}(X_{i_1}\mathrm{}X_{i_k})`$ are $`GL_n`$-invariant. A more sophisticated result than the previous theorem states:
###### Theorem 2.5.
(The First Fundamental Theorem of Matrix Invariants) The algebra $`\mathrm{\Omega }_{nd}^{GL_n}`$ is generated by all traces $`\text{tr}(X_{i_1}\mathrm{}X_{i_k})`$, where $`i_1,\mathrm{},i_k=1,\mathrm{},d`$.
It is difficult to judge about the priority in the proof of this theorem. (See the history of generic matrices in the paper by Formanek \[F5\].) The proof follows from general facts on invariant theory of $`GL_n`$ and can be found in the book by Gurevich \[Gu\] and in the papers by Sibirskii \[Si\] and Procesi \[P2\]. Although it was stated as well known by Kirillov \[Ki\], it seems that the understanding of the importance of the theorem is a result of its rediscovery by Procesi \[P2\].
Combining both theorems, we are looking for a concrete finite set of generators of $`\mathrm{\Omega }_{nd}^{GL_n}`$. We recall one of the most important theorems in the theory of algebras with polynomial identities.
###### Theorem 2.6.
(Nagata-Higman Theorem) If the (nonunitary) algebra $`R`$ is nil of bounded index $`n`$, i.e. $`r^n=0`$ for all $`rR`$, then $`R`$ is nilpotent, i.e. there exists an $`N=N(n)`$ such that $`r_1\mathrm{}r_N=0`$ for all $`r_1,\mathrm{},r_NR`$.
The Nagata-Higman theorem was established in 1953 by Nagata \[Na1\] for algebras over a field of characteristic 0 and then in 1956 generalized by Higman \[Hi\] when $`\text{char K}=p>n`$. Much later it was discovered that this theorem was first established in 1943 by Dubnov and Ivanov \[DI\] but their paper was overlooked by the mathematical community.
The class of nilpotency $`N(n)`$ in the Nagata-Higman theorem is related in the following nice way with invariant theory of matrices.
###### Theorem 2.7.
(Formanek \[F2\], Procesi \[P2, P3\], Razmyslov \[R2\]) Let $`N(n)`$ be the class of nilpotency in the Nagata-Higman theorem. Then the algebra of invariants $`\mathrm{\Omega }_{nd}^{GL_n}`$ is generated by the traces $`\text{tr}(X_{i_1}\mathrm{}X_{i_m})`$ of degree $`N(n)`$. For $`d`$ sufficiently large this bound is sharp.
It is important to know the exact value of the class of nilpotency $`N(n)`$ in the Nagata-Higman theorem. The upper bound given in the proof of Higman \[Hi\] is $`N(n)2^n1`$. The best known upper bound is due to Razmyslov \[R2\]. Applying trace polynomial identities of matrices, he obtained the bound $`N(n)n^2`$. The proof of the theorem of Razmyslov is given also in his book \[R3\] or in the book by Formanek \[F4\]. For a lower bound, Kuzmin \[Ku\] showed that $`N(n)\frac{1}{2}n(n+1)`$. A proof of the result of Kuzmin may be found also in the books by Drensky and Formanek \[DF\] or by Kanel-Belov and Rowen \[KBR\]. Hence
$$\frac{n(n+1)}{2}N(n)n^2.$$
###### Problem 2.8.
Find the exact value $`N(n)`$ of the class of nilpotency of nil algebras of index $`n`$ over a field of characteristic $`0`$.
###### Conjecture 2.9.
(Kuzmin \[Ku\]) The exact value $`N(n)`$ of the class of nilpotency of nil algebras of index $`n`$ over a field of characteristic $`0`$ is
$$N(n)=\frac{n(n+1)}{2}.$$
The only values of $`N(n)`$ are known for $`n4`$: Dubnov \[Du\] obtained in 1935
$$N(1)=1,N(2)=3,N(3)=6.$$
In 1993, Vaughan-Lee \[VL\] proved that
$$N(4)=10.$$
In this way the conjecture of Kuzmin is confirmed for $`n4`$. Recently, Shestakov and Zhukavets \[SZ\] have proved that the class of nilpotency of the two-generated algebras satisfying the identity $`x^5=0`$ is equal to 15, which agrees with the conjecture of Kuzmin for $`n=5`$. They have obtained the same result also in the more general setup of 2-generated superalgebras. Their proof is based on computer calculations with the GAP package.
Since $`N(2)=3`$, the algebra $`\mathrm{\Omega }_{2d}^{GL_2}`$ is generated by products of traces of degree $`3`$. The following result was established by Sibirskii \[Si\].
###### Theorem 2.10.
The elements
$$\text{tr}(X_i),\text{tr}(X_i^2),i=1,\mathrm{},d,\text{tr}(X_iX_j),1i<jd,$$
$$\text{tr}(X_iX_jX_k),1i<j<kd,$$
constitute a minimal set of generators of the algebra of $`2\times 2`$ matrix invariants $`\mathrm{\Omega }_{2d}^{GL_2}`$.
Now we give an idea about the Razmyslov-Procesi theory which is related with the second fundamental theorem of the matrix invariants, see Razmyslov \[R2\] and Procesi \[P2\], as well as the book by Razmyslov \[R3\] for other applications of his method. For simplicity we consider the case $`n=2`$ only. The Cayley-Hamilton theorem for $`2\times 2`$ matrices implies that
$$X^2\text{tr}(X)X+\text{det}(X)=0.$$
The Newton formulas give that
$$\text{det}(X)=\frac{1}{2}(\text{tr}^2(X)\text{tr}(X^2)).$$
This can be seen also directly. If $`\xi _1,\xi _2`$ are the eigenvalues of $`X`$, then
$$\text{tr}(X)=\xi _1+\xi _2,\text{tr}(X^2)=\xi _1^2+\xi _2^2,$$
$$\text{det}(X)=\xi _1\xi _2=\frac{1}{2}((\xi _1+\xi _2)^2(\xi _1^2+\xi _2^2))=\frac{1}{2}(\text{tr}^2(X)\text{tr}(X^2)).$$
In this way we obtain the mixed trace identity
$$c(X)=X^2\text{tr}(X)X+\frac{1}{2}(\text{tr}^2(X)\text{tr}(X^2))=0.$$
Now we consider the identity $`c(X_1+X_2)c(X_1)c(X_2)=0`$, i.e. we linearize the identity $`c(X)=0`$. In this way we obtain the mixed Cayley-Hamilton identity
$$\mathrm{\Psi }_2(X_1,X_2)=X_1X_2+X_2X_1\text{tr}(X_1)X_2\text{tr}(X_2)X_1+\text{tr}(X_1)\text{tr}(X_2)\text{tr}(X_1X_2)=0.$$
Since the trace is a nondegenerate bilinear form on $`M_2(K)`$, the vanishing of the polynomial $`\mathrm{\Psi }_2(X_1,X_2)`$ on $`M_2(K)`$ is equivalent to the vanishing of the pure Cayley-Hamilton identity
$$\mathrm{\Phi }_2(X_1,X_2,X_3)=\text{tr}(\mathrm{\Psi }_2(X_1,X_2)X_3)=0$$
on all $`2\times 2`$ matrices. Direct calculations show that
$$0=\mathrm{\Phi }_2(X_1,X_2,X_3)=\text{tr}(\mathrm{\Psi }_2(X_1,X_2)X_3)=\text{tr}(X_1X_2X_3)+\text{tr}(X_2X_1X_3)$$
$$\text{tr}(X_1)\text{tr}(X_2X_3)\text{tr}(X_2)\text{tr}(X_1X_3)+\text{tr}(X_1)\text{tr}(X_2)\text{tr}(X_3)\text{tr}(X_1X_2)\text{tr}(X_3).$$
If we delete the symbols of traces and the $`X`$’s in the above expression, we shall obtain the following linear combination of permutations
$$(123)+(213)(1)(23)(2)(13)+(1)(2)(3)(12)(3)=\underset{\sigma S_3}{}\text{sign}(\sigma ).$$
This suggests the following construction. We write the permutations in the symmetric group $`S_m`$ as products of disjoint cycles, including the cycles of length 1,
$$\sigma =(i_1\mathrm{}i_p)(j_1\mathrm{}j_q)\mathrm{}(k_1\mathrm{}k_r).$$
We define the associated trace function
$$\text{tr}_\sigma (x_1,\mathrm{},x_m)=\text{tr}(x_{i_1}\mathrm{}x_{i_p})\text{tr}(x_{j_1}\mathrm{}x_{j_q})\mathrm{}\text{tr}(x_{k_1}\mathrm{}x_{k_r}).$$
For every element
$$\underset{\sigma S_m}{}\alpha _\sigma \sigma KS_m,\alpha _\sigma K,\sigma S_m,$$
where $`KS_m`$ is the group algebra of $`S_m`$, we define the trace polynomial
$$f(x_1,\mathrm{},x_m)=\underset{\sigma S_m}{}\alpha _\sigma \text{tr}_\sigma (x_1,\mathrm{},x_m).$$
We also assume that for $`mk`$ the symmetric group $`S_m`$ acts on $`1,\mathrm{},m`$ and leaves invariant $`m+1,\mathrm{},k`$, i.e. $`S_m`$ is canonically embedded into $`S_k`$.
###### Theorem 2.11.
(The Second Fundamental Theorem of Matrix Invariants, Razmyslov \[R2\], Procesi \[P2\]) Let
$$f(x_1,\mathrm{},x_m)=\underset{\sigma S_m}{}\alpha _\sigma \text{tr}_\sigma (x_1,\mathrm{},x_m),\alpha _\sigma K,$$
be a multilinear trace polynomial of degree $`m`$. Then $`f=0`$ is a trace identity for the $`n\times n`$ matrix algebra, i. e. $`f(a_1,\mathrm{},a_m)=0`$ for all $`a_1,\mathrm{},a_mM_n(K)`$, if and only if
$$\underset{\sigma S_m}{}\alpha _\sigma \sigma $$
belongs to the two-sided ideal $`J(n,m)`$ of the group algebra $`KS_m`$ generated by the element
$$\underset{\sigma S_{n+1}}{}\text{sign}(\sigma )\sigma .$$
As in the case of $`2\times 2`$ matrices, the fundamental trace identity
$$\underset{\sigma S_{n+1}}{}(\text{sign }\sigma )\text{tr}_\sigma (x_1,\mathrm{},x_{n+1})=0$$
is actually the linearization of the Cayley-Hamilton polynomial.
There are several important objects related with invariant theory of matrices. As above, $`n,d2`$ are fixed integers and $`X_1,\mathrm{},X_d`$ are $`d`$ generic $`n\times n`$ matrices:
The algebra $`R_{nd}`$ generated by $`X_1,\mathrm{},X_d`$;
The pure (or commutative) trace algebra $`C_{nd}=\mathrm{\Omega }_{nd}^{GL_n}`$ generated by the traces of all products $`\text{tr}(X_{i_1}\mathrm{}X_{i_k})`$;
The mixed (or noncommutative) trace algebra $`T_{nd}`$ generated by $`R_{nd}`$ and $`C_{nd}`$ regarding the elements of $`C_{nd}`$ as scalar matrices;
The field of fractions $`Q(C_{nd})`$ of the algebra $`C_{nd}`$;
The algebra $`Q(C_{nd})R_{nd}`$.
All these algebras have no zero divisors and play important roles in mathematics. See the books \[DF, F4, J, P1\] for different aspects of the theory of algebras of matrix invariants and their applications to combinatorial and structure theory of PI-algebras, central division algebras, etc.
The algebra $`R_{nd}`$ is a well known object in the theory of PI-algebras, or algebras with polynomial identities. Let $`Kx_1,\mathrm{},x_d`$ be the free associative algebra (i.e. the algebra of polynomials in noncommuting variables). A polynomial $`f(x_1,\mathrm{},x_d)Kx_1,\mathrm{},x_d`$ is a polynomial identity in $`d`$ variables for $`M_n(K)`$ if $`f(a_1,\mathrm{},a_d)=0`$ for all $`a_1,\mathrm{},a_dM_n(K)`$. The set $`I(M_n(K))`$ of all polynomial identities in $`d`$ variables is a two-sided ideal of $`Kx_1,\mathrm{},x_d`$ and $`R_{nd}`$ is isomorphic to the factor algebra $`Kx_1,\mathrm{},x_d/I(M_n(K))`$.
Clearly, the algebra $`C_{nd}`$ is the algebra of matrix invariants.
The algebra $`T_{nd}`$ is also the algebra of invariant polynomial functions under a suitable action of $`GL_n`$. It is called the algebra of matrix concominants. It is a finitely generated $`C_{nd}`$-module and, as a $`C_{nd}`$-module, has a generating set consisting of products $`X_{j_1}\mathrm{}X_{j_k}`$, where $`k<N(n)`$, the class of nilpotency in the Nagata-Higman theorem.
The field of fractions $`Q(C_{nd})`$ appears naturally in field theory. One of the main problems related with $`Q(C_{nd})`$ is whether it is a purely transcendent extension of $`K`$.
Finally, $`Q(C_{nd})R_{nd}`$ is a central division algebra of dimension $`n^2`$ over its centre $`Q(C_{nd})`$ and serves as a source of counterexamples to the theory of central division algebras.
General invariant theory gives that $`C_{nd}`$ and $`T_{nd}`$ have nice algebraic properties.
###### Theorem 2.12.
(Van den Bergh, \[VB1\]) The algebra $`C_{nd}`$ is a Cohen-Macaulay and even Gorenstein unique factorization domain. The algebra $`T_{nd}`$ is a Cohen-Macaulay module over $`C_{nd}`$.
Recall that the Noether normalization theorem gives that $`C_{nd}`$ contains a homogeneous set of algebraically independent elements $`\{a_1,\mathrm{},a_k\}`$, where $`k=(d1)n^2+1`$ is the transcendence degree of the quotient field $`Q(C_{nd})`$ ($`k`$ is also equal to the Krull dimension of $`C_{nd}`$), such that $`C_{nd}`$ is integral over the polynomial algebra $`K[a_1,\mathrm{},a_k]`$. Such a set $`\{a_1,\mathrm{},a_k\}`$ is called a homogeneous system of parameters for $`C_{nd}`$. By a result of Stanley \[St1\], a graded $`C_{nd}`$-module is Cohen-Macaulay if and only if it is a free module with respect to some homogeneous system of parameters $`\{a_1,\mathrm{},a_k\}`$ of $`C_{nd}`$.
The algebras $`R_{nd},C_{nd}`$, and $`T_{nd}`$ are multigraded and their homogeneous components of degree $`(k_1,\mathrm{},k_d)`$ consist of all polynomials which are homogeneous of degree $`k_i`$ in the generic matrix $`X_i`$. Hence we may consider their Hilbert series in $`d`$ variables. For example,
$$H(C_{nd},t_1,\mathrm{},t_d)=\text{dim}\left(C_{nd}^{(k_1,\mathrm{},k_d)}\right)t_1^{k_1}\mathrm{}t_d^{k_d},$$
where $`C_{nd}^{(k_1,\mathrm{},k_d)}`$ is the homogeneous component of degree $`(k_1,\mathrm{},k_d)`$. Le Bruyn \[LB1\] for the case $`n=2`$, and Formanek \[F3\] and Teranishi \[T1, T3\] in the general case proved:
###### Theorem 2.13.
Let $`d2`$ for $`n3`$ and $`d>2`$ for $`n=2`$. Then $`H(C_{nd},t_1,\mathrm{},t_d)`$ and $`H(T_{nd},t_1,\mathrm{},t_d)`$ satisfy the functional equation
$$H(C_{nd},t_1^1,\mathrm{},t_d^1)=(1)^k(t_1\mathrm{}t_d)^{n^2}H(t_1,\mathrm{},t_d),$$
where $`k=(d1)n^2+1`$ is the Krull dimension of $`C_{nd}`$, and similarly for the Hilbert series of $`T_{nd}`$.
The proofs of this theorem given by Formanek and Teranishi are quite different and use, respectively, representation theory of general linear groups and the Molien-Weyl integral formula. Later, Van den Bergh paid attention that the proof can be considerably simplified using results of Stanley on Hilbert series of Cohen-Macaulay algebras.
We need some background on symmetric polynomials and representation theory of the general linear group, see e.g. the books by Weyl \[W2\] and Macdonald \[Mc\]. As in the case of polynomial algebras, the general linear group $`GL_d`$ acts canonically on the vector space with basis $`\{X_1,\mathrm{},X_d\}`$. If $`g=(g_{ij})`$, $`g_{ij}K`$, then the action of $`g`$ on $`X_j`$ is defined by
$$g(X_j)=g_{1j}X_1+\mathrm{}+g_{dj}X_d.$$
This action is extended diagonally on $`R_{nd},C_{nd},T_{nd}`$. If $`f(X_1,\mathrm{},X_d)`$ is any polynomial expression depending on $`X_1,\mathrm{},X_d`$ (maybe including also traces), then
$$g(f(X_1,\mathrm{},X_d))=f(g(X_1),\mathrm{},g(X_d)),gGL_d.$$
Representation theory of $`GL_d`$ says that every submodule of the $`GL_d`$-modules $`R_{nd},C_{nd},T_{nd}`$ is a direct sum of irreducible (or simple) submodules. The irreducible $`GL_d`$-submodules which appear in the decomposition are polynomial modules and are described in terms of partitions of integers. If
$$\lambda =(\lambda _1,\mathrm{},\lambda _d),\lambda _1\mathrm{}\lambda _d0,$$
is a partition of $`k`$ (notation $`\lambda k`$) in not more than $`d`$ parts, then we denote the related $`GL_d`$-module by $`W(\lambda )=W(\lambda _1,\mathrm{},\lambda _d)`$. To be explicit, below we consider the case of $`C_{nd}`$ only. If
$$C_{nd}=m(\lambda )W(\lambda ),m(\lambda )0,$$
i.e. there are $`m(\lambda )`$ direct summands isomorphic to $`W(\lambda )`$, then we say that $`W(\lambda )`$ appears with multiplicity $`m(\lambda )`$. The multiplicities $`m(\lambda )`$ for $`R_{nd},C_{nd}`$, and $`T_{nd}`$ play important role in the quantitative study of polynomial identities of matrices. (See the survey by Regev \[Re\] and the book of the author \[D2\] for applications of representation theory of $`S_n`$ and $`GL_d`$ to the theory of PI-algebras.) The Hilbert series of $`C_{nd}`$ has the form
$$H(C_{nd},t_1,\mathrm{},t_d)=m(\lambda )S_\lambda (t_1,\mathrm{},t_d),$$
where $`S_\lambda (t_1,\mathrm{},t_d)`$ is the Schur function associated to $`\lambda `$. Schur functions are important combinatorial objects and appear in many places in mathematics. For example, they form a basis of the vector space of all symmetric polynomials in $`d`$ variables. One of the possible ways to define Schur functions is via Vandermonde-like determinants. For a partition $`\mu =(\mu _1,\mathrm{},\mu _d)`$, define the determinant
$$V(\mu _1,\mathrm{},\mu _d)=\left|\begin{array}{cccc}t_1^{\mu _1}& t_2^{\mu _1}& \mathrm{}& t_d^{\mu _1}\\ t_1^{\mu _2}& t_2^{\mu _2}& \mathrm{}& t_d^{\mu _2}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ t_1^{\mu _d}& t_2^{\mu _d}& \mathrm{}& t_d^{\mu _d}\end{array}\right|.$$
Then the Schur function is
$$S_\lambda (t_1,\mathrm{},t_d)=\frac{V(\lambda _1+d1,\lambda _2+d2,\mathrm{},\lambda _{d1}+1,\lambda _d)}{V(d1,d2,\mathrm{},1,0)}.$$
The Schur functions play the role of characters of the corresponding representation of $`GL_d`$. If we know the Hilbert series $`H(C_{nd},t_1,\mathrm{},t_d)`$, then we can uniquely determine the $`GL_d`$-module structure of $`C_{nd}`$.
## 3. Concrete computations
We give a survey of some concrete results about generators, defining relations and Hilbert series of the algebras of matrix invariants. We start with $`2\times 2`$ matrices.
The first important reduction is the following. We take the generic $`2\times 2`$ matrix $`X`$ and present it in the form
$$X=\left(\begin{array}{cc}x_{11}& x_{12}\\ x_{21}& x_{22}\end{array}\right)=\left(\begin{array}{cc}\frac{x_{11}+x_{22}}{2}& 0\\ 0& \frac{x_{11}+x_{22}}{2}\end{array}\right)+\left(\begin{array}{cc}\frac{x_{11}x_{22}}{2}& x_{12}\\ x_{21}& \frac{x_{11}x_{22}}{2}\end{array}\right)$$
$$=\frac{\text{tr}(X)}{2}E+\left(\begin{array}{cc}y_{11}& y_{12}\\ y_{21}& y_{11}\end{array}\right)=\frac{\text{tr}(X)}{2}E+Y,$$
where $`E`$ is the identity matrix and
$$Y=\left(\begin{array}{cc}y_{11}& y_{12}\\ y_{21}& y_{11}\end{array}\right)$$
is a generic $`2\times 2`$ traceless matrix. This reduction implies that $`C_{2d}`$ is generated by $`\text{tr}(X_i)`$, $`i=1,\mathrm{},d`$, and $`\text{tr}(Y_{i_1}\mathrm{}Y_{i_k})`$, where $`Y_1,\mathrm{},Y_d`$ are generic traceless matrices and $`k2`$. Since the class of nilpotency in the Nagata-Higman theorem is $`N(2)=3`$ for $`n=2`$, we obtain that it is sufficient to consider the cases $`k=2,3`$ only. Similarly, $`T_{2d}`$ is generated by $`C_{2d}`$ and $`Y_1,\mathrm{},Y_d`$. It is also easy to see:
###### Proposition 3.1.
The algebra $`C_{2d}`$ has the presentation
$$C_{2d}=\left(K[\text{tr}(X_1),\mathrm{},\text{tr}(X_d)]\right)_KC_0,$$
where
$$C_0=K[\text{tr}(Y_iY_j),\text{tr}(Y_pY_qY_r)1ijd,1p<q<rd]/I$$
is the algebra generated by products of generic traceless $`2\times 2`$ matrices. Hence we may choose all defining relations as linear combinations of the traces of products of traceless matrices.
The description of $`T_{2d}`$ is easier than that of $`C_{2d}`$.
###### Theorem 3.2.
(i) (Procesi \[P4\]) The noncommutative trace algebra $`T_{2d}`$ is isomorphic to the tensor product of $`K`$-algebras
$$K[\text{tr}(X_1),\mathrm{},\text{tr}(X_m)]_KW_d,$$
where $`W_d`$ is the associative algebra generated by the generic traceless matrices $`Y_1,\mathrm{},Y_m`$.
(ii) (Razmyslov \[R1\]) The algebra of the generic traceless matrices $`W_d`$ is isomorphic to the factor algebra $`Ky_1,\mathrm{},y_d/I(M_2(K),sl_2(K))`$ of the free associative algebra $`Ky_1,\mathrm{},y_d`$, where $`sl_2(K)`$ is the Lie algebra of all traceless $`2\times 2`$ matrices and $`I(M_2(K),sl_2(K))`$ is the ideal of all polynomials in $`Ky_1,\mathrm{},y_d`$ which vanish under the substitutions of $`y_i`$ by elements in $`sl_2(K)`$. (Such polynomials are called weak polynomial identities for the pair $`(M_2(K),sl_2(K))`$.) The ideal $`I(M_2(K),sl_2(K))`$ is generated as a weak T-ideal by the weak polynomial identity $`[x_1^2,x_2]=0`$, i.e. $`I(M_2(K),sl_2(K))`$ is the minimal ideal of weak polynomial identities containing the element $`[x_1^2,x_2]`$. An equivalent description is that as an ideal of the free algebra, $`I(M_2(K),sl_2(K))`$ is generated by all elements $`[uv+vu,w]`$, where $`u,v,w`$ are all possible commutators $`[y_{j_1},\mathrm{},y_{j_k}]`$ in the variables $`y_1,\mathrm{},y_d`$.
(iii) (Drensky and Koshlukov \[DK\]) The ideal $`I(M_2(K),sl_2(K))`$ is the minimal ideal of the free associative algebra $`Ky_1,\mathrm{},y_d`$ which is invariant under the diagonal action of $`GL_d`$ and contains the elements $`[y_1^2,y_2]`$ and
$$s_4(y_1,y_2,y_3,y_4)=\underset{\sigma S_4}{}\text{sign}(\sigma )y_{\sigma (1)}y_{\sigma (2)}y_{\sigma (3)}y_{\sigma (4)},$$
the second polynomial appears for $`d4`$ only. Hence the algebra of $`2\times 2`$ generic traceless matrices has a uniform set of defining relations for any $`d2`$.
(iv) (Procesi \[P4\]) As a $`GL_d`$-module $`W_d`$ has the description
$$W_dW(\lambda _1,\lambda _2,\lambda _3),$$
where the sum is on all partitions $`\lambda `$ in at most three parts.
(v) (Formanek \[F1\]) The $`GL_d`$-module $`T_{2d}`$ has the decomposition
$$T_{2d}(\lambda _1\lambda _2+1)(\lambda _2\lambda _3+1)(\lambda _3\lambda _4+1)W(\lambda _1,\lambda _2,\lambda _3,\lambda _4).$$
(vi) (See \[D1, D2, F1, P4, LB1, LB2\] for other descriptions of the Hilbert series.) The Hilbert series of $`T_{2d}`$ is
$$H(T_{2d},t_1,\mathrm{},t_d)=\underset{i=1}{\overset{d}{}}\frac{1}{1t_i}S_{(\lambda _1,\lambda _2,\lambda _3)}(t_1,\mathrm{},t_d).$$
Now we give some results on $`C_{2d}`$. We follow the way we used in the previous theorem. For more details, especially for the Hilbert series and the $`GL_d`$-module decomposition of $`C_{2d}`$, see \[F1, LB2, P4\] or \[DF\].
###### Theorem 3.3.
(i) The commutative trace algebra $`C_{2d}`$ is isomorphic to the tensor product of $`K`$-algebras
$$K[\text{tr}(X_1),\mathrm{},\text{tr}(X_m)]_KC(W_d),$$
where $`C(W_d)`$ is the centre of the algebra $`W_d`$ defined above.
(ii) As a subalgebra of $`W_d`$, its centre $`C(W_d)`$ is generated by
$$Y_i^2,i=1,\mathrm{},d,Y_iY_j+Y_jY_i,1i<jd,$$
$$s_3(Y_i,Y_j,Y_k)=\underset{\sigma S_3}{}\text{sign}(\sigma )Y_{\sigma (1)}Y_{\sigma (2)}Y_{\sigma (3)},1i<j<kd,$$
where the symmetric group $`S_3`$ acts on $`\{i,j,k\}`$.
(iii) As a $`GL_d`$-module $`W_d`$ has the decomposition
$$W_dW(\lambda _1,\lambda _2,\lambda _3),$$
where the sum is on all partitions $`\lambda `$ in at most three parts such that both $`\lambda _1\lambda _2`$ and $`\lambda _2\lambda _3`$ are even.
Concerning the defining relations of $`C_{2d}`$, the case $`d=2`$ is trivial. Formanek, Halpin and Li \[FHL\] showed that $`C_{22}`$ is generated by the algebraically independent elements
$$\text{tr}(X_1),\text{tr}(X_2),\text{det}(X_1),\text{det}(X_2),\text{tr}(X_1X_2).$$
For $`d=3`$ Sibirskii \[Si\] found one relation between the generators of $`C_{23}`$ and, using the Hilbert series of $`C_{23}`$, Formanek \[F1\] proved that there are no more relations. In the general case, the description of the defining relations of $`C_{2d}`$ is reduced to a similar description of the defining relations of the subalgebra of $`C_{2d}`$ generated by
$$\text{tr}(Y_i^2),i=1,\mathrm{},d,\text{tr}(Y_iY_j),1i<jd,$$
$$\text{tr}(Y_iY_jY_k),1i<j<kd,$$
where $`Y_1,\mathrm{},Y_d`$ are generic traceless $`2\times 2`$ matrices. Since $`GL_2`$ acts on the generic matrices by conjugation, we may replace its action with the action of $`SL_2`$ and even with the action of $`PSL_2`$. Since $`PSL_2()SO_3()`$, the special orthogonal group, we may apply invariant theory of special linear groups. (The restriction $`K=`$ is not essential in the final version of the result.)
We consider the action of the special orthogonal group $`SO_3=SO_3(K)`$, i.e. the group of orthogonal $`3\times 3`$ matrices with determinant 1, on the polynomial algebra in $`3d`$ variables
$$K[u_i^{(1)},u_i^{(2)},u_i^{(3)}i=1,\mathrm{},d],$$
induced by the action of $`SO_3`$ on the three-dimensional vectors $`u_i=(u_i^{(1)},u_i^{(2)},u_i^{(3)})`$. It is a classical result that the algebra of invariants $`K[u_i^{(j)}]^{SO_3}`$ is generated by all scalar products
$$u_i,u_j=u_i^{(1)}u_j^{(1)}+u_i^{(2)}u_j^{(2)}+u_i^{(3)}u_j^{(3)},1ijd,$$
and all $`3\times 3`$ determinants of the coordinates
$$\mathrm{\Delta }(u_i,u_j,u_k)=\left|\begin{array}{ccc}u_i^{(1)}& u_j^{(1)}& u_k^{(1)}\\ u_i^{(2)}& u_j^{(2)}& u_k^{(2)}\\ u_i^{(3)}& u_j^{(3)}& u_k^{(3)}\end{array}\right|,1i<j<kd.$$
The defining relations express the fact that the underlying vector space is three-dimensional and every four vectors are linearly dependent. In particular, they use the properties of the Gram determinant:
$$\mathrm{\Gamma }_4(u_i,u_j,u_k,u_l;u_p,u_q,u_r,u_s)=\left|\begin{array}{cccc}u_i,u_p& u_i,u_q& u_i,u_r& u_i,u_s\\ u_j,u_p& u_j,u_q& u_j,u_r& u_j,u_s\\ u_k,u_p& u_k,u_q& u_k,u_r& u_k,u_s\\ u_l,u_p& u_l,u_q& u_l,u_r& u_l,u_s\end{array}\right|=0,$$
$$1i<j<k<ld,1p<q<r<sd,$$
$$\mathrm{\Delta }(u_i,u_j,u_k)\mathrm{\Delta }(u_p,u_q,u_r)\mathrm{\Gamma }_3(u_i,u_j,u_k;u_p,u_q,u_r)=0,$$
$$u_p,u_i\mathrm{\Delta }(u_j,u_k,u_l)u_p,u_j\mathrm{\Delta }(u_i,u_k,u_l)$$
$$+u_p,u_k\mathrm{\Delta }(u_i,u_j,u_l)u_p,u_l\mathrm{\Delta }(u_i,u_j,u_k)=0.$$
In order to apply invariant theory of $`SO_3`$ we need a scalar product (i.e. nondegenerate symmetric bilinear form) on $`sl_2(K)`$. We use the trace and define
$$u,v=\text{tr}(uv),u,vsl_2(K).$$
The following result gives the generators and the defining relations of $`C_{2d}`$ for $`d2`$. It is a translation of the description of the invariants of $`SO_3`$.
###### Theorem 3.4.
(i) The algebra $`C_{2d}`$ is generated by
$$\text{tr}(X_i),\text{tr}(Y_i^2),\text{tr}(Y_iY_j),\text{tr}(s_3(Y_i,Y_j,Y_k)),$$
where $`i,j,k=1,\mathrm{},d`$, and in the traces involving two or three traceless matrices we require $`i<j`$ or $`i<j<k`$, respectively.
(ii) Drensky \[D3\]) The defining relations of $`C_{2d}`$ with respect to the above generators are
$$\text{tr}(s_3(Y_i,Y_j,Y_k))\text{tr}(s_3(Y_p,Y_q,Y_r))+18\left|\begin{array}{ccc}\text{tr}(Y_iY_p)& \text{tr}(Y_iY_q)& \text{tr}(Y_iY_r)\\ \text{tr}(Y_jY_p)& \text{tr}(Y_jY_q)& \text{tr}(Y_jY_r)\\ \text{tr}(Y_kY_p)& \text{tr}(Y_kY_q)& \text{tr}(Y_kY_r)\end{array}\right|=0,$$
$$\text{tr}(Y_pY_i)\text{tr}(s_3(Y_j,Y_k,Y_l))\text{tr}(Y_pY_j)\text{tr}(s_3(Y_i,Y_k,Y_l))$$
$$+\text{tr}(Y_pY_k)\text{tr}(s_3(Y_i,Y_j,Y_l))\text{tr}(Y_pY_l)\text{tr}(s_3(Y_i,Y_j,Y_k))=0,$$
where, again, $`i,j,k,p,q,r=1,\mathrm{},d`$, and, where necessary, we require $`i<j<k<l`$ and $`p<q<r`$.
In order to work efficiently with an algebra $`R=K[x_1,\mathrm{},x_p]/I`$, it is not sufficient to know the generators of the ideal $`I`$. For computational purposes one needs also the Gröbner basis of $`I`$ with respect to some ordering on the monomials of $`K[x_1,\mathrm{},x_p]`$, see e.g. the book by Adams and Loustaunau \[AL\]. The Gröbner basis of $`C_{2d}`$ is given by Domokos and Drensky \[DD\], see their paper for more details.
Now we consider two generic $`3\times 3`$ matrices. Using the Molien-Weyl formula, Teranishi \[T1\] calculated the Hilbert series of $`C_{32}`$, namely,
$$H(C_{32},t_1,t_2)=\frac{1+t_1^3t_2^3}{(1t_1)(1t_2)q_2(t_1,t_2)q_3(t_1,t_2)(1t_1^2t_2^2)},$$
where
$$q_2(t_1,t_2)=(1t_1^2)(1t_1t_2)(1t_2^2),$$
$$q_3(t_1,t_2)=(1t_1^3)(1t_1^2t_2)(1t_1t_2^2)(1t_2^3).$$
He also found the following system of generators of $`C_{32}`$:
$$\begin{array}{c}\text{tr}(X_1),\text{tr}(X_2),\text{tr}(X_1^2),\text{tr}(X_1X_2),\text{tr}(X_2^2),\\ \\ \text{tr}(X_1^3),\text{tr}(X_1^2X_2),\text{tr}(X_1X_2^2),\text{tr}(X_2^3),\text{tr}(X_1^2X_2^2),\text{tr}(X_1^2X_2^2X_1X_2),\end{array}$$
where $`X_1,X_2`$ are generic $`3\times 3`$ matrices. He showed that the first ten of these generators form a homogeneous system of parameters of $`C_{32}`$ and $`C_{32}`$ is a free module with generators 1 and $`\text{tr}(X_1^2X_2^2X_1X_2)`$ over the polynomial algebra on these ten elements.
Abeasis and Pittaluga \[AP\] found a system of generators of $`C_{3d}`$ in terms of representation theory of the symmetric and general linear groups, in the spirit of its use in theory of PI-algebras. They showed that $`C_{3d}`$ has a minimal system of generators which spans a $`GL_d`$-module isomorphic to
$$G=W(1)W(2)W(3)W(1^3)W(2^2)W(2,1^2)$$
$$W(3,1^2)W(2^2,1)W(1^5)W(3^2)W(3,1^3).$$
(The partitions in \[AP\] are given in “Francophone” way, i.e., transposed to ours.)
It follows from the description of the generators of $`C_{32}`$ given by Teranishi \[T1\], that $`\text{tr}(X_1^2X_2^2X_1X_2)`$ satisfies a quadratic equation with coefficients depending on the other ten generators. The explicit (but very complicated) form of the equation was found by Nakamoto \[N\], over $``$, with respect to a slightly different system of generators. A much simpler description of $`C_{32}`$ was obtained by Aslaksen, Drensky, and Sadikova \[ADS\]. The following generators are in the spirit of the ideas of \[AP\].
###### Proposition 3.5.
Let $`X_1,X_2`$ and $`Y_1,Y_2`$ be, respectively, two generic and two generic traceless $`3\times 3`$ matrices. The algebra $`C_{32}`$ is generated by
$$\begin{array}{c}\text{tr}(X_1),\text{tr}(X_2),\text{tr}(Y_1^2),\text{tr}(Y_1Y_2),\text{tr}(Y_2^2),\\ \\ \text{tr}(Y_1^3),\text{tr}(Y_1^2Y_2),\text{tr}(Y_1Y_2^2),\text{tr}(Y_2^3),\\ \\ V=\text{tr}(Y_1^2Y_2^2)\text{tr}(Y_1Y_2Y_1Y_2),W=\text{tr}(Y_1^2Y_2^2Y_1Y_2)\text{tr}(Y_2^2Y_1^2Y_2Y_1).\end{array}$$
Now we define the following elements of $`C_{32}`$:
$$U=\left|\begin{array}{cc}\text{tr}(Y_1^2)& \text{tr}(Y_1Y_2)\\ \text{tr}(Y_1Y_2)& \text{tr}(Y_2^2)\end{array}\right|,$$
$$W_1=U^3,W_2=U^2V,W_4=UV^2,W_7=V^3,$$
$$W_5=V\left|\begin{array}{ccc}\text{tr}(Y_1^2)& \text{tr}(Y_1Y_2)& \text{tr}(Y_2^2)\\ \text{tr}(Y_1^3)& \text{tr}(Y_1^2Y_2)& \text{tr}(Y_1Y_2^2)\\ \text{tr}(Y_1^2Y_2)& \text{tr}(Y_1Y_2^2)& \text{tr}(Y_2^3)\end{array}\right|,$$
$$W_6=\left|\begin{array}{cc}\text{tr}(Y_1^3)& \text{tr}(Y_1Y_2^2)\\ \text{tr}(Y_1^2Y_2)& \text{tr}(Y_2^3)\end{array}\right|^24\left|\begin{array}{cc}\text{tr}(Y_2^3)& \text{tr}(Y_1Y_2^2)\\ \text{tr}(Y_1Y_2^2)& \text{tr}(Y_1^2Y_2)\end{array}\right|\left|\begin{array}{cc}\text{tr}(Y_1^3)& \text{tr}(Y_1^2Y_2)\\ \text{tr}(Y_1^2Y_2)& \text{tr}(Y_1Y_2^2)\end{array}\right|,$$
$$W_3^{}=U\left|\begin{array}{ccc}\text{tr}(Y_1^2)& \text{tr}(Y_1Y_2)& \text{tr}(Y_2^2)\\ \text{tr}(Y_1^3)& \text{tr}(Y_1^2Y_2)& \text{tr}(Y_1Y_2^2)\\ \text{tr}(Y_1^2Y_2)& \text{tr}(Y_1Y_2^2)& \text{tr}(Y_2^3)\end{array}\right|,$$
where $`U,V`$ are defined above. Finally, we define one more element $`W_3^{\prime \prime }`$ as follows. Recall that a linear mapping $`\delta `$ of an algebra $`R`$ is a derivation if $`\delta (rs)=\delta (r)s+r\delta (s)`$ for all $`r,sR`$. We consider the derivation $`\delta `$ of $`C_{32}`$ which commutes with the trace and satisfies the conditions
$$\delta (X_1)=0,\delta (X_2)=X_1,\delta (Y_1)=0,\delta (Y_2)=Y_1.$$
Then
$$W_3^{\prime \prime }=\frac{1}{144}\underset{i=0}{\overset{6}{}}(1)^i\delta ^i(\text{tr}^3(Y_2^2))\delta ^{6i}(\text{tr}^2(Y_2^3)).$$
The following theorem gives the defining relation of $`C_{32}`$. It uses representation theory of $`GL_2`$, combinatorics, computations by hand and easy computer calculations with standard functions of Maple.
###### Theorem 3.6.
(Aslaksen, Drensky, Sadikova \[ADS\]) The algebra of invariants $`C_{32}`$ of two $`3\times 3`$ matrices is generated by the elements from the previous theorem, subject to the defining relation
$$W^2\left(\frac{1}{27}W_1\frac{2}{9}W_2+\frac{4}{15}W_3^{}+\frac{1}{90}W_3^{\prime \prime }+\frac{1}{3}W_4\frac{2}{3}W_5\frac{1}{3}W_6\frac{4}{27}W_7\right)=0.$$
The calculation of the Hilbert series of $`C_{nd}`$ and $`T_{nd}`$ based on the Molien-Weyl formula is quite complicated because requires evaluations of multiple integrals. Van den Bergh \[VB2\] sujected a way which involves graph theory. As a consequence, he established important properties of $`H(C_{2d},t_1,\mathrm{},t_d)`$ and $`H(T_{2d},t_1,\mathrm{},t_d)`$.
Berele and Stembridge \[BS\] applied the method of van den Bergh \[VB2\] and calculated the Hilbert series of $`T_{32}`$. Using the above results of Aslaksen, Drensky, and Sadikova on $`C_{32}`$ and the explicit form of the Hilbert series of $`T_{32}`$, Benanti and Drensky \[BD\] found a polynomial subalgebra $`S`$ of $`C_{32}`$ and a finite set of generators of the free $`S`$-module $`T_{32}`$. They gave also a set of defining relations of $`T_{32}`$ as an algebra and a Gröbner basis of the corresponding ideal. (See the survey article by Ufnarovski \[U\] for a background on Gröbner bases in the noncommutative case, as well as the paper by Mikhalev and Zolotykh \[MZ\] which is closer to the situation in \[BD\].)
For two generic $`4\times 4`$ matrices, the Hilbert series of $`C_{42}`$ was calculated (with some typos) by Teranishi \[T1, T2\] and corrected by Berele and Stembridge \[BS\]. Teranishi found also a homogeneous system of parameters and a system of generators, in the spirit of the $`3\times 3`$ case. Recently, Drensky and Sadikova \[DS\] have found another system of generators of $`C_{42}`$ which is minimal and seems to be more convenient for concrete calculations.
The Hilbert series of a graded vector space with $`GL_d`$-module structure determines uniquely its decomposition into irreducible submodules. Hence, in principle, one may calculate the multiplicities $`m(\lambda )`$ if one knows the concrete form of the Hilbert series. Berele \[B1\] used the Hilbert series of $`C_{32}`$ found by Teranishi \[T1\] and described the asymptotics of $`m(\lambda _1,\lambda _2)`$. (Due to a technical error (an omitted summand) some of the coefficients of the polynomials in the asymptotics of Berele are slightly different from the real ones.) Another approach to the problem was suggested by Drensky and Genov \[DG1\]. Let
$$f(t_1,t_2)=\underset{i,j0}{}a_{ij}t_1^it_2^j,$$
$`a_{ij}K`$, $`a_{ij}=a_{ji}`$, be a symmetric function in two variables which is a formal power series from $`K[[t_1,t_2]]`$. We present it in the form
$$f(t_1,t_2)=\underset{\lambda _1\lambda _2}{}m(\lambda _1,\lambda _2)S_{(\lambda _1,\lambda _2)}(t_1,t_2)$$
and want to find the multiplicities $`m(\lambda _1,\lambda _2)`$. In most of the cases which we consider, $`f(t_1,t_2)`$ is given explicitly as a rational function. So, it is natural to express $`m(\lambda _1,\lambda _2)`$ not in terms of the coefficients $`a_{ij}`$ but in a more direct way. We introduce the generating function of the multiplicities
$$M(f,t,u)=\underset{\lambda _1\lambda _2}{}m(\lambda _1,\lambda _2)t^{\lambda _1}u^{\lambda _2}K[[t,u]]$$
and call it the multiplicity series of $`f(t_1,t_2)`$. It is more convenient to introduce a new variable $`v=tu`$ and to consider the series
$$M^{}(f,t,v)=\underset{\lambda _1\lambda _2}{}m(\lambda _1,\lambda _2)t^{\lambda _1\lambda _2}v^{\lambda _2}K[[t,v]],$$
because the mapping $`M^{}:K[[t_1,t_2]]^{S_2}K[[t,v]]`$ is a bijective linear mapping which is continuous with respect to the formal power series topology. It is easy to see that $`f(t_1,t_2)`$ and $`M^{}(f,t,v)`$ are related by
$$f(t_1,t_2)=\frac{t_1M^{}(f,t_1,t_1t_2)t_2M^{}(f,t_2,t_1t_2)}{t_1t_2}.$$
Hence, if we have a potential candidate $`h(t,v)`$ for $`M^{}(f,t,v)`$, it is easy to verify whether $`h(t,v)=M^{}(f,t,v)`$.
Also, the elementary symmetric function $`e_2=t_1t_2`$ behaves like a constant,
$$M^{}(g(t_1t_2)f(t_1,t_2),t,v)=g(v)M^{}(f,t,v),$$
and this simplifies the calculations. Applying quite complicated (also technically) arguments, Drensky and Genov \[DG1\] found the multiplicity series of the Hilbert series of $`C_{32}`$. They also corrected the technical errors in \[B1\].
###### Theorem 3.7.
(i) \[DG1\] The multiplicity series of the Hilbert series of the algebra $`C_{32}`$ of invariants of two $`3\times 3`$ matrices is
$$M^{}(H(C_{32},t_1,t_2),t,v)=\frac{1}{(1v^2)(1v^3)^2}\times $$
$$\times (\frac{(1+v^2+v^4)((1+v^2)(1t^2v)+2tv(1v))}{3(1v)(1v^2)^3(1t)^2(1t^2)}+$$
$$+\frac{(1v)(1+tv)}{3(1v^2)(1t)(1t^2)}+\frac{(1v^2)(1tv)}{3(1v^3)(1t^3)}$$
$$\frac{v^3((1v+v^2)(1t^2v^2)+tv(1v^2))}{(1v)(1v^2)^2(1v^4)(1t)(1t^2)(1tv)}).$$
(ii) \[B1, DG1\] Let $`\lambda =(p,q)`$ and let $`m(p,q)`$ be the multiplicity of $`S_{(p,q)}(t_1,t_2)`$ in $`H(C_{32},t_1,t_2)`$. Then for $`p>2q0`$
$$m(p,q)=\frac{q^7}{7!2^5.3^2}+\frac{(pq)q^6}{6!2^4.3^2}+\frac{(pq)^2q^5}{2!5!2^33^2}+𝒪((p+q)^6)$$
$$=\frac{p^2q^5}{17280}\frac{11pq^6}{103680}+\frac{71q^7}{1451520}+𝒪((p+q)^6);$$
for $`2qpq0`$
$$m(p,q)=\frac{q^7}{7!2^5.3^2}+\frac{(pq)q^6}{6!2^4.3^2}+\frac{(pq)^2q^5}{2!5!2^33^2}\frac{(2qp)^7}{7!2^5.3^2}+𝒪((p+q)^6)$$
$$=\frac{p^7}{1451520}\frac{p^6q}{103680}+\frac{p^5q^2}{17280}\frac{p^4q^3}{5184}+\frac{p^3q^4}{2592}\frac{7p^2q^5}{17280}+\frac{7pq^6}{34560}\frac{19q^7}{483840}+𝒪((p+q)^6).$$
Later the methods for calculating the multiplicity series of symmetric functions of special kinds were significantly improved \[DG2\]. The Hilbert series of $`C_{42}`$ calculated by Teranishi \[T1, T2\] (with some typos corrected in \[BS\]) and the Hilbert series of $`T_{32}`$ and $`T_{42}`$ calculated by Berele and Stembridge \[BS\] allowed to express their multiplicity series and to determine the asympotics of the multiplicities. We shall state simplified versions of the results:
###### Theorem 3.8.
(i) (Drensky, Genov, Valenti \[DGV\]) The multiplicities $`m_{(\lambda _1,\lambda _2)}(C_{32})`$ and $`m_{(\lambda _1,\lambda _2)}(T_{32})`$ of the Hilbert series of $`C_{32}`$ and $`T_{32}`$, respectively, are related by
$$m_{(\lambda _1,\lambda _2)}(T_{32})9m_{(\lambda _1,\lambda _2)}(C_{32}).$$
(ii)( Drensky and Genov \[DGV\]) Let $`\lambda =(\lambda _1,\lambda _2)`$. The multiplicities $`m_\lambda (C_{42})`$ of the Hilbert series of $`C_{42}`$ satisfy the condition
$$m_\lambda (C_{42})=\{\begin{array}{cc}m_1+𝒪((\lambda _1+\lambda _2)^{13}),\hfill & \text{if }\lambda _1>3\lambda _2\text{,}\hfill \\ m_1+m_2+𝒪((\lambda _1+\lambda _2)^{13}),\hfill & \text{if }3\lambda _2\lambda _1>2\lambda _2\text{,}\hfill \\ m_1+m_2+m_3+𝒪((\lambda _1+\lambda _2)^{13}),\hfill & \text{if }2\lambda _2\lambda _1\text{,}\hfill \end{array}$$
where
$$m_1=\frac{(\lambda _1\lambda _2)^3\lambda _2^{11}}{11!3!2^83^2}\frac{(\lambda _1\lambda _2)^2\lambda _2^{12}}{12!2!2^83^3}+\frac{127(\lambda _1\lambda _2)\lambda _2^{13}}{13!2^{10}3^4}\frac{305\lambda _2^{14}}{14!2^93^5},$$
$$m_2=\frac{(3\lambda _2\lambda _1)^{14}}{14!2^{10}3^55^2},$$
$$m_3=\frac{(\lambda _1\lambda _2)(2\lambda _2\lambda _1)^{13}}{13!2^{10}3^25}\frac{7(2\lambda _2\lambda _1)^{14}}{14!2^935^2}.$$
The multiplicities $`m_\lambda (T_{42})`$ satisfy
$$m_\lambda (T_{42})=16m_\lambda (C_{42})+𝒪((\lambda _1+\lambda _2)^{13}).$$
We want to mention that Berele and Stembridge \[BS\] computed also the Hilbert series of $`C_{33}`$ and $`T_{33}`$ but the methods of \[DG1, DG2, DG3, DGV\] do not work successfully for symmetric functions in three variables. One can introduce the multiplicity series of a symmetric function in any number of variables, generalizing in an obvious way the case of symmetric functions in two variables. A recent theorem of Berele \[B2\] gives the rationality of the multiplicity series of a class of rational symmetric functions in any number of variables, including the Hilbert series of $`C_{nd}`$ and $`T_{nd}`$. Unfortunately, it is not clear how to perform the concrete calculations, even for the Hilbert series of $`C_{33}`$ and $`T_{33}`$.
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# A class of integrable spin Calogero-Moser systems II: exact solvability
## 1. Introduction
In \[LX1\],\[LX2\], we introduce a class of spin Calogero-Moser (CM) systems associated with so-called classical dynamical r-matrices with spectral parameter, as defined and classified in \[EV\] for simple Lie algebras, following the pioneering work of Felder \[F\] and Felder and Wieczerkowski \[FW\] in which the classical dynamical Yang-Baxter equation (CDYBE) with spectral parameter was introduced and studied in the context of conformal field theory. The main purpose of this sequel is to show how to obtain the explicit solutions of the associated integrable systems in \[LX2\] by using the factorization method developed in \[L2\].
The spin CM systems constructed in the afore-mentioned papers are of three types-rational, trigonometric and elliptic. Indeed, for each of the canonical forms of the three types of $`z`$-dependent classical dynamical r-matrices in \[EV\], there is an intrinsic way to construct an associated spin CM system and its realization in the dual bundle of a corresponding coboundary dynamical Lie algebroid. In this way, we are led to a family of rational spin CM systems parametrized by subsets $`\mathrm{\Delta }^{}\mathrm{\Delta }`$ which are closed with respect to addition and multiplication by $`1`$. Here $`\mathrm{\Delta }`$ is the root system associated with a complex simple Lie algebra $`𝔤`$ with a fixed Cartan subalgebra $`𝔥.`$ In the trigonometric case, we also have a family but here the family is parametrized by subsets $`\pi ^{}`$ of a fixed simple system $`\pi \mathrm{\Delta }.`$ Finally, we have an elliptic spin Calogero-Moser system for every simple Lie algebra. Let us summarize a few key features of these Hamiltonian systems and their realization spaces as follows: (a) in each case the phase space $`P`$ is a Hamiltonian $`H`$-space (with equivariant momentum map $`J`$) which admits an $`H`$-equivariant realization in the dual bundle $`A\mathrm{\Gamma }`$ of an infinite dimensional coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }`$ and the Hamiltonian is the pullback of a natural invariant function on $`A\mathrm{\Gamma }`$ under the realization map $`\rho `$, (b) the coboundary dynamical Lie algebroids involved are associated with solutions of the modified dynamical Yang-Baxter equation (mDYBE), (c) the pullbacks of the natural invariant functions on $`A\mathrm{\Gamma }`$ by $`\rho `$ do not Poisson commute everywhere on $`P`$, but they do so on the fiber $`J^1(0)`$ in all cases, (d) the reduced Hamiltonian system on $`J^1(0)/H`$ admits a natural collection of Poisson commuting integrals. As a matter of fact, we now know from \[L2\] that $`A\mathrm{\Gamma }`$ is also a Hamiltonian $`H`$-space and it follows from the same work that the integrable flows on the reduced space are realized on the Poisson quotient $`\gamma ^1(0)/H`$, where $`\gamma :A\mathrm{\Gamma }𝔥^{},(q,\lambda ,X)\lambda `$ is the momentum map of the $`H`$-action on $`A\mathrm{\Gamma }.`$ Indeed we have $`\rho (J^1(0))\gamma ^1(0)`$ in each case. Since $`\rho `$ is $`H`$-equivariant, it therefore induces a Poisson map between the corresponding Poisson quotients.
We now turn to our approach on exact solvability. As we showed in \[L2\], the mDYBE is associated with factorization problems on trivial Lie groupoids the solutions of which provide an effective method to integrate the generalized Lax equations on $`\gamma ^1(0).`$ (See \[L1\] for the groupoid version.) In the context of our spin CM systems, the factors which appear in the factorization problems are elements of certain Lie subgroupoids of trivial Lie groupoids of the form $`\mathrm{\Gamma }=U\times LG\times U,`$ where $`LG`$ is the loop group associated with a simple Lie group $`G`$, and $`U`$ is an open subset of a Cartan subalgebra $`𝔥𝔤=Lie(G).`$ Consequently, it is essential to have a precise description of the elements which belong to these Lie subgroupoids. As a first step towards the solutions of the factorization problems, we begin by analyzing the corresponding decompositions on the infinitesimal level, i.e., at the level of Lie algebroids. Here, the $`r`$-matrix $`:A^{}\mathrm{\Gamma }A\mathrm{\Gamma }`$ of the coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }`$ plays an important role. Indeed, for associated bundle maps $`^\pm :A^{}\mathrm{\Gamma }A\mathrm{\Gamma }`$, $`Im^\pm `$ are Lie subalgebroids of the trivial Lie algebroid $`A\mathrm{\Gamma }`$ and the infinitesimal version of factorization is the decomposition
$$(0_q,X,0)=\frac{1}{2}^+(0_q,X,0)\frac{1}{2}^{}(0_q,X,0)$$
where the element $`(0_q,X,0)`$ on the left hand side of the above expression is in the adjoint bundle of $`A\mathrm{\Gamma }.`$ Thus it is essential to be able to describe the elements of $`Im^\pm .`$ As it turns out, $`Im^+=_{qU}\{0_q\}\times L^+𝔤\times 𝔥`$ in all three cases, where $`L^+𝔤`$ is the Lie subalgebra of the loop algebra $`L𝔤`$ consisting of convergent power series $`_0^{\mathrm{}}X_nz^n.`$ On the other hand, $`Im^{}`$ in each of the three cases is given by a matched product $`^{}𝒬`$ (in the sense of Mokri \[Mok\]) where the ideal $`^{}`$ coincides with the adjoint bundle of $`Im^{}`$ and $`𝒬`$ is a Lie subalgebroid of $`Im^{}`$ isomorphic to $`TU.`$ In spite of this, the method of solution of the factorization problems is quite different in the three cases under consideration. That this is so is due to the difference in the analyticity properties of the elements in the ideals $`^{}`$. In the rational case and trigonometric case, we can reduce the factorization problems to finite dimensional problems due to some special features of the Lax operators. (Of course, this is also a reflection of the underlying character of the flows.) However, this is not so in the elliptic case-here we will only do things for the classical Lie algebras and indeed we will only give details for $`𝔤=sl(N,)`$ as the arguments for the other classical Lie algebras are similar. In this case, the explicit solution of the factorization problem is obtained with the help of multi-point Baker-Akheizer functions connected with the spectral curve $`C`$. Thus the solution of the equations can be expressed in terms of Riemann theta functions.
The paper is organized as follows. In Section 2, we present a number of basic results which will be used throughout the paper. More specifically, we will begin with a review on the geometric scheme to construct integrable systems based on realization in the dual bundles of coboundary dynamical Lie algebroids and the factorization theory which we mentioned earlier. Then we will turn our attention to a subclass of coboundary dynamical Lie algebroids defined by the classical dynamical r-matrices with spectral parameter. We will also recall what we mean by spin Calogero-Moser systems associated with this subclass of coboundary dynamical Lie algebroids. In Section 3, we discuss the solution of the integrable rational spin Calogero-Moser systems by solving the corresponding factorization problem. In Section 4, we handle the trigonometric case. Finally in Section 5, we analyze the elliptic case.
To close, we would like to point out what was done in the paper \[KBBT\] so that the reader can better understand why a different method is required for our more general class of systems here. To cut the story short, what the authors considered in \[KBBT\] are the $`gl(N,)`$-rational spin Calogero-Moser system of Gibbons and Hermsen \[GH\], as well as their trigonometric and elliptic counterparts. From our point of view, these are special cases which can be obtained from more general $`gl(N,)`$-systems (see \[BAB\] and \[ABB\]) by restricting the matrix of ‘spin variables’ to some special coadjoint orbits of $`gl(N,)^{}gl(N,)`$ which can be parametrized by vectors $`a_j,b_j^l`$, $`l<N,j=1,\mathrm{},N.`$ The method for solving such systems in \[KBBT\] is based on the connection with the matrix KP equation and is specific to these special coadjoint orbits of $`gl(N,)^{}gl(N,)`$. For a sketch of this method in the elliptic case together with an explanation of its limitations, we refer the reader to Remark 5.2.8 (a).
## 2. Invariant Hamiltonian systems associated with coboundary dynami- fak cal Lie algebroids and the factorization method
In the first two subsections, we shall present a number of basic results from \[L2\] which will be used throughout the paper. In particular, we shall give a summary of the factorization theory. In the last subsection, we shall turn our attention to a subclass of coboundary dynamical Lie algebroids defined by the classical dynamical r-matrices with spectral parameter \[LX2\]. Since the paper is concerned with the solution of the class of integrable spin Calogero-Moser systems introduced in \[LX2\], we will recall its construction and its relation to this subclass of Lie algebroids in this last subsection.
## 2.1 Geometry of the modified dynamical Yang-Baxter equation
Let $`G`$ be a connected Lie group, $`HG`$ a connected Lie subgroup, and $`𝔤`$, $`𝔥`$ their Lie algebras. We shall denote by $`\iota :𝔥𝔤`$ the Lie inclusion. In what follows, the Lie groups and Lie algebras can be real or complex unless we specify otherwise.
If $`U𝔥^{}`$ is a connected $`Ad_H^{}`$-invariant open subset, we say that a smooth (resp. holomorphic) map $`R:UL(𝔤^{},𝔤)`$ (here and henceforth we denote by $`L(𝔤^{},𝔤)`$ the set of linear maps from $`𝔤^{}`$ to $`𝔤`$) is a classical dynamical r-matrix \[EV\] associated with the pair $`(𝔤,𝔥)`$ iff $`R`$ is pointwise skew symmetric
$$<R(q)(A),B>=<A,R(q)B>$$
$`(\mathrm{2.1.1})`$
and satisfies the classical dynamical Yang-Baxter condition
$$\begin{array}{cc}& [R(q)A,R(q)B]+R(q)(ad_{R(q)A}^{}Bad_{R(q)B}^{}A)\hfill \\ \hfill +& dR(q)\iota ^{}A(B)dR(q)\iota ^{}B(A)+d<R(A),B>(q)=\chi (A,B),\hfill \end{array}$$
$`(\mathrm{2.1.2})`$
for all $`qU`$, and all $`A,B𝔤^{}`$, where $`\chi :𝔤^{}\times 𝔤^{}𝔤`$ is $`G`$-equivariant. The dynamical $`r`$-matrix is said to be $`H`$-equivariant if and only if
$$R(Ad_{h^1}^{}q)=Ad_hR(q)Ad_h^{}$$
$`(\mathrm{2.1.3})`$
for all $`hH,qU.`$
We shall equip $`\mathrm{\Gamma }=U\times G\times U`$ with the trivial Lie groupoid structure over $`U`$ with structure maps (target, source, $`\mathrm{},`$ multiplication)
$$\begin{array}{cc}& \alpha (u,g,v)=u,\beta (u,g,v)=v,ϵ(u)=(u,1,u),i(u,g,v)=(u,g^1,v)\hfill \\ & m((u,g,v),(v,g^{},w))=(u,gg^{},w)\hfill \end{array}$$
$`(\mathrm{2.1.4})`$
and let $`A\mathrm{\Gamma }=KerT\alpha |_{ϵ(U)}=_{qU}\{0_q\}\times 𝔤\times 𝔥^{}TU\times 𝔤`$ be its Lie algebroid with anchor map denoted by $`a.`$ (See \[CdSW\] and \[M1\] for details.) Recall that associated with an $`H`$-equivariant classical dynamical r-matrix $`R`$ there is a natural Lie algebroid bracket $`[,]_{A^{}\mathrm{\Gamma }}`$ on the dual bundle $`A^{}\mathrm{\Gamma }`$ of $`A\mathrm{\Gamma }`$ \[BKS\],\[L2\] such that the pair $`(A\mathrm{\Gamma },A^{}\mathrm{\Gamma })`$ is a Lie bialgebroid in the sense of MacKenzie and Xu \[MX\].(Lie bialgebroids are infinitesimal versions of the Poisson groupoids of Weinstein \[W\].) Throughout the paper, the pair $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }})`$ together with the anchor map $`a_{}:A^{}\mathrm{\Gamma }TU`$ given by
$$a_{}(0_q,A,Z)=(q,\iota ^{}Aad_Z^{}q)$$
$`(\mathrm{2.1.5})`$
will be called the coboundary dynamical Lie algebroid associated to $`R`$.
A special case of (2.1.2) is the modified dynamical Yang-Baxter equation (mDYBE):
$$\begin{array}{cc}& [R(q)A,R(q)B]+R(q)(ad_{R(q)A}^{}Bad_{R(q)B}^{}A)\hfill \\ \hfill +& dR(q)\iota ^{}A(B)dR(q)\iota ^{}B(A)+d<R(A),B>(q)\hfill \\ \hfill =& [K(A),K(B)]\hfill \end{array}$$
$`(\mathrm{2.1.6})`$
where $`KL(𝔤^{},𝔤)`$ is a nonzero symmetric map which satisfies $`ad_XK+Kad_X^{}=0`$ for all $`X𝔤`$,i.e., $`K`$ is $`G`$-equivariant. In \[L2\], the class of coboundary dynamical Lie algebroids associated with mDYBE was singled out and was shown to have some rather remarkable properties. We will restrict to this class of $`(A^{}\mathrm{\Gamma },[,]_{A^{}\mathrm{\Gamma }},a_{})`$ in the rest of the subsection.
Following \[L1\] and \[L2\], we introduce the bundle map
$$:A^{}\mathrm{\Gamma }A\mathrm{\Gamma },(0_q,A,Z)(0_q,\iota Z+R(q)A,\iota ^{}Aad_Z^{}q)$$
$`(\mathrm{2.1.7})`$
and call it the r-matrix of the Lie algebroid $`A^{}\mathrm{\Gamma }`$. Also, define
$$𝒦:A^{}\mathrm{\Gamma }A\mathrm{\Gamma },(0_q,A,Z)(0_q,K(A),0),$$
$`(\mathrm{2.1.8})`$
and set
$$^\pm =\pm 𝒦,R^\pm (q)=R(q)\pm K.$$
$`(\mathrm{2.1.9})`$
###### Proposition 2.1.1
(a) $`^\pm `$ are morphisms of transitive Lie algebroids and, as morphisms of vector bundles over $`U`$, are of locally constant rank. (b) $`Im^\pm `$ are transitive Lie subalgebroids of $`A\mathrm{\Gamma }.`$
In the rest of the subsection,we shall assume $`𝔤`$ has a nondegenerate invariant pairing $`(,)`$ such that $`(,)_{𝔥\times 𝔥}`$ is also nondengenerate. Without loss of generality, we shall take the map $`K:𝔤^{}𝔤`$ in the above discussion to be the identification map induced by $`(,)`$. Indeed, with the identifications $`𝔤^{}𝔤`$, $`𝔥^{}𝔥`$, we have $`K=id_𝔤.`$ We shall regard $`R(q)`$ as taking values in $`End(𝔤)`$, and the derivatives as well as the dual maps are computed using $`(,)`$. Also, we have $`ad^{}ad`$, $`\iota ^{}\mathrm{\Pi }_𝔥`$, where $`\mathrm{\Pi }_𝔥`$ is the projection map to $`𝔥`$ relative to the direct sum decomposition $`𝔤=𝔥𝔥^{}`$. We shall keep, however, the notation $`A^{}\mathrm{\Gamma }`$ although as a set it can be identified with $`A\mathrm{\Gamma }.`$
We now introduce the following subbundles of the adjoint bundle $`Kera=\{(0_q,X,0)qU,X𝔤\}`$ of $`A\mathrm{\Gamma }`$:
$$^\pm =\{(0_q,X,0)KeraqU,^{}(0_q,X,Z)=0\text{for some}Z𝔥\}.$$
$`(\mathrm{2.1.10})`$
###### Proposition 2.1.2
(a) $`^\pm `$ are ideals of the transitive Lie algebroids $`Im^\pm `$. (b) Equip $`Im^+/^+`$ and $`Im^{}/^{}`$ with the quotient transitive Lie algebroid structures, then the map $`\theta :Im^+/^+Im^{}/^{}`$ defined by
$$\theta (^+(0_q,X,Z)+_q^+)=^{}(0_q,X,Z)+_q^{}$$
is an isomorphism of transitive Lie algebroids.
###### Theorem 2.1.3
(a) Every element $`(0_q,X,0)Kera`$ admits a unique decomposition
$$(0_q,X,0)=𝒳_+𝒳_{}$$
where $`(𝒳_+,𝒳_{})(^+,^{})(\{0_q\}\times 𝔤\times \{0\})`$ with $`\theta (𝒳_++_q^+)=𝒳_{}+_q^{}.`$ (b) The coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }`$ is isomorphic to the Lie subalgebroid
$$\{(𝒳_+,𝒳_{})(Im^+\underset{TU}{}Im^{})_qqU,\theta (𝒳_++_q^+)=𝒳_{}+_q^{}\}$$
of $`Im^+\underset{TU}{}Im^{}`$.
## 2.2 Invariant Hamiltonian systems and the factorization method
We shall continue to use the same assumptions on $`𝔤`$, $`K`$ and $`R`$ as in the latter part of Section 2.1. Although it is not entirely necessary to assume that $`R`$ satisfies the mDYBE for Theorem 2.2.1, however, it serves our purpose here.
Let $`P`$ be a Poisson manifold and suppose $`\rho =(m,\tau ,L):PU\times 𝔥\times 𝔤A\mathrm{\Gamma }`$ is a realization of $`P`$ in the dual bundle $`A\mathrm{\Gamma }`$ of the coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }`$, equipped with the Lie-Poisson structure, i.e., $`\rho `$ is a Poisson map. Recall that $`A\mathrm{\Gamma }`$ is a Hamiltonian $`H`$-space under the natural action $`h.(q,\lambda ,X)=(Ad_{h^1}^{}q,Ad_{h^1}^{}\lambda ,Ad_hX)`$ and the projection map $`\gamma :A\mathrm{\Gamma }𝔥^{},(q,\lambda ,X)\lambda `$ is an equivariant momentum map. We assume: A1. $`P`$ is a Hamiltonian $`H`$-space with an equivariant momentum map $`J:P𝔥`$,
A2. the realization map $`\rho `$ is $`H`$-equivariant,
A3. for some regular value $`\mu 𝔥`$ of $`J`$,
$$\rho (J^1(\mu ))\gamma ^1(0)=U\times \{0\}\times 𝔤.$$
$`(\mathrm{2.2.1})`$
Let $`I(𝔤)`$ be the ring of smooth ad-invariant functions on $`𝔤`$, $`i_\mu :J^1(\mu )P`$ the inclusion map, and $`\pi _\mu :J^1(\mu )J^1(\mu )/H_\mu `$ the canonical projection, where $`H_\mu `$ is the isotropy subgroup of $`\mu `$ for the $`H`$-action on $`P.`$ Also, let $`Pr_3`$ denote the projection map onto the third factor of $`U\times 𝔥\times 𝔤A\mathrm{\Gamma }.`$ We consider $`H`$-invariant Hamiltonian systems on $`P`$, generated by Hamiltonians of the form $`=L^{}f`$, where $`fI(𝔤)`$.
###### Theorem 2.2.1
Under the above assumptions, (a) The Hamiltonian $``$ descends to a unique function $`_\mu `$ on the reduced Poisson variety $`P_\mu =J^1(\mu )/H_\mu .`$ Moreover, $`_\mu `$ admits a natural family of Poisson commuting functions given by the reduction of functions in $`L^{}I(𝔤)`$, (b) If $`\psi _t`$ is the induced flow on $`\gamma ^1(0)`$ generated by the Hamiltonian $`Pr_3^{}f`$, and $`\varphi _t`$ is the Hamiltonian flow of $``$ on $`P`$, then under the flow $`\varphi _t`$, we have
$`{\displaystyle \frac{d}{dt}}m(\varphi _t)=\mathrm{\Pi }_𝔥df(L(\varphi _t)),`$
$`{\displaystyle \frac{d}{dt}}\tau (\varphi _t)=0,`$
$`{\displaystyle \frac{d}{dt}}L(\varphi _t)=[L(\varphi _t),R(m(\varphi _t))df(L(\varphi _t))]+dR(m(\varphi _t))(\tau (\varphi _t))df(L(\varphi _t))`$
where the term involving $`dR`$ drops out on $`J^1(\mu )`$. Moreover, the reduction $`\varphi _t^{red}`$ of $`\varphi _ti_\mu `$ on $`P_\mu `$ defined by $`\varphi _t^{red}\pi _\mu =\pi _\mu \varphi _ti_\mu `$ is a Hamiltonian flow of $`_\mu .`$
This theorem applies in particular to the special case where $`P=A\mathrm{\Gamma }`$, $`\rho =id_{A\mathrm{\Gamma }}`$ (see \[L2\] for the proof that $`A\mathrm{\Gamma }`$ is a Hamiltonian $`H`$-space). In fact, we have a factorization method for solving the generalized Lax equations
$$\begin{array}{cc}& \frac{d}{dt}(q,0,X)\hfill \\ \hfill =& (\mathrm{\Pi }_𝔥df(X),0,[X,R(q)df(X)])\hfill \end{array}$$
$`(\mathrm{2.2.2})`$
on the invariant manifold $`\gamma ^1(0).`$ In what follows, if $`(u,g,v)\mathrm{\Gamma },`$ the symbol $`𝚕_{(u,g,v)}`$ will stand for left translation in $`\mathrm{\Gamma }`$ by $`(u,g,v)`$, the lift of $`Im(^+,^{})A\mathrm{\Gamma }\underset{TU}{}A\mathrm{\Gamma }`$ to the groupoid level will be denoted by the same symbol, and $`𝔸𝕕`$ is the adjoint representation of $`\mathrm{\Gamma }`$ on its adjoint bundle $`Kera`$, defined by $`𝔸𝕕_\gamma (q,0,X)=(q^{},0,Ad_kX)`$, for $`\gamma =(q^{},k,q)\mathrm{\Gamma }`$. Lastly, if $`Y`$ is a section of $`A\mathrm{\Gamma }`$, the exponential $`exp(tY):U\mathrm{\Gamma }`$ is defined by the formula $`exp(tY)(q)=f_t(ϵ(q))`$, where $`f_t`$ is the local flow generated by the left-invariant vector field $`\stackrel{}{Y}`$: $`\stackrel{}{Y}(u,g,v)=T_{ϵ(v)}𝚕_{(u,g,v)}Y(v).`$ In particular, if we take $`Y`$ to be the constant section $`(0,0,\xi )`$ of $`A\mathrm{\Gamma },`$ $`\xi 𝔤`$, an easy computation shows that $`exp\{t(0,0,\xi )\}(q)=(q,e^{t\xi },q).`$
###### Theorem 2.2.2
Suppose that $`fI(𝔤)`$, $`F=Pr_3^{}f`$ and $`q_0U`$, where $`U`$ is simply connected. Then for some $`0<T\mathrm{}`$, there exists a unique element $`(\gamma _+(t),\gamma _{}(t))=((q_0,k_+(t),q(t)),(q_0,k_{}(t),q(t)))Im(^+,^{})\mathrm{\Gamma }\underset{U\times U}{\times }\mathrm{\Gamma }`$ for $`0t<T`$ which is smooth in t, solves the factorization problem
$$exp\{2t(0,0,df(X_0))\}(q_0)=\gamma _+(t)\gamma _{}(t)^1$$
$`(\mathrm{2.2.3})`$
and satisfies
$$(T_{\gamma _+(t)}𝚕_{\gamma _+(t)^1}\dot{\gamma }_+(t),T_{\gamma _{}(t)}𝚕_{\gamma _{}(t)^1}\dot{\gamma }_{}(t))(^+,^{})(\{q(t)\}\times \{0\}\times 𝔤)$$
$`(\mathrm{2.2.4}a)`$
with
$$\gamma _\pm (0)=(q_0,1,q_0).$$
$`(\mathrm{2.2.4}b)`$
Moreover, the solution of Eqn.(2.2.2) with initial data $`(q,0,X)(0)=(q_0,0,X_0)`$ is given by the formula
$$(q(t),0,X(t))=𝔸𝕕_{\gamma _\pm (t)^1}(q_0,0,X_0).$$
$`(\mathrm{2.2.5})`$
###### Corollary 2.2.3
Let $`\psi _t`$ be the induced flow on $`\gamma ^1(0)`$ as defined in (2.2.5) and let $`\varphi _t`$ be the Hamiltonian flow of $`=L^{}f`$ on $`P`$, where $`L=Pr_3\rho `$ for a realization map $`\rho :PA\mathrm{\Gamma }`$ satisfying A1-A3. If we can solve for $`\varphi _t(x)`$, $`xJ^1(\mu )`$ explicitly from the relation $`\rho (\varphi _t)(x)=\psi _t(\rho (x))`$, then the formula $`\varphi _t^{red}\pi _\mu =\pi _\mu \varphi _ti_\mu `$ gives an explicit expression for the flow of the reduced Hamiltonian $`_\mu `$.
## 2.3 Classical dynamical r-matrices with spectral parameter and the fake associated spin Calogero-Moser systems
From now onwards, we let $`𝔤`$ be a complex simple Lie algebra of rank $`N`$ with Killing form $`(,)`$ and let $`G`$ be the connected and simply-connected Lie group which integrates $`𝔤.`$ We fix a Cartan subalgebra $`𝔥`$ of $`𝔤`$ and let $`𝔤=𝔥_{\alpha \mathrm{\Delta }}𝔤_\alpha `$ be the root space decomposition of $`𝔤`$ with respect to $`𝔥`$. For each $`\alpha \mathrm{\Delta }`$, denote by $`H_\alpha `$ the element in $`𝔥`$ which corresponds to $`\alpha `$ under the isomorphism between $`𝔥`$ and $`𝔥^{}`$ induced by the Killing form $`(,)`$. We fix a simple system of roots $`\pi =\{\alpha _1,\mathrm{},\alpha _N\}`$ and denote by $`\mathrm{\Delta }^\pm `$ the corresponding positive/negative system. Also, for each $`\alpha \mathrm{\Delta }^+`$, we pick root vectors $`e_\alpha 𝔤_\alpha `$, $`e_\alpha 𝔤_\alpha `$ which are dual with respect to $`(,)`$ so that $`[e_\alpha ,e_\alpha ]=H_\alpha .`$
Let $`r:𝔥\times 𝔤𝔤`$ be a classical dynamical r-matrix with spectral parameter in the sense of \[EV\], with coupling constant equal to $`1`$. Then we can construct the associated classical dynamical r-matrix $`R`$ \[LX2\] for the pair $`(L𝔤,𝔥)`$, where $`L𝔤`$ is the loop algebra of $`𝔤`$. Indeed, if $`L𝔤^{}`$ denotes the restricted dual of $`L𝔤`$ and we make the identification $`L𝔤^{}L𝔤`$, then we have the following result.
###### Proposition 2.3.1
For $`XL𝔤`$, we have the formula
$$(R(q)X)(z)=\frac{1}{2}X(z)+\underset{k0}{}\frac{1}{k!}(\frac{^kr}{z^k}(q,z),X_{(k+1)}1).$$
$`(\mathrm{2.3.1})`$
Moreover, $`R`$ is a solution of the mDYBE with $`K=\frac{1}{2}id_{L𝔤}`$.
Remark 2.3.2. As the reader will see, the formula in (2.3.1) will be used to compute the explicit expressions for $`R`$ which play a critical role in characterizing the elements in the Lie subalgebroids $`Im^\pm `$ in our analysis in Section 3, 4 and 5 below. This formula has also been used in connection with symmetric coboundary dynamical Lie algebroids in \[L3\].
We now fix a simply-connected domain $`U𝔥`$ on which $`R`$ is holomorphic. Also, introduce the trivial Lie groupoids $`\mathrm{\Omega }=U\times G\times U`$, $`\mathrm{\Gamma }=U\times LG\times U,`$ where $`LG`$ is the loop group of the simple Lie group $`G.`$ Then we can use the map $`R:UL(L𝔤,L𝔤)`$ in Proposition 2.3.1 to construct the associated coboundary dynamical Lie algebroid $`A^{}\mathrm{\Gamma }=T^{}U\times L𝔤^{}TU\times L𝔤`$ so that its dual bundle has a Lie-Poisson structure. (See \[L2\] for explicit formulas.) On the other hand, we shall equip the dual bundle $`A^{}\mathrm{\Omega }TU\times 𝔤`$ of the trivial Lie algebroid $`A\mathrm{\Omega }`$ with the corresponding Lie-Poisson structure. Explicitly, $`\{\phi ,\psi \}_{A^{}\mathrm{\Omega }}(q,p,\xi )=(\delta _2\phi ,\delta _1\psi )(\delta _1\phi ,\delta _2\psi )+(\xi ,[\delta \phi ,\delta \psi ])`$ where $`\delta _1\phi ,`$ $`\delta _2\phi `$ and $`\delta \phi `$ are the partial derivatives of $`\phi `$ with respect to the variables in $`U`$, $`𝔥`$ and $`𝔤`$ respectively.
###### Theorem 2.3.3
The map $`\rho =(m,\tau ,L):A^{}\mathrm{\Omega }TU\times 𝔤TU\times L𝔤A\mathrm{\Gamma }`$ given by
$$\rho (q,p,\xi )=(q,\mathrm{\Pi }_𝔥\xi ,p+r_{}^\mathrm{\#}(q)\xi )$$
$`(\mathrm{2.3.2})`$
is an H-equivariant Poisson map, when the domain is equipped with the Lie-Poisson structure corresponding to the trivial Lie algebroid $`A\mathrm{\Omega }TU\times 𝔤`$, and the target is equipped with the Lie-Poisson structure corresponding to $`A^{}\mathrm{\Gamma }`$. Here, the map $`r_{}^\mathrm{\#}(q):𝔤L𝔤`$ is defined by
$$((r_{}^\mathrm{\#}(q)\xi )(z),\eta )=(r(q,z),\eta \xi )$$
$`(\mathrm{2.3.3})`$
for $`\xi `$, $`\eta 𝔤.`$
Let $`Q`$ be the quadratic function
$$Q(X)=\frac{1}{2}_c(X(z),X(z))\frac{dz}{2\pi iz}$$
$`(\mathrm{2.3.4})`$
where $`c`$ is a small circle around the origin. Clearly, $`Q`$ is an ad-invariant function on $`L𝔤`$.
###### Definition Definition 2.3.4
Let $`r`$ be a classical dynamical r-matrix with spectral parameter with coupling constant equal to 1. Then the Hamiltonian system on $`A^{}\mathrm{\Omega }TU\times 𝔤`$ (equipped with the Lie-Poisson structure as in Theorem 2.3.3) generated by the Hamiltonian
$$(q,p,\xi )=(L^{}Q)(q,p,\xi )=\frac{1}{2}_c(L(q,p,\xi )(z),L(q,p,\xi )(z))\frac{dz}{2\pi iz}$$
$`(\mathrm{2.3.5})`$
is called the spin Calogero-Moser system associated to $`r`$.
###### Proposition 2.3.5
The Hamiltonians of the spin Calogero-Moser systems are invariant under the Hamiltonian $`H`$-action on $`A^{}\mathrm{\Omega }TU\times 𝔤`$ given by $`h(q,p,\xi )=(q,p,Ad_h\xi )`$ with equivariant momentum map $`J(q,p,\xi )=\mathrm{\Pi }_𝔥\xi .`$ Moreover, under the Hamiltonian flow, we have
$$\begin{array}{cc}& \dot{q}=\mathrm{\Pi }_𝔥(M(q,p,\xi ))_1,\hfill \\ & \dot{L}(q,p,\xi )=[L(q,p,\xi ),R(q)M(q,p,\xi )]\hfill \end{array}$$
$`(\mathrm{2.3.6})`$
on the invariant manifold $`J^1(0)`$, where
$$M(q,p,\xi )(z)=L(q,p,\xi )(z)/z.$$
$`(\mathrm{2.3.7})`$
Clearly, the second part of the above proposition is a consequence of Theorem 2.2.1 (b) and Theorem 2.3.3.
In order to discuss the associated integrable models, we have to restrict to a smooth component of the reduced Poisson variety $`J^1(0)/H=U\times 𝔥\times (𝔥^{}/H)`$. For this purpose, we restrict to the following open submanifold of $`𝔤`$:
$$𝒰=\{\xi 𝔤\xi _{\alpha _i}=(\xi ,e_{\alpha _i})0,i=1,\mathrm{},N\}.$$
$`(\mathrm{2.3.8})`$
Then the $`H`$-action in Proposition 2.3.5 above induces a Hamiltonian $`H`$-action on $`TU\times 𝒰`$ and we denote the corresponding momentum map also by $`J`$ so that $`J^1(0)=TU\times (𝔥^{}𝒰)`$. Now, recall from \[LX2\] that the formula
$$g(\xi )=\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}\underset{j=1}{\overset{N}{}}(C_{ji}\mathrm{log}\xi _{\alpha _j})h_{\alpha _i}\right)$$
$`(\mathrm{2.3.9})`$
defines an $`H`$-equivariant map $`g:𝒰H`$ where $`C=(C_{ij})`$ is the inverse of the Cartan matrix and $`h_{\alpha _i}=\frac{2}{(\alpha _i,\alpha _i)}H_{\alpha _i}`$, $`i=1,\mathrm{},N`$. Using $`g`$, we can identify the reduced space $`J^1(0)/H=TU\times (𝔥^{}𝒰/H)`$ with $`TU\times 𝔤_{red}`$, where $`𝔤_{red}=ϵ+_{\alpha \mathrm{\Delta }\pi }e_\alpha `$, and $`ϵ=_{j=1}^Ne_{\alpha _j}`$. Thus the projection map $`\pi _0:J^1(0)TU\times 𝔤_{red}`$ is the map
$$(q,p,\xi )(q,p,Ad_{g(\xi )^1}\xi ).$$
$`(\mathrm{2.3.10})`$
We shall write $`s=_{\alpha \mathrm{\Delta }}s_\alpha e_\alpha `$ for $`s𝔤_{red}`$ (note that $`s_{\alpha _j}=1\text{for}j=1,\mathrm{},N`$). By Poisson reduction \[MR\], the reduced manifold $`TU\times 𝔤_{red}`$ has a unique Poisson structure which is a product structure where the second factor $`𝔤_{red}`$ is equipped with the reduction (at 0) of the Lie-Poisson structure on $`𝒰`$ by the $`H`$-action. If $``$ is the Hamiltonian defined in (2.3.5), we shall denote its reduction to $`TU\times 𝔤_{red}`$ by $`_0.`$
## 3. The rational spin Calogero-Moser systems
## 3.1. Lax operators, Hamiltonian equations and the Lie subalgebroids
Let $`𝔤`$ be a complex simple Lie algebra, as in Section 2.3. In addition to the basis $`\{e_\alpha \}_{\alpha \mathrm{\Delta }}`$ of $`_{\alpha \mathrm{\Delta }}g_\alpha `$ in that section, let us now fix an orthonormal basis $`(x_i)_{1iN}`$ of $`𝔥`$. Thus we will write $`p=_ip_ix_i`$, $`\xi =_i\xi _ix_i+_{\alpha \mathrm{\Delta }}\xi _\alpha e_\alpha `$ for $`p𝔥`$ and $`\xi 𝔤.`$
The rational spin Calogero-Moser systems are the Hamiltonian systems on $`TU\times 𝔤`$ (as defined in Definition 2.3.4) associated to the rational dynamical r-matrices with spectral parameter:
$$r(q,z)=\frac{\mathrm{\Omega }}{z}+\underset{\alpha \mathrm{\Delta }^{}}{}\frac{1}{\alpha (q)}e_\alpha e_\alpha ,$$
$`(\mathrm{3.1.1})`$
where $`\mathrm{\Delta }^{}\mathrm{\Delta }`$ is any set of roots which is closed with respect to addition and multiplication by $`1,`$ and $`\mathrm{\Omega }(S^2𝔤)^𝔤`$ is the Casimir element corresponding to the Killing form $`(,).`$ Accordingly, the Lax operators are given by
$$L(q,p,\xi )(z)=p+\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi _\alpha }{\alpha (q)}e_\alpha +\frac{\xi }{z}$$
$`(\mathrm{3.1.2})`$
and so we have a family of Hamiltonians parametrized by $`\mathrm{\Delta }^{}`$:
$$(q,p,\xi )=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{2}\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi _\alpha \xi _\alpha }{\alpha (q)^2}.$$
$`(\mathrm{3.1.3})`$
Note that in particular, we have
$$L(q,p,\xi )(\mathrm{})𝔤_\mathrm{\Delta }^{}$$
$`(\mathrm{3.1.4})`$
where
$$𝔤_\mathrm{\Delta }^{}=𝔥+\underset{\alpha \mathrm{\Delta }^{}}{}𝔤_\alpha $$
$`(\mathrm{3.1.5})`$
is a reductive Lie subalgebra of $`𝔤`$. As the reader will see, this fact is important later on, when we solve the factorization problem.
###### Proposition 3.1.1
The Hamiltonian equations of motion generated by $``$ on $`TU\times 𝔤`$ are given by
$$\begin{array}{cc}& \dot{q}=p,\hfill \\ & \dot{p}=\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi _\alpha \xi _\alpha }{\alpha (q)^3}H_\alpha ,\hfill \\ & \dot{\xi }=[\xi ,\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi _\alpha }{\alpha (q)^2}e_\alpha ].\hfill \end{array}$$
$`(\mathrm{3.1.6})`$
###### Demonstration Proof
From the expression
$$\{\phi ,\psi \}(q,p,\xi )=(\delta _2\phi ,\delta _1\psi )(\delta _1\phi ,\delta _2\psi )+(\xi ,[\delta \phi ,\delta \psi ])$$
for the Poisson bracket on $`TU\times 𝔤`$, the equations of motion are given by $`\dot{q}=\delta _2`$, $`\dot{p}=\delta _1`$, and $`\dot{\xi }=[\xi ,\delta ]`$. Therefore, (3.1.6) follows by a direct computation. $`\mathrm{}`$
We shall solve these equations on $`J^1(0)`$ by our factorization method. To do so, it is essential to have the explicit expression of the classical dynamical r-matrix $`R`$ associated to $`r`$. Before we make the computation, let us recall that the loop algebra $`L𝔤`$ admits a direct sum decomposition
$$L𝔤=L^+𝔤L_0^{}𝔤$$
$`(\mathrm{3.1.7})`$
into Lie subalgebras, where $`L^+𝔤`$ consists of convergent power series $`_0^{\mathrm{}}X_nz^n`$, while $`L_0^{}𝔤`$ consists of Laurent tails, of the form $`_{\mathrm{}}^1X_nz^n`$. We shall denote by $`\mathrm{\Pi }_\pm `$ the projection operators relative to this splitting.
###### Proposition 3.1.2
The classical dynamical r-matrix R associated with the meromorphic map $`r`$ in (3.1.1) is given by
$$(R(q)X)(z)=\frac{1}{2}\left(\mathrm{\Pi }_+X\mathrm{\Pi }_{}X\right)(z)\underset{\alpha \mathrm{\Delta }^{}}{}\frac{(X_1)_\alpha }{\alpha (q)}e_\alpha .$$
$`(\mathrm{3.1.8})`$
In particular,
$$(R(q)M(q,p,\xi ))(z)=\frac{1}{2}M(q,p,\xi )(z)\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi _\alpha }{\alpha (q)^2}e_\alpha $$
$`(\mathrm{3.1.9})`$
for $`M(q,p,\xi )(z)=L(q,p,\xi )(z)/z.`$ (See (2.3.6),(2.3.7).)
###### Demonstration Proof
By direct differentiation, we find that $`\frac{^kr}{z^k}(q,z)=k!\frac{\mathrm{\Omega }}{z^{k+1}}`$, $`k1`$. Substituting into (2.3.1), the formula follows. $`\mathrm{}`$
Remark 3.1.3. (a) The formula in (3.1.8) shows that the classical dynamical r-matirx $`R`$ is a perturbation of the standard r-matrix associated with the splitting in (3.1.7). (b) If we restrict ourselves to $`J^1(0)`$, then by equating the coefficients of $`z^0`$ and $`z^1`$ on both sides of the Lax equation $`\dot{L}(q,p,\xi )=[L(q,p,\xi ),R(q)M(q,p,\xi )]`$, we can recover the equations for $`p`$ and $`\xi `$ respectively in (3.1.6). However, the Lax equation only gives $`\alpha (\dot{q}p)`$ for all $`\alpha \mathrm{\Delta }^{}`$. Therefore, unless $`\mathrm{\Delta }^{}=\mathrm{\Delta }`$, otherwise, we cannot recover the equation for $`q`$ from that of $`L(q,p,\xi )`$. This remark shows that the full set of equations in (2.3.6) is important. We now give the equations of motion for the reduction of $``$ on $`TU\times 𝔤_{red}`$, with Hamiltonian given by
$$_0(q,p,s)=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{2}\underset{\alpha \mathrm{\Delta }^{}}{}\frac{s_\alpha s_\alpha }{\alpha (q)^2}.$$
$`(\mathrm{3.1.10})`$
###### Proposition 3.1.4
The Hamiltonian equations of motion generated by $`_0`$ on the reduced Poisson manifold $`TU\times 𝔤_{red}`$ are given by
$`\dot{q}=p,`$
$`\dot{p}={\displaystyle \underset{\alpha \mathrm{\Delta }^{}}{}}{\displaystyle \frac{s_\alpha s_\alpha }{\alpha (q)^3}}H_\alpha ,`$
$`\dot{s}=[s,]`$
where
$$\frac{=\underset{\alpha \mathrm{\Delta }^{}}{}\frac{s_\alpha }{\alpha (q)^2}e_\alpha +\underset{i,j}{}C_{ji}}{\alpha \mathrm{\Delta }^{}\alpha _j\alpha \mathrm{\Delta }N_{\alpha ,\alpha _j\alpha }\frac{s_\alpha s_{\alpha _j\alpha }}{\alpha (q)^2}h_{\alpha _i}.}$$
(Here we use the notation $`[e_\alpha ,e_\beta ]=N_{\alpha ,\beta }e_{\alpha +\beta }`$ if $`\alpha +\beta \mathrm{\Delta }`$.)
###### Demonstration Proof
The equations for $`q`$ and $`p`$ are obvious from Propostion 3.1.1 and the fact that $`s_\alpha =\xi _\alpha e^{\alpha (logg(\xi ))}.`$ To derive the equation for $`s`$, we differentiate $`s=Ad_{g(\xi )^1}\xi `$ with respect to $`t`$, assuming that $`\xi `$ satisfies the equation in Proposition 3.1.1 with $`\mathrm{\Pi }_𝔥\xi =0.`$ This gives
$$\dot{s}=[s,\underset{\alpha \mathrm{\Delta }^{}}{}\frac{s_\alpha }{\alpha (q)^2}e_\alpha T_{g(\xi )^1}r_{g(\xi )}\frac{d}{dt}g(\xi )^1].()$$
Now, from the expression for $`g(\xi )`$ in (2.3.9), we find
$$T_{g(\xi )^1}r_{g(\xi )}\frac{d}{dt}g(\xi )^1=\underset{i,j}{}C_{ji}\dot{\xi }_{\alpha _j}\xi _{\alpha _j}^{}{}_{}{}^{1}h_{\alpha _i}.$$
But from Proposition 3.1.1, we have
$`\dot{\xi }_{\alpha _j}=`$ $`([\xi ,{\displaystyle \underset{\alpha \mathrm{\Delta }^{}}{}}{\displaystyle \frac{\xi _\alpha }{\alpha (q)^2}}e_\alpha ],e_{\alpha _j})`$
$`=`$ $`{\displaystyle \frac{\xi _{\alpha _j}{\displaystyle }}{\alpha \mathrm{\Delta }^{}}}`$
$`\alpha _j\alpha \mathrm{\Delta }N_{\alpha ,\alpha _j\alpha }{\displaystyle \frac{s_\alpha s_{\alpha _j\alpha }}{\alpha (q)^2}}.`$
Therefore, on substituting the above expressions into (\*), we obtain the desired equation for $`s`$. $`\mathrm{}`$
In order to solve the equations in (2.3.6) by the factorization method, it is necessary to have precise description of the Lie algebroids and Lie groupoids which are involved. We now begin to describe these geometric objects. Let $`L^{}𝔤`$ be the Lie subalgebra of $`L𝔤`$ consisting of series of the form $`_{\mathrm{}}^0X_nz^n`$. From the explicit expression for $`R`$ in (3.1.8), we have
$$(R^\pm (q)X)(z)=\pm (\mathrm{\Pi }_\pm X)(z)\underset{\alpha \mathrm{\Delta }^{}}{}\frac{(X_1)_\alpha }{\alpha (q)}e_\alpha .$$
$`(\mathrm{3.1.11})`$
Therefore, $`R^+(q)XL^+𝔤`$, while $`R^{}(q)XL_\mathrm{\Delta }^{}^{}𝔤`$, where
$$L_\mathrm{\Delta }^{}^{}𝔤=\{XL^{}𝔤X(\mathrm{})𝔤_\mathrm{\Delta }^{}\}.$$
$`(\mathrm{3.1.12})`$
The proof of the next proposition is obvious and will be left to the reader.
###### Proposition 3.1.5
(a) $`Im^+=_{qU}\{0_q\}\times L^+𝔤\times 𝔥`$. (b) $`^+=_{qU}\{0_q\}\times L^+𝔤\times \{0\}`$ = adjoint bundle of $`Im^+`$.
Remark 3.1.6. Indeed, we also have $`\{^+(0_q,X,0)qU,XL𝔤\}=_{qU}\{0_q\}\times L^+𝔤\times 𝔥.`$ Before we turn to the characterization of $`Im^{}`$, let us recall the notion of a matched pair of Lie algebroids introduced in \[Mok\] as an infinitesimal version of the notion of a matched pair of Lie groupoids \[M2\].(These are generalizations of the corresponding notions for Lie algebras and Lie groups, see \[KSM\], \[LW\],\[Maj\].)
###### Definition Definition 3.1.7
Two Lie algebroids $`A_1`$, $`A_2`$ over the base $`B`$ is said to form a matched pair of Lie algebroids iff the Whitney sum $`W=A_1A_2`$ admits a Lie algebroid structure over the same base with $`A_1`$ and $`A_2`$ as Lie subalgebroids. In this case, the Lie algebroid $`W`$ is called the matched product of $`A_1`$ and $`A_2`$ and is denoted by $`A_1A_2`$.
###### Proposition 3.1.8
(a) The ideal $`^{}`$ is given by
$$^{}=\left\{(0_q,X,0)\right|qU,XL_\mathrm{\Delta }^{}^{}𝔤,X_1𝔥^{}\text{and}\underset{\alpha \mathrm{\Delta }^{}}{}(X_1)_\alpha e_\alpha =ad_q\mathrm{\Pi }_𝔥^{}X_0\}.$$
(b) $`Im^{}=^{}𝒬`$, where
$$𝒬=\{(0_q,\stackrel{~}{Z},Z)qU,Z𝔥,\text{and}\stackrel{~}{Z}(z)=Zz^1\}$$
is a Lie subalgebroid of $`Im^{}`$. Hence $`^{}`$ coincides with the adjoint bundle of the transitive Lie algebroid $`Im^{}`$. Moreover, $`^{}(\{0_q\}\times L𝔤\times \{0\})`$ can be characterized as the set
$$\{(0_q,X,Z)Z𝔥,XL_\mathrm{\Delta }^{}^{}𝔤,\mathrm{\Pi }_𝔥X_0=0\text{and}\mathrm{\Pi }_{𝔤_\mathrm{\Delta }^{}}X_1=Z+ad_qX_0\}$$
where $`\mathrm{\Pi }_{𝔤_\mathrm{\Delta }^{}}`$ is the projection map relative to the decomposition $`𝔤=𝔤_\mathrm{\Delta }^{}(𝔤_\mathrm{\Delta }^{})^{}.`$
###### Demonstration Proof
(a) From the definition of $`^{}`$ in (2.1.10) and the expression for $`R^+(q)`$ in (3.1.11), we have
$`(0_q,X,0)^{}`$
$``$ $`^+(0_q,X,Z)=0\text{for some}Z𝔥`$
$``$ $`\mathrm{\Pi }_𝔥X_1=0,\iota Z+(\mathrm{\Pi }_+X)(z){\displaystyle \underset{\alpha \mathrm{\Delta }^{}}{}}{\displaystyle \frac{(X_1)_\alpha }{\alpha (q)}}e_\alpha =0\text{ for some}Z𝔥`$
$``$ $`XL_\mathrm{\Delta }^{}^{}𝔤,X_1𝔥^{}\text{and}{\displaystyle \underset{\alpha \mathrm{\Delta }^{}}{}}(X_1)_\alpha e_\alpha =ad_q\mathrm{\Pi }_𝔥^{}X_0.`$
Hence the assertion.
(b) For an element $`Z𝔥`$, let $`\stackrel{~}{Z}`$ be the loop given by $`\stackrel{~}{Z}(z)=Zz^1`$. Consider an arbitrary element $`^{}(0_q,X,Z)`$ in $`Im^{}`$. Clearly, it admits the decomposition
$`^{}(0_q,X,Z)`$
$`=`$ $`(0_q,\iota Z\mathrm{\Pi }_{}X+\stackrel{~}{\mathrm{\Pi }_𝔥X_1}{\displaystyle \underset{\alpha \mathrm{\Delta }^{}}{}}{\displaystyle \frac{(X_1)_\alpha }{\alpha (q)}}e_\alpha ,0)`$
$`+(0_q,\stackrel{~}{\mathrm{\Pi }_𝔥X_1},\mathrm{\Pi }_𝔥X_1)`$
where the first term is in $`^{}`$ and the second term is in $`𝒬`$. This shows that $`Im^{}^{}𝒬.`$ Conversely, take an arbitrary element $`(0_q,X,0)+(0_q,\stackrel{~}{Z},Z)^{}𝒬`$ and let $`YL𝔤`$ be defined by $`Y=\mathrm{\Pi }_{}X+\stackrel{~}{Z}.`$ Then from the characterization of $`^{}`$ in part (a), we have
$`^{}(0_q,Y,\mathrm{\Pi }_𝔥X_0)`$
$`=`$ $`(0_q,\mathrm{\Pi }_𝔥X_0\mathrm{\Pi }_{}Y{\displaystyle \underset{\alpha \mathrm{\Delta }^{}}{}}{\displaystyle \frac{(Y_1)_\alpha }{\alpha (q)}}e_\alpha ,\mathrm{\Pi }_𝔥Y_1)`$
$`=`$ $`(0_q,\mathrm{\Pi }_𝔥X_0+\mathrm{\Pi }_{}X\stackrel{~}{Z}+{\displaystyle \underset{\alpha \mathrm{\Delta }^{}}{}}(X_0)_\alpha e_\alpha ,Z)`$
$`=`$ $`(0_q,X,0)+(0_q,\stackrel{~}{Z},Z)`$
and this shows $`^{}𝒬Im^{}.`$ Combining the two inclusions, we conclude that $`Im^{}=^{}𝒬.`$ We shall leave the details of the other assertions to the reader. $`\mathrm{}`$
As a consequence of Proposition 3.1.5, we have
$$\begin{array}{cc}\hfill Im^+/^+& =\underset{qU}{}\{0_q\}\times (L^+𝔤/L^+𝔤)\times 𝔥\hfill \\ & \underset{qU}{}\{0_q\}\times \{0\}\times 𝔥\hfill \end{array}$$
$`(\mathrm{3.1.13})`$
where the identification map is given by
$$(0_q,X+L^+𝔤,Z)(0_q,0,Z).$$
$`(\mathrm{3.1.14})`$
Similarly, it follows from Proposition 3.1.8 that
$$Im^{}/^{}𝒬$$
$`(\mathrm{3.1.15})`$
and the identification map is
$$(0_q,X,Z)+_q^{}(0_q,\stackrel{~}{Z},Z).$$
$`(\mathrm{3.1.16})`$
The following proposition is obvious.
###### Proposition 3.1.9
The isomorphism $`\theta :Im^+/^+Im^{}/^{}`$ defined in Proposition 2.1.2 (b) is given by
$$\theta (0_q,0,Z)=(0_q,\stackrel{~}{Z},Z).$$
Moreover, $`Im(^+,^{})=Im^+\underset{TU}{}Im^{}`$.
## 3.2. Solution of the integrable rational spin Calogero-Moser systems
We begin by solving the equation
$$\begin{array}{cc}& \frac{d}{dt}(q,0,L(q,p,\xi ))\hfill \\ \hfill =& (p,0,[L(q,p,\xi ),R(q)M(q,p,\xi )])\hfill \end{array}$$
$`(\mathrm{3.2.1})`$
where explicitly,
$$M(q,p,\xi )(z)=\frac{1}{z}\left(p+\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi _\alpha }{\alpha (q)}e_\alpha \right)+\frac{\xi }{z^2}.$$
$`(\mathrm{3.2.2})`$
To do so, we have to solve the factorization problem
$$exp\{t(0,0,M(q^0,p^0,\xi ^0))\}(q^0)=\gamma _+(t)\gamma _{}(t)^1$$
$`(\mathrm{3.2.3})`$
for $`(\gamma _+(t),\gamma _{}(t))=((q^0,k_+(t),q(t)),(q^0,k_{}(t),q(t)))Im(^+,^{})`$ satisfying the condition in (2.2.4), where $`(q^0,p^0,\xi ^0)J^1(0)=TU\times (𝒰𝔥^{})`$ is the initial value of $`(q,p,\xi )`$. In what follows, we shall denote by $`LG`$, $`L^+G`$, $`L_1^{}G`$ and $`L_\mathrm{\Delta }^{}^{}G`$ the loop groups corresponding to the Lie algebras $`L𝔤`$, $`L^+𝔤`$, $`L_0^{}𝔤`$ and $`L_\mathrm{\Delta }^{}^{}𝔤`$ respectively. We shall also denote $`k_\pm (t)(z)`$ by $`k_\pm (z,t)`$. Then $`k_+(,t)L^+G`$, while $`k_{}(,t)L_\mathrm{\Delta }^{}^{}G`$ and satisfies additional constraints. From the factorization problem on the Lie groupoid above, it follows that
$$e^{tM(q^0,p^0,\xi ^0)(z)}=k_+(z,t)k_{}(z,t)^1$$
$`(\mathrm{3.2.4})`$
where $`k_\pm (,t)`$ are to be determined. To do so, we recall from the Birkhoff factorization theorem \[PS\] that (at least for small values of $`t`$)
$$e^{tM(q^0,p^0,\xi ^0)(z)}=g_+(z,t)g_{}(z,t)^1$$
$`(\mathrm{3.2.5})`$
for unique $`g_+(,t)L^+G`$ and $`g_{}(,t)L_1^{}G`$. But from (3.2.2) above, it is clear that $`e^{tM(q^0,p^0,\xi ^0)}L_1^{}G`$, so we must have (i) $`g_+1`$, (ii) $`g_{}(z,t)=e^{tM(q^0,p^0,\xi ^0)}(z)`$ for all $`t`$ and consequently,
$$k_+(z,t)k_{}(\mathrm{},t).$$
$`(\mathrm{3.2.6})`$
Thus we have the relation
$$e^{tM(q^0,p^0,\xi ^0)(z)}=k_{}(\mathrm{},t)k_{}(z,t)^1$$
$`(\mathrm{3.2.7})`$
where $`\gamma _{}(t)=(q^0,k_{}(t),q(t))`$ is subject to the condition $`T_{\gamma _{}(t)}𝚕_{\gamma _{}(t)^1}\dot{\gamma }_{}(t)^{}(\{q(t)\}\times \{0\}\times L𝔤).`$ But from the characterization of $`^{}(\{q(t)\}\times \{0\}\times L𝔤)`$ in Proposition 3.1.8 (b), we have
$$\begin{array}{cc}& \mathrm{\Pi }_{𝔤_\mathrm{\Delta }^{}}Res_{z=0}T_{k_{}(z,t)}l_{k_{}(z,t)^1}\dot{k}_{}(z,t)\hfill \\ \hfill =& \dot{q}(t)+ad_{q(t)}T_{k_{}(\mathrm{},t)}l_{k_{}(\mathrm{},t)^1}\dot{k}_{}(\mathrm{},t).\hfill \end{array}$$
$`(\mathrm{3.2.8})`$
On the other hand, by differentiating (3.2.7) with respect to $`t`$, we find
$`T_{k_{}(z,t)}l_{k_{}(z,t)^1}\dot{k}_{}(z,t)T_{k_{}(\mathrm{},t)}l_{k_{}(\mathrm{},t)^1}\dot{k}_{}(\mathrm{},t)`$
$`=`$ $`Ad_{k_{}(\mathrm{},t)^1}M(q^0,p^0,\xi ^0)(z)`$
from which it follows that
$$Res_{z=0}T_{k_{}(z,t)}l_{k_{}(z,t)^1}\dot{k}_{}(z,t)=Ad_{k_{}(\mathrm{},t)^1}\left(p^0+\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi _\alpha ^0}{\alpha (q^0)}e_\alpha \right).$$
$`(\mathrm{3.2.9})`$
Therefore, upon substituting (3.2.9) into (3.2.8), we obtain
$`Ad_{k_{}(\mathrm{},t)}\dot{q}(t)+[T_{k_{}(\mathrm{},t)}r_{k_{}(\mathrm{},t)^1}\dot{k}_{}(\mathrm{},t),Ad_{k_{}(\mathrm{},t)}q(t)]`$
$`=`$ $`L(q^0,p^0,\xi ^0)(\mathrm{}),`$
that is,
$$\frac{d}{dt}Ad_{k_{}(\mathrm{},t)}q(t)=L(q^0,p^0,\xi ^0)(\mathrm{}).$$
$`(\mathrm{3.2.10})`$
Hence the factorization problem boils down to
$$q^0+tL(q^0,p^0,\xi ^0)(\mathrm{})=Ad_{k_{}(\mathrm{},t)}q(t)$$
$`(\mathrm{3.2.11})`$
where $`q(t)`$ and $`k_{}(\mathrm{},t)`$ are to be determined. But from (3.1.4) and the fact that $`𝔤_\mathrm{\Delta }^{}`$ is reductive, we can find (at least for small values of $`t`$) unique $`d(t)H`$ and $`g(t)G_\mathrm{\Delta }^{}`$ (unique up to $`g(t)g(t)\delta (t)`$, where $`\delta (t)H`$) such that
$$q^0+tL(q^0,p^0,\xi ^0)=Ad_{g(t)}d(t)$$
$`(\mathrm{3.2.12})`$
with $`g(0)=1`$, $`d(0)=q^0`$. Hence
$$q(t)=d(t).$$
$`(\mathrm{3.2.13})`$
On the other hand, let us fix one such $`g(t)`$. We shall seek $`k_{}(\mathrm{},t)`$ in the form
$$k_{}(\mathrm{},t)=g(t)h(t),h(t)H.$$
$`(\mathrm{3.2.14})`$
To determine $`h(t)`$, note that the characterization of $`^{}(\{q(t)\}\times \{0\}\times L𝔤)`$ in Proposition 3.1.8 (b) also gives
$$\mathrm{\Pi }_𝔥T_{k_{}(\mathrm{},t)}l_{k_{}(\mathrm{},t)^1}\dot{k}_{}(\mathrm{},t)=0.$$
$`(\mathrm{3.2.15})`$
Using this condition, we find that $`h(t)`$ satisfies the equation
$$\dot{h}(t)=T_el_{h(t)}\left(\mathrm{\Pi }_𝔥(T_{g(t)}l_{g(t)^1}\dot{g}(t))\right)$$
$`(\mathrm{3.2.16})`$
with $`h(0)=1`$. Solving the equation explicitly, we obtain
$$h(t)=exp\left\{_0^t\mathrm{\Pi }_𝔥(T_{g(\tau )}l_{g(\tau )^1}\dot{g}(\tau ))𝑑\tau \right\}.$$
$`(\mathrm{3.2.17})`$
Hence $`k_+(z,t)k_{}(\mathrm{},t)`$ and $`k_{}(z,t)e^{tM(q^0,p^0,\xi ^0)(z)}k_{}(\mathrm{},t)`$ satisfy (3.2.4).
###### Theorem 3.2.1
Let $`(q^0,p^0,\xi ^0)J^1(0)=TU\times (𝒰𝔥^{}).`$ Then the Hamiltonian flow on $`J^1(0)`$ generated by
$$(q,p,\xi )=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{2}\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi _\alpha \xi _\alpha }{\alpha (q)^2}$$
with initial condition $`(q(0),p(0),\xi (0))=(q^0,p^0,\xi ^0)`$ is given by
$$\begin{array}{cc}& q(t)=d(t),\hfill \\ & \xi (t)=Ad_{k_{}(\mathrm{},t)^1}\xi ^0,\hfill \\ & p(t)=\dot{d}(t)=Ad_{k_{}(\mathrm{},t)^1}L(q^0,p^0,\xi ^0)(\mathrm{})\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi (t)_\alpha }{\alpha (q(t))}e_\alpha \hfill \end{array}$$
$`(\mathrm{3.2.18})`$
where $`d(t)`$ and $`k_{}(\mathrm{},t)`$ are constructed from the above procedure.
###### Demonstration Proof
The formulas for $`p(t)`$, $`\xi (t)`$ are obtained by equating the coefficients of $`z^0`$ and $`z^1`$ on both sides of the expression $`L(q(t),p(t),\xi (t))(z)=Ad_{k_{}(\mathrm{},t)^1}L(q^0,p^0,\xi ^0)(z).`$ $`\mathrm{}`$
We now turn to the solution of the associated integrable model on $`TU\times 𝔤_{red}`$ with Hamiltonian $`_0(q,p,s)=\frac{1}{2}_ip^2\frac{1}{2}_{\alpha \mathrm{\Delta }^{}}\frac{s_\alpha s_\alpha }{\alpha (q)^2}`$ and with equations of motion given in Proposition 3.1.4.
###### Corollary 3.2.2
Let $`(q^0,p^0,s^0)TU\times 𝔤_{red}`$ and suppose $`s^0=Ad_{g(\xi ^0)^1}\xi ^0`$ where $`\xi ^0𝒰𝔥^{}.`$ Then the Hamiltonian flow generated by $`_0`$ with initial condition $`(q(0),p(0),s(0))=(q^0,p^0,s^0)`$ is given by
$$\begin{array}{cc}& q(t)=d(t),\hfill \\ & s(t)=Ad_{\left(\stackrel{~}{k}_{}(\mathrm{},t)g\left(Ad_{\stackrel{~}{k}_{}(\mathrm{},t)^1}s^o\right)\right)^1}s^0,\hfill \\ & p(t)=Ad_{\left(\stackrel{~}{k}_{}(\mathrm{},t)g\left(Ad_{\stackrel{~}{k}_{}(\mathrm{},t)^1}s^o\right)\right)^1}L(q^0,p^0,s^0)(\mathrm{})\underset{\alpha \mathrm{\Delta }^{}}{}\frac{s_\alpha (t)}{\alpha (q(t))}e_\alpha .\hfill \end{array}$$
$`(\mathrm{3.2.19})`$
where $`\stackrel{~}{k}_{}(\mathrm{},t)=g(\xi ^0)^1k_{}(\mathrm{},t)g(\xi ^0)`$ depends only on $`s^0`$ and $`k_{}(\mathrm{},t)`$, $`d(t)`$ are as in Theorem 3.2.1.
###### Demonstration proof
We shall obtain the Hamiltonian flow generated by $`_0`$ by Poisson reduction. Using the relation $`\varphi _t^{red}\pi _0=\pi _0\varphi _ti_0`$ from Corollary 2.2.3, we have
$$\varphi _t^{red}(q^0,p^0,s^0)=(q(t),p(t),Ad_{g(\xi (t))^1}\xi (t))$$
where $`q(t)`$ and $`p(t)`$ are given by the expressions in Theorem 3.2.1 above. Thus
$$s(t)=Ad_{g(\xi (t))^1}\xi (t)=Ad_{\left(\stackrel{~}{k}_{}(\mathrm{},t)g\left(Ad_{\stackrel{~}{k}_{}(\mathrm{},t)^1}s^o\right)\right)^1}s^0$$
where we have used the $`H`$-equivariance of the map $`g`$ to show that
$$g(\xi ^0)^1k_{}(\mathrm{},t)g(\xi (t))=\stackrel{~}{k}_{}(\mathrm{},t)g\left(Ad_{\stackrel{~}{k}_{}(\mathrm{},t)^1}s^o\right).$$
To express $`p(t)`$ in the desired form, simply apply $`Ad_{g(\xi (t))^1}`$ to both sides of the expression for $`p(t)`$ in the above theorem, this gives
$$p(t)=Ad_{g(\xi (t))^1k_{}(\mathrm{},t)^1g(\xi ^0)}L(q^0,p^0,s^0)(\mathrm{})\underset{\alpha \mathrm{\Delta }^{}}{}\frac{s_\alpha (t)}{\alpha (q(t))}e_\alpha $$
where we have used the relation $`s_\alpha (t)=e^{\alpha (g(\xi (t)))}\xi _\alpha (t)`$. Hence the desired expression for $`p(t)`$ follows. The assertion that $`\stackrel{~}{k}_{}(\mathrm{},t)`$ depends only on $`s^0`$ is clear. $`\mathrm{}`$
Remark 3.2.3. (a) The expression in (3.2.12) shows that the solution blows up precisely when the factorization fails. However, initial conditions do exist for which the solution exists for all time.
(b) The first example of a rational spin Calogero-Moser system is due to Gibbons and Hermsen \[GH\]. Analogous to what was done there, we can show that
$$\dot{q}=L(q,p,\xi )(\mathrm{})+[q,\underset{\alpha \mathrm{\Delta }^{}}{}\frac{\xi _\alpha }{\alpha (q)^2}e_\alpha ]$$
on $`J^1(0)`$ from which we can also deduce the relation (3.2.12). Thus on the surface, it appears that there is no need to use $`L(q,p,\xi )(z)`$. We remark, however, that our factorization problem (which involves $`L(q^0,p^0,\xi ^0)(z)`$) does carry more information and that we do need $`L(q,p,\xi )(z)`$ in order to establish the Liouville integrability of the reduced system on $`TU\times 𝔤_{red}.`$ In other words, our realization picture embraces both exact solvability and complete integrability. A unifying and representation independent method to establish the Liouville integrability of the integrable spin CM systems in \[LX2\] for all simple Lie algebras will be given in a forthcoming paper.
(c) We now explain the Poisson meaning of the limiting Lax operator $`L(q,p,\xi )(\mathrm{}).`$ To do so, we recall that $`r(q)=_\mathrm{\Delta }^{}\frac{1}{\alpha (q)}e_\alpha e_\alpha `$ is a classical dynamical r-matrix with zero coupling constant in the sense of \[EV\]. Therefore, if we define $`R:UL(𝔤,𝔤)`$ by
$$R(q)\xi =r^{\mathrm{}}(q)\xi =\underset{\alpha ^{}}{}\frac{\xi _\alpha }{\alpha (q)}e_\alpha ,$$
then $`R`$ is a solution of the CDYBE (i.e., (2.1.2) with $`\chi 0`$). Let $`A^{}\mathrm{\Omega }TU\times 𝔤`$ be the coboundary dynamical Lie algebroid associated with $`R`$ and let $`A\mathrm{\Omega }TU\times 𝔤`$ be the trivial Lie algebroid. Then according to \[L2\],
$$:A^{}\mathrm{\Omega }A\mathrm{\Omega },(q,p,\xi )(q,\mathrm{\Pi }_𝔥\xi ,p+R(q)\xi )$$
is a morphism of Lie algebroids. Consequently, the dual map $`^{}`$ is an $`H`$-equivariant Poisson map, when the domain and target are equipped with the corresponding Lie-Poisson structure. Explicitly,
$`^{}(q,p,\xi )=`$ $`(q,\mathrm{\Pi }_𝔥\xi ,pR(q)\xi )`$
$`=`$ $`(q,\mathrm{\Pi }_𝔥\xi ,L(q,p,\xi )(\mathrm{})).`$
Moreover, if we define
$$L^{\mathrm{}}(q,p,\xi )=L(q,p,\xi )(\mathrm{}),$$
then the Hamiltonian $``$ of the rational spin CM system in (3.1.3) is also given by
$$(q,p,\xi )=((L^{\mathrm{}})^{}E)(q,p,\xi )$$
where $`E`$ is the quadratic function on $`𝔤`$ defined by
$$E(\xi )=\frac{1}{2}(\xi ,\xi ).$$
This shows that the Hamiltonian system defined by $``$ admits a second realization in $`A\mathrm{\Omega }`$ and this clarifies the Poisson-geometric meaning of $`L(q,p,\xi )(\mathrm{})`$. We note, however, that the r-matrix $``$ introduced earlier in this remark is degenerate in the sense that it is not associated with a factorization problem.
## 4. The trigonometric spin Calogero-Moser systems
## 4.1. Lax operators, Hamiltonian equations and the Lie subalgebroids
In this section, we take the trigonometric spin Calogero-Moser systems to be the Hamiltonian systems in Definition 2.3.4 associated to the following trigonometric dynamical r-matrices with spectral parameter:
$$r(q,z)=c(z)\underset{i}{}x_ix_i\underset{\alpha \mathrm{\Delta }}{}\varphi _\alpha (q,z)e_\alpha e_\alpha $$
$`(\mathrm{4.1.1})`$
where
$$c(z)=\mathrm{cot}z$$
$`(\mathrm{4.1.2})`$
and
$$\varphi _\alpha (q,z)=\{\begin{array}{cc}\frac{\mathrm{sin}(\alpha (q)+z)}{\mathrm{sin}\alpha (q)\mathrm{sin}z},\hfill & \alpha <\pi ^{}>\hfill \\ \frac{e^{iz}}{\mathrm{sin}z},\hfill & \alpha \overline{\pi }_{}^{}{}_{}{}^{+}\hfill \\ \frac{e^{iz}}{\mathrm{sin}z},\hfill & \alpha \overline{\pi }_{}^{}{}_{}{}^{}.\hfill \end{array}$$
$`(\mathrm{4.1.3})`$
In (4.1.3) above, $`\pi ^{}`$ is an arbitrary subset of a fixed simple system $`\pi \mathrm{\Delta }`$, $`<\pi ^{}>`$ is the root span of $`\pi ^{}`$ and $`\overline{\pi }_{}^{}{}_{}{}^{\pm }=\mathrm{\Delta }^\pm <\pi ^{}>^\pm .`$ Accordingly, the Lax operators are given by
$$\begin{array}{cc}\hfill L(q,p,\xi )(z)=& p+c(z)\underset{i}{}\xi _ix_i\underset{\alpha \mathrm{\Delta }}{}\varphi _\alpha (q,z)\xi _\alpha e_\alpha \hfill \\ \hfill =& p+\underset{\alpha <\pi ^{}>}{}c(\alpha (q))\xi _\alpha i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{+}}{}\xi _\alpha e_\alpha \hfill \\ & +i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{}}{}\xi _\alpha e_\alpha +c(z)\xi .\hfill \end{array}$$
$`(\mathrm{4.1.4})`$
Hence we have a family of dynamical systems parametrized by subsets $`\pi ^{}`$ of $`\pi `$ with Hamiltonians of the form:
$$\begin{array}{cc}\hfill (q,p,\xi )=& \frac{1}{2}\underset{i}{}p_i^2\frac{1}{2}\underset{\alpha <\pi ^{}>}{}\left(\frac{1}{\mathrm{sin}^2\alpha (q)}\frac{1}{3}\right)\xi _\alpha \xi _\alpha \frac{5}{6}\underset{\alpha \mathrm{\Delta }<\pi ^{}>}{}\xi _\alpha \xi _\alpha \hfill \\ & \frac{1}{3}\underset{i}{}\xi _i^2.\hfill \end{array}$$
$`(\mathrm{4.1.5})`$
Remark 4.1.1. The trigonometric dynamical r-matrices in (4.1.1) are gauge equivalent to those used in \[LX2\]. Although the corresponding Hamiltonians in (4.1.5) above contain the additional term $`\frac{1}{3}_i\xi _i^2`$, however, the Hamiltonian flows on $`J^1(0)`$ and the reduced systems are the same as those in \[LX2\]. The reason why we use the dynamical r-matrices in (4.1.1) is due to the fact that the corresponding Lie subalgebroids $`Im^\pm `$ are simpler to analyze.
The next two propositions follow from direct calculation, as in Propositions 3.1.1 and 3.1.4. We shall leave the proof to the reader.
###### Proposition 4.1.2
The Hamiltonian equations of motion generated by $``$ on $`TU\times 𝔤`$ are given by
$$\begin{array}{cc}& \dot{q}=p,\hfill \\ & \dot{p}=\underset{\alpha <\pi ^{}>}{}\frac{\mathrm{cot}\alpha (q)}{\mathrm{sin}^2\alpha (q)}\xi _\alpha \xi _\alpha H_\alpha ,\hfill \\ & \dot{\xi }=[\xi ,\frac{2}{3}\mathrm{\Pi }_𝔥\xi \underset{\alpha <\pi ^{}>}{}\left(\frac{1}{\mathrm{sin}^2\alpha (q)}\frac{1}{3}\right)\xi _\alpha e_\alpha \frac{5}{3}\underset{\alpha \mathrm{\Delta }<\pi ^{}>}{}\xi _\alpha e_\alpha ].\hfill \end{array}$$
$`(\mathrm{4.1.6})`$
###### Proposition 4.1.3
The Hamiltonian equations of motion generated by
$$_0(q,p,s)=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{2}\underset{\alpha <\pi ^{}>}{}\left(\frac{1}{\mathrm{sin}^2\alpha (q)}\frac{1}{3}\right)s_\alpha s_\alpha \frac{5}{6}\underset{\alpha \mathrm{\Delta }<\pi ^{}>}{}s_\alpha s_\alpha $$
on the reduced Poisson manifold $`TU\times 𝔤_{red}`$ are given by
$`\dot{q}=p,`$
$`\dot{p}={\displaystyle \underset{\alpha \mathrm{\Delta }^{}}{}}{\displaystyle \frac{\mathrm{cot}\alpha (q)}{\mathrm{sin}^2\alpha (q)}}s_\alpha s_\alpha H_\alpha ,`$
$`\dot{s}=[s,]`$
where
$`=`$ $`{\displaystyle \underset{\alpha <\pi ^{}>}{}}\left({\displaystyle \frac{1}{\mathrm{sin}^2\alpha (q)}}{\displaystyle \frac{1}{3}}\right)s_\alpha e_\alpha {\displaystyle \frac{5}{3}}{\displaystyle \underset{\alpha \mathrm{\Delta }<\pi ^{}>}{}}s_\alpha e_\alpha `$
$`{\displaystyle \frac{+{\displaystyle \underset{i,j}{}}C_{ji}{\displaystyle }}{\alpha <\pi ^{}>\pi ^{}}}`$
$`\alpha _j\alpha \mathrm{\Delta }N_{\alpha ,\alpha _j\alpha }\left({\displaystyle \frac{1}{\mathrm{sin}^2\alpha (q)}}{\displaystyle \frac{1}{3}}\right)s_\alpha s_{\alpha _j\alpha }h_{\alpha _i}`$
$`{\displaystyle \frac{+{\displaystyle \frac{5}{3}}{\displaystyle \underset{i,j}{}}C_{ji}{\displaystyle }}{\alpha \mathrm{\Delta }<\pi ^{}>}}`$
$`\alpha _j\alpha \mathrm{\Delta }N_{\alpha ,\alpha _j\alpha }s_\alpha s_{\alpha _j\alpha }h_{\alpha _i}.`$
(Here the notation $`N_{\alpha ,\beta }`$ is as in Proposition 3.1.4.)
###### Proposition 4.1.4
The classical dynamical r-matrix $`R`$ associated with the trigonometric dynamical r-matrix with spectral parameter in (4.1.1) is given by
$$\begin{array}{cc}\hfill (R(q)X)(z)& =\frac{1}{2}X(z)+\underset{k=0}{\overset{\mathrm{}}{}}\frac{c^{(k)}(z)}{k!}X_{(k+1)}\underset{\alpha <\pi ^{}>}{}c(\alpha (q))(X_1)_\alpha e_\alpha \hfill \\ & +i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{+}}{}(X_1)_\alpha e_\alpha i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{}}{}(X_1)_\alpha e_\alpha .\hfill \end{array}$$
$`(\mathrm{4.1.7})`$
###### Demonstration Proof
The formula follows from (2.3.1) and (4.1.1) by a direct calculation where we have used the formula $`\frac{d^k}{dw^k}|_{w=0}\varphi _\alpha (q,zw)=c^{(k)}(z),k1`$. $`\mathrm{}`$
###### Corollary 4.1.5
On $`J^1(0)`$, we have
$$\begin{array}{cc}& (R(q)M(q,p,\xi ))(z)\hfill \\ \hfill =& \frac{1}{2}M(q,p,\xi )(z)c(z)p+\underset{\alpha \mathrm{\Delta }<\pi ^{}>}{}\varphi _\alpha (q,z)c(z)\xi _\alpha e_\alpha \hfill \\ & +\underset{\alpha <\pi ^{}>}{}\varphi _\alpha (q,z)(c(\alpha (q))+c(z)c(\alpha (q)+z))\xi _\alpha e_\alpha \hfill \end{array}$$
$`(\mathrm{4.1.8})`$
where $`M(q,p,\xi )(z)=L(q,p,\xi )(z)/z.`$
###### Demonstration Proof
The formula in (4.1.8) follows from (4.1.7) by algebra on using the following expansion in a deleted neighborhood of 0:
$$M(q,p,\xi )(z)=\frac{\xi }{z^2}+\frac{1}{z}M(q,p,\xi )_1+O(1),$$
where
$`M(q,p,\xi )_1`$
$`=`$ $`p+{\displaystyle \underset{\alpha <\pi ^{}>}{}}c(\alpha (q))\xi _\alpha e_\alpha i{\displaystyle \underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{+}}{}}\xi _\alpha e_\alpha +i{\displaystyle \underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{}}{}}\xi _\alpha e_\alpha .`$
$`\mathrm{}`$
Our next lemma is obvious from (4.1.7) and the expansion $`c^{(k)}(z)=k!z^{(k+1)}+O(1),k0,`$ in a deleted neighborhood of $`0.`$
###### Lemma 4.1.6
For $`XL𝔤`$, $`R^+(q)XL𝔤.`$
###### Lemma 4.1.7
(a) For $`XL𝔤`$, $`R^{}(q)X`$ has singularities at the points of the rank one lattice $`\pi `$, and is holomorphic in $`\pi .`$ Moreover, $`R^{}(q)X`$ is simply-periodic with period $`\pi `$. (b) The principal part of $`R^{}(q)X`$ at $`z=0`$ is $`(\mathrm{\Pi }_{}X)(z).`$ (c) $`R^{}(q)X`$ is bounded as $`z\mathrm{}`$ in a period strip with
$$\underset{y\mathrm{}}{lim}(R^{}(q)X)(x+iy)=i\mathrm{\Pi }_𝔥X_1+\underset{\alpha <\pi ^{}>}{}(ic(\alpha (q)))(X_1)_\alpha e_\alpha +2i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{+}}{}(X_1)_\alpha e_\alpha ,$$
$$\underset{y\mathrm{}}{lim}(R^{}(q)X)(xiy)=i\mathrm{\Pi }_𝔥X_1\underset{\alpha <\pi ^{}>}{}(i+c(\alpha (q)))(X_1)_\alpha e_\alpha 2i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{}}{}(X_1)_\alpha e_\alpha .$$
###### Demonstration Proof
(a) Clearly, $`c^{(k)}(z)`$ are periodic with period $`\pi `$ and meromorphic in $``$ with poles at the points of the rank one lattice $`\pi `$. Therefore, the assertion follows.
(b) This follows from the property that for $`k0`$, we have $`c^{(k)}(z)=k!z^{(k+1)}+O(1)`$ in a deleted neighborhood of $`z=0`$.
(c) First of all, note that $`lim_{y\pm \mathrm{}}\mathrm{cot}(x+iy)=i.`$ On the other hand, it is easy to check that the derivatives of $`\mathrm{cot}z`$ always contain $`\mathrm{csc}^2z`$ as a factor. Therefore, we have $`lim_{y\pm \mathrm{}}c^{(k)}(x+iy)=0`$ for $`k1.`$ The formulas for $`lim_y\mathrm{}(R^{}(q)X)(x\pm iy)`$ are now obvious from (4.1.7). $`\mathrm{}`$
In order to describe the membership of the elements $`(R^{}(q)X)(\pm i\mathrm{})`$ in Lemma 4.1.7 (c) and for subsequent analysis, we need to introduce a number of Lie subalgebras of $`𝔤`$ and their corresponding Lie groups. To begin with, let $`𝔟^{}=𝔥+_{\alpha \mathrm{\Delta }^{}}𝔤_\alpha `$ and $`𝔟^+=𝔥+_{\alpha \mathrm{\Delta }^+}𝔤_\alpha `$ be opposing Borel subalgebras of $`𝔤.`$ Then for each $`\pi ^{}\pi `$, we have the parabolic subalgebras
$$𝔭_\pi ^{}^\pm =𝔟^\pm +\underset{\alpha <\pi ^{}>^{}}{}𝔤_\alpha .$$
$`(\mathrm{4.1.9})`$
Recall that $`𝔭_\pi ^{}^\pm `$ admit the following direct sum decomposition \[Kn\]
$$𝔭_\pi ^{}^\pm =𝔤_\pi ^{}+𝔫_\pi ^{}^\pm $$
$`(\mathrm{4.1.10})`$
where
$$𝔤_\pi ^{}=𝔥+\underset{\alpha <\pi ^{}>}{}𝔤_\alpha $$
$`(\mathrm{4.1.11})`$
is the Levi factor of $`𝔭_\pi ^{}^\pm `$, and
$$𝔫_\pi ^{}^\pm =\underset{\alpha \overline{\pi }^^\pm }{}𝔤_\alpha $$
$`(\mathrm{4.1.12})`$
are the nilpotent radicals. We shall denote by $`\mathrm{\Pi }_{𝔤_\pi ^{}}^\pm `$ the projection maps onto $`𝔤_\pi ^{}`$ relative to the splitting $`𝔭_\pi ^{}^\pm =𝔤_\pi ^{}+𝔫_\pi ^{}^\pm `$. On the other hand, the connected and simply-connected Lie subgroups of $`G`$ with corresponding Lie subalgebras $`𝔭_\pi ^{}^\pm `$, $`𝔤_\pi ^{}`$, and $`𝔫_\pi ^{}^\pm `$ will be denoted respectively by $`P_\pi ^{}^\pm `$, $`G_\pi ^{}`$, and $`N_\pi ^{}^\pm `$ and we have $`P_\pi ^{}^\pm =N_\pi ^{}^\pm G_\pi ^{}`$.
Thus it follows from Lemma 4.1.7 (c) that $`(R^{}(q)X)(\pm i\mathrm{})𝔭_\pi ^{}^\pm `$. The proof of the our next proposition is obvious.
###### Proposition 4.1.8
(a) $`Im^+=_{qU}\{0_q\}\times L^+𝔤\times 𝔥`$. (b) $`^+=_{qU}\{0_q\}\times L^+𝔤\times \{0\}`$ = adjoint bundle of $`Im^+`$.
Remark 4.1.9. Indeed, in going through the proof of Proposition 4.1.8 (a) above, one can show that
$$\{^+(0_q,X,0)qU,XL𝔤\}=\underset{qU}{}\{0_q\}\times L^+𝔤\times 𝔥.$$
###### Proposition 4.1.10
$`Im^{}=^{}𝒬,`$ where
$$𝒬=\{(0_q,c()Z,Z)qU,Z𝔥\}$$
$`(\mathrm{4.1.13})`$
is a Lie subalgebroid of $`Im^{}`$ and the ideal $`^{}`$ coincides with the adjoint bundle of $`Im^{}`$ and admits the following characterization:
$$(0_q,X,0)_q^{}\text{if and only if}$$
(a) $`X`$ is holomorphic in $`\pi `$ with singularities at the points of the rank one lattice $`\pi `$,
(b) $`X(z)`$ is periodic with period $`\pi `$,
(c) $`\mathrm{\Pi }_𝔥X_1=0,`$
(d) $`X`$ is bounded as $`z\mathrm{}`$ in a period strip with
$$\underset{y\mathrm{}}{lim}X(x+iy)=\iota Z\underset{\alpha <\pi ^{}>}{}(ic(\alpha (q)))(X_1)_\alpha e_\alpha 2i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{+}}{}(X_1)_\alpha e_\alpha ,$$
$$\underset{y\mathrm{}}{lim}X(xiy)=\iota Z+\underset{\alpha <\pi ^{}>}{}(i+c(\alpha (q)))(X_1)_\alpha e_\alpha +2i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{}}{}(X_1)_\alpha e_\alpha ,$$
for some $`Z𝔥`$. Consequently, $`X(\pm i\mathrm{})𝔭_\pi ^{}^\pm `$ and
$$\mathrm{\Pi }_{𝔤_\pi ^{}}^{}X(i\mathrm{})=Ad_{e^{2iq}}\mathrm{\Pi }_{𝔤_\pi ^{}}^+X(i\mathrm{}).$$
$`(\mathrm{4.1.14})`$
###### Demonstration Proof
From the definition of $`^{}`$ , we have
$`(0_q,X,0)^{}`$
$``$ $`^+(0_q,X,Z)=0\text{for some}Z𝔥`$
$``$ $`\mathrm{\Pi }_𝔥X_1=0,\iota Z+R^+(q)X=0\text{for some}Z𝔥`$
$``$ $`\mathrm{\Pi }_𝔥X_1=0,X(z)=\iota Z(R^{}(q)X+c()\mathrm{\Pi }_𝔥X_1)(z)\text{for some}Z𝔥.`$
Therefore, by Lemma 4.1.7 above and the relation $`c(\alpha (q))+i=e^{2i\alpha (q)}(c(\alpha (q))i)`$, we conclude that $`X`$ satisfies the properties in (a)-(d). Conversely, suppose $`XL𝔤`$ satisfies the properties in (a)-(d). Consider
$$D(z)=X(z)+(R^{}(q)X)(z).$$
Then by the properties of $`X`$ and Lemma 4.1.7, $`D(z+\pi )=D(z)`$ and the principal part of $`D`$ at $`z=0`$ is zero. Therefore, $`D`$ extends to a holomorphic map from $``$ to $`𝔤`$. Moreover, $`D`$ is bounded as $`z\mathrm{}`$ in the period strip and $`lim_y\mathrm{}D(x\pm iy)=\iota Z`$. Write $`D(z)=_jd_j(z)x_j+_{\alpha \mathrm{\Delta }}d_\alpha (z)e_\alpha .`$ Then $`d_j`$ and $`d_\alpha `$ are entire functions for $`1jN,\alpha \mathrm{\Delta }`$ and are periodic with period $`\pi `$. Therefore, when we combine this with the boundedness of $`d_j`$ and $`d_\alpha `$ as $`z\mathrm{}`$ in the period strip, we conclude that $`d_j(z)=d_j(=\text{constant})`$ for each $`j`$ and $`d_\alpha (z)=d_\alpha (=\text{constant})`$ for each $`\alpha `$. But now it follows from $`lim_y\mathrm{}D(x\pm iy)=\iota Z`$ that we must have $`D(z)=_jd_jx_j=\iota Z`$. Consequently, $`X=\iota ZR^{}(q)X`$ and this in turn implies that $`\iota Z+R^+(q)X=0`$. As $`\mathrm{\Pi }_𝔥X_1=0`$, we have $`(0_q,X,0)_q^{}`$, as was to be proved. The proof of the assertion $`Im^{}=^{}𝒬`$ is similar to the one of Proposition 5.1.9 and so we will omit the details. $`\mathrm{}`$
## 4.2. Solution of the integrable trigonometric spin Calogero-Moserfakkesystems
In principle, we have to solve the factorization problem
$$exp\{t(0,0,M(q^0,p^0,\xi ^0))\}(q^0)=\gamma _+(t)\gamma _{}(t)^1$$
$`(\mathrm{4.2.1})`$
for $`(\gamma _+(t),\gamma _{}(t))=((q^0,k_+(t),q(t)),(q^0,k_{}(t),q(t)))Im(^+,^{})`$ satisfying the condition
$$(T_{\gamma _+(t)}𝚕_{\gamma _+(t)^1}\dot{\gamma }_+(t),T_{\gamma _{}(t)}𝚕_{\gamma _{}(t)^1}\dot{\gamma }_{}(t))(^+,^{})(\{q(t)\}\times \{0\}\times L𝔤),$$
$`(\mathrm{4.2.2})`$
where $`(q^0,p^0,\xi ^0)J^1(0)=TU\times (𝒰𝔥^{})`$ is the initial value of $`(q,p,\xi )`$ and $`M(q^0,p^0,\xi ^0)(z)=L(q^0,p^0,\xi ^0)(z)/z.`$ (We shall denote $`k_\pm (t)(z)`$ by $`k_\pm (z,t)`$.) However, as we shall see in the next two propositions and their corollary, it actually suffices to solve for $`q(t)`$, $`k_+(0,t)`$ and $`k_{}(\pm i\mathrm{},t)`$ and we will find the factorization problems for these quantities from (4.2.1) and (4.2.2) above. In what follows, we shall denote by $`(q(t),p(t),\xi (t))`$ the Hamiltonian flow on $`J^1(0)`$ generated by $``$ with initial condition $`(q(0),p(0),\xi (0))=(q^0,p^0,\xi ^0)`$.
###### Proposition 4.2.1
With the notations introduced above, (a) $`L(q(t),p(t),\xi (t))(\pm i\mathrm{})`$ exist. Explicitly,
$$\begin{array}{cc}& L(q(t),p(t),\xi (t))(\pm i\mathrm{})\hfill \\ \hfill =& L(q(t),p(t),\xi (t))(\pi /2)i\xi (t)\hfill \\ \hfill =& p(t)+\underset{\alpha <\pi ^{}>}{}(c(\alpha (q(t)))i)\xi _\alpha (t)e_\alpha 2i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{\pm }}{}\xi _\alpha (t)e_\alpha \hfill \end{array}$$
$`(\mathrm{4.2.3})`$
and therefore $`L(q(t),p(t),\xi (t))(\pm i\mathrm{})𝔭_\pi ^{}^\pm .`$ (b) $`L(q(t),p(t),\xi (t))(\pm i\mathrm{})`$ satisfy the Lax equations
$$\begin{array}{cc}& \frac{d}{dt}L(q(t),p(t),\xi (t))(\pm i\mathrm{})\hfill \\ \hfill =& [L(q(t),p(t),\xi (t))(\pm i\mathrm{}),\underset{\alpha <\pi ^{}>}{}\mathrm{csc}^2(\alpha (q(t)))\xi _\alpha (t)e_\alpha ].\hfill \end{array}$$
$`(\mathrm{4.2.4})`$
###### Demonstration Proof
(a) The existence of $`L(q(t),p(t),\xi (t))(\pm i\mathrm{})`$ and their explicit formulas are obtained from (4.1.4) by noting that $`lim_{y\pm \mathrm{}}\mathrm{cot}(iy)=i.`$
(b) According to Proposition 2.3.5, we have
$$\dot{L}(q(t),p(t),\xi (t))(z)=[L(q(t),p(t),\xi (t))(z),(R^{}(q(t))M(q(t),p(t),\xi (t)))(z)]$$
where
$`(R^{}(q(t))M(q(t),p(t),\xi (t)))(z)`$
$`=`$ $`c(z)p(t)+{\displaystyle \underset{\alpha \mathrm{\Delta }<\pi ^{}>}{}}\varphi _\alpha (q(t),z)c(z)\xi _\alpha (t)e_\alpha `$
$`+{\displaystyle \underset{\alpha <\pi ^{}>}{}}\varphi _\alpha (q(t),z)(c(\alpha (q(t)))+c(z)c(\alpha (q(t))+z))\xi _\alpha (t)e_\alpha `$
by (4.1.8). Now, it is easy to see from (4.1.3) that
$$\varphi _\alpha (q(t),i\mathrm{})=\{\begin{array}{cc}(i+\mathrm{cot}\alpha (q(t))),\hfill & \alpha <\pi ^{}>\hfill \\ 2i,\hfill & \alpha \overline{\pi }_{}^{}{}_{}{}^{+}\hfill \\ 0,\hfill & \alpha \overline{\pi }_{}^{}{}_{}{}^{}\hfill \end{array}$$
whereas
$$\varphi _\alpha (q(t),i\mathrm{})=\{\begin{array}{cc}(i+\mathrm{cot}\alpha (q(t))),\hfill & \alpha <\pi ^{}>\hfill \\ 0,\hfill & \alpha \overline{\pi }_{}^{}{}_{}{}^{+}\hfill \\ 2i,\hfill & \alpha \overline{\pi }_{}^{}{}_{}{}^{}.\hfill \end{array}$$
Therefore, upon taking the limit as $`z=iy\pm i\mathrm{}`$ in the above expression for $`(R^{}(q(t))M(q(t),p(t),\xi (t)))(z),`$ we find that
$`(R^{}(q(t))M(q(t),p(t),\xi (t)))(\pm i\mathrm{})`$
$`=`$ $`\pm iL(q(t),p(t),\xi (t))(\pm i\mathrm{}){\displaystyle \underset{\alpha <\pi ^{}>}{}}\mathrm{csc}^2(\alpha (q(t)))\xi _\alpha (t)e_\alpha `$
from which the assertion follows. $`\mathrm{}`$
Remark 4.2.2. Although $`L(q,p,\xi )(\pm i\mathrm{})`$ exist and satisfy Lax equations, however, they are deficient in the sense that they do not provide enough conserved quantities for complete integrability. In order to establish Liouville integrability, we must use the Lax operator with spectral parameter $`L(q,p,\xi )(z)`$.
We next spell out some of the consequences of the condition in (4.2.2) which will clarify the relation between the term $`_{\alpha <\pi ^{}>}\mathrm{csc}^2(\alpha (q(t)))\xi _\alpha (t)e_\alpha `$ which appears in the Lax equations above for $`L(q(t),p(t),\xi (t))(\pm i\mathrm{})`$ and the factors $`k_\pm (z,t).`$
###### Proposition 4.2.3
(a) $`k_+(0,t)G_\pi ^{}`$ and satisfies the equation
$$T_{k_+(0,t)}l_{k_+(0,t)^1}\dot{k}_+(0,t)=\underset{\alpha <\pi ^{}>}{}\mathrm{csc}^2(\alpha (q(t)))\xi _\alpha (t)e_\alpha .$$
(b) $`k_{}(\pm i\mathrm{},t)P_\pi ^{}^\pm `$ and satisfy the equations
$`T_{k_{}(\pm i\mathrm{},t)}l_{k_{}(\pm i\mathrm{},t)^1}\dot{k}_{}(\pm i\mathrm{},t)`$
$`=\pm iL(q(t),p(t),\xi (t))(\pm i\mathrm{}){\displaystyle \underset{\alpha <\pi ^{}>}{}}\mathrm{csc}^2(\alpha (q(t)))\xi _\alpha (t)e_\alpha .`$
###### Demonstration Proof
(a) It follows from (4.2.1) and (4.2.2) that (see the proof of Theorem 2.2.2 in \[L2\])
$$T_{\gamma _+(t)}𝚕_{\gamma _+(t)^1}\dot{\gamma }_+(t)=^+(q(t),0,M(q(t),p(t),\xi (t))).$$
Consequently, we have
$`T_{k_+(z,t)}l_{k_+(z,t)^1}\dot{k}_+(z,t)`$
$`=`$ $`(R^+(q(t))M(q(t),p(t),\xi (t)))(z)`$
$`=`$ $`M(q(t),p(t),\xi (t))(z)c(z)p(t)+{\displaystyle \underset{\alpha \mathrm{\Delta }<\pi ^{}>}{}}\varphi _\alpha (q(t),z)c(z)\xi _\alpha (t)e_\alpha `$
$`+{\displaystyle \underset{\alpha <\pi ^{}>}{}}\varphi _\alpha (q(t),z)(c(\alpha (q(t)))+c(z)c(\alpha (q(t))+z))\xi _\alpha (t)e_\alpha .()`$
Since $`\mathrm{cot}z=\frac{1}{z}+O(z^3)`$ in a deleted neighborhood of $`0`$, the $`z^0`$ term in the Laurent series expansion about $`0`$ of $`M(q(t),p(t),\xi (t))(z)`$, $`c(z)p(t)`$ and $`\varphi _\alpha (q(t),z)c(z)`$ for $`\alpha \mathrm{\Delta }<\pi ^{}>`$ is equal to zero in each case. On the other hand, for $`\alpha <\pi ^{}>`$, the $`z^0`$ term in the Laurent series expansion of $`\varphi _\alpha (q(t),z)(c(\alpha (q(t)))+c(z)c(\alpha (q(t))+z))`$ about $`0`$ is $`\mathrm{csc}^2(\alpha (q(t)))`$. The formula for $`T_{k_+(0,t)}l_{k_+(0,t)^1}\dot{k}_+(0,t)`$ thus follows when we let $`z0`$ in (\*) above.
(b) It also follows from (4.2.1) and (4.2.2) that (see the proof of Theorem 2.2.2 of \[L2\])
$$T_{\gamma _{}(t)}𝚕_{\gamma _{}(t)^1}\dot{\gamma }_{}(t)=^{}(q(t),0,M(q(t),p(t),\xi (t)))$$
and hence
$`T_{k_{}(z,t)}l_{k_{}(z,t)^1}\dot{k}_{}(z,t)`$
$`=`$ $`(R^{}(q(t))M(q(t),p(t),\xi (t)))(z).`$
The formulas for $`T_{k_{}(\pm i\mathrm{},t)}l_{k_{}(\pm i\mathrm{},t)^1}\dot{k}_{}(\pm i\mathrm{},t)`$ then follow from the proof of Proposition 4.2.1 (b). Finally the assertion that $`k_{}(\pm i\mathrm{},t)P_\pi ^{}^\pm `$ is a consequence of these formulas and Proposition 4.2.1 (a). $`\mathrm{}`$
###### Corollary 4.2.4
In terms of $`k_+(0,t)`$, we have
$$L(q(t),p(t),\xi (t))(\pm i\mathrm{})=Ad_{k_+(0,t)^1}L(q^0,p^0,\xi ^0)(\pm i\mathrm{}).$$
Consequently,
$$L(q(t),p(t),\xi (t))(z)=Ad_{k_+(0,t)^1}L(q^0,p^0,\xi ^0)(z).$$
###### Demonstration Proof
By using Proposition 4.2.3 (a) and Proposition 4.2.1 (b), we can check that $`Ad_{k_+(0,t)}L(q(t),p(t),\xi (t))(\pm i\mathrm{})`$ are constants, hence
$$Ad_{k_+(0,t)}L(q(t),p(t),\xi (t))(\pm i\mathrm{})=L(q^0,p^0,\xi ^0)(\pm i\mathrm{}).$$
Now it is clear from (4.2.3) that
$$2L(q(t),p(t),\xi (t))(\pi /2)=L(q(t),p(t),\xi (t))(i\mathrm{})+L(q(t),p(t),\xi (t))(i\mathrm{})$$
and
$$2i\xi (t)=L(q(t),p(t),\xi (t))(i\mathrm{})L(q(t),p(t),\xi (t))(i\mathrm{}).$$
As
$$L(q(t),p(t),\xi (t))(z)=L(q(t),p(t),\xi (t))(\pi /2)+c(z)\xi (t),$$
the second assertion is a consequence of the first one by virtue of the above relations. $`\mathrm{}`$
Combining the formulas in Proposition 4.2.3 (a) and (b), and the fact that
$$L(q(t),p(t),\xi (t))(\pm i\mathrm{})=Ad_{k_{}(\pm i\mathrm{},t)^1}L(q^0,p^0,\xi ^0)(\pm i\mathrm{}),$$
$`(\mathrm{4.2.5})`$
we obtain the following factorization problems on $`P_\pi ^{}^\pm `$ :
$$e^{itL(q^0,p^0,\xi ^0)(i\mathrm{})}=k_{}(i\mathrm{},t)k_+(0,t)^1,$$
$`(\mathrm{4.2.6})`$
$$e^{itL(q^0,p^0,\xi ^0)(i\mathrm{})}=k_{}(i\mathrm{},t)k_+(0,t)^1$$
$`(\mathrm{4.2.7})`$
where $`k_+(0,t)`$ and $`k_{}(\pm i\mathrm{})`$ are to be determined. The nature of these factorization problems are of course quite different from that of those in the well-known group-theoretic scheme for constant r-matrices.(Compare, for example, the factorization problems in \[RSTS\],\[STS\], \[DLT\] with our solution of (4.2.6), (4.2.7) below.)
We shall use the following notation: for $`g^\pm P_\pi ^{}^\pm `$, $`𝝂^\pm (g^\pm )N_\pi ^{}^\pm `$, $`𝝀^\pm (g^\pm )G_\pi ^{}`$ will denote the factors in the unique factorization $`g^\pm =𝝂^\pm (g^\pm )𝝀^\pm (g^\pm )`$. In order to solve (4.2.6) and (4.2.7), note that
$$(q^0,k_{}(t),q(t))=(q^0,\widehat{k}_{}(t),q^0)(q^0,e^{c()(q^0q(t))},q(t))$$
$`(\mathrm{4.2.8})`$
by the global version of Proposition 4.1.10 where $`(q^0,\widehat{k}_{}(t),q^0)`$ is in the Lie group bundle integrating $`^{}`$ and the second factor $`(q^0,e^{c()(q^0q(t))},q(t))`$ is in the Lie groupoid integrating $`𝒬`$. Consequently, the factorization problems on $`P_\pi ^{}^\pm `$ in (4.2.6) and (4.2.7) can be recast in the form
$$e^{itL(q^0,p^0,\xi ^0)(i\mathrm{})}=𝝂^+(\widehat{k}_{}(i\mathrm{},t))𝝀^+(\widehat{k}_{}(i\mathrm{},t))e^{i(q^0q(t))}k_+(0,t)^1,$$
$`(\mathrm{4.2.9})`$
$$e^{itL(q^0,p^0,\xi ^0)(i\mathrm{})}=𝝂^{}(\widehat{k}_{}(i\mathrm{},t))𝝀^{}(\widehat{k}_{}(i\mathrm{},t))e^{i(q^0q(t))}k_+(0,t)^1.$$
$`(\mathrm{4.2.10})`$
Now, from the fact that $`e^{itL(q^0,p^0,\xi ^0)(i\mathrm{})}P_\pi ^{}^+`$, we can find unique $`n_+(t)N_\pi ^{}^+`$, $`g_+(t)G_\pi ^{}`$ satisfying $`n_+(0)=g_+(0)=1`$ such that
$$e^{itL(q^0,p^0,\xi ^0)(i\mathrm{})}=n_+(t)g_+(t).$$
$`(\mathrm{4.2.11})`$
Similarly, we can find unique $`n_{}(t)N_\pi ^{}^{}`$, $`g_{}(t)G_\pi ^{}`$ satisfying $`n_{}(0)=g_{}(0)=1`$ such that
$$e^{itL(q^0,p^0,\xi ^0)(i\mathrm{})}=n_{}(t)g_{}(t).$$
$`(\mathrm{4.2.12})`$
By comparing (4.2.9) (resp. (4.2.10))with (4.2.11) (resp. (4.2.12)), we obtain
$$𝝂^+(\widehat{k}_{}(i\mathrm{},t))=n_+(t),𝝂^{}(\widehat{k}_{}(i\mathrm{},t))=n_{}(t).$$
$`(\mathrm{4.2.13})`$
Hence the factorization problems reduce to
$$g_+(t)=𝝀^+(\widehat{k}_{}(i\mathrm{},t))e^{i(q^0q(t))}k_+(0,t)^1,$$
$`(\mathrm{4.2.14})`$
$$g_{}(t)=𝝀^{}(\widehat{k}_{}(i\mathrm{},t))e^{i(q^0q(t))}k_+(0,t)^1.$$
$`(\mathrm{4.2.15})`$
But from the global version of (4.1.14), we have
$$𝝀^{}(\widehat{k}_{}(i\mathrm{},t))=e^{2iq^0}𝝀^+(\widehat{k}_{}(i\mathrm{},t))e^{2iq^0}.$$
$`(\mathrm{4.2.16})`$
Substitute this into (4.2.15) above, we find
$$e^{2iq^0}g_{}(t)=𝝀^+(\widehat{k}_{}(i\mathrm{},t))e^{i(q^0+q(t))}k_+(0,t)^1.$$
$`(\mathrm{4.2.17})`$
Consequently, when we eliminate $`𝝀^+(\widehat{k}_{}(i\mathrm{},t))`$ from (4.2.14) and (4.2.17), we obtain the following factorization problem on $`G_\pi ^{}`$:
$$g_{}(t)^1e^{2iq^0}g_+(t)=k_+(0,t)e^{2iq(t)}k_+(0,t)^1.$$
$`(\mathrm{4.2.18})`$
But $`G_\pi ^{}`$ is a reductive Lie group, hence we can find (for at least small values of $`t`$) $`x(t)G_\pi ^{}`$ (unique to transformations $`x(t)x(t)\delta (t)`$ where $`\delta (t)H`$) and unique $`d(t)H`$ such that
$$g_{}(t)^1e^{2iq^0}g_+(t)=x(t)d(t)x(t)^1$$
$`(\mathrm{4.2.19})`$
with $`x(0)=1,d(0)=e^{2iq^0}.`$ This determines $`q(t)`$ via the formula
$$q(t)=\frac{1}{2i}logd(t).$$
$`(\mathrm{4.2.20})`$
On the other hand, let us fix one such $`x(t)`$. We shall seek $`k_+(0,t)`$ in the form
$$k_+(0,t)=x(t)h(t),h(t)H.$$
$`(\mathrm{4.2.21})`$
To determine $`h(t)`$, we shall impose the following condition (which is a corollary of Proposition 4.2.3 (a)):
$$\mathrm{\Pi }_𝔥T_{k_+(0,t)}l_{k_+(0,t)^1}\dot{k}_+(0,t)=0$$
$`(\mathrm{4.2.22})`$
where $`\mathrm{\Pi }_𝔥`$ is the projection map to $`𝔥`$ relative to the direct sum decomposition $`𝔤_\pi ^{}=𝔥+_{\alpha <\pi ^{}>}𝔤_\alpha .`$ Substitute (4.2.21) into (4.2.22), we see that $`h(t)`$ satisfies the equation
$$\dot{h}(t)=T_el_{h(t)}(\mathrm{\Pi }_𝔥T_{x(t)}l_{x(t)^1}\dot{x}(t))$$
$`(\mathrm{4.2.23})`$
with $`h(0)=1.`$ Solving the equation explicitly, we find that
$$h(t)=exp\left\{_0^t\mathrm{\Pi }_𝔥(T_{x(\tau )}l_{x(\tau )^1}\dot{x}(\tau ))𝑑\tau \right\}.$$
$`(\mathrm{4.2.24})`$
###### Theorem 4.2.5
Let $`(q^0,p^0,\xi ^0)J^1(0)=TU\times (𝒰𝔥^{}).`$ Then the Hamiltonian flow on $`J^1(0)`$ generated by
$$\begin{array}{cc}\hfill (q,p,\xi )=& \frac{1}{2}\underset{i}{}p_i^2\frac{1}{2}\underset{\alpha <\pi ^{}>}{}\left(\frac{1}{\mathrm{sin}^2\alpha (q)}\frac{1}{3}\right)\xi _\alpha \xi _\alpha \frac{5}{6}\underset{\alpha \mathrm{\Delta }<\pi ^{}>}{}\xi _\alpha \xi _\alpha \hfill \\ & \frac{1}{3}\underset{i}{}\xi _i^2.\hfill \end{array}$$
with initial condition $`(q(0),p(0),\xi (0))=(q^0,p^0,\xi ^0)`$ is given by
$$\begin{array}{cc}& q(t)=\frac{1}{2i}logd(t),\hfill \\ & \xi (t)=Ad_{k_+(0,t)^1}\xi ^0,\hfill \\ & p(t)=Ad_{k_+(0,t)^1}L(q^0,p^0,\xi ^0)(\pm i\mathrm{})\underset{\alpha <\pi ^{}>}{}(c(\alpha (q(t)))i)\xi _\alpha (t)e_\alpha \hfill \\ & \pm 2i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{\pm }}{}\xi _\alpha (t)e_\alpha \hfill \end{array}$$
$`(\mathrm{4.2.25})`$
where $`d(t)`$ and $`k_+(0,t)`$ are constructed from the above procedure.
###### Demonstration Proof
The formula for $`\xi (t)`$ is a consequence of Corollary 4.2.4 and the relation
$$2i\xi (t)=L(q(t),p(t),\xi (t))(i\mathrm{})L(q(t),p(t),\xi (t))(i\mathrm{}).$$
On the other hand, the formula for $`p(t)`$ follows by equating the two different expressions for $`L(q(t),p(t),\xi (t))(\pm i\mathrm{})`$ in (4.2.3) and in Corollary 4.2.4. $`\mathrm{}`$
By Poisson reduction, we can now write down the solution of the associated integrable model on $`TU\times 𝔤_{red}`$ with Hamiltonian $`_0`$ whose equations of motion are given in Proposition 4.1.3, as in Corollary 3.2.2.
###### Corollary 4.2.6
Let $`(q^0,p^0,s^0)TU\times 𝔤_{red}`$ and suppose $`s^0=Ad_{g(\xi ^0)^1}\xi ^0`$ where $`\xi ^0𝒰𝔥^{}.`$ Then the Hamiltonian flow generated by $`_0`$ with initial condition $`(q(0),p(0),s(0))=(q^0,p^0,s^0)`$ is given by
$$\begin{array}{cc}& q(t)=d(t),\hfill \\ & s(t)=Ad_{\left(\stackrel{~}{k}_+(0,t)g\left(Ad_{\stackrel{~}{k}_+(0,t)^1}s^o\right)\right)^1}s^0,\hfill \\ & p(t)=Ad_{\left(\stackrel{~}{k}_+(0,t)g\left(Ad_{\stackrel{~}{k}_+(0,t)^1}s^o\right)\right)^1}L(q^0,p^0,s^0)(\pm i\mathrm{})\hfill \\ & \underset{\alpha <\pi ^{}>}{}(c(\alpha (q(t)))i)s_\alpha (t)e_\alpha \pm 2i\underset{\alpha \overline{\pi }_{}^{}{}_{}{}^{\pm }}{}s_\alpha (t)e_\alpha \hfill \end{array}$$
$`(\mathrm{4.2.26})`$
where $`\stackrel{~}{k}_+(0,t)=g(\xi ^0)^1k_+(0,t)g(\xi ^0)`$ and $`k_+(0,t)`$, $`d(t)`$ are as in Theorem 4.2.5.
Remark 4.2.7. The reader should contrast the factorization problems in this section with the ones in \[L2\]. Although the Hamiltonians are rather similar (we can transform the hyperbolic spin CM systems in \[L2\] to trigonometric ones), however, the factorization problems involved are quite different.
## 5. The elliptic spin Calogero-Moser systems
## 5.1. Lax operators, Hamiltonian equations and the Lie subalgebroids
In this section, $`\mathrm{}(z)`$ is the Weierstrass $`\mathrm{}`$-function with periods $`2\omega _1`$,$`2\omega _2`$, and $`\sigma (z)`$, $`\zeta (z)`$ are the related Weierstrass sigma-function and zeta-function.
We consider the following elliptic dynamical r-matrix with spectral parameter, given by
$$r(q,z)=\zeta (z)\underset{i}{}x_ix_i\underset{\alpha \mathrm{\Delta }}{}l(\alpha (q),z)e_\alpha e_\alpha $$
$`(\mathrm{5.1.1})`$
where
$$l(w,z)=\frac{\sigma (w+z)}{\sigma (w)\sigma (z)}.$$
$`(\mathrm{5.1.2})`$
Then the associated spin Calogero-Moser system on $`TU\times 𝔤`$ is called the elliptic spin Calogero-Moser system. Explicitly, the Hamiltonian is of the form
$$(q,p,\xi )=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{2}\underset{\alpha \mathrm{\Delta }}{}\mathrm{}(\alpha (q))\xi _\alpha \xi _\alpha $$
$`(\mathrm{5.1.3})`$
and the Lax operator is given by
$$L(q,p,\xi )(z)=p+\zeta (z)\underset{i}{}\xi _ix_i\underset{\alpha \mathrm{\Delta }}{}l(\alpha (q),z)\xi _\alpha e_\alpha .$$
$`(\mathrm{5.1.4})`$
Our next result gives the Hamiltonian equations of motion generated by $`.`$ Using the same method of calculation as in the proof of Proposition 3.1.4, we can also compute the corresponding equations generated by its reduction
$$_0(q,p,s)=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{2}\underset{\alpha \mathrm{\Delta }}{}\mathrm{}(\alpha (q))s_\alpha s_\alpha $$
$`(\mathrm{5.1.5})`$
on $`TU\times 𝔤_{red}.`$
###### Proposition 5.1.1
The Hamiltonian equations of motion generated by $``$ on $`TU\times 𝔤`$ are given by
$$\begin{array}{cc}& \dot{q}=p,\hfill \\ & \dot{p}=\frac{1}{2}\underset{\alpha \mathrm{\Delta }}{}\mathrm{}^{}(\alpha (q))\xi _\alpha \xi _\alpha H_\alpha ,\hfill \\ & \dot{\xi }=[\xi ,\underset{\alpha \mathrm{\Delta }}{}\mathrm{}(\alpha (q))\xi _\alpha e_\alpha ].\hfill \end{array}$$
$`(\mathrm{5.1.6})`$
###### Proposition 5.1.2
The Hamiltonian equations of motion generated by $`_0`$ on the reduced Poisson manifold $`TU\times 𝔤_{red}`$ are given by
$`\dot{q}=p,`$
$`\dot{p}={\displaystyle \frac{1}{2}}{\displaystyle \underset{\alpha \mathrm{\Delta }}{}}\mathrm{}^{}(\alpha (q))s_\alpha s_\alpha H_\alpha ,`$
$`\dot{s}=[s,]`$
where
$$\frac{=\underset{\alpha \mathrm{\Delta }}{}\mathrm{}(\alpha (q))s_\alpha e_\alpha +\underset{i,j}{}C_{ji}}{\alpha \mathrm{\Delta }\alpha _j\alpha \mathrm{\Delta }N_{\alpha ,\alpha _j\alpha }\mathrm{}(\alpha (q))s_\alpha s_{\alpha _j\alpha }h_{\alpha _i}.}$$
(Here the notation $`N_{\alpha ,\beta }`$ is as in Proposition 3.1.4.)
###### Proposition 5.1.3
The classical dynamical r-matrix $`R`$ associated with the elliptic dynamical r-matrix with spectral parameter in (5.1.1) is given by
$$\begin{array}{cc}\hfill (R(q)X)(z)& =\frac{1}{2}X(z)+\underset{k=0}{\overset{\mathrm{}}{}}\frac{\zeta ^{(k)}(z)}{k!}\mathrm{\Pi }_𝔥X_{(k+1)}\hfill \\ & +\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{k!}\underset{\alpha \mathrm{\Delta }}{}\frac{d^k}{dw^k}|_{w=0}l(\alpha (q),zw)(X_{(k+1)})_\alpha e_\alpha .\hfill \end{array}$$
$`(\mathrm{5.1.7})`$
###### Demonstration Proof
Here, we have used the formula $`l(w,z)=l(w,z).`$ Otherwise, the proof is similar to that of Proposition 4.1.4. $`\mathrm{}`$
###### Corollary 5.1.4
On $`J^1(0)`$, we have
$$\begin{array}{cc}& (R(q)M(q,p,\xi ))(z)\hfill \\ \hfill =& \frac{1}{2}M(q,p,\xi )(z)\zeta (z)p+\underset{\alpha \mathrm{\Delta }}{}l(\alpha (q),z)(\zeta (\alpha (q))+\zeta (z)\zeta (\alpha (q)+z))\xi _\alpha e_\alpha \hfill \end{array}$$
$`(\mathrm{5.1.8})`$
where $`M(q,p,\xi )(z)=L(q,p,\xi )(z)/z.`$
###### Demonstration Proof
In a deleted neighborhood of $`0`$, we have the expansion
$$M(q,p,\xi )(z)=\frac{\xi }{z^2}+\frac{1}{z}\left(p+\underset{\alpha }{}\zeta (\alpha (q))\xi _\alpha e_\alpha \right)+O(1).$$
On the other hand, by direct differentiation, we find
$$\frac{d}{dw}|_{w=0}l(\alpha (q),zw)=l(\alpha (q),z)(\zeta (\alpha (q)+z)\zeta (z)).$$
Therefore, on using (5.1.7), we obtain the desired formula. $`\mathrm{}`$
Remark 5.1.5. Using (5.1.8), we can check that in this case, the equations in (5.1.6) can be recovered from the Lax equation $`\dot{L}(q,p,\xi )=[L(q,p,\xi ),R(q)M(q,p,\xi )]`$ on $`J^1(0).`$ The computation makes use of the following identities:
(i) $`l(w,z)l(w,z)=\mathrm{}(z)\mathrm{}(w)`$,
(ii) $`l(x,z)l(y,z)[\zeta (x+z)\zeta (x)\zeta (y+z)+\zeta (y)]=l(x+y,z)[\mathrm{}(x)\mathrm{}(y)]`$,
(iii) $`\zeta (x+y)\zeta (x)\zeta (y)=\frac{1}{2}\frac{\mathrm{}^{}(x)\mathrm{}^{}(y)}{\mathrm{}(x)\mathrm{}(y)}.`$
We shall leave the details to the interested reader.
Our next lemma is a simple consequence of the fact that for $`k0`$, we have $`\zeta ^{(k)}(z)=k!z^{(k+1)}+O(1)`$, $`\frac{d^k}{dw^k}|_{w=0}l(\alpha (q),zw)=k!z^{(k+1)}+O(1)`$ in a deleted neighborhood of $`0`$.
###### Lemma 5.1.6
For $`XL𝔤`$, $`R^+(q)XL^+𝔤`$.
###### Lemma 5.1.7
For $`XL𝔤`$, $`R^{}(q)X+\zeta ()\mathrm{\Pi }_𝔥X_1`$ has singularities at the points of the rank $`2`$ lattice
$$\mathrm{\Lambda }=2\omega _1+2\omega _2$$
$`(\mathrm{5.1.9})`$
and is holomorphic in $`\backslash \mathrm{\Lambda }.`$ Moreover, the quasi-periodicity condition
$$(R^{}(q)X+\zeta ()\mathrm{\Pi }_𝔥X_1)(z+2\omega _i)=Ad_{e^{2\eta _iq}}(R^{}(q)X+\zeta ()\mathrm{\Pi }_𝔥X_1)(z)$$
$`(\mathrm{5.1.10})`$
holds, where $`\eta _i=\zeta (\omega _i),`$ $`i=1,2`$.
###### Demonstration Proof
From (5.1.7), we obtain
$`(R^{}(q)X+\zeta ()\mathrm{\Pi }_𝔥X_1)(z)`$
$`=`$ $`{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\zeta ^{(k)}(z)}{k!}}\mathrm{\Pi }_𝔥X_{(k+1)}+{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\alpha \mathrm{\Delta }}{}}{\displaystyle \frac{1}{k!}}{\displaystyle \frac{d^k}{dw^k}}|_{w=0}l(\alpha (q),zw)(X_{(k+1)})_\alpha e_\alpha `$
from which it is clear that $`R^{}(q)X+\zeta ()\mathrm{\Pi }_𝔥X_1`$ is holomorphic in $`\backslash \mathrm{\Lambda }`$ with singularities at the points of $`\mathrm{\Lambda }`$. On the other hand, it is easy to check that
$$\frac{d^k}{dw^k}|_{w=0}l(\alpha (q),z+2\omega _iw)=e^{2\eta _i\alpha (q)}\frac{d^k}{dw^k}|_{w=0}l(\alpha (q),zw),k0.$$
Hence the second assertion follows. $`\mathrm{}`$
The proof of our next proposition is obvious.
###### Proposition 5.1.8
(a) $`Im^+=_{qU}\{0_q\}\times L^+𝔤\times 𝔥`$. (b) $`^+=_{qU}\{0_q\}\times L^+𝔤\times \{0\}`$ = adjoint bundle of $`Im^+`$.
Indeed, we can show more, namely,
$$\{^+(0_q,X,0)qU,XL𝔤\}=\underset{qU}{}\{0_q\}\times L^+𝔤\times 𝔥.$$
$`(\mathrm{5.1.11})`$
###### Proposition 5.1.9
$`Im^{}=^{}𝒬,`$ where
$$𝒬=\{(0_q,\zeta ()Z,Z)qU,Z𝔥\}$$
$`(\mathrm{5.1.12})`$
is a Lie subalgebroid of $`Im^{}`$ and the ideal $`^{}`$ coincides with the adjoint bundle of $`Im^{}`$ and admits the following characterization:
$$(0_q,X,0)_q^{}\text{if and only if}$$
(a) $`X`$ is holomorphic in $`\backslash \mathrm{\Lambda }`$ with singularities at the points of $`\mathrm{\Lambda }`$,
(b) $`X(z+2\omega _i)=Ad_{e^{2\eta _iq}}X(z),i=1,2`$,
(c) $`\mathrm{\Pi }_𝔥X_1=0.`$
###### Demonstration Proof
As in the proof of Proposition 4.1.10, we have
$`(0_q,X,0)^{}`$
$``$ $`\mathrm{\Pi }_𝔥X_1=0,X(z)=\iota Z(R^{}(q)X+\zeta ()\mathrm{\Pi }_𝔥X_1)(z)\text{for some}Z𝔥.`$
Therefore, it follows from Lemma 5.1.7 that $`X`$ satisfies the properties in (a)-(c). Conversely, suppose $`XL𝔤`$ satisfies the properties in (a)-(c). Consider
$$D(z)=X(z)+(R^{}(q)X)(z).$$
Then $`D(z)`$ satisfies the quasi-periodicity condition $`D(z+2\omega _i)=Ad_{e^{2\eta _iq}}D(z),i=1,2`$ and the principal part of $`D`$ at $`0`$ is zero. Hence $`D`$ extends to a holomorphic map from $``$ to $`𝔤.`$ Write $`D(z)=_jd_j(z)x_j+_{\alpha \mathrm{\Delta }}d_\alpha (z)e_\alpha .`$ Then $`d_j`$ and $`d_\alpha `$ are entire functions for $`1jN,\alpha \mathrm{\Delta }`$ and the quasi-periodicity condition implies that for $`i=1,2`$, we have
$$d_j(z+2\omega _i)=d_j(z),1jN,$$
$`(\mathrm{5.1.13})`$
$$d_\alpha (z+2\omega _i)=d_\alpha (z)e^{2\eta _i\alpha (q)},\alpha \mathrm{\Delta }.$$
$`(\mathrm{5.1.14})`$
From (5.1.13) and Liouville’s theorem, it follows that $`d_j(z)=d_j(=\text{constant})`$ for each $`j`$. On the other hand, observe that the meromorphic function $`l(\alpha (q),z)`$ satisfies the same quasi-periodicity condition as $`d_\alpha `$, that is, $`l(\alpha (q),z+2\omega _i)=l(\alpha (q),z)e^{2\eta _i\alpha (q)}.`$ Hence we conclude from (5.1.14) and the above observation that $`d_\alpha (z)=f_\alpha (z)l(\alpha (q),z)`$, where $`f_\alpha `$ is an elliptic function. But the order of a non-constant elliptic function is never less than $`2`$. Therefore, as $`l(\alpha (q),z)`$ has simple poles at the points of $`\mathrm{\Lambda }`$, we must have $`f_\alpha d_\alpha 0`$ for each $`\alpha \mathrm{\Delta }.`$ Hence we have shown that $`X=\iota ZR^{}(q)X`$, where $`Z=_jd_jx_j.`$ This in turn implies that $`\iota Z+R^+(q)X=0`$ and so $`(0_q,X,0)^{}.`$
Next, we show $`Im^{}=^{}𝒬`$. Consider an arbitrary element $`^{}(0_q,X,Z)`$ in $`Im^{}`$. Clearly, we have the decomposition
$`^{}(0_q,X,Z)`$
$`=`$ $`(0_q,\iota Z+R^{}(q)X+\zeta ()\mathrm{\Pi }_𝔥X_1,0)`$
$`+(0_q,\zeta ()\mathrm{\Pi }_𝔥X_1,\mathrm{\Pi }_𝔥X_1)`$
where the first term is in $`^{}`$ (by Lemma 5.1.7 and the characterization of $`^{}`$ which we established above) and the second term is in $`𝒬`$. This shows that $`Im^{}^{}𝒬.`$ Conversely, take an arbitrary element $`(0_q,X,0)+(0_q,\zeta ()Z,Z)^{}𝒬`$. From the definition of $`^{}`$, we have $`^+(0_q,X,Z^{})=0`$ for some $`Z^{}𝔥`$. Let $`Y=X+\zeta ()Z.`$ Then
$`^{}(0_q,Y,Z^{})`$
$`=`$ $`(0_q,\iota Z^{}R^{}(q)X+R^{}(q)\zeta ()Z,Z)(\mathrm{\Pi }_𝔥Y_1=Z)`$
$`=`$ $`(0_q,\iota Z^{}R^+(q)X+X+R^{}(q)\zeta ()Z,Z)`$
$`=`$ $`(0_q,X,0)+(0_q,\zeta ()Z,Z)(\iota Z^{}R^+(q)X=0andR^{}(q)\zeta ()Z=\zeta ()Z)`$
and this establishes the reverse inclusion $`^{}𝒬Im^{}`$. The assertion that $`^{}`$ coincides with the adjoint bundle of $`Im^{}`$ is now clear. $`\mathrm{}`$
From Proposition 5.1.8, it follows that
$$\begin{array}{cc}\hfill Im^+/^+& =\underset{qU}{}\{0_q\}\times (L^+𝔤/L^+𝔤)\times 𝔥\hfill \\ & \underset{qU}{}\{0_q\}\times \{0\}\times 𝔥\hfill \end{array}$$
$`(\mathrm{5.1.15})`$
where the identification map is given by
$$(0_q,X+L^+𝔤,Z)(0_q,0,Z).$$
$`(\mathrm{5.1.16})`$
Similarly, as a consequence of Proposition 5.1.9, we obtain
$$Im^{}/^{}𝒬.$$
$`(\mathrm{5.1.17})`$
This time, the identification is given by the map
$$(0_q,X,Z)+_q^{}(0_q,\zeta ()Z,Z).$$
$`(\mathrm{5.1.18})`$
The following proposition is now obvious.
###### Proposition 5.1.10
The isomorphism $`\theta :Im^+/^+Im^{}/^{}`$ defined in Proposition 2.1.2 (b) is given by
$$\theta (0_q,0,Z)=(0_q,\zeta ()Z,Z).$$
## 5.2. Solution of the integrable elliptic spin Calogero-Moser systems
We are now ready to discuss the factorization problem
$$exp\{t(0,0,M(q^0,p^0,\xi ^0))\}(q^0)=\gamma _+(t)\gamma _{}(t)^1$$
$`(\mathrm{5.2.1})`$
where $`(\gamma _+(t),\gamma _{}(t))=((q^0,k_+(t),q(t)),(q^0,k_{}(t),q(t)))Im(^+,^{})`$ is to be determined subject to the constraint in (2.2.4) with $`𝔤`$ replaced by $`L𝔤`$ and where $`(q^0,p^0,\xi ^0)J^1(0)=TU\times (𝒰𝔥^{})`$ is the initial value of $`(q,p,\xi )`$. (Recall that $`M(q^0,p^0,\xi ^0)(z)=L(q^0,p^0,\xi ^0)(z)/z.`$) Note that by the global version of Proposition 5.1.9, we have the unique factorization
$$(q^0,k_{}(t),q(t))=(q^0,\widehat{k}_{}(t),q^0)(q^0,e^{\zeta ()(q^0q(t))},q(t))$$
$`(\mathrm{5.2.2})`$
where $`(q^0,\widehat{k}_{}(t),q^0)`$ is in the Lie group bundle integrating $`^{}`$ and the second factor $`(q^0,e^{\zeta ()(q^0q(t))},q(t))`$ is in the Lie groupoid integrating $`𝒬`$. As before, we denote $`k_\pm (t)(z)`$ by $`k_\pm (z,t)`$. Also, denote $`\widehat{k}_{}(t)(z)`$ by $`\widehat{k}_{}(z,t)`$. Then $`k_+(,t)L^+G`$, while $`\widehat{k}_{}(,t)=k_{}(,t)e^{\zeta ()(q(t)q^0)}`$ enjoys the following properties: (a) $`\widehat{k}_{}(,t)`$ is holomorphic in $`\backslash \mathrm{\Lambda }`$ with singularities at the points of $`\mathrm{\Lambda }`$,
(b) $`\widehat{k}_{}(z+2\omega _i,t)=e^{2\eta _iq^0}\widehat{k}_{}(z,t)e^{2\eta _iq^0},i=1,2,`$
(c) $`\left(\frac{d}{dt}|_{t=0}\widehat{k}_{}(z,t)\right)_1𝔥^{}`$. From the factorization problem on the Lie groupoid above and (5.2.2), it follows that
$$e^{tM(q^0,p^0,\xi ^0)(z)}=k_+(z,t)e^{\zeta (z)(q(t)q^0)}\widehat{k}_{}(z,t)^1$$
$`(\mathrm{5.2.3})`$
where $`q(t)`$, $`k_+(,t)`$ and $`\widehat{k}_{}(,t)`$ are to be determined. To do so, let us introduce the following gauge transformations of $`L(q,p,\xi )`$ and $`M(q,p,\xi )`$:
$$\begin{array}{cc}& L^e(q,p,\xi )(z):=Ad_{e^{\zeta (z)q}}L(q,p,\xi )(z),\hfill \\ & M^e(q,p,\xi )(z):=Ad_{e^{\zeta (z)q}}M(q,p,\xi )(z)\hfill \end{array}$$
$`(\mathrm{5.2.4})`$
for $`(q,p,\xi )J^1(0)=TU\times (𝒰𝔥^{})`$. Then the problem in (5.2.3) can be reformulated in the form
$$e^{tM^e(q^0,p^0,\xi ^0)(z)}=k_+^s(z,t)k_{}^e(z,t)^1$$
$`(\mathrm{5.2.5})`$
where
$$k_+^s(z,t)=e^{\zeta (z)q^0}k_+(z,t)e^{\zeta (z)q(t)}$$
$`(\mathrm{5.2.6})`$
and
$$k_{}^e(z,t)=e^{\zeta (z)q^0}\widehat{k}_{}(z,t)e^{\zeta (z)q^0}.$$
$`(\mathrm{5.2.7})`$
Now, on using the property of $`\widehat{k}_{}(z,t)`$ in property (b) above and the fact that $`l(w,z+2\omega _i)=e^{2\eta _iw}l(w,z)`$, it is straight forward to check that
$$k_{}^e(z+2\omega _i,t)=k_{}^e(z,t),$$
$`(\mathrm{5.2.8})`$
and
$$L^e(q,p,\xi )(z+2\omega _i)=L^e(q,p,\xi )(z)$$
$`(\mathrm{5.2.9})`$
for $`i=1,2`$. In view of this, it is natural to introduce the elliptic curve $`\mathrm{\Sigma }=/\mathrm{\Lambda }`$ where $`\mathrm{\Lambda }`$ is the rank 2 lattice in (5.1.9). Thus we can regard $`k_{}^e(,t)`$ as a holomorphic map on $`\mathrm{\Sigma }\backslash \{0\}`$ taking values in $`G`$. On the other hand, the factor $`k_+^s(,t)`$ in (5.2.6) is holomorphic in a deleted neighborhood of $`0\mathrm{\Sigma }`$. Hence we can think of (5.2.5) as a factorization problem on a small circular contour centered at $`0\mathrm{\Sigma }`$ where $`k_+^s(,t)`$ and $`k_{}^e(,t)`$ have analyticity properties as indicated above and satisfying additional constraints. Indeed, it follows from (5.2.6) and property (c) above for $`\widehat{k}_{}(,t)`$ that
$$k_+^s(z,t)e^{\frac{q^0}{z}}k_+(0,t)e^{\frac{q^0}{z}}\text{ in a deleted neighborhood of}0,$$
$`(\mathrm{5.2.10})`$
$$\left(\frac{d}{dt}|_{t=0}\widehat{k}_{}^e(z,t)\right)_1𝔥^{}.$$
$`(\mathrm{5.2.11})`$
Note that if $`(q(t),p(t),\xi (t))`$ is the solution of the Hamiltonian equations in (5.1.6) satisfying the initial condition $`(q(0),p(0),\xi (0))=(q^0,p^0,\xi ^0)`$, then by Theorem 2.2.2 and our discussion above, we have
$$\begin{array}{cc}& L^e(q(t),p(t),\xi (t))(z)\hfill \\ \hfill =& k_+^s(z,t)^1L^e(q^0,p^0,\xi ^0)(z)k_+^s(z,t)\hfill \\ \hfill =& k_{}^e(z,t)^1L^e(q^0,p^0,\xi ^0)(z)k_{}^e(z,t).\hfill \end{array}$$
$`(\mathrm{5.2.12})`$
In the following, we shall write down the solution of the factorization problem explicitly in terms of Riemann theta functions for the case where $`𝔤=sl(N,)`$ with $`𝔥`$ taken to be the Cartan subalgebra consisting of diagonal matrices in $`𝔤.`$ As similar procedures can also be carried out for other classical simple Lie algebras, we shall not give details here for the other cases.
For our purpose, we introduce the spectral curve $`C`$ as defined by the equation
$$det(L(q^0,p^0,\xi ^0)(z)wI)=0.$$
$`(\mathrm{5.2.13})`$
By (5.2.9), this defines an $`N`$-sheeted branched covering $`\pi :C\mathrm{\Sigma }`$ of the elliptic curve $`\mathrm{\Sigma }.`$ Let
$$I(q^0,p^0,\xi ^0;z,w):=det(L(q^0,p^0,\xi ^0)(z)wI)$$
$`(\mathrm{5.2.14})`$
for $`(z,w)^{}\times `$. We shall make the following genericity assumptions: (GA1) zero is a regular value of $`I(q^0,p^0,\xi ^0;,),`$
(GA2) the eigenvalues $`\lambda _1,\mathrm{},\lambda _N`$ of $`\xi ^0`$ are distinct. Then the curve $`C`$ is smooth. The points on $`C`$ corresponding to $`z=0`$ will be considered as points “at $`\mathrm{}`$”, we shall denote them by $`P_1,\mathrm{},P_N`$ respectively. Note that for $`P`$ on the finite part of $`C`$, we have $`dimker(L(z(P))w(P)I)=1`$, for otherwise we would obtain a contradiction to assumption (GA1). Consequently, there exists a unique eigenvector $`\widehat{v}(P)`$ of the matrix $`L^e(q^0,p^0,\xi ^0)(z(P))`$ corresponding to the eigenvalue $`w(P)`$ normalized by the condition $`\widehat{v}_1(P)=(e_1,\widehat{v}(P))=1`$. The next result gives a summary on the properties of the spectral curve and $`\widehat{v}(P)`$ and can be obtained by following the analysis in \[KBBT\].
###### Proposition 5.2.1
Under the genericity assumptions (GA1) and (GA2), the spectral curve $`C`$ has the following properties: (a) $`C`$ is smooth and is an $`N`$-sheeted branched cover of the elliptic curve $`\mathrm{\Sigma },`$
(b) in a deleted neighborhood of $`z=0`$, $`C`$ can be represented as
$$\underset{r=1}{\overset{N}{}}\left(\frac{\lambda _r}{z}+h_r(z)w\right)=0$$
where $`h_1,\mathrm{},h_N`$ are holomorphic in a neighborhood of $`z=0`$,
(c) the genus of $`C`$ is $`g=\frac{1}{2}(N^2N+2).`$
On the other hand, the components $`\widehat{v}_j(P)=(e_j,\widehat{v}(P))`$ of the eigenvector $`\widehat{v}(P)`$, $`j=2,\mathrm{},N`$, are meromorphic on the finite part of $`C`$ with polar divisor $`D=_{i=1}^{g1}\gamma _i.`$ Moreover, in a deleted neighborhood of $`P_k`$, we have
$$\widehat{v}_j(P)=e^{\zeta (z(P))(q_j^0q_1^0)}(\psi _j^{(k)}+O(z(P)))$$
$`(\mathrm{5.2.15})`$
where $`\psi ^{(k)}`$ is the eigenvector of $`\xi ^0`$ corresponding to $`\lambda _k`$ with $`\psi _1^{(k)}=1`$, $`k=1,\mathrm{},N.`$
Now, from the definition of $`\widehat{v}(P)`$ and $`M^e(q^0,p^0,\xi ^0)`$, we have
$$M^e(q^0,p^0,\xi ^0)(z(P))\widehat{v}(P)=(w(P)/z(P))\widehat{v}(P).$$
$`(\mathrm{5.2.16})`$
Hence it follows from (5.2.5) and (5.2.16) that
$$e^{t(w(P)/z(P))}(k_+^s(z(P),t)^1\widehat{v}(P))=k_{}^e(z(P),t)^1\widehat{v}(P)$$
$`(\mathrm{5.2.17})`$
for $`z(P)`$ in a deleted neighborhood of $`0\mathrm{\Sigma }.`$ Set
$$v_+(t,P)=k_+^s(z(P),t)^1\widehat{v}(P),$$
$`(\mathrm{5.2.18})`$
$$v_{}(t,P)=k_{}^e(z(P),t)^1\widehat{v}(P).$$
$`(\mathrm{5.2.19})`$
Then from (5.2.17), (5.2.12) and the definition of $`\widehat{v}(P)`$, we obtain
$$e^{t(w(P)/z(P))}v_+(t,P)=v_{}(t,P),$$
$`(\mathrm{5.2.20})`$
$$L^e(q(t),p(t),\xi (t))(z(P))v_\pm (t,P)=w(P)v_\pm (t,P).$$
$`(\mathrm{5.2.21})`$
In this way, we are led to scalar factorization problems for the components of a suitably normalized eigenvector of $`L^e(q(t),p(t),\xi (t))(z(P))`$.
###### Proposition 5.2.2
In a deleted neighborhood of $`P_k`$, $`k=1,\mathrm{},N`$, we have
$`v_+^j(t,P)`$ $`=e^{\zeta (z(P))(q_1^0q_j(t))}((k_+(0,t)^1\psi ^{(k)})_j+O(z(P)))`$
$`e^{(q_1^0q_j(t))z(P)^1}(k_+(0,t)^1\psi ^{(k)})_jasPP_k,`$
$`v_{}^j(t,P)`$ $`=e^{tw(P)z(P)^1+\zeta (z(P))(q_1^0q_j(t))}((k_+(0,t)^1\psi ^{(k)})_j+O(z(P)))`$
$`e^{t(\lambda _k/z(P)+h_k(0))z(P)^1+(q_1^0q_j(t))z(P)^1}(k_+(0,t)^1\psi ^{(k)})_jasPP_k.`$
###### Demonstration Proof
This is a consequence of (5.2.18)-(5.2.20),(5.2.6),(5.2.15) and the fact that in a deleted neighborhood of $`P_k`$, we have $`w(P)=\lambda _kz(P)^1+h_k(0)+O(z(P))`$ from Proposition 5.2.1 (b). $`\mathrm{}`$
In order to write down $`v_{}^j(t,P)`$, we will insert a fictitious pole together with a matching zero to this function at some point $`\gamma _0`$ on the finite part of $`C`$ distinct from $`\gamma _1,\mathrm{},\gamma _{g1}`$. By putting in an additional pole in this way, we would be able to construct $`v_{}^j(t,P)`$ as a multi-point Baker-Akheizer function. To do so, let us fix a canonical homology basis $`\{a_j,b_k\}_{1j,kg}`$ of the Riemann surface associated with $`C`$ and let $`\{\omega _i\}_{1ig}`$ be a cohomology basis dual to $`\{a_j,b_k\}_{1j,kg}`$, i.e. $`_{a_j}\omega _i=\delta _{ij}`$, $`_{b_j}\omega _i=\mathrm{\Omega }_{ij}`$. With respect to the Riemann matrix $`\mathrm{\Omega }=(\mathrm{\Omega }_{ij})`$, we construct the theta function
$$\theta (z_1,\mathrm{},z_g)=\underset{m^g}{}exp\{2\pi i(m,z)+\pi i(\mathrm{\Omega }m,m)\}.$$
$`(\mathrm{5.2.22})`$
We also introduce the Abel-Jacobi map
$$A:CJac(C),P(_{P_0}^P\omega _1,\mathrm{},_{P_0}^P\omega _g)$$
$`(\mathrm{5.2.23})`$
where $`P_0`$ is some fixed point on the finite part of $`C`$. Now, let $`d\mathrm{\Omega }^{(i)}`$, $`i=1,2`$, be the unique abelian differential of second kind with vanishing $`a`$-periods such that in a deleted neighborhood of $`P_k`$,
$$d\mathrm{\Omega }^{(1)}=d(z^1+\omega ^{(1)}(z)),$$
$`(\mathrm{5.2.24})`$
$$d\mathrm{\Omega }^{(2)}=d(\lambda _kz^2+h_k(0)z^1+\omega ^{(2)}(z)),$$
$`(\mathrm{5.2.25})`$
where $`\omega ^{(1)}(z),\omega ^{(2)}(z)`$ are regular at $`z=0`$, $`k=1,\mathrm{},N.`$ We shall denote by $`2\pi iU^{(i)}`$ the vector of $`b`$-periods of $`d\mathrm{\Omega }^{(i)}`$, $`i=1,2.`$ Then from Proposition 5.2.2 and the fact that $`v_{}^j(t,P)`$ has $`_{i=0}^{g1}\gamma _i`$ as a polar divisor, we obtain the following result from the standard construction of Baker-Akhiezer functions \[K\].
###### Proposition 5.2.3
For $`1jN`$,
$`v_{}^j(t,P)`$
$`=`$ $`f_j(t){\displaystyle \frac{\theta (A(P)+(q_1^0q_j(t))U^{(1)}+tU^{(2)}A(D)A(\gamma _0)K)}{\theta (A(P)A(D)A(\gamma _0)K)}}`$
$`\times exp[(q_1^0q_j(t))\mathrm{\Omega }^{(1)}(P)+t\mathrm{\Omega }^{(2)}(P)]`$
where
$$\mathrm{\Omega }^{(i)}(P)=_{P_0}^P𝑑\mathrm{\Omega }^{(i)},i=1,2$$
and $`K`$ is the vector of Riemann constants.
###### Corollary 5.2.4
$`\theta \left(q_j(t)U^{(1)}tU^{(2)}+V\right)=0`$ where $`V=A(D)q_1^0U^{(1)}+K,j=1,\mathrm{},N.`$
###### Demonstration Proof
This follows when we evaluate the expression for $`v_{}^j(t,P)`$ in Proposition 5.2.3 at the point $`\gamma _0`$ and equate the result to zero. $`\mathrm{}`$
Let
$$f(t)=diag(f_1(t),\mathrm{},f_N(t)).$$
$`(\mathrm{5.2.26})`$
In view of Proposition 5.2.3, we shall write
$$v_{}(t,P)=f(t)v_{}^\theta (t,P)$$
$`(\mathrm{5.2.27})`$
where the $`v_{}^\theta (t,P)`$ are known. Note that if we set $`t=0`$ in the above expression, we obtain $`\widehat{v}(P)=f(0)v_{}^\theta (0,P).`$ Clearly, $`f_1(0)=1`$; the other $`f_j(0)`$’s are then uniquely determined from the definition of $`\widehat{v}(P)`$. Now, for given $`z\mathrm{\Sigma }`$ which is not a branch point of the coordinate function $`z(P)`$, there exist $`N`$ points $`P_1(z),\mathrm{},P_N(z)`$ of $`C`$ lying over $`z`$. Hence we can define the matrices
$$\begin{array}{cc}\hfill \widehat{V}(z)=(\widehat{v}(P_1(z)),& \mathrm{},\widehat{v}(P_N(z))),V_{}(z,t)=(v_{}(t,P_1(z)),\mathrm{},v_{}(t,P_N)),\hfill \\ & V_{}^\theta (z,t)=(v_{}^\theta (t,P_1(z)),\mathrm{},v_{}^\theta (t,P_N(z))).\hfill \end{array}$$
$`(\mathrm{5.2.28})`$
With these definitions, if follows from (5.2.27) and (5.2.19) that
$$k_{}^e(z,t)=\widehat{V}(z)V_{}^\theta (z,t)^1f(t)^1$$
$`(\mathrm{5.2.29})`$
where $`f(t)`$ is still to be determined. To do this, we invoke the condition that $`T_{\gamma _+(t)}𝚕_{\gamma _+(t)^1}\dot{\gamma }_+(t)^+(\{q(t)\}\times \{0\}\times L𝔤).`$ Indeed, it follows from the proof of Theorem 2.2.2 in \[L2\] that
$$T_{\gamma _+(t)}𝚕_{\gamma _+(t)^1}\dot{\gamma }_+(t)=^+(q(t),0,M(q(t),p(t),\xi (t))).$$
$`(\mathrm{5.2.30})`$
Consequently, we have
$$\begin{array}{cc}& k_+(z,t)^1\dot{k}_+(z,t)\hfill \\ \hfill =& (R^+(q(t))M(q(t),p(t),\xi (t)))(z)\hfill \\ \hfill =& M(q(t),p(t),\xi (t))(z)\zeta (z)p(t)+\underset{ij}{}l(q_i(t)q_j(t),z)(\zeta (q_i(t)q_j(t))\hfill \\ & +\zeta (z)\zeta (q_i(t)q_j(t)+z))\xi _{ij}(t)e_{ij}\hfill \end{array}$$
$`(\mathrm{5.2.31})`$
where in the last step we have used (5.1.8). Therefore, when we expand the above expression about $`z=0`$ and compare the term in $`z^0`$, we find that
$$k_+(0,t)^1\dot{k}_+(0,t)=\underset{ij}{}\zeta ^{}(q_i(t)q_j(t))\xi _{ij}(t)e_{ij}$$
$`(\mathrm{5.2.32})`$
from which we obtain the condition
$$\mathrm{\Pi }_𝔥((k_+(0,t)^1\dot{k}_+(0,t))=0.$$
$`(\mathrm{5.2.33})`$
In order to state our next result, we introduce
$$\begin{array}{cc}& \omega _k^{(1)}=\underset{PP_k}{lim}\left(\mathrm{\Omega }^{(1)}(P)z(P)^1\right)\hfill \\ & \omega _k^{(2)}=\underset{PP_k}{lim}\left(\mathrm{\Omega }^{(2)}(P)(\lambda _kz(P)^2+h_k(0)z(P)^1)\right)\hfill \end{array}$$
$`(\mathrm{5.2.34})`$
for $`k=1,\mathrm{},N`$. We also introduce the matrix $`W^\theta (t)=(W_{jk}^\theta (t))`$ where
$$\begin{array}{cc}\hfill W_{jk}^\theta (t)& =\frac{\theta (A(P_k)+(q_1^0q_j(t))U^{(1)}+tU^{(2)}A(D)A(\gamma _0)K)}{\theta (A(P_k)A(D)A(\gamma _0)K)}\hfill \\ & \times exp\left[(q_1^0q_j(t))\omega _k^{(1)}+t\omega _k^{(2)}\right].\hfill \end{array}$$
$`(\mathrm{5.2.35})`$
###### Proposition 5.2.5
$`f(t)`$ satisfies the differential equation
$$\dot{f}(t)=f(t)\mathrm{\Pi }_𝔥\left(\dot{W}^\theta (t)W^\theta (t)^1\right)$$
and hence
$$f(t)=f(0)exp\left\{_0^t\mathrm{\Pi }_𝔥\left(\dot{W}^\theta (\tau )W^\theta (\tau )^1\right)𝑑\tau \right\}.$$
$`(\mathrm{5.2.36})`$
###### Demonstration Proof
From Proposition 5.2.2 and 5.2.3, we have
$`(k_+(0,t)^1\psi ^{(k)})_j`$
$`=`$ $`\underset{PP_k}{lim}v_{}^j(t,P)exp[t(\lambda _kz(P)^2+h_k(0)z(P)^1)(q_1^0q_j(t))z(P)^1]`$
$`=`$ $`f_j(t)W_{jk}^\theta (t)`$
which implies
$$k_+(0,t)=\mathrm{\Psi }W^\theta (t)^1f(t)^1$$
where $`\mathrm{\Psi }`$ is the $`N\times N`$ matrix whose $`k`$-th column is the vector $`\psi ^{(k)}`$, $`k=1,\mathrm{},N.`$ Differentiating the above expression with respect to $`t`$, we find
$$k_+(0,t)^1\dot{k}_+(0,t)=f(t)\dot{W}^\theta W^\theta (t)^1f(t)^1\dot{f}(t)f(t)^1.$$
Therefore, when we apply the condition in (5.2.33), the desired equation for $`f(t)`$ follows. Finally, the solution of the equation is obvious. $`\mathrm{}`$
Hence $`k_+(0,t)=\mathrm{\Psi }W^\theta (t)^1f(t)^1`$ and $`k_{}^e(z,t)=\widehat{V}(z)V_{}^\theta (z,t)^1f(t)^1`$ are completely determined. Therefore we have following result.
###### Theorem 5.2.6
Let $`(q^0,p^0,\xi ^0)J^1(0)=TU\times (𝒰𝔥^{})`$ satisfy the genericity assumptions (GA1), (GA2). Then the Hamiltonian flow on $`J^1(0)`$ generated by
$$(q,p,\xi )=\frac{1}{2}\underset{i}{}p_i^2\frac{1}{2}\underset{ij}{}\mathrm{}(q_iq_j)\xi _{ij}\xi _{ji}$$
with initial condition $`(q(0),p(0),\xi (0))=(q^0,p^0,\xi ^0)`$ is given by
$$\begin{array}{cc}& \theta \left(q_j(t)U^{(1)}tU^{(2)}+V\right)=0,j=1,\mathrm{},N,\hfill \\ & \xi (t)=k_+(0,t)^1\xi ^0k_+(0,t),\hfill \\ & p(t)=Ad_{k_{}^e(z,t)^1}L^e(q^0,p^0,\xi ^0)(z)\hfill \\ & +\underset{ij}{}l(q_i(t)q_j(t),z)e^{\zeta (z)(q_i(t)q_j(t))}\xi _{ij}(t)e_{ij}\hfill \end{array}$$
$`(\mathrm{5.2.37})`$
where $`k_+(0,t)`$, $`k_{}^e(z,t)`$ are given by the formulas above.
Finally we are ready to give the solutions of the associated integrable model on $`TU\times 𝔤_{red}`$ whose equations are given in Proposition 5.1.2.
###### Corollary 5.2.7
Let $`(q^0,p^0,s^0)TU\times 𝔤_{red}`$ and suppose $`s^0=Ad_{g(\xi ^0)^1}\xi ^0`$ where $`\xi ^0𝒰𝔥^{}.`$ Then the Hamiltonian flow generated by $`_0`$ with initial condition $`(q(0),p(0),s(0))=(q^0,p^0,s^0)`$ is given by
$$\begin{array}{cc}& \theta \left(q_j(t)U^{(1)}tU^{(2)}+V\right)=0,j=1,\mathrm{},N,\hfill \\ & s(t)=Ad_{\left(\stackrel{~}{k}_+(0,t)g\left(Ad_{\stackrel{~}{k}_+(0,t)^1}s^o\right)\right)^1}s^0,\hfill \\ & p(t)=Ad_{\left(\stackrel{~}{k}_+^e(z,t)g\left(Ad_{\stackrel{~}{k}_+(0,t)^1}s^0\right)\right)^1}L^e(q^0,p^0,s^0)(z)\hfill \\ & +\underset{ij}{}l(q_i(t)q_j(t),z)e^{\zeta (z)(q_i(t)q_j(t))}s_{ij}(t)e_{ij}\hfill \end{array}$$
$`(\mathrm{5.2.38})`$
where $`k_+(0,t)`$, $`k_{}^e(z,t)`$ are given by the formulas above.
Remark 5.2.8. (a) In \[KBBT\], the authors considered the $`gl(N,)`$-elliptic spin Calogero-Moser system with Hamiltonian
$$H(q,p,f)=\frac{1}{2}\underset{i=1}{\overset{N}{}}p_i^2+\frac{1}{2}\underset{ij}{}\mathrm{}(q_iq_j)\xi _{ij}\xi _{ji}$$
$`(\mathrm{5.2.39})`$
and they imposed the following restriction on $`\xi =(\xi _{ij})gl(N,)gl(N,)^{}`$, namely, they set
$$\xi _{ij}=b_i^Ta_j$$
$`(\mathrm{5.2.40})`$
for all $`i`$ and $`j`$ where $`a_j`$, $`b_j`$ are (column) vectors in $`^l,`$ $`l<N,`$ satisfying the nontrivial Poisson bracket relations $`\{a_{i,\alpha },b_{j,\beta }\}=\delta _{i,j}\delta _{\alpha ,\beta }.`$ Thus from the outset, it is clear that these authors were restricting themselves to special coadjoint orbits of $`gl(N,)^{}gl(N,)`$ which consist of matrices of the form $`B^TA`$, where
$$A=(a_1,\mathrm{},a_N),B=(b_1,\mathrm{},b_N)$$
$`(\mathrm{5.2.41})`$
are $`l\times N`$ matrices. However, it is only through the imposition of (5.2.40) that they were able to make the connection with the matrix KP equation. The precise relation is that the equations of motion for $`a_j`$, $`b_j`$ (up to gauge equivalence) and $`q_j`$ are the necessary and sufficient condition for the time-dependent matrix Schrödinger equation
$$\left(_t_x^2+\underset{j=1}{\overset{N}{}}a_j(t)b_j^T(t)\mathrm{}(xq_j(t))\right)\mathrm{\Psi }=0$$
$`(\mathrm{5.2.42})`$
and its adjoint
$$\stackrel{~}{\mathrm{\Psi }}^T\left(_t_x^2+\underset{j=1}{\overset{N}{}}a_j(t)b_j^T(t)\mathrm{}(xq_j(t))\right)=0$$
$`(\mathrm{5.2.43})`$
($`\stackrel{~}{\mathrm{\Psi }}^T\stackrel{~}{\mathrm{\Psi }}^T`$) to admit solutions of the form
$$\mathrm{\Psi }=\underset{j=1}{\overset{N}{}}s_j(t,k,z)\mathrm{\Phi }(xq_j(t),z)e^{kx+k^2t},$$
$`(\mathrm{5.2.44})`$
$$\stackrel{~}{\mathrm{\Psi }}=\underset{j=1}{\overset{N}{}}s_j^+(t,k,z)\mathrm{\Phi }(x+q_j(t),z)e^{kxk^2t},$$
$`(\mathrm{5.2.45})`$
where $`s_j`$,$`s_j^+`$ are functions which take values in $`^l`$ and
$$\mathrm{\Phi }(x,z)=\frac{\sigma (zx)}{\sigma (x)\sigma (z)}e^{\zeta (z)x}.$$
$`(\mathrm{5.2.46})`$
Furthermore, it is in the course of proving this result that the Lax pair as well as the constraint $`f_{jj}=b_j^Ta_j=2,j=1,\mathrm{},N,`$ emerge naturally. In short, the method of solution in \[KBBT\] of their elliptic spin Calogero-Moser system as defined in (5.2.39), (5.2.40) on the constraint manifold $`\{f_{jj}=b_j^Ta_j=2,j=1,\mathrm{},N\}`$ is based on the above correspondence.
In what follows, we shall give a sketch of this method to solve the equations of motion for $`a_j(t),b_j(t)`$ and $`q_j(t)`$ so that the reader can understand its limitations. To start with, one has a normalized Baker-Akhiezer vector function $`\psi (x,t,P)`$ (and its dual $`\psi ^+(x,t,P)`$) which is uniquely determined by the spectral curve, a divisor of degree $`g+l1`$, and prescribed behaviour in deleted neighborhoods of the punctures $`P_j`$, $`j=1,\mathrm{},l.`$ (These are in turn fixed by the initial data.) From a general result in \[K\], corresponding to $`\psi (x,t,P)`$ and $`\psi ^+(x,t,P)`$ is an algebro-geometric potential $`u(x,t)`$ satisfying
$$\left(_t_x^2+u(x,t)\right)\psi (x,t,P)=0$$
$`(\mathrm{5.2.47})`$
and
$$(\psi ^+(x,t,P))^T\left(_t_x^2+u(x,t)\right)=0.$$
$`(\mathrm{5.2.48})`$
Next, by analyzing $`\psi (x,t,P)`$ as a function of $`x`$, and on comparing $`\mathrm{\Psi }(x,t,P)`$ with $`\psi (x,t,P)`$, one can show that there exists a constant invertible matrix $`\chi _0`$ such that
$$\chi _0u(x,t)\chi _0^1=\underset{j=1}{\overset{N}{}}a_j(t)b_j^T(t)\mathrm{}(xq_j(t)).$$
$`(\mathrm{5.2.49})`$
Since $`u(x,t)`$ is determined by $`\psi (x,t,P)`$, it is in this fashion that the authors in \[KBBT\] were able to write down $`a_j(t),b_j(t)`$ and the equation satisfied by $`q_j(t)`$ in terms of theta functions. Now, let us examine the expression $`_{j=1}^Na_j(t)b_j^T(t)\mathrm{}(xq_j(t))`$ for the potential carefully. Clearly, it depends on the $`a_j`$’s and $`b_j`$’s through the rank one matrices $`a_jb_j^T`$ rather than on the entries of the matrix $`\xi =(b_j^Ta_j).`$ For this reason, the solution method sketched above is rather specific to these special coadjoint orbits of $`gl(N,)^{}.`$ Obviously, similar remarks also hold for the corresponding rational and trigonometric cases. Thus it is clear that this method is not applicable to our more general class of spin Calogero-Moser systems associated with simple Lie algebras here. (b) The first link between elliptic solutions of integrable PDEs and discrete particle systems was found in the paper of Airault, McKean and Moser \[AMM\]. Indeed, the PDE in \[AMM\] is KdV and the corresponding discrete particle system is the usual (spinless) Calogero-Moser system. In the context of the spin CM system defined by (5.2.39) and (5.2.40) above, we remark that its correspondence with matrix KP is related to some interesting algebraic geometry and we refer the reader to \[T\] for details. For our general class of spin CM systems which we address in this work, whether it has any connection with integrable PDEs is entirely open at this point.
Acknowledgments. The factorization method in \[L2\] and the present work has been the subject of several lectures in the last few years. The author would like to express his appreciation to his colleagues at various institutions for their interest and for their hospitality during his visits. He would also like to thank Alan Weinstein for an inspiring lecture at MSRI in 1989 which set him to learn about groupoids. Special thanks is due to Armando Treibich for an interesting explanation on tangential covers at IHP in the summer of 2004. Finally, he acknowledges the helpful advice of Percy Deift and an anonymous referee on the presentation of the material.
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# Competition of ferromagnetism and superconductivity in Sc3InB
## I Introduction
Recent discovery of superconductivity in the intermetallic perovskite MgCNi<sub>3</sub> nature1 with T$`{}_{c}{}^{}8`$ K was a great surprise due to large Ni contents and this compound was rather expected to be near ferromagnetic critical point pick1 . It was established that electron–phonon mechanism is responsible for superconductivity in this material, but few details are still not clear. NMR $`t_1`$ relaxation time experiments nmr or specific heat measurements spheat have resulted in typical behaviors, supporting s–wave type pairing with electron–phonon coupling constant $`\lambda 0.8`$. Conversely, rather unconventional behaviors have been observed in other experiments (e.g. increase of critical temperature with pressure press1 or unusual low–temperature behavior of London penetration depth $`\lambda _L(T)`$ penetr ), which can be partly connected with spin fluctuations likely appearing due to vicinity of ferromagnetism. Moreover, complex dynamical properties of this superconductor, as e.g. soft–mode behaviors and instability of Ni vibrations sav1 ph , make its theoretical analysis quite cumbersome.
In this paper we report on predictions of superconductivity in related Sc<sub>3</sub>InB compound, which also seems to be close to magnetism limit due to weak ferromagnetic properties of Sc<sub>3</sub>In. It was already revealed that Sc<sub>3</sub>InB crystallizes in a perovskite structure (space group $`P`$m-3m, CaTiO<sub>3</sub> type) with lattice constant $`a=4.56`$ Å cryst . However, the same number of valence electrons of boron and indium elements may give rise to lattice instabilities and In/B antisite defects are plausible.
## II Theoretical study
Electronic structure calculations were performed by Full Potential KKR method kkr , kkr2 within the LDA framework employing von Barth-Hedin formula for the exchange-correlation potential. Figure 1 presents electron density of states (DOS) in Sc<sub>3</sub>InB compound. The Fermi level (E<sub>F</sub>) is located on the decreasing slope of large DOS peak with the $`n(E_F)`$ value being as large as $`90`$ states/Ry. The main contributions can be attributed to Sc ($`d`$–states) and In ($`p`$–states) (see Tab. 1). Noteworthy, appearance of a large DOS peak close to $`E_F`$, coming essentially from transition metal Ni–$`d`$ states, was also characteristic of electronic structure of MgCNi<sub>3</sub> (Mg contribution was negligible) pick1 ,szajek . The Sc<sub>3</sub>InB case is not similar, since In plays more active role in formation of electronic states near $`E_F`$ due to one electron more on $`p`$–shell.
### Predictions of superconductivity in Sc<sub>3</sub>InB
The electron–phonon coupling strength was estimated by calculating McMillan–Hopfield parameters mcm from Gaspari–Györffy formulas gasgy – pick2 within Rigid Muffin Tin Approximation (RMTA). The electron–phonon coupling constant was then calculated from relation: $`\lambda =_i\eta _i/m_i\omega _i^2`$, where $`i`$ corresponds to atoms in the unit cell. The $`\omega _i^2`$ parameter was derived from phonon DOS $`F(\omega )`$ <sup>1</sup><sup>1</sup>1definition: $`\omega ^n=\omega ^{n1}\alpha ^2F(\omega )𝑑\omega /\omega ^1\alpha ^2F(\omega )𝑑\omega \omega ^{n1}F(\omega )𝑑\omega /\omega ^1F(\omega )𝑑\omega `$, if $`\alpha ^2(\omega )const.`$, where $`\alpha ^2(\omega )`$ is the electron–phonon interaction coefficient., computed for minimum–energy lattice constant $`a_0=8.610`$ $`a_B`$ (1 $`a_B`$ = 0.529 Å) within Density Functional Perturbation Theory, using the PWscf package pwscf .
In order to take into account the markedly different masses of Sc, In and B atoms, phonon spectrum was analyzed while focusing on two parts: low frequency region with predominantly In and Sc modes (peak at 13 meV comes from flat acoustic In modes) and high frequency region with B vibrations (above 50 meV). Noteworthy, similar separation of phonon DOS was earlier observed experimentally in MgCNi<sub>3</sub> ph . Values of $`\omega _i^2`$ used next to estimate $`\lambda `$, were calculated separately for both regions (with cut line at 50 meV): $`\omega _i^2`$ for Sc and In are taken to be equal, and represent the low–frequency part, while $`\omega _B^2`$ was taken from high frequency part of the spectrum. Results and estimation of $`\lambda _i`$ are presented in Table 1, total electron–phonon coupling constant is $`\lambda =0.94`$.
Using McMillan formula mcm for critical temperature and applying typical value of Coulomb pseudopotential $`\mu ^{}=0.13`$, we got $`T_c12`$ K (we used $`\omega /1.2`$ with $`\omega =19.5`$ meV instead of $`\omega _D/1.45`$ in McMillan formula in practical computations). Both $`\lambda `$ and $`T_c`$ values belong to moderate regime of superconducting parameters within the RMTA model, and are higher than in MgCNi<sub>3</sub>.
We can also notice advantageous trend of electron–phonon coupling in Sc<sub>3</sub>InB if comparing with MgCNi<sub>3</sub> superconductor. Our KKR calculations for MgCNi<sub>3</sub> showed that Ni has the largest McMillan–Hopfield parameter: $`\eta _{Ni}=20`$ mRy/$`a_B^2`$ with respect to other atoms ($`\eta _C=9`$ mRy/$`a_B^2`$, $`\eta _{Mg}`$ \- negligible). The value of $`\eta _{Sc}`$ in Sc<sub>3</sub>InB is similar to $`\eta _{Ni}`$ in MgCNi<sub>3</sub> but $`\eta _B`$ is over three times larger than $`\eta _C`$ (both In and Mg contributions are much smaller). In view of recent C isotope effect measurements isef (very large $`\alpha _C=0.54`$ coefficient) and bearing in mind a particular sensitivity of $`T_c`$ on the carbon concentration, one can conclude that both C and Ni sublattices are important in superconductivity of MgCNi<sub>3</sub>. In view of these criterion $`T_c`$ in Sc<sub>3</sub>InB can be expected higher than in MgCNi<sub>3</sub>.
### Magnetic properties of Sc<sub>3</sub>In
It seems that instability towards magnetism may also be present in Sc<sub>3</sub>InB<sub>1-x</sub> compound. First, we have studied electronic structure of both allotropic phases of Sc<sub>3</sub>In. The hexagonal compound (Ni<sub>3</sub>Sn-type, $`a`$ = 6.42 Å, $`c`$ = 5.18 Å cryst ) is well–known weak itinerant ferromagnet, while the cubic compound (Cu<sub>3</sub>Au–type, $`a=4.46`$ Å, synthesized under high pressure cryst ) has not been yet investigated to our knowledge. Present KKR calculations showed that both allotropic forms of Sc<sub>3</sub>In should exhibit magnetic ground state supported by magnetic moment on Sc atoms, i.e. 0.26 $`\mu _B`$ (hexagonal phase) and 0.27 $`\mu _B`$ (cubic phase). We should remind that experimentally observed magnetic moment is much weaker ($`0.05\mu _B`$/Sc in hexagonal phase), as already underlined in the previous LAPW calculations singh . Besides, the superconductor–to–ferromagnet transition can appear in Sc<sub>3</sub>InB<sub>1-x</sub> if varying B content (KKR-CPA computations are in progress). Thus, as already suggested in MgCNi<sub>3</sub> pick1 , sav1 one can expect that the proximity of ferromagnetic quantum critical point (resulting in enhanced spin fluctuations) may probably compete with superconductivity.
## III Experimental analyses
$``$1.5 g sample was first prepared by arc melting under high purity argon atmosphere (99.9995) starting from the appropriate proportions of the elements (purity $`>`$ 99.95) to obtain the Sc<sub>3</sub>InB formula. The resulting small ingot was melted several times in order to insure homogeneity. Pieces of the ingots were made by using a metal mortar, and inspection by using an optical microscope reveals the bright and homogeneous aspect of the fractured surfaces. Then, XRD patterns were recorded at $`\lambda _{K\alpha }(Cu)`$ using a Bragg-Brentano diffractometer equipped with a backscattering pyrolitic graphite monochromator. The diffraction pattern reveals the presence of a dominant amount of cubic phase with the addition of a minor impurity. Probably because no annealing procedure was applied, the crystallised state of the sample was not of the best and no effective crystal structure refinement was applied. However using the PowderCell code pcell , the two main phases were clearly identified and a rough determination of the cell parameters was made possible. The main phase ($``$70 % vol.) is simple cubic of perovskite type of structure and the second one ($``$30 % vol.) is hexagonal of Ni<sub>2</sub>In type of structure (space group $`P6_3/mmc`$). Accounting for the large difference in between the scattering lengths of the p-elements B and In, an estimate of the composition of the main phase was made. All the results are displayed on Table 2.
The best agreement for the main phase composition leads to consider that the p–elements B and In are neither fully ordered nor randomly distributed, with about 0.70 In and 0.30 B atom on the 1a site, then 0.65 In and 0.35 B atom on the 1b site of the space group $`P`$m-3m. Consideration for such a type of atomic disordering was already reported in literature d2 . Both the compositions of the two compounds (3-1-1 and 2-1) confirm that boron is a difficult element to combine, thus the remaining boron should be detected in the XRD pattern, unless it is difficult to evidence effectively as a very light and often poorly crystallised element.
Susceptibility and resistivity measurements were made using a a.c. susceptometer in temperature ranging from 50 to 2 K. Millimeter sized pieces of the ingot were measured successively and lead to the same results. A typical record is shown on Figure 3, thus revealing the onset of a superconducting state down to 4.4 K, the transition being very sharp with no detectable hysteresis loop. A change in the resistivity trace was also observed simultaneously.
Unexpectedly, the rest of the ingot and the small pieces rapidly change their brittle and bright aspects after several hours let in ambient atmosphere. Besides, a tentative to melt twice the main parts of the initial ingot was made, but at this time no evidence for any superconducting transition was found. Then, several new syntheses were undertaken, thus operating as possible similarly as for the first one procedure. Again, the new ingots we obtained do not display any transition down to 2 K. As revealed by the theoretical derivations (see below), the superconducting state and the related transition look fairly dependant on the p–element ordering and the stoichiometry in the 1a and 1b sites. So, we anticipate that during and after the different melts, the new samples do not exhibit the same atom ordering as it was resulting from the first attempt. New syntheses are now scheduled to be undertaken using different methods and techniques in order to achieve optimized compounds.
Because of experimental problems with synthesis, effect of boron site vacancy on electronic structure near E<sub>F</sub> and superconducting properties was simulated using KKR method with coherent potential approximation (CPA). We have found significant change in $`\eta `$ values not only for B, but also for Sc, i.e. for 7% boron deficiency (Sc<sub>3</sub>InB<sub>0.93</sub>), employing the same values of $`\omega _i^2`$, total $`\lambda `$ decreased over 30% to 0.62, which resulted in $`T_c4`$ K. In Sc<sub>3</sub>InB<sub>0.85</sub> coupling constant is so small ($`\lambda 0.4`$), $`T_c<0.5`$ K seems to be below the standard low–temperature measurements. Similar decrease (but less rapid than in the case of B vacancy) of critical parameters was detected from KKR-CPA analysis when antisite In/B disorder increased.
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# Nearsightedness of Electronic Matter in One Dimension
## I Introduction
This paper is based on a preliminary remark by W. Kohn,Kohn96 about a general concept called “nearsightedness of electronic matter (NEM)” and on a recent short report (PK),ProdanKohn which amplified that remark in various aspects but did not include detailed proofs. In the present paper, we select that part of PK dealing with non-interacting 1D electrons and provide a full discussion, including detailed proofs. Future publications will amplify other sections of PK.
By “electronic matter” we understand a system of many electrons with significant wavefunction overlap, in equilibrium under the action of a given external potential $`v(r)`$. We shall consider the change of a local electronic property, like the electron density $`n(r)`$, under the action of an arbitrarily strong potential perturbation $`w(r^{})`$, which vanishes inside a specified sphere, $`|rr^{}|=\text{R}`$. Note that we allow situations when the perturbation completely surrounds the point $`r`$. NEM states that the resulting density change at $`r`$, $`\mathrm{\Delta }n(r,\text{R})`$, is bounded by a function $`\overline{\mathrm{\Delta }n}(r,\text{R})`$,
$$\mathrm{\Delta }n(r,\text{R})\overline{\mathrm{\Delta }n}(r,\text{R}),$$
(1)
independent of the amplitude or shape of $`w(r^{})`$, and that
$$\overline{\mathrm{\Delta }n}(r,\text{R})0,\text{monotonically as}\text{R}\mathrm{},$$
(2)
provided only that the chemical potential $`\mu `$ is held fixed.
The essence of NEM is contained in Eq. (1). Although it may not look very special, consider our perturbing potential $`w(r^{})`$, confined outside the sphere of radius $`R`$. The common sense says that, as we increase its strength, we need to increase the radius $`R`$ if we want to maintain its effect at the center of the sphere below a certain level. The common sense also says that we need to increase $`R`$ to infinity as we make $`w(r^{})`$ stronger and stronger. In reality, this is not so: if the chemical potential is kept fixed, these effects will saturate and, in fact, no matter what $`w(r^{})`$ we put outside the sphere of radius $`R`$, they cannot exceed a certain upper bound, which we will determine in this paper.
For a given $`\mathrm{\Delta }n`$, we can solve for R in $`\overline{\mathrm{\Delta }n}(r,\text{R})=\mathrm{\Delta }n`$ and define the nearsightedness range $`\text{R}(r,\mathrm{\Delta }n)`$. The significance of $`\text{R}(r,\mathrm{\Delta }n)`$ is the following: any perturbation beyond $`\text{R}(r,\mathrm{\Delta }n)`$, of arbitrary shape and amplitude, cannot produce a density change at $`r`$ larger than $`\mathrm{\Delta }n`$. $`\text{R}(r,\mathrm{\Delta }n)`$ provides a simple, quantitative measure of nearsightedness.
The above formulation of NEM often reminds people of Thomas-Fermi screening, sometimes even when we discuss insulators. However, let us point a few facts. If one puts a charge inside the uniform charged electron gas and calculates the density response, he will find that the Thomas-Fermi exponential screening is valid only very near the impurity. Further away from the impurity, he will see Friedel oscillations, decaying as an inverse power law.Friedel These oscillations are not negligible. In fact, the Friedel oscillations were observed experimentally, thought not directly, one year after their theoretical prediction.Rowland1 However, this was not realized until much later, when Walter Kohn made the connection between the two results.Rowland2 ; Kohn60 He showed that there is a big discrepancy between the prediction of the Thomas-Fermi screening theory and these experimental results and that a self-consistent calculation along Friedel’s lines, of the density change due to impurities in coper, brings the theory and experiment to almost perfect agreement. The picture that emerged was that, in the asymptotic region, the screening only renormalizes the amplitude of the Friedel oscillations. The whole issue was considered very important at that time, because it clearly demonstrated the existence of a sharp Fermi surface in real metals. We also like to mention one very well known fact in surface physics, where self-consistent calculations of metallic surfaces showed that the effective potential goes to the bulk value extremely fast, typically within one or two layers. However, the density oscillations extend much further into the bulk and they can be viewed as the Friedel oscillations generated by the screened surface potential. With these being said, we hope that the reader will dissociate, right from the begenning, NEM from the Thomas-Fermi exponential screening and the nearsightedness range from the Thomas-Fermi screening length.
Quantum gases display non-local density responses to local perturbations because of two factors. First, the effect of any local perturbation propagates to further distances through inter-particle interactions. This effect can be regarded as classical, since it manifests, in the same way, in classical gasses. Secondly, there is a purely quantum effect, that steams directly from the uncertainity principle. This paper is concerned with this purely quantum effect, so it neglegts the inter-particle interation effects, entirely. The fact that NEM exists for non-interacting systems is extremely important. To understand why, let us go back in time and recall that, at the beginning of the electronic structure calculations, when the exact diagonalization was the method of choice, people were facing the so called “exponential wall” when trying to extend the calculations to larger systems: because the number of operations in such calculations scales exponentially with the number of atoms, $`N`$, their applicability was, and is still limited to systems containg a few tens of atoms. Density Functional Theory (DFT)KohnHohenberg ; KohnSham provided a powerful alternative: because the number of operations in DFT calculations scales as $`N^3`$, we can now solve the electronic structure for systems containing hundreds of atoms. However, electronic structure calculations for biological and nano systems, or for extremely complex materials, involve thousands of atoms. At this scale, we start feeling the “$`N^3`$ wall.” Ab-initio quantum calculation for such complex systems will require a new generation of DFT algorithms, scaling linearly with the number of atoms. It is now generally accepted that NEM is the physical basis for these algorithms.Giullia
The quest for linearly scaling algorithms was initiated by W.T. Yang, who was the first to argue that O(N) algorithms are possible.Yang91 The algorithm proposed by Yang is known by the name of Divide and Conquer (DC). There are now several reviews on the linear scaling electronic structure calculations. We will mention here the one by GoedeckerGoedecker and the one by Wu and Jayanthi,Wu which, at the time of their publication, gave an exhaustive discussion of O(N) methodologies. If we examine carefully these methodologies, they are all based on the same original idea, namely gluing together calculations done for smaller systems. What is different, is the representation and the way the size of these smaller systems (the truncation) is determined. For example, the real space approaches will use the decay of the density matrix while the localized basis set approaches will use the overlap of these functions to judge how large these subsystems should be.
Let us focus on the original implementation given by Yang.Yang91 Consider a self-consistent DFT iteration process, for a large quantum system. Each iterative step consists in calculating the density of a non-interacting electron gas in equilibrium under a given effective potential (known from the previous iteration). In traditional approaches, this requires a number of operations that scales as $`N^3`$. DC algorithm, if it works as it is supposed to, requires a number of operations that scales linearly with $`N`$. It goes like this: The large system is devided in non-overlaping sub-regiones, which are then surrounded with buffer zones. A global chemical potential, $`\mu `$, is fixed and the orbitals are calculated and populated up to $`\mu `$, for each sub-region + buffer zone. The density in the buffer zones is discarded so, at this point, one has calculated a density for each sub-region and, by putting together all these sub-densities, one can construct the global density. The charge neutrality condition is then checked for this global density and the chemical potential is adjusted, if necessary. Note that charge neutrality must be satisfied by the entire system, not by each subsystem. In this way we have completed the DFT iteration step in a number of operations that scales linearly with the size of the system. Now, the question is how accurate is this algorithm? In fact, the most important question is, can we obtain arbitrary accuracy with this algorithm? To answer, we need to compare the density calculated for a sub-region + buffer zone and the density calculated for the entire system at once, at the same chemical potential $`\mu `$. Now one can see why NEM can be regarded as the basis for this algorithm: the artificial termination at the outer boundary of the buffer zone, no matter how it is done, represents the change in the effective potential in our formulation of NEM. For example, such changes of the effective potential occure when one calculates the density matrix and ignores the points outside a sub-region, or when one calculates the density and ignores the elements of a localized basis set that are centered outside a sub-region. Since we have no control on how the effective potential is modified by such truncations, Eq. (1) is paramount: it tells us that the effects of any artificial termination cannot exceed an upper bound. This upper bound is an intrinsic characteristic of the system: is independent of the method of termination. Eq. (2) tells us that, if we take the buffer zones large enough, the difference between the sub-density and the real density can be made smaller than any desired accuracy.
Examples and estimates of the width of the buffer zones, together with a discussion of how to optimize this algorithm in 1, 2 and 3D and how the CPU time scales with the the desired accuracy can be found in Ref. ProdanKohn, . We also like to mention that DC has been recently implemented to systems containing as many as 65,000 atoms and shown that it can generate electronic structures of the same quality as a traditional approach will do for a system of, let us say, 10 atoms.Vashishta05 The tests performed in this numerical work agree qualitatively with our theoretical predictions.
There is another important issue related to DC. The ground energy in DFT is given by:
$$E=\underset{j}{}ϵ_j+E_{xc}[n]v_{xc}(r)n(r)\frac{1}{2}\frac{n(r)n(r^{})}{|rr^{}|}.$$
(3)
All the above terms can be calculated directly from the density, except the first one. However, this term is just the integral of the energy density,
$$ϵ(r)=\underset{ϵ_iϵ_F}{}ϵ_i|\psi _i(r)|^2,$$
(4)
and we will show that $`ϵ(r)`$ is also nearsighted. As a consequence, within the DC algorithm, $`ϵ(r)`$ can be calculated with arbitrary precision, like the density.
The goal of this paper is two-fold. We want to prove NEM, i.e. Eqs. (1) and (2), for a simple system, which is the 1D non-interacting, spin $`\frac{1}{2}`$ fermions, in periodic potentials, and also want to show that one can obtain exact estimates of the nearsightedness range, which is extremely important for DC. The case of 1D non-interacting electrons is important for several reasons. Inspite of its simplicity, it captures all the important aspects of nearsightedness. This allows for a thorough investigation of NEM, while keeping the technical aspects at a reasonable level. The 1D non-interacting case is relevant for linear molecular chains when treated within DFT.
## II The strategy
We shall first develop general tools that will allow us to compute, for arbitrary perturbations, the asymptotic behavior of the density change,
$$\mathrm{\Delta }n(x)=2\underset{ϵ_i\mu }{}|\psi _i(x)|^22\underset{ϵ_i^0\mu }{}|\psi _i^0(x)|^2,$$
(5)
where $`\psi _i^0(x),ϵ_i^0`$ and $`\psi _i(x),ϵ_i`$ are the wave functions and the corresponding energies of the unperturbed and perturbed systems, respectively. The factor 2 in front of the sums comes from the spin. The above expression is not very useful when dealing with the asymptotic behavior of $`\mathrm{\Delta }n(x)`$. Instead, we will work with an integral representation. Why an integral representation? For answer, we point to the theory of special functions, where the functions are most often defined and introduced as infinite series but, with no exception, their asymptotic behavior is derived from equivalent integral representations.
We can obtain an integral representation of $`\mathrm{\Delta }n(x)`$ by using the Green’s functions. Indeed, if $`G_E^0(EH_0)^1`$ and $`G_E(EH)^1`$ denote the Green’s functions of the unperturbed and perturbed systems, respectively, then
$$\mathrm{\Delta }n(x)=\frac{1}{\pi }_𝒞[G_EG_E^0](x,x)𝑑E,$$
(6)
where $`𝒞`$ is a contour in the complex energy plane, surrounding the occupied states. This can be seen from the eigenfunction expansions of $`G_E^0`$ and $`G_E`$ and the residue theorem. Now, the eigenfunction expansions of $`G_E^0`$ and $`G_E`$ are, again, not very useful. Instead, we will use the following compact representation:
$$G_E(x,x^{})=\frac{\psi _<(x_<)\psi _>(x_>)}{W(\psi _<,\psi _>)},$$
(7)
where $`x_<=\mathrm{min}(x,x^{})`$ and $`x_>=\mathrm{max}(x,x^{})`$; $`\psi _<(x)`$ and $`\psi _>(x)`$ are the solutions of the Schrodinger equation at energy $`E`$, satisfying the boundary condition to the left and right, respectively, and $`W(\psi _<,\psi _>)`$ is the Wronskian of the two solutions. For infinite systems, the case considered in this paper, $`\psi _<(x)`$ and $`\psi _>(x)`$ are the solutions decaying at $`\pm \mathrm{}`$, respectively. We will always assume that $`E`$ does not belong to the energy spectrum. When the contour $`𝒞`$ intersects the energy spectrum, such as for the case of metals, $`G_E`$ will have a discontinuity at the point of intersection. Strictly speaking, this point must be excluded from $`𝒞`$, which does not change the result of the integration. For all the other points of $`𝒞`$, $`G_E`$ is uniquely defined and given by Eq. (7). Later, we will use the reflection and transmission coefficients to construct extremely simple and compact expressions of the Green’s functions (see Eqs. (15), (22) and (IV.2)).
The last step of our strategy will be to identify the special point in the complex energy plane that determines the asymptotic behavior of the integral Eq. (6).
This strategy will require from us to go into the complex energy plane. We will make use of the analytic structure of the Bloch functions and band energies derived in Ref. Kohn59, , which is briefly discussed in the next section. These analytic structure results can be generalized to linear molecular chains and even to 3D crystals.Prodan05 Also, the above expression for the Green’s function, Eq. (7), can be generalized to linear molecular chains,Prodan05 or to 3D crystals.Allen In fact, the entire strategy can be applied in 2 and 3D, as it was already shown in Ref. ProdanKohn, .
## III The unperturbed system
Throughout this paper, $`v(x)`$ will be taken as a periodic, inversion-symmetric potential. Following is a brief discussion, largely taken from Ref. Kohn59, , of the periodic Schrodinger equation.
The solutions of the periodic Schrodinger equation,
$$[d^2/dx^2+v(x)]\psi =E\psi ,v(x+b)=v(x),$$
(8)
are the well known Bloch functions $`\psi _k`$ ($`k`$ the wave vector) which will be normalized as in Ref. Kohn59, . Their fundamental property is $`\psi _k(x+b)=e^{ikb}\psi _k(x)`$. When dealing with complex values of $`k`$, is much more convenient to work with the variable $`\lambda =e^{ikb}`$, instead of $`k`$. Thus, from now on, we will index the Bloch functions by $`\lambda `$; their fundamental property becomes:
$$\psi _\lambda (x+b)=\lambda \psi _\lambda (x).$$
(9)
The parameter $`\lambda `$ relates to the energy of $`\psi _\lambda `$ through the following equation:
$$\lambda ^22\mu (E)\lambda +1=0,$$
(10)
whith $`\mu (E)`$ the Kramers’ function.Kramers35 By examining the fundamental property Eq. (9), one can see that the physical states correspond to the case $`|\lambda |=1`$ ($`\lambda `$ on the unit circle), otherwise $`\psi _\lambda (x)`$ explodes either to $`x=\mathrm{}`$ or $`x=\mathrm{}`$. For $`\lambda `$ on the unit circle, the solutions behave like waves, thus it is appropriate to use the term Bloch waves. When discussing arbitrary values of $`\lambda `$, however, it is more appropriate to use the term Bloch functions.
The energy spectrum consists of energy bands, indexed here by $`n=1,2,\mathrm{}`$, which are separated by energy gaps. The energy bands can be computed by solving for $`E`$ in Eq. (10) for all $`\lambda `$ on the unit circle. Due to the symmetry $`\lambda 1/\lambda `$ in Eq. (10), we can and shall restrict $`\lambda `$ to $`|\lambda |1`$, and view $`\psi _\lambda (x)`$ and $`\psi _{1/\lambda }(x)`$ as two independent wavefunctions. For $`|\lambda |<1`$, it follows from the fundamental property, Eq. (9), that $`\psi _\lambda (x)`$ decays to zero as $`x\mathrm{}`$ and $`\psi _{1/\lambda }(x)`$ decays to zero as $`x\mathrm{}`$.
If we restrict $`\lambda `$ to the unit disk, $`E`$ uniquely determines $`\lambda `$. The opposite is not true, instead $`E(\lambda )`$ is a multi-valued complex function, with branch points of order one at $`\lambda _1`$, $`\lambda _2`$, …. Each $`\lambda _n`$ is real and $`0<(1)^n\lambda _n<1`$. The corresponding energy, $`\stackrel{~}{E}_nE(\lambda _n)`$, is also real and located in the $`n`$-th gap. Estimates of $`\lambda _n`$ in the small gap and tight binding limits are given in Appendix A. $`E(\lambda )`$ can be represented on a Riemann surface with one sheet corresponding to one band. The $`n`$-th sheet can be taken as the entire unit disk, except a cut extending from $`\lambda _{n1}`$ to $`\lambda _n`$, as shown in Fig. 1. Near a branch point, $`E(\lambda )`$ behaves as the square root,
$$E(\lambda )=\stackrel{~}{E}_n+2\alpha _n(\lambda /\lambda _n1)^{1/2}+\mathrm{}.$$
(11)
The function $`E(\lambda )`$ on the $`n`$-th Riemann sheet will be denoted by $`E_n(\lambda )`$.
The integral in Eq. (6) will be mapped into the complex $`\lambda `$-plane, by changing the variable from $`E`$ to $`\lambda `$. Thus, it will be important to understand how the contour $`𝒞`$ looks in this plane. It is more easy to understand how a given contour $`\gamma `$ in the $`\lambda `$-plane is mapped in the complex energy plane, i.e. to construct the points $`E(\lambda )`$ when $`\lambda `$ sweeps through the points of $`\gamma `$. The mapping $`E_n(\lambda )`$, from the unit disk to the complex $`E`$-plane, is generic in one dimension, in the sense that it is qualitatively independent of the periodic potential. The general picture is as follows: $`E_n(\lambda )`$ maps the unit disk into a domain $`𝒟_n`$ (see Fig. 2). The domains $`𝒟_n`$, $`n=1,2,\mathrm{}`$, are disjoint (with the exception of a possible common boundary), and, all together, they cover the entire complex $`E`$-plane. Now let us consider several contours. In Fig. 2, we used the same line-style for a contour and its image. The thick lines in Fig. 2b represent the energy bands. Now consider the contours $`𝒞_1`$ and $`𝒞_2`$ in the $`\lambda `$-plane, starting and ending at zero and surrounding the branch cuts infinitely close. They are mapped into $`𝒞_1^{}`$ and $`𝒞_2^{}`$, shown in Fig. 2b. This figure displays only a finite sector of the two contours, which extend from $`i\mathrm{}`$ to $`+i\mathrm{}`$. The unit circle is mapped in a loop that surrounds the $`n`$-th band, infinitely close. The segment from $`1`$ to $`\lambda _{n1}`$ is mapped on the real axis, from the lower edge of the $`n`$-th band down to $`\stackrel{~}{E}_{n1}`$. The segment from $`\lambda _n`$ to $`+1`$ is also mapped on the real axis, from $`\stackrel{~}{E}_n`$ down to the upper edge of the $`n`$-th band. The domain $`𝒟_n`$ mentioned above, lies between the curves $`𝒞_1^{}`$ and $`𝒞_2^{}`$. From the above information, one should be able to construct, qualitatively, the image of any other contour.
$`\psi _\lambda (x)`$ and $`\psi _{1/\lambda }(x)`$ are multi-valued analytic functions of $`\lambda `$, with branch points of order 3 at $`\lambda _n`$. For $`\lambda `$ on the $`n`$-th Riemann sheet, $`\psi _\lambda (x)`$ will be denoted by $`\psi _{n,\lambda }(x)`$. Both functions diverge at the branch points as,
$$\psi _{n,\lambda }(x)=\frac{u_n(x)e^{q_nx}}{(\lambda /\lambda _n1)^{1/4}}+\mathrm{}$$
(12)
and
$$\psi _{n,1/\lambda }(x)=\frac{u_n^{}(x)e^{q_nx}}{(\lambda /\lambda _n1)^{1/4}}+\mathrm{},$$
(13)
where $`u_n(x)`$ and $`u_n^{}(x)`$ are periodic (antiperiodic) functions for $`n`$ even (odd) and $`q_n`$ is defined by $`|\lambda _n|=e^{q_nb}`$. The Wronskian of the two independent Bloch functions is given by
$$W(\psi _{n,1/\lambda },\psi _{n,\lambda })=\frac{b\lambda }{2\pi }\frac{dE_n(\lambda )}{d\lambda }.$$
(14)
Consequently, the Green’s function $`G_E^0`$ satisfies the identity:
$$G_E^0(x,x^{})\frac{dE}{\pi i}=2i\psi _{n,1/\lambda }(x_<)\psi _{n,\lambda }(x_>)\frac{d\lambda }{b\lambda }.$$
(15)
## IV The Effect of Perturbations
### IV.1 One Sided Perturbations
We consider here perturbations that are either to the left or to the right of the point $`x`$, where we measure the density change $`\mathrm{\Delta }n(x)`$. Let us assume that $`w`$ is to the left of $`x`$. For convenience, we choose the origin of $`x`$ at the right edge of $`w`$, so that $`w`$ is confined in the interval $`[L,0]`$, with $`L`$ arbitrarily large but finite. We calculate the particle and energy density changes at $`x>0`$.
As already mentioned, the density change $`\mathrm{\Delta }n(x)`$ is given by
$$\mathrm{\Delta }n(x)=\frac{1}{\pi i}_𝒞[G_E(x,x)G_E^0(x,x)]𝑑E,$$
(16)
where $`G_E^0`$ and $`G_E`$ are the unperturbed and perturbed Green’s functions, respectively, and $`𝒞`$ is a contour in the complex energy plane, surrounding the eigenvalues below $`ϵ_F`$ (see for example Fig. 3a). Similarly, the change of the energy density is
$$\mathrm{\Delta }ϵ(x)=\frac{1}{\pi i}_𝒞E[G_E(x,x)G_E^0(x,x)]𝑑E.$$
(17)
We can focus on $`\mathrm{\Delta }n(x)`$ and give only the final results for $`\mathrm{\Delta }ϵ(x)`$.
We construct the perturbed Green’s function from two independent solutions of the Schrodinger equation:
$$[d^2/dx^2+v(x)+w(x)]\psi (x)=E\psi (x).$$
(18)
Outside the interval $`[L,0]`$, the solutions are linear combinations of $`\psi _\lambda `$ and $`\psi _{1/\lambda }`$. As already mentioned, we need the solutions decaying to $`\mathrm{}`$, which can be conveniently written in terms of the reflection and transmission coefficients:
$$\psi _{n,\lambda }^<(x)=\{\begin{array}{c}T_n(\lambda )\psi _{n,1/\lambda }(x),x<L\hfill \\ \psi _{n,1/\lambda }(x)+R_n^+(\lambda )\psi _{n,\lambda }(x),x>0\hfill \end{array}$$
(19)
and
$$\psi _{n,\lambda }^>(x)=\{\begin{array}{c}\psi _{n,\lambda }(x)+R_n^{}(\lambda )\psi _{n,1/\lambda }(x),x<L\hfill \\ T_n(\lambda )\psi _{n,\lambda }(x),x>0,\hfill \end{array}$$
(20)
where $`E`$ was taken to be in $`𝒟_n`$ (Fig. 2). The Wronskian of the two independent solutions is
$$W(\psi _{n,\lambda }^<,\psi _{n,\lambda }^>)=\frac{b}{2\pi }\lambda T_n(\lambda )\frac{dE_n(\lambda )}{d\lambda },$$
(21)
leading to the following useful identity, for $`x>0`$:
$$[G_EG_E^0](x,x)\frac{dE}{\pi i}=2iR_n^+(\lambda )\psi _{n,\lambda }(x)^2\frac{d\lambda }{b\lambda }.$$
(22)
This identity, together with Eqs. (16) and (17), shows that, for $`x>0`$, $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$ are completely determined by the unperturbed wavefunctions and reflection coefficient, in a simple and universal way.
Eq. (22) also provides the analytic structure of the reflection coefficient: $`R_n^+(\lambda )`$ has branch points of order 1 at $`\lambda _{n1}`$ and $`\lambda _n`$. If we go around these branch points, $`R_n^+`$ becomes $`R_{n1}^+`$ and $`R_{n+1}^+`$, respectively. In other words, $`R_n^+`$ are different branches of a multi-valued function $`R^+(\lambda )`$. Near the branch points,
$$R^+(\lambda )=R^+(\lambda _n)+r_n^+(\lambda /\lambda _n1)^{1/2}+\mathrm{}.$$
(23)
The poles of $`R^+(\lambda )`$, if any, are mapped by $`E(\lambda )`$ into the poles of $`G_E`$, i.e. the energies of the bound states. Since the bound states are located in the gaps, the poles of $`R^+(\lambda )`$ are always located on the real axis and away from the branch cuts.
By evaluating the Wronskian of $`\psi _\lambda ^<(x)`$ and $`\psi _{1/\lambda }^<(x)`$ for $`x<L`$ and for $`x>0`$ and equating the two results, one can derive the following identity:
$$T(\lambda )T(1/\lambda )+R^+(\lambda )R^+(1/\lambda )=1.$$
(24)
For $`|\lambda |=1`$, this identity becomes $`|T(\lambda )|^2+|R^+(\lambda )|^2=1`$, showing that $`|R^+(\lambda )|1`$ for all $`\lambda `$ on the unit circle. Similar conclusion holds for $`R^{}(\lambda )`$.
#### IV.1.1 Metals
We assume the $`n`$-th band partially occupied and choose $`𝒞`$ in Eq. (16) as in Fig. 3, where $`\lambda _F`$ is defined by $`ϵ_F=E_n(\lambda _F)`$. We write
$$x=y+mb,$$
(25)
with $`y`$ restricted to the first unit cell. Mapping into the $`\lambda `$-plane by using Eq. (22) and recalling the fundamental property of the Bloch functions, for $`x>0`$, we obtain
$$\mathrm{\Delta }n(x)=\frac{2i}{b}_C^{}R_n^+(\lambda )\psi _{n,\lambda }(y)^2\lambda ^{2m1}𝑑\lambda .$$
(26)
An integration by parts gives
$`\mathrm{\Delta }n(x)`$ $`=`$ $`2\text{Im}{\displaystyle \frac{R_n^+(\lambda _F)\psi _{n,\lambda _F}(x)^2}{mb}}`$ (27)
$``$ $`{\displaystyle \frac{i}{b}}{\displaystyle _C^{}}{\displaystyle \frac{\lambda ^{2m}}{2m}}{\displaystyle \frac{d}{d\lambda }}\left[R_n^+(\lambda )\psi _{n,\lambda }(y)^2\right]𝑑\lambda .`$
The integral is of order $`1/x^2`$, as it can be seen from another integration by parts. We can conclude
$$\mathrm{\Delta }n(x)\frac{2}{x}\text{Im}[R_n^+(\lambda _F)\psi _{n,\lambda _F}(x)^2],$$
(28)
for large $`x`$. Similarly,
$$\mathrm{\Delta }ϵ(x)\frac{2ϵ_F}{x}\text{Im}[R_n^+(\lambda _F)\psi _{n,\lambda _F}(x)^2].$$
(29)
Since $`|R_n^+(\lambda _F)|1`$, the amplitudes of $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$ cannot exceed, in the asymptotic limit, the upper bounds
$$\overline{\mathrm{\Delta }n}(x)\frac{2}{x}|\psi _{n,\lambda _F}(x)|^2,$$
(30)
and
$$\overline{\mathrm{\Delta }ϵ}(x)ϵ_F\overline{\mathrm{\Delta }n}(x),$$
(31)
independent of the shape and amplitude of $`w(x)`$. This proves NEM for metals and one sided perturbations.
A comparison between Eq. (28) and an exact calculation of $`\mathrm{\Delta }n(x)`$ for the perturbed Kronig-Penney modelKronig
$$v_t(x)=\{\begin{array}{c}v_0_{l=\mathrm{}}^{\mathrm{}}\delta \left(x+\frac{b}{2}lb\right),x>0\hfill \\ V_0,x0,\hfill \end{array}$$
(32)
is shown in Fig. 4. Note that the asymptotic regime starts from about two lattice constants away from the perturbation.
#### IV.1.2 Insulators
We assume the first $`n`$ bands completely filled. For insulators, the calculations are more involved since the perturbing potential $`w(x)`$ may generate bound states in the insulating gap ($``$ the gap above the $`n`$-th band). There can be a discrete or a continuum set of states inside the insulating gap. When the set is discrete, there are four distinct possibilities, as shown in Fig. 5. Let us analyze these cases first.
No bound states in the insulating gap: We can take $`𝒞_1^{}`$ (with opposite orientation, see Fig. 2b) as the contour of integration in Eq. (16). Mapping into the complex $`\lambda `$-plane and using the fundamental property of the Bloch functions, gives
$$\mathrm{\Delta }n(x)=\lambda _n^{2m}\frac{2i}{b}\underset{𝒞_1}{}R_n^+(\lambda )\psi _{n,\lambda }(y)^2\left(\frac{\lambda }{\lambda _n}\right)^{2m1}\frac{d\lambda }{\lambda _n},$$
(33)
with $`𝒞_1`$ shown in Fig. 2a and $`y`$ and $`m`$ defined in Eq. (25). The integrand diverges at $`\lambda _n`$ as $`(\lambda \lambda _n)^{1/2}`$ but this singularity is integrable. Away from the branch point, the integrand is finite and $`(\lambda /\lambda _n)^{2m}`$ becomes small as we follow the contour $`𝒞_1`$ towards $`\lambda =0`$. Thus, for large $`m`$, the main contribution to the integral comes from the region in the immediate vicinity of the branch point. Expanding the integrand near $`\lambda _n`$ and keeping the leading term,Olver97 we find
$$\mathrm{\Delta }n(x)R_n^+(\lambda _n)\frac{2}{b}\underset{𝒞_1}{}\frac{i(\frac{\lambda }{\lambda _n})^{2m1}}{\sqrt{\frac{\lambda }{\lambda _n}1}}\frac{d\lambda }{\lambda _n}u_n(x)^2e^{2q_nx}$$
(34)
The integral is equal to $`2B(2m,1/2)`$ ($`B`$ = Beta function) and behaves asymptotically as $`\sqrt{2\pi /m}`$. We conclude,
$$\mathrm{\Delta }n(x)2R_n^+(\lambda _n)\left(\frac{2\pi }{xb}\right)^{1/2}u_n(x)^2e^{2q_nx}.$$
(35)
Similarly,
$$\mathrm{\Delta }ϵ(x)2R_n^+(\lambda _n)\stackrel{~}{E}_n\left(\frac{2\pi }{xb}\right)^{1/2}u_n(x)^2e^{2q_nx}.$$
(36)
An implementation of Eq. (35) to the perturbed Kronig-Penney model described in Eq. (32) is given in Fig. 6. Note again that the asymptotic regime starts from one or two lattice sites from the perturbation.
To end the proof of nearsightedness, we need to show that $`|R^+(\lambda _n)|`$ cannot exceed an upper bound. Since the reflection coefficient is not evaluated on the unit circle, the inequality $`|R^+(\lambda )|1`$ is no longer guaranteed. However, if $`h_\lambda ^<`$ and $`h_\lambda `$ denote the logarithmic derivatives at $`x=0`$ of $`\psi _\lambda ^<(x)`$ defined in Eq. (19), and of $`\psi _\lambda (x)`$, respectively, then
$$R^+(\lambda )=\frac{h_\lambda ^<h_{1/\lambda }}{h_\lambda ^<h_\lambda }\frac{\psi _{1/\lambda }(0)}{\psi _\lambda (0)}.$$
(37)
As in Ref. Kohn59, , we choose the phase of the Bloch functions such that $`\psi _\lambda (0)=\psi _{1/\lambda }(0)`$, so we can eliminate the last factor in the above expression. For $`E`$ in the insulating gap, $`\psi _\lambda ^<(x)`$, $`\psi _\lambda (x)`$ and $`\psi _{1/\lambda }(x)`$ are real functions and since
$$dh_\lambda ^</dE=\psi _\lambda ^<(0)^2_{\mathrm{}}^0\psi _\lambda ^<(x)^2𝑑x$$
(38a)
$$dh_{1/\lambda }/dE=\psi _{1/\lambda }(0)^2_{\mathrm{}}^0\psi _{1/\lambda }(x)^2𝑑x$$
(38b)
$$dh_\lambda /dE=\psi _\lambda (0)^2_0^{\mathrm{}}\psi _\lambda (x)^2𝑑x,$$
(38c)
it follows that $`h_\lambda =h_{1/\lambda }`$ and $`h_\lambda ^<`$ and $`h_{1/\lambda }`$ are decreasing functions of $`E`$. The typical behavior of $`h_\lambda `$ and $`h_{1/\lambda }`$ is shown in Fig. 7. Now, if there are no bound states in the gap, $`h_\lambda ^<`$ and $`h_\lambda `$ cannot be equal for any $`E`$ in the gap. Then, since $`h_\lambda `$ increases while $`h_\lambda ^<`$ decreases with $`E`$, $`h_\lambda ^<`$ can take values only in the shaded area of Fig. 7, below 0. Consequently, the right side of Eq. (37) is smaller or equal to 1, i.e.
$$|R^+(\lambda )|1,$$
(39)
remains valid when $`E`$ is in the insulating gap.
We can then conclude that the amplitudes of $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$, in the asymptotic limit, cannot exceed the upper bound
$$\overline{\mathrm{\Delta }n}(x)2\left(\frac{2\pi }{xb}\right)^{1/2}u_n(x)^2e^{2q_nx},$$
(40)
and
$$\overline{\mathrm{\Delta }ϵ}(x)\stackrel{~}{E}_n\overline{\mathrm{\Delta }n}(x).$$
(41)
This completes the proof of NEM for insulators and one sided perturbing potentials that do not generate bound states in the insulating gap.
Bound states in the insulating gap: We show in Appendix B that if
$$w(x)\psi _\pm (x)^2𝑑x0,$$
(42)
where $`\psi _\pm (x)`$ denotes the Bloch function at the upper/lower edge of an energy band, then $`w`$ generates bound states above/below this band, even for infinitely small coupling constants. Thus, the presence of bound states in the gaps is not a rare occurrence in one dimension.
The asymptotic forms of $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$ depend on how the bound states in the insulating gap are occupied. When all bound states below the branch point $`\stackrel{~}{E}_n`$ are occupied and the ones above $`\stackrel{~}{E}_n`$ are unoccupied, i.e. the situation illustrated in Fig. 5b, the asymptotic behavior, Eqs. (35) and (36), of $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$ remains unchanged.
Consider now that there are unoccupied bound states below $`\stackrel{~}{E}_n`$, as illustrated in Fig. 5c. Let $`\phi _0`$, of energy $`E_0`$, be such a state. For $`x>0`$, $`\phi _0`$ is equal, up to a factorization constant, to the exponentially decaying Bloch function of energy $`E_0`$ ($`E_0=E_n(\lambda _0)`$):
$$\phi _0(x)=\left[\frac{(1\lambda _0^2)\mathrm{\Lambda }}{_0^b|\psi _{n,\lambda _0}(x)|^2𝑑x}\right]^{1/2}\psi _{n,\lambda _0}(x),$$
(43)
where
$$\mathrm{\Lambda }_0^{\mathrm{}}|\phi _0(x)|^2𝑑x(\mathrm{\Lambda }1).$$
(44)
We have $`|\lambda _0|=e^{q_0b}`$, with $`q_0`$ strictly larger than zero, and $`\psi _{\lambda _0}(x)=e^{q_0x}u_0(x)`$, with $`u_0(x+b)=(1)^nu_0(x)`$. Since $`q_0`$ decreases as $`E_0`$ moves away from $`\stackrel{~}{E}_n`$, the first unoccupied state will have the slowest exponential decay, among all unoccupied states below $`\stackrel{~}{E}_n`$. Thus, when the contribution of these states is subtracted from Eq. (35), one finds that the asymptotic form of $`\mathrm{\Delta }n(x)`$ is determined by the first unoccupied bound state:
$$\mathrm{\Delta }n(x)\frac{2(1\lambda _0^2)\mathrm{\Lambda }}{_0^b|\psi _{n,\lambda _0}(x)|^2𝑑x}u_0(x)^2e^{2q_0x},$$
(45)
with the index $`0`$ referring to the first unoccupied bound state in the insulating gap. Similarly,
$$\mathrm{\Delta }ϵ(x)\frac{2(1\lambda _0^2)E_0\mathrm{\Lambda }}{_0^b|\psi _{n,\lambda _0}(x)|^2𝑑x}u_0(x)^2e^{2q_0x}.$$
(46)
Since $`\mathrm{\Lambda }1`$, the amplitudes of $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$ cannot exceed, in the asymptotic limit, the upper bounds
$$\overline{\mathrm{\Delta }n}(x)\frac{2(1\lambda _0^2)}{_0^b|\psi _{n,\lambda _0}(x)|^2𝑑x}u_0(x)^2e^{2q_0x},$$
(47)
and
$$\overline{\mathrm{\Delta }ϵ}(x)E_0\overline{\mathrm{\Delta }n}(x).$$
(48)
The results remain the same if, instead of unoccupied bound states below $`\stackrel{~}{E}_n`$, there are occupied bound states above $`\stackrel{~}{E}_n`$, as in Fig. 5d. In this case, the index 0 will refer to the last occupied bound state.
Continuum states in the insulating gap: We consider the case when $`w(x)`$ fills the entire insulating gap with continuum spectrum, such as when the insulator is in contact with an infinite metal. In this case, $`R_n^+(\lambda )`$ has a branch cut on the real axis. The states are considered occupied up to a Fermi energy, $`ϵ_F`$, which is in the insulating gap of the unperturbed insulator. We consider only the generic case when $`ϵ_F\stackrel{~}{E}_n`$ and define $`\lambda _F`$ by $`ϵ_F=E_n(\lambda _F)`$. This $`\lambda _F`$ is located strictly inside the unit circle, as opposed to the case of metals. We also define $`q_F>0`$ so that $`|\lambda _F|=e^{q_Fb}`$.
In Eq. (16), we consider the contour of integration shown in Fig. 8a. Mapping into the complex $`\lambda `$-plane and using again the fundamental property of the Bloch functions,
$$\mathrm{\Delta }n(x)=\lambda _F^{2m}\frac{2i}{b}\underset{𝒞}{}R_n^+(\lambda )\psi _{n,\lambda }(y)^2\left(\frac{\lambda }{\lambda _F}\right)^{2m1}\frac{d\lambda }{\lambda _F}.$$
(49)
With the new variable $`q`$ defined by $`\lambda /\lambda _F=e^{qb}`$, we can write
$$\mathrm{\Delta }n(x)=4\lambda _F^{2m}\text{Im}\underset{q>0}{}R_n^+(\lambda ^+)\psi _{n,\lambda }(y)^2e^{2mqb}𝑑q,$$
(50)
where $`\lambda ^+\lambda +i0^+`$. The asymptotic behavior can be extracted from a simple integration by parts:
$$\mathrm{\Delta }n(x)2\text{Im}[R_n^+(\lambda _F^+)]\frac{\psi _{n,\lambda _F}(x)^2}{x}.$$
(51)
By writing $`\psi _{n,\lambda _F}(x)`$ as $`u_F(x)e^{q_Fx}`$, with $`u_F(x+b)=(1)^nu_F(x)`$, we can conclude:
$$\mathrm{\Delta }n(x)2\text{Im}[R_n^+(\lambda _F^+)]u_F(x)^2\frac{e^{2q_Fx}}{x}.$$
(52)
Similarly,
$$\mathrm{\Delta }ϵ(x)2\text{Im}[R_n^+(\lambda _F^+)]ϵ_Fu_F(x)^2\frac{e^{2q_Fx}}{x}.$$
(53)
An implementation of Eq. (52) to the perturbed Kronig-Penney model Eq. (32) is shown in Fig. 9. Notice again that the asymptotic regime starts from two lattice constants away from the perturbation.
To end the proof of NEM, we need to give an upper bound on the amplitudes of $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$. The imaginary part of $`R_n^+(\lambda _F^+)`$ is proportional to the local density of states, $`g(E,x)`$, at $`E=ϵ_F`$ and $`x=0`$. Indeed, Eq. (22) provides the following identity:
$$\text{Im}[R_n^+(\lambda _F^+)]\psi _{n,\lambda _F}(x)^2=\frac{1}{2\pi }\frac{dϵ_F}{dq_F}\text{Im}[G_{ϵ_F+i0}(x,x)],$$
(54)
leading to
$$\text{Im}[R_n^+(\lambda _F^+)]=\frac{dϵ_F/dq_F}{2\psi _{n,\lambda _F}(0)^2}g(ϵ_F,0).$$
(55)
Note that the coefficient in front of $`g(ϵ_F,0)`$ is determined by the unperturbed system. If we limit ourselves to the generic case of $`w^{}`$s that generate finite densities of states at $`ϵ_F`$, then NEM follows from Eqs. (52) and (55). For practical applications, we consider this argument sufficient. However, to achieve a full proof of NEM, we need to consider also the cases when $`g(E,0)`$ diverges (or becomes extremely large) as $`Eϵ_F`$. For this special cases, the asymptotic form of Eq. (50) cannot be extracted from a simple integration by parts and $`\mathrm{\Delta }n(x)`$ is no longer given by Eq. (52); its specific functional form will depend on the type of singularity of $`g(E,0)`$. This special situations will not be discussed here.
### IV.2 Two Sided Perturbations
We consider here the case when the point $`x`$, where we evaluate $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$, has perturbing potentials $`w_L`$ to the left and $`w_R`$ to the right, as schematically shown in Fig. 10. For one sided perturbations, $`\mathrm{\Delta }n(x)`$ decays as $`x`$ moves further and further away from the perturbation and, because of this simple picture, NEM is intuitive and simply to grasp. When left and right perturbing potentials are present, this simple picture is gone: there will be interference terms in $`\mathrm{\Delta }n(x)`$, whose amplitude remain constant in the region between the two perturbing potentials. In addition, $`w_L+w_R`$ can induce a strong, qualitative change of the system, namely, the energy bands may become quantized. In spite all of these, we will show the following: for metals, the interference terms are not negligible, but $`\mathrm{\Delta }n(x)`$ still remains bounded. For insulators, the interference terms are exponentially small and $`\mathrm{\Delta }n(x)`$ is given by the simple superposition of the left and right density changes. Similar conclusions hold for $`\mathrm{\Delta }ϵ(x)`$.
For convenience, we fix the origin in the middle of the interval that separates the two perturbing potentials and consider the distance, R, from the origin to the right/left edge of $`w_{L/R}`$ to be an integer of $`b`$, $`\text{R}=Nb`$. We are interested in the behavior, for $`x`$ near the origin, of $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$ when $`\text{R}\mathrm{}`$. We follow our general strategy and derive first an expression for the Green’s function on the interval $`[\text{R},\text{R}]`$. We look again for the solutions of the Schrodinger equation,
$$[d^2/dx^2+v(x)+w_L(x)+w_R(x)]\psi (x)=E\psi (x),$$
(56)
which decay at $`\mathrm{}`$. On the intervals from $`\mathrm{}`$ to the left edge of $`w_L`$ and from the right edge of $`w_R`$ to $`+\mathrm{}`$, these solutions can be expressed as in Eqs. (19) and (20), in terms of the total (corresponding to $`w_L+w_R`$) transmission and reflection coefficients. Then, one can use again the reflection and transmission coefficients to continue these solutions inside the interval $`\text{R}<x<\text{R}`$. On this interval, they take the following form:
$$\psi _\lambda ^>(x)=\frac{T(\lambda )}{\stackrel{~}{T}_R(\lambda )}[\psi _\lambda (x)+\stackrel{~}{R}_R^{}(\lambda )\psi _{1/\lambda }(x)]$$
(57)
and
$$\psi _\lambda ^<(x)=\frac{T(\lambda )}{\stackrel{~}{T}_L(\lambda )}[\psi _{1/\lambda }(x)+\stackrel{~}{R}_L^+(\lambda )\psi _\lambda (x)],$$
(58)
where $`\stackrel{~}{T}_{L,R}(\lambda )`$ and $`\stackrel{~}{R}_{L,R}^\pm (\lambda )`$ are the transmission and reflection coefficients of the left/right potentials, and $`T(\lambda )`$ is the total transmission coefficient,
$$T(\lambda )=\frac{\stackrel{~}{T}_L(\lambda )\stackrel{~}{T}_R(\lambda )}{1\stackrel{~}{R}_L^+(\lambda )\stackrel{~}{R}_R^+(\lambda )}.$$
(59)
All these coefficients depend on R: If $`R_{L,R}^\pm (\lambda )`$ denote the reflection coefficients when the right/left edge of $`w_{L/R}`$ is at $`x=0`$ (thus $`R_{L,R}^\pm (\lambda )`$ are independent of R), then Eq. (37) gives:
$$\stackrel{~}{R}_{L,R}^\pm (\lambda )=\lambda ^{2N}R_{L,R}^\pm (\lambda ).$$
(60)
The Wronskian of the two independent solutions is the same as given in Eq. (21). From the two independent solutions, Eqs. (57) and (58), and their Wronskian, we derive, for $`\text{R}<x<\text{R}`$, the following identity:
$`[G_E(x,x)G_E^0(x,x)]{\displaystyle \frac{dE}{\pi i}}`$
$`=`$ $`{\displaystyle \frac{2i\lambda ^{2N}R_L^+(\lambda )}{1\lambda ^{4N}R_L^+(\lambda )R_R^{}(\lambda )}}\psi _\lambda (x)^2{\displaystyle \frac{d\lambda }{b\lambda }}`$
$`+`$ $`{\displaystyle \frac{2i\lambda ^{2N}R_R^{}(\lambda )}{1\lambda ^{4N}R_L^+(\lambda )R_R^{}(\lambda )}}\psi _{1/\lambda }(x)^2{\displaystyle \frac{d\lambda }{b\lambda }}`$
$`+`$ $`{\displaystyle \frac{4i\lambda ^{4N}R_L^+(\lambda )R_R^{}(\lambda )}{1\lambda ^{4N}R_L^+(\lambda )R_R^{}(\lambda )}}\psi _\lambda (x)\psi _{1/\lambda }(x){\displaystyle \frac{d\lambda }{b\lambda }}.`$
We now can use Eqs. (16) and (17) to find $`\mathrm{\Delta }n(x)`$ and $`\mathrm{\Delta }ϵ(x)`$, for $`x`$ near the origin. By using the fundamental property of the Bloch functions, we can understand the behavior of each term in the above identity. For $`\lambda `$ not on the unit circle, as it is the case in our integrals, the first term decay exponentially as $`x`$ moves to the right; the second term decay exponentially as $`x`$ moves to the left, but the third term is periodic, with period $`b`$. Fortunately, the amplitude of this term becomes smaller and smaller as the two perturbing potentials are moved apart.
From Eq. (IV.2), one can easily obtain the energy spectrum of the perturbed system. So, let us discuss first how the simultaneous presence of $`w_{L/R}`$ affects the energy spectrum. We are interested in the last occupied band (indexed by $`n`$), so, from now on, $`\lambda `$ is considered on the $`n`$-th Riemann sheet. The energies of the discrete states are given by the poles of the Green’s function. From the identity Eq. (IV.2), one can see that these poles correspond to those $`\lambda `$ satisfying the equation
$$R_L^+(\lambda )R_R^{}(\lambda )=1/\lambda ^{4N}.$$
(62)
Since the discrete state energies are real, the solutions of Eq. (62) are always located on the unit circle or on the real axis, away from the branch cuts. Inside the unit disk we have $`|\lambda |<1`$; consequently, $`|1/\lambda ^{4N}|`$ becomes very large in the limit $`\text{R}\mathrm{}`$. Thus, if there are solutions of Eq. (62) inside the unit disk, they must be located very close to the poles of either $`R_L^+(\lambda )`$ or $`R_R^{}(\lambda )`$. In other words, these solutions are perturbations of the bound states generated by $`w_L`$ (or $`w_R`$) alone, already discussed in the previous subsection. Poles on the unit circle exist if and only if both $`|R_{L,R}^\pm (\lambda )|`$ are equal to 1. In this case, the energy band degenerates into discrete energy spectrum. Because the left hand side of Eq. (62) is slowly varying compared to the right hand side, one can get the qualitative picture by setting the left hand side constant. If the amplitudes of $`R_{L,R}^\pm (\lambda )`$ are equal to 1 on the whole unit circle, then Eq. (62) has $`4N`$ solutions, $`\lambda _k`$, distributed on the unit circle; the spacing between two consecutive solutions is $`2\pi /4N+O(1/N^2)`$. This is the picture in the $`\lambda `$-plane. In the $`E`$-plane, the discrete energies are given by $`E(\lambda _k)`$. The spacing between two consecutive energies is $`2\pi _\lambda E(\lambda _k)/4N+O(1/N^2)`$.
#### IV.2.1 Metals
The effects of band quantization will be the strongest for metals, since the Fermi energy lies inside the band. We calculate the density change by integrating Eq. (IV.2) along the contour of integration shown in Fig. 3a and map the integral into the complex $`\lambda `$-plane. The asymptotic behavior of $`\mathrm{\Delta }n(x)`$, for large R, is determined by the behavior of the integrand near $`\lambda _F`$ and $`1/\lambda _F`$ (see Fig. 3b). In the immediate vicinity of these points, we can replace the slowly varying functions in the right hand side of Eq. (IV.2) \[the reflection coefficients and the Bloch functions\] with their value at $`\lambda _F`$ and $`1/\lambda _F`$. The integral then can be explicitly calculated and the result is:
$`\mathrm{\Delta }n(x,\text{R})`$ $``$ $`{\displaystyle \frac{2}{\text{R}}}\text{Im}{\displaystyle \frac{\mathrm{tanh}^1[\lambda _F^{2N}(R_L^+R_R^{})^{1/2}]}{(R_L^+R_R^{})^{1/2}}}`$
$`\times `$ $`[R_L^+\psi _{\lambda _F}(x)^2+R_R^{}\psi _{1/\lambda _F}(x)^2]`$
$`+`$ $`{\displaystyle \frac{2}{\text{R}}}\text{Im}[\mathrm{ln}(1\lambda _F^{4N}R_L^+R_R^{})]|\psi _{\lambda _F}(x)|^2,`$
with the reflection coefficients evaluated at $`\lambda _F`$.
To prove NEM, we need to find an upper bound on the above expression. A simple analysis reveals that the largest density changes occur when $`|R_{L,R}^\pm |=1`$, i.e. when the band is quantized at the Fermi energy. In this case, we can rewrite Eq. (IV.2.1) as:
$`\mathrm{\Delta }n(x,\text{R})`$ $``$ $`{\displaystyle \frac{4}{\text{R}}}\text{Im}\left[\mathrm{tanh}^1[\lambda _F^{2N}(R_L^+R_R^{})^{1/2}]\right]`$
$`\times `$ $`\text{Re}[(R_L^+R_R^{})^{1/2}\psi _{\lambda _F}(x)^2]`$
$`+`$ $`{\displaystyle \frac{2}{\text{R}}}\text{Im}[\mathrm{ln}(1\lambda _F^{4N}R_L^+R_R^{})]|\psi _{\lambda _F}(x)|^2.`$
The density change, as function of R, has discontinuities every time when $`\lambda _F^{4N}R_L^+R_R^{}=1`$, i.e. when a discrete energy crosses the Fermi level. These discontinuities are finite: since $`|\text{Im}[\mathrm{tanh}^1(z)]|\pi /4`$ and $`|\text{Im}[\mathrm{ln}(1z)]|\pi /2`$, for $`|z|1`$, the amplitude of the asymptotic term of $`\mathrm{\Delta }n(x,\text{R})`$ cannot exceed the upper bound
$$\overline{\mathrm{\Delta }n}(x,\text{R})\frac{2\pi }{\text{R}}|\psi _{\lambda _F}(x)|^2,$$
(65)
independent of $`w_{L/R}`$ potentials and of the position of the Fermi energy relative to the discretized energies. Similarly, the amplitude of $`\mathrm{\Delta }ϵ(x,\text{R})`$ cannot exceed the upper bound
$$\overline{\mathrm{\Delta }ϵ}(x,\text{R})ϵ_F\overline{\mathrm{\Delta }n}(x,\text{R}),$$
(66)
and this completes our discussion of NEM for metals.
These upper bounds are optimal, in the sense that there are $`w_{L/R}`$ potentials (the worst scenario) that generate a density and an energy density that are equal to these upper bounds. By comparing with the results of the previous Section, one can see that interference has non-trivial effects: these upper bonds are not simply the superposition of the left and right upper bounds.
#### IV.2.2 Insulators
We consider first the situation when there are no bound states in the insulating gap. We can take $`𝒞_1^{}`$ of Fig. 2 as the contour of integration in Eq (16), which is mapped into $`𝒞_1`$ in the complex $`\lambda `$-plane. For any $`\lambda `$ on this curve, the denominators in the right side of Eq. (IV.2) goes exponentially to 1 as $`\text{R}\mathrm{}`$. Consequently, in this limit, the structure of the integrand becomes completely analogous with the one studied in the previous subsection. The asymptotic behavior can be extracted as previously and the result is
$`\mathrm{\Delta }n(x,\text{R})2R_L^+(\lambda _n)\left({\displaystyle \frac{2\pi }{b\text{R}}}\right)^{1/2}u_n(x)^2e^{2q_n\text{R}}`$
$`2R_R^{}(\lambda _n)\left({\displaystyle \frac{2\pi }{b\text{R}}}\right)^{1/2}u_n^{}(x)^2e^{2q_n\text{R}}`$
$`2R_L^+(\lambda _n)R_R^{}(\lambda _n)\left({\displaystyle \frac{\pi }{b\text{R}}}\right)^{1/2}u_n(x)u_n^{}(x)e^{4q_n\text{R}}.`$ (67)
Thus, in the limit of large R, $`\mathrm{\Delta }n(x,\text{R})`$ is just the sum of the independent density changes due to the left and right potentials, plus an exponentially small correction. From the previous subsection, we can conclude that the amplitude of $`\mathrm{\Delta }n(x,\text{R})`$ cannot exceed, for large R, the upper bound
$$\overline{\mathrm{\Delta }n}(x,\text{R})2\left(\frac{2\pi }{b\text{R}}\right)^{1/2}[u_n(x)^2+u_n^{}(x)^2]e^{2q_n\text{R}}.$$
(68)
Similarly,
$$\overline{\mathrm{\Delta }ϵ}(x,\text{R})\stackrel{~}{E}_n\overline{\mathrm{\Delta }n}(x,\text{R}).$$
(69)
For the case when there are bound or continuum states in the insulating gap, the conclusion is the same: for large R,, the density change near the origin is the sum of the independent changes induced by the left and right potentials. Upper bounds on $`\mathrm{\Delta }n(x,\text{R})`$ can be trivially derived from the previous subsection.
## V The Nearsightedness Range
The nearsightedness range $`\text{R}(x,\mathrm{\Delta }n)`$ was introduced as the range beyond which any perturbation, no matter how large, induces a density change at $`x`$ less than the given $`\mathrm{\Delta }n`$. The asymptotic $`\text{R}(x,\mathrm{\Delta }n)`$, in the limit $`\mathrm{\Delta }n0`$, can now be easily calculated from the upper bounds on $`\mathrm{\Delta }n(x)`$, derived in this paper. Since the periodic systems are macroscopically homogeneous, the asymptotic $`\text{R}(x,\mathrm{\Delta }n)`$ will be independent of $`x`$.
When solving for R in $`\overline{\mathrm{\Delta }n}(x,\text{R})=\mathrm{\Delta }n`$, we first average $`\overline{\mathrm{\Delta }n}(x,\text{R})`$ over one unit cell. For metals, Eq. (65) leads to the following asymptotic expression:
$$\text{R}(\mathrm{\Delta }n)1/\mathrm{\Delta }n.$$
(70)
Such universal behavior is characteristic only to 1 dimension; in higher dimensions, the nearsightedness range will depend on the average particle density.ProdanKohn
For insulators and $`w^{}s`$ that generate no bound states in the insulating gap, Eq. (68) leads to
$$\text{R}(\mathrm{\Delta }n)\frac{1}{2q_n}\mathrm{ln}\frac{\stackrel{~}{n}}{\mathrm{\Delta }n},$$
(71)
where
$$\stackrel{~}{n}=\frac{4\sqrt{2\pi q_n}}{b}_0^b[u_n(x)^2+u_n^{}(x)^2]𝑑x.$$
(72)
In the small gap and tight binding limits, $`\stackrel{~}{n}`$ is completely determined by the exponential decay constant $`q_n`$, $`\stackrel{~}{n}4q_n\sqrt{2/\pi }`$ and $`\stackrel{~}{n}4\sqrt{q_n/\pi b}`$, respectively.
It is important to notice that the nearsightedness range does not depend on the details of the underlying potential $`v(x)`$, but on some simple parameters that can be defined also for non periodic potentials. For example, $`q_n`$ can be identified with the exponential decay constant of the density matrix.
## VI Discussion
The above analysis provides a quantitative analysis of the nearsightedness of electronic matter for non-interacting fermions, moving in one dimension under the action of periodic potentials. Although the simplest case possible, it allowed us to understand the different mechanisms behind NEM. Although we cannot point to one simple and general physical explanation of NEM, it is now clear that NEM is due to a destructive interference not of the wave amplitudes but of density amplitudes $`n_j`$ associated with the single particle eigenstates $`\psi _j`$. The asymptotic behavior of $`\mathrm{\Delta }n(x)`$ was found to be determined by the reflection coefficient. For specific cases, the amplitude of $`\mathrm{\Delta }n(x)`$ cannot exceed an upper bound simply because, when evaluated at allowed energies, the reflection coefficients are always smaller than 1. More general, and now including 2 and 3D, one will find that asymptotic behavior of $`\mathrm{\Delta }n(x)`$ is determined by certain elements of the scattering matrix and, for specific cases, the unitarity of the scattering matrix imposes certain upper bounds. The situation is, however, more complicated when bound states appear in the insulating gap or when the bands become quantized.
We have introduced a new concept, the nearsightedness range, $`\text{R}(x,\mathrm{\Delta }n)`$, which is well defined only because there is this upper bound on $`\mathrm{\Delta }n(x)`$. $`\text{R}(x,\mathrm{\Delta }n)`$ is a characteristic of the unperturbed system and gives a simple and effective measure of nearsightedness. For periodic metals, we found $`\text{R}(x,\mathrm{\Delta }n)`$ to have, in the asymptotic limit $`\mathrm{\Delta }n0`$, a universal expression, namely $`1/\mathrm{\Delta }n`$. For insulators, $`\text{R}(x,\mathrm{\Delta }n)`$ is strongly dependent on the band structure but has a weak, logarithmic dependence on $`\mathrm{\Delta }n`$.
Although the estimates given in this paper can be applied only to 1D systems, we think we gain some knowledge that can be useful for more general situations. We are convinced that NEM exists in dimensions higher than 1, where it can be quantified in a similar way. In particular, we believe that a complete theoretical analysis and optimization of the O(N) divide and conquer algorithm is possible in all dimensions. Preliminary results in this direction have been already given in Ref. ProdanKohn, . The one dimension analysis proved to be extremely useful by providing a viable strategy and some understanding of the effects of the bound states in the in insulating gap and of the band quantization on NEM.
###### Acknowledgements.
I would like to thank Prof. Walter Kohn, who suggested and supervised this project. This work was completed when the author was visiting the Physics Department at the University of California at Santa Barbara and was supported by Grants No. NSF-DMR03-13980, NSF-DMR04-27188 and DOE-DE-FG02-04ER46130.
## Appendix A
We estimate here the exponential decay rate $`q_n`$, related to the branch point by $`|\lambda _n|=e^{q_nb}`$. According to Ref. Kohn59, :
$$q_n=\frac{1}{b}\mathrm{ln}[|\mu _n|+\sqrt{\mu _n^21}],$$
(73)
where $`\mu _n`$ is the Krammers function evaluated at the branch point $`\stackrel{~}{E}_n`$, defined by $`d\mu /dE|_{E=\stackrel{~}{E}_n}=0`$.
For small gaps, the behavior of $`\mu (E)`$ inside the entire gap is well approximated by a quadratic function of $`E`$:
$$\mu (E)(1)^n\left[1\frac{m_n^{}b^2(EE_n^+)(EE_{n+1}^{})}{2(E_{n+1}^{}E_n^+)}\right],$$
(74)
where $`E_n^\pm `$ is the upper/lower edge of the $`n`$-th band, and $`m_n^{}`$ is the effective mass at the upper edge of the $`n`$-th band. Since $`\mu _n1`$, $`q_n\frac{1}{b}\sqrt{2(|\mu _n|1)}`$, which, together with Eq. (74), lead to
$$q_n=\frac{1}{2}\sqrt{m_n^{}(E_{n+1}^{}E_n^+)}.$$
(75)
We consider now a periodic potential $`_lV_a(xlb)`$, where $`V_a(x)`$ vanishes for $`|x|>c`$ and has atomic levels $`E_nk_n^2`$, $`n=1,\mathrm{}`$. In the limit $`b\mathrm{}`$, we show that
$$q_n\frac{1}{b}\mathrm{ln}\frac{8\sqrt{E_n}}{ebW_n},$$
(76)
where $`W_n`$ is the width of the $`n`$-th energy band. For $`x`$ in $`[b/2,b/2]`$ and $`|x|>c`$, the solutions of the Schrodinger equation at an energy $`E=k^2`$ are of the general form
$$\psi (x)=\{\begin{array}{c}a_{}(k)e^{kx}+b_{}(k)e^{kx},x<c\\ a_+(k)e^{kx}+b_+(k)e^{kx},x>c,\end{array}$$
(77)
with
$$\left(\begin{array}{c}a_+(k)\\ b_+(k)\end{array}\right)=\widehat{T}(k)\left(\begin{array}{c}a_{}(k)\\ b_{}(k)\end{array}\right),$$
(78)
$`\widehat{T}(k)`$ being the transfer matrix of the potential $`V_a`$. The energy levels of $`V_a`$ correspond to the zeroes of $`T_{22}(k)`$, already denoted by $`k_n`$. The Kramers function is given by
$$\mu (k)=\frac{1}{2}\left[T_{11}(k)e^{kb}+T_{22}(k)e^{kb}\right].$$
(79)
We estimate first the bandwidths. We look for the solutions of $`\mu (k)=\pm 1`$, which give the band edges. For $`b`$ large, the solutions of this equation must be located very close to the zeros of $`T_{22}(k)`$ since, otherwise, the second term in Eq. (79) becomes very large. We can then linearize, $`T_{22}(k)(kk_n)T_{22}^{}(k_n)`$ and neglect the exponentially small term in Eq. (79), in which case the equation $`\mu (k)=\pm 1`$ can be trivially solved, leading to
$$W_n=\frac{8k_ne^{k_nb}}{|T_{22}^{}(k_n)|}.$$
(80)
We now calculate $`\stackrel{~}{E}_n=\stackrel{~}{k}_n^2`$, defined byKohn59
$$\left(\frac{d\mu }{dk}\right)_{k=\stackrel{~}{k}_n}=0T_{22}^{}(\stackrel{~}{k}_n)bT_{22}(\stackrel{~}{k}_n).$$
(81)
For $`b`$ large, the solutions of the above equation must also be close to the zeroes of $`T_{22}(k)`$. Linearizing $`T_{22}(k)`$, we find $`\stackrel{~}{k}_n=k_n1/b`$ and the Kramers function evaluated at $`\stackrel{~}{k}_n`$ is
$$\mu _n=\frac{T_{22}^{}(k_n)e^{k_nb}}{2eb}=\frac{4k_n}{ebW_n}.$$
(82)
Since $`\mu _n1`$, $`q_n\frac{1}{b}\mathrm{ln}[2\mu _n]`$ and Eq. (76) follows.
## Appendix B
Let $`w(x)`$ be a perturbing potential of finite support and such that
$$w(x)\psi _+(x)^2𝑑x>0$$
(83)
or
$$w(x)\psi _{}(x)^2𝑑x<0,$$
(84)
where $`\psi _\pm (x)`$ is the Bloch function at the upper/lower edge of an energy band. We show here that, even for infinitely small coupling constants, such potential will pull bound states out from the band.
Let $`H_0`$ denote the periodic Hamiltonian and $`HH_0+\gamma w`$. It is known that $`H`$ has a bound state at some energy $`E`$ if and only if the operatorSimonTr
$$\widehat{K}_E=\gamma w^{1/2}(EH_0)^1|w|^{1/2}$$
(85)
has an eigenvalue equal to 1.SimonTr Here, $`w^{1/2}=w/|w|^{1/2}`$. We show that, for any given energy $`E`$ below/above the band, $`\widehat{K}_E`$ has an eigenvalue equal to 1 for some positive $`\gamma `$, which decreases to zero as $`E`$ approaches the edges of the band, provided the condition Eq. (83)/(84) is satisfied.
If $`n`$ is odd, the lower edge of the band corresponds to $`\lambda =1`$. We take an energy $`E`$ below such band and let $`\lambda `$, which is real and less than 1, be its corresponding $`\lambda `$-parameter. Eq. (15) gives
$$\widehat{K}_E(x,x^{})=\frac{2\pi \gamma }{b}w(x)^{1/2}\frac{\psi _{1/\lambda }(x_<)\psi _\lambda (x_>)}{\lambda dE/d\lambda }|w(x^{})|^{1/2},$$
(86)
and we notice that $`dE/d\lambda \lambda 1`$, for $`\lambda 1`$, i.e. the kernel $`\widehat{K}_E(x,x^{})`$ diverges at $`\lambda =1`$. We can separate the diverging part by expanding
$`\psi _{\lambda ^1}(x_<)\psi _\lambda (x_>)`$ $`=`$ $`\psi _{}(x)\psi _{}(x^{})`$ (87)
$`+`$ $`(\lambda 1)W_\lambda (x,x^{}).`$
This provides the following decomposition,
$$\widehat{K}_E=\gamma \alpha (\lambda )|\phi _1\phi _2|+\gamma \widehat{A}(\lambda ),$$
(88)
where
$$\alpha \frac{2\pi }{b\lambda dE/d\lambda },$$
(89)
$$\widehat{A}(\lambda )\frac{2\pi }{b}\frac{1\lambda }{\lambda dE/d\lambda }w^{1/2}W_\lambda |w|^{1/2},$$
(90)
and
$$\{\begin{array}{c}\phi _1(x)w(x)^{1/2}\psi _{}(x)\hfill \\ \phi _2(x)|w(x)|^{1/2}\psi _{}(x).\hfill \end{array}$$
(91)
The first term of Eq. (88) diverges while the second one is analytic at $`\lambda =1`$. Now let $`\mathrm{\Psi }`$ be given by
$$\mathrm{\Psi }=(\gamma \widehat{A}(\lambda )1)^1\phi _1,$$
(92)
which is well defined for small $`\gamma `$. Then
$$\widehat{K}_E\mathrm{\Psi }=\mathrm{\Psi }+[1+\gamma \alpha (\lambda )\phi _2|(\gamma \widehat{A}(\lambda )1)^1|\phi _1]\phi _1.$$
(93)
In other words, $`\widehat{K}_E`$ has an eigenvalue at +1, if
$$1+\gamma \alpha (\lambda )\phi _2|(\gamma \widehat{A}(\lambda )1)^1|\phi _1=0.$$
(94)
We can rewrite this equation as
$$\gamma \phi _2|(\gamma \widehat{A}(\lambda )1)^1|\phi _1\frac{b\lambda }{2\pi }\frac{dE}{d\lambda }=0.$$
(95)
If we denote the left side with $`F(\gamma ,\lambda )`$, then $`F(0,1)=0`$ and $`_\gamma F(0,1)=w(x)\psi _{}(x)^2𝑑x>0`$, i.e. the conditions of the analytic implicit function theorem are satisfied, which means that, for any $`\lambda `$ near $`+1`$, there is always a solution $`\gamma (\lambda )`$ to the Eq. (94). Moreover,
$$\gamma (\lambda )=\frac{b\lambda }{2\pi }\frac{dE}{d\lambda }\left[w(x)\psi _{}(x)^2𝑑x\right]^1+\mathrm{},$$
(96)
where the dots indicate terms of order $`o[(1\lambda )^2]`$. $`\gamma `$ is real and positive, for $`E`$ below the band, and goes to zero as $`E`$ approaches the band edge.
The other possible cases, $`\lambda =1`$ and $`n`$ even, follow in the same way.
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# Further spectroscopic observations of the CSL-1 object
## 1 Introduction
As recently stated in the masterly review by Kibble (2004), the last two years have seen a renewal of interest in the cosmological role of cosmic strings, after more than a decade of relative quiescence: an interest which was mainly triggered by the discovery of the unusual object CSL-1 (R.A.$`{}_{2000}{}^{}=12^h23^m30\mathrm{"}5`$; $`\delta _{2000}=12^{}38^{}57\mathrm{"}0`$) (Sazhin et al., 2003, hereafter Paper I) in the OAC-DF (Osservatorio Astronomico di Capodimonte - Deep Field; Alcalá et al., 2004) and by the indirect evidence obtained by Schild (2004) from the luminosity fluctuations of the quasar $`Q0957+561A,B`$.
CSL-1 (see Fig.1) is a double source laying in a low density field. In the original images the components, 1.9 arcsec apart, appeared to be extended and with roundish and identical shapes. By low resolution spectroscopy we learned that both components are at a redshift of $`0.46\pm 0.008`$, and by photometry (both global properties and luminosity profiles) that they match the properties of two giant elliptical galaxies. Detailed analysis showed that the spectra of the two components were identical at a 99.9% level. Such a conclusion, however, was hindered both by the limited wavelength range spanned by the spectra, and by their relatively low signal-to-noise ratio.
As discussed in Paper I, the only possible explanation of the CSL-1 properties is: i) either an unlikely chance alignment of two giant ellipticals at the same redshift and with very similar spectra, or ii) a gravitational lensing phenomenon. But in this second case, due to the lack of asymmetry in the two images, the lens could not be modelled with the standard lensing by a massive compact source. Actually, the usual gravitational lenses, i.e. those formed by bound clumps of matter, always produce inhomogeneous gravitational fields which distort the images of extended background sources (cf. Schneider et al., 1992; Kochanek et al., 2004). The detailed modelling of CSL-1 proved that the two images were virtually undistorted (see Paper I for details). The only other explanation left in the framework of the gravitational lensing theory was that of a lensing by a cosmic string.
In Paper I we indicated two possible experimenta crucis: i) the detection of the sharp edges at faint light levels, since this is the signature expected from the lensing by a cosmic string, and ii) the detection of a small amplitude but sharp discontinuity in the local CMB (Gangui et al., 2001). The first test calls for high angular resolution and deep observations with HST; time has already been allocated but the observations have not yet been performed. The second test was attempted using the preliminary release of the WMAP data (Kaiser & Stebbins, 1984; Lo & Wright, 2005). By careful processing, a $`1\sigma `$ positive detection at the position of CSL-1 was found. Even though such a detection seems to imply the unrealistic speed of $`0.94c`$ for the string, the authors pointed out that both the low angular resolution and the low S/N ratio could prevent the detection of the expected signature.
Another test was proposed by several authors (Vilenkin & Shellard, 1994; Hindmarsh, 1990; Bernardeu & Uzan, 2001; de Laix et al., 1996), and more recently in Huterer & Vachaspati (2003). It is based on the fact that the alignment of the background object (a galaxy) inside the deficit angle of the string is a stochastic process determined by the area of the lensing strip and by the surface density distribution of the extragalactic objects which are laying behind the string. All the lensed objects will fall inside a narrow strip defined by the deficit angle computed along the string pattern. A preliminary investigation of the CSL-1 field showed a significant excess of gravitational lens candidates selected on the bases of photometric criteria only (Sazhin et al., 2005). Spectroscopic observations are being obtained for a first set of candidates and will be discussed elsewhere.
In this paper we address on firmer grounds the issue of the gravitational lens nature of CSL-1 using intermediate resolution spectra obtained at the ESO Very Large Telescope + FORS 1 spectrograph.
## 2 The data and data reduction
New spectra of CSL-1 were obtained in March 2005 at VLT<sup>1</sup><sup>1</sup>1The Very Large Telescope is operated by the European Southern Observatory and is located at Mount Paranal in Chile; http://www.eso.org/paranal/ using the FORS 1 spectrograph under Director’s Discretionary Time (proposal 274.A-5039).
The spectra were acquired on March 15-19 2005 in the FORS 1 long slit configuration 600V+GG435 ($`\lambda /\mathrm{\Delta }\lambda =990`$ at central wavelength), setting the slit of the spectrograph across the centers of the two components of CSL-1. The observations were split in several exposures, to prevent saturation and to allow for a better removal of the bad pixels and cosmic rays. The risk of cross contamination between the spectra of the two components was minimized by retaining only the 6 frames with an average Signal to Noise (S/N) ratio of $`12`$ and a PSF FWHM$`<1.0\mathrm{"}`$<sup>2</sup><sup>2</sup>2The Point Spread Function represents the image of a pointlike object. Its Full Width at Half Maximum measures the level of blurring due to atmospheric and instrumental factors., for a total exposure time of 4740 sec.
During observations a short $`R`$-band exposure was also obtained to check the pointing of the instrument. Owing to the excellent seeing conditions, this image could be used to test the extended nature and the similarity of the light profiles of the two images. The dot-dashed line (with no data symbols) in Figure 2 shows the light profile of an unresolved star present in the field, compared to the light profiles of the two images of CSL-1, which appear to be clearly resolved and identical within the errors. A de Vaucouler fit yields $`r_e1.6\mathrm{"}`$.
The spectral data were reduced through standard MIDAS procedures and, after bias subtraction, flat fielding, wavelength calibration, and sky subtraction, the spectra were re-aligned to correct for dithering and stacked. Since the shift values were set to an integer number of pixels, no re–sampling was applied, in order not to affect the noise statistics. The stacking was performed using a simple median filter to reject cosmic rays.
The spectra of each component was extracted using a 5 pixel strip (1 pix=0.21”) centered on the emission peak with the purpose of maximizing the S/N ratio. In a similar way, the background counts were extracted in two stripes located 40 pixels from each component so to measure the local background while minimizing the contribution from the source. The error on the spectral counts was calculated with the following expression:
$$\sigma (ADU)=\sqrt{\sigma _{bkg}^2+\frac{N(ADU)}{n_{exp}\times g}}$$
where: $`N(ADU)`$ are the source counts measured along the spectrum, $`n_{exp}`$ are the number of median averaged exposures, $`g`$ is the instrumental gain, and $`\sigma _{bkg}`$ represents the background r.m.s. measured over $`5\times 20`$ pixels centered at each wavelength. By folding the images we estimate an average cross-contamination of $`7\%\pm 1\%`$. The resulting spectra and their ratio are shown in Figure 3. No flux calibration was applied to the data. It has to be stressed that the narrow spikes visible in the Figure are residuals left after the removal of bright sky lines.
The spectra of the two components turned out to be identical at visual inspection. In fact Pearson’s, Spearman, and Kendall correlation tests indicate correlation coefficients of $`0.96`$, $`0.94`$ and $`0.94`$ respectively, with a significance of $`>99.9\%`$ in all cases. The degree of similarity was further quantified by running a $`\chi ^2`$ test both on the whole wavelength range<sup>3</sup><sup>3</sup>3The regions affected by sky lines subtraction residuals were excluded from the test. and on the most prominent spectral absorption features, namely Ca II ($`H`$, $`K`$), H$`\beta `$, $`H\delta `$, $`H\gamma `$ lines and $`G`$ band. The test yields $`\chi _\nu ^2=1.03`$ implying that the two spectra are consistent within $`<2\sigma `$ (80%); we also compared the distribution of the observed differences to the frequencies expected in the case of pure gaussian noise, finding that the two are consistent at the 95% level and that there is no deterministic part in the residuals. An even better agreement ($`1\sigma `$) is found for the individual absorption features. To further check if the observed consistency can be due to the known similarity of early-type galaxies we repeated the test on a sample of spectra extracted from the SDSS<sup>4</sup><sup>4</sup>4Sloan Digital Sky Survey, data release 4: http://www.sdss.org Luminous Red Galaxies, chosen to have a redshift difference and S/N comparable with the CSL1 data. We performed 2000 comparisons obtaining that $`<2\%`$ of the examined SDSS spectra are as consistent as the spectra of the two CSL1 components.
A cross correlation test based on the spectral lines mentioned above yelds a velocity difference between the two components of $`\mathrm{\Delta }v=14\pm 30`$ Km s<sup>-1</sup>. If however we exclude the $`H\beta `$ line which is affected by residual instrumental effects this figure reduces to $`0\pm 20`$ Km s<sup>-1</sup>.
## 3 Conclusion
The similarity of the spectra of the two components of CSL-1 and the zero velocity shift between them strongly support the interpretation of CSL-1 in terms of gravitational lensing. The data obtained so far do not allow to completely rule out the possibility of a chance alignment of two giant ellipticals but the new data presented in this paper make such an alignment very unlikely. In the case of chance alignment, in fact, the two images of CSL-1 would correspond to two giant ellipticals with identical shapes and spectra, placed at the same redshift. The probability of finding two ellipticals of $`M_R=22.3`$ within 2” (20 kpc) and with a radial distance $`<1`$ Mpc ($`2\sigma `$ upper limit) is $`P1.5\times 10^{15}`$, accounting for clustering effects (e.g. Zehavi et al., 2005)<sup>5</sup><sup>5</sup>5$`P=_{V_1,V_2}N_{gal}^2[1+\xi (r)]𝑑V_1𝑑V_2`$ were $`N_{gal}`$ is the space density of elliptical galaxies, $`\xi (r)`$ is the galaxy correlation function, $`V_1`$ is the volume enclosing the two galaxies and $`V_2`$ is the volume of the survey.. Integrating over the volume sampled by the OAC-Deep Field for a galaxy of the same magnitude as CSL-1 we calculate that we expect to find $`9\times 10^4`$ pairs in the whole survey. Including the spectral similarity we obtain an upper limit of $`P<2\times 10^5`$. As it was already mentioned in Paper I, this could still be explained if CSL-1 belonged to a rich cluster with two central dominant galaxies, but this is not the case. Careful inspection of the CSL-1 field shows in fact that it is a rather isolated object with no other nearby galaxies of comparable brightness; furthermore the velocity difference measured from the two spectra is much smaller that the one expected in a rich cluster ($`\mathrm{\Delta }v300`$ Km s<sup>-1</sup>).
In the gravitational lensing scenario, as already stated in Paper I, the observed phenomenology cannot be understood in terms of lensing by compact clumps of matter such as, for instance, a Singular Isothermal Sphere model or any other model listed in the C.R. Keeton’s Lens Modeling Software (Kochanek et al., 2004). The only possible type of lens which can produce a morphology similar to that observed in CSL-1 seems to be a cosmic string (Sazhin et al., 2003).
Lensing by a cosmic string seems capable to explain all the observational evidences gathered so far and deserves further investigation. Cosmic strings were predicted by Kibble (1976) and their role in cosmology has been extensively discussed by Zeldovich (1980) and Vilenkin (1981). Recent work (Kibble, 2004; Polchinsky, 2004; Davis & Kibble, 2005) has also shown their relevance for both fundamental physics and cosmology. In particular it has become apparent that the detection of a cosmic string would lead to a direct measure of the energy scale of symmetry breaking in GUT theories. If we assume that CSL-1 is produced by a cosmic string, its measured properties would imply a linear density of the string of order of $`G\mu 410^7`$ and a corresponding energy scale of GUT of $`10^{15}`$ GeV (Kibble, 1976; Particle data group, 2004).
Hopefully the question on the nature of CSL1 will soon be answered by our HST observations approved in Cycle 14 to carry out the test proposed by the authors in Paper I and which will allow to verify the cosmic string hypothesis on firmer grounds.
## Acknowledgments
M.V.Sazhin acknowledges the INAF-Capodimonte Astronomical Observatory for hospitality and financial support. O. Khovanskaya acknowledges the Department of Physics of the University Federico II in Naples for financial support. The authors wish to thank C. Cezarsky, Director General of the European Southern Observatory for allocating Director’s Discretionary Time to this project. The work was also supported by the Russian Fund of Fundamental Investigations No. 04-02-17288. We thank the anonymous referees for the helpful suggestions.
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# A Brief Introduction To Operator Quantum Error Correction Dedicated to John Holbrook on the occasion of his 65th birthday
## 1. Introduction
In this paper we give an introduction to some of the mathematical aspects of quantum error correction, with an emphasis on the unified approach – called operator quantum error correction – recently introduced in . The field of quantum error correction took flight during the mid 1990’s . A central goal of this young field is to help construct quantum computers via the development of schemes that allow for the protection of quantum information against the noise associated with evolution of quantum systems. As it turns out, many of the problems in quantum error correction have an operator theoretic flavour. Here we shall briefly discuss the fundamental error correction protocols in quantum computing. We also describe various conditions that characterize correction in Operator QEC, and provide a new operator proof for the main testable condition in this scheme.
First let us briefly discuss the basic setting for quantum computation. See as examples of more extensive introductions. To each quantum system, the postulates of quantum mechanics associate a Hilbert space $`=^n`$. The finite dimensional case is the current focus in quantum computing for experimental reasons. A two-level quantum system is represented on $`=^2`$. This could describe, for instance, the ground and excited energy states of an electron in an atom. These are the two classical states that we observe, corresponding to an orthonormal basis $`\{|0,|1\}`$ for $`^2`$. However, quantum mechanics dictates that any linear combination $`|\psi `$ of these classical states is an allowable state, even though we only observe either the $`|0`$ or $`|1`$ state. When it is a non-trivial linear combination, $`|\psi `$ is said to be in a superposition of the classical states. A unit vector $`|\psi `$ is the fundamental quantum bit of information, also called a “qubit”. Equivalently, we could consider the rank one projection $`|\psi \psi |`$. The corresponding $`n`$-qubit composite system is realized on $`=(^2)^n=^{2^n}`$ with orthonormal basis $`\left\{|i_1\mathrm{}i_n|i_1\mathrm{}|i_n:i_j\{0,1\}\right\}`$ determined by the underlying two-level system.
More generally, we will only know that our system is in one of several states with various possibilities. So the direct generalization of a classical probability distribution in quantum information theory is a positive matrix $`\rho `$ with trace equal to one, a so-called density matrix.
A fundamental problem in quantum computation is to physically manipulate the superpositions inherent in quantum systems, without collapsing or “decohering” them. To accomplish this, methods must be developed to correct the errors that occur as quantum information is transferred from one physical location to the next inside, for instance, a quantum computer. To deal with this problem we must discuss evolution of quantum systems, a subject to which we now turn.
## 2. Evolution of Quantum Systems and Error Correction
The reversibility postulate of quantum mechanics implies that evolution in a closed quantum system occurs via unitary maps. From the discrete perspective, if we take a snapshot of this evolution, then a density matrix $`\rho `$ will encode the possible states of the system with various probabilities at a given time. An evolution of the system corresponds to a map $`\rho U\rho U^{}`$ for some unitary operator $`U`$.
In the context of quantum computing, the quantum systems of interest are “open” as they are exposed to external environments during computations. In such cases, the open system is regarded as part of a larger closed quantum system given by the composite of the system and the environment. If $`_S`$ and $`_E`$ are the system and environment Hilbert spaces, then the closed system is represented on $`=_E_S`$.
The characterization of evolution in open quantum systems requires first that density operators are mapped to density operators; i.e. probability densities are mapped to probability densities. Thus, such a map must be positive and trace preserving. However, this property must be preserved when the system is exposed to all possible environments. In terms of the map, if $``$ describes an evolution of the system, then the map $`id_E:(_E_S)(_E_S)`$ must also be positive and trace preserving for all $`E`$. Hence, the widely accepted working definition of a quantum operation (or evolution, or channel) on a Hilbert space $``$, is a completely positive, trace preserving map $``$ on $`()`$ (CPTP for short).
Deriving from a theorem of Choi and Kraus , every CPTP map $`:()()`$ has an “operator-sum representation” of the form $`(\rho )=_aE_a\rho E_a^{}`$ for some set of (non-unique) operators $`\{E_a\}()`$ with $`_aE_a^{}E_a=1\mathrm{l}`$. The $`E_a`$ are called the noise operators or errors associated with $``$. In the context of quantum error correction, it is precisely the effects of these errors that must be mitigated. As a short hand, we write $`=\{E_a\}`$ when an error model for $``$ is known.
Error correction in quantum computing is a much more delicate problem in comparison to its classical counterpart. As a simple observation, consider that the only errors that occur classically are some version of bit flips; e.g., $`|0`$ goes to $`|1`$ or vice-versa. More generally, in quantum computing subtleties arise from the fact that a given qubit can be corrupted to an infinite number of possible superpositions. In terms of operators on single bits or qubits for instance, whereas the Pauli bit flip matrix $`X=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ is the fundamental classical error matrix, any unitary matrix is a possible error in quantum computing. Of course, there are many other issues, such as the fabled “No Cloning Theorem”. The linearity of quantum mechanics implies that the analogue of the classically well-used repetition code does not extend to arbitrary qubits $`|\psi |\psi |\psi `$. Fortunately, methods have been, and are being, developed to overcome these challenges.
## 3. Standard Model of Quantum Error Correction
The “Standard Model” of quantum error correction involves triples $`(,,𝒞)`$ where $`𝒞`$ is a subspace, a quantum code, of a Hilbert space $``$ associated with a given quantum system, and the error $``$ and recovery $``$ are quantum operations on $`()`$.
Recall from the discussions above that we are forced by quantum mechanics to consider subspaces $`𝒞`$ as sets of codes, as linear combinations of classical codewords are perfectly allowable codewords in this setting. In the trivial case, when $`=\{U\}`$ is implemented by a single unitary error operator, the recovery is just the reversal operation $`=\{U^{}\}`$; that is,
$$\rho \stackrel{}{}U\rho U^{}\stackrel{}{}U^{}(U\rho U^{})U=\rho .$$
Of course, here there is no need to restrict the input operators $`\rho `$.
More generally, the set $`(,,𝒞)`$ forms an “error triple” if $``$ undoes the effects of $``$ on $`𝒞`$ in the following sense:
(1) $`()(\sigma )=\sigma \sigma (𝒞),`$
where $`𝒞`$ is naturally regarded as embedded inside $``$.
When there exists such an $``$ for a given pair $`,𝒞`$, the subspace $`𝒞`$ is said to be correctable for $``$. The existence of a recovery operation $``$ of $`=\{E_a\}`$ on $`𝒞`$ is characterized by the following condition : There exists a scalar matrix $`\mathrm{\Lambda }=(\lambda _{ab})`$ such that
(2)
$$P_𝒞E_a^{}E_bP_𝒞=\lambda _{ab}P_𝒞a,b,$$
where $`P_𝒞`$ is the projection of $``$ onto $`𝒞`$. It is not hard to see that this condition is independent of the operator-sum representation for $``$. We note that Eq. (2) is a special case of Eq. (10) below.
The motivating case of an error model that satisfies Eq. (2) occurs when the restrictions $`E_a|_{P_𝒞}=E_a|_𝒞`$ of the noise operators to $`𝒞`$ are scalar multiples of unitary operators $`U_a`$, such that the subspaces $`U_a𝒞`$ are mutually orthogonal. In this situation the positive scalar matrix $`\mathrm{\Lambda }`$ is diagonal. A correction operation here may be constructed by an application of the measurement operation determined by the subspaces $`U_a𝒞`$, followed by the reversals of the corresponding restricted unitaries $`U_aP_𝒞`$. Specifically, if $`P_a`$ is the projection of $``$ onto $`U_a𝒞`$, then $`=\{U_a^{}P_a\}`$ satisfies Eq. (1) for $``$ on $`𝒞`$.
Let us discuss a simple example. Let $`𝒞`$ be the subspace of $`=^8`$ given by $`𝒞=\mathrm{span}\{|000,|111\}`$. Let $`=\{\frac{1}{\sqrt{3}}X_k:k=1,2,3\}`$ with the Pauli matrix $`X`$ and $`X_1=X1\mathrm{l}_21\mathrm{l}_2`$, and similarly for $`X_2`$, $`X_3`$. In this case, $`\mathrm{\Lambda }=\frac{1}{3}1\mathrm{l}_3`$. The correction operation $``$ may be constructed as above.
## 4. Noiseless Subsystems
To describe the notion of noiseless subsystems from , we begin with a decomposition of the system Hilbert space
$$=\underset{J}{}_J^A_J^B,$$
where the “noisy subsystems” $`_J^A`$ have dimension $`m_J`$ and the “noiseless subsystems” $`_J^B`$ have dimension $`n_J`$. We focus on the case where information is encoded in a single noiseless sector of $`()`$, so
$$=(^A^B)𝒦$$
with $`dim(^A)=m`$, $`dim(^B)=n`$ and $`dim𝒦=dimmn`$. We shall write $`\sigma ^A`$ for operators in $`(^A)`$ and $`\sigma ^B`$ for operators in $`(^B)`$.
Let $`\{|\alpha _k:1km\}`$ be an orthonormal basis for $`^A`$ and let
$$\{P_{kl}=|\alpha _k\alpha _l|1\mathrm{l}_n:1k,lm\}$$
be the corresponding family of matrix units in $`(^A)1\mathrm{l}^B`$. Recall that the partial trace over $`A`$ on $`^A^B`$ is the quantum operation defined on elementary tensors by $`\mathrm{Tr}_A(\sigma ^A\sigma ^B)=\mathrm{Tr}(\sigma ^A)\sigma ^B`$.
Define for a fixed decomposition $`=(^A^B)𝒦`$ the operator semigroup
(3) $`𝔄=\{\sigma ():\sigma =\sigma ^A\sigma ^B,\mathrm{for}\mathrm{some}\sigma ^A\mathrm{and}\sigma ^B\}.`$
For notational purposes, we assume that bases have been chosen and define the matrix units $`P_{kl}`$ as above, so that $`P_k=P_{kk}`$, $`P_𝔄P_1+\mathrm{}+P_m`$ and $`P_𝔄=^A^B`$. We also define a map $`𝒫_𝔄`$ by the action $`𝒫_𝔄()=P_𝔄()P_𝔄`$. The following result motivates the (generalized) definition of NS’s from . (See for a proof.)
###### Lemma 4.1.
Given a fixed decomposition $`=(^A^B)𝒦`$ and a quantum operation $``$ on $`()`$, the following three conditions are equivalent:
1. $`\sigma ^A\sigma ^B,\tau ^A:(\sigma ^A\sigma ^B)=\tau ^A\sigma ^B`$
2. $`\sigma ^B,\tau ^A:(1\mathrm{l}^A\sigma ^B)=\tau ^A\sigma ^B`$
3. $`\sigma 𝔄:\left(\mathrm{Tr}_A𝒫_𝔄\right)(\sigma )=\mathrm{Tr}_A(\sigma )`$.
###### Definition 4.2.
The $`^B`$ sector of the semigroup $`𝔄`$ encodes a noiseless subsystem for $``$ when it satisfies the equivalent conditions of Lemma 4.1.
The NS framework discussed here is a generalization of both the “Standard NS” and “Decoherence-Free Subspace” methods of passive error correction, both of which are used for unital quantum operations. The method described here applies to all CPTP maps. See for more discussions on this point.
As a simple example of how such subsystems naturally arise, let $`\mathrm{\Phi }:(^A)(^A)`$ be an arbitrary CPTP map and let $`\mathrm{\Psi }:(^B)(^B)`$ be CPTP with a Standard NS $`_0^B^B`$; i.e., $`\mathrm{\Psi }(\rho )=\rho `$ for all $`\rho (_0^B)`$. Then $`_0^B`$ encodes a noiseless subsystem inside $`^A^B`$ for the map $`=\mathrm{\Phi }\mathrm{\Psi }:(^A^B)(^A^B)`$.
To be of use in practical applications, we need testable conditions for a map $`=\{E_a\}`$ to admit a NS described by a semigroup $`𝔄`$. Towards this end, we have proved the following theorem.
###### Theorem 4.3.
Let $`=\{E_a\}`$ be a quantum operation on $`()`$ and let $`𝔄`$ be a semigroup in $`()`$ as above. Then the following three conditions are equivalent:
* The $`^B`$ sector of $`𝔄`$ encodes a noiseless subsystem for $``$.
* The subspace $`P_𝔄=^A^B`$ is invariant for the operators $`E_a`$ and the restrictions $`E_a|_{P_𝔄}`$ belong to the algebra $`(^A)1\mathrm{l}^B`$.
* The following two conditions hold:
(4)
$$P_kE_aP_l=\lambda _{akl}P_{kl}a,k,l$$
for some set of scalars $`(\lambda _{akl})`$ and
(5)
$$E_aP_𝔄=P_𝔄E_aP_𝔄a.$$
Proof. Since the matrix units $`\{P_{kl}\}`$ generate $`(^A)1\mathrm{l}^B`$, it follows that (3) is a restatement of (2). Here we sketch the proof of the equivalence of (1) and (3), see for details. To prove the necessity of Eqs. (4), (5) for (1), it follows from properties of the map $`\mathrm{\Gamma }=\{P_{kl}\}`$ and Lemma 4.1 that there exist scalars $`\mu _{kiajl,k^{}l^{}}`$ such that
(6)
$$P_{ki}E_aP_{jl}=\underset{k^{}l^{}}{}\mu _{kiajl,k^{}l^{}}P_{k^{}l^{}}.$$
Multiplying both sides of this equality on the right by $`P_l`$ and on the left by $`P_k`$, we see that $`\mu _{kiajl,k^{}l^{}}=0`$ when $`kk^{}`$ or $`ll^{}`$. This implies Eq. (4) with $`\lambda _{akl}=\mu _{kkall,kl}`$. Equation (5) follows from Lemma 4.1 and consideration of the operator-sum representation for $``$.
On the other hand, if Eqs. (4), (5) hold, then for all $`\sigma =P_𝔄\sigma 𝔄`$ we have
$$(\sigma )=\underset{a,k,k^{}}{}P_kE_a\sigma E_a^{}P_k^{}.$$
This implies that for all $`\sigma =\sigma ^A\sigma ^B𝔄`$,
$`(\sigma ^A\sigma ^B)`$ $`=`$ $`{\displaystyle \underset{a,k,k^{},l,l^{}}{}}P_kE_aP_l(\sigma ^A\sigma ^B)P_l^{}E_a^{}P_k^{}`$
$`=`$ $`{\displaystyle \underset{a,k,k^{},l,l^{}}{}}\lambda _{akl}\overline{\lambda }_{ak^{}l^{}}P_{kl}(\sigma ^A\sigma ^B)P_{l^{}k^{}}.`$
Condition (1) now follows from the fact that the matrix units $`P_{kl}`$ act trivially on the $`(^B)`$ sector. $`\mathrm{}`$
## 5. Operator Quantum Error Correction
The Operator QEC approach consists of triples $`(,,𝔄)`$ where $``$ and $``$ are quantum operations on some $`()`$, and $`𝔄`$ is a semigroup in $`()`$ defined as above with respect to a fixed decomposition $`=(^A^B)𝒦`$.
###### Definition 5.1.
Given a triple $`(,,𝔄)`$ we say that the $`^B`$ sector of $`𝔄`$ is correctable for $``$ if
(7) $`\left(\mathrm{Tr}_A𝒫_𝔄\right)(\sigma )=\mathrm{Tr}_A(\sigma )\text{for all}\sigma 𝔄.`$
Equivalently, $`(,,𝔄)`$ is a correctable triple if the $`^B`$ sector of the semigroup $`𝔄`$ encodes a noiseless subsystem for the error map $``$. Table 1 indicates the special cases captured by Operator QEC. Our choice of terminology here is motivated by the fact that correctable codes in this scheme take the form of operator algebras and operator semigroups. We point the reader to for examples of error triples on subsystems that require non-trivial recovery operations, and for other recent related work.
An important feature of Operator QEC is that a semigroup $`𝔄`$ is correctable exactly when the $`\mathrm{C}^{}`$-algebra $`𝔄_0=1\mathrm{l}^A(^B)`$ can be corrected precisely.
###### Theorem 5.2.
Let $`=\{E_a\}`$ be a quantum operation on $`()`$ and let $`𝔄`$ be a semigroup in $`()`$ as above. Then the $`^B`$ sector of $`𝔄`$ is correctable for $``$ if and only if there is a quantum operation $``$ on $`()`$ such that
(8) $`()(\sigma )=\sigma \sigma 1\mathrm{l}^A(^B).`$
Proof. If Eq. (8) holds, then condition (2) of Lemma 4.1 holds for $``$ with $`\tau ^A=1\mathrm{l}^A`$ and hence $`𝔄`$ is correctable for $``$. On the other hand, suppose that $`𝔄`$ is correctable for $``$ and condition (2) of Lemma 4.1 holds for $``$. Note that the map $`\mathrm{\Gamma }^{}=\{\frac{1}{\sqrt{m}}P_{kl}\}`$ is trace preserving on $`(^A^B)`$. Thus, from basic properties of the map $`\mathrm{\Gamma }=\{P_{kl}\}`$, we have for all $`\sigma ^B`$,
(9) $`(\mathrm{\Gamma }^{})(1\mathrm{l}^A\sigma ^B)=\mathrm{\Gamma }^{}(\tau ^A\sigma ^B)1\mathrm{l}^A\sigma ^B.`$
By trace preservation the proportionality factor must be one, and hence Eq. (8) is satisfied for $`(\mathrm{\Gamma }^{})`$. The map $`\mathrm{\Gamma }^{}`$ may be extended to a quantum operation on $`()`$ by including the projection $`P_𝔄^{}`$ onto $`𝒦`$ as a noise operator. As this does not effect the calculation (9), the result follows. $`\mathrm{}`$
We next give a testable condition, Eq. (10), that characterizes correction in the Operator QEC regime. Notice that this is a generalization of Eq. (2) for Standard QEC. This condition was introduced in and necessity was established. Sufficiency was proved in up to a set of technical conditions, and more recently in with full generality. (The work of also links this condition with an interesting information theoretic condition.) Here we include a sketch of the proof of necessity from , and a new operator theoretic version of the proof of sufficiency sketched in . We assume that matrix units $`\{P_{kl}=|\alpha _k\alpha _l|1\mathrm{l}^B\}`$ inside $`(^A)1\mathrm{l}^B`$ have been chosen as above.
###### Theorem 5.3.
Let $`=\{E_a\}`$ be a quantum operation on $`()`$ and let $`𝔄`$ be a semigroup in $`()`$ as above. For the $`^B`$ sector of $`𝔄`$ to be correctable for $``$, it is necessary and sufficient that there are scalars $`\mathrm{\Lambda }=(\lambda _{abkl})`$ such that
(10) $`P_kE_a^{}E_bP_l=\lambda _{abkl}P_{kl}a,b,k,l.`$
Proof. For necessity, note first that Theorem 5.2 gives us a CPTP map $``$ on $`()`$ such that $``$ acts as the identity channel on $`𝔄_0=1\mathrm{l}^A(^B)()`$.
Suppose that $`=\{R_b\}`$. The noise operators for the operation $``$ are $`\{R_bE_a\}`$, and using arguments similar to those of Theorem 4.3 (see for details) we may find scalars $`\mu _{abkl}`$ such that
$$P_kR_bE_aP_l=\mu _{abkl}P_{kl}a,b,k,l.$$
Consider the products
$`\left(P_kR_bE_aP_l\right)^{}\left(P_k^{}R_bE_a^{}P_l^{}\right)`$ $`=`$ $`\left(\overline{\mu _{abkl}}P_{lk}\right)\left(\mu _{a^{}bk^{}l^{}}P_{k^{}l^{}}\right)`$
$`=`$ $`\{\begin{array}{cc}(\overline{\mu _{abkl}}\mu _{a^{}bkl^{}})P_{ll^{}}& \text{if }k=k^{}\hfill \\ 0& \text{if }kk^{}\hfill \end{array}.`$
Now, the subspace $`𝒞`$ can be shown to be invariant for the noise operators $`R_bE_a`$. Hence for fixed $`a,a^{}`$ and $`l,l^{}`$ we use $`_bR_b^{}R_b=1\mathrm{l}`$ to obtain
$`\left({\displaystyle \underset{b,k}{}}\overline{\mu _{abkl}}\mu _{a^{}bkl^{}}\right)P_{ll^{}}`$ $`=`$ $`{\displaystyle \underset{b,k}{}}\left(P_lE_a^{}R_b^{}P_k\right)\left(P_kR_bE_a^{}P_l^{}\right)`$
$`=`$ $`{\displaystyle \underset{b}{}}P_lE_a^{}R_b^{}P_𝔄R_bE_a^{}P_l^{}`$
$`=`$ $`P_lE_a^{}\left({\displaystyle \underset{b}{}}R_b^{}R_b\right)E_a^{}P_l^{}`$
$`=`$ $`P_lE_a^{}E_a^{}P_l^{}`$
The proof of necessity is completed by setting $`\lambda _{aa^{}ll^{}}=_{b,k}\overline{\mu _{abkl}}\mu _{a^{}bkl^{}}`$ for all $`a,a^{}`$ and $`l,l^{}`$.
For sufficiency, let us assume that Eq. (10) holds. Let $`\sigma _k=|\alpha _k\alpha _k|(^A)`$, for $`1km`$, and define a CPTP map $`_k:(^B)()`$ by $`_k(\rho ^B)(\sigma _k\rho ^B).`$ With $`PP_𝔄`$ and $`E_{a,k}E_aP|\alpha _k`$, it follows that $`_k=\{E_{a,k}\}`$. We shall find a CPTP map that globally corrects all of the errors $`E_{a,k}`$.
To do this, first note that we may define a CPTP map $`_B:(^B)()`$ with error model
$$_B=\{\frac{1}{\sqrt{m}}E_{a,k}:a,1km\}.$$
Then Eq. (10) and $`P=_kP_k`$ give us
$`1\mathrm{l}^BE_{a,k}^{}E_{b,l}1\mathrm{l}^B`$ $`=`$ $`1\mathrm{l}^B\alpha _k|PE_a^{}E_bP|\alpha _l1\mathrm{l}^B`$
$`=`$ $`{\displaystyle \underset{k^{},l^{}}{}}1\mathrm{l}^B\alpha _k|P_k^{}E_a^{}E_bP_l^{}|\alpha _l1\mathrm{l}^B`$
$`=`$ $`{\displaystyle \underset{k^{},l^{}}{}}\lambda _{abk^{}l^{}}\mathrm{\hspace{0.17em}1}\mathrm{l}^B\alpha _k|P_{k^{}l^{}}|\alpha _l1\mathrm{l}^B`$
$`=`$ $`\lambda _{abkl}1\mathrm{l}^B.`$
In particular, Standard QEC implies the existence of a CPTP map $`:()(^B)`$ such that $`(_B)(\rho ^B)=\rho ^B`$ for all $`\rho ^B`$.
This implies that
$`()(1\mathrm{l}^A\rho ^B)`$ $`=`$ $`\left({\displaystyle \underset{k}{}}_k(\rho ^B)\right)`$
$`=`$ $`m\left({\displaystyle \underset{k,a}{}}{\displaystyle \frac{1}{m}}E_{a,k}\rho ^BE_{a,k}^{}\right)`$
$`=`$ $`m_B(\rho ^B)=m\rho ^B.`$
Hence we may define a CPTP ampliation map $`I_𝔄:(^B)()`$ via $`I_𝔄(\rho ^B)=\frac{1}{m}(1\mathrm{l}^A\rho ^B`$). Thus on defining $`^{}I_𝔄`$, we obtain
$$\left(^{}\right)(1\mathrm{l}^A\rho ^B)=1\mathrm{l}^A\rho ^B\rho ^B(^B).$$
The result now follows from an application of Theorem 5.2. $`\mathrm{}`$
## 6. Concluding Remark
The focus of research in quantum error correction has mainly been on finite dimensional problems to this point. Primarily this reflects the current status of experimental efforts to build quantum computers, and the fact that many scientists working in the area are closely linked with experimentalists. Thus, in the author’s opinion, there is an opportunity here for operator theorists. In particular, mathematicians working in the field have, for the most part, not had the luxury of exploring infinite dimensional aspects and extensions of the quantum error correction framework. It is expected that problems of this nature will eventually be of experimental interest, and we expect they would be of current mathematical interest.
Acknowledgements. This paper was prepared as part of the Proceedings of the 25th Anniversary Meeting of the Great Plains Operator Theory Symposium held at the University of Central Florida in June 2005. We thank the organizers and participants for a stimulating meeting. Thanks also to Palle Jorgensen for helpful comments on an early draft of this paper. This work was partially supported by NSERC.
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# Correspondence principle for the diffusive dynamics of a quartic oscillator: deterministic aspects and the role of temperature
## I Introduction
A widely accepted explanation for the emergence of classical behavior from the quantum substratum is provided by the environment-induced decoherence (EID) program zurek03 ; omnes ; joos ; isar : the interaction of the system with many uncontrollable external degrees of freedom (the environment) yields a dramatic destruction of quantum coherences prohibiting, as a consequence, the existence of quantum superpositions in the macroscopic world. Decoherence, an experimentally demonstrated effect brune ; cheng ; myatt to which real systems are inevitably submitted, is the essence of the arguments given by EID defenders to explain, for instance, why an astronomical body behaves classically zurek98 ; zurek96 and to address fundamental questions, as for instance the Schröedinger cat paradox and the measurement problem omnesB .
Recent works, however, have discussed the real effectiveness of decoherence as a necessary mechanism to guarantee the quantum-classical transition. In reference wiebe , e.g., the authors show that a small amount of smoothing (due to apparatus resolution) is sufficient to ensure the classical limit. In adelcio , it is argued that diffusive decoherence can produce only an attenuation of quantum coherences, not their complete destruction. The authors thus conclude that the classical limit is a matter of experimental resolution. Another important objection raised in references wiebe ; ballentine points out to the fact that decoherence is not able to reduce the spreading of the wave function. Yet, in reference petitjean relevant questions concerning the role of the high-temperature limit and the generality of environment models have been addressed. Such a present debate forbids one to believe that the understanding of the quantum-classical transition is an exhausted subject.
In this work the correspondence principle is discussed by means of an analytical analysis carried out for the the dynamics of a nonlinear oscillator coupled with a phase-damping reservoir composed by $`N`$ harmonic oscillators in thermal equilibrium. Such a thermal bath has been shown to be effective in modelling nondissipative decoherence for a long range of finite temperatures angelo06 . With the help of this model, the quantum-classical correspondence is analyzed in two main directions, as indicated below.
(i) On one hand, it is investigated the correspondence between the quantum centroid and the Newtonian trajectory. This analysis is motivated by Einstein’s paper to Born about the foundations of quantum mechanics einstein . Einstein concludes that the only possible interpretation for the wave function is the one based on statistical ensembles but, on the other hand, emphasizes the “inevitably conception that Physics must ferment a realistic description of only one system”. Indeed, “nature as a whole may be thought as an individual system (existing only one time, with no need for repetitions) and not as a ‘ensemble of systems’.”einstein . Einstein’s dissatisfaction is invoked here in terms of the following motivating question: what can quantum mechanics predict about the future of an individual system? Obviously, such a question intends to challenge the quantum formalism to realize a task which is supposed to be an exclusiveness of Newtonian mechanics. However, since physics is fundamentally concerned with experimental observations, which in turn are inevitably limited by the apparatus resolution, it is reasonable to conceive that neither the idea of an objective reality (usually associated with classical beliefs) nor a completely indeterministic scenario (usually associated with quantum world) can be maintained. In fact, as it is argued in Sec.IV, once a margin of error is defined the concept of deterministic behavior can be physically accepted, at least for the short-time dynamics of initially localized wave functions. Moreover, it is shown that within such a short-time regime (the Ehrenfest regime) decoherence is not effective, in agreement with the claims of ref.wiebe ; ballentine .
(ii) On the other hand, beyond this time scale within which deterministic behavior is expected to be experimentally verifiable, predictions for individual systems can be given only in terms of statistical averages and variances. In this case, the quantum-classical correspondence is investigated in terms of comparisons between quantum and Liouvillian results, for which the experimental resolution is also regarded. In Sec.V, it is shown that decoherence is indeed a necessary mechanism to promote the quantum-classical correspondence. Our results show in addition that the reestablishment of the quantum-classical correspondence can be ensured by purely diffusive decoherence only within the regime of high temperatures.
## II The model
Recently, a model of reservoir has been proposed which is able to describe nondissipative decoherence at finite temperatures angelo06 . Here, it is investigate the dynamics of a system composed by such a reservoir coupled to the well known milburn ; lili ; adelcio quartic oscillator. The Hamiltonian of the model reads
$`\widehat{H}`$ $`=`$ $`\widehat{H}_S+\widehat{H}_R+\widehat{H}_I,`$ (1a)
where
$`\widehat{H}_S`$ $`=`$ $`\mathrm{}\omega _s\left(\widehat{n}_s+{\displaystyle \frac{1}{2}}\right)+\mathrm{}^2g_s\left(\widehat{n}_s+{\displaystyle \frac{1}{2}}\right)^2,`$ (1b)
$`\widehat{H}_R`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}\mathrm{}\omega _k\left(\widehat{n}_k+{\displaystyle \frac{1}{2}}\right),`$ (1c)
$`\widehat{H}_I`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}\mathrm{}^2g_k\left(\widehat{n}_s+{\displaystyle \frac{1}{2}}\right)\left(\widehat{n}_k+{\displaystyle \frac{1}{2}}\right).`$ (1d)
The constants $`g_s`$ and $`g_k`$ denote, respectively, the nonlinearity parameter and the coupling parameter, whereas $`\omega _s`$ and $`\omega _k`$ stand for the harmonic frequency of the nonlinear oscillator and the harmonic frequencies of the bath oscillators. The initial joint quantum state is assumed to be
$`\widehat{\rho }(0)=|\alpha _0\alpha _0|{\displaystyle \frac{e^{\beta \widehat{H}_R}}{\text{Tr}\left[e^{\beta \widehat{H}_R}\right]}},`$ (2)
where $`\beta =(k_BT)^1`$ and $`T`$ is the equilibrium temperature. The assumptions of separability (for the initial joint state) and purity (for the system state) are both convenient for a suitable construction of the respective classical distribution. The corresponding classical dynamics is described in terms of the Hamiltonian
$`H=H_S+H_R+H_I,`$ (3a)
where
$`H_S`$ $`=`$ $`\omega _s\left({\displaystyle \frac{p_s^2+q_s^2}{2}}\right)+g_s\left({\displaystyle \frac{p_s^2+q_s^2}{2}}\right)^2,`$ (3b)
$`H_R`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}\omega _k\left({\displaystyle \frac{p_k^2+q_k^2}{2}}\right),`$ (3c)
$`H_I`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}g_k\left({\displaystyle \frac{p_s^2+q_s^2}{2}}\right)\left({\displaystyle \frac{p_k^2+q_k^2}{2}}\right).`$ (3d)
For convenience, new canonical variables $`q_k=X_k\sqrt{m_k\omega _k}`$ and $`p_k=P_k/\sqrt{m_k\omega _k}`$ have been used instead of the usual position-momentum pair $`(X_k,P_k)`$. The quantum and classical Hamiltonians are correctly related by the scheme of bosonic quantization proposed in ref.angelo03 . A Gaussian spectral distribution given by
$`g_k={\displaystyle \frac{\mathrm{\Omega }}{\sqrt{N}}}\mathrm{exp}\left[{\displaystyle \frac{\pi \left(kk_0\right)^2}{2N^2}}\right],`$ (4)
where $`k_0=N/2`$ and $`\mathrm{\Omega }`$ is constant, will be assumed in our calculations in the next sections. As has been shown in ref.angelo06 it allows for the appropriated implementation of the thermodynamic limit ($`N\mathrm{}`$). For simplicity and convenience, the assumption $`\omega _k=\omega _0`$ will be also adopted here. As discussed in ref.angelo06 , although restrictive, this choice does not change the main properties of the model, namely, the Gaussian decay for the quantum coherences and the structural form of the decoherence time.
For the Liouvillian dynamics the following phase space distribution is used:
$`\rho (0)={\displaystyle \frac{e^{\frac{(q_sq_0)^2}{\mathrm{}}}}{\sqrt{\pi \mathrm{}}}}{\displaystyle \frac{e^{\frac{(p_sp_0)^2}{\mathrm{}}}}{\sqrt{\pi \mathrm{}}}}{\displaystyle \underset{k=1}{\overset{N}{}}}\left({\displaystyle \frac{e^{\frac{p_k^2+q_k^2}{z\mathrm{}}}}{z\pi \mathrm{}}}\right),`$ (5)
where
$`z^1=\mathrm{tanh}\left({\displaystyle \frac{\beta \mathrm{}\omega }{2}}\right).`$ (6)
This distribution corresponds to the Wigner function of $`\widehat{\rho }(0)`$ given by (2). Note that the Wigner function for the reservoir state reduces to the traditional Boltzmanian function $`e^{\beta H_R}/\text{Tr}\left(e^{\beta H_R}\right)`$ only in the high-temperature regime $`(\beta \mathrm{}\omega 1)`$.
There are several reasons supporting the choice of this model, namely: (i) as it has been shown in angelo06 , the phase-damping reservoir induces nondissipative decoherence at finite temperature, this being one important difference from other purely diffusive environments adelcio ; toscano ; (ii) the nonlinear oscillator used as the system of interest constitutes a paradigm of nontrivial quantum dynamics which exhibits interference and cat formation milburn ; lili ; (iii) classical and quantum solutions are analytical and exact for all regimes of parameters, with no need for either Markovian or weak coupling assumptions.
## III Analytical results
### III.1 Newtonian results
The Newtonian solution is obtained by integrating the system of $`2(N+1)`$ differential equations defined by Hamilton’s equations, $`\dot{q}_k=_{p_k}H`$ and $`\dot{p}_k=_{q_k}H`$. The final result for the Newtonian trajectory of the system can be conveniently written in matrix notation as follows:
$`R_N(\tau )=𝕄_1\left[\varphi _N\right]R_0,`$ (7a)
where
$`\varphi _N(\tau )`$ $``$ $`{\displaystyle \frac{\omega _s\tau }{\mathrm{}g_s}}+\left({\displaystyle \frac{R_0^\text{T}R_0}{\mathrm{}}}\right)\tau +{\displaystyle \underset{k}{}}{\displaystyle \frac{g_k}{g_s}}\left({\displaystyle \frac{p_k^2+q_k^2}{\mathrm{}}}\right)\tau ,`$
$`𝕄_1[\chi ]`$ $``$ $`\left[\begin{array}{cc}\mathrm{cos}\chi & \mathrm{sin}\chi \\ & \\ \mathrm{sin}\chi & \mathrm{cos}\chi \end{array}\right],`$ (7f)
where $`R_N(\tau )\left(\genfrac{}{}{0pt}{}{q_s(\tau )}{p_s(\tau )}\right)`$, $`R_0\left(\genfrac{}{}{0pt}{}{q_0}{p_0}\right)`$, and $`R_0^\text{T}R_0=p_0^2+q_0^2`$. For convenience, all results have been written in terms of the following dimensionless parameter:
$`\tau =\mathrm{}g_st.`$ (8)
Note that the Newtonian solution does not depend at all on $`\mathrm{}`$. The effect of the nonlinearity becomes evident in equations above: the rotation matrix $`𝕄[\varphi _N]`$ becomes dependent on the initial condition $`R_0`$. This fact can be interpreted as the basic mechanism responsible for the twist found in the dynamics of classical distributions milburn ; adelcio and for the dynamical formation of cat states (quantum superposition of several coherent states), as shown in ref.lili .
### III.2 Quantum results
The analytical solution for the reduced density matrix associated with the system of interest is obtained by taking the trace of the global density matrix, $`\widehat{\rho }(t)=e^{ı\widehat{H}t/\mathrm{}}\widehat{\rho }(0)e^{ı\widehat{H}t/\mathrm{}}`$, over the reservoir degrees of freedom, i.e., $`\widehat{\rho }_S(t)=\text{Tr}_R\widehat{\rho }(t)`$. As shown in ref.angelo06 , the result can be written in Fock basis as
$`\widehat{\rho }_S(\tau )`$ $`={\displaystyle \underset{n,n^{}=0}{\overset{\mathrm{}}{}}}e^{|\alpha _0|^2}{\displaystyle \frac{\alpha _0^n}{\sqrt{n!}}}{\displaystyle \frac{(\alpha _0^{})^n^{}}{\sqrt{n^{}!}}}e^{ı(n^2n^2)\tau }\times `$ (9)
$`e^{ı(nn^{})\left(1+\frac{\omega _s}{\mathrm{}g_s}+\frac{G_1}{2g_s}\right)\tau }C_{nn^{}}(\tau )|nn^{}|,`$
where $`G_1_kg_k`$ and
$`C_{nn^{}}(\tau ){\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \frac{1e^{\beta \mathrm{}\omega _k}}{1e^{\beta \mathrm{}\omega _k}e^{ı(g_k/g_s)(nn^{})\tau }}}.`$ (10)
From (10) we see that
$`\left|C_{nn^{}}\right|={\displaystyle \underset{k}{}}\left\{1+{\displaystyle \frac{\mathrm{sin}^2\left[\frac{g_k}{g_s}\frac{(nn^{})\tau }{2}\right]}{\mathrm{sinh}^2\left[\frac{\beta \mathrm{}\omega _k}{2}\right]}}\right\}^{\frac{1}{2}}.`$ (11)
For short times $`(\tau 1)`$ this function behaves as
$`\left|C_{nn^{}}\right|e^{(nn^{})^2\tau ^2/\tau _{DQ}^2},`$ (12)
where $`\tau _{DQ}`$, which denotes the characteristic decoherence time of the phase-damping reservoir, is given by
$`\tau _{DQ}=\sqrt{8}{\displaystyle \frac{g_s}{G_2}}\mathrm{sinh}\left({\displaystyle \frac{\beta \mathrm{}\omega }{2}}\right),`$ (13)
with $`G_2^2_kg_k^2`$. (In the thermodynamic limit $`G_20.89\mathrm{\Omega }`$.) Result (12) shows the character of Gaussian decay induced by the phase-damping reservoir over the off-diagonal terms ($`nn^{}0`$) of the density operator given by (9). Hereafter the assumption $`\omega _k=\omega `$, as suggested in angelo06 , will be implemented for simplicity.
The column matrix $`R_Q(\tau )\left(\genfrac{}{}{0pt}{}{\widehat{q}_s}{\widehat{p}_s}\right)`$, defined in terms of the expectation values of the position and the momentum, is calculated with the help of the usual relations: $`\widehat{q}_s=\sqrt{\mathrm{}/2}(\widehat{a}_s^{}+\widehat{a}_s)`$ and $`\widehat{p}_s=ı\sqrt{\mathrm{}/2}(\widehat{a}_s^{}\widehat{a}_s)`$, where $`\widehat{a}_s=\text{Tr}[\widehat{a}_s\widehat{\rho }_S(\tau )]`$ and $`\widehat{a}_s^{}=\text{Tr}[\widehat{a}_s^{}\widehat{\rho }_S(\tau )]`$. The calculations of these terms are too lengthy (though straightforward) in virtue of the many degrees of freedom of the system, and so they will be omitted. The final result can be written as
$`R_Q(\tau )=\left(\mathrm{\Gamma }_{1Q}𝕄_1[\theta _{1Q}]\right)\left(\mathrm{\Lambda }_{1Q}𝕄_1[\varphi _{1Q}]\right)R_0,`$ (14a)
where
$`\mathrm{\Gamma }_{1Q}(\tau )`$ $``$ $`\left|C_1(\tau )\right|,`$ (14b)
$`\mathrm{\Lambda }_{1Q}(\tau )`$ $``$ $`\mathrm{exp}\left[\left({\displaystyle \frac{R_0^\text{T}R_0}{\mathrm{}}}\right)\mathrm{sin}^2\tau \right],`$ (14c)
$`\theta _{1Q}(\tau )`$ $``$ $`ı\mathrm{ln}\left[{\displaystyle \frac{|C_1(\tau )|}{C_1(\tau )}}\right]+{\displaystyle \frac{G_1\tau }{2g_s}},`$ (14d)
$`\varphi _{1Q}(\tau )`$ $``$ $`2\tau +{\displaystyle \frac{\omega _s\tau }{\mathrm{}g_s}}+\left({\displaystyle \frac{R_0^\text{T}R_0}{2\mathrm{}}}\right)\mathrm{sin}(2\tau ).`$ (14e)
The definition of the auxiliary dimensionless functions, $`\mathrm{\Gamma }_{1Q}`$, $`\mathrm{\Lambda }_{1Q}`$, $`\theta _{1Q}`$, and $`\varphi _{1Q}`$, in equations (14b)-(14e), allowed to write the final result (14a) in such a suitable compact form, which will be useful for further comparisons. Interestingly, in Eq.(14a) the term $`\mathrm{\Lambda }_{1Q}𝕄_1[\varphi _{1Q}]`$, depending solely on parameters associated with the quartic oscillator, and the term $`\mathrm{\Gamma }_{1Q}𝕄_1[\theta _{1Q}]`$, related to the interaction with the reservoir, have been factorized. Note that $`\mathrm{\Gamma }_{1Q}`$ is an attenuation factor describing the environmental decoherence (see Eq.(12)) whereas the factor $`\mathrm{\Lambda }_{1Q}`$ is the one responsible for the well known revivals occurring in the closed dynamics of the quartic oscillator lili ; milburn ; agarwal . In fact, in the absence of the reservoir ($`\mathrm{\Omega }=0`$ and $`\mathrm{\Gamma }_{1Q}𝕄_1[\theta _{1Q}]=\mathrm{𝟏}`$), one may verify from (14c) the occurrence of revivals at the instants
$`\tau _R=n\pi (n=1,2,3,\mathrm{}).`$ (15)
This phenomenon, which has been shown to be a signature of the self-interference mechanism lili , has no classical analogue, as can be seen in the corresponding Liouvillian result given by Eq.(22e). In references adelcio ; milburn , numerical simulations comparing quantum (Wigner and Husimi) with classical distributions indeed demonstrate the absence of interference in the classical dynamics.
Note by (7) and (14) that one may rigorously verify that
$`\underset{\mathrm{}0}{lim}R_Q=R_N`$ (16)
in the absence of environmental reservoir ($`\mathrm{\Omega }=0`$). In contrast, for $`\mathrm{\Omega }0`$ this relation no longer holds, what can be demonstrated from (11) by showing that $`\underset{\mathrm{}0}{lim}|C_{nn^{}}|<1`$, for $`t>0`$. It is important to realize that this is not an effect of the many degrees of freedom composing the reservoir. In fact, the inequality applies even for only one oscillator in the reservoir. Rather, it occurs because the entanglement among the subsystems does not vanish asymptotically as $`\mathrm{}`$ tends to zero (see, e.g., angeloK05 ).
For the quadratic terms, $`R_{2Q}\left(\genfrac{}{}{0pt}{}{\widehat{q}_s^2}{\widehat{p}_s^2}\right)`$, a really tedious but equally straightforward calculation yields
$`R_{2Q}(\tau )=\left({\displaystyle \frac{R_0^\text{T}R_0+\mathrm{}}{2}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right)+`$
$`\mathrm{\Gamma }_{2Q}\mathrm{\Lambda }_{2Q}\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right){\displaystyle \frac{R_0^\text{T}}{\sqrt{2}}}𝕄_2\left[\theta _{2Q}+\varphi _{2Q}\right]{\displaystyle \frac{R_0}{2}},`$ (17a)
where
$`\mathrm{\Gamma }_{2Q}(\tau )`$ $``$ $`\left|C_2(\tau )\right|,`$ (17b)
$`\mathrm{\Lambda }_{2Q}(\tau )`$ $``$ $`\mathrm{exp}\left[\left({\displaystyle \frac{R_0^\text{T}R_0}{\mathrm{}}}\right)\mathrm{sin}^2(2\tau )\right],`$ (17c)
$`\theta _{2Q}(\tau )`$ $``$ $`ı\mathrm{ln}\left[{\displaystyle \frac{|C_2(\tau )|}{C_2(\tau )}}\right]+{\displaystyle \frac{G_1\tau }{g_s}},`$ (17d)
$`\varphi _{2Q}(\tau )`$ $``$ $`6\tau +{\displaystyle \frac{2\omega _s\tau }{\mathrm{}g_s}}+{\displaystyle \frac{R_0^\text{T}R_0}{2\mathrm{}}}\mathrm{sin}(4\tau ),`$ (17e)
$`𝕄_2[\chi ]`$ $``$ $`\left[\begin{array}{cc}\mathrm{cos}\chi & \mathrm{sin}\chi \\ & \\ \mathrm{sin}\chi & \mathrm{cos}\chi \end{array}\right].`$ (17i)
Once again, several dimensionless functions have been defined to guarantee compactness of the main result (17a). Also, $`\mathrm{\Gamma }_{2Q}`$ and $`\mathrm{\Lambda }_{2Q}`$ are terms associated with the reservoir and with the quartic oscillator, respectively.
Then, for the quantum variances, written as
$`\left(\mathrm{\Delta }R_Q\right)^2\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{\Delta }q_Q^2}{\mathrm{\Delta }p_Q^2}}\right)=\left({\displaystyle \genfrac{}{}{0pt}{}{\widehat{q}_s^2\widehat{q}_s^2}{\widehat{p}_s^2\widehat{p}_s^2}}\right),`$ (18)
we have
$`\left(\mathrm{\Delta }R_Q\right)^2=\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right)\left[{\displaystyle \frac{\mathrm{}}{2}}+{\displaystyle \frac{R_0^\text{T}R_0}{2}}\left(1\mathrm{\Gamma }_{1Q}^2\mathrm{\Lambda }_{1Q}^2\right)\right]+`$
$`\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right){\displaystyle \frac{R_0^\text{T}}{\sqrt{2}}}\left\{\mathrm{\Gamma }_{2Q}\mathrm{\Lambda }_{2Q}𝕄_2[\psi _{2Q}]\mathrm{\Gamma }_{1Q}^2\mathrm{\Lambda }_{1Q}^2𝕄_2[2\psi _{1Q}]\right\}{\displaystyle \frac{R_0}{\sqrt{2}}},`$
where $`\psi _{iQ}\varphi _{iQ}+\theta _{iQ}`$, with $`i=1,2`$.
### III.3 Liouvillian results
The time evolved Liouvillian distribution is obtained by integrating the Liouville equation, as shown in ref.angelo04 . The formal solution is given by
$`\rho (q_k,p_k,t)=\rho (q_k(t),p_k(t),0).`$ (20)
This formula establishes that the result can be obtained by replacing the arguments $`q_k`$ and $`p_k`$ of the initial distribution $`\rho (q_k,p_k,0)`$ given by Eq.(5) by the Newtonian trajectories evolved backwards in time, $`q_k(t)`$ and $`p_k(t)`$, respectively.
Classical statistical averages are calculate as
$`F_L(t)={\displaystyle 𝑑VF\rho (q_k,p_k,t)},`$ (21)
where $`dV=dq_sdp_sdq_1dp_1\mathrm{}dq_Ndp_N`$.
With formulas (20) and (21), it is possible do calculate the Liouvillian averages $`q_s`$, $`q_s^2`$, $`p_s`$, $`p_s^2`$, and also the related variances, $`\mathrm{\Delta }q_L^2=q_s^2q_s^2`$ and $`\mathrm{\Delta }p_L^2=p_s^2p_s^2`$. As for the quantum analysis, although straightforward, the calculations are too lengthy and do not add any further information to the results, so they can be omitted. The Liouvillian results may be written as:
$`R_L(\tau )=\left(\mathrm{\Gamma }_{1L}𝕄_1[\theta _{1L}]\right)\left(\mathrm{\Lambda }_{1L}𝕄_1[\varphi _{1L}]\right)R_0,`$ (22a)
where
$`D_m(\tau )`$ $``$ $`{\displaystyle \underset{k}{}}\left[1ım\left({\displaystyle \frac{g_k}{g_s}}\right){\displaystyle \frac{z\tau }{2}}\right]^1,`$ (22b)
$`\mathrm{\Gamma }_{1L}(\tau )`$ $``$ $`\left|D_1(\tau )\right|,`$ (22c)
$`\eta _1(\tau )`$ $``$ $`1+\tau ^2,`$ (22d)
$`\mathrm{\Lambda }_{1L}(\tau )`$ $``$ $`\eta _1^1\mathrm{exp}\left\{\left({\displaystyle \frac{R_0^\text{T}R_0}{\mathrm{}}}\right){\displaystyle \frac{\tau ^2}{\eta _1(\tau )}}\right\},`$ (22e)
$`\theta _{1L}(\tau )`$ $``$ $`ı\mathrm{ln}\left[{\displaystyle \frac{|D_1(\tau )|}{D_1(\tau )}}\right],`$ (22f)
$`\varphi _{1L}(\tau )`$ $``$ $`{\displaystyle \frac{\omega _s\tau }{\mathrm{}g_s}}+{\displaystyle \frac{\tau }{\eta _1(\tau )}}\left({\displaystyle \frac{R_0^\text{T}R_0}{\mathrm{}}}\right),`$ (22g)
$`(\tau )`$ $``$ $`{\displaystyle \frac{1}{\eta _1}}\left[\begin{array}{cc}1\tau ^2& 2\tau \\ & \\ 2\tau & 1\tau ^2\end{array}\right],`$ (22k)
with $`^\text{T}=^\text{T}=\mathrm{𝟏}`$. $`𝕄_1`$ is the rotation matrix defined by (7f), $`R_L\left(\genfrac{}{}{0pt}{}{q_s}{p_s}\right)`$, and $`R_0=\left(\genfrac{}{}{0pt}{}{q_0}{p_0}\right)`$. In Eq.(22b), the index $`m`$ is such that $`m=1,2`$ (see, e.g., equations (22c) and (25b)).
In the short-time regime the classical attenuation factor reduces to
$`\left|D_m(\tau )\right|`$ $`=`$ $`e^{m^2\tau ^2/\tau _{DL}^2},`$ (23)
where $`\tau _{DL}`$, the Liouvillian analogue of the decoherence time (13), is given by
$`\tau _{DL}=\sqrt{8}{\displaystyle \frac{g_s}{G_2}}\mathrm{tanh}\left({\displaystyle \frac{\beta \mathrm{}\omega }{2}}\right).`$ (24)
Once again the reservoir contributions could be factorized in the analytical solution for the centroid. Here, however, an extra rotation associated with the nonlinearity of the oscillator is imposed upon the initial condition $`R_0`$ through the matrix $``$.
It is worth emphasizing that the attenuation factor $`\mathrm{\Lambda }_{1L}`$ does not exhibit any revival, as has been mentioned. While quantum distribution are able to suffer self-interference, which in turn produces relocalization and revival in the amplitude of the centroid, classical distributions do not interfere, as shown in adelcio ; milburn .
For the quadratic terms, $`R_{2L}\left(\genfrac{}{}{0pt}{}{q_s^2}{p_s^2}\right)`$, we have
$`R_{2L}(\tau )=\left({\displaystyle \frac{R_0^\text{T}R_0+\mathrm{}}{2}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right)+`$
$`\mathrm{\Gamma }_{2L}\mathrm{\Lambda }_{2L}\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right){\displaystyle \frac{R_0^\text{T}}{\sqrt{2}}}𝕄_2\left[\theta _{2L}+\varphi _{2L}\right]_2{\displaystyle \frac{R_0}{\sqrt{2}}},`$ (25a)
where
$`\mathrm{\Gamma }_{2L}(\tau )`$ $``$ $`\left|D_2(\tau )\right|,`$ (25b)
$`\mathrm{\Lambda }_{2L}(\tau )`$ $``$ $`\eta _2^1\mathrm{exp}\left[4\left({\displaystyle \frac{R_0^\text{T}R_0}{\mathrm{}}}\right){\displaystyle \frac{\tau ^2}{\eta _2(\tau )}}\right],`$ (25c)
$`\eta _2(\tau )`$ $``$ $`1+4\tau ^2,`$ (25d)
$`\theta _{2L}(\tau )`$ $``$ $`ı\mathrm{ln}\left[{\displaystyle \frac{|D_2|}{D_2}}\right],`$ (25e)
$`\varphi _{2L}(\tau )`$ $``$ $`{\displaystyle \frac{2\omega _s\tau }{\mathrm{}g_s}}+2\left({\displaystyle \frac{R_0^\text{T}R_0}{\mathrm{}}}\right){\displaystyle \frac{\tau }{\eta _2}},`$ (25f)
$`_2`$ $``$ $`{\displaystyle \frac{1}{\eta _2^2}}\left[\begin{array}{cc}112\tau ^2& 2\tau (34\tau ^2)\\ & \\ 2\tau (34\tau ^2)& 112\tau ^2\end{array}\right].`$ (25j)
For the Liouvillian variances, written as
$`\left(\mathrm{\Delta }R_L\right)^2\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{\Delta }q_L^2}{\mathrm{\Delta }p_L^2}}\right)=\left({\displaystyle \genfrac{}{}{0pt}{}{q_s^2q_s^2}{p_s^2p_s^2}}\right),`$ (26)
the following result has been obtained:
$`\left(\mathrm{\Delta }R_L\right)^2=\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right)[{\displaystyle \frac{\mathrm{}}{2}}+{\displaystyle \frac{R_0^\text{T}R_0}{2}}(1\mathrm{\Gamma }_{1L}^2\mathrm{\Lambda }_{1L}^2)]+\left({\displaystyle \genfrac{}{}{0pt}{}{1}{1}}\right)\times `$
$`{\displaystyle \frac{R_0^\text{T}}{\sqrt{2}}}\left\{\mathrm{\Gamma }_{2L}\mathrm{\Lambda }_{2L}𝕄_2[\psi _{2L}]_2\mathrm{\Gamma }_{1L}^2\mathrm{\Lambda }_{1L}^2𝕄_2[2\psi _{1L}]_1^2\right\}{\displaystyle \frac{R_0}{\sqrt{2}}},`$
where $`\psi _{iL}\varphi _{iL}+\theta _{iL}`$ with $`i=1,2`$.
Figures 1 and 2 illustrates the main characteristics of quantities such as $`R_N^\text{T}R_N=R_0^\text{T}R_0`$ (Newtonian), $`R_Q^\text{T}R_Q=\mathrm{\Gamma }_{1Q}^2\mathrm{\Lambda }_{1Q}^2R_0^\text{T}R_0`$ (quantum), and $`R_L^\text{T}R_L=\mathrm{\Gamma }_{1L}^2\mathrm{\Lambda }_{1L}^2R_0^\text{T}R_0`$ (Liouvillian) for the closed ($`\mathrm{\Omega }=0`$) and the open ($`\mathrm{\Omega }>0`$) dynamics of the quartic oscillator. Quadratic terms such as $`R^\text{T}R`$, constructed with the results (7a), (14a), and (22a), are convenient to plot because they do not display the oscillations due to the rotation matrix $`𝕄_1`$. The effect of diffusive decoherence in destructing quantum revivals becomes evident in the figure. Next, the correspondence principle is discussed with the help of these results.
## IV Quasi-Determinism
The concept of determinism is mainly concerned with philosophical beliefs, according to which there are laws (hidden or not) univocally connecting the past and the future of every dynamical system, including animals and human beings. There is no chance for free choice. Note that the success of Newtonian theory in making predictions for the macroscopic world gives scientific support for such a philosophical concept. On the other hand, any attempt to apply the concept to the microscopic world will fail, once the notion of a trajectory linking the past and the future does not exist in the quantum realm. In fact, variances are always present, as predicted by Heisenberg’s uncertainty principle. Despite this fact, it is still possible to defend the idea of determinism by speculating on the existence of some sort of hidden laws (or hidden variables).
This paper, however, does not intend to settle this controversial question from an ontological viewpoint, asserting at the end whether determinism exist or not. Rather, the aim is to use the tools of the physical science, namely, fundamental laws and experimental investigations, to evaluate under which conditions the concept of determinism is acceptable in physics. In particular, the following test will be employed here: if the position and the momentum of a particle can be correctly predicted at a given instant for an individual system, then the concept of determinism will be regarded as a meaningful one in physics. This is what we call a scenario of quasi-determinism, since the motion will be experimentally indistinguishable from the truly deterministic Newtonian one. It is important to keep in mind this specific connotation attributed here to the word determinism, since it will be used from now on to incite the discussion about the central problem: the correspondence principle.
The argument motivating the historical statement of the correspondence principle is the one according to which quantum results must match classical ones because macroscopic world is correctly described by classical physics. However, from a rigorous point of view this reasoning cannot be considered as completely correct. Actually, it strongly depends on our capability of making the experimental checking. Extremely accurate measurements on the center of mass position of a macroscopic projectile near the surface of the Earth, would indicate some dispersion around the Newtonian parabolic trajectory. In fact, once we admit that the classical world is governed by quantum laws at a microscopic level, we have to conceive the center of mass motion as being determined by the quantum dynamics of all quantum particles composing the projectile. Accordingly, the center of mass is expected to behave quantum mechanically, with variances such as $`\mathrm{\Delta }x_{cm}=(\widehat{x}_{cm}^2\widehat{x}_{cm}^2)^{\frac{1}{2}}`$ being associated with its motion. That is, in principle, the quantum nature of the center of mass, and thus the deviation from Newtonian predictions, could be verified by such an accurate measurements. In this context, the status of the Newtonian mechanics as the basis for enunciating the correspondence principle is questionable.
To circumvent these difficulties, the following approach is proposed. We assume that the quantum-classical correspondence is achieved when quantum predictions are consistent with the concept of determinism (such as defined above). That is, if quantum mechanics is capable of making correct predictions for an individual system (what is expected to occur in a macroscopic regime) then the correspondence principle will be satisfied. As a consequence, since in this case the wave function dispersion is expected not to be detectable in virtue of the insufficient experimental resolution, Newtonian predictions can be successfully applied as well. Then, the approach is the following: for insufficient experimental accuracy, quantum dispersion is not apparent, Newtonian predictions is adequate, and the concept of deterministic behavior survives in physics. In this case, in virtue of Ehrenfest’s theorem, quantum mechanics will also succeed in making predictions for individual systems, at least for the short-time dynamics of initially localized wave functions (see, e.g., the famous example of the chaotic moon of Saturn zurek98 ; zurek96 , for which the estimated Ehrenfest time is about two decades) and, at the end, Einstein’s requirements will be satisfied.
Let us consider, as a first example, the quantum dynamics of a free particle initially prepared in a minimal uncertainty state. A well known calculation cohen shows that the product of the quantum variances evolves as
$`\mathrm{\Delta }x(t)\mathrm{\Delta }p(t)={\displaystyle \frac{\mathrm{}}{2}}\sqrt{1+\left({\displaystyle \frac{\mathrm{}t}{2m\mathrm{\Delta }x_0^2}}\right)^2}.`$ (28)
Consider now an experience in which the phase space resolution, defined by the measurement device, is given by the action $`\delta S`$, which here will be given by $`\delta S=M\mathrm{}/2`$, $`M`$ being an arbitrary real number. The product given in (28) achieves the critical value $`\delta S`$ for
$`t_{det}=t_EM\sqrt{1{\displaystyle \frac{1}{M^2}}},`$ (29)
where
$`t_E={\displaystyle \frac{2m\mathrm{\Delta }x_0^2}{\mathrm{}}}.`$ (30)
$`t_E`$ is the time scale firstly derived by Ehrenfest ehrenfest for the critical spreading of the free particle wave function. (In his work, Ehrenfest showed that $`t_E`$ can reach large values for macroscopic particles.) Now, within the time scale defined by $`t_{det}`$ the quantum dispersion of the wave function is smaller than the experimental resolution, so it is not possible to distinguish the quantum character of the phenomenon. In this case, Newtonian mechanics can be successfully applied to make predictions about the future phase space state $`(x,p)`$ of the particle. In other words, up to the instant $`t_{det}`$ the idea of a deterministic behavior is physically acceptable (within the experimental precision considered). Then, within this time scale, quantum mechanics satisfies Einstein’s principle as a theory capable of making predictions for individual systems, and in turn, the quantum-classical correspondence (as we have defined it) has been achieved. Beyond this time scale, quantum mechanics indeed fails in making predictions for individual systems, but exactly the same thing occurs with Newtonian mechanics. The determinism has been lost and the correspondence breaks off.
Now these ideas are applied to the dynamics of the nonlinear oscillator in the presence of the reservoir. Is decoherence capable of extending the time scale of determinism? It has already been pointed out in some works wiebe ; ballentine that decoherence is not able to reduce the wave function spreading. Since the spreading is the basic mechanism responsible for the breakdown in the determinism, as it has been suggested by the example above, one can anticipate that decoherence is not able to realize such a task. In fact that can be demonstrated in our model as follows.
Firstly, it is easy to see that decoherence does not restrain spreading. Using $`g_s=0`$ (the system of interest is now a harmonic oscillator) one obtains, by (12) and (III.2), the following long-time ($`\tau >\tau _{DQ}`$) variances product:
$`\mathrm{\Delta }q_Q(\mathrm{})\mathrm{\Delta }p_Q(\mathrm{})={\displaystyle \frac{\mathrm{}}{2}}+{\displaystyle \frac{R_0^\text{T}R_0}{2}}.`$ (31)
This result is larger than $`\mathrm{}/2`$, which is the value obtained for the closed dynamics of a harmonic oscillator. Clearly, purely diffusive decoherence has induced diffusion.
Secondly, let us consider the nonlinear oscillator under decoherence in a situation in which the experimental resolution is written as $`\delta S=M\mathrm{}/2`$. Since $`q`$ and $`p`$ plays essentially the same role in our model we chose for convenience an average phase space variance $`\mathrm{\Delta }^2(\mathrm{\Delta }q_Q^2+\mathrm{\Delta }p_Q^2)/2`$ instead of the product $`\mathrm{\Delta }q_Q\mathrm{\Delta }p_Q`$. Defining the determinism break time, $`\tau _{det}`$, as the instant at which the mean variance becomes comparable to the experimental resolution ($`\mathrm{\Delta }^2=\delta S`$), one may use (III.2) to obtain
$`1\mathrm{\Gamma }_{1Q}^2(\tau _{det})\mathrm{\Lambda }_{1Q}^2(\tau _{det})={\displaystyle \frac{\mathrm{}}{R_0^\text{T}R_0}}(M1).`$ (32)
Interestingly, a quite similar result is obtained when one compares the quantum centroid (14a) with the Newtonian trajectory (7a) through the prescription $`\frac{1}{2}(R_N^\text{T}R_NR_Q^\text{T}R_Q)=\delta S`$. This attests the consistence of our approach: the determinism break time, $`\tau _{det}`$, which corresponds to the time of critical spreading, is directly related to the instant at which quantum centroid and Newtonian trajectory stop agreeing.
Expanding Eq.(32) for $`\tau _{det}1`$ one obtains
$`\tau _{det}={\displaystyle \frac{\tau _{det}^{(c)}}{\sqrt{1+\left[\frac{\tau _{det}^{(c)}}{\tau _{DQ}}\right]^2\frac{2R_0^\text{T}R_0}{(M1)\mathrm{}}}}},`$ (33)
where
$`\tau _{det}^{(c)}=\sqrt{{\displaystyle \frac{(M1)\mathrm{}^2}{2(R_0^\text{T}R_0)^2}}}.`$ (34)
In these expressions $`\tau _{det}^{(c)}`$ denotes the break time for the closed dynamics of the nonlinear oscillator and $`\tau _{DQ}`$ the decoherence time given by Eq.(13). Equation (33) demonstrates that decoherence does not extend the determinism time scale, since $`\tau _{det}<\tau _{det}^{(c)}`$. Therefore, in the context of Einstein’s requirements for individual systems, decoherence plays no essential role.
In addition, by (8) one may write $`\tau _{det}=\mathrm{}g_st_{det}`$ to show, with the help of the results above, that $`t_{det}\mathrm{}^{\frac{1}{2}}t_E`$, where $`t_E`$ is the Ehrenfest time for the quartic oscillator angelo03 ; berman81 . This is another indicative that the time scale of determinism is indeed equivalent to the Ehrenfest one.
## V The quantum to classical transition
Once we accept the validity of the uncertainty principle at the basis of the microscopic phenomena we are forced to regard the Newtonian mechanics as an incomplete theory. In fact, since the initial phase space state $`(q_0,p_0)`$ of a particle can never be determined with arbitrary precision, the future phase space state can be predicted only statistically. Hence, the feeling of an objective reality disappears. In this scenario, the Liouvillian theory turns out to be the more general classical theory, whereas the Newtonian theory works as an approximation for the short-time regime of narrow distributions, as suggested in the precedent section. Thus, the correspondence principle is now discussed by comparing quantum with Liouvillian results, for times beyond the determinism time scale defined in the precedent section.
Firstly one should note from our analytical and exact results that there is no regime of parameters (with or without the reservoir) capable of making quantum and classical results mathematically identical. Although apparently naive, this sentence emphasizes, in agreement with adelcio ; wiebe ; ballentine , the approximated character of the quantum-classical correspondence and, in turn, the role of the experimental resolution in defining it. Accordingly, once again the theoretical action $`\delta S=M\mathrm{}/2`$, with $`M`$ real, will play the role of the experimental phase space resolution in our calculations.
The future phase space state of an individual system subjected to either quantum or classical statistical fluctuations can be predicted only in terms of averages and variances as $`(q,p)\pm \frac{1}{2}(\mathrm{\Delta }q,\mathrm{\Delta }p)`$. Then, the quantum-classical correspondence will be investigated in terms of the following differences:
$`D(\tau )={\displaystyle \frac{1}{2}}\left|R_Q^\text{T}(\tau )R_Q(\tau )R_L^\text{T}(\tau )R_L(\tau )\right|`$ (35)
for the centroids and
$`D(\tau )=\left|{\displaystyle \frac{\mathrm{\Delta }q_Q^2(\tau )+\mathrm{\Delta }p_Q^2(\tau )}{2}}{\displaystyle \frac{\mathrm{\Delta }q_L^2(\tau )+\mathrm{\Delta }p_L^2(\tau )}{2}}\right|`$ (36)
for the variances.
These definitions disregard, for simplicity, the influence of the rotation matrices $`𝕄_{1,2}`$ and $``$ but capture the attenuation effects due to both the reservoir and the nonlinearity. Using results (14a), (22a), (III.2), and (III.3), its is possible to show that both equations above reduce to
$`D(\tau )=\left|\mathrm{\Gamma }_{1Q}^2(\tau )\mathrm{\Lambda }_{1Q}^2(\tau )\mathrm{\Gamma }_{1L}^2(\tau )\mathrm{\Lambda }_{1L}^2(\tau )\right|{\displaystyle \frac{R_0^\text{T}R_0}{2}}.`$ (37)
This result demonstrating the equivalence between (35) and (36) suggests the existence in this model of a common dynamics for the quantum-classical differences associated with the centroids and with the variances. As a consequence, one may conceive the existence of a single time scale describing the breakdown in the correspondence, this being a counter-example contributing to the discussion raised in ref.adelcio03 .
In Fig.3 the behavior of the function $`D`$ is shown as a function of $`\tau `$ for the closed and the open dynamics of the nonlinear oscillator in two different regimes of temperature. Note that diffusive decoherence, modelled in terms of our phase-damping reservoir, is indeed effective in restoring the quantum-classical correspondence. However, a word concerning the role of temperature is in order.
Results (13) and (24) indicate that the attenuation effects due to the reservoir may occur also for low temperatures. However, in such a regime the differences between quantum and Liouvillian decay times become more accentuated. Indeed, from (13) and (24) it is easy to show that
$`\tau _{DQ}=\tau _{DL}\mathrm{cosh}\left({\displaystyle \frac{\beta \mathrm{}\omega }{2}}\right).`$ (38)
Now, the differences in the decays for the low-temperature regime ($`\beta \mathrm{}\omega 1`$) become evident. This result turns out to be an indicative that the high-temperature regime is an additional necessary condition for the establishment of the quantum-classical transition. However, this is not the whole truth, since the condition $`\tau _{DQ}=\tau _{DL}`$ does not automatically imply that $`\tau _{DQ}<\tau _R`$, where $`\tau _R`$ is the revival time given by (15). In fact, since $`\tau _R`$ is the time scale which determines the appearance of purely quantum effects, such as quantum interference, the classical limit is guaranteed by imposing that decoherence occurs early. In general one may state that the quantum-classical transition is achieved as long as the decoherence time be smaller than the time scale for which quantum phenomenon (e.g., self-interference) takes place in the dynamics. In our model this condition implies, by (13) and (15), that
$`{\displaystyle \frac{k_BT}{\mathrm{}\omega }}>{\displaystyle \frac{g_s}{2\mathrm{\Omega }}},`$ (39)
where the high-temperature regime ($`\beta \mathrm{}\omega 1`$) and the thermodynamic limit ($`N\mathrm{}`$, $`G_20.89\mathrm{\Omega }`$) have already been considered. This simple relation establishes the lower bound for the value of the reservoir temperature capable of yielding the quantum-classical transition. In addition, this result shows that nondissipative decoherence inducing quantum-classical transition can occur at finite temperatures.
## VI Summary
In this paper, the exact quantum and Liouvillian dynamics of a quartic oscillator coupled with a purely diffusive reservoir at arbitrary temperature have been solved analytically in terms of expectation values and variances associated with the phase space variables. The results demonstrate the effectiveness of our reservoir model in describing nondissipative decoherence and inducing quantum-classical transition.
Concerning the correspondence principle, this paper intend to contribute in two main directions. On one hand, examples were given showing that the Ehrenfest time is intimately connected with the time scale within which the concept of determinism (suitably defined) is acceptable in physics. In fact, within the Ehrenfest time scale initially localized wave functions remain sufficiently narrow in such a way that its dispersion cannot overcome the limits imposed by the experimental resolution. In this case, Newtonian and quantum mechanics are both well succeeded in making predictions for individual systems, thus satisfying Einstein’s requirements. Beyond the Ehrenfest time scale quantum predictions deviate from Newtonian ones, but the latter are no longer correct. This analysis gives answer to Einstein’s insatisfaction and also attributes a new physical meaning to the Ehrenfest time. In addition, it has been shown that diffusive decoherence plays no essential role within such a time scale, this being a point of agreement with the claims of ref.wiebe ; ballentine .
On the other hand, in a second direction, this paper shows that although decoherence is indeed necessary to reestablish quantum-classical correspondence beyond the determinism time scale, a certain amount of coarse-graining (related to experimental resolution) is also required in order to characterize the correspondence in a more precise way. This is necessary since the quantum-classical correspondence cannot be fully achieved at a mathematical level. Also, our results points out to the fact that the high-temperature regime is another fundamental ingredient in the scenario of the quantum-classical correspondence. A quantification of how high the reservoir temperature must be, in order to guarantee the correspondence, is given for our model in terms of a lower bound for the temperature. Although the assumption of $`T\mathrm{}`$ appears in several approaches (see, e.g., adelcio ; toscano ) as an artifact to simulates the effects of diffusive decoherence from the traditional dissipative master equations, here it is justified in terms of physical conditions. Yet, our approach points out to the possibility of modelling diffusive decoherence at finite temperatures.
###### Acknowledgements.
The author would like to thank A. F. Gomes, A. D. Ribeiro, C. F. Woellner, F. R. L. Parisio and A. S. Sant’Anna for fruitful discussions and helpful suggestions.
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# Nonlinear spin relaxation in strongly nonequilibrium magnets
## I Introduction
The problem of spin relaxation from a state close to equilibrium has a long history and is well studied, being related to the description of spin motion in the vicinity of different magnetic resonances. This type of spin relaxation is usually characterized by linear differential equations, such as Bloch equations. The theory of spin motion close to equilibrium has been expounded in numerous literature, among which it would be possible to mention several good books \[1–7\].
Essentially nonlinear spin motion arises if the system is prepared in a strongly nonequilibrium initial state, e.g. with magnetization opposite to an external magnetic field, and, in addition, is coupled to a resonator. Such nonlinear dynamics are commonly treated by the Bloch equations supplemented by the Kirchhoff equation for a resonator electric circuit \[8–11\]. However, the phenomenological Bloch equations do not allow for the elucidation of different physical processes involved in the behaviour of the system and are not able to describe several, probably the most interesting, self-organized regimes of spin motion, as was demonstrated in Refs. \[12–14\]. Some physical models, based on microscopic spin Hamiltonians, have also been considered, whose survey can be found in recent reviews . But in each of these models one standardly studies only some particular substances and considers only a part of spin interactions, mainly secular dipole-dipole interactions, and one takes into account only some of the known attenuation processes. At the same time, it is evident that taking care of only particular model elements can easily lead to wrong physical conclusions, since real physical materials always include several different characteristics competing with each other. The study of nonlinear spin relaxation is of paramount importance not solely owing to its theoretical beauty but also because it can be employed in a variety of applications, such as the measurement of materials parameters, ultrafast repolarization of solid-state targets, creation of sensitive field detectors, usage in quantum computing and others, as is discussed in reviews . One of the major possible applications is in achieving the regime of superradiant operation by spin masers \[13,17–19\]. Punctuated nonlinear dynamics of spin assemblies can also be a new tool for information processing .
The aim of the present paper is to develop a general theory of nonlinear spin relaxation, being based on a realistic microscopic Hamiltonian including, in addition to the Zeeman terms, the main spin interactions, and taking account of the different major mechanisms of spin attenuation. By considering just some limited models, it is easy to come to false conclusions and to predict fictitious physical effects that by no means can exist in real materials. It is only by carefully treating different competing mechanisms that one can derive reliable physical implications.
## II Basic Spin Hamiltonian
Keeping in mind the applicability of the theory to a wide class of spin systems, we start with a rather general Hamiltonian including the major spin interactions the most often met in magnetic materials \[1–7,21–23\]. Let us consider a solid sample containing $`N`$ vector spins $`𝐒_i`$ enumerated by the index $`i=1,2,\mathrm{},N`$. The spin operators $`𝐒_i`$ can represent any particles of spin $`S`$, starting from $`S=1/2`$ to very high spin values. These can be nuclear or electronic spins, as in the standard problems of nuclear or electronic spin resonances \[1–7,15\]. Magnetic molecules, forming molecular magnets, can possess various spins ranging from $`S=1/2`$ up to $`S=27/2`$, as is reviewed in Refs. \[16,19,24–26\]. Bose-Einstein condensates of dilute gases (see reviews \[27–30\]), being placed in optical lattices can form localized clouds with an effective spin per site of order $`10^2`$ or $`10^3`$. Spin dynamics (mainly linear) is an intensively developing field of research, called spintronics .
The Hamiltonian of a spin system can, generally, be separated into two parts,
$$\widehat{H}=\underset{i}{}\widehat{H}_i+\frac{1}{2}\underset{ij}{}\widehat{H}_{ij},$$
(1)
the first term being related to individual spins, while the second representing spin interactions. The single-spin Hamiltonian
$$\widehat{H}_i=\mu _0𝐁𝐒_iD(S_i^z)^2$$
(2)
consists of the Zeeman energy and the energy of the single-site magnetic anisotropy. Here $`\mu _0\mathrm{}\gamma _S`$, with $`\gamma _S`$ being the gyromagnetic ratio of a particle with spin $`S`$. For electronic spins, $`\mu _0<0`$, while for nuclear spins $`\mu _0`$ can be either positive or negative. The total magnetic field
$$𝐁=B_0𝐞_z+(B_1+H)𝐞_x$$
(3)
contains external longitudinal, $`B_0`$, and transverse, $`B_1`$, magnetic fields, and also a feedback field $`H`$ of a resonator, if the sample is coupled to a resonant electric circuit. The anisotropy parameter $`D`$ is positive for an easy-axis anisotropy and negative in the case of an easy-plane anisotropy.
The interaction Hamiltonian
$$\widehat{H}_{ij}=\underset{\alpha \beta }{}D_{ij}^{\alpha \beta }S_i^\alpha S_j^\beta J_{ij}𝐒_i𝐒_j$$
(4)
includes dipole and exchange interactions. The dipolar tensor is
$$D_{ij}^{\alpha \beta }=\frac{\mu _0^2}{r_{ij}^3}\left(\delta _{\alpha \beta }3n_{ij}^\alpha n_{ij}^\beta \right),$$
(5)
where $`\alpha ,\beta =x,y,z`$ and
$$r_{ij}|𝐫_{ij}|,𝐧_{ij}\frac{𝐫_{ij}}{r_{ij}},𝐫_{ij}𝐫_i𝐫_j.$$
This tensor enjoys the properties
$$\underset{\alpha }{}D_{ij}^{\alpha \alpha }=0,\underset{j(i)}{}D_{ij}^{\alpha \beta }=0,$$
(6)
of which the first is exact and the second one is asymptotically exact for a macroscopic sample with a large number of spins $`N1`$. A positive exchange integral corresponds to ferromagnetic interactions and negative, to antiferromagnetic interactions.
It is convenient to represent the Hamiltonians through the ladder spin operators $`S_i^\pm S_i^x\pm iS_i^y`$. Then the single-spin term (2) writes as
$$\widehat{H}_i=\mu _0B_0S_i^z\frac{1}{2}\mu _0(B_1+H)\left(S_i^++S_i^{}\right)D\left(S_i^z\right)^2.$$
(7)
With the notation
$$a_{ij}D_{ij}^{zz},b_{ij}\frac{1}{4}\left(D_{ij}^{xx}D_{ij}^{yy}2iD_{ij}^{xy}\right),c_{ij}\frac{1}{2}\left(D_{ij}^{xz}iD_{ij}^{yz}\right),$$
(8)
the interaction Hamiltonian (4) transforms to
$$\widehat{H}_{ij}=a_{ij}\left(S_i^zS_j^z\frac{1}{2}S_i^+S_j^{}\right)+b_{ij}S_i^+S_j^++b_{ij}^{}S_i^{}S_j^{}+$$
$$+2c_{ij}S_i^+S_j^z+2c_{ij}^{}S_i^{}S_j^zJ_{ij}\left(S_i^+S_j^{}+S_i^zS_j^z\right).$$
(9)
The interaction parameters $`a_{ij}=a_{ji}`$, $`b_{ij}=b_{ji}`$, and $`c_{ij}=c_{ji}`$ are symmetric and have the property
$$\underset{j(i)}{}a_{ij}=\underset{j(i)}{}b_{ij}=\underset{j(i)}{}c_{ij}=0,$$
(10)
following from Eqs. (6).
The equations of motion for the spin operators are obtained from the Heisenberg equations and the commutation relations
$$[S_i^+,S_j^{}]=2\delta _{ij}S_i^z,[S_i^z,S_j^\pm ]=\pm \delta _{ij}S_i^\pm .$$
In order to represent the evolution equations in a compact form, it is convenient to introduce the local fields
$$\xi _0\frac{1}{\mathrm{}}\underset{j(i)}{}\left[a_{ij}S_j^z+c_{ij}^{}S_j^{}+c_{ij}S_j^++J_{ij}\left(S_i^zS_j^z\right)\right],$$
$$\xi \frac{i}{\mathrm{}}\underset{j(i)}{}\left[2c_{ij}S_j^z\frac{1}{2}a_{ij}S_j^{}+2b_{ij}S_j^++J_{ij}\left(S_i^{}S_j^{}\right)\right]$$
(11)
and the effective force
$$f\frac{i}{\mathrm{}}\mu _0(B_1+H)+\xi .$$
(12)
There is a characteristic frequency, the Zeeman frequency, which we denote as
$$\omega _0\frac{\mu _0}{\mathrm{}}B_0.$$
(13)
Then as the equations of motion for the spin operators, we obtain
$$\frac{dS_i^{}}{dt}=i(\omega _0+\xi _0)S_i^{}+fS_i^z+i\frac{D}{\mathrm{}}\left(S_i^{}S_i^z+S_i^zS_i^{}\right),$$
(14)
with its Hermitian conjugate, and
$$\frac{dS_i^z}{dt}=\frac{1}{2}\left(f^+S_i^{}+S_i^+f\right).$$
(15)
The following description of spin dynamics will be based on these equations.
## III Triggering Spin Fluctuations
Suppose that the spin system is prepared in a strongly nonequilibrium state, being polarized along the $`z`$-axis. What then could be the triggering mechanisms initiating spin motion and their relaxation to an equilibrium state? It is evident that imposing transverse magnetic fields would push the spins to move. But assume that there are no transverse magnetic fields at the initial time and no transverse coherence is imposed on the system. What then would initiate the spin motion? Here it is important to stress the role of local spin waves as of the triggering mechanism for starting the spin relaxation.
The appearance of spin waves is due to the local fields (11). In order to consider spin waves, or more generally, spin fluctuations that arise in a state which is not necessarily equilibrium, it is appropriate to work with the operator equations (14) and (15). Let us define the operator deviation
$$\delta S_i^\alpha S_i^\alpha <S_i^\alpha >$$
(16)
from an average $`<S_i^\alpha >`$, which is not necessarily an equilibrium average, but which can be an average over a nonequilibrium statistical operator, though such that $`<S_i^\alpha >`$ weakly depends on the index $`i`$, because of which it can be taken out of the sums in Eqs. (11). Then, owing to Eqs. (10), we have
$$\xi _0=\frac{1}{\mathrm{}}\underset{j(i)}{}\left[a_{ij}\delta S_j^z+c_{ij}^{}\delta S_j^{}+c_{ij}\delta S_j^++J_{ij}\left(\delta S_i^z\delta S_j^z\right)\right],$$
$$\xi =\frac{i}{\mathrm{}}\underset{j(i)}{}\left[2c_{ij}\delta S_j^z\frac{1}{2}a_{ij}\delta S_j^{}+2b_{ij}\delta S_j^++J_{ij}\left(\delta S_i^{}\delta S_j^{}\right)\right],$$
(17)
which demonstrates that these local fields really correspond to local spin fluctuations.
To emphasize the role of the spin fluctuations, let us set $`B_1=H=0`$, that is, looking at the case when the transverse fields do not initiate the spin motion. And, respectively, let $`<S_i^\pm >=0`$, but the longitudinal polarization be finite, $`<S_i^z>0`$. Then $`S_i^\pm =\delta S_i^\pm `$. The behaviour of spin fluctuations is characterized by linearizing Eqs. (14) and (15) with respect to the operator deviations (16). The linearization of the single-site anisotropy term in Eq. (14) has to be done so that to satisfy the known exact relations for $`S=1/2`$ and $`S\mathrm{}`$, which can be represented as
$$S_i^{}S_i^z+S_i^zS_i^{}=\left(2\frac{1}{S}\right)<S_i^z>S_i^{}.$$
(18)
Introduce the single-site anisotropy frequency
$$\omega _D(2S1)\frac{D}{\mathrm{}}$$
(19)
and the effective spin frequency
$$\omega _s\omega _0\omega _D\frac{<S_i^z>}{S},$$
(20)
where $`\omega _0`$ is defined in Eq. (13). Then, linearizing Eqs. (14) and (15), we find
$$\frac{d}{dt}S_i^{}=i\omega _sS_i^{}+<S_i^z>\xi ,\frac{d}{dt}\delta S_i^z=0.$$
(21)
The second of these equations, under the initial condition $`\delta S_i^z(0)=0`$, gives $`\delta S_i^z=0`$.
Now let us employ the Fourier transforms for the interactions
$$a_{ij}=\frac{1}{N}\underset{k}{}a_ke^{i𝐤𝐫_{ij}},a_k=\underset{j(i)}{}a_{ij}e^{i𝐤𝐫_{ij}},$$
with the analogous transforms for $`b_{ij}`$ and $`J_{ij}`$, and for the spin operators
$$S_j^\pm =\underset{k}{}S_k^\pm e^{i𝐤𝐫_j},S_k^\pm =\frac{1}{N}\underset{j}{}S_j^\pm e^{\pm i𝐤𝐫_j}.$$
Using the notation
$$\alpha _k\omega _s+\frac{1}{\mathrm{}}\left(\frac{a_k}{2}+J_kJ_0\right)<S_i^z>,\beta _k\frac{2}{\mathrm{}}b_k<S_i^z>,$$
(22)
from the first of Eqs. (14), we obtain
$$\frac{d}{dt}S_k^{}=i\alpha _kS_k^{}+i\beta _kS_k^+.$$
(23)
Looking for the solution of the latter equation in the form
$$S_k^{}=u_ke^{i\omega _kt}+v_k^{}e^{i\omega _kt},$$
we find the spectrum of spin waves
$$\omega _k=\sqrt{\alpha _k^2|\beta _k|^2}.$$
(24)
In the long-wave limit, one gets
$$\omega _k|\omega _s|\left[1<S_i^z>\underset{<j>}{}\frac{a_{ij}+2J_{ij}}{4\mathrm{}\omega _s}(𝐤𝐫_{ij})^2\right],$$
(25)
where $`k0`$, and the summation is over the nearest neighbours.
In this way, in the spin system there are always transverse fluctuations, which can be named spin waves. The latter, as they have been described, are not necessarily the spin waves in an equilibrium state, as they are usually understood , but are to be considered in a generalized sense. Under spin waves, we mean here just transverse spin fluctuations. It is these transverse fluctuations that are responsible for triggering the initial motion of polarized spins, when there are no external transverse magnetic fields. This is why these transverse spin fluctuations can be called triggering spin waves. Taking into account such quantum spin fluctuations makes it possible to describe the dynamical regimes of spin motion, which do not exist for classical Bloch equations. And it becomes possible to develop a detailed picture of how the transverse spin coherence arises from initially chaotic fluctuations. This self-organized process of coherence emerging from chaos is one of the most interesting and challenging problems of spin dynamics.
## IV Spin Evolution Equations
The equations of motion (14) and (15) for spin operators are highly nonlinear. The nonlinearity comes from two sources. One is caused by the spin interactions accumulated in the local fluctuating fields (11). Another kind of nonlinearity enters through the effective force (12) containing feedback fields included in the term $`H`$. The treatment of the nonlinear spin dynamics will be done here by means of the scale separation approach \[11–15,34\], which is a generalization of the averaging technique to stochastic differential equations.
Notice, first of all, that there are two different spatial scales. One of them is related to local fields (11) describing random spin fluctuations (17), which is characterized by a spatial length of the order of the mean interparticle distance $`a_0`$. At this length scale, chaotic quantum spin fluctuations prevail. Another length scale is the wavelength $`\lambda a_0`$ corresponding to coherent effects associated with the characteristic spin rotation frequency $`\omega _s`$. At the latter scale, coherent spin correlations are important. These two different length scales allow us to distinguish two types of operators. One type are the local fluctuating fields (11), that is, the variables $`\xi _0`$, $`\xi `$, and $`\xi ^+`$, and another type are the spin operators $`S_i^{}`$, $`S_i^+`$, and $`S_i^z`$. The former, responsible for local short-range fluctuations, can be represented by random variables , while the latter keep track of long-range coherent effects. Respectively, it is convenient to define two sorts of averaging with respect to the corresponding variables. Then the statistical averaging over spin operators will be denoted by the single angle brackets $`<\mathrm{}>`$, while the averaging over the random local fields will be denoted by the double angle brackets $`\mathrm{}`$. The latter, treating the chaotic local spin fluctuations as white noise, are defined as
$$\xi _0(t)=\xi (t)=0,\xi _0(t)\xi _0(t^{})=2\gamma _3\delta (tt^{}),$$
$$\xi _0(t)\xi (t^{})=\xi (t)\xi (t^{})=0,\xi ^{}(t)\xi (t^{})=2\gamma _3\delta (tt^{}),$$
(26)
where $`\gamma _3`$ is the width of inhomogeneous dynamic broadening.
It is worth stressing that the white-noise approximation (26) is not principal and could be generalized to taking into account a coloured noise by including finite relaxation times. This, however, would result in much more complicated and cumbersome equations. It is therefore more convenient, following the ideas of the scale separation approach \[11–15\], to separate in the temporal behaviour of spin correlations two parts, fast and slow. The fast part is connected to the local spin fluctuations described by the spectrum of local spin waves (24). The characteristic frequencies of these fluctuations are defined by the near-neighbour spin coupling as well as by the applied external magnetic field. Here and in what follows, we assume that this external field is sufficiently strong, so that the fluctuation spectrum (24) is characterized by the frequencies of the order of the Zeeman frequency $`\omega _0`$, which is essentially larger than the frequency terms due to spin interactions. With the time $`2\pi /\omega _0`$ being the shortest among all other characteristic times, the related fast spin fluctuations can be effectively treated as white noise, as is done in Eq. (26). The influence of spin correlations slowly decaying in time can be appropriately included into the transverse relaxation time $`T_2`$ determined by the strength of the spin-spin coupling allowing for dipolar as well as exchange interactions. This effective relaxation time will also be taken into account in the following consideration, together with the effect of line narrowing due to high spin polarization .
Averaging over spin operators, because of their long-range role, one can employ the decoupling
$$<S_i^\alpha S_j^\beta >=<S_i^\alpha ><S_j^\beta >(ij).$$
(27)
Though this looks like a mean-field approximation, one should not forget that the restricted averaging, denoted by the single angle brackets $`<\mathrm{}>`$, by definition, involves only the spin degrees of freedom, without touching the stochastic variables $`\xi _0`$ and $`\xi ^{}`$. Therefore the quantum fluctuations are not lost in decoupling (27) but are preserved because of the dependence of the spin averages $`<S_i^\alpha >`$ on the random variables $`\xi _0`$ and $`\xi `$. Then decoupling (27) is termed the stochastic mean-field approximation \[11–16\].
A special care is to be taken in considering the single-site term of Eq. (14). When averaging the latter, one has to preserve the exact limiting properties known for $`S=1/2`$ and $`S\mathrm{}`$. The corresponding decoupling, correctly interpolating between the exact limiting behaviours is
$$<S_i^{}S_i^z+S_i^zS_i^{}>=\left(2\frac{1}{S}\right)<S_i^{}><S_i^z>.$$
(28)
Thus, for $`S=1/2`$, expression (28) becomes zero, as it should be, and for $`S\mathrm{}`$, one has $`2<S_i^{}><S_i^z>`$, again in agreement with the correct asymptotic behaviour.
Let us average the equations of motion (14) and (15) over the spin degrees of freedom, not touching the fluctuating random fields $`\xi _0`$ and $`\xi `$. Our aim is to obtain the evolution equations for the following variables: The transition function
$$u\frac{1}{SN}\underset{i=1}{\overset{N}{}}<S_i^{}>,$$
(29)
describing the average rotation of transverse spin components; the coherence intensity
$$w\frac{1}{S^2N(N1)}\underset{ij}{\overset{N}{}}<S_i^+S_j^{}>,$$
(30)
showing the level of coherence in the spin motion, and the spin polarization
$$s\frac{1}{SN}\underset{i=1}{\overset{N}{}}<S_i^z>,$$
(31)
defining the average polarization per particle.
In order to have the evolution equations representing realistic spin systems, but not just some unreasonable models, an accurate account must be taken of the main relaxation mechanisms. Being based on unrealistic models, omitting important existing attenuation processes, it would be easy to fall into the sin of predicting physical effects that in reality can never occur. We shall consider the following basic relaxation rates.
(1) Spin-lattice longitudinal attenuation $`\gamma _1`$, caused by spin-lattice interactions. The corresponding longitudinal relaxation time is $`T_11/\gamma _1`$. For different materials, $`\gamma _1`$ can be of different order. At low temperature, when spin-phonon interactions are suppressed, the parameter $`\gamma _1`$ can be rather small. For instance, in polarized nuclear targets at temperature of 1 K, one has $`\gamma _110^5`$ s<sup>-1</sup>. In molecular crystals below the blocking temperature of the order of 1 K, the spin-lattice rate can be between $`\gamma _110^7`$ and $`10^5`$ s<sup>-1</sup> (see more details in Refs. \[16,24–26\]). Being small, this relaxation parameter may not play an essential role at the initial stage of spin motion, however, it always plays a principal role at the late stages of spin relaxation.
(2) Polarization pump rate $`\gamma _1^{}`$, which is added to $`\gamma _1`$ when the sample is subject to a permanent pump supporting a nonequilibrium level of the longitudinal spin polarization. This rate can be made much larger than $`\gamma _1`$. Thus, by means of dynamic nuclear polarization, the pump rate for nuclear spins in solids can be as large as $`\gamma _1^{}0.01`$ and $`10`$ s<sup>-1</sup> . The sum of $`\gamma _1`$ and $`\gamma _1^{}`$ will be denoted as
$$\mathrm{\Gamma }_1\gamma _1+\gamma _1^{}.$$
(32)
(3) Spin dephasing rate $`\gamma _2`$, due to spin-spin interactions. This rate has been calculated by many authors, and the generally accepted value \[1–7\] writes as
$$\gamma _2=n_0\rho \frac{\mu _0^2}{\mathrm{}}\sqrt{S(S+1)},$$
(33)
where $`\rho N/V`$ is density and $`n_0`$ is a coefficient approximately equal to the number of nearest neighbours. The process of spin dephasing is mainly due to dipolar forces. Exchange interactions slightly narrow the line width (33), yielding a factor of about 0.8. The coefficient in Eq. (33) also depends on the type of lattice, so that the numerical factor here is approximate. The value of $`\gamma _2`$ is usually larger than that of $`\gamma _1`$. For example, in polarized solid targets $`\gamma _210^5`$ s<sup>-1</sup>, in molecular magnets it is $`\gamma _210^{10}`$ s<sup>-1</sup>. Inverse of $`\gamma _2`$ defines the spin dephasing time $`T_21/\gamma _2`$.
(4) Effective homogeneous broadening $`\gamma _2(s)`$ takes into account a correction to the spin dephasing rate $`\gamma _2`$, appearing in the case of strongly polarized spin systems. Such a strong polarization can be achieved in magnetically ordered materials, by applying strong longitudinal magnetic fields, or by dynamic polarization techniques. This effective broadening reads as
$$\gamma _2(s)=\gamma _2(0)\left(1s^2\right),\gamma _2(0)\gamma _2,$$
(34)
where $`s`$ is an average spin polarization (31) and $`\gamma _2`$ is given by Eq. (33). The derivation of Eq. (34) is explained in Appendix A. Under weak polarization, when $`s^21`$, one has $`\gamma _2(s)\gamma _2`$.
(5) Static inhomogeneous broadening $`\gamma _2^{}`$ is due to various magnetic defects, crystalline defects, field gradients, and a variety of additional interactions always present in any real materials \[1–7,21,31\]. Very often the inhomogeneity develops in matter not because of externally incorporated defects, but being due to the internal properties, when a heterogeneous state is more thermodynamically stable than a homogeneous state . This, e.g., happens in many colossal-magnetoresistance materials \[39–41\] and in high-temperature superconductors \[42–46\], where there appears mesoscopic phase separation. In general, $`\gamma _2^{}`$ can be both smaller as well as larger than $`\gamma _2`$. However in the majority of cases, to a very good approximation $`\gamma _2^{}\gamma _2`$. Summarizing the homogeneous and inhomogeneous mechanisms, discussed above, we denote the overall transverse relaxation rate as
$$\mathrm{\Gamma }_2\gamma _2\left(1s^2\right)+\gamma _2^{}.$$
(35)
(6) Dynamic inhomogeneous broadening $`\gamma _3`$ is caused by fast dynamic spin fluctuations, or the local spin waves, discussed in Sec. III. It comes into play through the stochastic averaging (26). The value of the broadening, due to local spin waves, is of the order or smaller than $`\gamma _2`$ \[14–16,21\]. As is emphasized in Sec. III, this dynamic broadening is crucially important at the initial stage of spin relaxation, when there are no applied transverse fields.
(7) Cross relaxation rates arise when there are several spin species in the system. For example, if there are two types of spins, $`S`$ and $`F`$, then the dynamic broadening for spin $`S`$ becomes
$$\gamma _3=\sqrt{\gamma _{SS}^2+\gamma _{SF}^2}.$$
(36)
Cross correlations can influence other relaxation rates, especially if the Zeeman frequencies of the spins $`S`$ and $`F`$ are close to each other \[1–7,15,16\].
(8) Spin radiation rate $`\gamma _r`$ arises when there exist the so-called wave packets of strongly correlated spins interacting with each other through the common radiation field. The possibility of the appearance of such an electromagnetic friction was, first, noticed by Ginzburg and later discussed by many authors (see e.g. ). This collective radiation rate is
$$\gamma _r=\frac{2}{3\mathrm{}}\rho \mu _0^2S(kL_s)^3,$$
(37)
where $`k`$ is the wave vector of the radiating field and $`L_s`$ is an effective linear size of a spin packet radiating coherently. Rate (37) has earlier been obtained in the classical approximation. In Appendix B, we briefly sketch how this rate can be derived in a fully quantum-mechanical picture. It is important to stress that the existence of rate (37) presupposes the occurrence of monochromatic radiation with a well-defined constant spin frequency $`\omega _s`$ and wave vector $`k`$, and that the radiation wavelength is much larger than the linear size $`L_s`$ of a spin packet, so that
$$kL_s1\left(k\frac{\omega _s}{c}\right).$$
(38)
If these conditions do not hold, no noticeable relaxation rate arises. And under the validity of these conditions, one has
$$\frac{\gamma _r}{\gamma _2}0.1(kL_s)^31.$$
(39)
The rate $`\gamma _r`$ is so much smaller than $`\gamma _2`$, and usually much smaller than $`\gamma _2^{}`$, that it can be safely neglected, being absolutely unable to influence the motion of spins. Actually, Bloembergen has already analysed this problem and come to the conclusion that the interaction of spins through the magnetodipole radiation field is completely negligible. However, one may put the following question. Suppose that the considered sample is ideally homogeneous, so that $`\gamma _2^{}`$ is very small, and let the initial spin polarization be very high, such that $`s_0^21`$. Then the effective transverse rate (35) at the initial time $`t=0`$ can become rather small. Could then the radiation rate (37) play any noticeable role, at least at the very initial stage of spin motion? We study this problem below.
(9) Thermal noise attenuation $`\gamma _T`$ emerges when the spin system is coupled to a resonant electric circuit. The resonator Nyquist noise, due to the thermal fluctuations of current in the circuit creates a fluctuational magnetic field, which has to be included in the effective force (12). The magnitude of the thermal field, produced by the Nyquist noise, is well known . It was found \[12–16\] that the resulting thermal attenuation is
$$\gamma _T=\frac{\eta \rho \mu _0^2\omega }{4\mathrm{}\gamma N}\mathrm{coth}\frac{\omega }{2\omega _T},$$
(40)
where $`\eta `$ is a filling factor, $`\omega `$ is the natural frequency of the electric circuit, $`\gamma `$ is the resonator ringing width, and $`\omega _Tk_BT/\mathrm{}`$ is the thermal frequency. Bloembergen and Pound first mentioned that, because of the macroscopic number of spins $`N`$ entering the denominator of $`\gamma _T`$, the latter is unable to influence any spin motion in a macroscopic sample. This conclusion was confirmed by accurate calculations \[12–16\].
(10) Resonator relaxation rate arises when the sample is coupled to a resonant electric circuit. Then in the effective force (12) the magnetic field $`H`$ is the resonator feedback field. The role of this field will be thoroughly studied in what follows.
Summarizing all said above, for the spin averages (29) to (31), we obtain the evolution equations
$$\frac{du}{dt}=i(\omega _s+\xi _0i\mathrm{\Gamma }_2)u+fs,$$
(41)
$$\frac{dw}{dt}=2\mathrm{\Gamma }_2w+\left(u^{}f+f^{}u\right)s,$$
(42)
$$\frac{ds}{dt}=\frac{1}{2}\left(u^{}f+f^{}u\right)\mathrm{\Gamma }_1(s\zeta ),$$
(43)
supplemented by the initial conditions
$$u(0)=u_0,w(0)=w_0,s(0)=s_0.$$
In these equations, $`\zeta `$ is a stationary spin polarization, the characteristic spin frequency is
$$\omega _s=\omega _0\omega _Ds,$$
(44)
with $`\omega _0`$ given by Eq. (13) and $`\omega _D`$, by Eq. (19). The total longitudinal rate $`\mathrm{\Gamma }_1`$ is defined in Eq. (32) and the total transverse rate $`\mathrm{\Gamma }_2`$, in Eq. (35). The effective force is
$$f=\frac{i}{\mathrm{}}\mu _0(B_1+H)+\xi +\gamma _ru,$$
(45)
where the last term is the friction force due to the interaction through magnetodipole radiation, and $`\gamma _r`$ is the magnetodipole radiation rate (37). Equations (41) to (43) are stochastic differential equations, since they contain the random variables $`\xi _0`$ and $`\xi `$, whose stochastic averages are given in Eqs. (26). The external transverse field $`B_1`$ and the resonator feedback field $`H`$ need yet to be specified.
## V Resonator Feedback Field
The resonator feedback field $`H`$ is created by the electric current of the coil surrounding the spin sample. We assume that the coil axis is along the axis $`x`$. The electric circuit is characterized by resistance $`R`$, inductance $`L`$, and capacity $`C`$. The spin sample is inserted into a coil of $`n`$ turns, length $`l`$, cross-section area $`A_c`$, and volume $`V_c=A_cl`$. The electric current in the circuit is described by the Kirchhoff equation
$$L\frac{dj}{dt}+Rj+\frac{1}{C}_0^tj(t^{})𝑑t^{}=E_f\frac{d\mathrm{\Phi }}{dt},$$
(46)
in which $`E_f`$ is an electromotive force, if any, and the magnetic flux
$$\mathrm{\Phi }=\frac{4\pi }{c}nA_c\eta m_x,$$
(47)
where $`\eta V/V_c`$ is a filling factor, is formed by the $`x`$-component of the magnetization density
$$m_x\frac{\mu _0}{V}\underset{i}{}<S_i^x>.$$
(48)
The electric current, circulating over the coil, creates a magnetic field
$$H=\frac{4\pi n}{cl}j.$$
(49)
The circuit natural frequency is
$$\omega \frac{1}{\sqrt{LC}}\left(L4\pi \frac{n^2A_c}{c^2l}\right)$$
(50)
and the circuit damping is
$$\gamma \frac{1}{\tau }=\frac{R}{2L}=\frac{\omega }{2Q},$$
(51)
where $`\tau `$ is called the circuit ringing time and $`Q\omega L/R`$ is the quality factor. Also, let us define the reduced electromotive force
$$e_f\frac{cE_f}{nA_c\gamma }.$$
(52)
Then the Kirchhoff equation (46) can be transformed to the equation
$$\frac{dH}{dt}+2\gamma H+\omega ^2_0^tH(t^{})𝑑t^{}=\gamma e_f4\pi \eta \frac{dm_x}{dt}$$
(53)
for the feedback magnetic field created by the coil.
The feedback equation (53) can be represented in another equivalent form that proved to be very convenient for defining the feedback field \[12–15\]. For this purpose, we involve the method of Laplace transforms and introduce the transfer function
$$G(t)=\left(\mathrm{cos}\omega ^{}t\frac{\gamma }{\omega ^{}}\mathrm{sin}\omega ^{}t\right)e^{\gamma t},$$
(54)
where
$$\omega ^{}\sqrt{\omega ^2\gamma ^2}.$$
Thus, we transform the feedback-field equation (53) to the integral representation
$$H=_0^tG(tt^{})\left[\gamma e_f(t^{})4\pi \eta \dot{m}_x(t^{})\right]𝑑t^{},$$
(55)
in which
$$\dot{m}_x(t)\frac{1}{2}\rho \mu _0S\frac{d}{dt}(u^{}+u).$$
(56)
Let the resonant part of the reduced electromotive force (52) be
$$e_f(t)=h_2\mathrm{cos}\omega t.$$
(57)
And let us introduce the notation
$$\nu _2\frac{\mu _0h_2}{2\mathrm{}}.$$
(58)
As usual, we assume that all attenuation parameters are much smaller than the characteristic spin frequency $`\omega _s`$. Then Eq. (55) can be solved by an iteration procedure, which in first order gives
$$\frac{\mu _0H}{\mathrm{}}=i(\alpha u\alpha ^{}u^{})+2\beta \mathrm{cos}\omega t.$$
(59)
Here the coupling function
$$\alpha =\gamma _0\omega _s\left[\frac{1\mathrm{exp}\{i(\omega \omega _s)t\gamma t\}}{\gamma +i(\omega \omega _s)}+\frac{1\mathrm{exp}\{i(\omega +\omega _s)t\gamma t\}}{\gamma i(\omega +\omega _s)}\right]$$
(60)
describes the coupling of spins with the resonator and the function
$$\beta =\frac{\nu _2}{2}\left(1e^{\gamma t}\right)$$
(61)
characterizes the action of the resonator electromotive force on spins. In Eq. (60) the notation for the natural spin width
$$\gamma _0\frac{\pi }{\mathrm{}}\eta \rho \mu _0^2S$$
(62)
is employed.
The spin-resonator coupling can be characterized by the dimensionless coupling parameter
$$g\frac{\gamma \gamma _0\omega _s}{\gamma _2(\gamma ^2+\mathrm{\Delta }^2)},$$
(63)
in which $`\mathrm{\Delta }\omega |\omega _s|`$ is the detuning. As is evident from Eq. (60), an efficient spin-resonator coupling is possible only when the detuning from the resonance is small, such that
$$\frac{|\mathrm{\Delta }|}{\omega }1(\mathrm{\Delta }\omega |\omega _s|).$$
(64)
When the resonance is sufficiently sharp, so that $`|\mathrm{\Delta }|<\gamma `$, then the coupling function (60) reduces to
$$\alpha =g\gamma _2\left(1e^{\gamma t}\right).$$
(65)
Thus, the resonator feedback field $`H`$ is defined by Eq. (59), in which $`\alpha `$ is given by Eq. (65) and $`\beta `$, by Eq. (61).
## VI Averaged Evolution Equations
The resonator field, defined in Eq. (59), has to be substituted in the effective force (45) entering the evolution equations (41) to (43). In Eq. (45), we also need to specify the external magnetic field $`B_1`$. In general, the latter may contain a constant part and an alternating term. So, let us take this transverse field in the form
$$B_1=h_0+h_1\mathrm{cos}\omega t.$$
(66)
In what follows, we shall use the notation
$$\nu _0\frac{\mu _0h_0}{\mathrm{}},\nu _1\frac{\mu _0h_1}{2\mathrm{}}.$$
(67)
Equations (41) to (43) are stochastic differential equations, containing the random variables $`\xi _0`$ and $`\xi `$ describing local spin fluctuations. In order to derive the evolution equations in terms of ordinary differential equations, we have to accomplish the averaging over random fluctuations. This can be done by following the scale separation approach \[11–16\], the usage of the stochastic averages (26), and by invoking the known techniques of treating stochastic variables .
Keeping in mind that the attenuation parameters are substantially smaller than the characteristic spin frequency $`\omega _s`$, we notice from Eqs. (41) to (43) that the function $`u`$ can be classified as fast, being compared with the temporal behaviour of the functions $`w`$ and $`s`$. The latter play the role of temporal quasi-invariants with respect to $`u`$.
Fist, we substitute into Eqs. (41) to (43) the effective force (45), the resonator field (59), and the transverse magnetic field (66). This results in the equations
$$\frac{du}{dt}=i(\omega _s+\xi _0)u(\mathrm{\Gamma }_2\alpha s\gamma _rs)u+f_1s\alpha su^{},$$
(68)
$$\frac{dw}{dt}=2(\mathrm{\Gamma }_2\alpha s\gamma _rs)w+\left(u^{}f_1+f_1^{}u\right)s\alpha s\left(u^2+(u^{})^2\right),$$
(69)
$$\frac{ds}{dt}=(\alpha +\gamma _r)w\frac{1}{2}\left(u^{}f_1+f_1^{}u\right)\mathrm{\Gamma }_1(s\zeta )+\frac{1}{2}\alpha \left(u^2+(u^{})^2\right),$$
(70)
in which
$$f_1i\nu _02i(\nu _1+\beta )\mathrm{cos}\omega t+\xi .$$
(71)
Then we solve Eq. (68) for the fast variable $`u`$, keeping the quasi-invariants fixed, which yields
$$u=u_0\mathrm{exp}\left\{(i\omega _s+\mathrm{\Gamma }_2\alpha s\gamma _rs)ti_0^t\xi _0(t^{})𝑑t^{}\right\}+$$
$$+s_0^tf_1(t^{})\mathrm{exp}\left\{(i\omega _s+\mathrm{\Gamma }_2\alpha s\gamma _rs)(tt^{})i_t^{}^t\xi _0(t^{\prime \prime })𝑑t^{\prime \prime }\right\}𝑑t^{}.$$
(72)
Solution (72) must be substituted in Eqs. (69) and (70) for the slow functions $`w`$ and $`s`$. After this, the latter equations have to be averaged over time and over the stochastic variables $`\xi _0`$ and $`\xi `$, again keeping the quasi-invariants fixed. To slightly simplify the resulting equations, one can take the initial condition for the transition function $`u`$ in the real form, such that $`u_0^{}=u_0`$, which is not principal but just makes the equations less cumbersome.
To present the resulting equations in a compact form, we introduce the effective attenuation
$$\mathrm{\Gamma }_3\gamma _3+\frac{\nu _0^2\mathrm{\Gamma }}{\omega _s^2+\mathrm{\Gamma }^2}\frac{\nu _0(\nu _1+\beta )\mathrm{\Gamma }}{\omega _s^2+\mathrm{\Gamma }^2}e^{\mathrm{\Gamma }t}+\frac{(\nu _1+\beta )^2\mathrm{\Gamma }}{\mathrm{\Delta }^2+\mathrm{\Gamma }^2}\left(1e^{\mathrm{\Gamma }t}\right),$$
(73)
in which
$$\mathrm{\Gamma }\mathrm{\Gamma }_2+\gamma _3(\alpha +\gamma _r)s.$$
(74)
And finally, after the described averaging, we obtain the evolution equations
$$\frac{dw}{dt}=2(\mathrm{\Gamma }_2\alpha s\gamma _rs)w+2\mathrm{\Gamma }_3s^2,$$
(75)
$$\frac{ds}{dt}=(\alpha +\gamma _r)w\mathrm{\Gamma }_3s\mathrm{\Gamma }_1(s\zeta ).$$
(76)
These equations are very general. They include various attenuation processes, described in Sec. IV, and take into account transverse constant and alternating fields (66), as well as the resonator electromotive force (57) entering through function (61). The resonator feedback field is responsible for the appearance of the coupling function (65). Notice that the radiation relaxation rate $`\gamma _r`$, defined in Eq. (37), enters everywhere together with the spin-resonator coupling $`\alpha `$. However their values are drastically different. Since
$$\frac{\gamma _r}{\alpha }0.1\frac{\gamma }{\omega _s}(kL_s)^31,$$
the value of $`\gamma _r`$ is so incomparably smaller than $`\alpha g\gamma _2`$, that it is evident, in the presence of a resonator, the rate $`\gamma _r`$ must be forgotten.
Moreover, even when there is no resonator, so that $`\alpha =\beta =0`$, the radiation rate $`\gamma _r`$ plays no role, since it is much smaller than $`\gamma _2`$, $`\gamma _2^{}`$, and $`\gamma _3`$. One might think that $`\gamma _r`$ could play a role in the following unrealistic case. Let us imagine an absolutely ideal lattice with no inhomogeneous broadening, that is, let us set $`\gamma _2^{}=0`$, which is certainly a purely imaginary situation. Then, according to Eq. (35), one has $`\mathrm{\Gamma }_2=\gamma _2(1s^2)`$. Assume that the spin system is completely polarized, with $`s_0=1`$. Hence, at the initial time, $`\mathrm{\Gamma }_2=0`$. Could then the spin motion be started by the term with $`\gamma _r`$? The answer is evident: As far as the largest terms in both Eqs. (75) and (76) are those containing $`\mathrm{\Gamma }_3`$, the terms with $`\gamma _r`$ are always negligible, even if $`\mathrm{\Gamma }_2=0`$. Even more, functions (30) and (31), by their definition, satisfy the inequality
$$w+s^21.$$
(77)
Therefore, if one sets $`s_0=1`$, then $`w_0=0`$, and the term $`\gamma _rw`$ simply disappears from the equations. Vice versa, if one sets a noticeable $`w_01`$, then $`s^21`$, and $`\mathrm{\Gamma }_2\gamma _2\gamma _r`$. In this way, the radiation rate $`\gamma _r`$ never plays any role in the spin motion, which is in agreement with the estimates by Bloembergen .
Note that the situation in spin systems is principally different from that happening in atomic systems. In the latter, both the linewidth $`\gamma _2=2|𝐝|^2k^3/3`$ as well as the collective radiation rate $`\gamma _r=(2/3)|𝐝|^2k^3N_c`$, where $`N_c`$ is the number of correlated atoms, forming a wave packet, are caused by the same physical process, by the interaction of atoms with their radiation field. Hence $`\gamma _r/\gamma _2=N_c1`$, which results in the coherentization of the dipole transitions. This is possible even if $`kL1`$, but the number of atoms in a partial wave packet is $`N_c1`$, since $`\gamma _r/\gamma _2=N_c1`$. Contrary to this, in spin systems the linewidth $`\gamma _2`$, given in Eq. (33), is due to direct dipole-dipole interactions, while the radiation rate (37) is a result of the spin interactions with their radiation field. This is why in the latter case, one always has $`\gamma _r\gamma _2`$, and the radiation rate $`\gamma _r`$ plays no part in the motion of spins.
We may also notice that in the effective attenuation (73) the terms due to the presence of a constant transverse field are less important than the terms caused by the local spin fluctuations and by the alternating transverse fields. Therefore, omitting the terms corresponding to a permanent transverse magnetic field, we have
$$\mathrm{\Gamma }_3=\gamma _3+\frac{(\nu _1+\beta )^2\mathrm{\Gamma }}{\mathrm{\Gamma }^2+\mathrm{\Delta }^2}\left(1e^{\mathrm{\Gamma }t}\right).$$
(78)
Finally, we obtain the evolution equations
$$\frac{dw}{dt}=2(\mathrm{\Gamma }_2\alpha s)w+2\mathrm{\Gamma }_3s^2,$$
(79)
$$\frac{ds}{dt}=\alpha w\mathrm{\Gamma }_3s\mathrm{\Gamma }_1(s\zeta ),$$
(80)
describing the averaged motion of spins.
## VII Coherence Emerging from Chaos
One of the most intriguing questions is how the spin motion could become coherent if initially it was not. This is a particular case of the general physical problem of how coherence emerges from chaos.
Being interested in a self-organized process of arising coherence, let us consider the case, when there are no external transverse fields pushing spins, that is $`\nu _1=\beta =0`$. Then Eq. (78) yields $`\mathrm{\Gamma }_3=\gamma _3`$. Assume also that there is no pumping, so that $`\gamma _1^{}=0`$, hence $`\mathrm{\Gamma }_1=\gamma _1`$. Under these conditions, the initial spin motion, for the time $`t`$ such that
$$\gamma _1t1,\gamma _2t1,\gamma _3t1,$$
(81)
follows from Eqs. (79) and (80) in the form
$$ww_0+2\left[\gamma _3s_0^2\gamma _2\left(1s_0^2+\kappa \right)w_0\right]t,$$
$$ss_0[(\gamma _1+\gamma _3)s_0\gamma _1\zeta ]t,$$
(82)
where the inhomogeneity coefficient is introduced,
$$\kappa \gamma _2^{}/\gamma _2.$$
(83)
If at the initial time no transverse polarization is imposed on the system, and the initial coherence function is zero, $`w_0=0`$, nevertheless the coherent spin motion starts developing according to the law
$$w2\gamma _3s_0^2t(w_0=0),$$
(84)
provided there is an initial longitudinal polarization $`s_00`$. The initiation of the emerging coherent motion is caused by local spin fluctuations creating the effective rate $`\gamma _3`$. Recall that in the Bloch equations coherent motion never appears if it is not imposed by the initial conditions. Contrary to this, Eqs. (79) and (80) take into account the local spin fluctuations triggering the motion of spins. The second of Eqs. (82), keeping in mind that usually $`\gamma _1\gamma _3\gamma _2`$, can be simplified to
$$ss_0(1\gamma _3t).$$
(85)
At the initial stage of spin motion, their coherence is yet incipient, and the motion is mainly governed by quantum chaotic spin fluctuations. The coherentization of the transverse motion goes through the resonator feedback field and the growing coupling function (65). The quantitative change in the spin motion happens when the coupling function (65) becomes so large that the term $`(\mathrm{\Gamma }_2\alpha s)`$ in Eq. (79) goes negative, which means that an efficient generation of coherence has started in the system. This is analogous to the beginning of maser generation \[15–19\]. The moment of time, when the regime of mainly chaotic quantum fluctuations transforms into the regime of predominantly coherent spin motion, can be called the chaos time. This time $`t_c`$ is defined by the equality $`\alpha s=\mathrm{\Gamma }_2`$, that is by the equation
$$\alpha s=\gamma _2(1s^2)+\gamma _2^{}(t=t_c).$$
(86)
From here, the estimate for the chaos time is
$$t_c=\tau \mathrm{ln}\frac{gs_0}{gs_01+s_0^2\kappa },$$
(87)
where $`\tau `$ is the resonator ringing time defined in Eq. (51). The regime of chaotic spin fluctuations lasts till the chaos time (87), after which the coherent stage of spin motion comes into play. As is clear from the above equations, the transformation from the chaotic to coherent regime goes as a gradual crossover. Notice that the quantity $`1s_0^2+\kappa `$ is positive since $`s_0^21`$. Then, in order that the chaos time (87) be positive and finite, the inequality
$$gs_0>1s_0^2+\kappa >0$$
(88)
must hold. For a strong spin-resonator coupling, when $`gs_01`$, the chaos time (87) reduces to
$$t_c\frac{\tau }{gs_0}\left(1s_0^2+\kappa \right).$$
(89)
As is seen, there exists a well defined stage of chaotic spin fluctuations, with a finite chaos time $`t_c>0`$, after which the coherent regime develops, if $`gs_0>0`$. The coupling parameter $`g`$ is defined in Eq. (63), from which it follows that one should have $`\omega _ss_0>0`$. Assuming that the initial spin polarization is positive, $`s_0>0`$, one gets the requirement that $`\omega _s>0`$. The latter, by definition (44), is equivalent to the condition $`\omega _0>\omega _Ds`$. Moreover, the coupling function (65) is obtained under the resonance condition (64), which implies that $`\omega _s`$ has to be close to the resonator natural frequency $`\omega `$. There are two ways of preserving the resonance condition (64). First, one can impose a sufficiently strong external magnetic field $`B_0`$, such that the frequency $`\omega _0`$, given by Eq. (13), would be much larger than $`\omega _D`$, defined in Eq. (19). This becomes trivial for $`S=1/2`$, when $`\omega _D=0`$. If $`\omega _0\omega _D`$, then it is easy to realize the resonance condition (64), with $`\omega _s\omega `$ and slightly varying in time detuning $`\mathrm{\Delta }=\omega \omega _s`$.
The second way of keeping the resonance condition (64) is by means of the chirping effect . This requires to vary in time the external magnetic field $`B_0`$ so that to maintain the equality
$$\frac{\mu _0B_0}{\mathrm{}}+(\omega +\omega _Ds)=\mathrm{\Delta },$$
(90)
with a fixed detuning.
## VIII Coherent Spin Relaxation
After the chaos time (87), the motion of spins becomes more and more coherent, being collectivized by the resonator feedback field, with the coupling function $`\alpha `$ reaching the value $`g\gamma _2`$. At the transient stage, when $`t>t_c`$ but $`tT_1`$, we may neglect the term with $`\gamma _1`$ in Eq. (80). Assuming that there is no pumping, that is $`\gamma _1^{}=0`$, one has $`\mathrm{\Gamma }_1=\gamma _1`$. Let us continue studying the case of the self-organized coherent spin motion, when there are no transverse external fields, so that $`\nu _1=\beta =0`$, hence $`\mathrm{\Gamma }_3=\gamma _3`$. When the coherence is well developed, then the main term in Eq. (79) is the first one, while the term with $`\gamma _3`$ can be neglected. Under these conditions, and using expression (35) for the rate $`\mathrm{\Gamma }_2`$, Eqs. (79) and (80) reduce to the form
$$\frac{dw}{dt}=2\gamma _2\left(1s^2+\kappa gs\right)w,$$
(91)
$$\frac{ds}{dt}=g\gamma _2w.$$
(92)
The solution of these equations is explained in Appendix C and it yields
$$w=\left(\frac{\gamma _p}{g\gamma _2}\right)^2\mathrm{sech}^2\left(\frac{tt_0}{\tau _p}\right),$$
$$s=\frac{\gamma _p}{g\gamma _2}\mathrm{tanh}\left(\frac{tt_0}{\tau _p}\right)+\frac{1+\kappa }{g}.$$
(93)
Here
$$\tau _p1/\gamma _p$$
(94)
is the pulse time showing the duration of the coherent relaxation occurring as a fast pulse. The delay time
$$t_0=t_c+\frac{\tau _p}{2}\mathrm{ln}\left|\frac{\gamma _p+\gamma _g}{\gamma _p\gamma _g}\right|$$
(95)
defines the time of the maximal coherence. The pulse width is given by the relation
$$\gamma _p^2=\frac{1}{2}\gamma _g^2\left[1+\sqrt{1+4\left(\frac{g\gamma _2}{\gamma _g}\right)^2w_c}\right],$$
(96)
in which
$$\gamma _g\gamma _2(gs_c1\kappa ).$$
(97)
The boundary values $`w_c`$ and $`s_c`$ are
$$w_c=w_0+2\left[\gamma _3s_0^2\gamma _2\left(1s_0^2+\kappa \right)w_0\right]t_c,s_c=s_0(1\gamma _3t_c),$$
(98)
with the chaos time $`t_c`$ given in Eq. (87). Since we are interested in the self-organized collective process, when there is no large transverse polarization imposed on the system at the initial time, we may set $`w_0s_0^2`$. Then Eq. (96) simplifies to
$$\gamma _p^2=\gamma _g^2+(g\gamma _2)^2w_c.$$
(99)
The pulse time (94) reads as
$$\tau _p=\frac{T_2}{\sqrt{(gs_c1\kappa )^2+g^2w_c}}.$$
(100)
It is easy to notice that if the spin-resonator coupling is weak, $`g1`$, then $`\gamma _p\gamma _g\gamma _2`$ and $`\tau _pT_2`$. In that case, no self-organized coherence can arise in the system.
Collective coherent effects appear in the spin motion only if the pulse time $`\tau _p`$ is smaller than the dephasing time $`T_2`$. The inequality $`\tau _p<T_2`$, according to Eq. (100), requires that
$$(gs_c1\kappa )^2+g^2w_c>1.$$
(101)
Three different regimes can satisfy Eq. (101).
The regime of collective induction happens when
$$gs_0<1+\kappa ,g^2w_0>1.$$
(102)
Then, as is clear from Eq. (97), one has $`\gamma _g<0`$, because of which $`t_0<t_c`$. This means that there is no a noticeable maximum in the coherence function $`w`$, since, by definition, the delay time (95) should occur after the chaotic stage, so that $`t_0>t_c`$. But the latter implies that $`\gamma _g>0`$.
The triggered coherent relaxation corresponds to
$$gs_0>1+\kappa ,0<g^2w_0<1.$$
(103)
And the purely self-organized coherent relaxation takes place when
$$gs_0>2+\kappa ,w_0=0.$$
(104)
In this classification, we keep in mind the inequality $`\gamma _3t_c1`$, owing to which $`w_cw_0`$ and $`s_cs_0`$. The initial coherence is assumed to be weak, so that $`w_01`$.
For $`w_0s_0^2`$, the delay time (95) can be represented as
$$t_0=t_c+\frac{\tau _p}{2}\mathrm{ln}\frac{4(gs_c1\kappa )^2}{g^2w_c}.$$
(105)
In the case of the purely self-organized coherent relaxation, for sufficiently large coupling and initial polarization, such that $`gs_01`$, the delay time (105) reduces to
$$t_0=t_c+\frac{\tau _p}{2}\mathrm{ln}\left|\frac{2}{\gamma _3t_c}\right|,$$
(106)
where $`\tau _p=T_2/gs_0`$. From these formulas, one sees that if $`\gamma _30`$, then $`t_0\mathrm{}`$, and no coherent relaxation is possible. This emphasizes the crucial role of the local spin fluctuations, whose existence results in the relaxation rate $`\gamma _3`$.
At the delay time (95), solutions (93) are given by the expressions
$$w(t_0)=w_c+\left(s_c\frac{1+\kappa }{g}\right)^2,s(t_0)=\frac{1+\kappa }{g}.$$
(107)
And for $`tt_0`$, they exponentially decay to the values
$$w4w(t_0)\mathrm{exp}(2\gamma _pt),$$
$$ss_c+\frac{2}{g}(1+\kappa )+2\left(s_c\frac{1+\kappa }{g}\right)\mathrm{exp}(2\gamma _pt).$$
(108)
At very large times $`tT_1`$, the transient equations (91) and (92) are no longer valid. Then one has to return to the full equations (79) and (80). With increasing time, the solutions tend to the stationary points defined by the zeros of the right-hand sides of these equations. Among the relaxation regimes to the stationary solutions, one is especially interesting, going through a long series of coherent pulses. This pulsing coherent relaxation takes place under a permanent external pumping described by a large pumping rate $`\gamma _1^{}\gamma _1`$. Then $`\mathrm{\Gamma }_1=\gamma _1^{}`$. If also the coupling parameter is sufficiently large, such that $`g\zeta 1`$ and
$$\frac{\gamma _3}{g\zeta \gamma _1^{}}1,$$
then the fixed point of Eqs. (79) and (80) is given by the expressions
$$w^{}=\frac{\gamma _1^{}}{g(\gamma _2+\gamma _2^{})}\left(1\frac{\gamma _3}{g\zeta \gamma _1^{}}\right),s^{}=\frac{1}{g}\left(1\frac{\gamma _3}{g\zeta \gamma _1^{}}\right),$$
(109)
corresponding to a stable focus. The relaxation to the stationary solutions (109) realizes through a series of sharp coherent pulses, similar to the form of Eqs. (93), with the temporal interval between the pulses asymptotically defined by the separation time
$$T_{sep}=\frac{2\pi }{\sqrt{2g\zeta \gamma _1^{}(\gamma _2+\gamma _2^{})}}.$$
(110)
The number of the separate coherent pulses can be estimated as $`N_{sep}=1/\gamma _1^{}T_{sep}`$, which gives
$$N_{sep}=\sqrt{\frac{g\zeta (\gamma _2+\gamma _2^{})}{2\pi ^2\gamma _1^{}}}.$$
Such a highly nontrivial relaxation regime occurs only under a strong pumping and a sufficiently strong coupling with a resonator.
## IX Influence of Cross Correlations
When in the sample, in addition to the studied spins, there are spins of other nature, the presence of the latter can certainly influence the dynamics of the former. Let us consider the case of two types of coexisting spins, $`S`$ and $`F`$. The total Hamiltonian is the sum
$$\widehat{H}=\widehat{H}_S+\widehat{H}_F+\widehat{H}_{SF}$$
(111)
of the Hamiltonians for $`S`$-spins, $`F`$-spins and their interactions. The Hamiltonian $`\widehat{H}_S`$ of $`S`$-spins is the same as in Eqs. (1) to (4). Let us accept for the Hamiltonian $`\widehat{H}_F`$ of $`F`$-spins a similar general form
$$\widehat{H}_F=\underset{i}{}\widehat{H}_{iF}+\frac{1}{2}\underset{ij}{}\widehat{H}_{ijF}.$$
(112)
The single-spin terms are
$$\widehat{H}_{iF}=\mu _{0F}𝐁𝐅_iD_F(F_i^z)^2,$$
(113)
with the total magnetic field (3). And the interaction terms are given by
$$\widehat{H}_{ijF}=\underset{\alpha \beta }{}D_{ijF}^{\alpha \beta }F_i^\alpha F_j^\beta J_{ijF}𝐅_i𝐅_j,$$
(114)
with the dipolar tensor
$$D_{ijF}^{\alpha \beta }=\frac{\mu _{0F}^2}{r_{ij}^3}\left(\delta _{\alpha \beta }3n_{ij}^\alpha n_{ij}^\beta \right).$$
Assume that the interactions between the $`S`$-and $`F`$-spins are represented by the Hamiltonian
$$\widehat{H}_{SF}=\underset{i}{}A𝐒_i𝐅_i+\underset{ij}{}\underset{\alpha \beta }{}A_{ij}^{\alpha \beta }S_i^\alpha F_j^\beta ,$$
(115)
containing the part of the single-site interactions of intensity $`A`$ and the part of the dipole interactions, with the dipolar tensor
$$A_{ij}^{\alpha \beta }=\frac{\mu _0\mu _{0F}}{r_{ij}^3}\left(\delta _{\alpha \beta }3n_{ij}^\alpha n_{ij}^\beta \right).$$
In particular, these could be hyperfine interactions between nuclear and electron spins .
We employ notation (8) for the interaction parameters of $`S`$-spins and an equivalent notation for the interaction parameters $`a_{ijF}`$, $`b_{ijF}`$, and $`c_{ijF}`$ of $`F`$-spins. Similarly, we define the interaction parameters
$$\overline{a}_{ij}A_{ij}^{zz},\overline{b}_{ij}\frac{1}{4}\left(A_{ij}^{xx}A_{ij}^{yy}2iA_{ij}^{xy}\right),\overline{c}_{ij}\frac{1}{2}\left(A_{ij}^{xz}iA_{ij}^{yz}\right)$$
(116)
for the spin cross interactions.
The local fields (11), acting on $`S`$-spins, are generalized to the form
$$\xi _0\frac{1}{\mathrm{}}\underset{j(i)}{}\left[a_{ij}S_j^z+c_{ij}^{}S_j^{}+c_{ij}S_j^++J_{ij}(S_i^zS_j^z)+\overline{a}_{ij}F_j^z+\overline{c}_{ij}^{}F_j^{}+\overline{c}_{ij}F_j^+\right],$$
$$\xi \frac{i}{\mathrm{}}\underset{j(i)}{}[2c_{ij}S_j^z\frac{1}{2}a_{ij}S_j^{}+2b_{ij}S_j^++J_{ij}(S_i^{}S_j^{})+$$
$$+2\overline{c}_{ij}F_j^z\frac{1}{2}\overline{a}_{ij}F_j^{}+2\overline{b}_{ij}F_j^+].$$
(117)
Analogous local fields act on $`F`$-spins,
$$\xi _{0F}\frac{1}{\mathrm{}}\underset{j(i)}{}\left[a_{ijF}F_j^z+c_{ijF}^{}F_j^{}+c_{ijF}F_j^++J_{ijF}(F_i^zF_j^z)+\overline{a}_{ij}S_j^z+\overline{c}_{ij}^{}S_j^{}+\overline{c}_{ij}S_j^+\right],$$
$$\xi _F\frac{i}{\mathrm{}}\underset{j(i)}{}[2c_{ijF}F_j^z\frac{1}{2}a_{ijF}F_j^{}+2b_{ijF}F_j^++J_{ijF}(F_i^{}F_j^{})+$$
$$+2\overline{c}_{ijF}S_j^z\frac{1}{2}\overline{a}_{ij}S_j^{}+2\overline{b}_{ij}S_j^+].$$
(118)
Instead of one effective force (12), we have now two forces
$$f\frac{i}{\mathrm{}}\mu _0(B_1+H)+\frac{i}{\mathrm{}}A_iF_i^{}+\xi ,$$
$$f_F\frac{i}{\mathrm{}}\mu _{0F}(B_1+H)+\frac{i}{\mathrm{}}A_iS_i^{}+\xi _F.$$
(119)
In addition to frequency (13), let us introduce the effective frequencies
$$\omega _{0F}\frac{\mu _{0F}}{\mathrm{}}B_0,\epsilon \frac{A}{\mathrm{}}.$$
(120)
The Heisenberg equations of motion for the system with Hamiltonian (111) yield the equations for $`S`$-spins
$$\frac{dS_i^{}}{dt}=i\left(\omega _0+\epsilon F_i^z+\xi _0\right)S_i^{}+S_I^zf+\frac{i}{\mathrm{}}D\left(S_i^{}S_i^z+S_i^zS_i^{}\right),$$
$$\frac{dS_i^z}{dt}=\frac{1}{2}\left(f^+S_i^{}+S_i^+f\right),$$
(121)
and the equations for $`F`$-spins
$$\frac{dF_i^{}}{dt}=i\left(\omega _{0F}+\epsilon S_i^z+\xi _{0F}\right)F_i^{}+F_I^zf_F+\frac{i}{\mathrm{}}D_F\left(F_i^{}F_i^z+F_i^zF_i^{}\right),$$
$$\frac{dF_i^z}{dt}=\frac{1}{2}\left(f_F^+F_i^{}+F_i^+f_F\right).$$
(122)
Again we assume that the sample is inserted into the coil of a resonant electric circuit. The feedback field acting on the sample is given by Eq. (53) or (55), where now the magnetic-moment density is
$$m_x=\frac{\mu _0}{V}\underset{i=1}{\overset{N}{}}<S_i^x>+\frac{\mu _{0F}}{V}\underset{j=1}{\overset{N_F}{}}<F_j^x>,$$
(123)
with $`N_F`$ being the number of $`F`$-spins.
Averaging Eqs. (121) and (122), we derive the evolution equations for functions (29), (30), and (31), corresponding to $`S`$-spins, as well as the equations for the functions
$$u_F\frac{1}{FN_F}\underset{i=1}{\overset{N_F}{}}<F_i^{}>,$$
(124)
$$w_F\frac{1}{F^2N_F(N_F1)}\underset{ij}{\overset{N_F}{}}<F_i^+F_j^{}>,$$
(125)
$$s_F\frac{1}{FN_F}\underset{i=1}{\overset{N_F}{}}<F_i^z>,$$
(126)
describing $`F`$-spins. In this notation, the transverse magnetic-moment density (123) is
$$m_x=\frac{1}{2}\rho \mu _0S(u^{}+u)+\frac{1}{2}\rho _F\mu _{0F}F(u_F^{}+u_F),$$
where $`\rho _F`$ is the density of $`F`$-spins.
The analysis of the evolution equations for the combined system of $`S`$\- and $`F`$-spins is the same as has been given above for one type of spins $`S`$, with the difference that all expressions become much more cumbersome. Again it is possible to show that in the triggering of spin motion an important role is played by the coupled $`SF`$ spin fluctuations, which yield the dynamic relaxation rates $`\gamma _3`$ and $`\gamma _{3F}`$ defined by the relations
$$\gamma _3^2=\gamma _{SS}^2+\gamma _{SF}^2,\gamma _{3F}^2=\gamma _{FF}^2+\gamma _{FS}^2,$$
(127)
where
$$\gamma _{SS}\rho \frac{\mu _0^2}{\mathrm{}}\sqrt{S(S+1)},\gamma _{SF}\sqrt{\rho \rho _F}\frac{\mu _0\mu _{0F}}{\mathrm{}}F,$$
$$\gamma _{FF}\rho _F\frac{\mu _{0F}^2}{\mathrm{}}\sqrt{F(F+1)},\gamma _{FS}\sqrt{\rho \rho _F}\frac{\mu _{0F}\mu _0}{\mathrm{}}S.$$
The effective frequencies of $`S`$\- and $`F`$-spins, respectively, are
$$\omega _S=\omega _0\omega _Ds+\epsilon s_FS,\omega _F=\omega _{0F}\omega _{DF}s_F+\epsilon sF,$$
(128)
where $`\omega _D`$ is given by Eq. (19) and
$$\omega _{DF}(2F1)\frac{D_F}{\mathrm{}}.$$
(129)
We shall not overload this paper by a detailed exposition of various cross correlations resulting from the complicated system of the coupled evolution equations for $`S`$\- and $`F`$-spins. Let us only emphasize the existence of a rather nontrivial nonlinear effect of mutual spin interactions through the resonator feedback field. Calculating the latter from the integral representation (55), with the transverse magnetic density (123), and substituting this into the evolution equations results in an effective mutual influence of spins through the feedback field. If the resonator is tuned to the characteristic frequency $`\omega _S`$ of $`S`$-spins, then for the latter, we derive the evolution equations similar to Eqs. (79) and (80), but with the effective spin-resonator coupling
$$g=\frac{\gamma \gamma _0\omega _S}{\gamma _2(\gamma ^2+\mathrm{\Delta }^2)}\left(1+\frac{\rho _F\mu _{0F}\epsilon s_FF}{\rho \mu _0\omega _F}\right),$$
(130)
instead of Eq. (63), and with $`\gamma _3`$ given by Eq. (127). Depending on the spin characteristics, coupling (130) can substantially surpass the value of Eq. (63). This is because the subsystem of $`F`$-spins, coupled to a resonator, becomes itself a kind of an additional resonator for $`S`$-spins.
## X Conclusion
A general theory is developed for describing nonlinear spin relaxation, which occurs when the spin system is prepared in a strongly nonequilibrium state and when the sample is coupled to a resonator electric circuit. A strongly nonequilibrium initial state can be realized by placing a polarized sample into an external magnetic field, whose direction is opposite to the sample magnetization. Nonlinearity in spin relaxation comes from direct spin-spin interactions and from their effective interactions through the resonator feedback field. Direct spin interactions are responsible for the appearance of local spin fluctuations, playing a crucial role at the starting stage of relaxation. The resonator feedback field collectivizes the spin motion, leading to coherent collective relaxation. The developed theory is based on a realistic Hamiltonian containing the main spin interactions. The role of various relaxation rates is thoroughly analysed.
The aim of the present paper has been to develop a general theory providing an accurate and realistic description of nonlinear spin relaxation. This theory can be employed for a large class of polarized spin materials. Applications to particular substances require a special consideration and separate publications. There exists a large variety of materials that can be treated by the developed theory. Just to give an example, we may mention the class of molecular magnets \[16,19,24–26\]. For instance, the molecular crystal V<sub>15</sub> is made of molecules of spin $`1/2`$, so has no magnetic anisotropy. Its nonlinear spin relaxation can be realized in a rather weak external field $`B_01`$ G. The molecules Mn<sub>12</sub> and Fe<sub>8</sub> possess the spin $`S=10`$. They form crystals with density $`\rho 10^{21}`$ cm<sup>-3</sup>. The anisotropy frequency is $`\omega _D10^{12}`$ s<sup>-1</sup>. At low temperatures below about 1 K, the molecules can be well polarized, with the spin-lattice relaxation parameters $`\gamma _110^510^7`$ s<sup>-1</sup>. The line width is caused by rather strong dipole interactions, with $`\gamma _210^{10}`$ s<sup>-1</sup>. The condition $`\omega _0>\omega _D`$ can be reached for $`B_0>10^5`$ G. In the molecular magnet, formed by the molecules Mn<sub>6</sub>, whose spin is $`S=12`$, the magnetic anisotropy is much weaker, with $`\omega _D10^{10}`$ s<sup>-1</sup>, being of the same order as $`\gamma _210^{10}`$ s<sup>-1</sup>. Therefore the required magnetic field is not high, $`B_0>10^3`$ G. Coupling a molecular crystal to a resonant circuit with the natural width $`\gamma \omega /2Q`$, where $`Q`$ is the resonator quality factor, one can attain the values of the coupling parameter as large as $`gQ10^4`$. With such a strong coupling, the influence of the resonator feedback field outperforms other relaxation mechanisms, producing fast coherent relaxation, with relaxation times $`\tau _p10^{13}`$ s. Such a fast reorientation of the magnetic moment can result in the emission of radiation pulses of high intensity.
Acknowledgements
I am grateful to E.P. Yukalova for helpful discussions. I appreciate the Mercator Professorship of the German Research Foundation.
Appendix A: Effective Homogeneous Broadening
The homogeneous broadening, existing in spin systems, arises from spin-spin interactions and is usually expressed through the moments $`M_n`$, which may depend on the level of the longitudinal polarization $`s`$, provided the latter is sufficiently large. The moments have been calculated in a number of works \[1–7,21\]. The most general and exact formula, relating the effective broadening with the moments, can be found in Abragam and Goldman , which for the Gaussian line shape is
$$\gamma _2(s)=\sqrt{\frac{\pi M_2^3(s)}{2[M_4(s)M_2^2(s)]}}.$$
The Lorentzian line shape yields to practically the same expression, with a slightly different coefficient. The broadening $`\gamma _2(s)`$ for the Lorentzian line is $`\sqrt{\pi }`$ of the Gaussian broadening. The dependence of the moments on the polarization has been accurately calculated , yielding
$$M_2(s)=M_2(0)(1s^2),M_4(s)=2.18M_2^2(0)(1s^2)(10.42s^2).$$
Substituting this into $`\gamma _2(s)`$, and taking into account that $`s^21`$, we obtain Eq. (34).
Appendix B: Spin Radiation Rate
To get a fully quantum-mechanical microscopic picture of spin interactions with electromagnetic field they radiate, one has to add to the spin Hamiltonian (1) the field Hamiltonian
$$\widehat{H}_f=\frac{1}{8\pi }\left(𝐄^2+𝐇^2\right)𝑑𝐫,$$
where $`𝐄=𝐄(𝐫,t)`$ is electric field and $`𝐇=𝐇(𝐫,t)`$ is magnetic field, and the operator energy of spin-field interactions
$$\widehat{H}_{sf}=\mu _0\underset{i=1}{\overset{N}{}}𝐒_i𝐇_i,$$
where $`𝐇_i=𝐇(𝐫_i,t)`$. From the Heisenberg equations of motion for the field variables, one finds the vector potential
$$𝐀(𝐫,t)=\frac{1}{c}𝐣(𝐫^{},t\frac{|𝐫𝐫^{}|}{c})\frac{d𝐫}{|𝐫𝐫^{}|},$$
in which the current density is
$$𝐣=c\mu _0\underset{i=1}{\overset{N}{}}𝐒_i\times \stackrel{}{}\delta (𝐫𝐫_i).$$
The vector potential $`𝐀_i𝐀(𝐫_i,t)`$ can be represented as
$$𝐀_i=𝐀_i^{}+𝐀_i^++𝐀_i^{},$$
where
$$𝐀_i^{}=\underset{j}{}\left(1+\frac{1}{c}\frac{}{t}\right)\frac{𝐫_{ij}}{r_{ij}^3}\times \stackrel{}{\mu }^{}S_j^{}\left(t\frac{r_{ij}}{c}\right),$$
$$𝐀_i^{}=\underset{j}{}\frac{𝐫_{ij}}{r_{ij}}\times \stackrel{}{\mu }_0S_j^z\left(t\frac{r_{ij}}{c}\right),$$
with the notation
$$\stackrel{}{\mu }\frac{\mu _0}{2}(𝐞_xi𝐞_y),\stackrel{}{\mu }_0\mu _0𝐞_z.$$
From here, we get the magnetic field $`𝐇_i𝐇(𝐫_i,t)`$ acting on an $`i`$-th spin as $`𝐇_i=\stackrel{}{}_i\times 𝐀_i`$, which gives the field
$$𝐇_i=𝐇_i^{}+𝐇_i^++𝐇_i^{},$$
in which
$$𝐇_i^{}=\underset{j}{}\left[\frac{\stackrel{}{\mu }^{}(\stackrel{}{\mu }^{}𝐧_{ij})𝐧_{ij}}{c^2r_{ij}}\frac{^2}{t^2}+\frac{\stackrel{}{\mu }^{}3(\stackrel{}{\mu }^{}𝐧_{ij})𝐧_{ij}}{r_{ij}^3}\left(1+\frac{r_{ij}}{c}\frac{}{t}\right)\right]S_j^{}\left(t\frac{r_{ij}}{c}\right),$$
$$𝐇_i^{}=\underset{j}{}\frac{\stackrel{}{\mu }_03(\stackrel{}{\mu }_0𝐧_{ij})𝐧_{ij}}{r_{ij}^3}S_j^z\left(t\frac{r_{ij}}{c}\right).$$
If the spins on different sites move independently of each other, so that the single-spin terms in the above sums chaotically oscillate, then the average magnetic field acting on each spin from the radiation of other spins is zero. Noticeable action of other spins can arise only if there exist the groups of spins, the so-called spin packets, which are strongly correlated, moving together. A substantial mutual interaction between spins, caused by their electromagnetic radiation, can appear only when this radiation is monochromatic, with a well-defined spin frequency $`\omega _s`$, the related wavelength $`\lambda =2\pi c/\omega _s`$, and wave vector $`k=\omega _s/c`$. This radiation can collectivize spins in a spin packet of size $`L_s`$, provided that
$$kL_s1.$$
When the radiation wavelength $`\lambda `$ is much larger than the system length $`L`$, then $`L_s=L`$. This, however, is not compulsory, and the size of a spin packet can be much shorter than $`L`$, but it should be much larger than the mean interspin distance. Thus, inequality (38) is a necessary condition for the appearance of collective effects.
Under condition (38), the above magnetic fields can be simplified, averaging them over spherical angles. The resulting expressions have to be added to the magnetic field in the effective force (12), which acquires one more term, being the friction force
$$f^{}=(\gamma _ri\delta \omega )u,$$
in which the collective radiation rate and frequency shift are
$$\gamma _r\gamma _0\underset{j}{\overset{N_s}{}}\frac{\mathrm{sin}(kr_{ij})}{kr_{ij}}\mathrm{\Theta }(ctr_{ij}),$$
$$\delta \omega \gamma _0\underset{j}{\overset{N_s}{}}\frac{\mathrm{cos}(kr_{ij})}{kr_{ij}}\mathrm{\Theta }(ctr_{ij}),$$
where
$$\gamma _0\frac{2}{3\mathrm{}}\mu _0^2Sk^3$$
is the single-spin natural width, $`\mathrm{\Theta }()`$ is a unit-step function, and $`N_s=\rho L_s^3`$ is the number of spins in a spin packet. These formulas can be further simplified to
$$\gamma _r=\gamma _0N_s=\frac{2}{3\mathrm{}}\mu _0^2Sk^3N_s$$
and
$$\delta \omega =\frac{3\gamma _r}{2kL_s}=\frac{1}{\mathrm{}}\rho \mu _0^2S(kL_s)^2.$$
The frequency shift is very small, even as compared to $`\gamma _2`$, since
$$\frac{\delta \omega }{\gamma _2}0.1(kL_s)^21.$$
Of course, such a small shift can be omitted, being negligible as compared to $`\gamma _2`$ and the more so as compared to $`\omega _s`$. And for the radiation rate $`\gamma _r`$, substituting there $`N_s=\rho L_s^3`$, we obtain Eq. (37).
Appendix C: Transient Stage of Relaxation
After the chaotic stage of spin fluctuations, the transient stage comes into play, characterized by Eqs. (91) and (92). The latter, by introducing the function
$$y\gamma _2\left(1s^2+\kappa gs\right)$$
and keeping in mind a sufficiently large coupling parameter $`gs`$, rearrange to
$$\frac{dw}{dt}=2yw,\frac{dy}{dt}=(g\gamma _2)^2w.$$
Differentiating the second of these equations, we have
$$\frac{d^2y}{dt^2}+2y\frac{dy}{dt}=0,$$
which yields
$$\frac{dy}{dt}+y^2=\gamma _p^2,$$
with $`\gamma _p`$ being an integration parameter. This Riccati equation possesses the solution
$$y=\gamma _p\mathrm{tanh}\left(\frac{tt_0}{\tau _p}\right),$$
in which $`\gamma _p\tau _p1`$ and $`t_0`$ is another integration constant. Inverting the dependence of $`y`$ on $`s`$ for $`s^21`$, we get
$$s=\frac{y}{g\gamma _2}+\frac{1+\kappa }{g}.$$
This gives the second of Eqs. (93), while the first of solutions (93) follows from Eq. (92). The integration constants $`\gamma _p`$ and $`t_0`$ are defined by the initial conditions, which for the transient stage are $`w_c=w(t_c)`$ and $`s_c=s(t_c)`$.
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# A Single Atom Mirror for 1D Atomic Lattice Gases
## I Introduction
One of the fundamental models in quantum optics is the interaction of a spin-$`1/2`$ system with a bosonic mode Cohen . The most prominent example is cavity quantum electrodynamics (CQED), where a two level atom interacts with a single mode of the radiation field in a high-Q cavity. CQED has been the topic of a series of seminal experiments both in the microwave and optical regime, demonstrating quantum control on the level of single atoms and photons in an open quantum system Kimble ; Haroche ; Rempe ; Walther .
In the present paper we will consider a system with the same basic ingredients, however in the context of cold atoms and quantum degenerate gases. The key feature of these systems is there controllability and weak decoherence. In particular we employ two aspects of control, the confinement of atoms in optical lattices JakschZollerReview ; Greiner ; Mandel ; Stoeferle and (magnetic or optical) Feshbach resonances as a way to manipulate atomic interactions Bolda ; Julienne ; Holland ; Theis . According to the setup described in Fig. 1(a) we will study the dynamics of an atomic quantum gas in 1D (with a single internal atomic state), representing bosonic or fermionic “modes”, controlled by an atomic spin-1/2 impurity. The quantum gas is confined by tight trapping potentials (e.g. an optical or magnetic trap), so that only the motional degrees along the $`z`$-axis in Fig. 1(a) are relevant. In the $`z`$-direction the motion is confined to the left by a trapping potential (e.g. a blue sheet of light), while the atomic impurity restricts the motion of the gas to the right due to collisional interactions of the quantum gas with the impurity. The atom representing the impurity can, for example, be a different atomic species in a tight trapping potential, a configuration discussed in Refs. Recati ; Raizen as an atomic quantum dot ($`0`$D system). Thus the impurity atom plays the role of “single atom mirror” confining the quantum gas in an “atomic cavity”.
In our model system the impurity atom is an internal two level system, which we write as an effective spin-1/2. In the following we will also interpret this two-level system as a qubit with two logical states $`|0=|`$ and $`|1=|`$. Cold atom collision physics allows for a situation where the collisional properties (scattering length) of the impurity atom and atoms in the quantum gas are spin-dependent. As illustrated in Fig. 1(a), we assume that in one spin state, say $`|`$, the single impurity atom is completely transparent for the quantum gas, i.e. the gas will leak out through the “mirror”. In contrast, in the other spin state the mirror atom is “highly reflective” confining the gas. For an impurity atom (qubit) initially prepared in a spin superposition
$$|\psi _Q(t=0)=\alpha _{}|+\alpha _{}|$$
the combined system at a time $`t`$ will be in a macroscopic superposition state
$`|\mathrm{\Psi }(t)=\alpha _{}||\varphi _{}(t)+\alpha _{}||\varphi _{}(t).`$ (1)
with $`|\varphi _\sigma (t)`$ many body wave functions of the gas atoms. Thus $`|\mathrm{\Psi }(t)`$ represents a Schrödinger cat state of two entangled quantum phases of gas atoms, the first one corresponding to gas confined by the mirror (Fig. 1(a) upper figure) and the second one to the expanding gas (Fig. 1(a) lower figure). The entanglement of the spin with a macroscopic number of atoms can be interpreted as a macroscopic quantum gate, as explained in Fig. 2), implementing a quantum nondemolition interaction (QND) qnd . In this sense the setup represents a “amplifier” of the state of the qubit. This situation is reminiscent of a Single Electron Transistor (SET) in mesoscopic physics set , and has stimulated the name Single Atom Transistor (SAT) for the setup Fig. 1(a) in Ref. Micheli , with the essential difference that the dynamics underlying (1) is completely coherent. We finally remark that this setup also allows for a single shot QND measurement of the impurity atom (qubit) by observing in a single experiment the distinct properties of the $`|\varphi _{}(t)`$ or $`|\varphi _{}(t)`$ quantum phases.
As a variant of the configuration of Fig. 1(a) we will consider below in particular the case where the quantum gas is loaded in an optical lattice, as illustrated in Fig. 1(b). In this case the gas could be loaded initially, for example, in a Mott insulating state, i.e. where large repulsion of the gas atom leads to a filling of the lattice sites with exactly one atom per lattice site Jaksch ; Greiner ; Stoeferle . The cat state (1) will thus correspond to a superposition of the Mott phase and the melted Mott phase, i.e. a (quasi-) condensate of atoms obtained by expansion of the atomic gas:
$`|\mathrm{\Psi }(t)=\alpha _{}||\text{BEC}+\alpha _{}||\text{Mott}.`$ (2)
In this case the distinguishing features of the two entangled quantum phases are the observation / non-observation of interference fringes as signatures of the Mott and BEC phase, when the atomic gas in released in a single experiment.
Transport through an impurity is a well studied problem in mesoscopic condensed matter physics Datta ; Mahan ; Levitov ; Cazalilla , which typically focuses on conductance properties of a system attached to leads. In contrast, in the context of cold gases we have a full time-dependent coherent dynamics in an otherwise closed system.
A short summary of the present work including results from numerical studies was presented in Ref. Micheli . In this paper we will present details of our analytical calculations, while we refer to Ref. Daley on a complementary numerical treatment of these problems using time dependent DMRG techniques.
The paper is organized as follows: In Sec. II we introduce the model used for describing the implementation of the Single Atom Mirror using cold atoms in optical lattice. In Sec. III we consider the detailed scattered processes involved in the transport of a single particle through the mirror. We solve exactly the scattering problem in the lattice by integrating the Lippmann-Schwinger Equation (LSE) and discuss the obtained scattering amplitudes and spectrum of the bound states. Finally, in Sec. IV we generalize the discussion to the case of interacting many-systems including the cases of a 1D degenerate Fermi-gas, a 1D quasi-condensate and Tonks gas.
## II Model
In this section we introduce the our model system by specifying the Hamiltonian for a 1D lattice gas coupled to an impurity, and we explain the key idea behind our setup. We will start with a discussion of spin-dependent collisions between the gas and the impurity, and then present the central idea of quantum interference as a way to switch atomic transport.
### II.1 Effective Spin-Dependent Hamiltonian
We consider the dynamics of a spin-1/2 atomic impurity $`Q`$ coupled to a 1D quantum gas of either bosonic or fermionic probe atoms $`A`$. The Hamiltonian for system is split into three parts as
$`H`$ $`=`$ $`H_A+H_Q+H_{AQ}.`$ (3)
Here $`H_Q`$ ($`H_A`$) describes the uncoupled dynamics of the impurity atom $`Q`$ (the degenerate quantum gas of probe atoms $`A`$), while $`H_{AQ}`$ accounts for the interaction between the two atomic species, $`Q`$ and $`A`$.
A degenerate quantum gas of bosonic or fermionic atoms $`A`$ trapped in the lowest band of a 1D optical lattice is well described by a Hubbard model JakschZollerReview
$`H_A`$ $`=`$ $`{\displaystyle \underset{j}{}}E_{A,j}a_j^{}a_jJ{\displaystyle \underset{ij}{}}a_i^{}a_j+{\displaystyle \frac{U}{2}}{\displaystyle \underset{j}{}}a_j^{}a_j^{}a_ja_j,`$ (4)
where $`a_j^{}`$ ($`a_j`$) are the creation (annihilation) operators for an atom $`A`$ on the site $`j`$, which obey standard commutation (anticommutation) relations for the case of bosonic (fermionic) atoms $`A`$. Moreover, $`E_{A,j}`$ account for the shift of the bare energy of an atom localized on the site $`j`$ in the presence of an external (e.g. magnetic) shallow trap, $`J`$ is the tunneling matrix element for neighboring sites $`ij`$ and $`U`$ gives the collisional interaction, i.e. the onsite-shift for two atoms $`A`$ localized within the same well (which would be zero for the case of fermions in the same internal state). Denoting the scattering-length of the atoms $`A`$ by $`a_s`$, and their mass by $`m`$, we have $`U=4\pi \mathrm{}^2a_sd^3r|w_j(𝐫)|^4/m`$, where $`w_j(𝐫)`$ is the Wannier function for a particle localize on the site $`j`$.
In the present setup we regard the impurity atom $`Q`$ to be trapped within a tight one-dimensional lattice, as depicted in Fig. 1(b). Therefore, we may restrict ourselves to the lowest trap-state of the $`j=0`$ well for the internal states $`\sigma =,`$, respectively. The uncoupled dynamics of the impurity corresponds to spin-1/2 system, i.e.
$`H_Q`$ $`=`$ $`{\displaystyle \underset{\sigma }{}}E_{Q,\sigma }|Q_\sigma Q_\sigma |,`$ (5)
where $`|Q_\sigma `$ ($`E_{Q,\sigma }`$) denotes the state (energy) of the atom $`Q`$ with spin $`\sigma =,`$.
Given the tight trapping of the impurity atom, the interaction of probe and impurity atom is restricted to the site of the impurity, an in general has the form of an effective spin-dependent collisional interaction
$`H_{AQ}`$ $`=`$ $`W_{\mathrm{eff},}||a_0^{}a_0+W_{\mathrm{eff},}||a_0^{}a_0,`$ (6)
where $`a_0^{}`$ ($`a_0`$) is the creation (annihilation) operator for a probe atom on the site of the impurity, $`j=0`$. Here, $`W_{\mathrm{eff},\sigma }=4\pi \mathrm{}^2a_\sigma /\mu d^3r|w_0(𝐫)|^2|\psi _{Q,\sigma }(𝐫)|`$ denotes the effective interaction for a probe atom $`A`$ and the impurity atom $`Q`$ in state $`\sigma `$ in terms of their effective scattering length $`a_\sigma `$ and $`\mu `$ is the reduced mass for $`A`$ and $`Q`$. The effective tunneling rate of a probe atom with energy $`E`$ through the impurity is then given by $`J_{\mathrm{eff},\sigma }=J^2/(EW_{\mathrm{eff},\sigma })+𝒪(J^4)`$ for the qubit in state $`\sigma `$. An obvious way to provide for a spin-dependent single atom mirror is to have the effective interaction for one spin state as large as possible ($`|W_{\mathrm{eff},}|J`$), thus blocking the transport of the probe atoms through the impurity site, while for the other it is effectively not present, ($`|W_{\mathrm{eff},}|J`$). This can be achieved, for example, by tuning the internal state dependent scattering length $`a_\sigma `$ or by engineering the spin-dependent trapping JakschZollerReview . The quality of the qubit dependent switch then depends on the difference of the moduli of the effective interactions, $`|W_{\mathrm{eff},}||W_{\mathrm{eff},}|`$. Thus the goal an efficient scheme is to make $`|W_{\mathrm{eff},}||W_{\mathrm{eff},}|`$ as large as possible and obtain $`|W_{\mathrm{eff},}|J|W_{\mathrm{eff},}|`$.
### II.2 Controlling the transport by interference
In this section we will show now that with the help of quantum interference we can engineer an effectively infinite (zero) atomic repulsion, $`W_{\mathrm{eff},}\mathrm{}`$ ($`W_{\mathrm{eff},}0`$), for the qubit in state $`\sigma =`$ ($`\sigma =`$). This is equivalent to tuning the Feshbach resonance to the point of infinite (zero) scattering length.
The quantum interference mechanism required to engineer the described spin-dependence of $`W_{\mathrm{eff},\sigma }`$ is obtained by exploiting the properties of either an optical or a magnetic Feshbach resonance. In the case of an optical Feshbach resonance a Raman laser drives the transitions from the joint state of the two atoms on the impurity site, $`a_0^{}|Q_\sigma |Q_\sigma a_0^{}|\mathrm{vac}`$, via an off-resonant excited molecular state back to a bound hetero-nuclear molecular state $`|M_\sigma `$ in the lowest electronic manifold (see Fig. 3). The Raman processes is described by the effective two-photon Rabi frequency $`\mathrm{\Omega }_\sigma `$ and detuning $`\mathrm{\Delta }_\sigma `$ for each spin-component $`\sigma `$. For the case of a magnetic Feshbach resonance, the effective Hamiltonian has the same form, but with $`\mathrm{\Omega }_\sigma `$ being the coupling strength between the open and closed channels and $`\mathrm{\Delta }_\sigma `$ the detuning of the magnetic field. The Hamiltonian describing the interaction between the probe atoms and the impurity is Holland
$$H_{AQ}=\underset{\sigma }{}[E_{M,\sigma }|M_\sigma M_\sigma |+\mathrm{\Omega }_\sigma (|M_\sigma Q_\sigma |a_0+\mathrm{h}.\mathrm{c}.)$$
$$+W_{Q,\sigma }|Q_\sigma Q_\sigma |a_0^{}a_0+W_{M,\sigma }|M_\sigma M_\sigma |a_0^{}a_0],$$
(7)
where the bare energy of the molecular bound state is $`E_{M,\sigma }=E_{A,0}+E_{Q,\sigma }+\mathrm{\Delta }_\sigma `$. Here the first two terms describe the resonant coupling induced by the Feshbach mechanism, while the last two describe the off-resonant collisions between an atom $`A`$ and an atom $`Q`$ (a molecule $`M`$) in state $`\sigma `$ by means of their on-site shift $`W_{Q,\sigma }`$ ($`W_{M,\sigma }`$) for the impurity site. The Hamiltonian (3) conserves the spin-component of the impurity, $`S_\sigma |Q_\sigma Q_\sigma |+|M_\sigma M_\sigma |`$, i.e. $`[H,S_\sigma ]=0`$. Therefore, we can consider the dynamics for the two spin components of $`Q`$ separately, and in the following we will drop the spin index $`\sigma `$ and choose the reference energy as $`E_{A,0}=E_{Q,\sigma }=0`$.
For off-resonant laser driving ($`|\mathrm{\Delta }|\mathrm{\Omega }`$), the Feshbach resonance enhances the interaction between $`A`$ and $`Q`$ atoms, giving the familiar result $`W_{\mathrm{eff}}=W_Q+\mathrm{\Omega }^2/\mathrm{\Delta }`$. However, for resonant driving ($`\mathrm{\Delta }=0`$) the physical mechanism changes, and the effective tunneling $`J_{\mathrm{eff}}`$ of an atom $`A`$ past the impurity (Fig. 4, $`\mathrm{I}\mathrm{III}`$) is blocked by quantum interference. On the impurity site, laser driving mixes the states $`a_0^{}|Q`$ and $`|M`$, forming two dressed states with energies $`E_\pm =W_Q/2\pm (W_Q^2/4+\mathrm{\Omega }^2)^{1/2}`$ (Fig. 4, II). The two resulting paths for a particle of energy $`E`$ destructively interfere so that for large $`\mathrm{\Omega }J`$ and $`W_Q=0`$,
$`J_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{J^2}{E+\mathrm{\Omega }}}{\displaystyle \frac{J^2}{E\mathrm{\Omega }}}0.`$
This is analogous to the interference effect underlying Electromagnetically Induced Transparency (EIT) eit , and is equivalent to having an effective interaction $`W_{\mathrm{eff}}\mathrm{}`$. In addition, if we choose $`\mathrm{\Delta }=\mathrm{\Omega }^2/W_Q`$, the paths interfere constructively, screening the background interactions $`W_Q`$ to produce perfect transmission ($`W_{\mathrm{eff}}0`$). The insensitivity of the interference scheme to losses from the dressed states due to their large detuning has been argued in Ref. Micheli .
## III Single particle scattering from an impurity
In this section we will analyze the scattering of a single probe atom $`A`$ from an impurity atom $`Q`$. We will formulate the scattering problem, then solve the time-independent and time-dependent Schrödinger Equation, to finally obtain the dynamics of wave-packets in the lattice.
We consider a probe atom $`A`$ approaching the impurity from the left, as a plain Bloch-wave with quasi-momentum $`k`$. Hence the state of the system is given by
$$|k=\left(\frac{a}{2\pi }\right)^{1/2}\underset{j}{}e^{ikx_j}|j,$$
(8)
where $`|j=|Qa_j^{}|\mathrm{vac}`$ is the joint state of the atoms $`A`$ and $`Q`$, with $`A`$ ($`Q`$) localized in the lowest vibrational state of the well $`j`$ (the impurity well $`j=0`$) and $`a`$ is the lattice spacing.
The free evolution of the system is given by the hopping of the atoms $`A`$ between neighboring sites at the tunneling rate $`J`$, whereas the composite molecule $`M`$ is detuned by $`\mathrm{\Delta }`$ from the threshold for the joint state of $`A`$ and $`Q`$. Thus, with $`E_k=2J\mathrm{cos}ka`$ being the energy of a Bloch-wave in the first Bloch-band with quasi-momentum $`k`$, we have
$`H_0`$ $`=`$ $`J{\displaystyle \underset{j}{}}(|j+1j|+\mathrm{h}.\mathrm{c}.)+\mathrm{\Delta }|MM|`$ (9)
$`=`$ $`{\displaystyle _{\pi /a}^{\pi /a}}𝑑kE_k|kk|+\mathrm{\Delta }|MM|,`$
where $`|M`$ denotes the molecular bound state localized on the impurity site. From Eq. (9) we obtain the propagation of the incoming plane wave $`|k`$ at group-velocity $`v_k=E_k/k=2Ja\mathrm{sin}ka`$ in the first Bloch band.
Due to the strong confinement of the particles $`A`$ and $`Q`$ in the lattices, their interaction is restricted to the impurity site. There, their bare interaction induces an on-site-shift $`W`$ for the joint atomic state of $`A`$ and $`Q`$ on the impurity ($`|0`$). Moreover, the photo-association lasers effectively couple the latter state to the trapped molecular state ($`|M`$) at Rabi-frequency $`\mathrm{\Omega }`$, yielding
$`V`$ $`=`$ $`W|00|+\mathrm{\Omega }(|M0|+\mathrm{h}.\mathrm{c}.).`$ (10)
### III.1 Scattering solution
The scattering of a particle $`A`$ with energy $`E=E_k`$ in the first Bloch band by the impurity $`Q`$ is described by a solution of the Lippmann-Schwinger Equation (LSE). The scattering wave function $`|\varphi _+`$ obeys
$$|\varphi _+=|k+G_0(E+i0^+)V|\varphi _+,$$
(11)
with incident plane wave $`|k`$ with quasimomentum $`k`$ ($`0<k<\pi /a`$), and $`G_0(z)=1/(zH_0)`$ the free propagator. Expanding the scattering wave function
$$|\varphi _+=\left(\frac{a}{2\pi }\right)^{1/2}\left[\underset{j}{}\alpha _j|j+\beta |M\right]$$
(12)
the amplitudes $`\alpha _j`$ and $`\beta `$ satisfy
$`\alpha _j`$ $`=`$ $`e^{ikx_j}+𝒢_j(E_k)\left(W\alpha _0+\mathrm{\Omega }\beta \right),`$ (13a)
$`\beta `$ $`=`$ $`𝒢_M(E_k)\mathrm{\Omega }\alpha _0`$ (13b)
with atomic and molecular propagators
$`𝒢_j(E)`$ $`=`$ $`j|G_0(E+i0^+)|0`$
$``$ $`{\displaystyle \frac{a}{2\pi }}{\displaystyle _{\pi /a}^{+\pi /a}}𝑑k{\displaystyle \frac{e^{ikx_j}}{EE_k+i0^+}}={\displaystyle \frac{e^{ik|x_j|}}{iv_k/a}},`$
$`𝒢_M(E)`$ $`=`$ $`M|G_0(E+i0^+)|M{\displaystyle \frac{1}{E\mathrm{\Delta }+i0^+}}.`$
Solving Eqs.(13) we find
$`\alpha _j`$ $`=`$ $`e^{ikx_j}+{\displaystyle \frac{W_k}{iv_k/aW_k}}e^{ik|x_j|},`$ (15a)
$`\beta `$ $`=`$ $`{\displaystyle \frac{i\mathrm{\Omega }v_k/a}{\mathrm{\Omega }^2+\left(E_k\mathrm{\Delta }+i0^+\right)\left(Wiv_k/a\right)}}`$ (15b)
with effective energy dependent interaction
$$W_k=W+\frac{\mathrm{\Omega }^2}{E_k\mathrm{\Delta }+i0^+},$$
(16)
where we read off the transmission and reflection amplitudes
$`t_k`$ $`=`$ $`{\displaystyle \frac{1}{1+iaW_k/v_k}}`$ (17a)
$`r_k`$ $`=`$ $`{\displaystyle \frac{1}{1iv_k/aW_k}},`$ (17b)
respectively.
Note that the presence of the molecular state introduces an effective energy-dependent interaction $`W_k`$. This can be interpreted in terms of an effective atomic scattering length with background scattering length proportional to $`W`$ and a resonant term, corresponding to an optical Feshbach resonance at energy given by the detuning from the molecular state $`\mathrm{\Delta }`$, and width determined by the Rabi frequency $`\mathrm{\Omega }`$.
The scattering matrix
$$S(E_k)=\left(\begin{array}{cc}r_k& t_k\\ t_k& r_k\end{array}\right)$$
(18)
is unitary, as follows readily from the above expressions (17). This implies $`T_k+R_k|t_k|^2+|r_k|^2=1`$. We can assign phase shifts $`t_k\pm r_k=\mathrm{exp}(i\delta _k^\pm )`$ for the symmetric and antisymmetric states $`|k\pm |k`$, $`\delta _k^+=2\mathrm{arctan}(aW_k/v_k)`$ and $`\delta _k^{}=0`$, respectively, so that $`R_k=\mathrm{sin}^2(\delta _k^+/2)`$.
### III.2 Discussion of the Scattering Amplitudes
In the absence of molecular couplings ($`\mathrm{\Omega }=0`$) the on-site interaction $`W`$ between the B atom and the impurity Q always gives rise to *partial* reflection and transmission, see Fig. 5a,
$`R_k=1T_k={\displaystyle \frac{W^2}{W^2+4J^2\mathrm{sin}^2\left(ka\right)}}<1.`$ (19)
The significant new feature introduced by the optical Feshbach resonance is that we can achieve essentially *complete blocking* ($`R_k=1`$) and *complete transmission* ($`T_k=1`$). We obtain this in the limits $`\mathrm{\Omega }J`$ and $`\mathrm{\Delta }=0`$, and $`\mathrm{\Omega }J`$ and $`\mathrm{\Delta }=\mathrm{\Omega }^2/W`$, respectively. Physically, the first case corresponds to tuning to the point of “infinite” scattering length, while the second case corresponds to tuning to the point of “zero” scattering length, respectively.
In the general case the energy dependence of transmission and reflection has the form of a Fano-like profile (see Fig. 5). In the region $`|E_k|2J`$ we may neglect the dispersion effects, i.e. $`v_k2Ja`$, and obtain Fano-line-shapes for transmission and reflection as
$`T(\epsilon )`$ $`=`$ $`{\displaystyle \frac{1}{1+q_f^2}}{\displaystyle \frac{\left(\epsilon +q_f\right)^2}{\epsilon ^2+1}},`$ (20a)
$`R(\epsilon )`$ $`=`$ $`{\displaystyle \frac{1}{1+1/q_f^2}}{\displaystyle \frac{\left(\epsilon 1/q_f\right)^2}{\epsilon ^2+1}},`$ (20b)
where $`\epsilon (EE_R)/(\gamma /2)`$ is the dimensionless energy in units of the resonance width $`\gamma `$, $`E_R`$ is the resonance energy and $`q_f`$ is the Fano-$`q`$-parameter. These parameters of the Fano-profile are related to $`J,W,\mathrm{\Delta },\mathrm{\Omega }`$ by
$`\gamma `$ $`={\displaystyle \frac{4J\mathrm{\Omega }^2}{W^2+4J^2}},`$ $`\mathrm{\Omega }^2`$ $`=J\gamma \left(q_f^2+1\right),`$ (21a)
$`E_R`$ $`=\mathrm{\Delta }{\displaystyle \frac{W\mathrm{\Omega }^2}{W^2+4J^2}},`$ $`\mathrm{\Delta }`$ $`=E_R{\displaystyle \frac{q_f\gamma }{2}},`$ (21b)
$`q_f`$ $`={\displaystyle \frac{W}{2J}},`$ $`W`$ $`=2Jq_f.`$ (21c)
For $`W=0`$ the asymmetry parameter $`q_f`$ vanishes, the reflection profile is symmetric, and for $`\mathrm{\Omega }>0`$ resembles a Breit-Wigner-profile, see Fig 5(b). The maximum $`R_k=1`$ is attained at $`\epsilon =q_f`$ ($`E_k=\mathrm{\Delta }`$, $`|\mathrm{\Delta }|<2J`$) and has a width $`\gamma =\mathrm{\Omega }^2/J`$.
For finite background collisions $`q_f0`$ ($`W0`$) the transmission profile is asymmetric, and shows an additional minimum $`R_k=0`$ at $`\epsilon =1/q_f`$ ($`E_k=\mathrm{\Delta }\mathrm{\Omega }^2/W=\mathrm{\Delta }_{}`$, $`|\mathrm{\Delta }_{}|<2J`$), see Fig 5(c).
Near the edges of the Bloch band, $`k_\pm =(\pi \pi )/2a`$ ($`E_k=\pm 2J`$), transmission and reflection deviate from the Fano line-shape Eq. (20). There the group-velocity $`v_k0`$ and thus also the transmission $`T_k`$ vanishes, unless the dressed resonance $`\mathrm{\Delta }_{}`$ is tuned to respective edge of the Bloch band. The transmission coefficient are given by
$`T_{kk_\pm }{\displaystyle \frac{4J^2a^2\left(kk_\pm \right)^2}{\left[W\mathrm{\Omega }^2/\left(\mathrm{\Delta }2J\right)\right]^2}}`$ for $`\mathrm{\Delta }_{}\pm 2J,`$
$`T_{kk_\pm }1{\displaystyle \frac{W^4a^2\left(kk_\pm \right)^2}{2\mathrm{\Omega }^4}}`$ for $`\mathrm{\Delta }_{}=\pm 2J.`$
The reflection coefficient $`R_k`$ as a function of energy $`2J<E_k<+2J`$ is shown in Fig. 5.
In the absence of molecular coupling, $`\gamma =0`$ ($`\mathrm{\Omega }=0`$), the reflection is unity at the band-edges, $`k=0,\pi /a`$, and decreases within the Bloch-band due to the increase of the group-velocity (see Fig.5(a)). The profile is symmetric about the middle of the Bloch band, $`k=\pi /2a`$, where it attains its *minimum*,
$`R_k={\displaystyle \frac{1}{1+q_f^2}}={\displaystyle \frac{1}{1+\left(2J/W\right)^2}}.`$
In the presence of molecular couplings, $`\gamma 0`$ ($`\mathrm{\Omega }>0`$), and for $`E_R=q_f=0`$ the reflection profile is still symmetric about $`k=\pi /2a`$ (see Fig.5(b)). However, now it approaches its *maximum*, $`R_k=1`$, at $`k=\pi /2a`$, and has now two minima at $`k\pi /4a`$ and $`k3\pi /4a`$, given by
$`R_k={\displaystyle \frac{1}{1+\left(2J/\gamma \right)^2}}={\displaystyle \frac{1}{1+\left(\sqrt{2}J/\mathrm{\Omega }\right)^4}}.`$
For $`\gamma ,q_f0`$ we obtain an asymmetric Fano-profile (see Fig.5(d)), which for $`|E_R+q_f\gamma /2|=|\mathrm{\Delta }|<2J`$ shows complete reflection, $`R_k=1`$, at $`\epsilon =q_f`$, while for $`|E_R\gamma /2q_f|<2J`$ one has perfect transmission, $`R_k=0`$, at $`\epsilon =+1/q_f`$.
The reflective and transmissive resonance are present regardless of the magnitude of $`q_f`$, and their width is $`\gamma `$. However, for $`\gamma 8J|q_f+1/q_f|`$ they may both occur within the physical energy range of the Bloch-band (see Fig.5(d)), while for $`\gamma >8J|q_f+1/q_f|`$ only one resonance appears (see Fig.5(b,c)). Thus in the limit $`\gamma 8J|q_f+1/q_f|`$ we achieve complete blocking , $`R_k=1`$ for all $`k`$, by tuning $`E_Rq_f\gamma /2`$ (see full-line in Fig.5(b)). Within the same limit we can also efficiently screen any background-interaction $`W`$ and achieve complete transparency, $`T_k=1`$ for all $`k`$, by tuning $`E_R\gamma /2q_f`$ (see full-line in Fig. 5(c)).
### III.3 Interference mechanism
Physically, the features of complete blocking ($`R_k=0`$) and complete transmission ($`T_k=0`$) are induced by an interference mechanism, as the probe atom may tunnel via two interfering paths of dressed atomic \+ molecular states, as depicted in Fig. 4.
For simplicity we start by elucidating the underlying interference mechanism (present for $`\mathrm{\Omega }0`$) in the regime of strong coupling, $`\mathrm{\Omega }^2(2J)^2+|W\mathrm{\Delta }|`$. In this regime we can consider the local dynamics within the individual sites, and treat the tunneling $`J`$ by means of perturbation theory.
To zeroth order in $`J`$ the Hamiltonian $`H`$ decouples the dynamics of the individual sites $`j`$ as
$`H^{(0)}=`$ $`W|00|`$ $`+\mathrm{\Omega }|0M|+`$
$`+`$ $`\mathrm{\Omega }|M0|`$ $`+\mathrm{\Delta }|MM|.`$ (23)
Outside the impurity ($`j0`$) its eigenstates are the joint states of the atoms $`A`$ and $`Q`$, $`|j`$, with energy $`E_0=0`$, whereas on the impurity ($`j=0`$) the strong coupling $`\mathrm{\Omega }`$ between the atomic state $`|0`$ and the molecular state $`|M`$ induces the two states to split into two dressed state $`|E_\pm `$ of atoms + molecules with energy $`E_\pm `$, see Fig. 4. By diagonalizing the $`2\times 2`$-matrix in Eq. (23) we obtain the amplitudes and energy of the dressed states as
$`|E_\pm `$ $`=`$ $`\left[{\displaystyle \frac{1\pm \xi }{2}}\right]^{1/2}|0\pm \left[{\displaystyle \frac{1\xi }{2}}\right]^{1/2}|M,`$ (24a)
$`E_\pm `$ $`=`$ $`{\displaystyle \frac{W+\mathrm{\Delta }}{2}}\pm \left[\mathrm{\Omega }^2+\left({\displaystyle \frac{W\mathrm{\Delta }}{2}}\right)^2\right]^{1/2},`$ (24b)
$`\xi `$ $`=`$ $`{\displaystyle \frac{W\mathrm{\Delta }}{\left[\left(W\mathrm{\Delta }\right)^2+4\mathrm{\Omega }^2\right]^{1/2}}},`$ (24c)
where $`\xi `$ characterizes the asymmetry of the amplitudes, i.e. for $`\xi =0`$ ($`W=\mathrm{\Delta }`$) the dressing is completely symmetric while for $`|\xi |=1`$ ($`\mathrm{\Omega }=0`$) the atomic and molecular state decouple. From Eq. (24b) we see that for $`\mathrm{\Omega }|W\mathrm{\Delta }|`$ the dressed states are far off-resonant from $`E=0`$, and hence will be only virtually populated.
The effects of the hopping of the atom $`A`$ on the $`E0`$ modes $`|j0`$ can be accounted by means of an effective Hamiltonian $`H_{\mathrm{eff}}`$. Following Ref. Cohen we obtain the dynamics as a perturbative series in the hopping amplitude $`J`$, $`H_{\mathrm{eff}}=H^{(0)}+H^{(1)}+H^{(2)}+\mathrm{}`$. To first order in $`J`$ one obtains
$`H^{(1)}`$ $`=`$ $`J{\displaystyle \underset{j<0}{}}|jj1|J{\displaystyle \underset{j>0}{}}|jj+1|+\mathrm{h}.\mathrm{c}.`$ (25)
$`=`$ $`{\displaystyle \underset{\alpha =L,R}{}}{\displaystyle 𝑑kE_k|k_\alpha k_\alpha |},`$
where $`|k_L`$ ($`|k_R`$) are the Bloch-waves with quasi-momentum $`k`$ on the left (right) side of the impurity site,
$`|k_{L,R}`$ $`=`$ $`\left({\displaystyle \frac{a}{2\pi }}\right)^{1/2}{\displaystyle \underset{\pm j>0}{}}e^{+ikx_j}|j.`$
This the flat dispersion relation $`E^{(0)}=0`$ on the left and right side of the impurity is bent to $`\epsilon (k)=2J\mathrm{cos}(ka)`$, i.e. we recover the Bloch-band(s).
To second order in $`J`$ we obtain
$$H^{(2)}=J_{\mathrm{eff}}\underset{i,j=\pm 1}{}|ij|,$$
(26)
$$J_{\mathrm{eff}}=\frac{J^2}{2}\left[\frac{1+\xi }{E_0E_+}+\frac{1\xi }{E_0E_{}}\right]=\frac{J^2\mathrm{\Delta }}{\mathrm{\Omega }^2W\mathrm{\Delta }}.$$
(27)
We see that tuning on resonance $`\mathrm{\Delta }=0`$ the two contributions in Eq. (27) cancel each other as
$`J_{\mathrm{eff}}={\displaystyle \frac{J^2}{\sqrt{W^2+\mathrm{\Omega }^2}}}{\displaystyle \frac{J^2}{\sqrt{W^2+\mathrm{\Omega }^2}}}=0,`$
which gives perfect blocking by the impurity. Furthermore, from Eq. (24b) we obtain that for $`\mathrm{\Delta }_{}=\mathrm{\Delta }\mathrm{\Omega }^2/W=0`$ one of the dressed states $`|E_B`$ becomes a resonance for an incoming particle ($`E_B=E_0=0`$) and for $`\mathrm{\Omega }0`$ provides for complete transmission by means of photo-assisted tunneling. The described interference mechanism induced by the optical Feshbach resonance is in marked contrast to the situation where one has background collisions. There the particle $`A`$ can tunnel only via one path through the impurity ($`|\xi |=1`$), and therefore the effective hopping rate is always finite, i.e. $`J_{\mathrm{eff}}=J^2/W0`$.
### III.4 Discussion of Bound-states
For completeness we here derive the exact bound-state spectrum of $`H`$. For the exact scattering solution, detailed in Sec. III, the bound-states take the role of dressed states $`|E_\pm `$, which are responsible for the interference mechanism. We will show that for arbitrary $`\mathrm{\Omega }`$ there are always *two* bound-states, provided $`|\mathrm{\Delta }\mathrm{\Omega }^2/W|>2J`$. For $`|\mathrm{\Delta }\mathrm{\Omega }^2/W|2J`$ one of the bound-states turns into a resonance, which makes the impurity completely transparent for the atom $`A`$, $`T_k=1`$. Furthermore, we will show that the bound-states for extends over several lattice sites for $`\mathrm{\Omega },W,\mathrm{\Delta }J`$. This is in marked contrast to the perturbative result, where the dressed states were localized on the impurity site, cf. Eq. (24a).
We obtain the bound states wavefunctions $`|\varphi _B`$ from the homogeneous Lippmann-Schwinger equation
$`|\varphi _B`$ $`=`$ $`G_0(E)V|\varphi _B,`$ (29)
where $`\varphi _B`$ denotes the bound-state with energy $`EE_B`$ ($`|E_B|>2J`$). Using the ansatz
$`|\varphi _B`$ $`=`$ $`{\displaystyle \underset{j}{}}\alpha _j|j+\beta |M`$ (30)
we find that the atomic and molecular amplitudes, $`\alpha _j`$ and $`\beta `$, satisfy
$`a_j`$ $`=`$ $`𝒢_j(E_B)\left[W\alpha _0+\mathrm{\Omega }\beta \right],`$ (31a)
$`\beta `$ $`=`$ $`𝒢_M(E_B)\mathrm{\Omega }\alpha _0.`$ (31b)
The atomic and molecular propagators are given by
$`𝒢_j(E_B)`$ $`=`$ $`j|G_0(E_B)|0={\displaystyle \frac{e^{|x_j|/r_B}\left[\mathrm{sign}\left(E_B\right)\right]^j}{2J\mathrm{sinh}\left(\kappa _Ba\right)}},`$
$`𝒢_M(E_B)`$ $`=`$ $`M|G_0(E_B)|0={\displaystyle \frac{1}{E_B\mathrm{\Delta }}},`$
and $`r_B=a/\mathrm{acosh}(|E_B|/2J)>0`$ denotes the size of the bound-state.
For convenience we first consider the case $`\mathrm{\Omega }=0`$. There the molecular state decouples from the atomic ones, and we have one bound-state $`|\varphi _1=|M`$ with energy $`E_1=\mathrm{\Delta }`$. Moreover, for $`W0`$ we have another bound-state $`|\varphi _2`$ with energy $`E_2=\mathrm{sign}(W)\sqrt{W^2+(2J)^2}`$. Its amplitudes are given by $`\beta =0`$ and
$`\alpha _j=\left[\mathrm{tanh}\left({\displaystyle \frac{a}{r_B}}\right)\right]^{1/2}e^{|x_j|/r_B}\left[\mathrm{sign}(W)\right]^j,`$
with the size $`r_B=a/\mathrm{arccos}\sqrt{1+(W/2J)^2}`$. The spectrum of the system is plotted in Fig. 6(b) as a function $`W`$. For attractive (repulsive) interaction $`W<0`$ ($`W>0`$) the energy of the bound-state $`\varphi _2`$, lies below (above) the Bloch-band, i.e. $`E_2<2J`$ ($`E_2>2J`$), respectively. For $`|W|<2J`$ bound state $`\varphi _2`$ extends over several lattice sites, while for $`|W|2J`$ it is localized on the impurity.
In the case $`\mathrm{\Omega }0`$ a nontrivial solution of Eq. (31) requires
$`E_B\left[1\left({\displaystyle \frac{2J}{E_B}}\right)^2\right]^{1/2}`$ $`=`$ $`W+{\displaystyle \frac{\mathrm{\Omega }^2}{E_B\mathrm{\Delta }}},`$ (33)
which determines the bound-state spectrum $`E_B(W,\mathrm{\Delta },\mathrm{\Omega })`$. From Eq. (31) we obtain the atomic and molecular amplitudes as
$`\alpha _j`$ $`=`$ $`\beta {\displaystyle \frac{E_B\mathrm{\Delta }}{\mathrm{\Omega }}}e^{|x_j|/r_B}\left[\mathrm{sign}\left(E_B\right)\right]^j,`$ (34a)
$`\beta `$ $`=`$ $`\left[1+{\displaystyle \frac{\left(E_B\mathrm{\Delta }\right)^2}{\mathrm{\Omega }^2\sqrt{1\frac{4J^2}{E_B^2}}}}\right]^{1/2}.`$ (34b)
We solve Eq. (33) by expressing one of the parameters, either $`W`$, $`\mathrm{\Omega }`$ or $`\mathrm{\Delta }`$, in terms of the bound-state energy $`E_B`$,
$`\mathrm{\Delta }\left(E_B\right)`$ $`=`$ $`E_B+{\displaystyle \frac{\mathrm{\Omega }^2}{E_B\sqrt{14J^2/E_B^2}+W}},`$ (35a)
$`\mathrm{\Omega }\left(E_B\right)`$ $`=`$ $`\left[\left(\mathrm{\Delta }E_B\right)\left(E_B\sqrt{1{\displaystyle \frac{4J^2}{E_B^2}}}+W\right)\right]^{1/2},`$ (35b)
$`W\left(E_B\right)`$ $`=`$ $`E_B\sqrt{1+{\displaystyle \frac{4J^2}{E_B^2}}}+{\displaystyle \frac{\mathrm{\Omega }^2}{E_B\mathrm{\Delta }}}.`$ (35c)
Inverting the functions Eq. (35) for $`E_B`$ yields $`E_B`$ as a function of $`\mathrm{\Delta }`$, $`\mathrm{\Omega }`$ and $`W`$, respectively. For fixed $`\mathrm{\Delta }`$, $`\mathrm{\Omega }`$, $`W`$ we carry out the inversion by plotting in Fig. 6 the r.h.s. of Eq. (35)(a,b,c) as a function of $`E_B`$. In particular in Fig. 6(b) we plot detuning $`\mathrm{\Delta }`$ as a function of $`E_B`$ for constant $`\mathrm{\Omega }`$ and $`W`$, in Fig. 6(c) we plot Rabi-frequency $`Omega`$ as a function of $`E_B`$ for constant $`\mathrm{\Delta }`$ and $`W`$, and in Fig. 6(d) we plot on-site shift $`W`$ as a function of $`E_B`$ for constant $`\mathrm{\Delta }`$ and $`\mathrm{\Omega }`$. In the following we will give a detailed discussion of Fig. 6.
For no background collisions, $`W=0`$, and arbitrary detuning $`\mathrm{\Delta }`$, one always has two bound-states $`\varphi _{1,2}`$ with energy $`E_1<2J`$ and $`E_2>2J`$, respectively, see solid line in Fig. 6(c) corresponding to $`\mathrm{\Omega }=4J`$. For $`|\mathrm{\Delta }|\mathrm{\Omega }`$ the energy one bound-state approaches $`\mathrm{\Delta }`$ and is wavefunction becomes localized on the impurity, while the energy the other approaches the Bloch-band and its wavefunction extends over several lattice-sites. For $`\mathrm{\Delta }=0`$ the two bound-states are split symmetrically, and their energies are given by (see solid line in Fig. 6(d))
$`E_{1,2}=\pm \sqrt{2J^2+\sqrt{4J^2+\mathrm{\Omega }^2}}.`$ (36)
The symmetric splitting allows for complete reflection at $`E_k=0`$. In the limit $`\mathrm{\Omega }2J`$ we recover the perturbative result Eq. (24b), as the energies of the bound states approach $`E_\pm =\pm \mathrm{\Omega }`$, with their wavefunctions given by the dressed states $`|E_\pm `$, cf. Eq. (24a).
For finite onsite shift, $`W0`$, we also have two bound-states provided $`|\mathrm{\Delta }_{}||\mathrm{\Delta }\mathrm{\Omega }^2/W|>2J`$, see dashed lines in Fig. 6(c,d). With increasing detuning $`\mathrm{\Delta }`$ the energy of the bound-state $`\varphi _1`$ approaches the Bloch-band from below until crossing it for $`2J+\mathrm{\Omega }^2/W<\mathrm{\Delta }<+2J+\mathrm{\Omega }^2/W`$. In this parameter regime there is merely one bound state, while the other develops a resonance. This allows for perfect transparency at $`E_k=\mathrm{\Delta }\mathrm{\Omega }^2/W`$.
Finally we remark that the size $`r_B`$ of the bound-states is inversely proportional to the separation of their energy $`E_B`$ from their Bloch-band. Thus for $`|E_B|2J`$ the wave-function of the bound-states extends over several lattice sites. For $`|E_B|2J`$ the bound-states are localized on the impurity and we recover the results of the previous section.
### III.5 Wave-packet dynamics
As an illustration of the time-dependence of the interference mechanism we simulate the evolution of a gaussian wave-packet $`\psi (t)`$ with mean quasi-momentum $`k=\pi /2a`$ incident from the left of the impurity.
These wave-packets are obtained as superposition of the scattering solutions $`\varphi _+`$ of Sec. III, i.e. their atomic and molecular amplitudes, $`\alpha _j(t)=j|\psi (t)`$ and $`\beta (t)=M|\psi (t)`$, are obtained as
$`\alpha _j(t)`$ $`=`$ $`{\displaystyle \underset{j^{}}{}}U_{j,j^{}}(t)\alpha _j^{}(0),`$ (37a)
$`\beta (t)`$ $`=`$ $`{\displaystyle \underset{j^{}}{}}U_{M,j^{}}(t)\alpha _j^{}(0),`$ (37b)
with the full propagator for the system given by
$`U_{j^{},j}(t)={\displaystyle \underset{B}{}}{\displaystyle \frac{e^{(|x_j^{}|+|x_j|)/r_BiE_Bt}\mathrm{sign}(E_B)^{j^{}+j}}{\frac{|E_B|}{\sqrt{E_B^2+4J^2}}+\left[\frac{\mathrm{\Omega }}{E_B\mathrm{\Delta }}\right]^2}}+`$
$`+{\displaystyle \frac{a}{2\pi }}{\displaystyle 𝑑ke^{iE_kt}\left[e^{+ik|x_j^{}x_j|}+r_ke^{+ik(|x_j^{}|+|x_j|)}\right]},`$
$`U_{M,j}(t)={\displaystyle \underset{B}{}}{\displaystyle \frac{e^{|x_j|/r_BiE_Bt}\mathrm{sign}(E_B)^j}{\sqrt{1+\frac{|E_B|(E_B\mathrm{\Delta })^2}{\mathrm{\Omega }^2\sqrt{E_B^24J^2}}}}}+`$
$`+{\displaystyle \frac{a}{2\pi }}{\displaystyle 𝑑ke^{iE_kt}\left[\frac{\mathrm{\Omega }t_k}{E_k\mathrm{\Delta }+i0^+}e^{+ik(|x_j|)}\right]}.`$ (38)
On the left side of Fig. 7 we plot the atomic populations of the individual sites, $`n(x_j,t)=|\alpha _j(t)|^2`$. The right side shows the corresponding atomic populations of the atom $`A`$ on the left, $`n_L=_{j<0}n(x_j,t)`$ (dashed line), and on the right side of the impurity, $`n_R=_{j>0}n(x_j,t)`$ (dashed-dotted line). We also plot the population on the impurity, i.e. the atomic population, $`n_0=n(x_0,t)`$ (solid line), and the molecular population, $`n_M(t)=|M|\psi (t)|^2`$ (dotted line). The three different sets in Fig. 7 correspond to different coupling strengths, $`\mathrm{\Omega }`$, and detunings, $`\mathrm{\Delta }`$. For all cases we choose $`W=2J`$. In Fig. 7(a) we have $`\mathrm{\Omega }=0`$: the atom is partially reflected from the impurity with $`R_k=1`$. In Fig. 7(b) we set $`\mathrm{\Omega }=2J`$ and $`\mathrm{\Delta }=0`$, which gives rise to complete reflection of the wavepacket, $`R_k=1`$. In Fig. 7(c) we have $`\mathrm{\Omega }=2J`$, but now $`\mathrm{\Delta }=2J`$. We have complete transmission of the atom through the impurity, $`T_k=1`$. All this is consistent with the results of Sec. III.
## IV Many body scattering from an impurity
In this section we will analyze the evolution of a 1D lattice gas of *many* atoms $`A`$ interacting with an impurity atom $`Q`$. Since the statistics of the atoms $`A`$ plays a dominant role, we will consider the cases of fermionic and bosonic atoms, separately. In this context we will study analytically the limiting cases of an ideal Fermi-gas, an ideal Bose-gas and a Tonks-gas. An exact numerical treatment of the dynamics for the lattice-gas having arbitrary interaction $`U`$ is given in Ref. Daley .
### IV.1 Ideal Fermi-gas
We first consider the case, where the probe atoms $`A`$ are spin-polarized fermions. The Hamiltonian for the system is given by
$`H`$ $`=`$ $`J{\displaystyle \underset{j}{}}\left(a_j^{}a_{j+1}+a_{j+1}^{}a_j\right)+\mathrm{\Delta }|MM|+`$ (39)
$`+\mathrm{\Omega }\left(|MQ|a_0+a_0^{}|QM|\right)+`$
$`+W_Q|QQ|a_0^{}a_0+W_M|MM|a_0^{}a_0,`$
where the operators $`a_j^{}`$ ($`a_j`$) create (annihilate) an atom $`A`$ on site $`j`$, and obey the canonical anti-commutation relations $`\{a_i,a_j^{}\}=\delta _{ij}`$ and $`\{a_i,a_j\}=\{a_i^{},a_j^{}\}=0`$. Moreover, $`|Q`$ ($`|M`$) denote the states with an atom $`Q`$ (a molecule in state $`M`$) on the impurity, and $`W_Q`$ ($`W_M`$) is the onsite-shift for an atom $`A`$ and an atom $`Q`$ (a molecule $`M`$) on the impurity.
For simplicity henceforth we will restrict ourselves to the case of equal on-site shifts $`W_M=W_QW`$. In this case we may rewrite the Hamiltonian as
$`H`$ $`=`$ $`J{\displaystyle \underset{j}{}}\left(a_j^{}a_{j+1}+a_{j+1}^{}a_j\right)+Wa_0^{}a_0+`$ (40)
$`\mathrm{\Delta }f^{}f+\mathrm{\Omega }\left(f^{}a_0+a_0^{}f\right),`$
where the ladder operators $`f^{}|MQ|`$ and $`f|QM|`$ obey standard fermionic anti-commutation relations and anti-commute with $`a_j`$ and $`a_j^{}`$. The corresponding equations of motions for $`a_j`$ and $`f`$ are linear, provided $`W_M=W_Q`$. Thus for a Fermi-gas of $`N`$ atoms $`A`$ the scattering off the impurity atom $`Q`$ will occur independently for each fermion $`a_k`$ with scattering amplitudes $`t_k`$ and $`r_k`$, according to their quasi-momentum $`k`$, cf. Eq. (17). The details of this calculation will be given below.
We will here detail the time-dependent scattering for a Fermi-gas of $`N`$ atoms $`A`$. For concreteness we assume the fermions to be initially trapped in a box of $`M`$ sites to the left the impurity atom $`Q`$. The corresponding wavefunction of the system is given by
$`|\mathrm{\Psi }(t=0)={\displaystyle \underset{n=1}{\overset{N}{}}}\left[{\displaystyle \underset{j}{}}\alpha _j(k_n)a_j^{}\right]|Q,`$ (41)
$`\alpha _j(k_n)=\sqrt{{\displaystyle \frac{2}{M}}}\{\begin{array}{cc}\mathrm{sin}(k_nx_j)& \text{for }Mj1,\\ 0& \text{else},\end{array}`$ (44)
where the quasi-momenta $`k_n=n\pi /(M+1)a`$. This corresponds to a Fermi-sea filled up to $`E_F=4J\mathrm{sin}^2(k_Fa)`$, where $`k_F=\nu \pi /a`$ is the Fermi-momentum and the initial filling-factor $`\nu =N/(M+1)`$.
At time $`t=0`$ we open the impurity (cf. Fig. 1(b)), and from Eq. (40) we obtain the evolution of the system as
$$|\mathrm{\Psi }(t)=\underset{n=1}{\overset{N}{}}\left[\underset{j}{}\alpha _j(k_n,t)a_j^{}+\beta (k_n,t)f^{}\right]|Q,$$
(45a)
$$\alpha _j(k_n,t)=\underset{j^{}}{}U_{j,j^{}}(t)\alpha _j^{}(k_n),$$
(45b)
$$\beta (k_n,t)=\underset{j^{}}{}U_{M,j^{}}(t)\alpha _j^{}(k_n),$$
(45c)
where $`U_{\alpha ,j}(t)`$ are the single-particle propagators, cf. Eq. (38). According to Eq. (45a) the scattering from the impurity occurs independently for each particle in the initial Fermi sea, with scattering amplitudes $`t_k`$ and $`r_k`$ given in Eq. (17) for $`0<kk_F`$. The atomic and molecular densities are thus given by the sum of the probabilities for the single fermions in the Fermi-gas,
$`n(x_j,t)`$ $`=`$ $`a_j^{}a_j_t={\displaystyle \underset{n=1}{\overset{N}{}}}|\alpha _j(k_n,t)|^2,`$ (46a)
$`n_M(t)`$ $`=`$ $`1n_Q(t)={\displaystyle \underset{n=1}{\overset{N}{}}}|\beta (k_n,t)|^2.`$ (46b)
Moreover, after opening the switch the atomic quasi-momentum distribution in the Fermi-gas for the semi-infinite system is given by
$`n(k,t)`$ $`=`$ $`{\displaystyle \frac{a}{2\pi }}{\displaystyle \underset{j,j^{}}{}}e^{ik(x_j^{}x_j)}a_j^{}^{}a_j^{}_t`$ (47)
$`=`$ $`{\displaystyle \frac{a}{2\pi }}{\displaystyle \underset{n=1}{\overset{N}{}}}\left|{\displaystyle \underset{j}{}}e^{ikx_j}\alpha _j(k_n,t)\right|^2.`$
In Fig. 8(a,b,c) we show the evolution for a Fermi-sea with $`\nu 3/4`$, i.e. $`N=38`$ particles initially on $`M=50`$ sites. For each simulation we have $`W=\mathrm{\Delta }=0`$, but the driving varies as $`\mathrm{\Omega }/J=0,1,2`$ in Fig. 8(a,b,c) respectively. On the left side we plot the atomic density $`n(x_j,t)`$ (darker regions correspond to higher density). To the right we plot the respective momentum profiles $`n(k,t)`$ for the Fermi-gas at times $`t=0,M/2J,M/J`$ (from bottom to top). This times are indicated by arrows each figure.
In Fig. 8(a) we see the evolution of the noninteracting system, $`\mathrm{\Omega }=0`$. The atomic cloud expands freely to the right after opening the switch at $`t=0`$. The corresponding momentum distribution is initially given by $`n(k,t=0)\theta (k_F|k|)Ma/2\pi `$ (see profile at the bottom). With progressing time $`0<t<M/2J`$ the gas develops a forward peak at $`k=\pi /2a`$ (see profile in the middle) until becoming asymmetric as $`n(k,t=M/J)(M/2\pi )\theta (k_Fk)2\theta (k)`$ for $`tM/J`$ (see profile at the top).
In Fig. 8(b) we show the behavior for for weak laser driving, $`\mathrm{\Omega }=J`$. We notice that there is already substantial blocking by the impurity. The corresponding momentum profiles show that the blocking is mainly due to the complete reflection of fermions with quasi-momentum near $`k=\pi /2a`$ to $`k^{}=\pi /2a`$.
In Fig. 8(c) we plot the densities for resonant driving with $`\mathrm{\Omega }=2J`$. The transport through the impurity is efficiently blocked by the impurity atom, as the initial densities $`n(x_j,t=0)`$ and $`n(k,t=0)`$ are almost completely preserved.
In the following we will consider the number of particles on the right side of the impurity,
$`N_R(t)`$ $`=`$ $`{\displaystyle \underset{j>0}{}}n(x_j,t),`$ (48)
and the corresponding particle current through the impurity, $`I_R`$.
$`I_R(t)`$ $`=`$ $`{\displaystyle \frac{dN_R}{dt}}(t).`$ (49)
They characterize the behavior of the switch.
In Fig. 9(a) we show the number of particles $`N_R(t)`$, for the same parameters as in Fig. 8, i.e. for each the initial filling factor $`\nu 3/4`$ and $`W=\mathrm{\Delta }=0`$. The solid line shows $`N_R`$ for the no coupling to the impurity, $`\mathrm{\Omega }=0`$, and corresponds to the densities shown in Fig. 8(a). The dashed-dotted line shows the behavior for $`\mathrm{\Omega }=J`$, cf. Fig.8(b), and the dashed line corresponds to $`\mathrm{\Omega }=2J`$, i.e. Fig.8(c). Moreover, in Fig. 9(b) we show the number of particles $`N_R(t)`$, for initial filling factor $`\nu 3/4`$ and $`\mathrm{\Omega }=\mathrm{\Delta }=0`$, but for an onsite shift $`W=0,J,2J,4J`$ (see solid, dash-dotted, dashed, dotted line), respectively. In general, after a short transient period, of the order of the inverse tunneling rate $`1/J`$, the number of particles on the right side of the impurity increases linearly with $`t`$. Thereby the system establishes a roughly constant flux of particles $`I_R`$ through the impurity. The flux persists up to $`tM/J`$, which is indicated by a vertical dotted line in Fig. 9(a,b). Then the population on the left side of the impurity is substantially depleted and therefore $`N_R(t)`$ saturates until all particles tunneled through the impurity, yielding $`N_R(t)=N`$ and $`I_R(t)=0`$ for $`t\mathrm{}`$.
We are interested in the linear regime. From Eq. (46a) we obtain the constant average current as (cf. the Landauer-formula)
$`I_R={\displaystyle \frac{dN_R}{dt}}={\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{k_F}}𝑑kv_kT_k,`$ (50)
where $`v_k=2Ja\mathrm{sin}(ka)`$ is the group-velocity of the quasi-particles with quasimomentum $`k`$, $`k_F=\nu \pi /a=N\pi /(M+1)a`$ the Fermi-quasimomentum and $`T_k`$ the corresponding transmission coefficients, cf. Eq. (17). Thus the average current is obtained by integrating the Fano-profiles $`T_k=1R_k`$, see e.g. Fig 5, up to the Fermi-momentum.
For an uncoupled impurity, $`W=\mathrm{\Omega }=0`$, we have $`T_k=1`$. Thus the current is given up to a constant by the Fermi-energy,
$`I_R^{(0)}`$ $`=`$ $`{\displaystyle \frac{E_F}{2\pi }}={\displaystyle \frac{2J}{\pi }}\mathrm{sin}^2\left({\displaystyle \frac{\nu \pi }{2}}\right).`$ (51)
In Fig. 8(b) we plot the dependence of $`I_R^{(0)}`$ as a function of the filling factor $`\nu `$ as a solid line.
For a finite on-site shift, $`|W|>0`$, but no laser-driving, $`\mathrm{\Omega }=0`$, we have $`T_k<1`$, see Eq. (19) and Fig. 5(a). Thus the current through the switch decreases as
$`I_R^{(W)}`$ $`=`$ $`{\displaystyle \frac{E_F}{2\pi }}{\displaystyle \frac{W^2}{2\pi E_0}}\mathrm{arccoth}\left({\displaystyle \frac{W^2+2JE_F}{E_0E_F}}\right)`$ (52)
where $`E_0=\sqrt{W^2+4J^2}`$ is the modulus of the bound-state energy. The exponential decay of $`I_R^{(W)}`$ with increasing coupling strength $`W`$, i.e. the arccoth-term in Eq. (52), is characteristic for a system with one bound-state. The dependence of the current $`I_R^{(W)}`$ on the filling-factor $`\nu `$ is shown in Fig. 10(a). The solid line shows the non-interacting value, $`W=0`$, while the dash-dotted, dashed, dotted line correspond to $`W=J,2J,4J`$, respectively.
However, for resonant driving at $`\mathrm{\Omega }>0`$ and couplings $`W=\mathrm{\Delta }=0`$ we obtain a symmetric Fano-profile for $`R_k`$ with respect to $`k=\pi /2a`$, see Fig. 5(b). Therefore, by integrating the latter profiles we obtain the current as
$`I_R^{(\mathrm{\Omega })}`$ $`=`$ $`{\displaystyle \frac{E_F}{2\pi }}{\displaystyle \frac{E_+E_{}^2}{2\pi \left(E_+^2E_{}^2\right)}}\mathrm{arccoth}\left({\displaystyle \frac{E_{}^2+2JE_F}{E_+E_F}}\right)+`$ (53)
$`{\displaystyle \frac{E_+^2E_{}}{2\pi \left(E_+^2E_{}^2\right)}}\mathrm{arccot}\left({\displaystyle \frac{E_+^2+2JE_F}{E_{}E_F}}\right),`$
with $`E_\pm =\sqrt{\sqrt{\mathrm{\Omega }^4+4J^4}\pm 2J^2}`$. The arccoth-term in Eq. (53) gives the mean effect of the reflection as that typical of a system with one bound-state, while the oscillating arccot-term is induced by the presence of two interfering poles in the scattering matrix. The current Eq. (53) is plotted in Fig. 10(b) as a function of the initial filling factor $`\nu `$ for the uncoupled impurity, $`\mathrm{\Omega }=0`$ (solid line), for $`\mathrm{\Omega }=J`$ (dashed-dotted line), and for $`\mathrm{\Omega }=2J`$ (dashed line). For finite driving the current $`I_R^{(\mathrm{\Omega })}`$ shows a plateau at $`\nu =1/2`$, as all the Bloch-waves near $`k=\pi /2a`$ are completely reflected from the impurity (see Fig. 5(b)). From Eq. (53) we obtain that already for $`\mathrm{\Omega }4J`$ the current of particles through the impurity is completely suppressed for arbitrary filling $`\nu `$, i.e. up to Fermi-energy $`E_F=4J`$.
In the following we discuss the dependence of the current for $`\mathrm{\Omega }>0`$ on the detuning $`\mathrm{\Delta }`$. In Fig. 10(c) we show the current $`I_R`$ for $`\mathrm{\Omega }=2J`$ but still $`W=0`$ as a function of the detuning $`\mathrm{\Delta }`$ for several initial densities $`\nu `$. The solid line corresponds to commensurate initial filling, $`\nu =1`$, the dash-dotted line to half-filling, $`\nu =1/2`$, and the dotted line to a dilute Fermi-gas with $`\nu =1/4`$. The current shows a symmetric profile with a minimum at $`\mathrm{\Delta }2J\mathrm{cos}^3(\nu \pi /2)`$ and approaches its threshold value $`I_R^{(0)}`$ for $`|\mathrm{\Delta }|\mathrm{\Omega }`$. Notice that the resonance for the many-body Fermi-gas with increasing density from the bottom of the Bloch-band, $`\mathrm{\Delta }2J`$, toward the middle of the band, $`\mathrm{\Delta }=0`$.
For finite $`W=2J`$ (see Fig. 10(d)) the dependence shows an asymmetric profile and reaches its threshold value $`I_R^{(W)}`$ (cf. Eq. (52)) for large detuning, $`|\mathrm{\Delta }|\mathrm{\Omega }^2/|W|`$. We notice that although the single fermions in the Fermi-sea scatter independently, we obtain a finite current for $`\mathrm{\Omega }<J`$, even on resonance. This is caused by the fact that the various fermionic modes $`a_k`$ see the resonance a different (energy-dependent) detuning $`\mathrm{\Delta }E_k`$, which leads to a shift of the minimum (and maximum) of the transmitted current proportional to the density of the Fermi-gas, see Fig. 10(c,d). However, in the limit of strong driving $`\mathrm{\Omega }J`$ we recover the features of perfect blocking for $`\mathrm{\Delta }0`$ and of perfect transmission for $`\mathrm{\Delta }\mathrm{\Omega }^2/W`$.
### IV.2 Ideal Bose-gas
We now consider the case, where the probe atoms $`A`$ are spin-less non-interacting bosons. The Hamiltonian for the system is given by
$`H`$ $`=`$ $`J{\displaystyle \underset{j}{}}(a_j^{}a_{j+1}+\mathrm{h}.\mathrm{c}.)+\mathrm{\Delta }|MM|+`$ (54)
$`+W_Q|QQ|a_0^{}a_0+W_M|MM|a_0^{}a_0+`$
$`+\mathrm{\Omega }\left(|MQ|a_0+a_0^{}|QM|\right),`$
where the operators $`a_j^{}`$ ($`a_j`$) create (annihilate) an atom $`A`$ on site $`j`$, and obey the canonical commutation relations $`[a_i,a_j^{}]=\delta _{ij}`$ and $`[a_i,a_j]=[a_i^{},a_j^{}]=0`$. Moreover, $`|Q`$ ($`|M`$) denote the states with an atom $`Q`$ (a molecule in state $`M`$) on the impurity, and $`W_Q`$ ($`W_M`$) is the onsite-shift for an atom $`A`$ and an atom $`Q`$ (a molecule $`M`$) on the impurity. As in the previous section we will henceforth restrict ourselves to the case of equal on-site shifts $`W_M=W_QW`$. In this case we may rewrite the Hamiltonian as
$`H`$ $`=`$ $`J{\displaystyle \underset{j}{}}\left(a_j^{}a_{j+1}+a_{j+1}^{}a_j\right)+Wa_0^{}a_0+`$ (55)
$`+\mathrm{\Delta }\sigma ^+\sigma ^{}+\mathrm{\Omega }\left(\sigma ^+a_0+a_0^{}\sigma ^{}\right),`$
where the pauli operators $`\sigma ^+|MQ|`$ and $`\sigma ^{}|QM|`$ obey canonical anti-commutation relations and commute with $`a_j`$ and $`a_j^{}`$. The Hamiltonian (55) corresponds a multi-mode Jaynes-Cummings model.
In the following we will consider the scattering of a gaussian wavepacket of $`N`$ bosons $`A`$, all initially occupying the same single particle state, $`\alpha (x_j,t=0)`$, approaching the impurity atom $`Q`$ with mean quasi-momentum $`\pi /a>k_0>0`$ and width $`\delta k_0\pi /a`$. The corresponding wavefunction for the system is given by
$`|\mathrm{\Psi }(t=0)`$ $`=`$ $`|Q{\displaystyle \frac{1}{\sqrt{N!}}}\left[{\displaystyle \underset{j}{}}\alpha (x_j,t=0)a_j^{}\right]^N|\mathrm{vac},`$ (56)
$`\alpha _0(x_j,0)`$ $`=`$ $`𝒩e^{\delta k_0^2(x_jx(0))+ik_0x_j},`$ (57)
where for $`\delta k_0\pi /a`$ the normalization is given by $`𝒩^2=(2\delta k_0^2a^2/\pi )^{1/2}=n_0/N`$ in terms of the peak density of the gaussian wavepacket, $`n_0n(x(0),t=0)`$, and $`x(0)0`$ denotes the mean position of the particles $`A`$ at $`t=0`$.
For $`\mathrm{\Omega }=0`$ the equations of motion for $`a_j`$ decouple from $`\sigma ^{}`$. Therefore, we obtain the scattering of the bosons $`A`$ by the impurity, as
$`|\mathrm{\Psi }(t)`$ $`=`$ $`|Q{\displaystyle \frac{1}{\sqrt{N!}}}\left[{\displaystyle \underset{j}{}}\alpha (x_j,t)a_j^{}\right]^N|\mathrm{vac},`$ (58)
where the single-particle wavefunction for finite $`W`$ was already obtained in Sec. III. For this case all the results obtained in Sec. III hold, and we obtain e.g. the density as $`N`$ times the single particle result, $`n(x_j,t)=N|\alpha (x_j,t)|^2|`$.
#### IV.2.1 Linearization of the impurity
For $`\mathrm{\Omega }J|\mathrm{\Delta }|,|W|`$, we obtained in Sec. III that the population of the molecular state was strongly suppressed, i.e. as $`(J/\mathrm{\Omega })^2`$, and thus we approximate $`\sigma ^z1`$. Thus we linearize the spin, i.e. set $`\sigma ^+b^{}`$ and $`\sigma ^{}b`$, where $`b`$ and $`b^{}`$ obey canonical commutation relations. The scattering of the bosons $`A`$ by the impurity is given by
$`|\mathrm{\Psi }(t)`$ $``$ $`{\displaystyle \frac{1}{\sqrt{N!}}}\left[{\displaystyle \underset{j}{}}\alpha _j(t)a_j^{}+\beta (t)b^{}\right]^N|\mathrm{vac}_A,`$ (59)
where the single-particle wavefunction $`\alpha _j(t)`$ and the amplitude $`\beta (t)`$ of the molecular state $`b^{}|\mathrm{vac}\sigma ^+|Q=|M`$, were obtained in Sec. III (see Eq. (38)).
Self-consistency of the replacement $`\sigma _z1`$ requires that the obtained molecular population $`n_m(t)1`$. From the linearization we obtain the molecular population as (see App. A)
$`n_m(t)`$ $`=`$ $`{\displaystyle \frac{n_0}{4\delta k_0^2\pi }}\left|{\displaystyle 𝑑k\frac{\mathrm{\Omega }t_ke^{iE_kt}}{E_k\mathrm{\Delta }}e^{\stackrel{~}{k}^2/4\delta k_0^2i\stackrel{~}{k}x(0)}}\right|^2,`$ (60)
where $`\stackrel{~}{k}kk_0`$, and the Fourier-integral is obtained analytically e.g. by using a saddle-point method, see App. A. We find that the maximal attained molecular population, $`n_m^{}`$, is proportional to the initial density $`n_0`$ of the gas. In the case of a broad resonance, $`\mathrm{\Omega }>\delta k_0`$, we find
$`n_m^{}`$ $``$ $`{\displaystyle \frac{n_0\mathrm{\Omega }^2v_0^2/a^2}{\left(E_0\mathrm{\Delta }\right)^2v_0^2/a^2+\left[\mathrm{\Omega }^2+W\left(E_0\mathrm{\Delta }\right)^2\right]^2}},`$ (61)
with $`E_0=2J\mathrm{cos}(k_0a)`$ and $`v_0=2J\mathrm{sin}(k_0a)`$. Moreover, for a extremely narrow resonance, $`\mathrm{\Omega }<\delta k_0`$ and $`|\mathrm{\Delta }|<J`$, one obtains
$`n_m^{}`$ $``$ $`{\displaystyle \frac{n_0}{2\delta k_0^2a}}{\displaystyle \frac{\mathrm{\Omega }^2}{W^2+v_{}^2/a^2}}e^{\left(k_0k_{}\right)^2/2\delta k_0^2},`$ (62)
where $`v_{}=2J\mathrm{sin}(k_{}a)`$ and $`k_{}\mathrm{arccos}(\mathrm{\Delta }/2J)`$ denotes the position of the maximum of $`T_k/(E_k\mathrm{\Delta })^2`$.
In Fig. 11(a) we show the molecular population, $`n_m`$, as obtained by the replacement $`\sigma _z1`$. In Fig. 11(a) we plot the $`n_m(t)/n_0`$ for an incoming gaussian wavepacket with $`k_0=\pi /2a`$ and $`\delta k_0=0.01\pi /a`$ for driving $`\mathrm{\Omega }=J/4`$ (dashed line), $`\mathrm{\Omega }=J/2`$ (dashed-dotted line) and $`\mathrm{\Omega }=J`$ (solid line). In all three cases we have $`\mathrm{\Delta }=W=0`$. We see that with increasing Rabi-frequency $`\mathrm{\Omega }`$ the attained molecular-population quickly drop as $`J^2/\mathrm{\Omega }^2`$, and that the molecular-population closely resembles the density distribution $`n(x_j,t=0)`$ of the atomic-cloud $`A`$. In Fig. 11(b) we plot the maximal attained population, $`n_m^{}`$, as a function of the incoming momentum of the gas, $`k_0`$, for $`\mathrm{\Omega }=J`$ and $`W=\mathrm{\Delta }=0`$. The four lines correspond to different width $`\delta k_0`$ of the wavepacket, i.e. $`\delta k_0=0.001\pi /a`$ (solid line), $`\delta k_0=0.01\pi /a`$ (dashed line), $`\delta k_0=0.1\pi /a`$ (dash-dotted line) and $`\delta k_0=0.2\pi /a`$ (dashed line). For a given width $`\delta k_0`$ the molecular population attains its maximum for $`k_0=\pi /2a`$, i.e. where we have complete reflection of the wave-packet (see Fig. 5(b)). At the point of complete-reflection, $`k_0=\pi /2a`$ the population attains its overall maximum for a narrow momentum-distribution, i.e. for $`\delta k_00`$ we have $`n_m4n_0`$ for $`\mathrm{\Omega }=J`$ and $`\mathrm{\Delta }=W=0`$. The dependence of $`n_m^{}`$ on the detuning $`\mathrm{\Delta }`$ and the Rabi-frequency $`\mathrm{\Omega }`$ is shown in Fig. 11(c,d) for $`W=0`$ and $`W=2J`$, respectively. In both figures the gaussian wavepacket has $`k_0=\pi /2a`$ and $`\delta k_0=0.02\pi /a`$, i.e. initially extends about $`\pi /a\delta k_0=50`$ lattice sites. For $`W=0`$ we have complete reflection of the wavepacket for $`\mathrm{\Delta }=0`$, and the attained molecular population is maximal, $`n_m3n_0`$, for $`\mathrm{\Omega }2J`$ and $`\mathrm{\Delta }0`$ (see Fig. 11(c)). However, for finite $`W=2J`$ we have also complete transmission of the wavepacket, i.e. for $`\mathrm{\Delta }=\mathrm{\Omega }/W`$. From Fig. 11(d) we notice that for a given $`Omega`$ the maximal population is shifted from $`\mathrm{\Delta }0`$ towards the point, where one has complete transmission of the wavepacket, i.e. $`\mathrm{\Delta }\mathrm{\Omega }/W`$. However, the overall maximum of $`n_m^{}`$ in both cases, $`W=0`$ and $`W=2J`$, is attained for $`\mathrm{\Omega }2J`$ and for stronger driving quickly drops as $`n_m^{}n_0(2J/\mathrm{\Omega })^2`$. As the replacement $`\sigma ^{}b`$ is self-consistent for a dilute gas with densities $`n_0(n_m^{}/n_0)^1`$, we see from 11 that the approximation holds, even on resonance $`\mathrm{\Delta }=0`$, for strong driving $`\mathrm{\Omega }>4J`$ up to densities as high as $`n_05`$ and for small densities $`n_0<1`$ only fails for $`\mathrm{\Delta }0`$ and $`\mathrm{\Omega }2J`$.
#### IV.2.2 Time-dependent Variational Ansatz
In the following we use a time-dependent variational Ansatz to describe the behavior of the many-body wavefunction in near resonance $`\mathrm{\Delta }0`$ for $`\mathrm{\Omega }J`$, i.e. in the regime where the approximation $`\sigma _z1`$ fails already for small densities, $`n_0>0.05`$. As a generalization of Eq. (30) for $`N1`$ bosonic atoms $`A`$ we choose as an number-conserving Ansatz for the state of the system
$`|\mathrm{\Psi }(t)`$ $`=`$ $`c_Q(t)|Q{\displaystyle \frac{\left[a_Q(t)^{}\right]^N}{\sqrt{N!}}}|\mathrm{vac}+`$ (63)
$`c_M(t)|M{\displaystyle \frac{\left[a_M(t)^{}\right]^{N1}}{\sqrt{(N1)!}}}|\mathrm{vac},`$
where $`a_\sigma (t)=_j\alpha _{j,\sigma }(t)a_j^{}`$ represent two non-orthogonal time-dependent modes for the field of the bosonic atoms $`A`$ given that the impurity is in state $`\sigma =Q,M`$, and the amplitudes for the impurity, $`c_\sigma (t)`$, and for the bosonic wavepackets, $`\alpha _{j,\sigma }(t)`$, are normalized as $`_\sigma |c_\sigma (t)|^2=_j|\alpha _{j,\sigma }(t)|^2`$. The equation of motion for variational parameters, $`c_\sigma (t)`$ and $`\alpha _{j,\sigma }(t)`$, are obtained by minimizing the corresponding action (cf. App. B)
$`S(t)`$ $`=`$ $`{\displaystyle \frac{\dot{\mathrm{\Psi }}_t|\mathrm{\Psi }_t\mathrm{\Psi }_t|\dot{\mathrm{\Psi }}_t}{2i}}\mathrm{\Psi }_t|H|\mathrm{\Psi }_t`$ (64)
with respect to $`c_\sigma (t)^{}`$ and $`\alpha _{j,\sigma }(t)^{}`$, as a set of coupled non-linear differential equations, which we integrate numerically. Thus we obtain the dynamics of the system.
In Fig. 12 we show the obtained reflection-coefficient
$`R=\underset{t\mathrm{}}{lim}{\displaystyle \frac{1}{N}}{\displaystyle \underset{j<0}{}}n(x_j,t)`$ (65)
and the attained peak molecular population $`n_m^{}`$ for a Gaussian-wavepacket with narrow momentum $`k_0=\pi /2a`$, i.e. $`\delta k_0=0.02\pi /a\pi /a`$ for initial density $`n_0=0.05,0.1,0.15,0.20,0.25`$ (values indicated in plots). In Fig. 12(a) (Fig. 12(b)) we plot $`R`$ ($`n_m^{}`$) as a function of the Rabi frequency $`\mathrm{\Omega }`$ for $`\mathrm{\Delta }=W=0`$, i.e. when complete reflection of the wavepacket was predicted by the bosononic approximation of the spin. The reflection shows a non-linear behavior in the density for $`\mathrm{\Omega }>.1`$, i.e. decreases as $`1/n_0`$ with increasing density $`n_0`$. While for $`n_0=0.05`$ (see dotted line) we have complete reflection of the wavepacket for $`\mathrm{\Omega }>2J`$, for higher densities we still have a finite transmission at $`\mathrm{\Omega }=2J`$. From Fig. 12(a) we see that with increasing density the transmission coefficient rapidly deviates from the low(zero)-density result and approaches a linear behavior in $`\mathrm{\Omega }`$ already for $`n_0=1/4`$. In Fig. 12(a) we plot the dependence of $`n_m^{}`$ on $`\mathrm{\Omega }`$, which shows that the maximal population is attained for $`\mathrm{\Omega }J/2`$ and decreases as $`(J/\mathrm{\Omega })^2`$ for $`\mathrm{\Omega }Jn_0`$. Moreover, we see that with increasing density, $`n_0>0.10`$, the peak in the molecular population, $`n_m^{}`$, is no longer linear in the density as was predicted by linearization, cf. Eq. (61) and cf. Eq. (61), but saturates toward the unitary limit $`n_m^{}1`$. The dependence of the reflection coefficient $`R`$ on the detuning $`\mathrm{\Delta }`$ is plotted in Fig. 12(c) for $`\mathrm{\Omega }=J`$ and $`W=0`$ and in Fig. 12(c) for $`\mathrm{\Omega }=J`$ and $`W=2J`$. In the limit of a very dilute Bose-gas (see dashed lines for $`n_0=0.05`$) we obtain the single-particle result $`T=|t(k_0)|^2`$ given by Eq. (17), showing a symmetric Fano-profile for $`W=0`$ and an asymmetric Fano-profile for $`W=2J`$, (see also Fig. 5(b,d)). We notice that at such weak-driving as $`\mathrm{\Omega }=J`$ the peak (and the asymmetry) in the Fano-profiles are suppressed with increasing density. However, for strong driving $`\mathrm{\Omega }4Jn_0`$ we recover the features of complete reflection (complete transmission) through the impurity site as was already predicted by the linearization of the impurity.
### IV.3 Hard-core Bosons
We now consider the limit of a strongly interacting Bose-gas. Its Hamiltonian is given by
$`H`$ $`=`$ $`J{\displaystyle \underset{j}{}}(a_{j+1}^{}a_j+\mathrm{h}.\mathrm{c}.)+{\displaystyle \frac{U}{2}}{\displaystyle \underset{j}{}}a_j^{}a_j^{}a_ja_j+`$ (66)
$`+\mathrm{\Delta }|MM|+\mathrm{\Omega }\left(|MQ|a_0+a_0^{}|QM|\right)+`$
$`+W_Qa_0^{}|QQ|a_0+W_Ma_0^{}a_0|MM|`$
where the onsite-shift for two-bosons $`A`$ on the same site, $`U`$, by far exceeds the tunneling rate $`J`$, i.e. $`UJ`$. Since double occupation of a site by two atoms $`A`$ is strongly suppressed, we may eliminate those excitation from $`H`$, e.g. by imposing $`a_j^20`$. In the following we will focus on the limiting case $`U/J\mathrm{}`$, i.e. that of a Tonks gas. In this limit we may fermionize the Hamiltonian (66) via a Jordan-Wigner transformation (JWT) Sachdev , which maps the commuting fields for the hard-core bosons, $`a_j`$, and the pseudo-spin of the impurity, $`\sigma ^{}|MQ|`$, onto anticommuting fields, $`c_j`$ and $`f`$, respectively. The JWT is given by
$`a_j=c_j{\displaystyle \underset{l<j}{}}\left(12c_l^{}c_l\right),a_j^{}=c_j^{}{\displaystyle \underset{l<j}{}}\left(12c_l^{}c_l\right),`$ (67a)
$`\sigma ^{}=f{\displaystyle \underset{l}{}}\left(12c_l^{}c_l\right),\sigma ^+=f^{}{\displaystyle \underset{l}{}}\left(12c_l^{}c_l\right).`$ (67b)
The fields $`c_j`$ and $`f`$ describe fermionic excitations for the new joint vacuum state of the system, $`|\mathrm{vac}_{CF}|Q|\mathrm{vac}`$. We rewrite the Hamiltonian (66) in terms of the fermionic excitations, $`c_j`$ and $`f`$, and obtain
$`H`$ $`=`$ $`J{\displaystyle \underset{j<M}{}}(c_{j+1}^{}c_j+\mathrm{h}.\mathrm{c}.)+`$ (68)
$`+\mathrm{\Delta }f^{}f+\mathrm{\Omega }(1)^{\widehat{N}_R}(f^{}c_0+\mathrm{h}.\mathrm{c}.)+`$
$`+W_Qff^{}c_0^{}c_0+W_Mf^{}fc_0^{}c_0,`$
where $`\widehat{N}_R=_{j>0}c_j^{}c_j=_{j>0}a_j^{}a_j`$ is the number of particles to the right of the impurity site. The Hamiltonian for the fermionic excitations, $`c_j`$ and $`f`$, is the same as the one obtained for the Fermi-gas, cf. Eq. (40), except for the appearance of the phase-factor $`(1)^{\widehat{N}_R}`$ for the coupling $`\mathrm{\Omega }`$ to the impurity.
We proceed by detailing the time-dependent scattering of a tonks gas with $`N`$ atoms $`A`$ off the impurity atom $`Q`$. We assume that at time $`t=0`$ the atoms $`A`$ are trapped within a box of $`M`$ sites to the left on impurity site. This corresponds to a fermi-sea of the fermionic modes $`c_j`$, and the state of the system at $`t=0`$ is given by (see also Sec.IV.1)
$$|\mathrm{\Psi }(0)=\underset{n=1}{\overset{N}{}}\left[\sqrt{\frac{2}{M}}\underset{j<0}{}\mathrm{sin}\left(k_nx_j\right)c_j^{}\right]|\mathrm{vac}_{CF}$$
$$=|Q\underset{𝐣<0}{}(1)^{S(𝐣)}\underset{n=1}{\overset{N}{}}\left[\sqrt{\frac{2}{M}}\mathrm{sin}\left(k_nx_{j_n}\right)a_{j_n}^{}\right]|\mathrm{vac}_A,$$
(69)
where $`k_n=n\pi /(M+1)a`$ are the quasi-momenta of the fermionic excitations, $`𝐣=(j_1,j_2,\mathrm{},j_N)`$ (with $`j_pj_q`$) denotes the position of the $`N`$ bosons and $`(1)^{S(𝐣)}`$ the permutational sign of $`𝐣`$, i.e. $`S(𝐣)=_{j_p>j_q}1`$. Due to the cumbersomeness of the many-body wavefunction in terms of the bosonic operators $`a_j`$, it is preferable to deal within the fermionic picture and extract the quantities of interest from the correlations for the fermions. The density $`n(x_j,t)`$ of the hardcore bosons $`A`$ corresponds to the density for the fermions $`c`$, while the correlations, $`\rho (x_i,x_j,t)`$, and the momentum distribution of the Tonks gas, $`n(k,t)`$, differ from those of a Fermi-gas, as
$`n(x_j,t)`$ $`=`$ $`a_j^{}a_j=c_j^{}c_j,`$ (70a)
$`\rho (x_i,x_j,t)`$ $`=`$ $`a_i^{}a_j=c_i^{}{\displaystyle \underset{l=i}{\overset{j}{}}}\left[12c_l^{}c_l\right]c_j,`$ (70b)
$`n(k,t)`$ $`=`$ $`{\displaystyle \frac{a}{2\pi }}{\displaystyle \underset{m,n}{}}e^{ik\left(x_mx_n\right)}a_m^{}a_n`$
$`=`$ $`{\displaystyle \frac{a}{2\pi }}{\displaystyle \underset{m,n}{}}e^{ik\left(x_mx_n\right)}c_m^{}{\displaystyle \underset{l=m}{\overset{n}{}}}\left[12c_l^{}c_l\right]c_n.`$
Diagonalizing the single-particle density matrix $`\rho (x_i,x_j,t)`$ one obtain the condensate fraction $`N_0(t)`$ as the largest eigenvalue and the wavefunction of the quasi-condensate $`\psi _0(x_j,t)`$ as the corresponding eigenmode. In the following we denote density of the quasi-condensate as $`n_0(x_j,t)`$ and its momentum distribution as $`n_0(k,t)`$.
In Fig. 13 we plot the initial density for a Tonks-gas trapped on $`M=50`$ sites for various filling factors $`\nu =1/4,1/2,3/4,1`$, i.e. we have $`N=13,25,38,50`$ particles for Fig. 13(a,b,c,d), respectively. The solid lines in the plot on the left shows the density in position space, $`n(x_j)`$, and the dotted lines show the contribution of the largest eigenmode of the single-particle-density matrix $`\rho (x_i,x_j)`$, $`n_0(x_j)`$. To the right we plot the corresponding quasi-momentum distributions of the gas, $`n(k)`$ (solid line), and for the largest eigenmode of $`\rho (x_i,x_j)`$, $`n_0(k)`$ (dotted line). While the density merely resembles that of a homogeneous Fermi-gas with local filling factor $`\nu `$, the momentum distributions strongly differs from the typical Fermi-sea (cf. Fig. 8) as it shows a sharp peak at $`k=0`$, as one would expect from a condensate. However, the condensed fraction $`N_0_jn_0(x_j)`$ (see dashed lines within the same figures) is not macroscopic, as it amounts only to $`\sqrt{N}`$ particles in the gas and thus the behavior of the Tonks-gas significantly differs from that of a true BEC. In fact we notice that the momentum distribution shows considerable wings, which account for the depletion in the quasi-condensate. Moreover, with increasing density $`\nu `$ the number of particle in the quasi-condensate depletes considerably until vanishing completely for $`\nu =1`$ (see Fig. 13(d)), where the system attains a Mott-insulator with no phase-correlations.
We now detail the free ($`\mathrm{\Omega }=W_Q=0`$) evolution of the Tonks gas after having opened the switch at $`t=0`$. From Eq. (68) we obtain the state of the system at time $`t`$ to be given by (cf. Sec.IV.1)
$`|\mathrm{\Psi }_t`$ $`=`$ $`|Q{\displaystyle \underset{n=1}{\overset{N}{}}}\left[\alpha _j(k_n,t)c_j^{}\right]|\mathrm{vac},`$ (71a)
$`\alpha _j(k_n,t)`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{M}}}{\displaystyle \underset{j^{}<0}{}}U_{j,j^{}}(t)\mathrm{sin}(k_nx_j),`$ (71b)
where $`U_{j,j^{}}(t)`$ denotes the free single-particle propagator, cf. Eq. (38) with $`W=\mathrm{\Omega }=0`$. The density $`n(x_j,t)`$ of the hardcore-bosons $`A`$ corresponds to that of the fermions, given by Eq. (46a). In Fig. 14 we show to the left the densities $`n(x,t)`$ of the Tonks-gas (solid lines) and of the condensate mode $`n_0(x_j,t)`$ (dashed line) at time $`t=M/J=50/J`$ after opening the switch. To the right we plot the corresponding momentum distributions, $`n(k,t)`$ and $`n_0(k,t)`$. As in Fig. 13 the subplots Fig. 14(a,b,c,d) correspond to an initial filling factor $`\nu =1/4,1/2,3/4,1`$, respectively. The density of the Tonks-gas, $`n(x_j,t)`$, corresponds to the one obtained for a Fermi-gas (see e.g. Fig. 8(a) for $`\nu =1/2`$), and shows the spreading of the gas through the impurity. The corresponding momentum distribution of the Tonks-gas, $`n(k,t)`$, is shifted away from $`k=0`$ towards $`k>0`$ and spread in momentum space, due to the tunneling of the particles through the impurity site. Moreover, after a brief transient period starts developing an new additional peak at $`kk_F/2=\nu \pi /2`$, which gains in magnitude until reaching its maximum value at $`t=M/J`$. This corresponds to the dynamical formation of a quasi-condensate $`n_0(k,t)`$ Rigol , which now propagates as a wave-packet through the impurity with $`k\nu \pi /2`$ (see dotted lines in Fig. 14). The number of particle is the quasi-condensate, $`N_0(t)=_jn_0(x_j,t)`$, gives the magnitude of the peak and grows with progressing time until it saturates at $`tM/J`$. At this time the Tonks gas on the right side of the impurity is significantly depleted and one has $`N_0(t)\sqrt{N}`$. The dynamical formation of a quasi-condensate propagating with $`k>0`$ is clearly seen for filling factors $`\nu >1/2`$, where it momentum peak at $`k=k_F/2`$ exceeds the magnitude of the initial $`k=0`$ coherent fraction. For the Mott-insulating state (commensurate filling $`\nu =1`$) the situation is particularly clear: the mott-insulator melts through the impurity (see left side of Fig. 14(d)) and forms a quasi-condensate propagating as a coherent wavepacket with velocity $`v_k=2Ja`$ through the lattice. Thereby the initially flat develops a sharp peak at $`k=\pi /a`$ (see right side of Fig. 14(d)). The number of particles in the quasi-condensate increases until reaching $`\sqrt{N}`$ for $`t=M/J`$. It saturates as the initial Mott-insulating state is significantly depleted. Moreover, from the evolution of the density distributions in position space (see plots on the right of Fig. 14(a-d)) we see that, although the quasi-condensate only amounts to a fraction of atoms in the system, its profile $`n_0(x_j,t)`$ describes the *transmitted* part of the Tonks-gas through the impurity accurately, i.e. up to a multiplicative factor.
In the presence of a finite on-site shift, $`W_Q0`$, but with driving $`\mathrm{\Omega }=0`$, the Hamiltonian is still bilinear in the fermionic modes $`c_i^{}`$ and $`c_j`$. The evolution of the system is given by Eq. (71), but now with $`U_{j^{},j}(t)`$ being the single-particle propagator for $`W0`$. In Fig. 15(a,b,c) we show the evolution of a Tonks gas with $`\nu =1/2`$ for $`W=0`$, $`W=J`$ and $`W=2J`$, respectively. On the left side we plot the density distribution, $`n(x_j,t)`$, in the middle the momentum distribution, $`n(k,t)`$, and to the right the quasi-condensate $`n_0(x_j,t)`$. Darker regions correspond to regions of higher density. For no interaction, $`W=0`$, the particles freely tunnel through the impurity. The momentum distribution is shifted towards $`k>0`$ as well as being broadened. At $`tM/2J=25/M`$ an additional peak in the momentum distribution is formed at $`k\pi /3a`$, which corresponds to the mode of the quasi-condensate, $`n_0(x_j,t)`$, tunneling through the impurity (see Fig. 15(a)). For $`W=J`$ the tunneling of particles through the impurity site is partially blocked and the dynamical formation of a monochromatic mode with $`k>0`$ is suppressed, as is see from the. Thereby the condensate mode $`n_0(x_j,t)`$ remains mainly localize to the left side of the impurity. For $`W=2J`$ we see that only a small fraction of atoms passes through the impurity and the momentum distribution remains centered at $`k=0`$, although becoming slighty broader. We also see that the condensed fraction of atoms in the condensate is efficiently hindered from passing the impurity and remains localized to the left of the impurity site.
In the presence of a finite coupling $`\mathrm{\Omega }>0`$, the Hamiltonian is no more bilinear in $`c_j^{}`$ and $`c_j`$, due the appearance of the nonlinear factor $`(1)^{\widehat{N}_R}=_{j>0}(12c_j^{}c_j)`$, which makes a description of the time-dependent scattering in terms of the fermionic modes $`c_j`$ in the general case as difficult as integrating out the full many-body Schrödinger equation (66) for the hard-core bosons $`A`$. However, the contribution of the nonlinear factor $`(1)^{\widehat{N}_R}`$ to the dynamics is negligible for strong driving, $`\mathrm{\Omega }J,|\mathrm{\Delta }|,W`$, and also for low densities, $`\nu 1`$. In this case the number of atoms on the right $`N_R(t)1`$, and those we set $`\widehat{N}_R0`$ in Eq. (68). The state of the system is given by
$`|\mathrm{\Psi }_t`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{N}{}}}\left[\alpha _j(k_n,t)c_j^{}+\beta (k_n,t)f^{}\right]|\mathrm{vac},`$ (72a)
$`\alpha _i(k_n,t)`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{M}}}{\displaystyle \underset{j<0}{}}U_{i,j}(t)\mathrm{sin}(k_nx_j),`$ (72b)
$`\beta (k_n,t)`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{M}}}{\displaystyle \underset{j<0}{}}U_{M,j}(t)\mathrm{sin}(k_nx_j),`$ (72c)
where $`U_{\alpha ,j}(t)`$ denotes the full single-particle propagator, cf. Eq. (38). In this regime the density distribution of the Tonks-gas, $`n(x_j,t)`$, corresponds to that of a Fermi-gas (see e.g. the left side in Fig. 8).
For the general case of many bosons $`A`$, arbitrary coupling strengths, and even finite $`U`$, we refer to the exact numerical simulation given in Ref. Daley . These simulations allow to test the behavior of the gas for essentially arbitrary repulsion $`U`$ and density $`\nu `$, i.e. for the full crossover regime from a weakly interacting dilute Bose-gas up to a dense Tonks gas.
## V Conclusion
We have studied a scheme utilizing quantum interference to control the transport of atoms in a 1D optical lattice by a single impurity atom. The two internal state represent a qubit (spin-1/2), which in one spin state is perfectly transparent to the lattice gas, and in the other spin state acts as a single atom mirror, confining the lattice gas. This allows to “amplify” the state of the qubit, and provides a single-shot quantum non-demolition measurement of the state of the qubit. We have derived exact analytical expression for the scattering of a single atom by the impurity, and gave approximate expressions for the dynamics a gas of many interacting bosonic of fermionic atoms.
A numerical study of this dynamics based on time-dependent DMRG techniques, which complements the present discussion, will be presented in Ref. Daley .
###### Acknowledgements.
The authors acknowledge helpful discussion with A. J. Daley and D. Jaksch. Work in Innsbruck is supported by the Austrian Science Foundation, EU Networks, and the Institute for Quantum Information.
## Appendix A Scattering of Gaussian Wavepackets
We consider the dynamics of a bosonic $`N`$-particle state of the form
$`|\mathrm{\Psi }_t`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N!}}}\left[{\displaystyle \underset{j}{}}\alpha _t(x_j)a_j^{}+\beta _tb^{}\right]^N|\mathrm{vac},`$
where at $`t=0`$ the gaussian wave-packet is given by $`\beta _0=0`$ and
$`\alpha _0(x_j)`$ $`=`$ $`𝒩_0e^{\delta k_0^2\left(x_jx_0\right)+ik_0x_j},`$
where $`x_00`$ ($`k_0>0`$) is the mean position (momentum) and $`\delta k_0<\pi /a`$ the width in momentum space at $`t=0`$. For $`\delta k_0\pi /a`$ we may take the continuum limit $`_j𝑑x/a`$ and obtain $`𝒩_0\left(2\delta k_0^2a^2/\pi \right)^{1/4}`$, and $`\delta x_01/2\delta k_0a/\pi `$. The momentum representation of the is given by
$`\stackrel{~}{\alpha }_0(k)`$ $`=`$ $`\left(a/2\pi \right)^{1/2}{\displaystyle \underset{j}{}}e^{ikx_j}\alpha _0(x_j)`$
$``$ $`\left(2\pi \delta k_0^2\right)^{1/4}e^{\left(kk_0\right)^2/4\delta k_0^2i\left(kk_0\right)x_0}.`$
### A.1 Molecular density
For the evolution we are interested at the population of the molecular state, which follows as
$`n_m(t)`$ $`=`$ $`b^{}b_t=N\left|{\displaystyle \underset{j}{}}U_{M,j}(t)\alpha _0(x_j)\right|^2`$
$`=`$ $`N𝒩_0^2\left|{\displaystyle \underset{j}{}}U_{M,j}(t)e^{\delta k_0^2\left(x_jx_0\right)+ik_0x_j}\right|^2,`$
where $`U_{M,j}(t)=M|e^{iHt}|j`$ is the single particle propagator. We introduce the peak of the initial atomic density $`n_0n(x_0,t=0)=N𝒩_0^2N\left(2/\pi \right)^{1/2}\delta k_0a`$.
For the scattering off the particles, i.e. for $`|x_0|r_B`$, we may neglect the finite range of the bound-state and have
$`n_m(t)`$ $``$ $`N\left|{\displaystyle \underset{j}{}}{\displaystyle \frac{a}{2\pi }}{\displaystyle 𝑑ke^{iE_kt}\frac{\mathrm{\Omega }t_k}{E_k\mathrm{\Delta }}e^{ik|x_j|}\alpha _0(x_j)}\right|^2`$
$`=`$ $`N\left({\displaystyle \frac{a}{2\pi }}\right)\left|{\displaystyle 𝑑ke^{iE_kt}\frac{\mathrm{\Omega }t_k}{E_k\mathrm{\Delta }}\stackrel{~}{\alpha }_0(k)}\right|^2.`$
Thus we have
$`n_m(t)`$ $`=`$ $`n_0\left|f_t\right|^2,`$
$`f_t`$ $`=`$ $`{\displaystyle \frac{1}{2\delta k_0\sqrt{\pi }}}{\displaystyle 𝑑k\frac{\mathrm{\Omega }t_k}{E_k\mathrm{\Delta }}e^{iE_kt\left(\frac{kk_0}{2\delta k_0}\right)^2i\left(kk_0\right)x_0}}.`$
For the Fourier transform we use a saddle-point method, however, we have to distinguish which one is narrower, either the width of the wave-packet, $`\delta k_0`$, or the width of dressing profile.
### A.2 Broad Fano-profile
For a broad resonance, i.e. $`\mathrm{\Omega }t_k/(E_k\mathrm{\Delta })`$ slowly varying on the Bloch band we expand the integral around $`kk_0`$, i.e. with
$`E_kE_0+v_0\left(kk_0\right)+{\displaystyle \frac{1}{2m_0}}\left(kk_0\right)^2,`$
$`{\displaystyle \frac{\mathrm{\Omega }t_k}{E_k\mathrm{\Delta }}}{\displaystyle \frac{\mathrm{\Omega }t_0}{E_0\mathrm{\Delta }}},`$
where energy $`E_0E_k_{k=k_0}`$, velocity $`v_0E_k/k_{k=k_0}`$ and effective mass $`m_01/(^2E_k/^2k)_{k=k_0}`$. Notice that $`m_k=1/a^2E_k`$ for the explicit shape of the Bloch-band and choice of the origin in the band-middle.
Thereby we obtain
$`n_m(t)`$ $`=`$ $`n_0{\displaystyle \frac{\delta k_t}{\delta k_0}}e^{2\delta k_tx_t^2}𝒟_0,`$
where the linear propagation, spreading and the dressing are
$`x_t`$ $`=`$ $`x_0+v_0t,`$
$`\delta k_t`$ $`=`$ $`{\displaystyle \frac{\delta k_0}{\sqrt{1+\left(2\delta k_0^2a^2E_0t\right)^2}}},`$
$`𝒟_0`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }^2|t_k|^2}{\left(E_k\mathrm{\Delta }\right)^2}}_{k=k_0}`$
$`=`$ $`{\displaystyle \frac{\mathrm{\Omega }^2}{\left(E_0\mathrm{\Delta }\right)^2+\left[\mathrm{\Omega }^2+W\left(E_0\mathrm{\Delta }\right)\right]^2a^2/v_0^2}}.`$
We notice that at $`t_{}=x_0/v_0=|x_0/v_0|`$ we attain a maximal molecular density of
$`n_m(t_{})`$ $`=`$ $`n_0{\displaystyle \frac{𝒟_0}{\sqrt{1+\left(2\delta k_0^2a^2E_0t\right)^2}}}.`$
For $`E_00`$ we might neglect the broadening/spreading. We recognize that $`𝒟_\mathcal{0}`$ is maximal for such detuning $`\mathrm{\Delta }`$ and initial momentum $`k_0`$ where
$`E_{}`$ $``$ $`\mathrm{\Delta }W{\displaystyle \frac{\mathrm{\Omega }^2}{W^2+v_{}^2/a^2}},`$
is on the bloch-band. This corresponds to the position of the Fano-profile. As $`v_{}^2/a^2=4J^2E_{}^2`$ the Equation for $`E_{}`$ is implicitly cubic gives the same recursive equation. Thus the maximal density for a broad resonance $`\mathrm{\Omega }|J\mathrm{\Delta }|`$ is suppressed as
$`n_m(t_{})n_0{\displaystyle \frac{v_0^2/a^2}{\mathrm{\Omega }^2}},`$ (73)
while far off resonance $`|\mathrm{\Delta }|\mathrm{\Omega }`$ we have
$`n_m(t_{})n_0{\displaystyle \frac{(\mathrm{\Omega }/\mathrm{\Delta })^2}{1+(aW/v_0)^2}}.`$ (74)
### A.3 Narrow Fano-profile
In the second case, that of a narrow resonance, i.e. a sharp Fano profile, we have that the dressing factor is narrower than the Gaussian wavepacket, hence we approximate via a Saddle-point method. We expand the dressing function,
$`𝒟_k={\displaystyle \frac{\mathrm{\Omega }t_k}{E_k\mathrm{\Delta }}}=\left[{\displaystyle \frac{ia\mathrm{\Omega }}{v_k}}+{\displaystyle \frac{E_k\mathrm{\Delta }}{\mathrm{\Omega }}}\left(1+{\displaystyle \frac{iaW}{v_k}}\right)\right]^1,`$
around the momentum $`k_{}`$ where $`|𝒟_k|`$ is maximal,
$`𝒟_kC_0e^{+iC_1(kk)a\frac{1}{2}C_2(kk_{})^2a^2},`$ (75)
with
$`C_0`$ $`=`$ $`{\displaystyle \frac{i\mathrm{\Omega }}{\zeta }},`$
$`C_1`$ $`=`$ $`{\displaystyle \frac{\gamma _{}+iW}{\zeta }}{\displaystyle \frac{im_{}}{\gamma _{}}}\left(1+i{\displaystyle \frac{E_{}\mathrm{\Delta }}{\zeta }}\right),`$
$`C_2`$ $`=`$ $`\left({\displaystyle \frac{\gamma _{}+iW}{\zeta }}\right)^2{\displaystyle \frac{im_{}}{\gamma _{}}}\left({\displaystyle \frac{\gamma +iW}{\zeta }}+{\displaystyle \frac{2\mathrm{\Omega }^2}{\zeta ^2}}\right)`$
$`+{\displaystyle \frac{m_{}^2}{\gamma _{}^2}}\left[1+{\displaystyle \frac{\left(E_{}\mathrm{\Delta }\right)^2}{\zeta ^2}}\right],`$
$`\zeta `$ $``$ $`{\displaystyle \frac{\mathrm{\Omega }^2+\left(E_{}\mathrm{\Delta }\right)\left(Wi\gamma _{}\right)}{\gamma _{}}},`$
where $`E_{}=E(k_{})`$, $`\gamma _{}=E(k)/ka_{k=k_{}}=v(k_{})/a`$, $`m_{}=^2E(k)/(ka)^2=E(k_{})`$ are the lowest expansion coefficient of the dispersion relation $`E(k)=2J\mathrm{cos}(ka)`$ at $`k=k_{}`$.
The position of its maximum $`k_{}`$ is obtained from the expansion (75)by requiring $`\mathrm{}[C_1]=0`$ with
$`\mathrm{}[C_1]=\gamma _{}{\displaystyle \frac{\mathrm{\Omega }^2W+\left(E_{}\mathrm{\Delta }\right)\left(\gamma _{}^2+W^2\right)}{\mathrm{\Omega }^4+2\left(E_{}\mathrm{\Delta }\right)\mathrm{\Omega }^2W+\left(E_{}\mathrm{\Delta }\right)\left(\gamma _{}^2+W^2\right)}}`$
$`{\displaystyle \frac{m_{}}{\gamma _{}}}{\displaystyle \frac{\left[\mathrm{\Omega }^2+\left(E_{}\mathrm{\Delta }\right)W\right]^2}{\mathrm{\Omega }^4+2\left(E_{}\mathrm{\Delta }\right)\mathrm{\Omega }^2W+\left(E_{}\mathrm{\Delta }\right)\left(\gamma _{}^2+W^2\right)}}.`$
In the limit of interest (i.e. near the middle of the Bloch-band), we have $`|m_k/\gamma _k|=|\mathrm{cot}(ka)|1`$ thus we obtain the energy of the maximum as a series
$`E_{}\mathrm{\Delta }{\displaystyle \frac{W\mathrm{\Omega }^2}{\gamma _{}^2+W^2}}+{\displaystyle \frac{m_{}\gamma _{}^2\mathrm{\Omega }^4}{\left(\gamma _{}^2+W^2\right)^3}}+𝒪^2\left(m_{}\right).`$ (76)
Since $`\gamma _{}`$, $`m_{}=E_{}`$ all depend on $`k_{}`$ the series gives implicitly the value of $`k_{}`$. Moreover, we notice that the truncation to first order in $`m_{}`$ yields the exact result for $`W=0`$.
Then with
$`\zeta `$ $``$ $`{\displaystyle \frac{\mathrm{\Omega }^2}{\gamma _{}iW}}\left[1i{\displaystyle \frac{m_{}}{\gamma _{}}}\left({\displaystyle \frac{\mathrm{\Omega }\gamma _{}}{W^2+\gamma _{}^2}}\right)^2\right],`$
$`C_0`$ $``$ $`{\displaystyle \frac{\left(W+i\gamma _{}\right)/\mathrm{\Omega }}{1im_{}\gamma _{}\mathrm{\Omega }^2/\left(\gamma _{}^2+W^2\right)^2}},`$
$`C_1`$ $``$ $`{\displaystyle \frac{\gamma _{}^2+W^2}{\mathrm{\Omega }^2}}{\displaystyle \frac{m_{}W}{\gamma _{}^2+W^2}},`$ (77)
$`C_2`$ $``$ $`\left({\displaystyle \frac{\gamma _{}^2+W^2}{\mathrm{\Omega }^2}}\right)^2i{\displaystyle \frac{m_{}}{\gamma _{}}}{\displaystyle \frac{\gamma _{}^2W^24i\gamma _{}W}{\mathrm{\Omega }^2}},`$
we can compute the Fourier-integral and obtain
$`|f_t|^2`$ $`=`$ $`|C_0|^2|A(t)|e^{\frac{1}{2}\mathrm{}\left[A(t)\left(2\delta k_0x(t)+i\frac{k_0k_{}}{\delta k_0}\right)^2+\left(\frac{k_0k_{}}{\delta k_0}\right)^2\right]},`$
$`A(t)`$ $`=`$ $`\left[1+2\delta k_0^2a^2\left(C_2+im_{}t\right)\right]^1,`$ (78)
$`x(t)`$ $`=`$ $`x_0+v_{}taC_1.`$
For $`E_{}0`$, i.e. near the middle of the Bloch-band, we might neglect the spreading of the wave-packet, i.e. $`m_{}0`$. Then we have
$`|C_0|^2`$ $`=`$ $`C_1={\displaystyle \frac{W^2+v_{}^2/a^2}{\mathrm{\Omega }^2}}=\sqrt{C_2},`$
$`|f(t)|^2`$ $`=`$ $`{\displaystyle \frac{C_1\mathrm{exp}\left[\frac{2\delta k_0^2x(t)^2+C_1^2\left(k_0k_{}\right)^2a^2}{1+2\delta k_0^2a^2C_1^2}\right]}{1+2\delta k_{0}^{}{}_{}{}^{2}a^2C_1^2}}.`$ (80)
In the limit $`\mathrm{\Omega }0^+`$ the probability $`|f(t)|^2`$ vanishes as
$`\mathrm{max}_t|f(t)|^2`$ $``$ $`{\displaystyle \frac{\mathrm{\Omega }^2}{W^2+v_{}^2/a^2}}{\displaystyle \frac{\mathrm{exp}\left[\frac{\left(k_0k_{}\right)^2}{2\delta k_0^2}\right]}{2\delta k_{0}^{}{}_{}{}^{2}a^2}}.`$ (81)
## Appendix B Variational Ansatz
For the Ansatz (63) we obtain the action (64) as
$`\begin{array}{c}S={\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}[c_\sigma ^{}{\displaystyle \frac{c_\sigma }{t}}\mathrm{\Delta }\delta _{\sigma M}|c_\sigma |^2+N_\sigma |c_\sigma |^2{\displaystyle \underset{j}{}}\alpha _{j,\sigma }\\ \times (i{\displaystyle \frac{\alpha _{j,\sigma }}{t}}+J\alpha _{j+1,\sigma }+J\alpha _{j1,\sigma }W\delta _{j0}\alpha _{j,\sigma })]\\ \mathrm{\Omega }\sqrt{N}c_M^{}c_Q\left({\displaystyle \underset{j}{}}\alpha _{j,M}^{}\alpha _{j,Q}\right)^N\alpha _{j,Q}+\mathrm{c}.\mathrm{c}.,\end{array}`$
with $`c_\sigma c_\sigma (t)`$, $`\alpha _{j,\sigma }\alpha _{j,\sigma }(t)`$ and $`N_Q=N_M+1=N`$.
Minimizing the action $`S`$ with respect to $`c_\sigma ^{}`$ and $`\alpha _{j,\sigma }`$, we obtain after some algebra,
$`ic_M^{}{\displaystyle \frac{c_M}{t}}`$ $`=`$ $`\mathrm{\Delta }|c_M|^2+\lambda (N1){\displaystyle \frac{\lambda +\lambda ^{}}{2}},`$
$`ic_Q^{}{\displaystyle \frac{c_Q}{t}}`$ $`=`$ $`\lambda ^{}N{\displaystyle \frac{\lambda +\lambda ^{}}{2}},`$
$`i|c_M|^2{\displaystyle \frac{\alpha _{j,M}}{t}}`$ $`=`$ $`|c_M|^2\left(J{\displaystyle \underset{\pm }{}}\alpha _{j\pm 1,M}+W\delta _{j0}\alpha _{j,M}\right)`$
$`{\displaystyle \frac{\lambda \lambda ^{}}{2}}\alpha _{j,M}+{\displaystyle \frac{\lambda }{s}}\alpha _{j,Q},`$
$`i|c_Q|^2{\displaystyle \frac{\alpha _{j,Q}}{t}}`$ $`=`$ $`|c_Q|^2\left(J{\displaystyle \underset{\pm }{}}\alpha _{j\pm 1,Q}+W\delta _{j0}\alpha _{j,Q}\right)`$
$`+{\displaystyle \frac{\lambda \lambda ^{}}{2}}\alpha _{j,Q}+{\displaystyle \frac{N1}{N}}{\displaystyle \frac{\lambda ^{}}{s^{}}}\alpha _{j,M}+{\displaystyle \frac{\lambda ^{}\delta _{j0}}{N\alpha _{0,Q}^{}}}`$
where the overlap $`s`$ and the effective coupling $`\lambda `$ are given by
$`s`$ $`=`$ $`{\displaystyle \underset{j}{}}\alpha _{j,M}^{}\alpha _{j,Q},`$
$`\lambda `$ $`=`$ $`\mathrm{\Omega }\sqrt{N}c_M^{}c_Qs^{N1}\alpha _{0,Q}.`$
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# Algorithm for Lang’s Theorem
## 1. Introduction
A finite group of Lie type can be described as the rational points of a connected reductive algebraic group over a finite field. Given a structure in the algebraic group, such as a conjugacy class or a maximal torus, we want to find the corresponding structures in the finite group of Lie type. This can often be achieved with Lang’s Theorem. We provide a computationally efficient algorithm for Lang’s Theorem in split connected reductive groups. Our algorithm is randomised but guaranteed to return a correct answer, ie, it is Las Vegas in the sense of \[Bab97\]. Glasby and Howlett \[GH97\] have already solved this problem in a special case; our algorithm is inspired by their work and the proof of Lang’s Theorem given in \[Mül03\].
Throughout this paper $`k`$ is a finite field of size $`q`$ and characteristic $`p`$, and $`k_r`$ is the unique degree $`r`$ extension of $`k`$ in the algebraic closure $`\overline{k}`$. The affine space of dimension $`N`$ can be identified with $`\overline{k}^N`$. An *affine variety* $`X`$ is a subset of $`\overline{k}^N`$ that consists of the zeroes of a collection of polynomials. The variety is *defined over $`k`$* if it is closed under the action of the map $`F:\overline{k}^N\overline{k}^N`$ that takes the $`q`$th power of each component. The restriction of $`F`$ to $`X`$ is called the *(standard) Frobenius endomorphism* of $`X`$. The set of *rational points* of $`X`$ over $`k_r`$, denoted by $`X\left(k_r\right)`$, consists of those elements of $`X`$ fixed by $`F^r`$. A *nonstandard Frobenius endomorphism* is a morphism $`F^{}:XX`$ such that $`(F^{})^s=F^s`$ for some positive integer $`s`$. The elements of $`X`$ fixed by $`F^{}`$ are the rational points of a *$`k`$-form* of $`X`$. In this paper, Frobenius endomorphisms are standard unless otherwise stated.
A *linear algebraic group* is an affine variety with group multiplication and inversion given by rational functions. See, for example, \[Spr98\] for more details including the definitions of reductive and connected groups. Every linear algebraic group contains a maximal connected subgroup $`G^{}`$, the component of the identity. This subgroup is normal and $`G/G^{}`$ is finite, so for many purposes it suffices to study connected groups. An important result on linear algebraic groups over finite fields is:
###### Theorem 1.1 (Lang’s Theorem).
If $`G`$ is a connected linear algebraic group defined over the finite field $`k`$ with Frobenius map $`F`$, then the map
$$GG,aa^Fa$$
is onto.
This is equivalent to the statement that the first Galois cohomology of $`G`$ is trivial.
In this paper, we give an algorithm for Lang’s Theorem in $`k`$-split connected reductive groups described by the Steinberg presentation as in \[CMT04\]. In particular, the root datum, and hence the Cartan type, of $`G`$ is known. Reductive groups are likely to be the critical case, since the problem for an arbitrary connected linear algebraic group could be solved by working down a composition series (see Section 3) and all simple connected groups are reductive. Our main result is:
###### Theorem 1.2.
Let $`k`$ be a finite field of characteristic greater than $`3`$. Let $`G`$ be a $`k`$-split connected reductive linear algebraic group. Let $`c`$ be in $`G(k_r)`$, and suppose we are given $`s`$, the order of $`c^{F^{r1}}\mathrm{}c^Fc`$. Then we can find $`aG(k_{rs})`$ such that $`c=a^Fa`$ in Las Vegas time $`O(n^9r^2s^2\mathrm{log}^2(n)\mathrm{log}^2(q))`$ where $`n`$ is the reductive rank of $`G`$.
We can improve significantly on this result for the classical groups:
###### Theorem 1.3.
Let $`G`$ be a $`k`$-split simple connected classical group defined over $`k`$. Let $`c`$ be in $`G(k_r)`$ and suppose we are given $`s`$, the order of $`c^{F^{r1}}\mathrm{}c^Fc`$. Then we can find $`aG(k_{rs})`$ such that $`c=a^Fa`$ in Las Vegas time $`O(n^5r^2s^2\mathrm{log}^2(q))`$ where $`n`$ is the reductive rank of $`G`$.
The parameter $`s`$ measures the size of the field extension required, as explained in Section 2. In Section 3, we use the concept of $`F`$-eigenvectors to reduce to a problem involving forms of $`G`$-modules. A solution to this problem and a proof of Theorem 1.3 is given in Section 4. This solution uses the algorithm for computing a standard Chevalley basis in the Lie algebra of $`G`$ described in Section 5. The running time of this algorithm is analysed in Section 6, leading to a proof of Theorem 1.2.
## 2. Minimum field degree
Computation in large finite fields is a challenging problem (see, for example, \[LN97\]). So we start with an easy result giving the size of the field extension needed for Lang’s Theorem. We define the *minimum field degree* of $`gG`$ as the smallest $`r`$ such that $`g^{F^r}=g`$. Note that $`g`$ has minimum field degree $`r`$ if, and only if, $`k_r`$ is the smallest extension of $`k`$ such that $`g`$ is in $`G\left(k_r\right)`$.
###### Proposition 2.1.
Let $`G`$ be a connected linear algebraic group defined over $`k`$. Let $`c`$ be an element of $`G`$ with minimum field degree $`r`$ and let $`s`$ be the order of $`c^{F^{r1}}\mathrm{}c^Fc`$. If $`c=a^Fa`$ for some $`a`$ in $`G`$, then the minimum field degree of $`a`$ is $`rs`$.
###### Proof.
Let $`m`$ be the minimum field degree of $`a`$. Clearly $`k_r`$ is a subfield of $`k_m`$, so $`r`$ is a divisor of $`m`$, say $`ru=m`$. Since $`c^{F^r}=c`$, we have
$$\left(c^{F^{r1}}\mathrm{}c^Fc\right)^u=c^{F^{m1}}\mathrm{}c^Fc=a^{F^m}a^{F^{m1}}\mathrm{}a^{F^2}a^Fa^Fa=a^{F^m}a.$$
Hence $`a^{F^m}=a`$ if, and only if, $`u`$ is a multiple of $`s`$. ∎
The most important consequence of this proposition is that the minimum field degree is independent of the particular choice of $`a`$ and can be computed beforehand. In all our timings we consider $`s`$, the order of $`c^{F^{r1}}\mathrm{}c^Fc`$, to be an input of our algorithm. While it is straightforward to compute $`s`$, no polynomial time algorithm is known. The best known method for computing $`s`$ is to convert from the Steinberg presentation of $`G`$ to a faithful representation \[CMT04\] and then compute the order of the corresponding matrix using the algorithm of \[CLG97\]. If the representation has degree $`d`$, this takes Las Vegas time $`O(d^3\mathrm{log}(q)\mathrm{log}\mathrm{log}(q^d))`$ plus the time required to factor a collection of integers of the form $`q^{d_i}1`$ with $`_id_id`$.
Suppose now that $`G`$ is a $`k`$-split reductive group with reductive rank $`n`$ and semisimple rank $`\mathrm{}`$. The element $`c`$, which is the input to our algorithm, has size $`O((n+\mathrm{}^2)r\mathrm{log}(q))`$; while the element $`a`$, which is the output, has size $`O((n+\mathrm{}^2)rs\mathrm{log}(q))`$. Since $`s`$ need not be bounded by a polynomial in $`n`$, $`\mathrm{}`$, $`r`$, and $`\mathrm{log}(q)`$, there is no algorithm for Lang’s Theorem that is polynomial in the size of the input. The best we can hope for is an algorithm which is polynomial in the size of the output.
## 3. Twisted eigenvectors
We can now give an outline of our main algorithm. Let $`G=G\left(\overline{k}\right)`$ be a connected linear algebraic group defined over $`k`$. Suppose that $`V`$ is a $`G`$-module of dimension $`d`$ defined over $`k`$, so that $`F`$ acts on $`V=\overline{k}^d`$ by taking the $`q`$th power of each component. We say that $`vV`$ is an *$`F`$-eigenvector* of $`c`$ if $`v^Fc=v`$ (note that the “$`F`$-eigenvalue” is always one). The set $`E(k)`$ of all $`F`$-eigenvectors in $`V`$ is a $`k`$-space of dimension $`d`$. By Lang’s theorem, the $`k_{rs}`$-span of $`E(k)`$ must be equal to $`V(k_{rs})`$. There is a variety $`E`$ defined over $`k`$ such that $`E(k_t)`$ is the $`k_t`$-span of $`E(k)`$ for every positive integer $`t`$. Such a variety is called a *$`k`$-form* of $`V`$ \[Spr98, Section 11.1\].
The following easy lemma is the key to our recursive approach.
###### Lemma 3.1.
Let $`G`$ be a connected linear algebraic group defined over $`k`$ and let $`V`$ be a $`G`$-module defined over $`k`$ with kernel $`KG`$. Let $`c`$ be an element of $`G`$. Suppose that $`E(k)`$ is the set of $`F`$-eigenvectors of $`c`$ in $`V`$. Then $`aG`$ satisfies $`ca^FKa`$ if, and only if, $`V(k)a=E(k)`$.
###### Proof.
If $`a^Fza=c`$ for $`zK`$, then, for all $`vV(k)`$, $`va=vza=va^Fc=(va)^Fc`$ and so $`vaE(k)`$. Conversely, if $`V(k)a=E(k)`$, then, for all $`vV(k)`$, $`va=(va)^Fc=va^Fc`$ and so $`a^Fca^1K`$. ∎
We call an element $`aG(k_{rs})`$ such that $`V(k)a=E(k)`$ a *transformer* in $`G`$ for the $`k`$-form $`E`$. Our approach to solving Lang’s Theorem is outlined in Algorithm 1.
Note that $`s`$ is taken to be the order of $`c^{F^{r1}}\mathrm{}c^Fc`$ in the top-level function call. It is not necessary to recompute $`s`$ for the recursive calls since a multiple of the element order works just as well.
Suppose that $`G`$ is split reductive and let $`T_0`$ be a $`k`$-split maximal torus of $`G`$. Using the methods of \[CMT04\], we can construct a module $`V`$ which is *projectively faithful*, that is, the kernel $`K`$ is contained in the centre of $`G`$. We can now take $`H=T_0`$ in Algorithm 1, since $`Z(G)`$ is contained in every maximal torus of $`G`$. Since a split torus has an easily constructed faithful module, there is at most one recursive call for reductive groups. The same algorithm could, in principle, be used for a nonreductive connected group $`G`$: construct a simple quotient $`G/N`$, take $`V`$ to be the $`G`$-module induced by a projectively faithful module for $`G/N`$, and take $`H`$ to be the preimage in $`G`$ of the maximal torus in $`G/N`$. However, finding the normal subgroup $`N`$ and constructing the quotient $`G/N`$ are nontrivial problems which lie beyond the scope of this paper.
Algorithms for finding transformers are discussed in the next section. We now give two algorithms for computing the $`F`$-eigenspace. The most straightforward method is given in Algorithm 2.
The key is to consider $`k_{rs}`$ as a $`k`$-space of dimension $`rs`$ and to consider $`V(k_{rs})=k_{rs}^d`$ as a $`k`$-space of dimension $`drs`$. The solution is then found by linear algebra over $`k`$. Computing $`S`$ takes time $`O(r^2s^2\mathrm{log}^2(q))`$, where the second factor of $`\mathrm{log}(q)`$ is for computing $`q`$th powers. Finding $`C`$ and the fixed space takes time $`O(d^3r^3s^3\mathrm{log}(q))`$. So the overall time is $`O(d^3r^3s^3\mathrm{log}^2(q))`$.
An alternative method, due to Glasby and Howlett \[GH97\], is given in Algorithm 3.
It takes time $`O(d^2r^2s^2\mathrm{log}^2(q))`$ to apply $`F`$ to a $`d\times d`$ matrix over $`k_{rs}`$, so computing $`a`$ takes time $`O(d^3r^2s^2\mathrm{log}^2(q))`$. Each randomly chosen $`x`$ has a probability of at least $`1/4`$ of yielding an invertible element $`a`$. Since this probability is bounded away from zero as $`q`$, $`r`$, $`s`$, and $`d`$ become large, the algorithm is Las Vegas. Note that we have an algorithm for Lang’s theorem in $`\text{GL}\left(V\right)`$ if the function returns $`a`$ instead of $`V(k)a^1`$.
We now have:
###### Theorem 3.2.
Let $`G`$ be a connected linear algebraic group defined over $`k`$ and let $`V`$ be a $`G`$-module defined over $`k`$ with dimension $`d`$. Let $`c`$ be an element of $`G`$ with minimum field degree $`r`$ and let $`s`$ be the order of $`c^{F^{r1}}\mathrm{}c^Fc`$. Then we can compute a basis for the $`k`$-space $`E(k)`$ of $`F`$-eigenvectors of $`c`$ in deterministic time $`O(d^3r^3s^3\mathrm{log}^2(q))`$ or Las Vegas time $`O(d^3r^2s^2\mathrm{log}^2(q))`$.
## 4. Finding transformers
Let $`G`$ be a $`k`$-split connected reductive linear algebraic group defined over $`k`$, let $`c`$ be in $`G(k_r)`$, and let $`s`$ be the order of $`c^{F^{r1}}\mathrm{}c^Fc`$. Let $`T_0`$ be the standard $`k`$-split maximal torus of $`G`$ determined by the Steinberg presentation. Let $`V`$ be a projectively faithful $`G`$-module and compute $`E`$, the $`k`$-form of $`F`$-eigenvectors of $`c`$. In this section, we show how to find a transformer $`aG(k_{rs})`$ such that $`E(k)a=V(k)`$. First we consider two special cases: split tori and classical groups. Then we give an algorithm for an arbitrary $`k`$-split connected reductive group. The key is to consider $`k`$-bases with some additional structure that ensures that $`G`$ is transitive on all such bases (or $`G_{\mathrm{ad}}`$ is transitive in Subsection 4.3).
### 4.1. Split tori and isogeny
A *$`k`$-split torus* $`T`$ of dimension $`n`$ is just the direct product of $`n`$ copies of $`\overline{k}^\times `$ with the Frobenius endomorphism taking the $`q`$th power of each component. The standard module $`V`$ is just $`\overline{k}^n`$ with the componentwise action. Suppose $`c=(c_1,\mathrm{},c_n)T(k_r)`$ and $`E`$ is the variety of $`F`$-eigenvectors of $`c`$ in $`V`$. Splitting $`T`$ into $`n`$ copies of $`\overline{k}^\times `$ and using Theorem 3.2, we can compute $`E`$ in Las Vegas time $`O(nr^2s^2\mathrm{log}^2(q))`$. Now $`E`$ has a basis of the form $`a_1e_1,\mathrm{},a_ne_n`$ where each $`a_ik_{rs}^\times `$ and $`e_i`$ is the $`i`$th standard basis vector in $`V`$. Now $`(a_1,\mathrm{},a_n)T(k_{rs})`$ is a transformer for $`E`$. Hence we have proved:
###### Proposition 4.1.
Let $`T`$ be a $`k`$-split torus of dimension $`n`$. Let $`c`$ be in $`T(k_r)`$, and suppose we are given $`s`$, the order of $`c^{F^{r1}}\mathrm{}c^Fc`$. Then we can find an element $`a`$ in $`T(k_{rs})`$ such that $`c=a^Fa`$ in Las Vegas time $`O(nr^2s^2\mathrm{log}^2(q)).`$
Consider two connected linear algebraic groups $`G`$ and $`H`$ defined over $`k`$. Let $`\iota `$ be a homomorphism $`GH`$ defined over $`k`$ which is onto with finite kernel $`K`$. Such a map is called an *isogeny*. Now suppose $`G`$ and $`H`$ are reductive and described by a Steinberg presentation with unipotent, Weyl, and toral generators as in \[CMT04\]. If there is an isogeny $`\iota :GH`$, then both groups have the same Cartan type. Furthermore, we can assume (after composing with an automorphism) that $`\iota `$ leaves unipotent and Weyl generators unchanged, and acts by a change of basis on the toral generators. We denote the standard tori generated by the toral generators of $`G`$ (resp. $`H`$) by $`T_0`$ (resp. $`U_0`$). Note that $`KZ(G)T_0`$. An important invariant of $`\iota `$ is the exponent of the finite group $`K`$, which we denote by $`m`$.
For $`gT_0(k_r)`$, we have $`\iota (g)^{F^r}=\iota (g^{F^r})=\iota (g)`$, so $`\iota (g)U_0(k_r)`$. This image can be computed in time $`O(n^3r\mathrm{log}(q))`$ by linear algebra in $`T_0(k_r)`$.
For $`hU_0(k_r)`$, we can find $`gT_0`$ such that $`\iota (g)=h`$. Then $`\iota (g^{F^r}g)=h^{F^r}h=1`$, ie, $`g^{F^r}gK`$. Hence $`(g^{F^r}g)^m=1`$ and so $`g^mT_0(k_r)`$. Using the fact that $`T_0`$ is a direct sum of copies of $`\overline{k}^\times `$, such a $`g`$ must be in $`T_0(k_{rm})`$. This preimage can be computed in time $`O(n^3rm\mathrm{log}(q))`$ by linear algebra in $`T_0(k_{rm})`$.
###### Proposition 4.2.
Let $`G`$ and $`H`$ be $`k`$-split connected reductive linear algebraic groups defined over $`k`$ with reductive rank $`n`$. Suppose we have an isogeny $`\iota :GH`$ defined over $`k`$ and that $`m`$ is the exponent of the kernel of $`\iota `$. For $`c`$ in $`G`$ or $`H`$, let $`s(c)`$ denote the order of $`c^{F^{r1}}\mathrm{}c^Fc`$, where $`r`$ is the minimal field degree of $`c`$.
1. Lang’s theorem can be solved for $`cG(k_r)`$ in time $`O(n^3r^2s(c)^2m^2\mathrm{log}^2(q))`$ plus the time needed to solve it for some $`c^{}H(k_r)`$ with $`s(c^{})s(c)`$.
2. Lang’s theorem can be solved for $`cH(k_r)`$ in time $`O(n^3rs(c)m^2\mathrm{log}(q))`$ plus the time needed to solve it for some $`c^{}G(k_{rm})`$ with $`s(c^{})ms(c)`$.
###### Proof.
If $`cG(k_r)`$, then $`c^{}=\iota (c)`$ can be found in time $`O(n^3r\mathrm{log}(q))`$. Clearly $`s(c^{})s(c)`$. Now we can find $`a^{}H_{rs(c)}`$ such that $`a^Fa^{}=c^{}`$. Let $`aG_{rsm}`$ be a preimage of $`a^{}`$ computed in time $`O(n^3rm\mathrm{log}(q))`$. Consider $`a^Fca^1K(k_{rs(c)m})T_0(k_{rs(c)m})`$. Now
$$(a^Fca^1)^{F^{rs(c)m1}}\mathrm{}(a^Fca^1)^F(a^Fca^1)=a^{F^{rs(c)m}}(c^{F^{rs(c)m1}}\mathrm{}c^Fc)a^1=1,$$
So by Proposition 4.1, we can find $`bT_0(k_{rs(c)m})`$ such that $`a^Fca^1=b^Fb`$ in Las Vegas time $`O(nr^2s(c)^2m^2\mathrm{log}^2(q))`$. Now $`(ab)^Fab=c`$ and Part (1) follows.
If $`cH(k_r)`$, we can find an element $`c^{}G(k_{rm})`$ such that $`\iota (c^{})=c`$ in time $`O(n^3rm\mathrm{log}(q))`$. Since $`(c^{F^{r1}}\mathrm{}c^Fc^{})^{s(c)}K`$, we get $`s(c^{})ms(c)`$. We can now find $`a^{}H(k_{rs(c)m^2})`$ such that $`c^{}=a^Fa^{}`$. Then $`a=\iota (a^{})`$ can be computed in time $`O(n^3rs(c)m^2\mathrm{log}(q))`$ and $`a^Fa=c`$ and Part (2) is proved. ∎
### 4.2. Classical groups
We now show how to find transformers for the classical groups, using the standard representations. Throughout this subsection we take $`V=\overline{k}^d`$ and $`B_0`$ to be the standard basis $`e_1,\mathrm{},e_d`$ of $`V(k)`$.
The simplest case is $`G=\text{GL}_d\left(\overline{k}\right)`$. Let $`B`$ be a $`k`$-basis of $`E(k)`$. Let $`a`$ be the matrix whose rows are the elements of $`B`$. Then $`B_0a=B`$, and so $`a`$ is a transformer for $`E`$.
Now suppose $`G=\text{SL}_d\left(\overline{k}\right)`$. Given a basis $`B`$ of $`V`$, define its *volume*, denoted $`\mathrm{vol}(B)`$, to be the determinant of the matrix whose rows are the elements of $`B`$. Then $`B_0`$ has volume one and $`G`$ is transitive on all bases of volume one. Now suppose $`B`$ is a basis of $`E(k)`$, the set of $`F`$-eigenvectors of $`cG`$. Then $`B^Fc=B`$, and so
$$\mathrm{vol}(B)^F=\mathrm{vol}(B^F)=\mathrm{vol}(Bc^1)=\mathrm{vol}(B)det(c)^1=\mathrm{vol}(B),$$
and $`\mathrm{vol}(B)k`$. We can now construct a basis $`B^{}`$ of $`E(k)`$ with volume one by dividing the first element of $`B`$ by the scalar $`\mathrm{vol}(B)`$. So the matrix that takes $`B_0`$ to $`B^{}`$ is a transformer in $`G`$.
Now suppose that $`q`$ is odd and $`M`$ is a nondegenerate orthogonal or symplectic form on $`V`$ written $`(u,v)`$ for $`u,vV`$. Further suppose that $`M`$ is defined over $`k`$. Then the invariant group
$$G=\{x\text{GL}_d(\overline{k})(ux,vx)=(u,v)\}$$
is a reductive linear algebraic group defined over $`k`$. Note that $`G`$ is not necessarily split or connected however. Define the $`m\times m`$ matrix
$$A_m=\left(\begin{array}{ccc}0& & 1\\ & \mathrm{}& \\ 1& & 0\end{array}\right)$$
and let $`\delta `$ be a fixed nonsquare in $`k`$. Then the form $`M`$ has precisely one of the following Gram matrices $`M_B`$ with respect to some basis $`B`$:
* If $`M`$ is orthogonal and $`d=2\mathrm{}+1`$, then
$$M_B=A_d\text{or}\left(\begin{array}{ccc}& & A_{\mathrm{}}\\ & \delta & \\ A_{\mathrm{}}& & \end{array}\right).$$
* If $`M`$ is orthogonal and $`d=2\mathrm{}`$, then
$$M_B=A_d\text{or}\left(\begin{array}{cccc}& & & A_\mathrm{}1\\ & 1& & \\ & & \delta & \\ A_\mathrm{}1& & \end{array}\right).$$
* If $`M`$ is symplectic, then $`d=2\mathrm{}`$ and
$$M_B=\left(\begin{array}{cc}& A_{\mathrm{}}\\ & \\ A_{\mathrm{}}& \end{array}\right).$$
A *normal basis* for $`M`$ is a basis of $`V`$ such that $`M_B`$ is one of these matrices.
Given a nondegenerate symplectic or orthogonal form $`M`$ on the $`k`$-space $`U`$, Algorithm 4 constructs a normal basis for $`U`$.
The quadratic equations involved always have solutions by the standard classification theory of bilinear forms over finite fields (see \[Gro02\] for more details). Each of these equations can be solved by standard techniques in time $`\mathrm{log}^2(q)`$. Note that this construction is rational (that is, it does not use extensions of $`k`$) and takes time $`O(d^3\mathrm{log}(q)+d\mathrm{log}^2(q))`$.
If the form $`M`$ is symplectic, we are done: our transformer is simply the matrix taking a normal basis of $`V(k)`$ to a normal basis of $`E(k)`$.
If $`M`$ is orthogonal, the two normal bases may have different Gram matrices, in which case the equation in Lang’s Theorem has no solution. This is to be expected, since $`G`$ is not connected in this case. If we take $`G=\mathrm{SO}(V,M)`$, then this problem can be avoided. Without loss of generality, the standard basis $`B_0`$ is normal. Suppose that $`B`$ is a normal basis of $`E(k)`$. As with the special linear group, $`\mathrm{vol}(B)`$ is in $`k`$. Also
$$det(M_B)=det(BM_{B_0}B^T)=\mathrm{vol}(B)^2det(M_{B_0}).$$
But the two choices given above for the Gram matrix have determinants in different cosets of $`k^{\times 2}`$, hence $`M_{B_0}=M_B`$. It now remains to ensure that $`\mathrm{vol}(B)`$ is one. Now $`\mathrm{vol}(B)^2=det(M_{B_0})/det(M_B)=1`$, so suppose $`\mathrm{vol}(B)=1`$. If $`M_B=A_d`$, then exchanging the first and last vectors in $`B`$ results in a new normal basis with volume one. Otherwise, negating the $`(\mathrm{}+1)`$st vector in $`B`$ has the same effect.
Similar methods work for quadratic forms in characteristic two.
Now suppose we have a split simple classical group $`G`$ of (reductive and semisimple) rank $`n`$. Then $`G`$ is isogenous to one of the groups considered above, with $`d=O(n)`$.
If $`G`$ has type $`\mathrm{A}_n`$, then there is an isogeny map $`\iota :\text{SL}_n\left(\overline{k}\right)G`$ with $`mn+1`$. By Proposition 4.2(2), we can solve Lang’s Theorem in $`G`$ in Las Vegas time $`O(n^5r^2s^2\mathrm{log}^2(q))`$.
If $`G`$ is not of type $`\mathrm{A}_{\mathrm{}}`$, then there is a series of at most 2 isogeny maps connecting $`G`$ with one of the groups considered above. For each of these maps, $`m`$ is at most 4. By Proposition 4.2, we can solve Lang’s Theorem in Las Vegas time $`O(n^3r^2s^2\mathrm{log}^2(q))`$.
We have now proved Theorem 1.3. For groups of Cartan type $`\mathrm{G}_2`$ and $`\mathrm{F}_4`$, a similar result can probably be obtained by exploiting the structure of composition and Jordan algebras, respectively \[SV00\].
### 4.3. Adjoint representation
Now consider an arbitrary $`k`$-split connected reductive linear algebraic group $`G`$, with reductive rank $`n`$ and semisimple rank $`\mathrm{}`$. Then $`G`$ has a root datum $`(X,\mathrm{\Phi },Y,\mathrm{\Phi }^{})`$ with respect to a $`k`$-split maximal torus $`T_0`$. Here $`X`$ and $`Y`$ are free $``$-modules of dimension $`n`$ with a bilinear pairing $`,:X\times Y`$ putting them in duality. We fix dual bases $`e_1,\mathrm{},e_n`$ for $`X`$ and $`f_1,\mathrm{},f_n`$ for $`Y`$. The roots $`\mathrm{\Phi }`$ are a finite subset of $`X`$ and the coroots $`\mathrm{\Phi }^{}`$ are a finite subset of $`Y`$. There is a one-to-one correspondence $`:\mathrm{\Phi }\mathrm{\Phi }^{}`$ such that $`\alpha ,\alpha ^{}=2`$ for every $`\alpha \mathrm{\Phi }`$ . For more details see \[CMT04\].
The Lie algebra $`L=L(G)`$ is a $`G`$-module defined over $`k`$. This is called the *adjoint representation* of $`G`$ and it is projectively faithful. Now $`L(k)`$ has basis elements $`e_\alpha `$ for $`\alpha \mathrm{\Phi }`$ and $`h_iL(T_0)`$ for $`i=1,\mathrm{},n`$ with structure constants:
(1) $`[h_i,h_j]`$ $`=0,`$
(2) $`[e_\alpha ,h_i]`$ $`=\alpha ,f_ie_\alpha ,`$
(3) $`[e_\alpha ,e_\alpha ]`$ $`={\displaystyle \underset{i=1}{\overset{n}{}}}e_i,\alpha ^{}h_i,`$
(4) $`[e_\alpha ,e_\beta ]`$ $`=\{\begin{array}{cc}N_{\alpha \beta }e_{\alpha +\beta }\hfill & \text{for }\alpha +\beta \mathrm{\Phi },\hfill \\ 0\hfill & \text{for }\alpha +\beta \mathrm{\Phi }\text{}\beta \alpha \text{,}\hfill \end{array}`$
where the integral constants $`N_{\alpha \beta }`$ are defined in \[Car72\]. Such a basis is called a *Chevalley basis*.
Choose simple roots $`\alpha _1,\mathrm{},\alpha _{\mathrm{}}`$, and fix a linear ordering $`<`$ on $`\mathrm{\Phi }^+`$ which is compatible with height, ie, $`\mathrm{ht}(\alpha )<\mathrm{ht}(\beta )`$ implies that $`\alpha <\beta `$. Given a nonsimple positive root $`\xi `$, take the positive roots $`\alpha ,\beta `$ such that $`\xi =\alpha +\beta `$ and $`\alpha `$ is as small as possible with respect to the ordering on $`\mathrm{\Phi }^+`$. We call $`(\alpha ,\beta )`$ the *extraspecial pair* of $`\xi `$. We can choose a Chevalley basis of $`L`$ so that the integers $`N_{\alpha \beta }`$ are positive on extraspecial pairs by \[Car72\]. We call such a basis a *standard Chevalley basis*. Note that, as with the normal bases in Subsection 4.2, the problem of finding a standard Chevalley basis is rational over $`k`$.
The linear map $`a`$ taking the standard Chevalley basis of $`L(k)`$ to a standard Chevalley basis of $`E(k)`$ must be an automorphism of $`L(k_{rs})`$. We now need to find a transformer in $`G`$. Let $`G_{\mathrm{ad}}`$ be the adjoint group with the same Cartan type as $`G`$ and let $`\mathrm{\Gamma }`$ be the automorphism group of the Dynkin diagram of $`G`$. For each element of $`\mathrm{\Gamma }`$, fix a graph automorphism normalising $`T_0`$ and the Borel subgroup determined by $`\mathrm{\Phi }^+`$, as in \[Car72\]. Take $`Z`$ to be a complement to $`Z(L)[L,L]`$ in $`Z(L)`$; by construction, our graph automorphisms fix $`Z`$ pointwise. If the characteristic of $`k`$ is greater than $`3`$, then it follows from \[Hog82\] that
$$\mathrm{Aut}(L)=\mathrm{Aut}(Z)\times (\mathrm{\Gamma }G_{\mathrm{ad}}).$$
We can compute a decomposition $`a=z\gamma b`$ with $`z\mathrm{Aut}(Z)`$, $`\gamma `$ a graph automorphism, and $`bG_{\mathrm{ad}}(k_{rs})`$ in time $`O(d^3rs\mathrm{log}(q))`$ using a slight modification of Algorithm 6 of \[CMT04\]. Since $`L(k)z\gamma =L(k)`$, the element $`b`$ is a transformer in $`G_{\text{ad}}`$. It is easily checked on a case-by-case basis that the number of roots of $`G`$ is $`O(\mathrm{}^2)`$ and so the dimension of $`L`$ is $`O(n+\mathrm{}^2)`$. Hence Lang’s Theorem can be solved for $`cG_{\mathrm{ad}}(k_r)`$ in time $`O((n+\mathrm{}^2)^3r^2s^2\mathrm{log}^2(q))`$, once we have a standard Chevalley basis for $`E(k)`$.
Suppose now that $`G`$ is simple, so that $`\mathrm{}=n`$. Then there is an isogeny map $`GG_{\mathrm{ad}}`$ with $`m`$ at most $`n+1`$. We can now apply Proposition 4.2(2) and obtain:
###### Proposition 4.3.
Suppose that $`k`$ has characteristic greater than $`3`$. Let $`G`$ be a $`k`$-split connected simple linear algebraic group and let $`L`$ be the Lie algebra of $`G`$. Let $`c`$ be an element of $`G(k_r)`$ and suppose we are given $`s`$, the order of $`c^{F^{r1}}\mathrm{}c^Fc`$. Let $`E`$ be the variety of $`F`$-eigenvectors of $`c`$ in $`L`$. We can find $`aG(k_{rs})`$ in Las Vegas time $`O(n^8r^2s^2\mathrm{log}^2(q))`$ plus the time needed to find a standard Chevalley basis of $`E(k)`$.
We give an algorithm for finding a standard Chevalley basis in the next section. The timing of this algorithm is analysed in Section 6, leading to a proof of Theorem 1.2.
## 5. Computing a standard Chevalley basis
We now give an algorithm for constructing a Chevalley basis of the Lie algebra $`L`$ of a $`k`$-split connected reductive group $`G`$. Recall that $`L`$ is a $`p`$-Lie algebra \[Jac62, Section V.7\]. The first and hardest step is finding a $`k`$-split maximal toral $`p`$-subalgebra. This is similar to the algorithm of \[dGIR96\] for finding a Cartan subalgebra, but ensuring that the subalgebra is $`k`$-split makes things considerably more complex. Once we have a split maximal toral $`p`$-subalgebra, a Chevalley basis can be constructed using \[Car72, Section 4.2\].
Our algorithm only works for fields of characteristic $`p>3`$. Whenever possible we state results for characteristics $`2`$ and $`3`$, in the hope that the gaps in our argument for small $`p`$ can be filled later.
We assume that $`L`$ is given as a structure constant algebra, but we frequently compute in the adjoint representation. Throughout this section $`n`$ denotes the reductive rank of $`G`$, $`\mathrm{}`$ denotes the semisimple rank of $`G`$, and $`d`$ denotes the dimension of the Lie algebra $`L`$. Recall that our Steinberg presentation of $`G`$ determines a $`k`$-split maximal torus $`T_0`$.
### 5.1. Toral subalgebras
A Lie algebra $`L`$ over a field of characteristic $`p`$ is called a *$`p`$-Lie algebra* if it is equipped with a map $`p:LL`$ satisfying the axioms
(5) $`(x+y)^p`$ $`=x^p+y^p+{\displaystyle \underset{i=1}{\overset{p1}{}}}s_i(x,y),`$
(6) $`(ax)^p`$ $`=a^px^p,`$
(7) $`[xy^p]`$ $`=x(\mathrm{ad}y)^p`$
where $`x,yL`$, $`a\overline{k}`$, $`s_i`$ is defined in \[Jac62, Section V.7\], and $`a^p`$ and $`(\mathrm{ad}y)^p`$ are the usual $`p`$th powers.
Given values of the $`p`$-map on a basis of $`L`$, we can compute the values on an arbitrary element using Equations (5) and (6). But $`s_{p1}`$ involves Lie products of length $`p`$, so the time taken for this computation is not polynomial in $`\mathrm{log}(p)`$. Given $`xL`$, we can use (7) to compute the coset $`x^p+Z(L)`$ in time $`O(\mathrm{}^6\mathrm{log}(q)\mathrm{log}(p))`$, since $`dim(L/Z(L))`$ is $`O(\mathrm{}^2)`$. We also define the *$`q`$-map* by applying the $`p`$-map $`e`$ times, where $`q=p^e`$; values of this map modulo $`Z(L)`$ can be computed in time $`O(\mathrm{}^6\mathrm{log}^2(q))`$.
We say that $`xL`$ is *semisimple* if it is contained in the $`p`$-subalgebra generated by $`x^p`$. A *toral subalgebra* of $`L`$ is a subalgebra defined over $`k`$ consisting entirely of semisimple elements. Note that a toral subalgebra need not be a $`p`$-subalgebra. However every subalgebra $`H`$ of $`L`$ is contained in a minimal $`p`$-subalgebra called the *$`p`$-closure* of $`H`$ in $`L`$. The $`p`$-closure of a toral subalgebra is toral, and so a maximal toral subalgebra is automatically a $`p`$-subalgebra. An $`n`$-dimensional toral $`p`$-subalgebra $`H`$ is *$`k`$-split* if $`H(k)`$ is isomorphic, as a $`p`$-Lie algebra, to the vector space $`k^n`$ with trivial Lie product and the $`p`$-map acting componentwise.
If $`L`$ is the Lie algebra of a $`k`$-split connected reductive group $`G`$, then the values of the $`p`$-map on a Chevalley basis are
$$h_{i}^{}{}_{}{}^{p}=h_i\text{and}e_{\alpha }^{}{}_{}{}^{p}=0,$$
provided that $`p>3`$. Clearly $`H_0:=L(T_0)=h_1,\mathrm{},h_n`$ is a $`k`$-split toral subalgebra.
We say that the Lie algebra $`L`$ is *$`k`$-split* if it contains a maximal toral subalgebra which is $`k`$-split. The following theorem collects together the properties of toral subalgebras which we need.
###### Theorem 5.1.
Let $`L`$ be the $`p`$-Lie algebra of a $`k`$-split connected reductive group $`G`$.
1. $`L`$ is $`k`$-split with split maximal toral subalgebra $`H_0`$.
2. The centre of $`L`$ is a $`k`$-split toral subalgebra when $`p>2`$.
3. Every toral subalgebra of $`L`$ is abelian.
4. Every ($`k`$-split) maximal toral subalgebra of $`L`$ is the Lie algebra of a ($`k`$-split) maximal torus of $`G`$ (when $`p>2`$).
5. The maximal toral subalgebras of $`L`$ are $`G`$-conjugate.
6. The $`k`$-split maximal toral subalgebras of $`L`$ are $`G(k)`$-conjugate when $`p>2`$.
###### Proof.
In Part (1) it only remains to prove maximality, which follows from \[Hum67, Proposition 13.2\]. Part (3) is given in \[Hum78, Lemma 8.1\] for characteristic zero, but the same proof works for positive characteristic. Part (5) is Corollary 13.5 of \[Hum67\].
We now prove Part (2). Let $`\{e_\alpha ,h_i\}`$ be a Chevalley basis with respect to $`H_0`$. Suppose that $`zZ(L)`$ and write
$$z=\underset{i=1}{\overset{n}{}}t_ih_i+\underset{\alpha \mathrm{\Phi }}{}a_\alpha e_\alpha .$$
Let $`h_\alpha =_{i=1}^ne_i,\alpha ^{}h_i`$, then the coefficient of $`e_\alpha `$ in $`[z,h_\alpha ]`$ is $`2a_\alpha `$. Since $`[z,h_\alpha ]=0`$ and $`p>2`$, we get $`a_\alpha =0`$. Hence $`z`$ is in $`H_0=h_1,\mathrm{},h_n`$. Since $`H_0`$ is a split toral subalgebra, $`Z(L)`$ is also. (The idea for this proof is from \[Hog82, Lemma 6.10\].)
Every maximal toral subalgebra of $`L`$ is the Lie algebra of a maximal torus of $`G`$ by \[Hum67, Proposition 13.2\]. For $`p>2`$, split tori correspond to split toral subalgebras by \[Sel67, Theorem 9\]. Hence Part (4) is proved.
From now on we assume $`p>2`$. By \[Hum67, Proposition 13.6\], $`TL(T)`$ gives a one-to-one correspondence between maximal tori of $`G`$ and maximal toral subalgebras of $`L`$. Once again, split tori correspond to split toral $`p`$-algebras. Part (6) now follows from the corresponding result for tori. ∎
###### Corollary 5.2.
Given a subalgebra $`H`$ of $`L`$ defined over $`k`$, we can determine if $`H`$ is $`k`$-split maximal toral in time $`O(\mathrm{}^7\mathrm{log}^2(q))`$.
###### Proof.
First check that $`H`$ is abelian of dimension $`n`$, and compute $`H/Z`$. Note that $`H/Z`$ has dimension at most $`\mathrm{}`$. By Theorem 5.1(2), it suffices to determine if $`H/Z(L)`$ is a split toral algebra. This is done by testing whether $`b^q+Z(L)=b+Z(L)`$ where $`b+Z(L)`$ runs over a basis of $`L/Z(L)`$. As we argued at the beginning of this section, this takes time $`O(\mathrm{}^6\mathrm{log}^2(q))`$ for each basis element. ∎
Since semisimple elements are common in $`L(k)`$ (see Section 6) and the centraliser of such an element is reductive of rank $`n`$, we can find a maximal toral subalgebra by Algorithm 5.
The basic idea of our algorithm is to randomly select a series of increasingly split maximal toral subalgebras. We now assign a conjugacy class of $`W`$ to every maximal toral subalgebra $`H`$, which measures how split $`H`$ is. See \[Leh92\] for a more detailed version of this construction. There exists $`gG\left(\overline{k}\right)`$ such that $`H=H_{0}^{}{}_{}{}^{g}`$, by Theorem 5.1(5). Note that $`H_{0}^{}{}_{}{}^{F}=H_0`$ and $`H^F=H`$, since both are defined over $`k`$. Now
$$H_{0}^{}{}_{}{}^{g^Fg^1}=((H_{0}^{}{}_{}{}^{g})^F)^{g^1}=(H^F)^{g^1}=H^{g^1}=H_0,$$
so $`g^Fg^1`$ is in $`N_G(H_0)=N_G(T_0)`$. Let $`w`$ be the image of $`g^Fg^1`$ under projection onto the Weyl group $`W=N_G(T_0)/T_0`$. The element $`w`$ is uniquely determined by $`H`$ up to conjugacy in $`W`$.
### 5.2. Root decompositions of $`L`$
The *root decomposition* of $`L`$ with respect to $`H_0`$ is
$$L=H_0\underset{\alpha \mathrm{\Phi }}{}L_\alpha $$
where the *root space* $`L_\alpha =\{bL[b,h]=\alpha (h)b\text{ for all }hH_0\}`$ and each root $`\alpha \mathrm{\Phi }`$ is a linear functional $`H_0\overline{k}`$ defined over $`k`$. This decomposition is defined over $`k`$ by \[Sel67, Theorem 6\]. If the characteristic of $`k`$ is greater than 3, every root space has dimension one.
Let $`H`$ be a maximal toral subalgebra of $`L`$, fix $`gG\left(\overline{k}\right)`$ such that $`H=H_{0}^{}{}_{}{}^{g}`$ and let $`w`$ be the image of $`g^Fg^1`$ in $`W`$. For $`\alpha \mathrm{\Phi }`$, define $`\alpha ^g:H\overline{k}`$ by $`\alpha ^g(h)=\alpha (h^{g^1})`$. Then the root decomposition with respect to $`H`$ is
$$L=H\underset{\alpha \mathrm{\Phi }}{}L_{\alpha ^g},$$
where $`L_{\alpha ^g}=\{bL[b,h]=\alpha ^g(h)b\text{ for all }hH\}=L_{\alpha }^{}{}_{}{}^{g}`$. This decomposition is not defined over $`k`$ in general.
Fix a basis $`h_1,\mathrm{},h_n`$ of $`H`$ and let $`f=(f_1,\mathrm{},f_n)`$ be a sequence of irreducible polynomials in $`k[X]`$ with $`f_i(X)X`$ for at least one $`i`$. Define
$$L_f=\{yLyf_i(\mathrm{ad}(h_i))=0\text{ for }i=1,\mathrm{},\mathrm{}\}.$$
If $`L_f0`$, we call $`f`$ a *generalised root* and $`L_f`$ a *generalised root space*. The *generalised root decomposition* of $`L`$ with respect to $`H`$ is
$$L=H\underset{f}{}L_f,$$
where $`=(L,H)`$ is the set of generalised roots of $`L`$ with respect to $`H`$. This decomposition is defined over $`k`$. The generalised roots are computed by Algorithm 6.
Complete factorisation of a polynomial of degree $`d`$ over $`k`$ takes time at most $`O(d^3(\mathrm{log}(d)+\mathrm{log}(q))\mathrm{log}(q))`$, as shown in \[vzGG03\]. Factoring the characteristic polynomials $`g`$ is the dominant contribution to the running time of this algorithm. Since each $`g`$ has degree at most $`d`$, and the sum of the degrees of all the $`g`$s is at most $`nd`$, the algorithm takes time $`O(nd^3(\mathrm{log}(d)+\mathrm{log}(q))\mathrm{log}(q))`$.
In fact, we do not apply this algorithm directly to $`L`$, since we want our time to depend on $`\mathrm{}`$ but not on $`n`$ (this is necessary for analysing the recursion in Algorithm 8). By Theorem 5.1(2), the centre $`Z(L)`$ is contained in $`H`$. So we can construct a basis $`h_1,\mathrm{},h_n`$ for $`H(k)`$ with $`h_1,\mathrm{},h_m=Z(L)`$ central for some $`mn`$. Extend this to a basis $`B`$ of $`L(k)`$. Let $`\varphi `$ be the pullback map $`L/Z(L)L`$ which takes $`b+Z`$ to $`b`$ for all $`bB`$. Note that $`\varphi `$ is a linear map, but need not be a Lie algebra map. We compute in $`L/Z(L)`$, since it has dimension $`O(\mathrm{}^2)`$ independent of $`n`$, and the results are then transfered into $`L`$ via $`\varphi `$. However, $`L/Z(L)`$ need not be the Lie algebra of a group of Lie type, so most of our theoretical results do not apply to this quotient.
Let $``$ be the set of generalised roots of $`L/Z(L)`$ with respect to $`H/Z(L)`$. Given $`f=(f_1,\mathrm{},f_m)`$ define the sequence $`f^{}=(f_1,\mathrm{},f_m,X,\mathrm{},X)`$ of length $`n`$. It is now easy to see that $`\varphi ((L/Z(L))_f)=L_f^{}`$. Hence the generalised root decomposition of $`L`$ with respect to $`H`$ follows immediately once we have the decomposition of $`L/Z(L)`$ with respect to $`H/Z`$. Since the dimension of $`L/Z(L)`$ is $`O(\mathrm{}^2)`$, the decomposition of $`L/Z(L)`$ can be computed in time $`O(\mathrm{}^7(\mathrm{log}(\mathrm{})+\mathrm{log}(q))\mathrm{log}(q))`$.
Given a generalised root $`f(L,H)`$, the subspace $`L_f`$ is a direct sum of components $`L_{\alpha ^g}`$ of the root decomposition with respect to $`H`$. So we can partition $`\mathrm{\Phi }`$ into subsets $`\mathrm{\Phi }_f`$ such that $`L_f=_{\alpha \mathrm{\Phi }_f}L_{\alpha ^g}`$. Define the *degree* of $`f`$ to be the lowest common multiple of the degrees of the $`f_i`$. Given a generalised root $`f`$, we define
$$f_{}=((1)^{\mathrm{deg}(f_1)}f(X),\mathrm{},(1)^{\mathrm{deg}(f_n)}f(X)).$$
Clearly $`\mathrm{\Phi }_f_{}=\mathrm{\Phi }_f`$. Note that we can have $`f=f_{}`$ when the degree of $`f`$ is greater than one.
We now prove some properties of the sets $`\mathrm{\Phi }_f`$.
###### Lemma 5.3.
Let $`f`$ be a generalised root of $`L=L(G)`$ with respect to $`H`$.
1. The action of $`F`$ on $`\{L_{\alpha ^g}\alpha \mathrm{\Phi }_f\}`$ is equivalent to the action of $`w`$ on $`\mathrm{\Phi }_f`$.
2. $`\mathrm{\Phi }_f`$ is a union of orbits of $`w`$ on $`\mathrm{\Phi }`$.
3. If $`\mathrm{deg}(f)=1`$, then $`w`$ acts trivially of $`\mathrm{\Phi }_f`$. If in addition $`q>3`$, then $`\mathrm{\Phi }_f`$ contains a single root.
4. If $`\mathrm{deg}(f)=2`$ and $`f=f_{}`$, then $`w`$ acts on $`\mathrm{\Phi }_f`$ by negation.
###### Proof.
Write $`g^Fg^1=t\dot{w}`$ for some $`tT_0`$. Now
$$L_{\alpha ^g}^{}{}_{}{}^{F}=L_{\alpha }^{}{}_{}{}^{gF}=L_{\alpha }^{}{}_{}{}^{F^1gF}=L_{\alpha }^{}{}_{}{}^{g^F}=L_{\alpha }^{}{}_{}{}^{t\dot{w}g}=L_{\alpha w}^{}{}_{}{}^{g}=L_{(\alpha w)^g},$$
and so Part (1) is proved. Part (2) follows since $`L_{f}^{}{}_{}{}^{F}=L_f`$.
Part (3) holds because $`L_f`$ is a root space when $`\mathrm{deg}(f)=1`$.
Suppose $`\mathrm{deg}(f)=2`$ and $`f=f_{}`$. Let $`\alpha \mathrm{\Phi }_f`$. Then $`L_{\alpha ^g}^{}{}_{}{}^{F}=L_{(\alpha w)^g}`$ and so $`(\alpha w)^g(h_i)`$ and $`\alpha ^g(h_i)`$ are conjugate roots of $`f_i`$. But if $`\mathrm{deg}(f_i)=2`$, then $`f=f_{}`$ implies that the conjugate roots are negatives of each other. And if $`\mathrm{deg}(f_i)=1`$, then $`f=f_{}`$ implies that the only root of $`f_i`$ is zero. In either case $`(\alpha w)^g(h_i)=\alpha ^g(h_i)`$ and so $`w`$ acts by negation. ∎
### 5.3. Fundamental subalgebras
Now that we have the generalised root decomposition of $`L`$ with respect to $`H`$, we consider the subalgebra $`M_f`$ generated by a generalised root space $`L_f`$. Such subalgebras often turn out to be fundamental: We define a *(split) fundamental subgroup* of $`G`$ as a connected reductive subgroup normalised by a (split) maximal torus. A subalgebra $`M`$ of $`L`$ is *(split) fundamental* if it is the Lie algebra of a (split) fundamental group. This subgroup is denoted $`G_M`$. Fundamental subalgebras clearly normalise a maximal toral subalgebra.
The most important properties of such algebras for our purposes are:
###### Theorem 5.4.
Suppose that $`k`$ has characteristic greater than $`3`$. Let $`M`$ be a fundamental subalgebra of $`L`$ normalised by the maximal toral subalgebra $`H`$.
1. $`MH`$ is a maximal toral subalgebra of $`M`$.
2. If $`H`$ is a split maximal toral subalgebra, then $`M`$ is split fundamental.
###### Proof.
We have $`HC_{M+H}(H)C_L(H)=H`$, so $`M+H`$ has root decomposition
(8) $`M+H`$ $`=H{\displaystyle \underset{\beta }{}}M_\beta ,`$
where $`\beta `$ runs over $`\mathrm{\Phi }(M+H,H)`$, the set of roots of $`M+H`$ with respect to $`H`$. Suppose $`m+hM_\beta `$ where $`mM`$ and $`hH`$. Then, for all $`h^{}H`$, $`[m,h^{}]=[m+h,h^{}]=\alpha (h^{})(m+h)`$. But $`[m,h^{}]M`$, and so $`hM`$ and $`M_\beta M`$. By intersecting (8) with $`M`$, using the fact that each $`M_\beta `$ has dimension one, we obtain the root decomposition
$$M=(HM)\underset{\beta }{}M_\beta ,$$
and so $`HM`$ is a maximal toral subalgebra of $`M`$.
Finally if $`M`$ normalises $`H`$ and $`H`$ is split, then, by Theorem 5.1(4), $`H=L(T)`$ for some split maximal torus $`T`$ of $`G`$. Then $`G_M`$ normalises $`T`$ and Part (2) is proved. ∎
Recall that the closure $`\overline{\mathrm{\Psi }}`$ of $`\mathrm{\Psi }\mathrm{\Phi }`$ is just the set of all roots that can be written as a sum of elements of $`\mathrm{\Psi }`$. Note that if $`\overline{\mathrm{\Psi }}`$ is also closed under negation, it is a *subsystem*. If $`\overline{\mathrm{\Psi }}`$ is a subsystem, we say $`wW`$ is *inner* on $`\overline{\mathrm{\Psi }}`$ if the action of $`w`$ on $`\overline{\mathrm{\Psi }}`$ is induced by an element of $`W(\overline{\mathrm{\Psi }})`$.
###### Lemma 5.5.
Suppose that $`k`$ has odd characteristic. Let $`M`$ be the subalgebra generated by $`_{\alpha \mathrm{\Psi }}L_{\alpha ^g}`$, where $`\mathrm{\Psi }`$ is an orbit in $`\mathrm{\Phi }`$ under the action of $`wW`$.
1. $`M`$ is fundamental or soluble.
2. If $`M`$ is fundamental, then $`G_M`$ is semisimple.
3. $`M`$ is fundamental if, and only if, $`\overline{\mathrm{\Psi }}`$ is a subsystem.
4. If $`\overline{\mathrm{\Psi }}`$ is a subsystem and $`w`$ is inner on $`\overline{\mathrm{\Psi }}`$, then $`M`$ is split fundamental.
###### Proof.
Since $`[_{\alpha \mathrm{\Psi }}L_{\alpha ^g},H]_{\alpha \mathrm{\Psi }}L_{\alpha ^g}`$, we have $`[M,H]M`$ and so $`M`$ normalises $`H`$. Since $`[L_\alpha ,L_\beta ]L_{\alpha +\beta }`$ (recalling that $`L_0=H`$), we have
$$M=(HM)\underset{\alpha \overline{\mathrm{\Psi }}}{}L_{\alpha ^g}.$$
Let $`\mathrm{\Psi }=\mathrm{\Psi }_1\mathrm{}\mathrm{\Psi }_m`$ be the finest decomposition of $`\mathrm{\Psi }`$ into pairwise orthogonal subsets. Then $`\overline{\mathrm{\Psi }}=\overline{\mathrm{\Psi }}_1\mathrm{}\overline{\mathrm{\Psi }}_m`$ is also an orthogonal decomposition. Clearly $`w`$ permutes the sets $`\mathrm{\Psi }_i`$ and, since $`w`$ is transitive on $`\mathrm{\Psi }`$, it must be transitive on them. Since $`\overline{\mathrm{\Psi }}_1`$ is never orthogonal to $`\overline{\mathrm{\Psi }}_1`$, we either have $`\overline{\mathrm{\Psi }}_1=\overline{\mathrm{\Psi }}_1`$ or $`\overline{\mathrm{\Psi }}_1`$ is disjoint from $`\overline{\mathrm{\Psi }}_1`$. By the transitivity of $`w`$, whichever of these cases holds for $`\overline{\mathrm{\Psi }}_1`$, also holds for all $`\overline{\mathrm{\Psi }}_i`$. In particular, $`\overline{\mathrm{\Psi }}`$ is closed under negation iff $`\overline{\mathrm{\Psi }}_1`$ is. Let $`\psi `$ be the sum of all the elements of $`\overline{\mathrm{\Psi }}_1`$. Now $`\overline{\mathrm{\Psi }}_1`$ is closed under negation iff $`\psi =0`$ (since $`\psi =0`$ implies $`\alpha =_{\beta \overline{\mathrm{\Psi }}_1,\beta \alpha }\beta \overline{\mathrm{\Psi }}_1`$ for all $`\alpha \overline{\mathrm{\Psi }}_1`$, and the converse is trivial). We define $`M_i=H_i_{\alpha \overline{\mathrm{\Psi }}_i}L_{\alpha ^g}`$, where $`H_i`$ is the subalgebra of $`H`$ generated by $`[L_{\alpha ^g},L_{\alpha ^g}]`$ for all $`\alpha \overline{\mathrm{\Psi }}_i`$. Note that $`M=_iM_i`$.
Suppose first that $`\psi 0`$. Then the root subsystem generated by $`\mathrm{\Psi }_1`$ is just $`\overline{\mathrm{\Psi }}_1\overline{\mathrm{\Psi }}_1`$. Since this root subsystem is irreducible, $`\psi `$ induces an ordering on it which makes $`\overline{\mathrm{\Psi }_1}`$ the set of positive roots. Hence $`M_1`$ is just the Borel subalgebra of the Lie algebra of a simple group, and so must be soluble. The transitivity of $`w`$ on the sets $`\overline{\mathrm{\Psi }}_i`$ implies that $`M_i`$ is soluble for every $`i`$, and so $`M=_iM_i`$ is soluble.
If $`\psi =0`$, then $`\overline{\mathrm{\Psi }}_1`$ is an irreducible root subsystem and so $`M_1`$ is fundamental with $`G_{M_1}`$ a simple group. Hence $`M=_iM_i`$ is fundamental with $`G_M`$ a semisimple group. Parts (1), (2) and (3) are now proved.
Now suppose that $`\overline{\mathrm{\Psi }}`$ is closed under negation and $`w`$ is inner on $`\overline{\mathrm{\Psi }}`$. By Lang’s theorem in $`G_M`$, we can find $`hG_M`$ such that $`h^Fh^1=\dot{w}`$. On the other hand, $`g`$ satisfies $`g^Fg^1=t\dot{w}`$ for some $`tT_0`$. Now the map $`\dot{w}F`$ is a nonstandard Frobenius endomorphism since $`\dot{w}^F=\dot{w}`$ and so $`(\dot{w}F)^m=F^m`$, where $`m`$ is the order of $`\dot{w}`$. Furthermore $`T_{0}^{}{}_{}{}^{\dot{w}F}=T_0`$. So, by Lang’s theorem in $`T_0`$, there is a $`uT_0`$ such that $`t=u^{\dot{w}F}u^1`$. Set $`\stackrel{~}{g}=u^{\dot{w}}g`$, so that
$$\stackrel{~}{g}^F\stackrel{~}{g}^1=u^{\dot{w}F}g^Fg^1u^{\dot{w}}=u^{\dot{w}F}tu\dot{w}=\dot{w}$$
and $`H_{0}^{}{}_{}{}^{\stackrel{~}{g}}=H_{0}^{}{}_{}{}^{\dot{w}^1u^1\dot{w}g}=H_{0}^{}{}_{}{}^{g}=H`$. Hence $`h^Fh^1=\stackrel{~}{g}^F\stackrel{~}{g}^1`$, that is $`\stackrel{~}{g}h^1`$ is defined over $`k`$ and so $`H^{h^1}=H_{0}^{}{}_{}{}^{\stackrel{~}{g}h^1}`$ is split. So
$$[M,H^{h^1}]=[M^h,H]^{h^1}=[M,H]^{h^1}M^{h^1}=M$$
and $`M`$ is split fundamental by Theorem 5.4(2). ∎
An immediate application is Algorithm 7 for computing the direct sum decomposition of a the Lie algebra $`L`$. Although more that one $`w`$-orbit of $`\mathrm{\Phi }`$ can have the same generalised root, this clearly is not possible for orbits in different components of $`\mathrm{\Phi }`$.
The components returned are fundamental subalgebras.
### 5.4. Finding a split maximal toral subalgebra
Suppose now that we have found a nontrivial split fundamental subalgebra $`M`$. The following proposition shows that we can use recursion to find a split maximal toral subalgebra of $`L`$.
###### Proposition 5.6.
Suppose that the characteristic of $`k`$ is greater than $`3`$. Let $`M`$ be a split fundamental subalgebra of $`L`$. Let $`K`$ be a split maximal toral subalgebra of $`M`$. Then $`C_L(K)`$ is a split fundamental subalgebra of $`L`$ with full rank $`n`$. Hence a split maximal toral subalgebra of $`C_L(K)`$ is also a split maximal toral subalgebra of $`L`$.
###### Proof.
Let $`G_M`$ be the split fundamental subgroup of $`G`$ such that $`L(G_M)=M`$. Let $`T`$ be a split maximal torus of $`G`$ which normalises $`G_M`$ and let $`H=L(T)`$. Then $`U=G_MT`$ is a split maximal torus of $`G_M`$. By Theorem 5.1(6), we can assume without loss of generality that $`K=L(U)`$.
Let $`C=C_G(U)`$. Clearly $`C`$ is normalised by $`T`$ and it is reductive by \[Hum75, Corollary 26.2A\], hence it is split fundamental.
Let $`\mathrm{\Psi }`$ be the subset of $`\mathrm{\Phi }=\mathrm{\Phi }(G,T)`$ consisting of roots of $`G_M`$, or equivalently of $`M`$. Let $`\mathrm{\Psi }^{}`$ be the elements of $`\mathrm{\Phi }`$ which are orthogonal to all elements of $`\mathrm{\Psi }`$. Then $`\mathrm{\Psi }^{}`$ is the root system of $`C`$ and
$$L(C)=H\underset{\alpha \mathrm{\Psi }^{}}{}L_\alpha .$$
Clearly $`L(C)C_L(K)`$.
Conversely, suppose $`xC_L(K)`$. Let $`\{e_\alpha ,h_i\}`$ be a Chevalley basis of $`L`$ with respect to $`H`$ and write
$$x=\underset{i=1}{\overset{n}{}}t_ih_i+\underset{\alpha \mathrm{\Phi }}{}a_\alpha e_\alpha .$$
If $`\alpha \mathrm{\Psi }^{}`$ then there exists $`\beta \mathrm{\Psi }`$ such that $`\alpha ,\beta ^{}0`$. By the basic properties of root data $`|\alpha ,\beta ^{}|3`$, so $`\alpha ,\beta ^{}`$ is still nonzero considered as an element of $`k`$. Now $`h_\beta =[e_\beta ,e_\beta ]`$ is in $`K`$, and so $`[x,h_\beta ]=0`$. But the coefficient of $`e_\alpha `$ in $`[x,h_\beta ]`$ is $`a_\alpha \alpha ,\beta ^{}`$ by (2). Hence $`a_\alpha =0`$ for all $`\alpha \mathrm{\Psi }^{}`$ and so $`xL(C)`$.
The second conclusion is an immediate consequence of the first. ∎
We now have a method for finding split maximal toral subalgebras of $`L`$: Find a maximal toral subalgebra $`H`$, and compute its generalised roots. For each generalised root $`f`$, construct the subalgebra $`M_f`$ generated by $`L_f`$. Now, assuming that we can find $`A`$ for which $`M_A=_{fA}M_f`$ is known to be split fundamental and strictly contained in $`L`$, find a split maximal toral subalgebra $`H_A`$ of $`M_A`$. By Proposition 5.6, a split maximal toral subalgebra of $`C_L(H_A)`$ is a split maximal toral subalgebra of $`L`$. Since $`M_A`$ and $`C_L(H_A)`$ are split fundamental subalgebras of $`L`$, Theorem 5.4(2) ensures that they are also the Lie algebras of $`k`$-split connected reductive algebraic groups and so this recursion is valid. Algorithm 8 gives the precise method we use.
Note that the second argument $`Z`$ passed to this function is intended to indicate that we have a basis of $`L(k)`$ extending a basis of $`Z(k)`$, and the pullback map $`\varphi :L/ZL`$. We take $`Z=Z(L)`$ initially. In Section 6, we give a method for ensuring that $`M_A`$ is known to be split fundamental.
### 5.5. Finding a Chevalley basis
We start by giving a recognition theorem for a standard Chevalley basis.
###### Theorem 5.7.
Suppose the finite field $`k`$ has characteristicgreater than $`3`$. Let $`G`$ be a $`k`$-split connected reductive linear algebraic group defined over $`k`$ and let $`L`$ be the Lie algebra of $`G`$. Let $`H`$ be a $`k`$-split maximal toral subalgebra of $`L`$ and let $`L=H_\alpha L_\alpha `$ be the root decomposition of $`L`$. Suppose we have a basis of $`L`$ consisting of $`h_iH`$ for $`i=1,\mathrm{},n`$ and $`e_\alpha L_\alpha `$ for $`\alpha \mathrm{\Phi }`$. Further suppose this basis satisfies the equation
$$[e_\alpha ,e_\alpha ]=\underset{i=1}{\overset{n}{}}e_i,\alpha ^{}h_i$$
for every simple root $`\alpha `$ and the equations
$$[e_\alpha ,e_\beta ]=N_{\alpha \beta }e_{\alpha +\beta }\text{and}[e_\alpha ,e_\beta ]=N_{\alpha ,\beta }e_{\alpha \beta }$$
for every extraspecial pair $`(\alpha ,\beta )`$. Then this is a standard Chevalley basis.
###### Proof.
We need to prove that this basis satisfies the defining equations given in Subsection 4.3. Equation (1) follow from the fact that a toral subalgebra is abelian, Equation (1) is given, and the Equation (4) follows from \[Car72, Theorem 4.2.1\]. It remains to prove Equation (2).
For $`yY`$, define $`h_y=_{i=1}^ne_i,yh_iH(k)`$. It suffices to prove that
(9) $`[e_\alpha ,h_y]=\alpha ,ye_\alpha ,`$
for some collection of elements $`y`$ generating $`Y`$.
Now (9) is true for all $`y\mathrm{\Phi }^{}`$ by \[Car72, Theorem 4.2.1\]. If $`\alpha ,y=0`$ for all $`\alpha \mathrm{\Phi }`$, then $`h_y`$ is central and so (9) is trivially true. Together, these two kinds of element generate $`Y`$ and so we are done. ∎
A consequence of this theorem is Algorithm 9 for finding a Chevalley basis of $`L`$.
Note that for an extraspecial pair $`(\alpha ,\beta )`$, we have $`0<N_{\alpha \beta }3`$, so division by $`N_{\alpha \beta }`$ is not a problem. The basis $`\{h_i\}`$ can be computed by elementary linear algebra. Note that in the second for-loop, the roots are taken in the linear order $`<`$ of Subsection 4.3, thus ensuring that $`e_\alpha `$ and $`e_\beta `$ are already known when we compute $`e_\gamma `$.
## 6. Time analysis
Let $`L`$ be the Lie algebra of the $`k`$-split connected reductive linear algebraic group $`G`$. We now find bounds on the probability of finding a maximal toral subalgebra $`HL`$ and a set $`A`$ of generalised roots such that $`M_A`$ is known to be split fundamental. To simplify our analysis, we just bound the probability that Algorithm 5 finds a maximal toral subalgebra in a single step, or equivalently that the random element chosen is regular semisimple. Subsection 6.1 gives bounds on the frequencies of regular semisimple elements corresponding to Weyl group elements. In Section 6.2, we bound the proportion of suitable Weyl group elements. We give the proof of Theorem 1.2 in Section 6.3.
Throughout this section, $`n`$ is the reductive rank of $`G`$, $`\mathrm{}`$ is the semisimple rank of $`G`$, $`d`$ is the dimension of $`L`$, and $`d_1,\mathrm{},d_{\mathrm{}}`$ are the invariant degrees of $`G`$ as defined in \[Car72, Section 9.3\].
### 6.1. Regular semisimple elements
An element of $`L`$ is *regular semisimple* if its centraliser is a maximal toral subalgebra. For any subvariety $`S`$ of $`L`$, let $`S_{\text{rss}}`$ be the variety of regular semisimple elements in $`S`$. Recall from Subsection 5.1 that the maximal toral subalgebras of $`L`$ are classified up to $`G(k)`$-conjugacy by the conjugacy classes of $`W`$. Fix $`w`$ in $`W`$ and let $`L_{\text{rss},w}`$ be the set of elements $`xL`$ which are regular semisimple and such that there exists $`gG`$ with $`C_L(x)=H_{0}^{}{}_{}{}^{g}`$ and $`g^Fg^1T_0\dot{w}`$. Although we give direct proofs, many results in this section also follow from Gus Lehrer’s analysis of hyperplane complements \[Leh92, Leh98\].
The following result bounds our chances of finding a regular semisimple element in $`L(k)`$ whose centraliser corresponds to the $`W`$-class of a given $`w`$.
###### Proposition 6.1.
Let $`L`$ be the Lie algebra of a $`k`$-split connected reductive group $`G`$ with root datum $`(X,\mathrm{\Phi },Y,\mathrm{\Phi }^{})`$. Let $`w`$ be an element of the Weyl group $`W`$. Define
$$Q_w(X)=\frac{_{i=1}^{\mathrm{}}(1X^{d_i})}{det_Y(1wX)}.$$
Then
$$\left(1\underset{i=1}{\overset{\mathrm{}}{}}\frac{c_i}{q^i}\right)Q_w(1/q)\frac{|w^W|}{|W|}\frac{|L_{\mathrm{rss},w}(k)|}{|L(k)|}Q_w(1/q)\frac{|w^W|}{|W|}.$$
where $`c_i=c_i(w)`$ is the number of $`w`$-orbits in $`\mathrm{\Phi }`$ consisting of roots $`\alpha `$ with the property that $`i`$ is the largest integer for which $`\alpha ,\alpha w,\mathrm{},\alpha w^{i1}`$ are $`\overline{k}`$-linearly independent.
###### Proof.
Fix some $`gG`$ such that $`g^Fg^1=\dot{w}`$ and define $`H_w=H_{0}^{}{}_{}{}^{g}`$. Let $`T_w=T_{0}^{}{}_{}{}^{g}`$ so that $`L(T_w)=H_w`$. Then
$$L_{\text{rss},w}(k)=\{xL_{\text{rss}}(k)xH_w(k)^h\text{ for some }hG(k)\},$$
which is in one-to-one correspondence with
$$\{(x,H)L_{\text{rss}}(k)\times H_w(k)^{G(k)}xH\}.$$
Since $`N_{G(k)}(H_w(k))/T_w(k)C_W(w)`$, we have $`|H_w(k)^{G(k)}|=\frac{|G(k)|}{|T_w(k)||C_W(w)|}`$. Hence
$$\frac{|L_{\text{rss},w}(k)|}{|L(k)|}=|(H_w)_{\text{rss}}(k)|\frac{|G(k)|}{|L(k)||T_w(k)|}\frac{|w^W|}{|W|}.$$
Given a root $`\alpha \mathrm{\Phi }`$, define
$$H_\alpha =\{hH_w\alpha ^g(h)=0\}.$$
Then $`H_\alpha `$ is a hyperplane in $`H_w`$ and $`(H_w)_{\text{rss}}=H_w_{\alpha \mathrm{\Phi }}H_\alpha `$. Now $`H_{\alpha }^{}{}_{}{}^{F}=H_{\alpha w}`$, so $`H_\alpha (k)=\left(_jH_{\alpha w^j}\right)(k)`$. This space has codimension $`i`$, the largest integer such that $`\alpha ,\alpha w,\mathrm{},\alpha w^{i1}`$ are linearly independent. So, for each $`i=1,\mathrm{},\mathrm{}`$, we are removing $`c_i`$ subspaces of codimension $`i`$ from a $`k`$-space of dimension $`n`$. Hence
$$q^n\left(1\underset{i=1}{\overset{\mathrm{}}{}}\frac{c_i}{q^i}\right)\left|(H_w)_{\text{rss}}(k)\right|q^n.$$
Using Theorem 9.4.10 of \[Car72\] and the fact that our group is untwisted, we get
$$|G(k)|=q^d\underset{i=1}{\overset{\mathrm{}}{}}\left(1\frac{1}{q^{d_i}}\right).$$
Using Proposition 3.3.5 of \[Car93\] and the fact that $`F`$ is the standard Frobenius, we find that $`T_w(k)`$ has order $`det_Y(qIw)`$. Hence
$$\frac{|G(k)|}{|L(k)||T_w(k)|}=\frac{q^d_i(11/q^{d_i})}{q^ddet_Y(qIw)}=\frac{Q_w(1/q)}{q^n}.$$
The following useful lemma can be proved by elementary calculus.
###### Lemma 6.2.
Let $`a_1,\mathrm{},a_m`$ be a sequence of nonnegative integers and suppose that no integer appears more than $`a`$ times in this sequence. Then
$$\underset{i}{}\left(1\frac{1}{q^{a_i}}\right)\left(1\frac{1}{q}\right)^{2a}.$$
### 6.2. Reflection derangements
Recall from Subsection 5.2 that there is a relationship between between the generalised roots $`f`$ with respect to a toral subalgebra and the orbits of the corresponding Weyl group element $`w`$ on $`\mathrm{\Phi }`$. This relationship need not be a one-to-one correspondence. As we saw in Lemma 5.3(3) and (4), this relationship is almost a one-to-one correspondence when the degree of $`f`$ is one, or the degree is two and $`f=f_{}`$. This happens when there is a root $`\alpha `$ such that $`\alpha w=\pm \alpha `$. In other words, when a reflection $`s_\alpha `$ is fixed under conjugation by $`w`$.
In this section, we count the number of Weyl group elements of this kind. Given a permutation representation of a group, an element of the group is called a *derangement* with respect to the representation if it fixes no points at all. The proportion of derangements of the symmetric group $`\mathrm{Sym}_m`$ acting on $`m`$ letters is known to approach $`1/e`$ as $`m\mathrm{}`$. We give similar results for a Weyl group acting on its reflections by conjugation. We refer to these elements as *reflection derangements*. We are grateful to Anthony Henderson for helping us with the proof of this proposition.
###### Proposition 6.3.
If $`W`$ is an irreducible Coxeter group of classical type $`\mathrm{A}_{\mathrm{}}`$, $`\mathrm{B}_{\mathrm{}}/\mathrm{C}_{\mathrm{}}`$, or $`\mathrm{D}_{\mathrm{}}`$, then the proportion of its reflection derangements approaches $`2e^{3/2}`$, $`e^{5/4}`$, $`2e^{5/4}+(4e)^1`$, respectively, as $`\mathrm{}\mathrm{}`$. For exceptional types, the proportions are as listed below:
$`\mathrm{G}_2`$ $`\mathrm{F}_4`$ $`\mathrm{E}_6`$ $`\mathrm{E}_7`$ $`\mathrm{E}_8`$ $`1/3`$ $`1/4`$ $`1409/2592`$ $`1646/2835`$ $`3385549/6220800`$
###### Proof.
Denote by $`f`$ the number of reflection derangements of $`W`$. We wish to determine $`f/|W|`$.
Type $`\mathrm{A}_{\mathrm{}}`$: The Weyl group $`W(\mathrm{A}_{\mathrm{}})`$ can be identified with the symmetric group $`\mathrm{Sym}_{\mathrm{}+1}`$ on $`\mathrm{}+1`$ letters. Write $`m=\mathrm{}+1`$ and write $`d_m`$ for the proportion of permutations in $`\mathrm{Sym}_m`$ without fixed points in $`\{1,\mathrm{},m\}`$. Denote by $`R_m`$ the set of all permutations in $`\mathrm{Sym}_m`$ with at most one fixed point in $`\{1,\mathrm{},m\}`$.
An element of $`\mathrm{Sym}_m`$ does not fix a reflection if, and only if, it belongs to $`R_m`$ and does not contain a transposition $`(i,j)`$ in its cycle decomposition. So
$$f=\left|R_m\underset{1i<jm}{}R_m^{i,j}\right|$$
where
$$R_m^{ij}=\{wR_mw\text{ contains }(i,j)\}.$$
We compute $`f`$ by inclusion/exclusion. As $`R_m^{ij}`$ and $`R_m^{ij^{}}`$ intersect trivially for $`jj^{}`$ we can find $`f`$ as an alternating sum over $`h`$-tuples of commuting transpositions:
$$\underset{h=0}{\overset{m/2}{}}(1)^h\left(\genfrac{}{}{0pt}{}{m}{2h}\right)\frac{(2h)!}{2^hh!}|R_{m2h}|.$$
Since, clearly, $`|R_m|=d_m+md_{m1}/m`$,
$$f=m!\underset{h=0}{\overset{m/2}{}}\left(\frac{1}{2}\right)^h\frac{1}{h!}\left(d_{m2h}+d_{m2h1}\right).$$
As $`lim_m\mathrm{}d_m=1/e`$, the required proportion tends to
$$\underset{m\mathrm{}}{lim}\frac{f}{m!}=\underset{h=0}{\overset{\mathrm{}}{}}\left(\frac{1}{2}\right)^h\frac{1}{h!}\frac{2}{e}=e^{\frac{1}{2}}\frac{2}{e}=2e^{\frac{3}{2}}.$$
Types $`\mathrm{B}_{\mathrm{}}`$ and $`\mathrm{C}_{\mathrm{}}`$: The Weyl group $`W=W(\mathrm{B}_{\mathrm{}})=W(\mathrm{C}_{\mathrm{}})`$ can be identified with the group of all permutations $`w`$ of $`\{\pm 1,\mathrm{},\pm \mathrm{}\}`$ such that $`(i)w=(iw)`$. Define the homomorphism $`\varphi :W\mathrm{Sym}_{\mathrm{}}`$ by $`iw^\varphi =|iw|`$. Then $`wW`$ fixes no reflections if, and only if, $`w^\varphi `$ is a derangement of $`\mathrm{Sym}_{\mathrm{}}`$ and, for every transposition $`(i,j)`$ contained in the cycle decomposition of $`w^\varphi `$, either $`(i,j,i,j)`$ or $`(j,i,j,i)`$ is contained in the cycle decomposition of $`w`$.
Writing $`S_{\mathrm{}}`$ for elements of $`W`$ such that $`w^\varphi `$ is a derangement and
$$S_{\mathrm{}}^{ij}=\{wS_{\mathrm{}}w\text{ contains }(i,j)(i,j)\text{ or }(i,j)(i,j)\},$$
we find that
$$f=\left|S_{\mathrm{}}\underset{1i<j\mathrm{}}{}S_{\mathrm{}}^{ij}\right|.$$
Again, we can count $`f`$ by taking alternating sums over $`h`$-tuples of commuting transpositions in $`W^\varphi `$. As each transposition in the decomposition of an element of $`w^\varphi `$ corresponds to two 4-cycles as indicated above, we find an extra factor $`2^h`$ compared to the $`A_{\mathrm{}}`$ case:
$$\underset{h0,2h\mathrm{}}{}(1)^h\left(\genfrac{}{}{0pt}{}{\mathrm{}}{2h}\right)\frac{(2h)!}{2^hh!}2^h|S_{\mathrm{}2h}|.$$
As $`|S_{\mathrm{}}|=2^{\mathrm{}}\mathrm{}!d_{\mathrm{}}`$,
$$f=\underset{h0,2h\mathrm{}}{}\left(1\right)^h\frac{\mathrm{}!}{h!}2^{\mathrm{}2h}d_{\mathrm{}2h}.$$
As $`lim_m\mathrm{}d_m=1/e`$ and $`|W(\mathrm{B}_{\mathrm{}})|=2^{\mathrm{}}\mathrm{}!`$, the required proportion tends to
$$\underset{m\mathrm{}}{lim}\frac{f}{2^{\mathrm{}}\mathrm{}!}=\underset{h=0}{\overset{\mathrm{}}{}}\left(\frac{1}{4}\right)^h\frac{1}{h!}\frac{1}{e}=e^{\frac{1}{4}}e^1=e^{\frac{5}{4}}.$$
Type $`\mathrm{D}_{\mathrm{}}`$: The Weyl group $`W(\mathrm{D}_{\mathrm{}})`$ is the subgroup of $`W(\mathrm{B}_{\mathrm{}})`$ consisting of all elements $`w`$ such that $`_{i=1}^{\mathrm{}}iw`$ is positive. In cycle notation, this means that $`w`$ has an even number of negative cycles (that is, cycles in which both positive and negative numbers occur).
Define $`\varphi :W\mathrm{Sym}_{\mathrm{}}`$ as the restriction of the map for type $`\mathrm{B}_{\mathrm{}}`$. Then $`wW`$ does not commute with any reflection if, and only if,
1. $`w^\varphi `$ fixes at most one element of $`\{1,\mathrm{},\mathrm{}\}`$ and, for every transposition $`(i,j)`$ contained in the cycle decomposition of $`\varphi (w)`$, the cycle occurring in $`w`$ is $`(i,j,i,j)`$ or $`(j,i,j,i)`$; or
2. $`w^\varphi `$ has exactly two fixed points, say $`i`$ and $`j`$, and the cycle decomposition of $`w`$ contains $`(i,i)(j)(j)`$ or $`(i)(i)(j,j)`$.
The number of elements of the type (ii) is clearly $`\left(\genfrac{}{}{0pt}{}{\mathrm{}}{2}\right)d_\mathrm{}22^\mathrm{}2(\mathrm{}2)!`$, contributing
$$\underset{\mathrm{}\mathrm{}}{lim}\frac{\left(\genfrac{}{}{0pt}{}{\mathrm{}}{2}\right)2^\mathrm{}2d_\mathrm{}2(\mathrm{}2)!}{|W(\mathrm{D}_{\mathrm{}})|}=\underset{\mathrm{}\mathrm{}}{lim}2^2d_\mathrm{}2=\frac{1}{4e}$$
to the required asymptotic proportion.
Writing $`T_{\mathrm{}}`$ for elements of $`W`$ such that $`w^\varphi `$ fixes at most two elements and
$$T_{\mathrm{}}^{i,j}=\{wT_{\mathrm{}}w\text{ contains }(i,j)(i,j)\text{ or }(i,j)(j,i)\},$$
we find that the set of elements of type (i) is
$$T_{\mathrm{}}\underset{1i<j\mathrm{}}{}T_{\mathrm{}}^{i,j}.$$
Again, we take alternating sums over $`h`$-tuples of commuting transpositions in $`\varphi (W)`$. As each transposition in the decomposition of an element of $`\varphi (w)`$ corresponds to two 4-cycles as indicated above, we find the same factor $`2^h`$ as for the $`B_{\mathrm{}}`$ case:
$$\underset{h=0}{\overset{\mathrm{}/2}{}}(1)^h\left(\genfrac{}{}{0pt}{}{\mathrm{}}{2h}\right)\frac{(2h)!}{2^hh!}2^h|T_{\mathrm{}2h}|.$$
As $`|T_{\mathrm{}}|=2^\mathrm{}1\mathrm{}!(d_{\mathrm{}}+d_\mathrm{}1)`$, the result is
$$\underset{h=0}{\overset{\mathrm{}/2}{}}\left(\frac{1}{4}\right)^h\left(d_{\mathrm{}2h}+d_{\mathrm{}2h1}\right),$$
which contributes
$$\underset{m\mathrm{}}{lim}\frac{f}{2^{\mathrm{}}\mathrm{}!}=\underset{h=0}{\overset{\mathrm{}}{}}\left(\frac{1}{4}\right)^h\frac{1}{h!}\frac{2}{e}=2e^{\frac{1}{4}}e^1=2e^{\frac{5}{4}}$$
to the required proportion. Hence, the asymptotic proportion is $`(4e)^1+2e^{5/4}`$.
The exceptional types: These were computed by machine. ∎
###### Corollary 6.4.
The proportion of reflection derangements in a Weyl group is less than $`\frac{2}{3}`$.
###### Proof.
Recall that if $`a_n>0`$ converges monotonically to zero, then $`_{i=0}^{\mathrm{}}(1)^ia_i`$ is called an *alternating series*. The maximum value of the partial sums $`s_n=_{i=0}^n(1)^ia_i`$ of such a series is one of the first two partial sums. Since the series in the previous proposition are sums of alternating sequences, it is always possible to find a constant $`M`$ such that the maximum value of the partial sums is one of $`s_1,\mathrm{},s_M`$. It is now easy to show on a case-by-case basis that the proportion of reflection derangements in an irreducible Weyl group is at most $`\frac{2}{3}`$.
If $`W`$ is a direct product decomposition into $`s`$ irreducible Weyl groups, then an element of $`W`$ is a reflection derangement if and only if each component of $`w`$ is a reflection arrangement, and so their proportion is at most $`\left(\frac{2}{3}\right)^s\frac{2}{3}`$. ∎
Together with Proposition 6.1, this shows that the chance of finding a regular semisimple element of $`L`$ corresponding to a reflection nonderangement in the Weyl group is at least one third, provided $`q`$ is large enough. To complete the analysis, we need a more precise bound on the probability of finding certain regular semisimple elements.
### 6.3. Time analysis
We start by looking at the Coxeter class in the Weyl group. The Coxeter element is actually a reflection derangement, but this proof is the model for our next result.
###### Proposition 6.5.
Suppose that $`W`$ is an irreducible Weyl group. If $`w_c`$ is a Coxeter element of $`W`$, then
$$\frac{|L_{\mathrm{rss},w_c}(k)|}{|L(k)|}\left(1\frac{\mathrm{}}{q^{\mathrm{}/2}}\right)\left(1\frac{1}{q}\right)^4\frac{1}{h}$$
where $`h`$ is the order of $`w_c`$.
###### Proof.
Suppose $`\alpha `$ is a root and $`\alpha w_{c}^{}{}_{}{}^{m}`$ is a linear combination of $`\alpha ,\alpha w_c,\mathrm{},\alpha w_{c}^{}{}_{}{}^{m1}`$. We prove that $`m\mathrm{}/2`$ on a case-by-case basis:
Type $`\mathrm{A}_{\mathrm{}}`$: Identify $`W`$ with $`\mathrm{Sym}_{\mathrm{}+1}`$ and consider $`\mathrm{\Phi }`$ to consist of roots $`e_ie_j`$ with $`ij`$. We can take $`w_c=(1,2,\mathrm{},\mathrm{}+1)`$ and $`\alpha =e_ie_j`$. So $`\alpha w_{c}^{}{}_{}{}^{m}=e_{i+m}e_{j+m}`$ with the subscripts taken modulo $`\mathrm{}+1`$. Hence $`\alpha w_{c}^{}{}_{}{}^{m}`$ is a linear combination of $`\alpha ,\mathrm{},\alpha w_{c}^{}{}_{}{}^{m}`$ iff $`i+m`$ and $`j+m`$ are both in $`[i,i+m1][j,j+m1]`$ modulo $`\mathrm{}+1`$. By the pigeon hole principle, this can only happen if $`m(\mathrm{}+1)/2`$.
Type $`\mathrm{C}_{\mathrm{}}`$: Identify $`W`$ with the set of permutations $`w`$ of $`\{\pm 1,\mathrm{},\pm \mathrm{}\}`$ such that $`(i)w=(iw)`$ for $`i=1,\mathrm{},\mathrm{}`$. Consider $`\mathrm{\Phi }=\mathrm{\Phi }(C_{\mathrm{}})`$ to consist of roots $`\epsilon e_i\delta e_j`$ with $`\epsilon ,\delta \{\pm 1\}`$, $`i,j=1,\mathrm{},\mathrm{}`$ and $`\epsilon i\delta j`$. We can take
$$w_c=(1,2,\mathrm{},\mathrm{},1,2,\mathrm{},\mathrm{})$$
and $`\alpha =\epsilon e_i\delta e_j`$. The same argument used in type $`\mathrm{A}_{\mathrm{}}`$ now shows that $`m\mathrm{}/2`$.
Type $`\mathrm{B}_{\mathrm{}}`$: The permutation action of $`W(\mathrm{B}_{\mathrm{}})`$ on its roots is isomorphic to the action of $`W(\mathrm{C}_{\mathrm{}})`$ on its roots, so the same argument works.
Types $`\mathrm{D}_{\mathrm{}}`$: Identify $`W`$ with the elements of $`W(\mathrm{C}_{\mathrm{}})`$ such that $`_{i=1}^{\mathrm{}}(iw)>0`$ and consider $`\mathrm{\Phi }`$ to consist of the roots $`\epsilon e_i\delta e_j`$ with $`\epsilon ,\delta \{\pm 1\}`$, $`i,j=1,\mathrm{},\mathrm{}`$ and $`ij`$. We can take $`w_c=(1,2,\mathrm{},\mathrm{}1,1,2,\mathrm{},\mathrm{}+1)(\mathrm{},\mathrm{})`$ and $`\alpha =\epsilon e_i\delta e_j`$. Once again $`m\mathrm{}/2`$ if $`i,j\mathrm{}`$. If $`i=\mathrm{}`$, $`j\mathrm{}`$, then $`\alpha w_{c}^{}{}_{}{}^{m}=(1)^m\epsilon e_{\mathrm{}}\delta e_{j+m}`$ with the second subscript taken modulo $`\mathrm{}1`$, and so $`m\mathrm{}1`$.
Exceptional types: These are easily checked by computer.
It is well known that every orbit of $`w_c`$ on $`\mathrm{\Phi }`$ has size $`h`$, so $`_ic_i(w_c)=2N/h=\mathrm{}`$. We have shown that $`c_i(w_c)=0`$ for $`i<\mathrm{}/2`$, so
$$1\underset{i=1}{\overset{\mathrm{}}{}}\frac{c_i(w_c)}{q^i}1\frac{\mathrm{}}{q^{\mathrm{}/2}}.$$
The functions $`Q_{w_c}(X)`$ are straightforward to compute and are given in Table 1.
The terms in which every coefficient is positive can be ignored, since they are bounded below by 1 when we set $`X=1/q`$. Since no term $`1X^a`$ appears more than twice in these polynomials and $`q3`$, it follows by Lemma 6.2 that $`Q_{w_c}(1/q)(11/q)^4`$.
The required inequality now follows from the first inequality of Proposition 6.1 and the fact that the centraliser of $`w_c`$ has order $`h`$. ∎
We now consider reflection nonderangements that are, in some sense, close to being Coxeter elements.
###### Proposition 6.6.
Suppose that $`W`$ is an irreducible Weyl group of rank greater than one. If $`W`$ is classical with rank at least $`7`$ then there is a reflection nonderangement $`w`$ such that
$$\frac{|L_{\mathrm{rss},w}(k)|}{|L(k)|}\left(1\frac{3}{q}\frac{4}{q^2}\frac{\mathrm{}+5}{q^{(\mathrm{}2)/2}}\right)\left(1\frac{1}{q}\right)^6\frac{1}{4\mathrm{}}.$$
For other Cartan types there is a reflection nonderangement $`w`$ such that
$$\frac{|L_{\mathrm{rss},w}(k)|}{|L(k)|}\left(1\underset{i=1}{}\frac{c_i}{q^i}\right)\left(1\frac{1}{q}\right)^6\frac{1}{c}.$$
with the constants $`c`$ and $`c_i`$ listed in Table 2.
###### Proof.
Fix a root $`\beta `$. Assume $`\beta `$ is short (resp. long) for Cartan type $`\mathrm{B}_{\mathrm{}}`$ (resp. $`\mathrm{C}_{\mathrm{}}`$). Let $`\mathrm{\Phi }_\beta =\{\gamma \mathrm{\Phi }\gamma ,\beta ^{}=0\}.`$ Then $`\mathrm{\Phi }_\beta `$ is a subsystem of $`W`$ and, except in type $`\mathrm{D}_4`$, it has at most two irreducible components. Let $`\mathrm{\Phi }_\beta ^{}`$ be the irreducible summand of $`\mathrm{\Phi }_\beta `$ of maximal rank. Let $`s_\beta `$ be the reflection in $`\beta `$ and let $`w_\beta `$ be the Coxeter element of $`W(\mathrm{\Phi }_\beta ^{})`$. We take $`w=s_\beta w_\beta `$, except for type $`\mathrm{A}_1`$ where we use $`w=1`$, type $`\mathrm{G}_2`$ where we use $`w=s_\beta `$, and type $`\mathrm{D}_4`$ where we use $`s_1s_2s_1s_3s_2s_1s_4s_2s_1s_3s_2`$. (Here $`s_i`$ is the $`i`$th simple reflection, with the numbering given in \[Bou75\].) These elements are all reflection nonderangements.
First we prove that
$$\underset{i=1}{}\frac{c_i}{q^i}\frac{3}{q}+\frac{4}{q^2}+\frac{\mathrm{}+5}{q^{(\mathrm{}2)/2}}$$
for the classical types of rank at least 7.
Type $`\mathrm{A}_{\mathrm{}}`$: Assume $`\beta =e_1e_2`$. Then $`\mathrm{\Phi }_\beta `$ has type $`\mathrm{A}_\mathrm{}2`$, and so orbits within $`\mathrm{\Phi }_\beta `$ contribute at most $`\frac{\mathrm{}2}{q^{(\mathrm{}2)/2}}`$ to the sum, as in the previous proof. If $`\alpha \mathrm{\Phi }_\beta `$, then $`\alpha =\pm (e_ie_j)`$ where $`i=1`$ or $`2`$ and $`j>2`$. These roots form one orbit of size $`2`$ and either two orbits of size $`\mathrm{}1`$ or one orbit of size $`2(\mathrm{}1)`$. So these orbits contribute at most $`1/q+2/q^{\mathrm{}}`$.
Type $`\mathrm{B}_{\mathrm{}}`$ with $`\beta `$ short: Assume $`\beta =e_1e_2`$. Then $`\mathrm{\Phi }_\beta `$ has type $`\mathrm{B}_\mathrm{}1`$, and so the orbits within $`\mathrm{\Phi }_\beta `$ contribute at most $`\frac{\mathrm{}1}{q^{(\mathrm{}1)/2}}`$. If $`\alpha \mathrm{\Phi }_\beta `$, then $`\alpha =\epsilon e_i\delta e_j`$ where $`i=1`$ or $`2`$ and $`j>2`$. These roots form four orbits of size two with $`m=1`$ and four orbits with $`m=\mathrm{}2`$.
Type $`\mathrm{C}_{\mathrm{}}`$ with $`\beta `$ long: This is similar to type $`\mathrm{B}_{\mathrm{}}`$, with the short roots and long roots exchanged.
Type $`\mathrm{D}_{\mathrm{}}`$: Assume $`\beta =e_1e_2`$. Then $`\mathrm{\Phi }_\beta `$ has type $`\mathrm{D}_\mathrm{}2\mathrm{A}_1`$ and $`\mathrm{\Phi }_\beta ^{}`$ is the subsystem of type $`\mathrm{D}_\mathrm{}2`$. So the orbits within $`\mathrm{\Phi }_\beta ^{}`$ contribute at most $`\frac{\mathrm{}2}{q^{(\mathrm{}1)/2}}`$ to the sum. If $`\alpha \mathrm{\Phi }_\beta ^{}`$, then $`\alpha =\epsilon e_i\delta e_j`$ where $`i=1`$ or $`2`$ and $`j>2`$. These roots form at most four orbits with $`m=\mathrm{}2`$.
The values of the constants in Table 2 are easily computed in Magma. The constant $`c`$ is just $`|C_W(w)|`$. The functions $`Q_w(X)`$ are given in Table 3.
Applying Lemma 6.2, we get $`Q_w(1/q)(11/q)^6`$.
For groups not covered in Table 2, let $`h_\beta `$ be the Coxeter number of $`\mathrm{\Phi }_\beta ^{}`$. Then the centraliser of $`w_\beta `$ in $`W(\mathrm{\Phi }_\beta ^{})`$ has order $`h_\beta `$, and the centraliser of $`w`$ in $`W`$ has order $`2h_\beta 4\mathrm{}`$. The required result now follows from the first inequality of Proposition 6.1. ∎
Finally we are in a position to give an analysis of our algorithm. We refer to Algorithm 10, a version of Algorithm 8 which searches for maximal toral subalgebras corresponding to reflection nonderangements.
As discussed in Subsection 6.2, finding $`f`$ with $`\mathrm{deg}(f)=1`$, or $`\mathrm{deg}(f)=2`$ and $`f=f_{}`$ is equivalent to the corresponding Weyl group element being a reflection nonderangement. When $`\mathrm{deg}(f)=2`$ and $`f=f_{}`$, we have always found in practice that $`M/Z`$ is of type $`\mathrm{A}_1`$, and so has dimension $`3`$. We do not have a proof of this however, so it is necessary to check and then decompose over the field extension $`k_2`$ in the unlikely event that we get a larger subalgebra.
###### Theorem 6.7.
Suppose that the characteristic of $`k`$ is greater than $`3`$. Let $`G`$ be a $`k`$-split connected reductive group and let $`L`$ be the Lie algebra of $`G`$. We can find a split maximal toral subalgebra of $`L`$ in Las Vegas time $`O(n^3\mathrm{}^6\mathrm{log}^2(\mathrm{})\mathrm{log}^2(q))`$.
###### Proof.
Before calling Algorithm 10, we compute the centre of $`L`$, which takes time $`O((n+\mathrm{}^2)^3\mathrm{log}(q))`$. Using Algorithm 7, we can assume $`G`$ is simple. As indicated in Algorithm 10, the computations within the main loop take time $`O(\mathrm{}^7\mathrm{log}(\mathrm{})\mathrm{log}^2(q))`$.
By Proposition 6.6, if $`G`$ is classical with rank at least $`7`$, we obtain a split toral subalgebra $`M`$ with probability at least
$$\left(1\frac{1}{q}\right)^6\left(1\frac{3}{q}\frac{4}{q^2}\frac{\mathrm{}+5}{q^{(\mathrm{}2)/2}}\right)\frac{1}{4\mathrm{}}.$$
For $`q5`$ and $`\mathrm{}7`$, this is at least
$$\left(\frac{4}{5}\right)^6\left(1\frac{3}{5}\frac{4}{25}\frac{12}{5^{5/2}}\right)\frac{1}{4\mathrm{}}>0.$$
Similarly for the Cartan types in Table 2, except for $`\mathrm{D}_4`$,
$$\left(1\underset{i=1}{}\frac{c_i}{q^i}\right)\left(1\frac{1}{q}\right)^6\frac{1}{c}\left(1\underset{i=1}{}\frac{c_i}{5^i}\right)\left(1\frac{1}{5}\right)^6\frac{1}{c}>0.$$
For type $`\mathrm{D}_4`$, the bound is negative for $`q=5`$, but positive for $`q7`$. So it remains to consider the Lie algebra $`\mathrm{D}_4(5)`$. But for any fixed Lie algebra, it is easily seen that there is a nonzero chance of the algorithm working, since there is a chance that the toral subalgebra found by Algorithm 5 is already split. We have now shown that there is a constant $`C>0`$ such that the probability of success after one iteration of the main loop is at least $`C/\mathrm{}`$.
Since
$$\underset{\mathrm{}\mathrm{}}{lim}\left(1\frac{C}{\mathrm{}}\right)^a\mathrm{}=e^{aC},$$
we can choose $`a`$ such that
$$\left(1\frac{C}{\mathrm{}}\right)^a\mathrm{}\frac{1}{e^4}$$
for all $`\mathrm{}`$. Hence the probability of failure after $`a\mathrm{}\mathrm{log}(\mathrm{})`$ repetitions of the loop is at most
$$\left(1\frac{C}{\mathrm{}}\right)^{a\mathrm{}\mathrm{log}(\mathrm{})}\left(\frac{1}{e^4}\right)^{\mathrm{log}(\mathrm{})}=\frac{1}{\mathrm{}^4}.$$
Clearly the depth of recursion is at most $`\mathrm{}`$, which contributes a further factor of $`\mathrm{}`$ to our timing. The ranks of all the subalgebras in all the calls at a particular depth sum to at most $`\mathrm{}`$, so the total number of recursive calls is at most $`\mathrm{}^2`$. Hence the overall probability of success is at least
$$\left(1\frac{1}{\mathrm{}^4}\right)^\mathrm{}^2\left(1\frac{1}{2\mathrm{}^2}\right)^\mathrm{}^2\frac{1}{2}.$$
Hence Algorithm 10 takes Las Vegas time $`O(\mathrm{}^9\mathrm{log}^2(\mathrm{})\mathrm{log}^2(q))`$. Combining this with the preprocessing time of $`O((n+\mathrm{}^2)^3\mathrm{log}(q))`$, and using the fact that $`n\mathrm{}`$ we get the desired result. ∎
###### Corollary 6.8.
Suppose that the characteristic of $`k`$ is greater than $`3`$. Let $`G`$ be a $`k`$-split connected reductive group and let $`L`$ be the Lie algebra of $`G`$. We can find a Chevalley basis of $`L`$ in Las Vegas time $`O(n^3\mathrm{}^6\mathrm{log}^2(\mathrm{})\mathrm{log}^2(q))`$.
###### Proof.
The time taken to find a split maximal toral subalgebra clearly dominates the time for Algorithm 9. ∎
We can easily decompose $`G`$ into simple subgroups, since we know its root datum. Hence, combining this corollary with Proposition 4.3, we see that the algorithm for Lang’s Theorem takes Las Vegas time
$$O(n^3\mathrm{}^6\mathrm{log}^2(\mathrm{})\mathrm{log}^2(q)+n^8r^2s^2\mathrm{log}^2(q)),$$
which is easily simplified to the expression in Theorem 1.2.
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# Using 𝑏-tagging to enhance the SUSY reach of the CERN Large Hadron Collider
## Abstract
Assuming that supersymmetry is realized with parameters in the hyperbolic branch/focus point (HB/FP) region of the minimal supergravity (mSUGRA) model, we show that by searching for multijet + $`E_T^{\mathrm{miss}}`$ events with tagged $`b`$ jets the reach of experiments at the LHC may be extended by as much as 20% from current projections. The reason for this is that gluino decays to third generation quarks are enhanced because the lightest neutralino has substantial higgsino components. Although we were motivated to perform this analysis because the HB/FP region is compatible with the recent determination of the relic density of cold dark matter, our considerations may well have a wider applicability since decays of gluinos to third generation quarks are favoured in a wide variety of models.
preprint: preprint: IFT-P.017/2005, IFUSP-1612/2005, UH-511-1072-05
The WMAP collaboration wmap has determined the cosmological density of cold dark matter (CDM) to be,
$$\mathrm{\Omega }_{CDM}h^2=0.1126_{0.009}^{+0.008}.$$
(1)
This measured value is of the same order of magnitude as the density expected from the production in the Big Bang of a stable, weakly interacting particle with a mass $`𝒪`$(100 GeV), assuming only that it was in thermal equilibrium at some point in the past dodd . Hence, the precise WMAP measurement provides a stringent constraint on all models that include heavy, stable weakly interacting particles. However, since the dark matter may well be made of several components, strictly speaking (1) only implies an upper bound on the density of such particles. In particular this bound applies to the stable lightest supersymmetric particle (LSP), frequently the lightest neutralino $`\stackrel{~}{Z}_1`$, of $`R`$-parity conserving supersymmetric models that have been the focus of much attention during the last twenty-five years rev , and leads us to conclude that,
$$\mathrm{\Omega }_{\stackrel{~}{Z}_1}h^20.129(2\sigma ).$$
(2)
Compatibility with (2) is possible only if the neutralinos can annihilate efficiently which, in turn, is possible only if one of the following holds:
1. The lightest neutralino is hypercharge gaugino-like and annihilates via $`t`$-channel exchanges of relatively light ($`300`$ GeV) sfermions bulk .
2. The neutralino mass $`m_{\stackrel{~}{Z}_1}\frac{1}{2}m_{A,H}`$ so that it annihilates resonantly via the exchange of the neutral Higgs bosons $`A`$ or $`H`$ in the $`s`$-channel funnel . Since these heavier Higgs bosons are typically quite broad, and because the neutralinos have thermal motions, resonant Higgs annihilation occurs over a rather large range of parameters. There is also a small range of parameters where resonant annihilation via the lightest scalar Higgs boson $`h`$ leads to efficient neutralino annihilation lhiggs .
3. The neutralino mass is close to that of a charged or coloured particle; since this latter particle can annihilate efficiently, the neutralino density is correspondingly reduced as long as its interactions maintain it in thermal equilibrium with the co-annihilating charged particle coannih .
4. The parameter $`|\mu |`$ is small compared to the gaugino masses so that the lightest neutralino has significant higgsino components and annihilates effectively via couplings to electroweak gauge bosons matchevfp .
5. The neutralino has significant $`SU(2)`$ gaugino components, and so annihilates to $`W^+W^{}`$ via large isotriplet $`SU(2)`$ couplings. In this case, the lighter chargino tends to be close in mass to $`\stackrel{~}{Z}_1`$, so co-annihilations may also be important hansen .
Most SUSY analyses of the implications of the WMAP measurement have been performed within the framework of the mSUGRA model msugra which, assuming radiative electroweak symmetry breaking, is specified by the well known parameter set,
$$m_0,m_{1/2},A_0,\mathrm{tan}\beta ,\mathrm{sign}(\mu ).$$
Within this framework, which has also been the paradigm for many phenomenological analyses of supersymmetry (SUSY), neutralino annihilation as in item i. occurs in the so-called bulk region with small values of $`m_0`$ and $`m_{1/2}`$; as in ii. in the $`A`$ or $`H`$ funnels which occur only if $`\mathrm{tan}\beta `$ is large; as in iii. only close to the boundary of parameter space where $`\stackrel{~}{\tau }_1`$ becomes the LSP stau , or for special values of $`A_0`$, where $`\stackrel{~}{t}_1`$ becomes the LSP stop ; and as in iv. in the hyperbolic branch/focus point (HB/FP) region with large values of $`m_0`$ and modest to large values of $`m_{1/2}`$ fppapers . The last option v. is not realized within mSUGRA or for that matter in any model with unification of gaugino masses, or in any SUSY Grand Unified Theory (GUT) unless the field that breaks SUSY also breaks the GUT gauge symmetry anderson . In the recently studied non-universal Higgs masses (NUHM) extensions of mSUGRA nuhm , the Higgs funnel occurs for any value of $`\mathrm{tan}\beta `$ while the low $`|\mu |`$ region is accessible even for relatively low values of $`m_0`$ ournuhm .
Within the mSUGRA framework that we adopt for this study, it has been shown that, with an integrated luminosity of 100 fb<sup>-1</sup>, experiments at the Large Hadron Collider (LHC) will probe all of the bulk region, most of the Higgs funnel and, except for the largest values of $`\mathrm{tan}\beta `$, all of the stau co-annihilation region allowed by the WMAP data sasha . The reach in the low $`|\mu |`$ HB/FP region, however, cuts off around $`m_{\stackrel{~}{g}}1.61.8`$ TeV, where the signal from gluino pair production becomes rate limited. Although $`\stackrel{~}{W}_1`$ and $`\stackrel{~}{Z}_2`$ are relatively light and will be abundantly produced at the LHC, if $`|\mu |`$ is indeed small the efficiency for the well studied trilepton signal from $`\stackrel{~}{W}_1\stackrel{~}{Z}_2`$ production trilep is reduced, especially if we are deep in the HB/FP region where the mass gap between $`\stackrel{~}{W}_1`$ (or $`\stackrel{~}{Z}_2`$) and $`\stackrel{~}{Z}_1`$ becomes rather small and the daughter leptons as well as $`E_T^{\mathrm{miss}}`$ from charginos and neutralinos are soft. In this case, it has been shown tadas that by implementing specially designed cuts to separate the chargino pair production soft decay products from Standard Model (SM) background, experiments at an $`e^+e^{}`$ linear collider operating at $`\sqrt{s}=5001000`$ GeV will be able to probe portions of the HB/FP region not accessible at the LHC. Since the LHC is scheduled to commence operations in 2007, while a linear collider is even very optimistically at least a decade away, it is clearly worthwhile to explore all strategies that can potentially expand the LHC reach, especially in this low $`|\mu |`$ region favoured by the WMAP measurements.
An obvious option would be to re-examine the trilepton signal to see whether it is possible to separate it from SM background processes.<sup>1</sup><sup>1</sup>1While we were preparing this manuscript, we learnt that the reach via the trilepton channel has recently been re-examined in Ref. newbaer , and found to be smaller or comparable to the reach in other leptonic and the $`E_T^{\mathrm{miss}}`$ channels even when $`|\mu |<M_{1,2}`$. In this paper, however, we follow a completely different strategy and focus on the signal from gluino pair production (since squarks are very heavy). Our starting point is the observation that since the lighter higgsino-like neutralinos and charginos couple much more strongly to the third generation than to the first two generations, decays of gluino into third generation fermions will be strongly enhanced so that the signal may be expected to be rich in high $`E_T`$ $`b`$-jets dp . In contrast, the dominant SM backgrounds to multijet + $`E_T^{\mathrm{miss}}`$ channels which give the largest SUSY reach at the LHC come from $`t\overline{t}`$ production, from $`V+\mathrm{jet}`$ production ($`V=W,Z`$) and from QCD processes. Since the latter two backgrounds are not expected to be especially rich in hard bottom quark jets, and because experiments at the LHC are expected to have good $`b`$ tagging capability, we explore whether requiring the presence of tagged $`b`$-jets in the signal allows us to probe portions of the hyperbolic branch that are inaccessible using the by now standard analyses bcpt ; atlas ; cms of the various multijet + $`E_T^{\mathrm{miss}}`$ channels at the LHC. While $`b`$-tagging has been suggested before to explore the nature of the underlying model btag , to our knowledge it has never been proposed as a tool for SUSY discovery. We remark that although the $`b`$-tagged signal will be rate limited, unlike the trilepton signal from $`\stackrel{~}{W}_1\stackrel{~}{Z}_2`$ production, this signal will be relatively insensitive to how deep we are in the HB/FP region.
We use the program ISAJET 7.69 isajet with a toy calorimeter described in Ref. bcpt for our analysis. Jets are found using a cone algorithm with a cone size $`\mathrm{\Delta }R=\sqrt{\mathrm{\Delta }\eta ^2+\mathrm{\Delta }\varphi ^2}=0.7`$. Clusters with $`E_T>40`$ GeV and $`|\eta (\mathrm{jet})|<`$ 3 are defined to be jets. Muons (electrons) are classified as isolated if they have $`E_T>10`$ GeV (20 GeV) and visible activity in a cone with $`\mathrm{\Delta }R=0.3`$ about the lepton direction smaller than $`E_T<5`$ GeV. We identify a hadronic cluster with $`E_T40`$ GeV and $`|\eta (j)|<1.5`$ as a $`b`$ jet if it also has a $`B`$ hadron, with $`p_T(B)>15`$ GeV and $`|\eta (B)|<3`$, within a cone with $`\mathrm{\Delta }R=0.5`$ of the jet axis. We take the tagging efficiency $`ϵ_b=0.5`$, and assume that gluon and other quark jets can be rejected as $`b`$ jets by a factor $`R_b=150`$ (50) if $`E_T<100`$ GeV ($`E_T>250`$ GeV) and a linear interpolation in between brej . While we make no representation about the tagging efficiency and rejection against light quark and gluon jets that will be finally achieved, especially at high LHC luminosity and in the jetty environment of SUSY events, we felt that an exploratory study of just how much $`b`$-jet tagging helps with the LHC reach would be worthwhile.
Rather than perform time consuming scans over the WMAP favoured HP/FB regions of the $`m_0m_{1/2}`$ planes of the mSUGRA model for several values of $`\mathrm{tan}\beta `$, we have chosen three diverse model lines for our analysis. For each of these model lines, we take $`\mu >0`$ (the sign favoured by the result of experiment E821 at Brookhaven g-2 ) and fix $`A_0=0`$ — our results are largely insensitive to this choice — and
* $`m_{1/2}=0.295m_0477.5`$ GeV with $`\mathrm{tan}\beta =30`$ for Model Line 1,
* $`m_{1/2}=0.295m_0401.25`$ GeV with $`\mathrm{tan}\beta =30`$ for Model Line 2, and
* $`m_{1/2}=\frac{17}{60}m_0390`$ GeV with $`\mathrm{tan}\beta =52`$ for Model Line 3.
For values of $`m_01500`$ GeV, these model lines all lie in the WMAP allowed HB/FP region of the mSUGRA parameter space delineated in Ref. sasha . The first two model lines have an intermediate value of $`\mathrm{tan}\beta `$ with Model Line 1 being deep in the HB/FP region while Model Line 2 closer to the periphery of the corresponding WMAP region. We choose Model Line 3 again deep in the HB/FP region, but with a very large value of $`\mathrm{tan}\beta `$ to examine any effects from a very large bottom quark Yukawa coupling. We take $`m_t=175`$ GeV throughout this analysis.
The branching fractions for the decays of the gluino with a mass $`1650`$ GeV (close to the limit that can be probed at the LHC via the usual multijet \+ multilepton +$`E_T^{\mathrm{miss}}`$ analyses) into third generation fermions are shown in Table LABEL:tab:decays for the three model lines introduced above. The following features of Table LABEL:tab:decays are worth noting.
1. In all three cases, almost 90% of the gluino decays are to the third generation, so that we expect very hard top and bottom quark jets in SUSY events.
2. In Model Lines 1 and 3 that are deep in the HB/FP region, the gluino mainly decays to the higgsino-like lighter chargino and the two lightest neutralinos; in Model Line 2, $`\stackrel{~}{Z}_1`$ is dominantly the hypercharge gaugino, so that gluino decays to $`\stackrel{~}{Z}_2`$ and to $`\stackrel{~}{Z}_3`$ are favoured.
3. The main difference due to the large $`\mathrm{tan}\beta `$ value for Model Line 3 is the increased branching ratio for the decays $`\stackrel{~}{g}b\overline{b}\stackrel{~}{Z}_i`$ relative to $`\stackrel{~}{g}t\overline{t}\stackrel{~}{Z}_i`$. Although we have not separated these out in the Table, we have checked that while the direct decays to bottom comprise just about 10% of all gluino decays to neutralinos for Model Lines 1 and 2, these decays constitute about a third of all gluino to neutralino decays for Model Line 3. This is, of course, due to the increased Yukawa coupling of the bottom quark dreesprl .
4. In each of these cases, we see that decays of the gluino to the wino-like charginos and neutralinos have relatively small branching fractions despite their large $`SU(2)`$ gauge couplings. This is because decays mediated by lighter third generation squarks that have large Yukawa couplings dominate because these are dynamically as well as kinematically favoured.
The main message of this Table is that in the WMAP favoured HB/FP regions of the mSUGRA model, decays to third generation quarks dominate gluino decays. Moreover, although detailed decay patterns depend on both $`\mathrm{tan}\beta `$ and the value of $`\mu /M_1`$ ( i.e. on how deep we are in the HB/FP region), the total branching fraction for decays to third generation quarks is relatively insensitive to these details. Motivated by these observations we begin our examination of the inclusive $`b`$-jet signal for each of the model lines introduced above.
The major SM backgrounds to the multijet plus $`E_T^{\mathrm{miss}}`$ signal, with or without $`b`$-jets, come from $`W\mathrm{or}Z+\mathrm{jet}`$ production, from $`t\overline{t}`$ production and from QCD production of light quarks and gluons. In the last case, the $`E_T^{\mathrm{miss}}`$ arises from neutrinos from $`c`$ and $`b`$ quarks, from showering of $`W,Z`$ bosons and their subsequent decays to neutrinos, and from energy mismeasurement. The backgrounds from these sources are shown in Table 2 for two representative choices of cuts discussed below. Here $`S_T`$ is the transverse sphericity, $`m_{\mathrm{eff}}`$ is the scalar sum of the $`E_T`$ of the four hardest jets and $`E_T^{\mathrm{miss}}`$, $`\mathrm{\Delta }\varphi `$ is the azimuthal angle between $`E_T^{\mathrm{miss}}`$ and the hardest jet, and $`\mathrm{\Delta }\varphi _b`$ is the azimuthal opening angle between the two hardest $`b`$ jets in events with $`N_b2`$. Since we do not require multileptons in our analysis, backgrounds from $`WW,WZ`$ and $`ZZ`$ as well as three vector bosons processes are expected to be much smaller bcpt .
We see that for both sets of cuts the SM background is dominated by QCD. This differs from earlier results where it is argued that the QCD background can be reduced to negligible levels by requiring the signal to be sufficiently hard. Indeed, to track the reason for this difference, we have done a very high statistics analysis of the QCD background that was not possible in Ref. bcpt . Specifically, we took particular care to divide the event generation into a large number (50 bins for QCD, fewer for other backgrounds) of finely spaced hard scattering $`p_T`$ bins, especially for lower values of the hard scattering $`p_T`$ where the event weights are very large, even if the efficiency for passing the cuts is low.<sup>2</sup><sup>2</sup>2If, for any set of cuts, we find no events in our simulation of a particular background, we set this background cross section to a value corresponding to the single event level in the smallest weight bin in our simulation. We have checked that our backgrounds levels from $`W+\mathrm{jets}`$, $`Z+\mathrm{jets}`$ and $`t\overline{t}`$ processes are in agreement with those in Ref. bcpt and attribute the difference in the QCD background to statistics of the background simulation.<sup>3</sup><sup>3</sup>3In passing, we mention that we also found a significant QCD contribution to the background in the multijets + $`1\mathrm{}`$ \+ $`E_T^{\mathrm{miss}}`$ channel. This is in contrast to the result in Ref. tovey which was obtained using PYTHIA. We attribute this difference to the fact that PYTHIA does not include showering of $`W`$ and $`Z`$ bosons in QCD events. While this may result in some double counting of “Drell-Yan” $`W`$ and $`Z`$ production, we note that showering of vector bosons from the final state quarks will, in general tend to populate a different region of phase space. Since QCD typically contributes about half the background in Table 2, we may expect a small degradation of the LHC reach from these earlier projections.<sup>4</sup><sup>4</sup>4A more significant background issue may be that the showering algorithms appear to obtain a significantly softer distribution of the variable $`m_{\mathrm{eff}}`$ (introduced below) relative to evaluations using exact multijet plus $`W,Z`$ production matrix elements mangano . We have nothing to say about this, except that requiring additional $`b`$’s (and in the case of $`W`$, also requiring a transverse mass cut) will reduce this background considerably. At the very least, our analysis will indicate the extent to which the presence of tagged $`b`$’s increases the LHC reach. We may add that the $`E_T^c`$ analysis of Ref.bcpt that requires $`E_T^{\mathrm{miss}}`$ together with just two (rather than four) additional hard jets, with just one of these jets from showering, may be more robust to these matrix element corrections.
Also shown in Table 2 is the SUSY signal for three points along model lines. In our analysis, we consider a signal to be observable with a given integrated luminosity if, (a) its statistical significance $`N_S/\sqrt{N_B}5`$, (b$`N_S/N_B0.25`$, and (c$`N_S10`$. We see that Case 1 with $`m_{\stackrel{~}{g}}1`$ TeV which by early analyses should be easily observable at the LHC is also observable using cuts 1 with an integrated luminosity of just 10 fb<sup>-1</sup>, and moreover, the significance of the signal improves with increasing $`b`$ multiplicity. For Case 2, the SM background with the softer cuts 1 is too large except in the $`2b`$ channel, where the signal is observable for an integrated luminosity $`26`$ fb<sup>-1</sup>. However with the harder cuts of set 2 the signal, though unobservable without $`b`$-tagging, should be observable in the $`1b`$ ($`2b`$) channel for an integrated luminosity exceeding 26 fb<sup>-1</sup> (41 fb<sup>-1</sup>). For Case 3, both $`b`$-tagging capability and an integrated luminosity of at least 72 fb<sup>-1</sup> are essential for the observability of the signal. It is clear that $`b`$-tagging improves the reach of the LHC for points in the HB/FP region.
To quantify the improvement $`b`$-tagging makes to the capabilities of the LHC for the detection of SUSY, we have re-computed the reach of the LHC for each of the three model lines introduced above. Towards this end, for every mSUGRA parameter point that we examined, we generated a set of SUSY events using ISAJET. We also generated large samples of SM background events. We passed these events through the toy calorimeter simulation mentioned previously, and then analysed both the signal and the background for the entire set of cuts ($`5\times 5\times 6\times 5\times 3^4=60750`$ choices in all) listed in Table 3. We regard the signal as observable if it satisfies the observability criteria (a)–(c) listed above for at least one choice of cuts.
Notice that the cuts in this Table are much harder than those used in a recent analysis of the LHC SUSY signal sasha . This is because, unlike in Ref. sasha where the cuts were designed to extract the signal for a wide range of squark and gluino masses, here we focus on the optimization of the signal in the portion of the HB/FP region with heavy gluinos where the previous strategy fails. Note also that cut 1 in Table 2 corresponds to the softest of these cuts. Since, as we saw earlier, the signal for Case 1 (with $`m_{\stackrel{~}{g}}1`$ TeV) is comfortably observable with the present set of cuts as well as those in Ref. sasha , there is no danger that there will be a gap in the HB/FP region of parameter space where the signal is unobservable with either strategy.
Our results for the LHC reach are shown in Fig. 1, where we plot the largest statistical significance of the signal as we run over the cuts in Table 3 for a) Model Line 1, b) Model Line 2, and c) Model Line 3. The solid curves, from lowest to highest, denote this maximum statistical significance without any $`b`$ tagging requirement, requiring $`1`$ tagged $`b`$-jet, and $`2`$ tagged $`b`$-jets respectively, for an integrated luminosity of 100 fb<sup>-1</sup>, while the dotted curves show the corresponding results for 300 fb<sup>-1</sup> of integrated luminosity that may be expected from three years of LHC operation with the high design luminosity. While, without $`b`$-tagging, our reach for gluinos is $`200`$ GeV smaller than earlier projections sasha ; atlas ; cms , presumably because of differences in the background levels discussed above, we see that $`b`$-tagging will improve the mass reach of gluinos by 15-20%, provided that LHC experiments can accumulate an integrated luminosity of 100-300 fb<sup>-1</sup> and that $`b`$-tagging with an efficiency of $`50`$% remains possible even in the high luminosity environment.
A few remarks appear to be in order at this point:
* The statistical significance in Fig. 1 is not significantly improved if the $`b`$-tagging efficiency improves to 60%. The reason is that before tagging the signal typically contains (on average) 3-4 $`b`$ quark jets while the background typically contains (at most) just two $`b`$ quark jets. As a result, the increased efficiency enhances the $`b`$-tagged background more than the signal, and the statistical significance is essentially unchanged.
* We have also checked that with a $`b`$-tagging efficiency of 50% and an integrated luminosity of 100 fb<sup>-1</sup>, the signal with $`3`$ tagged $`b`$-jets is rate limited and no increase in the reach is obtained from that shown in the Figure.<sup>5</sup><sup>5</sup>5It is possible that a slightly increased reach may be obtained if the tagging efficiency is significantly larger than 50% or if the integrated luminosity is considerably higher than 100 fb<sup>-1</sup>. It would, however, be necessary to evaluate backgrounds from $`4b`$, $`4t`$ and $`t\overline{t}b\overline{b}`$ production processes before a definitive conclusion can be made.
* The search strategy proposed here does not use lepton information at all. We have checked that a transverse mass cut $`M_T(\mathrm{},E_T^{\mathrm{miss}})100`$ GeV on events with at least one isolated lepton does not increase the significance of the signal because the fraction of signal events (after our cuts) with an isolated lepton is not especially large.
* Although we have not shown this explicitly, we have checked that requiring the presence of additional isolated leptons does not lead to an increase in the reach relative to the $`2b`$ channel.
In summary, we have shown that in the HB/FP region of the mSUGRA model the reach of the LHC as measured in terms of gluino masses may be increased by 15-20% by requiring the presence of hard, tagged $`b`$-jets in SUSY events. While we were mainly motivated in our investigation by the fact that this part of parameter space is one of the regions compatible with the WMAP data, our considerations may have wider applicability since decays of heavy gluinos to third generation fermions are favoured in all models with common masses for sfermions with the same gauge quantum numbers. This is in part because the large top Yukawa coupling and, if $`\mathrm{tan}\beta `$ is large, also the bottom quark Yukawa coupling, cause the third generation squarks to be lighter than their siblings in the first two generations, and in part because of new contributions to gluino decay amplitudes from these large Yukawa couplings dreesprl . Specifically, we may expect that signals with tagged $`b`$-jets may also be useful in models with an inverted squark mass hierarchy imh , in models with unification of Yukawa couplings (because they require large $`\mathrm{tan}\beta `$), and possibly also in models with non-universal Higgs mass parameters that have recently been re-examined in light of the WMAP data ournuhm . Since the CMS and ATLAS experiments are expected to ultimately have good $`b`$-tagging capability, we urge that it be utilised to maximize the SUSY reach of the LHC.
###### Acknowledgements.
We thank F. Gianotto and F. Paige for discussions about $`b`$ tagging at the LHC, and H. Baer and M. Drees for comments on the manuscript. This research was supported in part by the U. S. Department of Energy under contract number DE-FG-03-94ER40833 and by Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP).
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# Remarks on quasi-isometric non-embeddability into uniformly convex Banach spaces†
## 1. Construction of the graph $`\mathrm{\Gamma }`$
The main idea is to build a graph with a fractal-like structure in such way that the deformed geometry is reflected on the large scale. We will first inductively construct a sequence $`\left\{\mathrm{\Gamma }\right\}_n`$ of finite graphs which will be the building blocks of $`\mathrm{\Gamma }`$ (compare \[NR\]).
Define $`\mathrm{\Gamma }_0`$ to be a single edge of length 1. To construct $`\mathrm{\Gamma }_1`$ take four copies of $`\mathrm{\Gamma }_0`$ and denote by $`p_i`$ and $`q_i`$ the vertices of $`i`$-th copy, $`i=0,1,2,3`$. The graph $`\mathrm{\Gamma }_1`$ is constructed by identifying $`q_i`$ with $`p_{i+1}`$ with $`imod4`$. Equip $`\mathrm{\Gamma }_1`$ with the path metric.
Similarly to construct $`\mathrm{\Gamma }_{n+1}`$ take four copies of $`\mathrm{\Gamma }_n`$, denote two vertices of valence $`2^n`$ that are distance $`2^n`$ away from each other in the $`i`$-th copy by $`p_i`$ and $`q_i`$, $`i=0,1,2,3`$ (there are two pairs of such vertices in $`\mathrm{\Gamma }_n`$), and identify $`q_i`$ with $`p_{i+1}`$ with $`imod4`$. Equip each $`\mathrm{\Gamma }_n`$ with a path metric. For each $`\mathrm{\Gamma }_n`$ a pair of vertices of valence $`2^n`$ and distance $`2^n`$ from each other will be called *a pair of primary vertices*. In every $`\mathrm{\Gamma }_n`$ there are exactly two pairs of primary vertices. Note that $`\mathrm{diam}\mathrm{\Gamma }_n=2^n`$.
For each $`n`$ denote by $`p_n`$ and $`q_n`$ a pair of primary vertices in $`\mathrm{\Gamma }_n`$. Construct the graph $`\mathrm{\Gamma }`$ by identifying $`q_n`$ with $`p_{n+1}`$ for every $`n`$ (see Fig. 2) and extending the metric in the obvious way. Note that one of the pairs of primary vertices in $`\mathrm{\Gamma }_n`$ changes valence under the isometric embeddings $`\mathrm{\Gamma }_n\mathrm{\Gamma }_{n+1}`$ and $`\mathrm{\Gamma }_n\mathrm{\Gamma }`$, however we will refer to them without change as primary vertices of an isometric copy of $`\mathrm{\Gamma }_n`$.
Actually, for our purposes it would be enough to consider just the set of all vertices of $`\mathrm{\Gamma }`$, however we found the quote from the Introduction to \[Gr\]: *‘Given a discrete metric space $`\mathrm{\Gamma }`$, one can make it more palatable by adding some meat to $`\mathrm{\Gamma }`$ in the form of edges and higher dimensional simplices with vertices in $`\mathrm{\Gamma }`$, without changing the quasi-isometry type’*, quite applicable in this situation.
Recall that a metric space $`X`$ is called locally finite if there exists a discrete subset $`𝒩X`$ and a constant $`C>0`$ such that for every $`xX`$ there is a $`y𝒩`$ satisfying $`d(x,y)C`$ and for every $`y𝒩`$ the number of elements in every ball around $`y`$ in $`𝒩`$ is finite. The graph $`\mathrm{\Gamma }`$ is a locally finite metric space.
## 2. Round ball spaces and nonexistence of the embedding
In this section we prove that the graph $`\mathrm{\Gamma }`$ constructed above does not admit a quasi-isometric embedding into any *round ball metric space*, which we define below.
###### Definition 1.
Let $`X,Y`$ be metric spaces. We say that $`f:XY`$ is a *quasi-isometry* if there are constants $`L>0`$, $`C0`$ such that
$$L^1d_X(x,y)Cd_Y(f(x),f(y))Ld_X(x,y)+C$$
for all $`x,yX`$.
###### Definition 2 (\[La\]).
A metric space $`X`$ is called a *round ball space* if for every $`\epsilon >0`$ there exists $`\delta _\epsilon >0`$ such that
$$\mathrm{diam}\left(B(x,\frac{1+\delta _\epsilon }{2}d(x,y))B(y,\frac{1+\delta _\epsilon }{2}d(x,y))\right)<\epsilon d(x,y)$$
for all $`x,yX`$.
The round ball condition generalizes to metric spaces the classical notion of uniform convexity for Banach spaces. This is made precise in the following
###### Proposition 3 (\[La\]).
A Banach space is a round ball space if and only if it is uniformly convex.
Given a map $`f:XY`$ between metric spaces and points $`x,yX`$, by $`L_{x,y}`$ we denote the Lipschitz constant of $`f|_{\{x,y\}}`$, i.e.
$$L_{x,y}=\frac{d_Y(f(x),f(y))}{d_X(x,y)}.$$
The following lemma can be extracted from \[La\].
###### Lemma 4.
Let $`x_1,x_2,x_3,x_4X`$ satisfy
1. $`d_X(x_1,x_3)=d_X(x_2,x_4)=C`$
2. $`d_X(x_1,x_2)=d_X(x_1,x_3)=d_X(x_2,x_4)=d_X(x_3,x_4)=C/2`$
for some constant $`C>0`$. If $`f:XY`$ is a map to a round ball metric space $`Y`$ and the restriction $`f|_{\{x_1,x_2,x_3,x_4\}}`$ is biLipschitz with constant $`L`$ then
$$\mathrm{max}\{L_{x_1,x_2},L_{x_1,x_4},L_{x_3,x_2},L_{x_3,x_4}\}(1+\delta _{L^2})L_{x_1,x_3}.$$
###### Proof.
Condition (1) implies that
$$d_Y(f(x_2),f(x_4))L^2d_Y(f(x_1),f(x_3)).$$
If $`B(f(x_1),R)B(f(x_3,R)`$ contains $`f(x_2)`$ and $`f(x_4)`$ then we must have $`R\frac{1}{2}\left(1+\delta _{L^2}\right)d_Y(f(x_1),f(x_3))`$, by the round ball condition. Since this is the case when we take
$$R=\mathrm{max}\{d_Y(f(x_1),f(x_2)),d_Y(f(x_1),f(x_4)),d_Y(f(x_3),f(x_2)),d_Y(f(x_3),f(x_4))\},$$
the assertion follows. ∎
###### Theorem 5.
The space $`\mathrm{\Gamma }`$ does not admit a quasi-isometric embedding into any round ball space.
###### Proof.
Assume that there exists a quasi-isometric embedding of the metric space $`\mathrm{\Gamma }`$ into a round ball space $`Y`$. Observe that such an embedding is *biLipschitz for large distances*, i.e. there are constants $`L>0`$ and $`S>0`$ such that
$$L^1d_\mathrm{\Gamma }(x,y)d_Y(f(x),f(y))Ld_\mathrm{\Gamma }(x,y)$$
for all $`x,y\mathrm{\Gamma }`$ satisfying $`d_\mathrm{\Gamma }(x,y)S`$. Choose $`n`$ such that $`\left(1+\delta _{L^2}\right)^nL^1>L`$. In $`\mathrm{\Gamma }`$ choose points $`x_1^{(0)}`$ and $`x_3^{(0)}`$ that are images of primary vertices in an isometric copy of $`\mathrm{\Gamma }_k`$ for some $`k`$ and satisfy $`d_\mathrm{\Gamma }(x_1^{(0)},x_3^{(0)})>2^{n+1}S`$. Denote by $`x_2^{(0)}`$ and $`x_4^{(0)}`$ the second pair of primary vertices in this copy of $`\mathrm{\Gamma }_k`$. The quadruple of points $`x_1^{(0)}`$, $`x_2^{(0)}`$, $`x_3^{(0)}`$, $`x_4^{(0)}`$ satisfies the hypothesis of lemma 4 and we get
(1)
$$\mathrm{max}\{L_{x_1^{(0)},x_2^{(0)}},L_{x_1^{(0)},x_4^{(0)}},L_{x_3^{(0)},x_2^{(0)}},L_{x_3^{(0)},x_4^{(0)}}\}>(1+\delta _{L^2})L_{x_1^{(0)},x_3^{(0)}}$$
Rename the pair for which the $`\mathrm{max}`$ on the left side is attained to $`x_1^{(1)}`$ and $`x_3^{(1)}`$. These two points constitute a pair of primary vertices of an isometric copy of $`\mathrm{\Gamma }_{k1}`$. Denote by $`x_2^{(1)}`$ and $`x_4^{(1)}`$ the remaining pair of primary vertices in $`\mathrm{\Gamma }_{k1}`$ . Applying lemma 4 to the quadruple $`x_1^{(1)}`$, $`x_2^{(1)}`$, $`x_3^{(1)}`$, $`x_4^{(1)}`$ and the renaming procedure together with equation (1) we get a pair of points $`x_1^{(2)}`$ and $`x_3^{(2)}`$ satisfying
$$L_{x_1^{(2)},x_3^{(2)}}(1+\delta _{L^2})^2L_{x_1^{(0)},x_3^{(0)}}.$$
The points $`x_1^{(2)}`$ and $`x_3^{(2)}`$ are again a pair of primary vertices of an isometric copy of $`\mathrm{\Gamma }_{k2}`$.
Continuing in this way after $`n`$ steps we will get a pair of points $`x_1^{(n)}`$, $`x_3^{(n)}`$ satisfying $`d_\mathrm{\Gamma }(x_1^{(n)},x_3^{(n)})S`$ and
$$L_{x_1^{(n)},x_3^{(n)}}(1+\delta _{L^2})^nL_{x_1^{(0)},x_3^{(0)}}.$$
By the choice of $`n`$ we get a contradiction to the Lipschitz condition for large distances. ∎
We remark that what is essential for the proof of Theorem 5 is the existence of an isometric copy of $`\mathrm{\Gamma }_n`$ in $`\mathrm{\Gamma }`$ for arbitrarily large $`n`$. Given the sequence $`\left\{\mathrm{\Gamma }_n\right\}_n`$ one can thus create similar examples using constructions like e.g. disjoint union with an appropriate metric.
*Bounded geometry example.* We also want to indicate that it is now easy to modify the construction of $`\mathrm{\Gamma }`$ to get a bounded geometry metric space with the property of non-embeddability. Simply consider for each $`n`$ the space $`V_n`$, the set of vertices of $`\mathrm{\Gamma }_n`$ with the metric multiplied by $`n`$ and glue them together as in the construction of $`\mathrm{\Gamma }`$. The resulting space however is not quasi-geodesic.
## 3. The geometry of $`c_0`$ and quasi-isometric embeddings
The fact that a 1-net in $`c_0`$ does not admit a biLipschitz for large distances embedding into any uniformly convex Banach space is well known in Banach space geometry, one can prove it using e.g. ultrapowers. We will show however how the methods from the previous section can be implemented in the $`\mathrm{}_{\mathrm{}}^n`$’s and $`c_0`$ and we will give an explicit geometric obstruction to uniform quasi-isometric embeddability (i.e. with constants $`L`$ and $`C`$ independent of $`n`$)<sup>1</sup><sup>1</sup>1We do not use here the standard notion of uniform containment of a sequence $`X_n`$ of finite dimensional spaces in a Banach space, namely that for every $`ϵ>0`$ there is a isomorphic embedding of $`X_n`$ with distortion less than $`1+ϵ`$ for every $`n`$, since it emphasizes the infinitesimal aspect of uniformity. Our definition is suitable for the purposes of large scale geometric behavior. of $`\mathrm{}_{\mathrm{}}^n`$’s into uniformly convex Banach spaces.
Consider the embedding of $`\mathrm{}_{\mathrm{}}^n`$ into $`\mathrm{}_{\mathrm{}}^{n+1}`$ given by adding $`(n+1)`$-st coordinate 0 to each vector in $`\mathrm{}_{\mathrm{}}^n`$. For vectors $`v,w\mathrm{}_{\mathrm{}}^n`$ take their images under this embedding $`\stackrel{~}{v},\stackrel{~}{w}\mathrm{}_{\mathrm{}}^{n+1}`$ and the vectors $`x=\frac{\stackrel{~}{v}+\stackrel{~}{w}}{2}+(\underset{n}{\underset{}{0,0,\mathrm{},0}},\frac{\stackrel{~}{v}\stackrel{~}{w}}{2})`$, $`y=\frac{\stackrel{~}{v}+\stackrel{~}{w}}{2}+(\underset{n}{\underset{}{0,0,\mathrm{},0}},\frac{\stackrel{~}{v}\stackrel{~}{w}}{2})`$. It is easy to check that the quadruple of points $`\stackrel{~}{v},\stackrel{~}{w},x,y`$ satisfies the hypothesis of lemma 4, thus we can apply with no change the procedure from the proof of Theorem 5 and recover
###### Proposition 6.
Let $`X`$ be a Banach space containing $`\mathrm{}_{\mathrm{}}^n`$’s quasi-isometrically uniformly. Then $`X`$ does not admit a quasi-isometric embedding into any uniformly convex Banach space.
The same argument gives a purely metric proof of the fact that the unit ball in $`c_0`$ does not admit a biLipschitz embedding into any uniformly convex Banach space, which is again obvious once we appeal to the linear structure of $`c_0`$.
## 4. Some remarks on a paper of Bourgain
For the purposes of the coarse Novikov Conjecture \[KY\] one can consider coarse embeddings into superreflexive Banach spaces, since these are exactly the ones that admit an equivalent uniformly convex norm, due to a theorem of Enflo \[En\]. It might be thus interesting to confront the above observations with a paper of J. Bourgain in which a metric characterization of superreflexivity was given \[Bo\]. The necessary and sufficient condition for a Banach space to be superreflexive was shown to be biLipschitz uniform non-embeddability of a sequence of trees $`T_j`$ defined below (see \[Bo\] for a precise formulation).
Denote $`\mathrm{\Omega }_n=\{1,1\}^n`$, $`T_n=_{in}\mathrm{\Omega }_i`$, $`T=_{j=1}^{\mathrm{}}T_j`$ and again add ”some meat” to $`T`$ in the form of edges in the obvious way and denote the resulting space by $`\overline{T}`$. First note that quasi-isometric embeddability of $`T`$ (or equivalently $`\overline{T}`$) into a superreflexive Banach space implies the existence of a biLipschitz embedding of $`T`$ into such a space.
###### Lemma 7.
Let $`X`$ be a discrete metric space and assume that $`X`$ embeds quasi-isometrically into a superreflexive Banach space. Then $`X`$ admits a biLipschitz embedding into a superreflexive Banach space.
###### Proof.
Let $`f:XE`$ be the quasi-isometric embedding. Define<sup>2</sup><sup>2</sup>2Although it is not of great importance, it is convenient to take the direct sum with the $`\mathrm{}_1`$-norm $`\stackrel{~}{f}:XE\mathrm{}_2(X)`$ by the formula $`\stackrel{~}{f}(x)=f(x)\delta _x`$. $`E\mathrm{}_2(X)`$ is super-reflexive and it is easy to verify that $`\stackrel{~}{f}`$ is a biLipschitz embedding. ∎
Thus $`\overline{T}`$ also does not admit a quasi-isometric embedding into any superreflexive Banach space, the argument in \[Bo\] is however of probabilistic nature, as mentioned earlier.
The geometry of the graph $`\mathrm{\Gamma }`$ is intuitively the very opposite of the hyperbolic geometry on a tree. This is indicated already by topological invariants, but also the geometries of these spaces are very different on the large scale. The next proposition shows that Theorem 5 and Bourgain’s result cannot be deduced from each other.
###### Proposition 8.
1. $`\overline{T}`$ does not admit a coarse embedding into $`\mathrm{\Gamma }`$
2. $`\mathrm{\Gamma }`$ does not admit a coarse embedding into $`\overline{T}`$.
###### Proof.
To see (1) note the obvious fact that since the graph is infinite in just one direction, a coarse embedding from $``$ to $`\mathrm{\Gamma }`$ must map both infinite ends of $``$ in the same direction in $`\mathrm{\Gamma }`$ and distant points on the line have to cross close to the joints of $`\mathrm{\Gamma }`$, i.e. the points in which $`\mathrm{\Gamma }_n`$ is glued with $`\mathrm{\Gamma }_{n+1}`$. Since $``$ is isometrically embedded in the tree $`\overline{T}`$, the assertion follows.
Similarly for (2) take two infinite geodesic rays in $`\mathrm{\Gamma }`$ and observe that they have to pass through the joints so that arbitrarily far on these geodesic rays some points are identified. In an isometric copy of $`\mathrm{\Gamma }_n\mathrm{\Gamma }`$ the pair of primary vertices that are not joints is at distance $`2^n`$ apart. Thus a coarse embedding of $`\mathrm{\Gamma }`$ into $`\overline{T}`$ from some point on cannot map these vertices into the same branch in $`\overline{T}`$. But that means that the joints as points on two geodesic rays will be mapped to points whose distance grows to infinity, which gives a contradiction. ∎
Note that to get quasi-isometric non-embeddability of $`\mathrm{\Gamma }`$ into $`\overline{T}`$ it is enough to use Theorem 5, since any $``$-tree is a round ball space.
The intuition behind the facts presented above might be the following: the tree $`T`$ carries an $`\mathrm{}_1`$-geometry, while the geometry of $`\mathrm{\Gamma }`$ resembles the one of $`c_0`$. Recall that $`\mathrm{}_1`$ and $`c_0`$ are the standard examples of non-uniformly convex Banach spaces, thus the presence of any of these geometries should be an obstruction to quasi-isometric embeddability into uniformly convex Banach spaces.
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# Singularity dynamics: Action and Reaction
## 1 Introduction
The method presented builds on the action principle which offers a general tool to describe both particle and field dynamics. This principle represents one of the great unification in physics and applies equally well to Newtonian mechanics , general relativity and electromagnetic fields. Using this single principle it is possible to deduce the Euler-Lagrange equations describing the evolution of one system or of many interacting systems. However describing these interactions requires the assumption of an additional action term. In this paper, we show that it is possible to eliminate this additional assumption and deduce the reaction dynamics of the singularity. In the case of scalar field singularities we derive Newton’s equation of motion together with the equivalence between Newtonian gravitational mass and inertial mass. For charges and magnetic monopoles we deduce their associated Lorentz force.
Singularities, vortices and topological defects have received considerable attention over the years; they determine the properties of many interesting materials such as vortices in superconductors , superfluids and two-dimensional condensate , dislocations and defects in solid and liquid crystals , optical phase singularities and optical vortices and even topological defect in cosmology . In this paper, we use singularities together with the field action principle to deduce their dynamical behaviour. The method presented here is general and can be applied to all of the fields mentioned above.
For clarity reasons we treat first the simplest case possible that being a Dirac source singularity in the Lorentz-invariant scalar wave equation. This leads to equations of motion similar to Newtonian mechanics (section 2). The method is generalised for other singularities and other fields such as a charges (section 3) where the electromagnetic field implies the Lorentz force. In the last section, we apply this principle to deduce the equations of motion of hypothetical magnetic monopoles where the dynamics concurs with the one obtained through charge-monopole symmetry considerations .
## 2 Scalar field singularity: Inertia
Before considering the anatomy of the field singularity, let us start by deducing the Euler-Lagrange equation for a scalar mass-less field $`U`$ described by the following action
$$𝒮_U=(U)d^3x𝑑t=\frac{1}{2}\frac{1}{c^2}(_tU)^2(\mathbf{}U)^2d^3xdt$$
(1)
where $`(U)`$ corresponds to the Lagrangian density, $`\mathbf{}=(_x,_y,_z)`$ to the nabla operator and $`c`$ to the speed of light. The action principle implies that the scalar field $`U`$ is solution of the standard homogenous wave equation
$$\mathbf{}^2U\frac{1}{c^2}_t^2U=0$$
(2)
where $`\mathbf{}^2`$ stands for the Laplacian operator. The next step in the singularity dynamics theory is to define a singular source term that maintains the Lorentz-invariance of the supporting field equation. Starting from the wave equation (2), we can define the stationary singular field $`U_s`$ as the time independent solution in the presence of a Dirac source at the origin of the coordinate system
$$\mathbf{}^2U_s\frac{1}{c^2}_t^2U_s=m_s\delta (𝐫)$$
(3)
where $`m_s`$ corresponds for the source amplitude. To define the Lorentz-invariant singularity, we consider equation (3) in a reference frame having the constant velocity $`𝐯`$ with respect to the stationary singularity. In this reference frame the singularity has the velocity $`𝐯`$. The transformation between the two reference frames is :
$`𝐫_{}^{}{}_{}{}^{}`$ $`=`$ $`\gamma (𝐫_{}+𝐯t)`$
$`𝐫_{}^{}{}_{}{}^{}`$ $`=`$ $`𝐫_{}`$ (4)
$`t^{}`$ $`=`$ $`\gamma \left(t+{\displaystyle \frac{𝐯𝐫}{c^2}}\right).`$
where $`\gamma =(1𝐯^2/c^2)^{1/2}`$ and the prime stands for the coordinates in the moving reference frame. The subscripts $``$ and $``$ represent the vectorial components parallel and perpendicular to the velocity $`𝐯`$. Applying these transformations to equation (3) leads us to a constant velocity moving singularity. Beforehand, we remark that the singularity can be decomposed into two components, $`\delta (𝐫)=\delta (𝐫_{})\delta (𝐫_{})`$. Further, the wave equation and the scalar field $`U_s`$ is invariant with respect to the transformation that is it remains unchanged in the transition from one reference frame to another. It is only the singularity which changes in the moving reference system
$`\mathbf{}_{}^{\mathbf{}}{}_{}{}^{2}U_s{\displaystyle \frac{1}{c^2}}_t^{}^2U_s`$ $`=`$ $`m_s\delta (\gamma (𝐫_{}^{}{}_{}{}^{}𝐯t^{}))\delta (𝐫_{}^{}{}_{}{}^{})`$
$`=`$ $`m_s\sqrt{1{\displaystyle \frac{𝐯^\mathrm{𝟐}}{c^2}}}\delta (𝐫_{}^{}{}_{}{}^{}𝐯t^{}))\delta (𝐫_{}^{}{}_{}{}^{})`$
$`=`$ $`m_s\sqrt{1{\displaystyle \frac{𝐯^\mathrm{𝟐}}{c^2}}}\delta (𝐫^{}𝐫_{}^{}{}_{s}{}^{})`$
where $`𝐫_{}^{}{}_{s}{}^{}=𝐯t^{}`$ is the position of the singularity in the moving reference system. We have used the scaling property of the single variable Dirac distribution $`|a|\delta (ax)=\delta (x)`$.
We can now write the equation of a moving singularity on an arbitrary path $`𝐫_s(t)`$ with the velocity $`𝐯=\dot{𝐫}_s`$. The singular field $`U_s`$ is solution of
$$\mathbf{}^2U_s\frac{1}{c^2}_t^2U_s=m_s\sqrt{1𝐯^2/c^2}\delta (𝐫𝐫_s).$$
(5)
and corresponds to a retarded potential.
In the following, we are considering the action principle in the case of the linear superposition of the singular field as defined in equation (5) and a given external field $`U_0`$ solution of (2). We remark here that the only free parameter of the system is the position of the Dirac singularity and its velocity. The Euler-Lagrange equation of motion is deduced when the action is stationary with respect path of the singularity. The total action integral may be expressed as
$`𝒮_{U_0+U_s}(𝐫_s(t))`$ $`=`$ $`{\displaystyle d^3x𝑑t(U_0+U_s)}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3x𝑑t\frac{1}{c^2}(_tU_0)^2}(\mathbf{}U_0)^2+{\displaystyle \frac{1}{c^2}}(_tU_s)^2(\mathbf{}U_s)^2`$
$`+{\displaystyle d^3x𝑑t\frac{1}{c^2}(_tU_0)(_tU_s)}(\mathbf{}U_0)(\mathbf{}U_s).`$
This action integral can be decomposed into two parts. The first integral has no cross-terms between the external field and the singular field. This integral is independent with respect to the position of the singularity and its variation is zero. This leaves only the second integral $`\mathrm{\Delta }𝒮_s`$ for the action principle
$`\mathrm{\Delta }𝒮_s`$ $`=`$ $`{\displaystyle d^3x𝑑tU_0\mathbf{}^2U_s}{\displaystyle \frac{1}{c^2}}U_0_t^2U_s`$
where we have integrated the action by parts and assumed the singular field to be zero on the boundary of the integration region. Next, we use equation (5) to replace the singular field by the Dirac singularity and finally we integrate in Lagrangian density in the spacial domain.
$`\mathrm{\Delta }𝒮_s`$ $`=`$ $`{\displaystyle d^3x𝑑tU_0m_s\sqrt{1𝐯^2/^2}\delta (𝐫𝐫_s)}`$
$`=`$ $`{\displaystyle 𝑑tU_0m_s\sqrt{1𝐯^2/c^2}}`$
The equations of movement can be obtained by using the conventional particle Euler-Lagrange differential equation where the position and velocity correspond to those of the singularity. Indeed we have
$`{\displaystyle \frac{d}{dt}}m_i𝐯`$ $`=`$ $`m_g\mathbf{}U_0`$ (6)
where we have defined an effective inertial $`m_i`$ and gravitational $`m_g`$ mass
$$\begin{array}{cc}\hfill m_i& =\frac{U_0m_s}{c^2\sqrt{1𝐯^2/c^2}}\hfill \\ \hfill m_g& =m_s\sqrt{1𝐯^2/c^2}.\hfill \end{array}$$
(7)
One can also define a differential relationship between the two masses by taking the scalar product between the velocity and equation (6). This leads to a energy balance equation that generalises the famous equations $`E=mc^2`$ to the dynamic domain.
$`{\displaystyle \frac{d}{dt}}m_ic^2`$ $`=`$ $`m_g{\displaystyle \frac{}{t}}U_0`$ (8)
Indeed, the energy balance equation shows that the variation of the inertial mass is proportional to its potential energy variation.
Finally, we can deduce Newtons equation of motion for non-relativistic velocities and far from any other singularities
$`{\displaystyle \frac{d}{dt}}m_s𝐯`$ $`=`$ $`0`$ (9)
where the scalar field gauge is chosen to be $`U_0(\mathrm{})=c^2`$.
## 3 Charge singularity: Lorentz force
In order to deduce the laws of motion of a charge we have to consider a singularity of the electromagnetic field described by Maxwell’s equations. We proceed in a similar way to the scalar singularity case. In a first instance, we deduce the homogenous Maxwell’s equations using the action principle. In a second step, we introduce a singular source to these equations making sure that their symmetries are maintained (i.e. Lorentz-invariance). Then, we write down the Lagrangian integral of the total field composed of an external field and the singular field. Finally, the action principle is applied to the total Lagrangian integral and the dynamics of the singularity is deduced.
The electromagnetic field action is expressed as:
$$𝒮_{em}=\left(\frac{ϵ_0}{2}𝐄𝐄\frac{1}{2\mu _0}𝐁𝐁\right)d^3x𝑑t$$
(10)
where the electric and magnetic fields are defined by
$$\begin{array}{cc}\hfill 𝐄& =\mathbf{}V_t𝐀\hfill \\ \hfill 𝐁& =\mathbf{}\times 𝐀\hfill \\ \hfill 0& =\mathbf{}𝐀+ϵ_0\mu _0_tV\hfill \end{array}$$
(11)
and where $`V`$ and $`𝐀`$ are the electric scalar and magnetic vector potential respectively. The vacuum permittivity and permeability is given by $`ϵ_0`$ and $`\mu _0`$. The third equation corresponds to Lorentz gauge. Applying the action principle with respect to the potentials lets us deduce Maxwell’s equations in absence of charges:
$`\mathbf{}\times 𝐄_t𝐁`$ $`=`$ $`0`$ (12a)
$`{\displaystyle \frac{1}{\mu _0}}\mathbf{}\times 𝐁ϵ_0_t𝐄`$ $`=`$ $`0`$ (12b)
$`\mathbf{}ϵ_0𝐄`$ $`=`$ $`0`$ (12c)
$`\mathbf{}𝐁`$ $`=`$ $`0.`$ (12d)
Like in the scalar case, we consider a single stationary Dirac singularity $`q\delta (𝐫)`$ in equation of the divergence of the electrical field (12c). The effect on the electromagnetic field of this source is equivalent to a stationary point charge $`q`$ at the origin. The singular charge fields $`𝐄_𝐪`$ and $`𝐁_𝐪`$ are solutions of the in-homogenous Maxwell’s equations:
$$\begin{array}{cc}\hfill \mathbf{}\times 𝐄_𝐪_t𝐁_𝐪& =0\hfill \\ \hfill \frac{1}{\mu _0}\mathbf{}\times 𝐁_𝐪ϵ_0_t𝐄_𝐪& =0\hfill \\ \hfill \mathbf{}ϵ_0𝐄_𝐪& =q\delta (𝐫)\hfill \\ \hfill \mathbf{}𝐁_𝐪& =0\hfill \end{array}$$
(13)
The main difference with the scalar case is in that the potential fields $`V_q`$ and $`𝐀_𝐪`$ change in the transition from one reference frame to another. Indeed, using the moving frame defined by (3) we have for the potentials:
$$\begin{array}{cc}\hfill 𝐀_{𝐪}^{}{}_{}{}^{}& =\gamma \left(𝐀_{𝐪}^{}{}_{}{}^{}+\frac{𝐯V_q}{c^2}\right)\hfill \\ \hfill 𝐀_{𝐪}^{}{}_{}{}^{}& =𝐀_{𝐪}^{}{}_{}{}^{}\hfill \\ \hfill V_q^{}& =\gamma (V_q+𝐯𝐀_𝐪)\hfill \end{array}$$
(14)
where the speed of light is given by $`1=c^2ϵ_0\mu _0`$. Combining the transformations (3) and (14) together with the definition of the electric and magnetic fields (11) we can deduce the transformation relationship between these fields
$$\begin{array}{cc}\hfill 𝐄_{𝐪}^{}{}_{}{}^{}& =𝐄_{𝐪}^{}{}_{}{}^{}\hfill \\ \hfill 𝐄_{𝐪}^{}{}_{}{}^{}& =\gamma (𝐄_{𝐪}^{}{}_{}{}^{}𝐯\times 𝐁_𝐪)\hfill \\ \hfill 𝐁_{𝐪}^{}{}_{}{}^{}& =𝐁_{𝐪}^{}{}_{}{}^{}\hfill \\ \hfill 𝐁_{𝐪}^{}{}_{}{}^{}& =\gamma \left(𝐁_{𝐪}^{}{}_{}{}^{}+\frac{𝐯\times 𝐄_𝐪}{c^2}\right).\hfill \end{array}$$
(15)
This implies the following transformation for equations (13):
$`\mathbf{}^{\mathbf{}}\times 𝐄_𝐪^{}_t^{}𝐁_𝐪^{}`$ $`=`$ $`0`$
$`{\displaystyle \frac{1}{\mu _0}}\mathbf{}^{\mathbf{}}\times 𝐁_𝐪^{}ϵ_0_t^{}𝐄_𝐪^{}`$ $`=`$ $`q𝐯\gamma \delta (\gamma (𝐫_{}^{}{}_{}{}^{}𝐯t^{}))\delta (𝐫_{}^{}{}_{}{}^{})`$
$`\mathbf{}^{\mathbf{}}ϵ_0𝐄_𝐪^{}`$ $`=`$ $`q\gamma \delta (\gamma (𝐫_{}^{}{}_{}{}^{}𝐯t^{}))\delta (𝐫_{}^{}{}_{}{}^{})`$
$`\mathbf{}^{\mathbf{}}𝐁_𝐪^{}`$ $`=`$ $`0`$
Applying the scaling property of the single variable Dirac distribution $`|a|\delta (ax)=\delta (x)`$ and introducing the definition of the position of the singularity $`𝐫_{}^{}{}_{q}{}^{}=𝐯t^{}`$ in the moving reference system we get:
$`\mathbf{}^{\mathbf{}}\times 𝐄_𝐪^{}+_t^{}𝐁_𝐪^{}`$ $`=`$ $`0`$
$`{\displaystyle \frac{1}{\mu _0}}\mathbf{}^{\mathbf{}}\times 𝐁_𝐪^{}ϵ_0_t^{}𝐄_𝐪^{}`$ $`=`$ $`q𝐯\delta (𝐫^{}𝐫_{}^{}{}_{q}{}^{})`$
$`\mathbf{}^{\mathbf{}}ϵ_0𝐄_𝐪^{}`$ $`=`$ $`q\delta (𝐫^{}𝐫_{}^{}{}_{q}{}^{})`$
$`\mathbf{}^{\mathbf{}}𝐁_𝐪^{}`$ $`=`$ $`0`$
where the relativistic coefficient $`\gamma `$ in the source terms cancelled out. The reason for this is the covariant behaviour of the electromagnetic potentials as opposed to the invariant scalar potential in equation (5).
This procedure gives us the form of the singularity following an arbitrary path $`𝐫_q(t)`$ with the velocity $`𝐯=\dot{𝐫}_q`$. The resulting vector singularity is equivalent to a charge and its associated current. The singular charge fields $`𝐄_q`$ and $`𝐁_q`$ are solutions of the in-homogenous Maxwell’s equations:
$$\begin{array}{cc}\hfill \mathbf{}\times 𝐄_𝐪_t𝐁_𝐪& =0\hfill \\ \hfill \frac{1}{\mu _0}\mathbf{}\times 𝐁_𝐪ϵ_0_t𝐄_𝐪& =q𝐯\delta (𝐫𝐫_q)\hfill \\ \hfill \mathbf{}ϵ_0𝐄_𝐪& =q\delta (𝐫𝐫_q)\hfill \\ \hfill \mathbf{}𝐁_𝐪& =0.\hfill \end{array}$$
(16)
To deduce the equations of motion of this moving charge singularity we define the total field as the linear superposition of the singular fields $`𝐄_𝐪`$ and $`𝐁_𝐪`$ and the external fields $`𝐄_\mathrm{𝟎}`$ and $`𝐁_\mathrm{𝟎}`$:
$`𝐄`$ $`=`$ $`𝐄_𝐪+𝐄_\mathrm{𝟎}`$
$`𝐁`$ $`=`$ $`𝐁_𝐪+𝐁_\mathrm{𝟎}`$
where the external fields are decomposed onto the external scalar and vector potentials
$`𝐄_\mathrm{𝟎}`$ $`=`$ $`\mathbf{}V_0_t𝐀_\mathrm{𝟎}`$
$`𝐁_\mathrm{𝟎}`$ $`=`$ $`\mathbf{}\times 𝐀_\mathrm{𝟎}`$
$`0`$ $`=`$ $`\mathbf{}𝐀_\mathrm{𝟎}+ϵ_0\mu _0_tV_0.`$
Introducing the total fields $`𝐄`$ and $`𝐁`$ into the electromagnetic Lagrangian (10) infers the following singularity action:
$$\mathrm{\Delta }𝒮_{sq}=\left(ϵ_0𝐄_\mathrm{𝟎}𝐄_𝐪\frac{1}{\mu _0}𝐁_\mathrm{𝟎}𝐁_𝐪\right)d^3x𝑑t$$
(17)
where we have kept only the terms that vary with respect to the trajectory of the singularity. Indeed, the terms in $`𝐄_𝐢𝐄_𝐢`$ and $`𝐁_𝐢𝐁_𝐢`$, where $`(i=q,0)`$, are not dependent on the position of the charge singularity. Rewriting the external field with the help of the scalar $`V_0`$ and vector potential $`𝐀_\mathrm{𝟎}`$ allows us to integrate by parts:
$`\mathrm{\Delta }𝒮_{sq}`$ $`=`$ $`{\displaystyle \left(ϵ_0\left(\mathbf{}V_0_t𝐀_\mathrm{𝟎}\right)𝐄_𝐪\frac{1}{\mu _0}\left(\mathbf{}\times 𝐀_\mathrm{𝟎}\right)𝐁_𝐪\right)d^3x𝑑t}`$
$`=`$ $`{\displaystyle \left(ϵ_0\left(V_0\mathbf{}𝐄_𝐪+𝐀_\mathrm{𝟎}_t𝐄_𝐪\right)\frac{1}{\mu _0}\left(\mathbf{}\times 𝐁_𝐪\right)𝐀_\mathrm{𝟎}\right)d^3x𝑑t}`$
where we have assumed the fields $`𝐄_𝐪`$ and $`𝐁_𝐪`$ to be zero on the boundary of the integration domain. Using equations (16) we can further simplify the action of the singularity
$`\mathrm{\Delta }𝒮_{sq}`$ $`=`$ $`{\displaystyle \left(qV_0\delta (𝐫𝐫_q)q𝐀_\mathrm{𝟎}𝐯\delta (𝐫𝐫_q)\right)d^3x𝑑t}`$
$`=`$ $`{\displaystyle \left(qV_0q𝐀_\mathrm{𝟎}𝐯\right)𝑑t}`$
which leads directly to Lorentz forces for a charge in the external field defined by $`𝐄_\mathrm{𝟎}`$ and $`𝐁_\mathrm{𝟎}`$. To convince ourselves of this we can rigidly link a scalar mass singularity introduced in previous section to a charge singularity. The total action of this double singularity is:
$`\mathrm{\Delta }𝒮_s+\mathrm{\Delta }𝒮_{sq}`$ $`=`$ $`{\displaystyle \left(U_0m_s\sqrt{1\frac{𝐯^2}{c^2}}+qV_0q𝐀_\mathrm{𝟎}𝐯\right)𝑑t}`$ (18)
which lead to the following equation of movement
$`{\displaystyle \frac{d}{dt}}m_i𝐯`$ $`=`$ $`m_g\mathbf{}U_0+q𝐄_\mathrm{𝟎}+q𝐯\times 𝐁_\mathrm{𝟎}`$ (19)
where $`m_i`$ and $`m_g`$ are the defined by (7). The energy balance equation is
$`{\displaystyle \frac{d}{dt}}m_ic^2`$ $`=`$ $`m_g{\displaystyle \frac{}{t}}U_0+q𝐯𝐄_\mathrm{𝟎}.`$ (20)
## 4 Magnetic Monopole singularity
Magnetic monopoles are hypothetical particles that behave in a similarly to charged particles with magnetic and electric fields inverted (). The monopoles correspond to sources of the magnetic field in equation (12d). In the following we are deducing the equations of motion of magnetic source singularities by considering a Dirac source $`q_m\delta (𝐫)`$ in equations (12d) where $`q_m`$ corresponds to the monopole strength.
Before proceeding further, we have to generalise the potential decomposition of the electromagnetic field. Indeed, the potentials that we used for the charged singularities does not allow the presence of sources of the magnetic field. Conventionally the magnetic field is by definition a rotational field. Therefore, we define a symmetric electromagnetic potential decomposition that includes an additional magnetic scalar and electric vector potential.
$$\begin{array}{cc}\hfill 𝐄& =\mathbf{}V_q_t𝐀_𝐪+\frac{1}{ϵ_0}\mathbf{}\times 𝐀_𝐦\hfill \\ \hfill 𝐁& =\mathbf{}\times 𝐀_𝐪+\mu _0\mathbf{}V_m+\mu _0_t𝐀_𝐦\hfill \\ \hfill 0& =\mathbf{}𝐀_𝐪+ϵ_0\mu _0_tV_q\hfill \\ \hfill 0& =\mathbf{}𝐀_𝐦+ϵ_0\mu _0_tV_m\hfill \end{array}$$
(21)
where the subscripts $`(q,m)`$ stand respectively for the charge and monopole potentials or fields. Using all the newly defined potential definitions we can apply again the action principle the electromagnetic action defined by equation (10). The resulting evolution equations are identical to equations (12). The only difference is the possibility to define a monopole source density in equation (22d).
$`\mathbf{}\times 𝐄_t𝐁`$ $`=`$ $`0`$ (22a)
$`{\displaystyle \frac{1}{\mu _0}}\mathbf{}\times 𝐁ϵ_0_t𝐄`$ $`=`$ $`0`$ (22b)
$`\mathbf{}ϵ_0𝐄`$ $`=`$ $`0`$ (22c)
$`\mathbf{}𝐁`$ $`=`$ $`0.`$ (22d)
For the monopole singularity, we consider a single stationary Dirac singularity at the origin. The singularity constitutes the free term of equation (22d) and acts as a source of magnetic field. The singular fields $`𝐄_𝐦`$ and $`𝐁_𝐦`$ are solution of:
$$\begin{array}{cc}\hfill \mathbf{}\times 𝐄_𝐦_t𝐁_𝐦& =0\hfill \\ \hfill \frac{1}{\mu _0}\mathbf{}\times 𝐁_𝐦ϵ_0_t𝐄_𝐦& =0\hfill \\ \hfill \mathbf{}ϵ_0𝐄_𝐦& =0\hfill \\ \hfill \mathbf{}𝐁_𝐦& =q_m\delta (𝐫)\hfill \end{array}$$
(23)
In this case there are four potential fields $`V_{sq}`$, $`V_{sm}`$, $`𝐀_{\mathrm{𝐬𝐪}}`$ and $`𝐀_{\mathrm{𝐬𝐦}}`$ and the electric and magnetic field is defined by equations (21). Taking the same frame transformation defined by (3) we have for the potentials:
$$\begin{array}{cc}\hfill 𝐀_{\mathrm{𝐬𝐪}}^{}{}_{}{}^{}& =\gamma (𝐀_{\mathrm{𝐬𝐪}}^{}{}_{}{}^{}+\frac{𝐯V_{sq}}{c^2})\hfill \\ \hfill 𝐀_{\mathrm{𝐬𝐪}}^{}{}_{}{}^{}& =𝐀_{\mathrm{𝐬𝐪}}^{}{}_{}{}^{}\hfill \\ \hfill V_{sq}^{}& =\gamma (V_{sq}+𝐯𝐀_{\mathrm{𝐬𝐪}})\hfill \\ \hfill 𝐀_{\mathrm{𝐬𝐦}}^{}{}_{}{}^{}& =\gamma (𝐀_{\mathrm{𝐬𝐦}}^{}{}_{}{}^{}+\frac{𝐯V_{sm}}{c^2})\hfill \\ \hfill 𝐀_{\mathrm{𝐬𝐦}}^{}{}_{}{}^{}& =𝐀_{\mathrm{𝐬𝐦}}^{}{}_{}{}^{}\hfill \\ \hfill V_{sm}^{}& =\gamma (V_{sm}+𝐯𝐀_{\mathrm{𝐬𝐦}}).\hfill \end{array}$$
(24)
The implied relationship between the electromagnetic fields in the two frames is identical to the transformation for the charges only field.
$$\begin{array}{cc}\hfill 𝐄_{𝐦}^{}{}_{}{}^{}& =𝐄_{𝐦}^{}{}_{}{}^{}\hfill \\ \hfill 𝐄_{𝐦}^{}{}_{}{}^{}& =\gamma (𝐄_{𝐦}^{}{}_{}{}^{}𝐯\times 𝐁_𝐦)\hfill \\ \hfill 𝐁_{𝐦}^{}{}_{}{}^{}& =𝐁_{𝐦}^{}{}_{}{}^{}\hfill \\ \hfill 𝐁_{𝐦}^{}{}_{}{}^{}& =\gamma (𝐁_{𝐦}^{}{}_{}{}^{}+\frac{𝐯\times 𝐄_𝐦}{c^2})\hfill \end{array}$$
(25)
Combining all the above transformations implies the singularity in equations (23) to transform as:
$`\mathbf{}^{\mathbf{}}\times 𝐄_𝐦^{}_t^{}𝐁_𝐦^{}`$ $`=`$ $`q_m𝐯\gamma \delta (\gamma (𝐫_{}^{}{}_{}{}^{}𝐯t^{}))\delta (𝐫_{}^{}{}_{}{}^{})`$
$`{\displaystyle \frac{1}{\mu _0}}\mathbf{}^{\mathbf{}}\times 𝐁_𝐦^{}ϵ_0_t^{}𝐄_𝐦^{}`$ $`=`$ $`0`$
$`\mathbf{}^{\mathbf{}}ϵ_0𝐄_𝐦^{}`$ $`=`$ $`0`$
$`\mathbf{}^{\mathbf{}}𝐁_𝐦^{}`$ $`=`$ $`q_m\gamma \delta (\gamma (𝐫_{}^{}{}_{}{}^{}𝐯t^{}))\delta (𝐫_{}^{}{}_{}{}^{})`$
Applying the scaling property of the single variable Dirac distribution $`|a|\delta (ax)=\delta (x)`$ and introducing the definition of the position of the singularity $`𝐫_{}^{}{}_{m}{}^{}=𝐯t^{}`$ in the moving reference system we get:
$`\mathbf{}^{\mathbf{}}\times 𝐄_𝐦^{}_t^{}𝐁_𝐦^{}`$ $`=`$ $`q_m𝐯\delta (𝐫^{}𝐫_{}^{}{}_{m}{}^{})`$
$`{\displaystyle \frac{1}{\mu _0}}\mathbf{}^{\mathbf{}}\times 𝐁_𝐦^{}ϵ_0_t^{}𝐄_𝐦^{}`$ $`=`$ $`0`$
$`\mathbf{}^{\mathbf{}}ϵ_0𝐄_𝐦^{}`$ $`=`$ $`0`$
$`\mathbf{}^{\mathbf{}}𝐁_𝐦^{}`$ $`=`$ $`q_m\delta (𝐫^{}𝐫_{}^{}{}_{m}{}^{})`$
where the relativistic coefficient $`\gamma `$ in the source terms cancelled out just like in the charge singularity case.
Consequently, the electromagnetic singular fields $`𝐄_𝐦`$ and $`𝐁_𝐦`$ are solution of
$$\begin{array}{cc}\hfill ϵ_0\mathbf{}\times 𝐄_𝐦\frac{1}{\mu _0}_t𝐁_𝐦& =q_m𝐯\delta (𝐫𝐫_m)\hfill \\ \hfill \frac{1}{\mu _0}\mathbf{}\times 𝐁_𝐦ϵ_0_t𝐄_𝐦& =0\hfill \\ \hfill \mathbf{}ϵ_0𝐄_𝐦& =0\hfill \\ \hfill \mathbf{}\frac{1}{\mu _0}𝐁_𝐦& =q_m\delta (𝐫𝐫_m).\hfill \end{array}$$
(26)
for a magnetic monopole singularity on an arbitrary path $`𝐫_m(t)`$ with a velocity $`𝐯=\dot{𝐫}_𝐦`$.
The equation of motion of the monopole singularity is deduced by defining the total field as the linear superposition of the singular fields $`𝐄_m`$ and $`𝐁_m`$ and the external fields $`𝐄_\mathrm{𝟎}`$ and $`𝐁_\mathrm{𝟎}`$:
$`𝐄`$ $`=`$ $`𝐄_𝐦+𝐄_\mathrm{𝟎}`$
$`𝐁`$ $`=`$ $`𝐁_𝐦+𝐁_\mathrm{𝟎}`$
where the external fields can be decomposed onto four potential fields $`V_{0q}`$, $`V_{0m}`$, $`𝐀_{\mathrm{𝟎}𝐪}`$ and $`𝐀_{\mathrm{𝟎}𝐦}`$
$`𝐄_\mathrm{𝟎}`$ $`=`$ $`\mathbf{}V_{0q}_t𝐀_{\mathrm{𝟎}𝐪}+{\displaystyle \frac{1}{ϵ_0}}\mathbf{}\times 𝐀_{\mathrm{𝟎}𝐦}`$ (27a)
$`𝐁_\mathrm{𝟎}`$ $`=`$ $`\mathbf{}\times 𝐀_{\mathrm{𝟎}𝐪}+\mu _0\mathbf{}V_{0m}+\mu _0_t𝐀_{\mathrm{𝟎}𝐦}`$ (27b)
$`0`$ $`=`$ $`\mathbf{}𝐀_{\mathrm{𝟎}𝐪}+ϵ_0\mu _0_tV_{0q}`$ (27c)
$`0`$ $`=`$ $`\mathbf{}𝐀_{\mathrm{𝟎}𝐦}+ϵ_0\mu _0_tV_{0m}.`$ (27d)
Introducing the total field in the electromagnetic Lagrangian (eq. 10) infers the following singularity action:
$$\mathrm{\Delta }𝒮_{sm}=\left(ϵ_0𝐄_\mathrm{𝟎}𝐄_𝐦\frac{1}{\mu _0}𝐁_\mathrm{𝟎}𝐁_𝐦\right)d^3x𝑑t$$
(28)
where we have kept only the terms that vary with respect to the trajectory of the singularity. Indeed, the terms in $`𝐄_𝐢𝐄_𝐢`$ and $`𝐁_𝐢𝐁_𝐢`$, where $`(i=m,0)`$, are not dependent on the position of the charge singularity. Rewriting the external field with the help of the potentials defined in equations (27a) and (27b) allows us to perform an integration by part
$`\mathrm{\Delta }𝒮_{sm}`$ $`=`$ $`{\displaystyle }(ϵ_0(\mathbf{}V_{0q}_t𝐀_{\mathrm{𝟎}𝐪}+\mathbf{}\times 𝐀_{\mathrm{𝟎}𝐦})𝐄_𝐦`$
$`{\displaystyle \frac{1}{\mu _0}}(\mathbf{}\times 𝐀_{\mathrm{𝟎}𝐪}+\mathbf{}V_{0m}+_t𝐀_{\mathrm{𝟎}𝐦})𝐁_𝐦)d^3xdt`$
$`=`$ $`{\displaystyle }(ϵ_0(V_{0q}\mathbf{}𝐄_𝐦+𝐀_{\mathrm{𝟎}𝐪}_t𝐄_𝐦+𝐀_{\mathrm{𝟎}𝐦}\mathbf{}\times 𝐄_𝐦)`$
$`{\displaystyle \frac{1}{\mu _0}}(𝐀_{\mathrm{𝟎}𝐪}\mathbf{}\times 𝐁_𝐦V_{0m}\mathbf{}𝐁_𝐦𝐀_{\mathrm{𝟎}𝐦}_t𝐁_𝐦))d^3xdt`$
where we have assumed the fields $`𝐄_𝐦`$ and $`𝐁_𝐦`$ to be zero at infinity i.e. the boundary of integration. Using equations (26) we can further simplify the action of the singularity
$`\mathrm{\Delta }𝒮_{sm}`$ $`=`$ $`{\displaystyle \left(q_mV_{0m}\delta (𝐫𝐫_q)q_m𝐀_{\mathrm{𝟎}𝐦}𝐯\delta (𝐫𝐫_q)\right)d^3x𝑑t}`$
$`=`$ $`{\displaystyle \left(q_mV_{0m}q_m𝐀_{\mathrm{𝟎}𝐦}𝐯\right)𝑑t}`$
which leads directly to a Lorentz type forces. Like in the case of the charge singularity, we can rigidly link a scalar mass singularity introduced in previous section to a monopole singularity. The total action of the double singularity is:
$`\mathrm{\Delta }𝒮_s+\mathrm{\Delta }𝒮_{sm}`$ $`=`$ $`{\displaystyle \left(U_0m_s\sqrt{1\frac{𝐯^2}{c^2}}+q_mV_{0m}q_m𝐀_{\mathrm{𝟎}𝐦}𝐯\right)𝑑t}`$ (29)
which lead to the following equation of movement
$`{\displaystyle \frac{d}{dt}}m_i𝐯`$ $`=`$ $`m_g\mathbf{}U_0+q_m𝐁_{\mathrm{𝟎}𝐦}+q_m𝐯\times 𝐄_{\mathrm{𝟎}𝐦}`$ (30)
where $`m_i`$ and $`m_g`$ are the defined by (7) and
$`𝐄_{\mathrm{𝟎}𝐦}`$ $`=`$ $`{\displaystyle \frac{1}{ϵ_0}}\mathbf{}\times 𝐀_{\mathrm{𝟎}𝐦}`$
$`𝐁_{\mathrm{𝟎}𝐦}`$ $`=`$ $`\mu _0\mathbf{}V_{0m}+\mu _0_t𝐀_{\mathrm{𝟎}𝐦}`$
correspond to the magnetic monopole part of the electromagnetic field. We remark here that the force acting on the monopole is only due to the monopole part of the electromagnetic field. Consequently, there is no direct interaction between charges and magnetic monopoles.
The energy balance equation is
$`{\displaystyle \frac{d}{dt}}m_ic^2`$ $`=`$ $`m_g{\displaystyle \frac{}{t}}U_0+q_m𝐯𝐁_{\mathrm{𝟎}𝐦}.`$ (32)
## 5 Conclusion
The field generated by a singularity gives rise to the laws of motion of the singularity. In other words, the action of the singularity on the field implies its dynamic reaction to the field. We have shown this to be the case for neutral bodies, charges and hypothetical magnetic monopoles. Using our method we have deduced not only the dynamics of masses, but also the relationship between the ”gravitational” source strength and the associated inertial mass. The charge singularity was shown to be the origin of the Lorentz force. The magnetic monopoles have an equivalent Lorentz force.
The conventional Lorentz force does not involve the magnetic monopole field. Reciprocally, the monopole Lorentz force does not involve fields originating from charged particles. Consequently, charges and magnetic monopoles do not interact, making the discovery of monopoles using charges difficult. A major question to be addressed concerns the role played by electromagnetic radiation. Can it provide the link between charges and magnetic monopoles?
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# Can the total angular momentum of 𝑢-quarks in the nucleon be accessed at Hermes?
## 1 Introduction
Over more than two decades, inclusive and semi-inclusive charged lepton scattering has been used as a powerful tool to successfully study the longitudinal momentum structure of the nucleon, which was parameterized in terms of parton distribution functions (PDFs). Hard exclusive reactions can be described in the theoretical framework of generalized parton distributions (GPDs) gpds1 ; gpds2 ; gpds3 ; gpds4 ; gpds5 . Their application became apparent after it had been shown jirule that measurements of the second moment of the sum of the ‘unpolarized’ GPDs $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ open, for the first time, access to the total angular momentum of partons in the nucleon:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{lim}$}}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(1)
In this relation $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ denote parton spin non-flip and spin flip GPDs ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}`$), respectively<sup>1</sup><sup>1</sup>1Throughout this paper the GPD definitions of a recent review markus are used.. GPDs depend on the fractions $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}`$ of longitudinal momentum of the proton carried by the parton and on $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$, the square of the 4-momentum transfer between initial and final protons (see. Fig. 1). As ordinary PDFs, also GPDs are subject to QCD evolution. Their $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ dependence has been perturbatively calculated up to next-to-leading order in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}`$ belevol and is omitted in the notations throughout the paper.
Recently, a simultaneous description of the transverse spatial and the longitudinal momentum structure of the nucleon was shown to be an appealing interpretation of GPDs impact1 ; impact3 ; impact4 ; impact2 . The concept of GPDs covers several types of processes, ranging from inclusive deeply inelastic lepton scattering to hard exclusive Compton scattering and meson production. Measurements of GPDs are expected to shed light especially on the hitherto theoretically uncharted territory of long-range (‘soft’) phenomena where parton-parton correlations are known to play an important role.
First steps towards the extraction of the GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}`$ have already been performed by scattering leptons off unpolarized protons through measurements of either cross sections h1dvcs ; zeusdvcs , or cross section asymmetries with respect to beam charge hermesdvcs or beam spin clasdvcs1 ; clasdvcs2 . Future measurements of the transverse target-spin asymmetry (TTSA) in hard exclusive electroproduction of a real photon (deeply virtual Compton scattering, DVCS) or a vector meson offer the possibility to acquire information on the spin-flip GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$. The most promising experiments to access it are those running at intermediate energy, where the spin-flip amplitude is expected to be sizable, while at higher energies it is suppressed due to $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}`$-channel helicity conservation. Thus at present a realistic program may be envisaged for Hermes, Clas and possibly Compass. In this paper the prospects are discussed for Hermes measurements of TTSAs in DVCS and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ electroproduction, and in particular their sensitivity to the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$-quark total angular momentum.
## 2 Modeling Generalized Parton Distributions
GPDs are most commonly parameterized using an ansatz based on double distributions radyushkin ; musatov complemented with the D-term dterm1 . Factorizing out the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-dependence, the non-forward GPDs can be related to the ordinary PDFs and the proton elastic form factors. In this framework gpv , the spin non-flip GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}`$ is given by
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.71}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(2)
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}`$=1.793 is the proton anomalous magnetic moment and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}`$ is the proton mass. The neutron Dirac form factor is neglected compared to the one of the proton.
For quarks, the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-independent part of the GPDs $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ is written as
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\theta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(3)
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}`$ is the D-term, and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}}`$ is the part of the GPD that is obtained from the double distribution (DD) $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\delta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(4)
For the double distributions the suggestion of Ref. radyushkin is used,
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{h}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(5)
where the profile function is given by musatov :
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{h}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Gamma }}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(6)
For $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{>}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is the ordinary quark density for the flavor $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$. The negative $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}`$ range corresponds to the antiquark density: $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$. The parameter $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}`$ characterizes to what extent the GPD depends on the skewness $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}`$. In the limit $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}`$ the GPD is independent on $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}`$, i.e., $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$. Note that $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}`$ is a free parameter for valence quarks ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}`$) or sea quarks ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$) and thus can be used as a fit parameter in the extraction of GPDs from hard electroproduction data frank .
For gluons, the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-independent part of the GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ is directly given by the double distribution,
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$
$`{\displaystyle \underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}{\displaystyle \underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\delta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ (7)
with the same form of the profile function in the double distribution
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{h}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(8)
The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-dependence for gluons is taken to be the same as that for quarks.
The factorized ansatz (2) is the simplest way of modeling GPDs. However, experimental studies of elastic diffractive processes indicate that the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-dependence of the cross section is entangled with its dependence on the photon-nucleon invariant mass collins . Recent evidence comes from lattice QCD calculations latticereg1 ; latticereg2 and phenomenological considerations markusreg ; vandreg . The non-factorized ansatz can be based on soft Regge-type parameterizations. In this case, the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-dependence is not factorized out and not controlled by a form factor as in Eq. (2). Instead, it is kept in Eqs. (3), (4) and (7). The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-dependence of double distributions is then modeled as gpv :
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(9)
which is referred to as Regge ansatz in the following. Here $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}`$ is the slope of the Regge trajectory, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.8}$}`$ GeV<sup>-2</sup> for quarks and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.25}$}`$ GeV<sup>-2</sup> for gluons.
In the factorized ansatz the spin-flip quark GPDs $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ are given by gpv :
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.71}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(10)
In the Regge ansatz the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$–dependence is modeled in analogy to Eq. (9).
The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-independent part is parameterized using the double distribution ansatz:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\theta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(11)
Note that the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}`$-term has the same size, but the opposite sign in Eqs. (11) and (3). Therefore, it drops out when calculating $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ according to Eq. (1).
The double distribution has a form analogous to the spin-nonflip case:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\delta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{K}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}$$
(12)
with:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{K}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{h}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\beta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(13)
The spin-flip parton densities $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ can not be extracted from deep-inelastic scattering (DIS) data, unlike the case of spin non-flip ones. Based on the chiral quark soliton model gpv , the spin-flip density is taken as a sum of valence and sea quarks contributions. Since in this model the sea part was found to be very narrowly peaked around $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$, the whole density is written as:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\delta }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(14)
In this expression, the shape of the valence quark part is given by that of the spin non-flip density. The coefficients $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ are constrained by the total angular momentum sum rule (1) and the normalization condition
$$\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(15)
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ is the anomalous magnetic moment of quarks of flavor $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1.67}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2.03}$}`$). The constraints yield:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$ (16)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$ (17)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\kappa }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$ (18)
Here $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ are the parton momentum contributions to the proton momentum:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(19)
In the given scenario the total angular momenta carried by $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$\- and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$-quarks, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$, enter directly as free parameters in the parameterization of the spin-flip GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$. Hence the parameterization (14) can be used to investigate the sensitivity of hard electroproduction observables to variations in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$.
As to the gluons, there exists no hint how the spin-flip GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ could be described. There is an expectation that $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ is not large compared to $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ ournew . Hence, for simplicity throughout the present study $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ is neglected (“passive” gluons, i.e. $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$).
As an example, Fig. 2 shows the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-independent part of various GPDs at $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.1}$}`$, based on the MRST98 mrst98 parameterization of PDFs at $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$=4 GeV<sup>2</sup>. Using instead CTEQ6L PDFs cteq as input, the results for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ quark GPDs are changed by less than 3%(10%); the GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ is up to 40% larger at $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$. Because of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$-quark dominance in electroproduction, uncertainties originating from $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$-quark PDFs can be safely neglected. Since gluons are absent in leading-order DVCS, uncertainties resulting from gluon PDFs are of little influence for DVCS asymmetries and have been found to lead to a fractional change of up to 15% for the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ asymmetries. For the following calculations the MRST98 PDF set is taken.
## 3 Sensitivity of DVCS to the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$-quark Total Angular Momentum
### 3.1 Cross Section and Asymmetries
The 5-fold cross section for the process $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is given by:
$$\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{16}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left|}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(20)
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ is the negative squared 4-momentum of the virtual photon, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is the Bjorken variable, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}`$ denotes the photon production amplitude and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}`$ is the electron charge. Since the DVCS and Bethe-Heitler (BH) processes have an identical final state, in which the photon is radiated either from a parton or from a lepton, respectively, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}`$ is given by the coherent sum of the BH amplitude $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}`$ and the DVCS amplitude $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}}`$:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(21)
in which
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{D}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{V}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}$$
(22)
describes the interference between both processes.
The coordinate system is defined in the target rest frame, as explained in Fig. 3. The theoretical formulae used below refer to the target being transversely polarized w.r.t. the virtual photon direction, while in the experiment the target polarization is transverse w.r.t. the incident lepton direction. At Hermes kinematics, these two directions are approximately parallel and the small longitudinal component ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{<}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{10}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}`$) of the target polarization along the virtual photon direction can be neglected. Thus the reasonable approximation
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}}$$
(23)
is used, where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}`$ ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}}`$) denotes the cross section for the unpolarized (transversely polarized) component.
Since in the kinematic region of the Hermes experiment the DVCS cross section is typically much smaller than the BH cross section kornow , the contribution of the DVCS term to the total cross section is neglected in the following. The contributions of the BH term for an unpolarized beam are:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\times }$}`$ (24)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$
The full expressions for the BH propagators $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ and for the Fourier coefficients $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}`$ can be found in Ref. belitsky <sup>2</sup><sup>2</sup>2The azimuthal angles defined in this work are different from those used in Ref. belitsky : $`\colorbox[rgb]{1,1,1}{$\varphi $}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$\pi $}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$\varphi $}_{\colorbox[rgb]{1,1,1}{$\text{belitsky}$}}`$ and $`\colorbox[rgb]{1,1,1}{$\varphi $}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$\varphi $}_\colorbox[rgb]{1,1,1}{$S$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$\pi $}\colorbox[rgb]{1,1,1}{$+$}\colorbox[rgb]{1,1,1}{$\phi $}_{\colorbox[rgb]{1,1,1}{$\text{belitsky}$}}`$..
The leading twist and leading order $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}`$ contributions of the DVCS-BH interference term to the total cross section can be written as:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\times }$}`$ (25)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$
Here $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}`$ ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$) stands for a negatively (positively) charged lepton beam and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is a kinematic pre-factor independent of azimuthal angles. The full expressions for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}`$ can be found in Eqs. (53-56) of Ref. belitsky . $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$ and $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}`$ are certain linear combinations of the Compton form factors $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$, $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}`$ and $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}`$, which are convolutions of the respective twist-2 GPDs $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$, $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}`$ and $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}`$ with the hard-scattering kernels as defined in Eq. (9) of Ref. belitsky .
The full expressions for $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$ and $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}`$ can be found in Eq. (71) in Ref. belitsky or in Eq. (60) in Ref. markusapet . Since $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is small in a wide range of experimentally relevant kinematics, terms with pre-factor $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}`$ or $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}`$ can be neglected, except for the GPD $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}`$ because the pion pole contribution to $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}`$ scales like $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$, so that $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$ and $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}`$ can be approximated as:
$`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$ (26)
Here $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ are the Dirac and Pauli form factors of the proton, respectively.
In order to constrain the GPDs involved in Eq. (26), the transverse polarization component of the interference term, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}}`$, has to be singled out. This can be accomplished by forming the transverse (T) target-spin asymmetry with unpolarized (U) beam:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}`$ (28)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$
As $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}`$ are independent on $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}`$, they do not appear in the numerators of Eq. (28). Since their dominant contribution to the denominator in Eq. (28) is given by $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}`$, the two amplitudes of the TTSA, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$, can be approximated as:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{c}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$ (29)
Note that the approximations used in this section are for illustrative purposes only and are not used in the numerical calculations described below.
### 3.2 Expected Value of TTSA and Projected Statistical Uncertainty
Since in the DVCS process the gluons enter only in NLO in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}`$, their contributions to cross section and TTSA are neglected. For the quarks, it can be seen from Eq. (26) that, besides the GPDs $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ which have been discussed in Sect. 2, there are two other GPDs, $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}`$ and $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}`$, involved in the TTSA for DVCS. Since they are not the main interest of this paper, in the calculations below they are always included and kept unchanged. In their model description, the forward limit of the GPD $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}`$ is fixed by the quark helicity distributions $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, while the GPD $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}`$ is evaluated from the pion pole which only provides a real part to $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}`$ in Eq. (26).
At present, there exists a code vgg designed to calculate observables in the exclusive reaction $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}`$. It has been used (see App. Appendix: TTSA Calculation in DVCS) to evaluate the TTSA arising from the DVCS-BH interference. The TTSA is calculated at the average kinematic values per bin in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$ taken from a measurement of the beam-spin asymmetry in DVCS at Hermes frank (see Tab. 1).
The statistical error of an asymmetry is independent on its size if the asymmetry itself is small. For a single beam (target) spin asymmetry it is obtained as:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(30)
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$ is the total number of events that is proportional to the integrated luminosity, and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ is the beam (target) polarization. The following projection is based on a future Hermes data set of 8 million DIS events to be taken with an unpolarized positron beam and a transversely polarized hydrogen target. Using the known statistical errors of the beam-spin asymmetry measurement at Hermes on an unpolarized hydrogen target (7 million DIS events, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{50}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\%}$}`$) frank , the projected statistical error for the TTSA is obtained.
The projections for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ are calculated for different values of the total angular momentum $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$. Since the contributions of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$-quark and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$-quark are proportional to the corresponding squared charge, the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$-quark contribution is suppressed and hence in the calculations a fixed value is used for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$. The latter was chosen to be $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$, inspired by the results of recent lattice calculations (see e.g. Ref. lattice\_prl\_qcdsf ). Using both Regge and factorized ansätze, the asymmetries are calculated for the four possible cases setting the profile parameters $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ to either one or infinity. Comparing all sets of projections to each other, the amplitudes of the TTSA appear to be sensitive only to the change in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ from one to infinity. The resulting differences are small and can be seen by comparing Figs. 4 and 5, where the amplitudes are shown in dependence on $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$ together with the projected statistical errors. In order to study the contributions of the GPDs $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}`$, $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}}`$ and $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{~}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}}`$ alone, calculations are done for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ as well.
As expected from Eqs. (26) and (3.1), variations in the parameter settings for the GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ become manifest in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ while $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ shows only minor modifications. The latter are apparent only in the kinematic regime of large $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}`$ or correspondingly large $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ since the contribution of the GPDs $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ to $`\widehat{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}`$ is suppressed by $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}`$ and thus has been neglected in Eq. (26). Within these model calculations $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ turns out to be sizable even when the calculation is done for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$. Thus a solid knowledge about the GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ is needed in order to constrain $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$. It has been shown dis04 that the model parameters for the GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$, in particular the size of the profile parameters $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$, can be well constrained by the envisaged Hermes DVCS measurements until 2007, using an unpolarized hydrogen target. Since in addition the profile parameters are assumed to be the same for the GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$, the only remaining free parameter is $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$. Hence the projected measurement of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ has a clear potential to constrain $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$, as can be seen from the left panels of Figs. 4 and 5.
The discriminative power of the envisaged TTSA measurement can be enhanced by combining all data into one point, since all considered models show the same kinematic dependences. The corresponding statistical power of a Hermes data set based on 8 million DIS events is shown in Fig. 6, for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ equal to one or infinity, and for three different values of the total $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$-quark angular momentum $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ plus the special case $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$.
It appears that for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}`$) the amplitude ranges between values of -0.17 and -0.27 (-0.19 and -0.29) when $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ ranges between 0.4 and 0. The projected statistical error for these integrated TTSA amplitudes is 0.017. Extrapolating the knowledge on the systematic uncertainty from the analysis of 2000 Hermes data frank , its size can be expected to not exceed the statistical error, such that a total experimental uncertainty below 0.025 appears as a realistic estimate. Altogether, the difference in the size of the TTSA due to a change of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ between zero and 0.4 corresponds to a 4$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}`$ effect, where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}`$ denotes the total experimental uncertainty. Thus, based on the GPD model used it can be expected that the upcoming DVCS results from Hermes<sup>3</sup><sup>3</sup>3The recent switch of the accelerator (HERA) to an electron beam will require to also perform the above calculations for the negative beam charge. However, the sensitivity to $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ of the combined electron and positron measurements is expected to be similar to the one calculated here for a positron beam only. will provide a constraint on the size of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$.
## 4 Sensitivity of Elastic $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ Electroproduction to the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$-quark Total Angular Momentum
Also a measurement of the TTSA in elastic vector meson electroproduction can be a source of information about the spin-flip generalized parton distribution $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$. An estimate for the asymmetry was obtained in Ref. gpv , using the factorized model of GPDs described in Sect. 2 without inclusion of gluons. The scope of this section is to also include the Regge ansatz, to check the assumption that the gluon contribution to the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ electroproduction cross section is small, and to eventually calculate the size of the TTSA at Hermes kinematics. The issue is raised since in contrast to DVCS, in vector meson elastic electroproduction gluons enter at the same order of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}`$ as quarks, namely at order $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}`$ to the power one. Hence this channel appears as one of the rare cases where gluon GPDs may be accessed through Hermes data.
### 4.1 Cross Section and Gluonic Contribution
It was shown frankstrik97 that the leading twist contribution to exclusive electroproduction of vector mesons requires both the virtual photon and the vector meson to be longitudinal, i.e. transversely polarized. Therefore the present calculations cover only the longitudinal part of the cross section.
The cross section of the reaction $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}`$ is given by gpv
$$\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{L}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{8}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{W}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}$$
(31)
where $`\stackrel{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ is the momentum of the virtual photon in the center of mass system of this photon and the initial proton, while $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{W}$}`$ is their invariant mass. The spin-flip amplitude reads
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{9}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒜}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}{\displaystyle \underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Phi }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}}}`$ (32)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒜}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{8}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{9}$}}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}}}{\displaystyle \underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Phi }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
and the spin-flip one is<sup>4</sup><sup>4</sup>4in the subsequent calculations the exact formulae were used:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}{\displaystyle \frac{\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{9}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\nu }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}{\displaystyle \underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Phi }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}}}`$ (33)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{9}$}}}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}}}{\displaystyle \underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝑑}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Phi }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$
Here $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ is a light-like vector along the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}`$-axis, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ is the 4-momentum transfer ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$). The modulus of its transverse component is given by $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}`$. The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$-meson wave function is taken in the form
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Phi }}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{6}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}$$
(34)
with $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{f}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}`$=0.216 GeV and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{z}$}`$ being the meson longitudinal momentum fraction carried by a parton. The complex factors $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒜}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ are given by:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒜}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}}{\displaystyle \underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$
$`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{8}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\times }$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{ϵ}$}}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{ϵ}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$ (35)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}}{\displaystyle \underset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\overset{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}`$
$`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{8}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\times }$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\{}$}{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{ϵ}$}}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{i}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{ϵ}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$ (36)
The TTSA is defined as
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}`$ (37)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$
The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ amplitude of the TTSA can be expressed in terms of the spin flip and spin non-flip amplitudes as gpv :
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$
$`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{I}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒜}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{\Delta }}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒜}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\xi }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{R}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒜}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$
Note that using the Trento convention Trento the sign of this equation is opposite to that in Ref. gpv and the normalization is larger by a factor of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$.
The cross section is calculated using both factorized and Regge ansätze for GPDs<sup>5</sup><sup>5</sup>5The principal values of the integrals in Eqs. (35) and (36) are calculated in the following way: $`\underset{\colorbox[rgb]{1,1,1}{$a$}\colorbox[rgb]{1,1,1}{$<$}\colorbox[rgb]{1,1,1}{$0$}}{\overset{\colorbox[rgb]{1,1,1}{$b$}\colorbox[rgb]{1,1,1}{$>$}\colorbox[rgb]{1,1,1}{$0$}}{\colorbox[rgb]{1,1,1}{$$}}}\frac{\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$x$}\colorbox[rgb]{1,1,1}{$)$}}{\colorbox[rgb]{1,1,1}{$x$}}\colorbox[rgb]{1,1,1}{$𝑑$}\colorbox[rgb]{1,1,1}{$x$}\colorbox[rgb]{1,1,1}{$=$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$b$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$\mathrm{ln}$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$b$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$$}\colorbox[rgb]{1,1,1}{$f$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$a$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$\mathrm{ln}$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$a$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$$}\underset{\colorbox[rgb]{1,1,1}{$a$}}{\overset{\colorbox[rgb]{1,1,1}{$0$}}{\colorbox[rgb]{1,1,1}{$$}}}\colorbox[rgb]{1,1,1}{$f$}^{\colorbox[rgb]{1,1,1}{$$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$x$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$\mathrm{ln}$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$x$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$𝑑$}\colorbox[rgb]{1,1,1}{$x$}\colorbox[rgb]{1,1,1}{$+$}\underset{\colorbox[rgb]{1,1,1}{$0$}}{\overset{\colorbox[rgb]{1,1,1}{$b$}}{\colorbox[rgb]{1,1,1}{$$}}}\colorbox[rgb]{1,1,1}{$f$}^{\colorbox[rgb]{1,1,1}{$$}}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$x$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$\mathrm{ln}$}\colorbox[rgb]{1,1,1}{$($}\colorbox[rgb]{1,1,1}{$x$}\colorbox[rgb]{1,1,1}{$)$}\colorbox[rgb]{1,1,1}{$𝑑$}\colorbox[rgb]{1,1,1}{$x$}`$. In this way the non-integrable singularity is exchanged by an integrable one.. The value $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ is taken for the profile parameter both for sea and valence quarks. It is found that using $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}`$ instead of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ leads to a rise of the cross section by a factor of about $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1.15}$}`$. The value $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{v}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{l}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ is chosen to provide a direct comparison to previous calculations gpd ; gpv . The value of the profile parameter for gluons is chosen as $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}`$ = 2 and it has been checked that choosing $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}`$ = 1 or $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}`$ does not change the cross section by more than 20%. The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{W}$}`$-dependence of the cross section for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$=4 GeV<sup>2</sup> is shown in Fig. 7. For both ansätze the calculations overshoot considerably the experimental data from Hermes rhoherm . However, a significant reduction of the calculated cross section might be expected if transverse motion effects are taken into account gpv ; gpd . On the other hand, also the double distribution based calculations of the DVCS cross section have been found to overshoot the data from H1 freundmac1 ; freundmac2 .
An unexpected result of the calculation shown in Fig. 7 is a quite small (15-20%) pure quark contribution to the cross section, while in Refs. gpv ; gpd the quark contribution was found to be dominant. Comparing the calculated quark contribution to experimental data (Fig. 7) it could also be concluded that the gluon contribution in the present calculation is substantially overestimated, while the quark contribution itself is reasonable and can explain alone (in the factorized ansatz) the value of the measured cross section. However, there exists experimental evidence that the gluon spin non-flip part is indeed large ournew .
On the amplitude level, the cross sections of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}`$ mesons are given as:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ (39)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\phi }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{9}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$
Here $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}`$ is the quark amplitude and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ the gluon amplitude (the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}`$-quark contribution to the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}`$ production amplitude is neglected) and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\phi }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}`$ is the effective phase between the quark and gluon amplitudes. In the existing GPD-based calculations gpv ; gpd , both quark and gluon contributions are dominated by the imaginary parts which have the same sign, i.e. $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\phi }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$. In the present calculation $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\phi }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{30}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}`$ is obtained. Considering the wave functions of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}`$ mesons to be similar (as it is supported by the measured values of their decay widths), $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}`$ follows, and the ratio of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}`$ to $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ cross sections reads:
$$\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{9}$}}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\phi }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(40)
At Hermes, the ratio of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}`$ was measured phitorho . The experimental value was 0.08$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}`$0.01, slightly increasing with $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$. Inserting it to the l.h.s of Eq. (40) and taking $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\phi }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}}`$ = 0(30) yields $`\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{r}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{m}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}}`$ = 0.7 (0.78). This value is in good agreement with the results of the present calculation, where $`\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}`$ ranges between 0.8 and 0.5 (0.5 and 0.3) for the factorized (Regge) ansatz when $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{W}$}`$ increases from 4 to 6 GeV. This is in contrast to the above mentioned result of a dominant quark contribution, $`\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{𝒯}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{|}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}`$ gpd ; gpv . Hence, it is concluded that $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ can not be neglected, i.e. to arrive at the measured cross section both quark and gluon amplitudes have to be scaled down in a similar proportion.
### 4.2 Expected Value of TTSA and Projected Statistical Uncertainty
The $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ amplitude of the TTSA in Eq. (4.1) is calculated at Hermes kinematics. The statistical error is extrapolated from a preliminary analysis of the Hermes longitudinal target-spin asymmetry measured on the deuteron ulrike that is based on 8 million DIS events. In the latter analysis the data is not split into parts corresponding to longitudinal and transverse virtual photons, while the present calculation is related to longitudinal photons only. At Hermes kinematics ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Q}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ GeV<sup>2</sup>), longitudinal photons constitute about 50% of all virtual photons. Also, the transverse target polarization is 0.75 while the longitudinal one is 0.85. The projected statistical error for 8 million DIS events taken on a transversely polarized target is then larger by a factor $`\sqrt{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.85}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.75}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1.6}$}`$ compared to that of Ref. ulrike . Note that this error estimate may be considered ‘optimistic’, since it assumes that the contribution to the asymmetry from longitudinal and transverse photons can be completely disentangled.
The calculated $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-dependences of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ are shown in Fig. 8 for different values of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$. As in the case of DVCS, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$ is fixed inspired by the fact that the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}`$-quark contribution is still suppressed, although the suppression in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ production is half as strong as in DVCS. Again, the choice of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ is based on the results of recent lattice calculations (see e.g. Ref. lattice\_prl\_qcdsf ). Note that in contrast to DVCS, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ results in a vanishing asymmetry. As it can be seen from comparing Fig. 8 to Figs. 4 and 5, the expected magnitude of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ production is much smaller than that in DVCS. This is due to a large gluonic contribution to the amplitude, which is considered as “passive” ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$), i.e. the gluons dilute the asymmetry in this case. It was found that the difference in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ between the factorized and Regge ansätze is negligible. Also the variation of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ only leads to a small difference as can be seen when comparing the left and right panels of Fig. 8, where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{B}$}`$\- and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$-dependences of the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ amplitude of the asymmetry are shown for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}`$, respectively. The amplitude of the integrated TTSA is shown in Fig. 9, for the same two cases. It is essentially independent of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{a}$}}`$ and ranges between values of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.10}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.01}$}`$ when $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ ranges between zero and 0.4.
The projected statistical error for the integrated TTSA amplitudes is 0.034. Extrapolating the knowledge on the systematic uncertainty from Ref. ulrike , its size can be expected to be about 0.02, such that a total experimental uncertainty below 0.04 appears as a realistic estimate. Altogether, the difference in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}`$ due to a change of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ between zero and 0.4 corresponds to an about 2$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}`$ effect, where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}`$ denotes the total experimental uncertainty. Thus it can be expected that the upcoming $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ electroproduction measurements performed at Hermes will provide an additional constraint on the size of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$.
A tempting possibility provided by $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ production is related to an estimate of the gluonic content of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$. Strongly simplifying, Eq. (4.1) represents the ratio $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{q}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$. Hence, when comparing the earlier calculations gpv where gluons have been neglected ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$) to the case of “passive” gluons presented above ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$), the asymmetry gets smaller (‘diluted’) by the presence of the term containing $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ in the denominator. On the other hand, if the measured asymmetry would be found large, this could imply that the gluons are “active” ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$), so that their contribution to the spin-flip amplitude can not be neglected.
## 5 Summary and Outlook
Transverse target-spin asymmetries (TTSAs) in DVCS and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ elastic electroproduction are the only candidates known by now to access the GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ on a proton target, in which $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ comes as a leading term. A code vgg based on the model developed in Ref. gpd ; gpv is used to calculate the expected TTSAs to be measured in DVCS on the Hermes transversely polarized hydrogen target. To check the accessibility of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ at Hermes, different parameterization ansätze and parameters of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{H}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ are chosen. As the model for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ depends on the total angular momentum of the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$-quarks in the proton, the possibility arises to check the sensitivity of the data to different values chosen as $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.4}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0.0}$}`$, while on the basis of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$-quark dominance and recent lattice calculations (see e.g. Ref. lattice\_prl\_qcdsf ) a fixed value $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ is used. The calculations are performed at the Hermes average kinematic values frank . The results show that the DVCS TTSA amplitude $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ is sensitive to the GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}`$ and thus to the total $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$-quark angular momentum $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$, while $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ is not. It was found that aside from $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ the amplitude $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ is largely independent on different parameterization ansätze and model parameters. Projected statistical errors for the asymmetries are evaluated by converting the ones from Ref. frank to a data set corresponding to 8 million DIS events taken on a transversely polarized hydrogen target.
The same parameterizations are used to calculate the TTSA in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ electroproduction by longitudinal virtual photons. The main difference to the DVCS case is a large gluonic contribution to the amplitude. At present, only the spin-nonflip part of the gluonic amplitude can be reasonably described, while the spin-flip gluonic GPD $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ is totally unknown. Therefore, throughout the calculation $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ is set to zero (“passive” gluons). Under this assumption, the situation in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ electroproduction appears less favorable concerning the sensitivity of the expected TTSA amplitude to the total angular momentum $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$. However, should the value of the amplitude be measured larger than that predicted by these calculations, this would imply that $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{E}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{g}$}`$ can not be neglected, and thus indicate that gluons inside the proton carry significant orbital angular momentum.
Altogether, transverse target-spin asymmetries in both DVCS and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ electroproduction are studied to evaluate projected uncertainties for extracting the value of $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$ from future data. Considering all anticipated Hermes data to be taken for DVCS ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\rho }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$-production), the projected total experimental $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}`$-uncertainty is estimated to correspond to a range of about 0.1 (0.2) in $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{J}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}`$.
## Acknowledgments
This study would have been impossible without the permanent advice of M. Diehl. The support of J. Volmer is highly appreciated by Z.Y. A.V. is grateful to A. Borissov for useful discussions. The advice of E. Aschenauer is appreciated. A.V. was supported by the Alexander von Humboldt foundation, RFBR grants 04-02-16445, 03-02-17291 and Heisenberg-Landau program. This work was supported in part by the US Department of Energy.
## Appendix: TTSA Calculation in DVCS
A code vgg is used to estimate the TTSA related to DVCS. The coordinate system and angles defined in the code are the same as depicted in Fig. 3. The polarization of the target in the code is defined according to the virtual photon direction. For a transversely polarized target, the target polarization direction can be chosen either in the lepton plane ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}`$ direction) or perpendicular to it ($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}`$ direction). The former corresponds to $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}`$ or $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}`$, the latter to $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ or $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$. Therefore the following intermediate asymmetries can be calculated:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`{\displaystyle \frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\frac{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{3}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pi }$}}{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$ (41)
Defining the following functions
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$ (42)
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{x}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{y}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
the contribution of the transverse target polarization component of the interference term $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}}`$ to the total cross section in Eq. (25) can be expressed as:
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{P}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{d}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\sigma }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{u}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{n}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{p}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\left[}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\right]}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}$$
(43)
Therefore the asymmetries defined in Eq. (28) can be computed as:
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}`$
$`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{U}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{T}$}}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{cos}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}`$ $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{A}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{sin}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\varphi }$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{.}$}`$
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# Matrix product representation of gauge invariant states in a 𝑍₂ lattice gauge theory
## 1 Introduction
The importance of the first-principle study in quantum chromodynamics is increasing largely because RHIC experiment has started and LHC is also coming. For precise description of high-energy heavy ion collisions, gauge theory needs to be studied at finite temperature and density in a systematic way. Ideally, we should also have a methodology for tracing time-evolution of quantum states based on the Schrödinger equation because heavy ion collisions should be treated as non-equilibrium evolving systems rather than static. Lattice gauge theory is the most useful method for studying the quark-gluon systems at zero and finite temperature. However, Monte Carlo integration does not work for lattice gauge theory with large chemical potential because of the severe sign problem. It would be worthwhile to pursue a systematic variational approach to gauge theory . In this paper, we apply the matrix product ansatz to a $`Z_2`$ hamiltonian lattice gauge theory on a spatial ladder lattice based on the previous work in a U(1) lattice gauge theory .
The matrix product ansatz is a variational method that originates in DMRG (density matrix renormalization group) . DMRG has been developed as the method that gives the most accurate results for spin and fermion chain models such as one-dimensional quantum Heisenberg and Hubbard models at zero and finite temperature . <sup>1</sup><sup>1</sup>1By “$`d`$-dimensional”, we mean ($`1+d`$)-dimensional spacetime. DMRG is also useful for diagonalization of transfer matrices in two-dimensional classical statistical systems . DMRG has been extended to two-dimensional quantum systems and can work for bosonic degrees of freedom . Since DMRG is a variational method based on diagonalization of hamiltonian and transfer matrices, it is free from the sign problem. Actually, DMRG has been successful in accurate study of $`\theta `$-vacuum in the massive Schwinger model . There is an old prediction by S. Coleman that quarks deconfine at $`\theta =\pi `$ . However, the model has not been analyzed accurately with the Monte Carlo method because the topological terms give complex action in the Euclidean theory.
The matrix product ansatz is a result of large simplification of trial wavefunction based on the knowledge established in the past DMRG studies. Therefore, the ansatz takes over the good points of DMRG. Besides, the matrix product ansatz has advantage to DMRG because the former can treat periodic one-dimensional systems accurately . Recent interesting progress of the matrix product ansatz is its application to non-equilibrium quantum physics and quantum information theory . The matrix product ansatz is a promising approach to further refinement and extension of the past DMRG studies.
In general, calculation of expectation values of hamiltonian becomes difficult exponentially as the system size increases. If the matrix product ansatz is introduced, the energy function has a simple matrix product form, which can be evaluated easily using a computer if the matrix size is small. It is expected that exponentially difficult problems are reduced into small tractable ones by the matrix product ansatz. Actually, the ansatz has been successful in giving accurate results in many one-dimensional quantum systems, where calculation errors can be controlled systematically. We can say that the matrix product ansatz is a first-principle variational method.
Lattice gauge hamiltonian is obtained by choosing temporal gauge in partition function of Euclidean lattice gauge theory . In hamiltonian formalism, gauge invariance needs to be maintained explicitly by imposing the Gauss law on the Hilbert space. It is a hard task to construct gauge-singlet variational space for general gauge group . On the other hand, Euclidean lattice gauge theory can keep gauge invariance manifestly by construction. This is one of the reasons why hamiltonian version of lattice gauge theory is not popular. In addition, no systematic methods had been known for diagonalization of gauge hamiltonian before the matrix product ansatz was applied to lattice gauge theory in ref. . If trial wavefunction is constrained directly with the Gauss law, the advantage of the matrix product ansatz is completely spoiled because calculation of energy function becomes impossible in a practical sense. If the hamiltonian is diagonalized without the Gauss law, all possible states are obtained including gauge variant states. However, it must be possible to extract gauge invariant states because all eigenstates of the hamiltonian can be classified using generators of the considered gauge group. Therefore, if the matrix product ansatz is used, we better start from the whole Hilbert space and then identify gauge invariant states using the Gauss law operator after all calculations.
In hamiltonian lattice gauge theory, the dimension of spatial lattice needs to be two or larger because the plaquette operator is a two-dimensional object. In this paper, we study a $`Z_2`$ lattice gauge theory on a spatial ladder chain for simplicity. The length of the chain needs to be sufficiently long for the matrix product ansatz to work well. The ladder chain is squashed into a chain because one-dimensional structure needs to be found in order to use the matrix product ansatz.
In the previous work , we have studied one- and two-dimensional $`S=1/2`$ Heisenberg models and a U(1) lattice gauge theory on a ladder chain using the matrix product ansatz in a similar way, where energy function is minimized using the Powell method. When the number of parameters is very large, such naive minimization is not useful because it takes long time to reach the bottom of the energy function. In this work, we use a diagonalization method introduced in ref. to obtain sufficient accuracy.
This paper is organized as follows. In section 2, hamiltonian lattice formulation of the $`Z_2`$ lattice gauge theory is briefly reviewed. In section 3, the matrix product ansatz is introduced and applied to a $`Z_2`$ lattice gauge hamiltonian. In the original construction, the matrix product states is assumed to have translational invariance. In this work, that condition is not imposed on variational space before diagonalization of the gauge hamiltonian. In section 4, numerical results are given. Section 5 is devoted to summary.
## 2 Quantum hamiltonian in the $`Z_2`$ lattice gauge theory
We are going to introduce the $`Z_2`$ lattice gauge theory, which was invented by F. Wegner . As seen in the literature, the simplicity of the model is useful for testing a new idea . The model cannot have non-vanishing magnetization because local gauge symmetry cannot break spontaneously, which is known as the Elitzur’s theorem . However, the model can have nontrivial phases depending on dimensionality.
We are interested in quantum hamiltonian of the model. Statistical mechanics and quantum hamiltonian are connected through the transfer matrix formalism. In the $`Z_2`$ lattice gauge theory, quantum hamiltonian is obtained by choosing temporal gauge in the partition function
$$H=\underset{n,i}{}\sigma _x(n,i)\lambda \underset{n,i,j}{}P(n,i,j),$$
(1)
where $`\sigma _x`$ and $`\sigma _z`$ are spin operators
$$\sigma _x=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _z=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
and $`P`$ is a plaquette operator
$$P(n,i,j)\sigma _z(n,i)\sigma _z(n+i,j)\sigma _z(n+i+j,i)\sigma _z(n+j,i).$$
(2)
In eq. (1), the first and second summations are taken on the spatial lattice for all possible link and plaquette operators, respectively. In general, arbitrary states can be represented as a superposition of products of $`|\pm _{n,i}`$, which are eigenstates of the spin operator $`\sigma _z(n,i)`$
$$\sigma _z(n,i)|\pm _{n,i}=\pm |\pm _{n,i}.$$
Let us introduce time-independent operators $`G(n)`$, each of which flips spins on all the links emerging from a site $`n`$
$$G(n)=\underset{\pm i}{}\sigma _x(n,i).$$
(3)
We have
$`G^1(n)\sigma _x(m,i)G(n)`$ $`=`$ $`\sigma _x(m,i),`$
$`G^1(n)\sigma _z(n,i)G(n)`$ $`=`$ $`\sigma _z(n,i),`$
$`G^1(n)\sigma _z(m,i)G(n)`$ $`=`$ $`\sigma _z(m,i),`$
where the last formula applies only if the link $`(m,i)`$ is not contained in $`G(n)`$. The operator $`G(n)`$ defines local gauge transformation
$$G(n)^1HG(n)=H.$$
(4)
In order for physical quantities to be gauge invariant, quantum states need to be invariant under gauge transformation
$$G(n)|\mathrm{\Psi }=|\mathrm{\Psi }.$$
(5)
We need to impose the Gauss law (5) on the wavefunction to keep gauge invariance. Otherwise, unphysical states may be obtained because gauge invariance is not guaranteed. When a state $`|\mathrm{\Psi }`$ satisfies the Gauss law (5), magnetization vanishes because a relation $`\mathrm{\Psi }|\sigma _z(n)|\mathrm{\Psi }=\mathrm{\Psi }|\sigma _z(n)|\mathrm{\Psi }`$ holds.
## 3 Matrix product ansatz on a ladder lattice
We are going to introduce the matrix product ansatz, which is a variational method inspired from density matrix renormalization group (DMRG). DMRG is a variational method that can reproduces very accurate results in one-dimensional quantum systems . In DMRG, wavefunction is represented as a product of orthogonal matrices because basis states are rotated for optimization with orthogonal matrices that diagonalize density matrices. The success of DMRG allows us to parametrize wavefunction as a product of finite-dimensional matrices from the beginning. This simplification of wavefunction is called the matrix product ansatz . Although DMRG has slow convergence in one-dimensional quantum systems with periodic boundary conditions, the matrix product ansatz gives much better accuracy .
In the original construction, the matrix product state is parametrized as follows :
$$|\mathrm{\Psi }=\mathrm{tr}\left(\underset{n=1}{\overset{L}{}}\underset{s_n}{}A[s_n]|s_n\right),$$
(6)
where $`\{|s_n\}`$ is a complete set of basis states for the $`n`$-th site (see figure 1). Hamiltonian needs to have periodicity for consistency with the trace operation in the variational state. As a result, energy becomes a function of the matrices $`A[s]`$. The minimum of the energy function corresponds to the ground state. We are going to apply the ansatz to a $`Z_2`$ gauge theory and see its compatibility with gauge symmetry.
Since this work is the first application of the matrix product ansatz to $`Z_2`$ gauge theory, we would like to consider a simple model. The simplest one is a $`Z_2`$ hamiltonian lattice gauge theory on a spatial ladder lattice (see figure 2). We assume periodicity in the horizontal direction on the ladder for later convenience. In figure 2, periodicity is denoted with the open circles.
In the free case $`\lambda =0`$, the hamiltonian (1) can be diagonalized analytically. The vacuum state is given by
$$|\mathrm{vac}=\underset{n,i}{}\frac{1}{\sqrt{2}}(|+_{n,i}+|_{n,i}).$$
(7)
Vacuum expectation values of the hamiltonian and plaquette operators are $`H/L=3`$ and $`P=0`$, respectively. In this case, one-dimensional matrices are sufficient to represent the vacuum state (7)
$$A[+]=A[]=\frac{1}{\sqrt{2}}.$$
(8)
When the coupling constant is very large $`\lambda \mathrm{}`$, the first term of the hamiltonian (1) can be neglected. As a result, the hamiltonian has a diagonal form. The vacuum state is given by
$$|\mathrm{vac}=\frac{1}{\sqrt{2^{2L+1}}}\underset{m=0}{\overset{2L}{}}G^{(m)}\underset{n,i}{}|\pm _{n,i},$$
(9)
where the operator $`G^{(m)}`$ represents a summation of all possible products that are composed of $`m`$ pieces of different Gauss-law operators $`G`$. The operator $`G^{(m)}`$ has $`{}_{2L}{}^{}C_{m}^{}`$ terms ($`G^{(0)}1`$). In eq. (9), the two definitions are identical. Then, we have $`H/L\lambda `$ and $`P1`$ in the limit $`\lambda \mathrm{}`$. The states (7) and (9) are both gauge invariant, $`G(n)|\mathrm{vac}=|\mathrm{vac}`$.
The $`Z_2`$ lattice gauge model has only link variables. In our construction, each link is assigned a different set of matrices $`A_n`$, $`B_n`$, and $`C_n`$ for parametrization of wavefunction (see figure 2). The index $`n`$ represents the $`n`$-th square on the ladder chain and runs from $`1`$ to $`L`$. The dimension of the matrices is $`M`$. Then, our matrix product state is
$$|\mathrm{\Psi }=\mathrm{tr}\left(\underset{n=1}{\overset{L}{}}\underset{s_n=\pm }{}\underset{t_n=\pm }{}\underset{u_n=\pm }{}A_n[s_n]B_n[t_n]C_n[u_n]|s_n_n|t_n_n|u_n_n\right),$$
(10)
where the matrices are multiplied in ascending order keeping the order of $`A_nB_nC_n`$, and the basis states $`|s_n`$, $`|t_n`$, and $`|u_n`$ are eigenstates of the spin operator $`\sigma _z`$ as before. In this expression, the variables $`s`$, $`t`$, and $`u`$ are used instead of the index $`i`$ to denote the position of the links. The implementation of the matrix product ansatz means that a ladder lattice has been represented as a one-dimensional system with non-nearest neighbor interactions. Gauge invariance of matrix product states will be discussed in the next section.
If we require orthogonality of optimum basis states according to ref. , we have
$`{\displaystyle \underset{j=1}{\overset{M}{}}}{\displaystyle \underset{s=\pm }{}}(X_n[s])_{ij}(X_n[s])_{i^{}j}`$ $`=`$ $`\delta _{ii^{}},`$ (11)
$`{\displaystyle \underset{i=1}{\overset{M}{}}}{\displaystyle \underset{s=\pm }{}}(X_n[s])_{ij}(X_n[s])_{ij^{}}`$ $`=`$ $`\delta _{jj^{}},`$ (12)
where $`X`$ stands for $`A,B`$, and $`C`$. If these conditions are not imposed, norm of the matrix product state (10) may becomes very small, which results in numerical instability.
Energy
$$E=\frac{\mathrm{\Psi }|H|\mathrm{\Psi }}{\mathrm{\Psi }|\mathrm{\Psi }},$$
(13)
is a function of the matrices $`A_n[s]`$, $`B_n[t]`$, and $`C_n[u]`$. The numerator and denominator can be calculated by evaluating trace of a product of $`3L`$ matrices numerically:
$`\mathrm{\Psi }|H|\mathrm{\Psi }=`$ $`{\displaystyle \underset{n=1}{\overset{L}{}}}\mathrm{tr}((a_ny_nz_n+x_nb_nz_n+x_ny_nc_n){\displaystyle \underset{m=n+1modL}{\overset{n1+LmodL}{}}}w_m+`$ (14)
$`+\lambda \alpha _n\beta _n\gamma _n\alpha _{n+1modL}y_{n+1modL}z_{n+1modL}{\displaystyle \underset{m=n+2modL}{\overset{n1+LmodL}{}}}w_m),`$
$$\mathrm{\Psi }|\mathrm{\Psi }=\mathrm{tr}\left(\underset{m=1}{\overset{L}{}}w_m\right),$$
(15)
where
$`a_n{\displaystyle \underset{s,s^{}}{}}(\sigma _x)_{ss^{}}A_n^{}[s]A_n[s^{}],b_n{\displaystyle \underset{t,t^{}}{}}(\sigma _x)_{tt^{}}B_n^{}[t]B_n[t^{}],c_n{\displaystyle \underset{u,u^{}}{}}(\sigma _x)_{uu^{}}C_n^{}[u]C_n[u^{}],`$
$`\alpha _n{\displaystyle \underset{s,s^{}}{}}(\sigma _z)_{ss^{}}A_n^{}[s]A_n[s^{}],\beta _n{\displaystyle \underset{t,t^{}}{}}(\sigma _z)_{tt^{}}B_n^{}[t]B_n[t^{}],\gamma _n{\displaystyle \underset{u,u^{}}{}}(\sigma _z)_{uu^{}}C_n^{}[u]C_n[u^{}],`$
$`x_n{\displaystyle \underset{s}{}}A_n^{}[s]A_n[s],y_n{\displaystyle \underset{t}{}}B_n^{}[t]B_n[t],z_n{\displaystyle \underset{u}{}}C_n^{}[u]C_n[u],`$
$`w_nx_ny_nz_n.`$
The dimension of the matrices on the left hand side is $`M^2`$. By the outer product symbol $``$, we mean
$$a_{(i,k),(j,l)}=\underset{s,s^{}}{}\sigma _{ss^{}}A_{ij}[s]A_{kl}[s^{}].$$
The minimum of the energy function (13) corresponds to the ground state, which can be obtained based on matrix diagonalization as explained below. We can reduce the minimization problem (13) into a generalized eigenvalue problem
$$v^{}\overline{H}v=Ev^{}Nv,$$
(16)
where $`\overline{H}`$ and $`N`$ are $`2M^2`$ by $`2M^2`$ matrices. To understand what is going here, let us consider how energy can be minimized by varying $`A_n[s]`$ when other matrices are fixed. Note that equations (14) and (15) are bilinear of the matrix $`A_n[s]`$
$`\mathrm{\Psi }|H|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{i,j,k,l}{}}{\displaystyle \underset{s,t}{}}(A_n^{}[s])_{ij}\overline{H}_{(i,j,s),(k,l,t)}(A_n[t])_{kl},`$ (17)
$`\mathrm{\Psi }|\mathrm{\Psi }`$ $`=`$ $`{\displaystyle \underset{i,j,k,l}{}}{\displaystyle \underset{s}{}}(A_n^{}[s])_{ij}N_{(i,j,s),(k,l,t)}(A_n[t])_{kl},`$ (18)
where the matrix $`N`$ is diagonal for the indices $`s`$ and $`t`$. Once these expressions are obtained and the variational parameters $`A_n[s]`$ are regarded as a vector $`v`$, the minimization problem (13) reduces to (16).
There is one more trick that needs to be implemented. As explained in equations (11) and (12), we encounter numerical instability on the right hand side in eq. (16) if the above procedure is used as it is. This is because the matrix $`N`$ may have very small eigenvalues if the matrices $`A_n[s]`$, $`B_n[s]`$, and $`C_n[s]`$ are varied freely. For this reason, we need to impose one of the conditions (11) or (12) on the matrices.
A matrix can be decomposed into a product of three matrices, which is called singular value decomposition. For example, if we regard a tensor $`A[s]`$ as a matrix $`A_{i,(j,s)}=(A[s])_{ij}=A_{IJ}`$, singular value decomposition of $`A_{IJ}`$ is given by
$$A_{IJ}=\underset{K=1}{\overset{M}{}}U_{IK}D_KV_{KJ}.$$
(19)
where $`D_K`$ are the singular values of the matrix $`A`$ . The matrices $`U`$ and $`V`$ are orthonormal
$$\underset{I=1}{\overset{M}{}}U_{IK}^{}U_{IK^{}}=\delta _{KK^{}},\underset{J=1}{\overset{2M}{}}V_{KJ}^{}V_{K^{}J}=\delta _{KK^{}}.$$
(20)
The decomposition $`A=U^{}V`$, where $`U^{}=UD`$ is a square matrix, is the key of the trick. Consider a part of the matrix product wavefunction
$$C[u]A[s]=C^{}[u]A^{}[s],$$
(21)
where $`(A^{}[s])_{ij}=V_{i,(j,s)}`$ and $`C^{}[u]=C[u]U^{}`$. If the decomposition (19) is accurate, the both representations give the same result. Based on this trick, the eigenvalue problem (16) is solved successively for all the sets of the matrices starting from the right end in figure 2
$$C_L[u_L]B_L[t_L]A_L[s_L]\mathrm{}C_1[u_1]B_1[t_1]A_1[s_1],$$
(22)
which we call sweep. In one sweep process, eq. (16) is solved $`3L`$ times. If this sweep process is repeated several times, energy of low-lying states converges to some value. Calculation error can be controlled systematically by increasing the matrix dimension $`M`$. See Appendix A for generation of initial matrices.
## 4 Numerical results
The matrix product ansatz assumes large lattice. Our lattice size $`L=500`$ is sufficiently large. We solve the generalized eigenvalue problem (16) using LAPACK . For steady states, real matrices are sufficient for parameterizing the matrix product state (10). Convergence of energy needs to be checked for the number of sweeps and the matrix dimension $`M`$. Energy density $`E/L`$ converges in accuracy of five digits or higher after two sweeps when the matrix size $`M`$ is fixed.
Table 1 shows energy spectra of six low-lying states for three values of the coupling constant: $`\lambda =0.1,1`$, and $`10`$. The sweep process has been repeated twice. In this model, convergence of energy is very fast in contrast to Heisenberg chains . Small matrix dimension is sufficient for good convergence. Since we have obtained low-lying states without imposing the Gauss law on the variational space, gauge variant states are contained. In table 1, gauge invariant states are denoted with underlines. The other states are gauge variant. As we will see, gauge invariant physical states can be identified by calculating expectation values of the Gauss law operator.
In the ladder chain model, the Gauss law operator $`G(n)`$ is a product of three $`\sigma _z`$ operators (two horizontal and one vertical). We evaluate expectation values of $`G(n)`$ on the upper lattice sites shown in figure 2. Then, the number of the Gauss law operators to be evaluated is $`L`$. Expectation values on the lower sites are same as the upper ones because of reflection symmetry. Figures 3 plots expectation values of the Gauss law operator $`G(n)`$ in the case of $`\lambda =10`$ for the states (a) $`E_0`$, (b) $`E_1`$, (c) $`E_2`$, and (d) $`E_3`$. In figures 3 (a) and (d), the Gauss law $`G(n)=1`$ is satisfied uniformly on every lattice sites. Therefore, the obtained states $`E_0`$ and $`E_3`$ are gauge invariant. On the other hand, in figures 3 (b) and (c), the states $`E_1`$ and $`E_2`$ are gauge variant because gauge symmetry is definitely broken at the site $`n=500`$. The position of this special lattice site depends on where the sweep process ends. The relation $`G(n)=1`$ or $`1`$ holds for the obtained low-lying states in accuracy of seven digits or higher when $`M=4`$.
According to the Elizur’s theorem, gauge variant operators have vanishing expectation values. We have checked that expectation values of single spin operators $`\sigma _z(n)`$ vanish for the gauge invariant states in accuracy of ten digits or higher when $`M=4`$. We also have checked that the gauge variant states have vanishing expectation values of $`\sigma _z(n)`$ in the same accuracy. These statements apply to the low-lying states shown in table 1.
In this way, we can classify the obtained states into gauge invariant states and others. Figure 4 plots vacuum energy density as a function of the coupling constant $`\lambda `$.
In order to see how spatial distribution look like, we evaluate expectation values of the plaquette operator for the vacuum and the excited states. Figure 5 plots expectation values of the plaquette as functions of the spatial lattice coordinate $`n`$ for the gauge invariant low-lying states (a) $`E_0`$, (b) $`E_3`$, and (c) $`E_5`$ in the case of $`\lambda =10`$. As expected, the vacuum state $`E_0`$ has complete uniformity. On the other hand, the excited states $`E_3`$ and $`E_5`$ have lumps around $`n=500`$. The higher excited states have similar lumps around the boundary. The obtained solutions have periodicity because the plaquette distributions are continuous around the boundary. The position of the lumps can be moved without changing energy in the sweep process as explained in figure 3. Figure 6 plots vacuum expectation values of the plaquette operator as a function of the coupling constant $`\lambda `$. For small and large $`\lambda `$, the tendency of energy and plaquette is consistent with the exact values given in section 3.
## 5 Summary
We have extracted gauge invariant physical states in a $`Z_2`$ hamiltonian lattice gauge theory on a spatial ladder chain. The calculations are based on the matrix product ansatz, which gives sufficiently convergent energy and wavefunction for low-lying states. In the future studies, similar calculations should be tested in higher dimensional lattice gauge theory including supersymmetric cases . The proposed method will be useful especially for non-perturbative analysis of vacuum structure at the amplitude level.
###### Acknowledgments.
The author would like to thank T. Nishino for useful communications. The numerical calculations were carried on the RIKEN RSCC system. This work has been partially supported by RIKEN BNL.
## Appendix A Generation of initial matrices
Before starting sweep process, we need to prepare initial values of matrices for the variational state (10). As explained in section 3, the matrices need to satisfy one of the orthogonality conditions (11) or (12) for numerical stability. The algorithm shown below is also useful for minimization of energy function with the Powell method .
Let us consider real matrices $`A[s]`$ that satisfy the following normalization condition
$$\underset{j=1}{\overset{M}{}}\underset{s=1}{\overset{K}{}}A[s]_{ij}A[s]_{i^{}j}=\delta _{ii^{}},$$
(23)
where $`M`$ is the dimension of the matrices $`A[s]`$ and $`K`$ is the degrees of freedom of each site or link. To parametrize the matrices $`A[s]`$, we introduce $`KM`$-dimensional vectors $`v^{(n)}`$
$$v^{(n)}=(v_1^{(n)},\mathrm{},v_{KMn+1}^{(n)},\underset{n1}{\underset{}{0,\mathrm{},0}}),n=1,\mathrm{},M.$$
(24)
These vectors are linearly independent and can be orthonormalized using the Gram-Schmidt method
$`a^{(1)}`$ $`=`$ $`{\displaystyle \frac{v^{(1)}}{|v^{(1)}|}}`$ (25)
$`b^{(k+1)}`$ $`=`$ $`v^{(k+1)}{\displaystyle \underset{n=1}{\overset{k}{}}}v^{(k+1)},a^{(n)}a^{(n)}`$ (26)
$`a^{(k+1)}`$ $`=`$ $`{\displaystyle \frac{b^{(k+1)}}{|b_{(k+1)}|}}`$ (27)
where $`k=1,\mathrm{},M1`$ and the brackets represent inner product. The orthonormalized vectors $`a^{(n)}`$ are used to parametrize the matrices $`A[s]_{ij}=A_{i,(j,s)}`$
$$A_{i,(j,s)}=\left(\begin{array}{c}a^{(1)}\\ \mathrm{}\\ a^{(M)}\end{array}\right),$$
(28)
which satisfy the conditions (23). When eq. (23) is satisfied, the number of independent degrees of freedom associated with the matrices $`A[s]`$ is given by
$$KM^2\frac{M(M+1)}{2},$$
(29)
which is equal to that of the vectors $`v^{(n)}`$
$$\underset{n=0}{\overset{M1}{}}(KMn)M.$$
(30)
In eq. (30), the first term counts the number of the parameters $`v_i^{(n)}`$ and the second term comes from the normalization conditions (25) and (27).
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# Lorentz contraction, geometry and range in antiproton-proton annihilation into two pions
## 1 Quark Model and intrinsic meson wavefunctions
In a quark model description of the $`\overline{p}p`$ annihilation process, one usually seeks guidance from Feynman diagrams in order to deduce appropriate transition operators (i.e. the operators linking the initial $`\overline{p}p`$ quarks and antiquarks so they form the right $`\overline{q}q`$ pairs). These so-called quark-line diagrams (QLD) are classified according to their flavor-flux topology into rearrangement or annihilation diagrams. In the former case a $`\overline{q}q`$ pair is annihilated and a quark and antiquark are rearranged to produce the final mesons, while in the latter case two $`\overline{q}q`$ pairs are annihilated and an $`\overline{q}q`$ state is created in the final state. The $`\overline{q}q`$ pairs annihilate into states with distinct quantum numbers $`J^P`$ and $`I`$. In a way, one can conceive of the model as an expansion of the annihilation amplitude in terms of increasing total angular momentum $`J^P`$. Parity requires the $`\overline{q}q`$ pairs be annihilated/created in an $`S=1`$ state. Thus the spin-multiplicity is fixed and for $`J^P=0^+,1^{}\mathrm{}`$ one gets $`{}_{}{}^{3}P_{0}^{}`$, $`{}_{}{}^{3}S_{1}^{}`$… annihilation operators. For a concise review of the QLD see, for instance, Ref. dover . The resulting transition operators may be written in a Hamiltonian form for both the $`{}_{}{}^{3}P_{0}^{}`$ and $`{}_{}{}^{3}S_{1}^{}`$ cases as
$`(^3P_0)`$ $`=`$ $`\gamma {\displaystyle \underset{ijmn}{}}a_m^{}(𝐤^{})a_n(𝐤)a_i(𝐩)b_j(𝐩^{})𝝈(𝐩𝐩^{})(2\pi )^3\delta (𝐤^{}𝐤𝐩^{}𝐩)+\mathrm{h}.\mathrm{c}.`$ (1)
$`(^3S_1)`$ $`=`$ $`\kappa {\displaystyle \underset{ijmn}{}}a_m^{}(𝐤^{})a_n(𝐤)a_i(𝐩)b_j(𝐩^{})[2𝝈_{ij}𝐤+i(𝝈_{mn}\times 𝝈_{ij})(𝐤𝐤^{})]\times `$ (2)
$`\times (2\pi )^3\delta (𝐤^{}𝐤𝐩^{}𝐩)+\mathrm{h}.\mathrm{c}.,`$
where color indices have been suppressed and $`ijmn`$ sums over flavor states. The delta-function imposes momentum transfer from the annihilated $`\overline{q}q`$ pair to one of the remaining (anti)quarks. Additionally, color is exchanged in the $`{}_{}{}^{3}S_{1}^{}`$ mechanism since $`J^P=1^{}`$ is the quantum number of the gluon. The relative strength $`\lambda =|\gamma /\kappa |`$ is a fit parameter.
Confinement is simulated in a harmonic oscillator basis. The single (anti)quark wavefunctions are therefore bundled in Gaussian wave packets in which the c.m. motion of the (anti)proton and pions can be separated from the relative motion of their respective quarks and antiquarks. The total wavefunction of the pion, for example, writes in momentum space as
$`\varphi _\pi (𝐩_i,𝐩_j)`$ $`=`$ $`N_\pi \left({\displaystyle \frac{4\pi }{\beta }}\right)^{3/2}u(𝐩_i)v(𝐩_j)\chi _{\mathrm{iso}}\chi _{\mathrm{color}}\mathrm{exp}\left\{{\displaystyle \frac{1}{2\beta }}{\displaystyle \underset{\mu =i,j}{}}\left[𝐩_\mu \frac{1}{2}𝐐_\pi \right]^2\right\}.`$ (3)
Here, $`𝐩_i`$ and $`𝐩_j`$ are the $`\overline{q}q`$ momenta and $`𝐐_\pi =𝐩_i+𝐩_j`$ is the pion momentum in the c.m. frame whereas $`u(𝐩_i)`$ and $`v(𝐩_j)`$ are the usual quark and antiquark Dirac spinors. The size parameter $`\beta `$ is chosen so as to give the pion a radius $`r_\pi =0.48`$ fm.
## 2 Lorentz Effects in Final States
So far, we have dealt with the wave representation of the pion in its rest frame where it is spherical (meson in an $`s`$-wave). In $`\overline{p}p\pi ^{}\pi ^+`$, on the other hand, the pions are produced at considerable c.m. energies already at very low $`\overline{p}`$-beam energies ( $`\sqrt{s}=1.92.2`$ GeV in Ref. hasan1 ). The large mass difference between the pion and the proton causes the final state to be highly relativistic, as can be simply seen from the relativistic factor $`\gamma =\sqrt{s}/2m_\pi c^26.88.0`$. In computing the annihilation amplitudes (and observables), one puts oneself in the c.m. frame in order to compare with experiment. Therefore, all ingredients must be transformed to the c.m. For the pions, this means they are not spherical anymore, since they are boosted from their rest frame to the c.m. frame. Eq. (3) cannot be used as such and instead the general Lorentz transformation $`𝐩_i^{}l^1𝐩_i`$ (where $`𝐩_i`$ is a quark momentum in the pion rest frame) yields
$`\varphi _\pi (l^1𝐩_i,l^1𝐩_j)`$ $`=`$ $`{\displaystyle \frac{N_\pi }{\sqrt{\gamma }}}\left({\displaystyle \frac{4\pi }{\beta }}\right)^{3/2}d_{\lambda \lambda ^{}}^{1/2}(\theta )u(l^1𝐩_i,\lambda )d_{\lambda \lambda ^{}}^{1/2}(\theta )v(l^1𝐩_j,\lambda )\chi _{\mathrm{iso}}\chi _{\mathrm{color}}\times `$ (4)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{4\beta }}\left[(𝐩_i𝐩_j)^2+({\displaystyle \frac{1}{\gamma ^2}}1)\left([𝐩_i𝐩_j]\widehat{𝐐}_\pi \right)^2\right]\right\}.`$
In the exponent, the momentum components along the boost direction $`\widehat{𝐐}_\pi =𝐐_\pi /|𝐐_\pi |`$ have been separated from the perpendicular ones. Naturally, for $`\gamma =1`$ one retrieves the original wavefunction in Eq. (3). The additional factor $`1/\sqrt{\gamma }`$ is due to proper normalization to one of the wavefunction. The Lorentz transformations are effected on the Dirac spinors via the Wick rotations $`d_{\lambda \lambda ^{}}^{1/2}(\theta )`$ in spin space, where the Wick angle $`\theta `$ depends on the relativistic factor $`\beta =v/c`$ rather than on $`\gamma `$. The indices $`\lambda `$ and $`\lambda ^{}`$ are the (anti)quark helicities in the pion rest frame and in the c.m. frame, respectively. We have applied the c.m. equal-time condition $`t_i^{}=t_j^{}`$ for the time components of the quark and antiquark, which is also frequently used in Bethe-Salpeter approaches. The effect of the time components is currently under investigation. While the boosting of spinors gives rise to subtle interference effects in the annihilation amplitudes (see discussion in Ref. paper2 ), its magnitude is negligible compared with the $`\gamma `$-components. Once the wavefunctions of Eq. (4) are employed to derive the annihilation amplitudes, it becomes clear that the boost effects are dramatic. In summary, the annihilation amplitudes, being of Gaussian form, acquire new relativistic terms. Their magnitude can be, depending on the $`\gamma `$ (and hence $`\sqrt{s}`$) value, up to two orders of magnitude larger than the one of non-relativistic terms. Furthermore, this strong $`\gamma `$-dependence is not constant. The boosts also introduce an angular dependence by means of the projection on the boost direction $`\widehat{𝐐}_\pi `$ in the exponential. It is this dependence (which at the hadronic level equates to the c.m. angle between the antiproton beam and an outgoing pion) that is novel and crucial to any improvement in a fit to the LEAR data.
## 3 Annihilation and Geometry
In the following, we propose a somewhat different geometric interpretation. The strong modification of both the $`{}_{}{}^{3}P_{0}^{}`$ and $`{}_{}{}^{3}S_{1}^{}`$ annihilation amplitudes is entirely due to the Lorentz contractions in the pions. In the overlap integrals this implies a narrowed overlap of the initial $`\overline{p}p`$ pair and the final pion wave functions. At first sight, this looks discouraging since our introductory discussion revolved around how to increase the annihilation range in a quark model. Yet, this decrease in overlap is merely one piece of the puzzle. The other part is illustrated in Fig. 1 where tentatively the angle dependence of the amplitudes due to relativistic effects is shown for two cases. The smaller blobs represent the pions and are round for $`\gamma =1`$ and oval (Lorentz contracted) otherwise. The distances between the two pions and the proton and antiproton are the same in each picture. However, in the first case the pions are produced perpendicularly to the $`\overline{p}`$ beam direction and all particles overlap. They do so somewhat less when the pions are contracted. In the second case, where the pions are emitted at a different angle, the four spheres overlap if relativity is neglected while the two pion-ellipsoids overlap much less with the the $`\overline{p}p`$ spheres if relativity is included. This angle-dependent overlap is responsible for a richer angular dependence of the annihilation amplitudes. We therefore obtain significant contributions to $`J1`$ partial waves not present previously and a fit to the LEAR data hasan1 is quite successful paper3 . On the other hand, it is clear that one cannot define an annihilation range per se since this range depends on the pion direction with respect to the initial $`\overline{p}`$-beam. Even so, one can conclude that in order for the $`\overline{p}p`$ pair to annihilate, the quarks and antiquarks must have a considerable overlap whose size is less than two nucleon radii.
We appreciated helpful discussions with Mary Alberg, Thomas Gutsche, Benoît Loiseau, Johann Haidenbauer, Fred Myhrer and Slawomir Wycech.
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# Instabilities of multi-hump vector solitons in coupled nonlinear Schrödinger equations
## 1 Introduction
The coupled NLS equations have wide applications in the modeling of physical processes. For instance, such equations with the cubic nonlinearity govern the nonlinear interaction of two wave packets and optical pulse propagation in birefringent fibers or wavelength-division-multiplexed optical systems . Similar equations with the saturable nonlinearity describe the propagation of several mutually-incoherent laser beams in biased photorefractive crystals . Various types of vector solitons including single-hump and multi-hump ones have been known to exist in these coupled NLS equations , and they have been observed in photorefractive crystals as well .
Linear stability of vector solitons in the coupled NLS equations is an important issue. Fundamental single-hump vector solitons are known to be stable . Stability of multi-hump vector solitons (which have one or more nodal points in one or more components) is more subtle. For the cubic nonlinearity, it was conjectured in based on the numerical evidence that multi-hump vector solitons were all linearly unstable. If the multi-hump solitons are pieced together by a few fundamental vector solitons, then their linear instability has been proven both analytically and numerically in . The linear instability for other types of multi-hump vector solitons has not been proven yet. For the saturable nonlinearity, multi-hump solitons have been shown to be stable in certain parameter regions , but the origins of their stability and instability have not yet been fully analyzed.
¿From a broader point of view, the theory of linear stability of vector solitons in coupled NLS equations was recently developed with the use of the closure theorem for the negative index of the linearized Hamiltonian . However, there are not many applications of the general theory to particular bifurcations of unstable eigenvalues , because the general theory excludes non-generic bifurcations. It is desirable to further develop a perturbation theory to the eigenvalue bifurcations, so that the origin of instability becomes more apparent in the context of the closure theorem.
In this paper, we investigate the linear stability of multi-hump vector solitons in the general Hamiltonian system of coupled NLS equations both analytically and numerically. Using the closure theorem for the negative index of the linearized Hamiltonian as well as the perturbation technique, we classify all possible bifurcations of unstable eigenvalues in two physical models with cubic or saturable nonlinearities. In the first model, we show that multi-hump vector solitons near points of local bifurcations are always linearly unstable, in agreement with numerical results in . In the second model, the situation is more complicated. Our results show that the 1st family of multi-hump vector solitons is indeed linearly stable near the local-bifurcation boundary, in agreement with numerical results in . However, for the 2nd family, we discovered a new oscillatory instability near the local-bifurcation boundary, which was missed in . Due to this oscillatory instability, the stability region of vector solitons for the 2nd family is drastically reduced from that reported in . Numerically, we track the unstable eigenvalues of multi-hump solitons and reveal various scenarios of eigenvalue bifurcations away from the local-bifurcation boundaries. We also map out the correct stability regions of multi-hump vector solitons in the entire parameter space. Furthermore, the number of numerically-obtained unstable eigenvalues agrees completely with that predicted by the negative index of the linearized Hamiltonian.
Our paper is structured as follows. The main formalism and the closure theorem for the negative index of the linearized Hamiltonian are described in Section 2. Analysis of unstable eigenvalues in the coupled NLS equations with cubic and saturable nonlinearities is developed in Sections 3 and 4, respectively. Section 5 summarizes our results and open questions. Appendix A reviews bifurcations of unstable eigenvalues by the perturbation method.
## 2 Main formalism
We consider a general Hamiltonian system of coupled NLS equations in the form:
$$i\frac{\psi _n}{z}+d_n\frac{^2\psi _n}{x^2}+\frac{U}{|\psi _n|^2}\psi _n=0,n=1,\mathrm{},N,$$
(2.1)
where $`z_+`$, $`x`$, $`\psi _n`$, $`d_n`$, and $`U=U(|\psi _1|^2,\mathrm{},|\psi _N|^2)`$. We assume that $`U(0)=U^{}(0)=0`$, and $`d_n>0`$ for all $`n`$. In optical fibers (photorefractive crystals), the function $`\psi _n(z,x)`$ is the envelope amplitude of the $`n^{\mathrm{th}}`$ channel (beam), $`z`$ is the propagation distance along the fiber (waveguide), and $`x`$ is the retarded time (the transverse coordinate) .
Following the recent work in , we study the linear stability of vector solitons:
$$\psi _n(z,x)=\mathrm{\Phi }_n(x)e^{i\beta _nz},$$
(2.2)
where $`\mathrm{\Phi }_n:`$, and $`\beta _n>0`$ for all $`n`$. We assume that none of the components $`\mathrm{\Phi }_n(x)`$ vanish identically on $`x`$. Linearization of the coupled NLS equations (2.1) follows from the expansion:
$$\psi _n(z,x)=\left\{\mathrm{\Phi }_n(x)+\left[u_n(x)+iw_n(x)\right]e^{\lambda z}+\left[\overline{u}_n(x)+i\overline{w}_n(x)\right]e^{\overline{\lambda }z}\right\}e^{i\beta _nz},$$
(2.3)
where $`u_n,w_n1`$, and the overline denotes the complex conjugation. The linearized equations for $`(u_n,w_n)`$ are the following non-self-adjoint problem in $`L^2(,^{2N})`$:
$$_1𝐮=\lambda 𝐰,_0𝐰=\lambda 𝐮,$$
(2.4)
where $`\lambda `$ is an eigenvalue, $`(𝐮,𝐰)^T:^{2N}`$ is the eigenvector, and $`_0`$ and $`_1`$ are the matrix Schrödinger operators with elements:
$`(_0)_{n,m}`$ $`=`$ $`\left(d_n{\displaystyle \frac{d^2}{dx^2}}+\beta _n{\displaystyle \frac{U}{\mathrm{\Phi }_n^2}}\right)\delta _{n,m},`$ (2.5)
$`(_1)_{n,m}`$ $`=`$ $`\left(d_n{\displaystyle \frac{d^2}{dx^2}}+\beta _n{\displaystyle \frac{U}{\mathrm{\Phi }_n^2}}\right)\delta _{n,m}2{\displaystyle \frac{^2U}{\mathrm{\Phi }_n^2\mathrm{\Phi }_m^2}}\mathrm{\Phi }_n\mathrm{\Phi }_m.`$ (2.6)
Since the eigenvalue problem (2.4) is a linearization of the Hamiltonian system, the values of $`\lambda `$ occur as pairs of real or purely imaginary eigenvalues, or as quadruplets of complex eigenvalues. Eigenvalues with $`\mathrm{Re}(\lambda )>0`$ lead to spectral instability of vector solitons (2.2). We denote the number of eigenvalues in the first open quadrant as $`N_{\mathrm{comp}}`$, the number of positive real eigenvalues as $`N_{\mathrm{real}}`$, and the number of purely imaginary eigenvalues with positive $`\mathrm{Im}(\lambda )`$ as $`N_{\mathrm{imag}}`$. The continuous spectrum has $`N`$ branches, located at the positive imaginary axis for $`\mathrm{Im}(\lambda )\beta _n`$, $`n=1,\mathrm{},N`$. Zero eigenvalue $`\lambda =0`$ has geometric multiplicity of at least $`(N+1)`$ and algebraic multiplicity of at least $`(2N+2)`$, in the assumption that none of the components $`\mathrm{\Phi }_n(x)`$ vanishes identically on $`x`$ .
Furthermore, we denote the number of negative and zero eigenvalues of operators $`_{0,1}`$ in $`L^2(,^N)`$ as $`n(_{0,1})`$ and $`z(_{0,1})`$, respectively. We also assume that the solution $`\mathrm{\Phi }_n(x)`$ depends smoothly on $`(\beta _1,\mathrm{},\beta _N)`$ in an open non-empty set of $`^N`$ and introduce the Hessian matrix $`𝒰`$ with elements:
$$𝒰_{n,m}=\frac{Q_n}{\beta _m},$$
(2.7)
where $`Q_n=Q_n(\beta _1,\mathrm{},\beta _N)=_{}\mathrm{\Phi }_n^2𝑑x`$. We denote the number of positive and zero eigenvalues of matrix $`𝒰`$ as $`p(𝒰)`$ and $`z(𝒰)`$, respectively. Finally, we introduce the linearized Hamiltonian (”energy”) of the eigenvalues $`\lambda `$ in $`H^1(,^{2N})`$:
$$h[𝐮,𝐰]=𝐮,_1𝐮+𝐰,_0𝐰,$$
(2.8)
where $`,`$ is the standard inner product in $`L^2(,^{2N})`$. The negative index of the linearized Hamiltonian is the number of negative eigenvalues of $`_1`$ and $`_0`$ in $`L^2(,^N)`$.
Several assumptions are imposed on the linearized problem (2.4) in a general case :
(i) $`z(_1)=1`$, $`z(_0)=N`$;
(ii) $`z(𝒰)=0`$;
(iii) no eigenvalues $`\lambda i`$ exist with $`h[𝐮,𝐰]=0`$;
(iv) no embedded eigenvalues $`\lambda i`$ exist with $`|\mathrm{Im}(\lambda )|\beta _{\mathrm{min}}`$, where $`\beta _{\mathrm{min}}=\mathrm{min}(\beta _1,\mathrm{},\beta _N).`$
Closure Theorem Assume that (i)–(iv) be satisfied. Let $`N_{\mathrm{imag}}^{}`$ be the number of eigenvalues $`\lambda i_+`$ with $`h[𝐮,𝐰]<0`$. Then
$`(i)N_{\mathrm{real}}+2N_{\mathrm{comp}}+2N_{\mathrm{imag}}^{}=n(_1)p(𝒰)+n(_0),`$ (2.9)
$`(ii)N_{\mathrm{real}}|n(_1)p(𝒰)n(_0)|,`$ (2.10)
$`(iii)N_{\mathrm{comp}}\mathrm{min}(n(_0),n(_1)p(𝒰)),`$ (2.11)
such that
$$|n(_1)p(𝒰)n(_0)|N_{\mathrm{unst}}n(_1)p(𝒰)+n(_0),$$
(2.12)
where $`N_{\mathrm{unst}}=N_{\mathrm{real}}+2N_{\mathrm{comp}}`$ is the total number of unstable eigenvalues in the problem (2.4).
This theorem was originally proved for the coupled NLS equations (2.1) in one dimension and then generalized to a three-dimensional NLS equation and to an abstract Hamiltonian dynamical system . It allows us to analytically trace unstable eigenvalues under parameter continuations, starting with the particular limits, where all eigenvalues $`\lambda `$ of negative energy $`h[𝐮,𝐰]`$ are known. Examples of such parameter continuation are recently reported in in the context of the coupled NLS equations.
Bifurcations of unstable eigenvalues may occur in the linearized problem (2.4), when operators $`_1`$ and $`_0`$ change according to a continuous deformation and one of the assumptions (i)–(iv) of the Closure Theorem fails. Bifurcations are reviewed in Appendix A. In what follows, we apply parameter continuation and bifurcation analysis to the system of coupled NLS equations (2.1) with cubic and saturable nonlinearities.
## 3 The coupled cubic NLS equations
We consider the system of coupled cubic NLS equations :
$`i\psi _{1z}+\psi _{1xx}+\left(|\psi _1|^2+\chi |\psi _2|^2\right)\psi _1`$ $`=`$ $`0,`$
$`i\psi _{2z}+\psi _{2xx}+\left(\chi |\psi _1|^2+|\psi _2|^2\right)\psi _2`$ $`=`$ $`0,`$ (3.1)
where $`\chi >0`$. The system is a particular example of (2.1) with $`N=2`$, $`d_1=d_2=1`$, and
$$U=\frac{1}{2}|\psi _1|^4+\chi |\psi _1|^2|\psi _2|^2+\frac{1}{2}|\psi _2|^4.$$
(3.2)
The system (3.1) has a countable infinite set of families of vector solitons $`𝚽(x)=(\mathrm{\Phi }_1,\mathrm{\Phi }_2)^T`$, classified by different nodal index $`𝐢=(i_1,i_2)^T`$, where $`i_n`$ is the number of zeros of $`\mathrm{\Phi }_n(x)`$ on $`x`$ . We consider here families of vector solitons with nodal index $`𝐢=(0,n)^T`$, $`n`$, which are locally close to the NLS soliton, $`𝚽_{\mathrm{NLS}}(x)=(\mathrm{\Phi }^{(0)},0)^T`$, where $`\mathrm{\Phi }^{(0)}(x)=\sqrt{2\beta _1}\mathrm{sech}(\sqrt{\beta _1}x)`$. We let $`\beta _1=1`$ and $`\beta _2=\beta `$ for convenience and introduce scalar Schrödinger operators:
$`L_0`$ $`=`$ $`{\displaystyle \frac{d^2}{dx^2}}+12\mathrm{sech}^2(x),`$ (3.3)
$`L_1`$ $`=`$ $`{\displaystyle \frac{d^2}{dx^2}}+16\mathrm{sech}^2(x),`$ (3.4)
$`L_s`$ $`=`$ $`{\displaystyle \frac{d^2}{dx^2}}+\lambda _n^2(\chi )2\chi \mathrm{sech}^2(x),`$ (3.5)
where
$$\lambda _n(\chi )=\frac{\sqrt{1+8\chi }(2n+1)}{2}.$$
(3.6)
The scalar operators $`L_0`$,$`L_1`$,$`L_s`$ define the matrix operators $`_0`$ and $`_1`$ at $`ϵ=0`$:
$$_0=\mathrm{diag}(L_0,L_s),_1=\mathrm{diag}(L_1,L_s).$$
We define the perturbation series expansions of vector solitons $`𝚽=(\mathrm{\Phi }_1,\mathrm{\Phi }_2)^T`$:
$`\mathrm{\Phi }_1(x)`$ $`=`$ $`\mathrm{\Phi }^{(0)}(x)+ϵ^2\mathrm{\Phi }^{(2)}(x)+\mathrm{O}(ϵ^4),`$
$`\mathrm{\Phi }_2(x)`$ $`=`$ $`ϵ\mathrm{\Phi }^{(1)}(x)+ϵ^3\mathrm{\Phi }^{(3)}(x)+\mathrm{O}(ϵ^5),`$ (3.7)
and
$$\beta =\lambda _n^2(\chi )+ϵ^2C_n(\chi )+\mathrm{O}(ϵ^4).$$
(3.8)
Corrections of the perturbation series (3.7)–(3.8) satisfy the linear equations:
$`L_s\mathrm{\Phi }^{(1)}`$ $`=`$ $`0,`$ (3.9)
$`L_1\mathrm{\Phi }^{(2)}`$ $`=`$ $`\chi \mathrm{\Phi }^{(0)}\left(\mathrm{\Phi }^{(1)}\right)^2,`$ (3.10)
$`L_s\mathrm{\Phi }^{(3)}`$ $`=`$ $`C_n(\chi )\mathrm{\Phi }^{(1)}+2\chi \mathrm{\Phi }^{(0)}\mathrm{\Phi }^{(1)}\mathrm{\Phi }^{(2)}+\left(\mathrm{\Phi }^{(1)}\right)^3.`$ (3.11)
The problem (3.9) has a decaying solution $`\mathrm{\Phi }^{(1)}\mathrm{\Phi }_n^{(1)}(x)`$ (see ). When $`n=0`$ and $`\chi >0`$, the solution $`\mathrm{\Phi }_0^{(1)}=\mathrm{sech}^s(x)`$, $`s=\lambda _0(\chi )`$ is a ground state. When $`n>0`$ and $`\chi >\chi _n=n(n+1)/2`$, the solution $`\mathrm{\Phi }_n^{(1)}(x)`$ is an excited state with exactly $`n`$ nodes on $`x`$. The problem (3.10) also has a decaying solution $`\mathrm{\Phi }^{(2)}(x)`$, since the right-hand-side $`\chi \mathrm{\Phi }^{(0)}\left(\mathrm{\Phi }_n^{(1)}\right)^2`$ is orthogonal to the kernel of the operator $`_1`$, which is $`\mathrm{\Phi }^{(0)}(x)`$. By the Fredholm Alternative Theorem, the problem (3.11) has a decaying solution if and only if the right-hand-side is orthogonal to the kernel of $`_s`$, which is $`\mathrm{\Phi }_n^{(1)}(x)`$. The orthogonality condition defines the parameter $`C_n(\chi )`$ in the form:
$$C_n(\chi )=\frac{\left(\mathrm{\Phi }_n^{(1)}\right)^2,\left(2\chi \mathrm{\Phi }^{(0)}\mathrm{\Phi }^{(2)}+\left(\mathrm{\Phi }_n^{(1)}\right)^2\right)}{\mathrm{\Phi }_n^{(1)},\mathrm{\Phi }_n^{(1)}}.$$
(3.12)
The condition $`C_n(\chi )0`$ gives the sufficient condition of continuation of the perturbation series expansions (3.7)–(3.8). Thus, for $`\chi >\chi _n`$ and $`C_n(\chi )0`$, there exists some $`R_n>0`$, such that the $`n`$-th family of vector solitons $`𝚽(x)=(\mathrm{\Phi }_1,\mathrm{\Phi }_2)^T`$ with the nodal index $`𝐢=(0,n)^T`$ bifurcates from $`𝚽_{\mathrm{NLS}}=(\mathrm{\Phi }^{(0)},0)^T`$ in the one-sided domain $`_n`$:
$$_n=\{\beta :0<\left|\beta \lambda _n^2(\chi )\right|<R_n,\mathrm{sign}\left(\beta \lambda _n^2(\chi )\right)=\mathrm{sign}(C_n(\chi ))\}.$$
(3.13)
These results for the first three families $`n=0,1,2`$ were analytically obtained and numerically verified in . We investigate stability of the $`n`$-th family of vector solitons in the one-sided domain $`_n`$ below.
### 3.1 Analytical results
We trace unstable eigenvalues using the Closure Theorem. We consider a generic case $`C_n(\chi )0`$ in the one-sided open domain $`_n`$ and show that the left-hand and right-hand sides of the closure relation (2.9) are equal to $`2n`$ for small $`ϵ0`$.
Operator $`L_0`$ in (3.3) has one bound state for zero eigenvalue, operator $`L_1`$ in (3.4) has two bound states for negative and zero eigenvalues, and operator $`L_s`$ in (3.5) has $`(n+1)`$ bound states with $`n`$ negative and one zero eigenvalues. Therefore, at $`ϵ=0`$, we have $`n(_0)=0+n=n`$, $`z(_0)=1+1=2`$, $`n(_1)=1+n`$ and $`z(_1)=1+1=2`$. It follows from Sturm Nodal Theorem that
$$n(_0)=n,z(_0)=2,ϵ0.$$
Since $`z(_1)=2>1`$, we have the bifurcation case $`z(_1)>1`$ for $`ϵ=0`$ (see Appendix A.1). It is however a degenerate bifurcation case, since it occurs on the boundary of the existence domain $`_n`$, such that $`\beta _n`$. We trace the zero eigenvalue of $`_1`$ for $`ϵ0`$ by the regular perturbation series,
$`𝐮(x)=\left[\begin{array}{c}0\\ \mathrm{\Phi }_n^{(1)}(x)\end{array}\right]+ϵ\left[\begin{array}{c}u^{(1)}(x)\\ 0\end{array}\right]+ϵ^2\left[\begin{array}{c}0\\ u^{(2)}(x)\end{array}\right]+\mathrm{O}(ϵ^3)`$ (3.20)
and
$`\lambda =ϵ^2\lambda _2+\mathrm{O}(ϵ^4).`$ (3.21)
Corrections of the perturbation series (3.20) satisfy a set of linear non-homogeneous equations:
$`L_1u^{(1)}`$ $`=`$ $`2\chi \mathrm{\Phi }^{(0)}\left(\mathrm{\Phi }_n^{(1)}\right)^2,`$ (3.22)
$`L_su^{(2)}`$ $`=`$ $`(\lambda _2C_n(\chi ))\mathrm{\Phi }_n^{(1)}+2\chi \mathrm{\Phi }^{(0)}\mathrm{\Phi }_n^{(1)}\left(u^{(1)}+\mathrm{\Phi }^{(2)}\right)+3\left(\mathrm{\Phi }_n^{(1)}\right)^3.`$ (3.23)
It follows from (3.10) and (3.22) that $`u^{(1)}=2\mathrm{\Phi }^{(2)}`$. By the Fredholm Alternative Theorem, decaying solutions of (3.23) exist if and only if the right-hand-side of (3.23) is orthogonal to $`\mathrm{\Phi }_n^{(1)}(x)`$. Using (3.12), we find that $`\lambda _2=2C_n(\chi )`$. Therefore, we have:
$$n(_1)=1+\mathrm{\Theta }(C_n(\chi ))+n,z(_1)=1,ϵ>0,$$
where $`\mathrm{\Theta }(z)`$ is the Heaviside step-function. We trace the zero eigenvalue of $`𝒰`$ from (2.7) and (3.7):
$`𝒰_{1,1}`$ $`=`$ $`{\displaystyle \frac{Q_1}{\beta _1}}|_{\beta _1=1}=2+2\mathrm{\Phi }^{(0)},\mathrm{\Phi }^{(2)}{\displaystyle \frac{ϵ^2}{\beta _1}}|_{\beta _1=1}+\mathrm{O}(ϵ^2),`$
$`𝒰_{1,2}`$ $`=`$ $`{\displaystyle \frac{Q_1}{\beta _2}}|_{\beta _1=1}=2\mathrm{\Phi }^{(0)},\mathrm{\Phi }^{(2)}{\displaystyle \frac{ϵ^2}{\beta _2}}|_{\beta _1=1}+\mathrm{O}(ϵ^2),`$
$`𝒰_{2,1}`$ $`=`$ $`{\displaystyle \frac{Q_2}{\beta _1}}|_{\beta _1=1}=\mathrm{\Phi }_n^{(1)},\mathrm{\Phi }_n^{(1)}{\displaystyle \frac{ϵ^2}{\beta _1}}|_{\beta _1=1}+\mathrm{O}(ϵ^2),`$
$`𝒰_{2,2}`$ $`=`$ $`{\displaystyle \frac{Q_2}{\beta _2}}|_{\beta _1=1}=\mathrm{\Phi }_n^{(1)},\mathrm{\Phi }_n^{(1)}{\displaystyle \frac{ϵ^2}{\beta _2}}|_{\beta _1=1}+\mathrm{O}(ϵ^2).`$
It follows from (3.8) that
$$\frac{ϵ^2}{\beta _2}|_{\beta _1=1}=\frac{1}{C_n(\chi )}+\mathrm{O}(ϵ^2)$$
(3.24)
and, due to the symmetry of $`𝒰`$,
$$det(𝒰)=\frac{2\mathrm{\Phi }_n^{(1)},\mathrm{\Phi }_n^{(1)}}{C_n(\chi )}+\mathrm{O}(ϵ^2).$$
(3.25)
Therefore, we have:
$$p(𝒰)=1+\mathrm{\Theta }(C_n(\chi )),z(𝒰)=0,ϵ>0.$$
We conclude that the bifurcation case $`z(_1)>1`$ on the boundary of the existence domain $`\beta _n`$ does not result in bifurcation of any eigenvalue $`\lambda `$ of the stability problem (2.4), such that $`n(_1)p(𝒰)+n(_0)=2n`$ is valid everywhere in $`\beta _n_n`$. It follows from the Closure Theorem that the ground state with $`n=0`$ is spectrally stable in $`\beta _n`$, while the $`n`$-th excited state with $`n1`$ may have at most $`N_{\mathrm{unst}}`$ unstable eigenvalues, where $`0N_{\mathrm{unst}}2n`$. We show that $`N_{\mathrm{unst}}=2N_{\mathrm{comp}}=2n`$ in $`\beta _n`$ in a generic case.
At $`ϵ=0`$, the stability problem (2.4) can be decoupled as follows:
$$L_1u_1=\lambda w_1,L_0w_1=\lambda u_1$$
(3.26)
and
$$L_s(u_2\pm iw_2)=\pm i\lambda (u_2\pm iw_2).$$
(3.27)
The first problem (3.26) has the continuous spectrum for $`\mathrm{Re}(\lambda )=0`$ and $`|\mathrm{Im}(\lambda )|1`$ and the zero eigenvalue $`\lambda =0`$ of algebraic multiplicity 4 and geometric multiplicity 2. The second problem (3.27) has the continuous spectrum for $`\mathrm{Re}(\lambda )=0`$ and $`|\mathrm{Im}(\lambda )|\lambda _n^2(\chi )`$, zero eigenvalue $`\lambda =0`$ of geometric and algebraic multiplicity 2, and $`2n`$ isolated eigenvalues in the points $`\lambda =\pm i\left(\lambda _k^2\lambda _n^2\right)`$, where $`k=0,1,\mathrm{},n1`$. It follows from (3.6) that for $`\chi >\chi _n`$:
$$\lambda _k^2\lambda _n^2=(nk)\left[2\lambda _n(\chi )+(nk)\right]>(nk)^21,0k<n.$$
(3.28)
Therefore, $`2n`$ isolated eigenvalues of the problem (3.27) are embedded in the continuous spectrum of the problem (3.26). These embedded eigenvalues have negative energy $`h[𝐮,𝐰]`$, since at $`ϵ=0`$:
$$𝐮_k,_1𝐮_k=𝐰_k,_0𝐰_k=(\lambda _k^2\lambda _n^2)\mathrm{\Phi }_k^{(1)},\mathrm{\Phi }_k^{(1)},0k<n,$$
(3.29)
where $`𝐮_k=(0,\mathrm{\Phi }_k^{(1)})^T`$ and $`𝐰_k=(0,i\mathrm{\Phi }_k^{(1)})^T`$ at $`ϵ=0`$. By Appendix A.4, all $`2n`$ embedded eigenvalues of negative energy $`h[𝐮,𝐰]`$ bifurcate in a general case of non-zero $`\mathrm{\Gamma }`$, see Eq. (A.29), to complex unstable eigenvalues $`\lambda `$, $`\mathrm{Re}(\lambda )>0`$ for $`ϵ0`$, such that $`N_{\mathrm{real}}+2N_{\mathrm{comp}}+2N_{\mathrm{imag}}^{}=2N_{\mathrm{comp}}=2n`$ in $`\beta _n`$.
### 3.2 Numerical results
For $`ϵ0`$, the linearized problem (2.4) satisfies the assumptions of the Closure Theorem. Therefore, all unstable eigenvalues $`N_{\mathrm{unst}}=2N_{\mathrm{comp}}=2n`$ are structurally stable for larger values of $`ϵ`$, until new bifurcations occur in the parameter continuations. We study numerically locations of unstable eigenvalues in the linearized problem (2.4) related to the vector solitons $`𝚽=(\mathrm{\Phi }_1,\mathrm{\Phi }_2)^T`$ with nodal index $`𝐢=(0,n)^T`$, $`n=1,2`$. Our numerical algorithm is based on the shooting technique in the complex $`\lambda `$-plane. We also determine the indices $`n(_0)`$, $`n(_1)`$ and $`p(𝒰)`$ by a numerics-assisted procedure as described in , and relate them to the number of unstable eigenvalues by using the closure relation (2.9).
Figure 1 shows the 1st-family of multi-hump vector solitons with the correspondence: $`u=\mathrm{\Phi }_1(x)`$, $`v=\mathrm{\Phi }_2(x)`$, and $`\omega =\sqrt{\beta }`$. For $`\omega <1`$, this family exists between $`\chi _1(\beta )<\chi <\chi _2(\beta )`$, where $`\chi =\chi _2(\beta )`$ is the local bifurcation boundary, and $`\chi =\chi _1(\beta )`$ is the nonlocal bifurcation boundary . Hence the one-sided domain $`\beta _1`$ is located to the left of the local bifurcation curve, and $`\mathrm{sign}(C_1)=1`$ in (3.13). When parameter $`\omega =0.6`$ is fixed, we readily find that $`\chi _1=0.28`$ and $`\chi _2=2.08`$. When $`\chi `$ moves from $`\chi _2`$ to $`\chi _1`$, the distance between the two pulses in the $`v`$ component grows. It diverges to infinity at the nonlocal bifurcation boundary $`\chi _2`$ (near point $`a`$ in Fig. 1).
Figure 2 shows unstable eigenvalues of the linearized problem (2.4) for $`\omega =0.6`$ and $`\chi _1<\chi <\chi _2`$. In the domain $`\beta _1`$, there is a pair of unstable complex eigenvalues $`\lambda =\mathrm{Re}(\sigma _2)\pm i\mathrm{Im}(\sigma _2)`$, which bifurcate from the embedded eigenvalues $`\lambda =\pm i(\lambda _0^2\lambda _1^2)`$, see Eq. (3.28).
When $`\chi 1^+`$, the complex eigenvalues $`\lambda =\mathrm{Re}(\sigma _2)\pm i\mathrm{Im}(\sigma _2)`$ approach the imaginary axis and become embedded eigenvalues $`\lambda =\pm i(1\omega ^2)`$. The case $`\chi =1`$ corresponds to the integrable Manakov system, when the linearized problem (2.4) has the following exact solution:
$$𝐮_0=\left(\begin{array}{cc}\mathrm{\Phi }_2& \\ \mathrm{\Phi }_1& \end{array}\right)𝐰_0=i\left(\begin{array}{cc}\mathrm{\Phi }_2& \\ \mathrm{\Phi }_1& \end{array}\right),\lambda =\pm i(1\omega ^2).$$
(3.30)
This exact solution is generated by the additional polarizational-rotation symmetry in the potential function $`U=U(|\psi _1|^2+|\psi _2|^2)`$ at $`\chi =1`$. Embedded eigenvalues $`\lambda =\pm i(1\omega ^2)`$ have negative energy $`h[𝐮,𝐰]`$, since
$$𝐮_0,_1𝐮_0=𝐰_0,_0𝐰_0=(1\omega ^2)\left(\mathrm{\Phi }_1,\mathrm{\Phi }_1\mathrm{\Phi }_2,\mathrm{\Phi }_2\right)<0,$$
(3.31)
where the last inequality is confirmed numerically from integration of the exact solutions :
$`\mathrm{\Phi }_1(x)`$ $`=`$ $`{\displaystyle \frac{\sqrt{1\omega ^2}\mathrm{cosh}\omega x}{\mathrm{cosh}x\mathrm{cosh}\omega x\omega \mathrm{sinh}x\mathrm{sinh}\omega x}},`$ (3.32)
$`\mathrm{\Phi }_2(x)`$ $`=`$ $`{\displaystyle \frac{\omega \sqrt{1\omega ^2}\mathrm{sinh}x}{\mathrm{cosh}x\mathrm{cosh}\omega x\omega \mathrm{sinh}x\mathrm{sinh}\omega x}},`$ (3.33)
in the entire domain of existence: $`0<\omega <1`$. Since embedded eigenvalues $`\lambda =\pm i(1\omega ^2)`$ at $`\chi =1`$ have negative energy $`h[𝐮,𝐰]`$, they bifurcate to the complex plane for $`\chi 1`$ when the polarizational symmetry is destroyed. This is indeed the case as shown in Fig. 2, in agreement with Appendix A.4.
There is an additional instability bifurcation at $`\chi =1`$. This bifurcation comes about because at this $`\chi `$ value, the 1st-family of vector solitons $`𝚽=(\mathrm{\Phi }_1,\mathrm{\Phi }_2)^T`$ can be generalized to asymmetric solutions with an additional free parameter . As a result, the derivative of the asymmetric vector solitons with respect to the free parameter is an eigenvector in the kernel of the operator $`_1`$, such that $`z(_1)=2`$ at $`\chi =1`$. When $`\chi 1`$, the integrability of the Manakov system is destroyed, and a pair of real or purely imaginary eigenvalues is generated, in agreement with Appendix A.1. Indeed, Fig. 2 shows a pair of purely imaginary eigenvalues $`\lambda =\pm i\sigma _1`$ for $`\chi >1`$, which merge to the end points $`\lambda =\pm i\omega ^2`$ of the continuous spectrum at $`\chi =1.185`$, and a pair of real eigenvalues $`\lambda =\pm \sigma _1`$ for $`\chi <1`$.
Now we relate results of Fig. 2 to the closure relation (2.9). For this purpose, we have determined the indices $`n(_1)`$ and $`p(𝒰)`$ by the numerics-assisted procedure in (we note that $`n(_0)=1`$ everywhere in the existence domain of the 1st-family of vector solitons). For $`\omega =0.6`$, we have found numerically that
$$n(_1)=\{\begin{array}{c}4,\chi _1<\chi <1,\\ 3,1<\chi <\chi _2,\end{array}p(𝒰)=2,\text{for all }\chi _1<\chi <\chi _2,$$
(3.34)
such that
$$n(_1)+n(_0)p(𝒰)=\{\begin{array}{c}3,\chi _1<\chi <1,\\ 2,1<\chi <\chi _2.\end{array}$$
(3.35)
On the other hand, Fig. 2 shows that
$$N_{\mathrm{comp}}=1,\text{for all }\chi _1<\chi <\chi _2,N_{\mathrm{real}}=\{\begin{array}{c}1,\chi _1<\chi <1\\ 0,1<\chi <\chi _2\end{array}$$
(3.36)
and the closure relation (2.9) is thus satisfied.
Figures 3 and 4 show similar results for the 2nd family of multi-hump vector solitons. When parameter $`\omega =0.6`$ is fixed, the local bifurcation boundary is $`\chi _2=4.68`$ and the nonlocal bifurcation boundary is $`\chi _1=1.68`$. Again, the one-sided domain $`\beta _2`$ is located to the left of the local bifurcation boundary, such that $`\mathrm{sign}(C_2)=1`$. However, such solitons exist on both sides of the nonlocal bifurcation boundary $`\chi =\chi _1`$. As $`\chi `$ moves leftward from $`\chi =\chi _2`$, it first crosses the nonlocal bifurcation boundary $`\chi =\chi _1`$, then turns around, and approaches the nonlocal bifurcation boundary from the left side (see point $`a`$ in Fig. 3). This behavior has been explained both analytically and numerically in .
Using the same shooting algorithm, we have obtained the unstable eigenvalues of the linearized problem (2.4) and displayed them in Fig. 4. In the domain $`\beta _2`$, there exist two pairs of unstable complex eigenvalues, such that the pair $`\lambda =\mathrm{Re}(\sigma _4)\pm i\mathrm{Im}(\sigma _4)`$ bifurcates from the pair of embedded eigenvalues $`\lambda =\pm i(\lambda _0^2\lambda _2^2)`$, while the pair $`\lambda =\mathrm{Re}(\sigma _3)\pm i\mathrm{Im}(\sigma _3)`$ bifurcates from the pair of embedded eigenvalues $`\lambda =\pm i(\lambda _1^2\lambda _2^2)`$. Fig. 4 also shows that the pair $`\lambda =\mathrm{Re}(\sigma _3)\pm i\mathrm{Im}(\sigma _3)`$ approach the imaginary axis and become a pair of embedded eigenvalues at $`\chi =\chi _a2.49`$, but then reappear as a pair of complex unstable eigenvalues for $`\chi <\chi _a`$, in agreement with Appendix A.4.
There are two more instability bifurcations in Fig. 4. At $`\chi =\chi _b2.44`$, the zero eigenvalue bifurcates into a pair of imaginary eigenvalues $`\lambda =\pm i\sigma _2`$ for $`\chi >\chi _b`$, which then merge into the end points $`\lambda =\pm i\omega ^2`$ of the continuous spectrum at $`\chi 2.72`$. When $`\chi <\chi _b`$, this zero eigenvalue bifurcates into a pair of real unstable eigenvalues $`\lambda =\pm \sigma _2`$. At yet another point $`\chi =\chi _c=2.17`$, the zero eigenvalue bifurcates into a pair of imaginary eigenvalues $`\lambda =\pm i\sigma _1`$ for $`\chi >\chi _c`$, which merge into the end points $`\lambda =\pm i\omega ^2`$ of the continuous spectrum at $`\chi 3.01`$. When $`\chi <\chi _c`$, this zero eigenvalue bifurcates into a pair of real unstable eigenvalues $`\lambda =\pm \sigma _1`$.
To relate numerical results of Fig. 4 to the closure relation (2.9), we have again determined the indices $`n(_1)`$ and $`p(𝒰)`$ by the numerical algorithm, while $`n(_0)=2`$ everywhere in the existence domain of the 2nd family of vector solitons. For $`\omega =0.6`$, we have found numerically that
$$n(_1)=\{\begin{array}{c}5,\chi _1<\chi <\chi _b,\\ 4,\chi _b<\chi <\chi _2,\end{array}p(𝒰)=\{\begin{array}{c}1,\chi _1<\chi <\chi _c,\\ 2,\chi _c<\chi <\chi _2,\end{array}$$
(3.37)
such that
$$n(_1)+n(_0)p(𝒰)=\{\begin{array}{c}6,\chi _1<\chi <\chi _c,\hfill \\ 5,\chi _c<\chi <\chi _b,\hfill \\ 4,\chi _b<\chi <\chi _2.\hfill \end{array}$$
(3.38)
On the other hand, Fig. 4 shows that
$$N_{\mathrm{comp}}=2,\text{for all }\chi _1<\chi <\chi _2,N_{\mathrm{real}}=\{\begin{array}{c}2,\chi _1<\chi <\chi _c,\hfill \\ 1,\chi _c<\chi <\chi _b,\hfill \\ 0,\chi _b<\chi <\chi _2.\hfill \end{array}$$
(3.39)
Hence the closure relation (2.9) is satisfied. In particular, the instability bifurcation at $`\chi =\chi _b`$ is due to a jump in the index $`n(_1)`$ (similar to the 1st family), while the instability bifurcation at $`\chi =\chi _c`$ is due to a jump in the index $`p(𝒰)`$. These two bifurcations occur in agreement with Appendices A.1 and A.2.
Although our results were obtained here for a particular value $`\omega =0.6`$, we expect that similar results hold for other values of $`\omega `$, when $`0<\omega <1`$. We conclude that the 1st and 2nd families of multi-hump vector solitons in the coupled cubic NLS equations are all linearly unstable (except for the integrable Manakov system $`\chi =1`$, where the 1st family is neutrally stable).
It is remarkable that the main features of instability bifurcations for the 1st-family of multi-hump vector solitons are repeated for the 2nd-family of vector solitons, irrelevant whether the coupled NLS equations are integrable or not. This allows us to conjecture that a similar pattern of unstable eigenvalues persists for a general $`n`$-th family of multi-hump vector solitons, with more unstable eigenvalues and additional instability bifurcations appearing as $`n`$ increases.
## 4 Two coupled saturable NLS equations
We consider the system of two coupled saturable NLS equations :
$`i\psi _{1z}+\psi _{1xx}+{\displaystyle \frac{|\psi _1|^2+|\psi _2|^2}{1+s(|\psi _1|^2+|\psi _2|^2)}}\psi _1`$ $`=`$ $`0,`$
$`i\psi _{2z}+\psi _{2xx}+{\displaystyle \frac{|\psi _1|^2+|\psi _2|^2}{1+s(|\psi _1|^2+|\psi _2|^2)}}\psi _2`$ $`=`$ $`0,`$ (4.1)
where $`s>0`$. This system is a particular example of (2.1) with $`N=2`$, $`d_1=d_2=1`$, and
$$U=\frac{1}{s}\left(|\psi _1|^2+|\psi _2|^2\right)\frac{1}{s^2}\mathrm{log}\left(1+s(|\psi _1|^2+|\psi _2|^2)\right).$$
(4.2)
We consider again the $`n`$-th family of vector solitons $`𝚽=(\mathrm{\Phi }_1,\mathrm{\Phi }_2)^T`$ with nodal index $`𝐢=(0,n)^T`$, $`n`$ and convenient parametrization $`\beta _1=1`$ and $`\beta _2=\beta `$. The $`n`$-th family bifurcates from the scalar solution $`𝚽=(\mathrm{\Phi }_0,0)^T`$, where $`\mathrm{\Phi }_0(x)`$ satisfies the ODE:
$$\mathrm{\Phi }_0^{\prime \prime }\mathrm{\Phi }_0+\frac{\mathrm{\Phi }_0^3}{1+s\mathrm{\Phi }_0^2}=0.$$
(4.3)
The local bifurcation occurs at $`\beta =\beta _n(s)`$, when there exists a $`n`$-nodal bound state $`\mathrm{\Phi }_n(x)`$ in the linear eigenvalue problem:
$$\mathrm{\Phi }_n^{\prime \prime }\beta _n\mathrm{\Phi }_n+\frac{\mathrm{\Phi }_0^2\mathrm{\Phi }_n}{1+s\mathrm{\Phi }_0^2}=0.$$
(4.4)
Vector solitons in the $`n`$-th family disappear at the nonlocal bifurcation boundary. The domain of existence for the first three families $`n=0,1,2`$ has been obtained numerically in . We trace analytically unstable eigenvalues and show that $`N_{\mathrm{real}}+2N_{\mathrm{comp}}+2N_{\mathrm{imag}}^{}=2n`$ for small $`|\beta \beta _n(s)|1`$. At $`\beta =\beta _n(s)`$, the stability problem (2.4) can be decoupled as follows:
$$L_1u_1=\lambda w_1,L_0w_1=\lambda u_1$$
(4.5)
and
$$L_s(u_2\pm iw_2)=\pm i\lambda (u_2\pm iw_2),$$
(4.6)
where
$`L_0`$ $`=`$ $`{\displaystyle \frac{d^2}{dx^2}}+1{\displaystyle \frac{\mathrm{\Phi }_0^2}{1+s\mathrm{\Phi }_0^2}},`$ (4.7)
$`L_1`$ $`=`$ $`{\displaystyle \frac{d^2}{dx^2}}+1{\displaystyle \frac{\mathrm{\Phi }_0^2(3+s\mathrm{\Phi }_0^2)}{(1+s\mathrm{\Phi }_0^2)^2}},`$ (4.8)
$`L_s`$ $`=`$ $`{\displaystyle \frac{d^2}{dx^2}}+\beta _n(s){\displaystyle \frac{\mathrm{\Phi }_0^2}{1+s\mathrm{\Phi }_0^2}}.`$ (4.9)
The first problem (4.5) is the linearized stability problem in the scalar saturable NLS equation (4.1) for $`\mathrm{\Phi }_0(x)`$. Based on numerical data in , we assume that the bound state $`\mathrm{\Phi }_0(x)`$ is spectrally stable in the scalar saturable NLS equation and the problem (4.5) does not have any eigenvalues of negative energy $`h[𝐮,𝐰]`$. It has the continuous spectrum at $`\mathrm{Re}(\lambda )=0`$ and $`|\mathrm{Im}(\lambda )|1`$, the zero eigenvalue $`\lambda =0`$ of algebraic multiplicity 4 and geometric multiplicity 2, and possibly isolated eigenvalues $`\lambda i`$ of positive energy.
The second problem (4.6) has the continuous spectrum at $`\mathrm{Re}(\lambda )=0`$ and $`|\mathrm{Im}(\lambda )|\beta _n(s)`$, zero eigenvalue $`\lambda =0`$ of geometric and algebraic multiplicity 2, and $`2n`$ isolated eigenvalues $`\lambda =\pm i[\beta _k(s)\beta _n(s)]`$, where $`k=0,1,\mathrm{},n1`$. By the Sturm Nodal Theorem, eigenvalues $`\beta _k(s)`$ ($`k=0,\mathrm{},n`$) are ordered in the decreasing order and are characterized by eigenfunctions $`\mathrm{\Phi }_k(x)`$ with $`k`$ nodes, such that the ground state $`\mathrm{\Phi }_0(x)`$ corresponds to $`\beta _0(s)=1`$ and the $`n`$-th excited state $`\mathrm{\Phi }_n(x)`$ corresponds to $`\beta _n(s)`$. The $`2n`$ eigenvalues have negative energy $`h[𝐮,𝐰]`$, since
$$𝐮_k,_1𝐮_k=𝐰_k,_0𝐰_k=\left[\beta _k(s)\beta _n(s)\right]\mathrm{\Phi }_k,\mathrm{\Phi }_k<0,0k<n,$$
(4.10)
where $`𝐮_k=(0,\mathrm{\Phi }_k)^T`$ and $`𝐰_k=(0,i\mathrm{\Phi }_k)^T`$ at $`\beta =\beta _n(s)`$. Therefore, $`2N_{\mathrm{imag}}^{}=2n`$ at $`\beta =\beta _n(s)`$. Since the zero eigenvalue has the generic algebraic multiplicity six, no negative eigenvalues of $`_1`$ and $`_0`$ arise from the zero eigenvalue, such that we have $`N_{\mathrm{real}}+2N_{\mathrm{comp}}+2N_{\mathrm{imag}}^{}=2n`$ for $`|\beta \beta _n(s)|1`$ by continuity of the negative index $`n(h)`$. The ground state with $`n=0`$ is therefore spectrally stable in the existence domain (which is $`\beta =1`$, $`s>0`$ for $`n=0`$).
We show that the nodal bound state with $`n1`$ may have at most $`N_{\mathrm{unst}}`$ unstable eigenvalues, where $`0N_{\mathrm{unst}}2n2`$. The number of unstable eigenvalues is reduced by two, since there are two eigenvalues $`\lambda =\pm i(1\beta )`$ of negative energy, which exist for all vector solitons with $`n1`$ due to the polarizational-rotation symmetry in the potential function $`U=U(|\psi _1|^2+|\psi _2|^2)`$. This pair of eigenvalues is similar to the one which occurs in the coupled NLS equations (3.1) with $`\chi =1`$, such that the stability problem (2.4) has exactly the same solution (3.30) for $`\beta =\omega ^2`$. These eigenvalues have negative energy $`h[𝐮,𝐰]`$ due to (3.31), where the inequality remains true in the entire existence domain, as follows from numerical data in . If $`\beta <1/2`$, these eigenvalues $`\lambda =\pm i(1\beta )`$ are embedded in the continuous spectrum of the problem (2.4), but never bifurcate off the continuous spectrum as the parameters vary. Therefore, $`2N_{\mathrm{imag}}^{}2`$ and $`0N_{\mathrm{unst}}2n2`$ for $`n1`$. As a result, the 1st family of vector solitons is spectrally stable when $`\beta `$ is near the local bifurcation boundary $`\beta =\beta _n(s)`$.
When $`\beta =\beta _2(s)`$, the eigenvalues $`\lambda =\pm i\left[\beta _1(s)\beta _2(s)\right]`$ are embedded into the continuous spectrum if $`\beta _1(s)>2\beta _2(s)`$, and are isolated from the continuous spectrum if $`\beta _1(s)<2\beta _2(s)`$. In the first case, the embedded eigenvalues bifurcate generally to complex unstable eigenvalues for $`0<|\beta \beta _2(s)|1`$, according to Appendix A.4. In the second case, the isolated eigenvalues do not bifurcate to complex unstable eigenvalues for $`0<|\beta \beta _2(s)|1`$. Since these eigenvalues have negative energy \[see Eq. (4.10)\], the number of unstable eigenvalues $`N_{\mathrm{unst}}`$ is then zero. In other words, vector solitons of the 2nd family near the local bifurcation boundary with $`\beta _1(s)<2\beta _2(s)`$ are spectrally stable.
We compare the above theoretical results with numerical results in , where unstable eigenvalues in the linearized problem (2.4) were obtained for the 1st and 2nd families \[see Figs. 2 and 3 in , where $`(\lambda ,\beta )`$ correspond to our parameters $`(\beta ,\sigma )`$\]. It was shown in that the 1st-family of vector solitons is stable in the domain $`\beta _1(s)<\beta <\beta _{\mathrm{stab}}^{(1)}(s)`$, in agreement with our prediction of linear stability near the first local bifurcation boundary $`\beta =\beta _1(s)`$. At $`\beta =\beta _{\mathrm{stab}}^{(1)}(s)`$, the bifurcation $`z(𝒰)=1`$ occurs, which generates a pair of real eigenvalues for $`\beta >\beta _{\mathrm{stab}}^{(1)}(s)`$ and a pair of imaginary eigenvalues for $`\beta <\beta _{\mathrm{stab}}^{(1)}(s)`$, in agreement with Appendix A.2. Therefore,
$$n(_1)+n(_0)p(𝒰)=\{\begin{array}{c}3,\beta _{\mathrm{stab}}^{(1)}<\beta <1,\hfill \\ 2,\beta _1<\beta <\beta _{\mathrm{stab}}^{(1)},\hfill \end{array}$$
(4.11)
and
$$N_{\mathrm{imag}}^{}=2,\text{for all }\beta ,N_{\mathrm{real}}=\{\begin{array}{c}1,\beta _{\mathrm{stab}}^{(1)}<\beta <1,\hfill \\ 0,\beta _1<\beta <\beta _{\mathrm{stab}}^{(1)}.\hfill \end{array}$$
(4.12)
For the 2nd family, the eigenvalues $`\lambda =\pm i\left[\beta _1(s)\beta _2(s)\right]`$ are embedded when $`0.646<s<0.857`$ and isolated when $`s>0.857`$. According to our prediction, the embedded eigenvalues should bifurcate to the complex unstable eigenvalues for $`0<|\beta \beta _2(s)|1`$. However, it was claimed in that vector solitons in the 2nd family were stable near the local bifurcation boundary $`0<|\beta \beta _2(s)|1`$ for all values of $`s`$. In order to resolve this discrepancy, we have numerically determined the eigenvalue spectrum for vector solitons in the 2nd family by the shooting method. Figures 5 and 6 present numerical results for $`\beta =0.16`$ and $`\beta =0.49`$, respectively.
For $`\beta =0.16`$, the local bifurcation boundary of the 2nd family occurs at $`s=0.785`$ and the eigenvalues $`\lambda =\pm i\left[\beta _1(s)\beta _2(s)\right]`$ are embedded, such that a pair of unstable complex eigenvalues $`\lambda =\mathrm{Re}(\sigma _2)\pm i\mathrm{Im}(\sigma _2)`$ bifurcates for $`s<0.785`$. These complex eigenvalues indicate the oscillatory instability of vector solitons near the local bifurcation boundary $`|\beta \beta _2(s)|1`$ for $`\beta =0.16`$. Thus the claim in on the stability of 2nd-family vector solitons anywhere near the local bifurcation boundary is incorrect. Note that the complex eigenvalues $`\lambda =\mathrm{Re}(\sigma _2)\pm i\mathrm{Im}(\sigma _2)`$ persist throughout the entire existence domain of the 2nd family, which is $`0.513<s<0.785`$ at $`\beta =0.16`$.
Furthermore, a pair of purely imaginary eigenvalues $`\lambda =\pm i|\sigma _1|`$ bifurcate from the end points $`\lambda =\pm i\beta `$ of the continuous spectrum at $`s=0.763`$, merge into the origin at $`s=0.696`$, and then bifurcate into a pair of real eigenvalues $`\lambda =\pm \sigma _1`$ for $`s<0.696`$. This exponential instability induced by the real eigenvalue $`\lambda =\sigma _1`$ has been reported in . Since it was claimed in that $`z(𝒰)=0`$ in the entire existence domain of the 2nd family of vector solitons, we conclude that the instability bifurcation at $`s=0.696`$ falls into the scenario of Appendix A.1 with $`z(_1)=2`$. In the interval $`0.513<s<0.696`$, the oscillatory instability is overshadowed by the exponential instability from the eigenvalue $`\lambda =\sigma _1`$. However, in the interval $`0.696<s<0.785`$, this oscillatory instability is the only instability experienced by vector solitons. Since $`\mathrm{Re}(\sigma _2)`$ is less than $`0.011`$ in the interval $`0.696<s<0.785`$, it may explain why this oscillatory instability was missed in the numerical results of .
For $`\beta =0.49`$, the local bifurcation boundary of the 2nd family occurs at $`s=0.894`$, and the eigenvalues $`\lambda =\pm i\left[\beta _1(s)\beta _2(s)\right]`$ are isolated, such that vector solitons near the local bifurcation boundary are spectrally stable. This is indeed confirmed in Fig. 6, where the spectrum diagram is displayed for all values of $`s`$. It is seen that at the local bifurcation boundary $`s=0.894`$, there are two pairs of isolated imaginary eigenvalues. The pair $`\lambda =\pm i\left[\beta _1(s)\beta _2(s)\right]=\pm 0.292i`$ has the negative energy, while the pair $`\lambda =\pm 0.310i`$ has the positive energy. The second pair corresponds to the positive eigenvalue of $`_s`$ in (4.9). As $`s`$ moves leftward from the boundary point 0.894, these two imaginary eigenvalues move toward each other. At $`s=0.891`$, they coalesce and create a quadruple of complex eigenvalues, in agreement with Appendix A.3. However, this instability persists only in a tiny interval $`0.882<s<0.891`$, and it is very weak, with growth rates below 0.01. At $`s=0.882`$, these complex eigenvalues coalesce and bifurcate back into two pairs of purely imaginary eigenvalues again. When $`s`$ decreases further, one pair of these imaginary eigenvalues (denoted as $`\pm \sigma _1`$ in Fig. 6) always remain imaginary, but the other pair of imaginary eigenvalues ($`\pm \sigma _2`$ in Fig. 6) move toward zero and become real for $`s<0.80`$, in agreement with Appendix A.1. There is one more eigenvalue ($`\sigma _3`$) in Fig. 6, which bifurcates from the edge of the continuous spectrum at $`s=0.718`$, and always stays imaginary. The pattern of Fig. 6 differs from that of Fig. 5 in that complex instability is not set in at the local bifurcation boundary $`\beta =\beta _2(s)`$, and, once it is set in, it is confined in a narrow interval of $`s`$. Similar to the $`\beta =0.16`$ case above, this narrow interval of oscillatory instability was missed in , but the exponential instability was captured there.
Finally, we have mapped out the regions of exponential and oscillatory instabilities in the entire domain of existence for the 2nd family of vector solitons. The results are shown in Fig. 7. The almost-straight boundary lines show the local and nonlocal bifurcation boundaries . The large domain of exponential instability away from the local bifurcation boundary corresponds to the one computed in Fig. 2 of . The domain of oscillatory instability consists of two sub-domains. The larger sub-domain is at $`s<0.857`$, where the eigenvalues $`\lambda =\pm i\left[\beta _1(s)\beta _2(s)\right]`$ at the local bifurcation boundary $`\beta =\beta _2(s)`$ are embedded. The smaller sub-domain is at $`s>0.857`$, where the eigenvalues $`\lambda =\pm i\left[\beta _1(s)\beta _2(s)\right]`$ are isolated but bifurcate to complex eigenvalues away from the local bifurcation boundary $`\beta =\beta _2(s)`$. Both sub-domains of oscillatory instabily were missed in .
## 5 Summary
We have applied the Closure Theorem for the negative index of the linearized Hamiltonian to multi-hump vector solitons in the general coupled NLS equations (2.1). Unstable eigenvalues of the linearized problem (2.4) are approximated with the perturbation series expansions and found numerically with the shooting method. Not only the numerical results are in excellent agreement with the closure relation (2.9), but also the closure relation (2.9) shows that all unstable eigenvalues are recovered with the numerical shooting method. These analytical and numerical results establish that all multi-hump vector solitons in the non-integrable coupled cubic NLS equations are linearly unstable, while multi-hump vector solitons in the coupled saturable NLS equations can be linearly stable in certain regions of the parameter space. In the latter case, we have also discovered a new oscillatory instability which was missed before. This oscillatory instability significantly reduces the stability domains of vector solitons.
We note that the Closure Theorem is applied differently in Sections 3 and 4. For coupled cubic NLS equations, we compute the closure relation (2.9) from the right-hand side $`n(_1)p(𝒰)+n(_0)`$ at the local bifurcation boundary and then match it with the number of unstable eigenvalues of the linearized problem (2.4). For coupled saturable NLS equations, we compute the closure relation (2.9) directly from the left-hand side $`N_{\mathrm{real}}+2N_{\mathrm{comp}}+2N_{\mathrm{imag}}^{}`$ at the local bifurcation boundary and continue it in the entire existence domain.
While the negative index theory is well understood in and well illustrated in and this paper, the following problem still remains a challenge: how do we understand the stability of vector solitons in sign-indefinite coupled-mode equations, such as the nonlinear Dirac equations? The Closure Theorem is obviously invalid for the Dirac equations as the continuous spectrum has both positive and negative energies. Thus a generalization of the Closure Theorem to such systems is highly desirable for future advances.
Acknowledgements. The work of D.E.P. was supported in part by NSERC grant No. 5-36694. The work of J.Y. was supported in part by NASA EPSCoR and AFOSR grants.
## Appendix A Appendix: Bifurcations of unstable eigenvalues
Here we classify the four special cases, when one of the assumptions (i)–(iv) of the Closure Theorem fails. We derive sufficient conditions, when new unstable eigenvalues with $`\mathrm{Re}(\lambda )>0`$ bifurcate in the linearized problem (2.4) from eigenvalues with $`\mathrm{Re}(\lambda )=0`$. For clarity of notations, we use an equivalent form of the spectral problem (2.4):
$$_1𝐮=\gamma _0^1𝐮,𝐮X_c^{(u)}(),$$
(A.1)
where $`\gamma =\lambda ^2`$ and $`X_c^{(u)}()`$ is the constrained subspace of $`L^2()`$:
$`X_c^{(u)}=\{𝐮L^2:\mathrm{\Phi }_n𝐞_n,𝐮=0,n=1,\mathrm{},N\}.`$ (A.2)
Here $`𝐞_1`$,…,$`𝐞_N`$ are unit vectors in $`^N`$ and none of the components $`\mathrm{\Phi }_n(x)`$ is assumed to vanish identically on $`x`$. Operator $`_0`$ is always invertible in $`X_c^{(u)}()`$, since eigenvectors $`\{\mathrm{\Phi }_n(x)𝐞_n\}_{n=1}^N`$ form a basis in the kernel of $`_0`$. In the domain $`𝒟_ϵ=\{\lambda :|\lambda |>ϵ\}`$ for any $`ϵ>0`$, there exists a relation between $`𝐮(x)`$ and $`𝐰(x)`$:
$$𝐰=\lambda _0^1𝐮,\lambda 𝒟_ϵ,$$
(A.3)
such that two eigenvectors $`(𝐮,𝐰)^T`$ and $`(𝐮,𝐰)^T`$ of the linearized problem (2.4) for $`\lambda `$ and $`(\lambda )`$ correspond to a single eigenvector $`𝐮(x)`$ of the problem (A.1) for $`\gamma =\lambda ^2`$.
### A.1 Zero eigenvalues of $`_1`$ and $`_0`$
The kernel of $`_0`$ has a basis of $`N`$ eigenvectors $`\{\mathrm{\Phi }_n(x)𝐞_n\}_{n=1}^N`$. Therefore, the existence domain of the vector solitons (2.2) is confined by the boundaries, where $`\mathrm{\Phi }_n(x)0`$ for some $`n=1,\mathrm{},N`$. Let $``$ be an open simply-connected domain in the parameter space $`(\beta _1,\mathrm{},\beta _N)`$, where the vector solitons (2.2) exist. No bifurcations of zero eigenvalues of $`_0`$ occur in $`𝜷`$.
The kernel of $`_1`$ has always the eigenvector $`𝚽^{}(x)`$. When $`z(_1)=1`$, this is the only eigenvector in the kernel of $`_1`$. When $`z(_1)>1`$, additional linearly independent eigenvectors $`𝐮_0(x)`$ exist in the kernel of $`_1`$, such that the geometric multiplicity of $`\lambda =0`$ in the linearization problem (2.4) exceeds $`(N+1)`$. Under parameter continuation, the zero eigenvalue $`\lambda =0`$ generally moves either to the real or purely imaginary axes of $`\lambda `$. We study this bifurcation in the case, when $`z(_1)=2`$, $`z(𝒰)=0`$, and $`𝜷`$.
###### Proposition A.1
Let $`ϵ`$ be the bifurcation parameter and, at $`ϵ=0`$, there exists a non-zero eigenvector $`𝐮_0X_c^{(u)}()`$, such that $`_1𝐮_0=\mathrm{𝟎}`$ and $`𝐮_0`$ is linearly independent of $`𝚽^{}(x)`$. Assume that $`_1(ϵ)`$ and $`_0(ϵ)`$ are $`C^1`$-functions at $`ϵ=0`$, such that $`l_0=𝐮_0,_0^1(0)𝐮_00`$ and $`\delta l_1=𝐮_0,_1^{}(0)𝐮_00`$. Then, there exists $`ϵ_0>0`$ such that the linearized problem (2.4) has a real positive eigenvalue $`\lambda `$ in the domain:
$$𝒟_ϵ=\{ϵ:0<|ϵ|<ϵ_0,\mathrm{sign}(ϵ)=\mathrm{sign}(l_0\delta l_1)\}.$$
* Proof. We expand solutions of (A.1) in power series of $`ϵ`$:
$$𝐮(x)=𝐮_0(x)+ϵ𝐮_1(x)+\mathrm{O}(ϵ^2),\gamma =ϵ\gamma _1+\mathrm{O}(ϵ^2).$$
(A.4)
The function $`𝐮_1(x)`$ solves the non-homogeneous problem in $`X_c^{(u)}()`$:
$$_1(0)𝐮_1+_1^{}(0)𝐮_0=\gamma _1_0^1(0)𝐮_0.$$
(A.5)
Using the Fredholm Alternative Theorem, we find from (A.5) that $`\delta l_1=\gamma _1l_0`$. If $`\delta l_10`$, $`l_00`$, and $`\mathrm{sign}(ϵ)=\mathrm{sign}(l_0\delta l_1)`$, the eigenvalue $`\gamma `$ is negative in the first order of $`ϵ`$, such that $`\lambda =\pm \sqrt{\gamma }`$ are real.
###### Corollary A.1
Let $`_0`$ be positive definite, such that $`l_0>0`$. A new negative eigenvalue $`\mu (ϵ)`$ of $`_1`$ as $`ϵ0`$ results in a new negative eigenvalue $`\gamma (ϵ)`$ of the problem (A.1), such that
$$\underset{ϵ0}{lim}\frac{\gamma }{\mu }=\frac{𝐮_0,𝐮_0}{𝐮_0,_0^1(0)𝐮_0}.$$
(A.6)
Bifurcation $`z(_1)>1`$ may occur only if $`N>1`$ in the system (2.1). Analysis of this bifurcation with the Lyapunov-Schmidt reduction method are reported in a similar content in .
### A.2 Zero eigenvalues of $`𝒰`$
When $`z(𝒰)>0`$, algebraic multiplicity of $`\lambda =0`$ exceeds $`(2N+2)`$. Under a parameter continuation, the zero eigenvalue $`\lambda =0`$ generally moves either to the real or purely imaginary axis of $`\lambda `$. We study this bifurcation in the case, when $`z(𝒰)=1`$, $`z(_1)=1`$, and $`𝜷`$.
###### Proposition A.2
Let $`ϵ`$ be the bifurcation parameter and, at $`ϵ=0`$, there exists a non-zero eigenvector $`𝛎^N`$, $`𝒰𝛎=\mathrm{𝟎}`$, such that the eigenvector $`𝐮_0X_c^{(u)}()`$ solves the problem:
$$_1𝐮_0=\underset{n=1}{\overset{N}{}}\nu _n\mathrm{\Phi }_n(x)𝐞_n.$$
(A.7)
Assume that $`_1(ϵ)`$, $`_0(ϵ)`$, and $`𝒰(ϵ)`$ are $`C^1`$-functions at $`ϵ=0`$, such that $`l_0=𝐮_0,_0^1(0)𝐮_00`$ and $`\delta u=𝛎,𝒰^{}(0)𝛎0`$. Then, there exists $`ϵ_0>0`$ such that the linearized problem (2.4) has a real positive eigenvalue $`\lambda `$ in the domain:
$$𝒟_ϵ=\{ϵ:0<|ϵ|<ϵ_0,\mathrm{sign}(ϵ)=\mathrm{sign}(l_0\delta u)\}.$$
* Proof. If there exists $`𝝂^N`$, such that $`𝒰(0)𝝂=\mathrm{𝟎}`$, then the eigenvector $`𝐮_0(x)`$ for the problem (A.7) is given explicitly as
$$𝐮_0(x)=\underset{n=1}{\overset{N}{}}\nu _n\frac{𝚽(x)}{\beta _n}.$$
(A.8)
Using $`_1𝚽/\beta _n=\mathrm{\Phi }_n𝐞_n`$ and $`_0\mathrm{\Phi }_n𝐞_n=\mathrm{𝟎}`$ for any $`ϵ`$, we obtain the following derivative relations:
$`_1(0){\displaystyle \frac{𝚽^{}(0)}{\beta _n}}+_1^{}(0){\displaystyle \frac{𝚽}{\beta _n}}`$ $`=`$ $`\mathrm{\Phi }_n^{}(0)𝐞_n,`$ (A.9)
$`_0(0)\mathrm{\Phi }_n^{}(0)𝐞_n+_0^{}(0)\mathrm{\Phi }_n𝐞_n`$ $`=`$ $`0,`$ (A.10)
where $`𝚽^{}(0)`$ stands for derivative of $`𝚽(ϵ)`$ in $`ϵ`$. We expand solutions of (A.1) in power series of $`ϵ`$:
$$𝐮(x)=𝐮_0(x)+ϵ𝐮_1(x)+\mathrm{O}(ϵ^2),\gamma =ϵ\gamma _1+\mathrm{O}(ϵ^2).$$
(A.11)
Since the eigenvector $`𝐮_0(x)`$ solves the non-homogeneous problem (A.7), the relation between $`𝐮(x)`$ and $`𝐰(x)`$ is modified as follows:
$$𝐰=\lambda _0^1𝐮+\frac{1}{\lambda }\underset{n=1}{\overset{N}{}}\nu _n\mathrm{\Phi }_n(x)𝐞_n.$$
(A.12)
As a result, the function $`𝐮_1(x)`$ solves the non-homogeneous problem:
$$_1(0)𝐮_1+_1^{}(0)𝐮_0=\gamma _1_0^1(0)𝐮_0\underset{n=1}{\overset{N}{}}\nu _n\mathrm{\Phi }_n^{}(0)𝐞_n,$$
(A.13)
subject to the constraints:
$$\mathrm{\Phi }_n(0)𝐞_n,𝐮_1+\mathrm{\Phi }_n^{}(0)𝐞_n,𝐮_0=0.$$
(A.14)
Using the Fredholm Alternative Theorem, we find from (A.13) and (A.14) that
$$𝐮_0,_1^{}(0)𝐮_0=\gamma _1𝐮_0,_0^1(0)𝐮_02\underset{n=1}{\overset{N}{}}\nu _n\mathrm{\Phi }_n^{}(0)𝐞_n,𝐮_0.$$
(A.15)
As a result,
$$\gamma _1𝐮_0|_0^1(0)𝐮_0=\underset{n=1}{\overset{N}{}}\underset{m=1}{\overset{N}{}}\nu _n\nu _m\frac{}{ϵ}\mathrm{\Phi }_m(ϵ)𝐞_m,\frac{𝚽(ϵ)}{\beta _n}|_{ϵ=0}=\frac{1}{2}𝝂,𝒰^{}(0)𝝂,$$
(A.16)
such that $`\delta u=2\gamma _1l_0`$. When $`\delta u0`$, $`l_00`$, and $`\mathrm{sign}(ϵ)=\mathrm{sign}(l_0\delta u)`$, the eigenvalue $`\gamma `$ is negative in the first order of $`ϵ`$, such that $`\lambda =\pm \sqrt{\gamma }`$ are real.
###### Corollary A.2
Let $`_0`$ be positive definite, such that $`l_0>0`$. A new negative eigenvalue $`\mu (ϵ)`$ of $`𝒰`$ as $`ϵ0`$ results in a new negative eigenvalue $`\gamma (ϵ)`$ of the problem (A.1).
Bifurcation $`z(𝒰)>0`$ was analyzed in with power series expansions of the problem (2.4) near $`\lambda =0`$ and in with Taylor series expansions of the Evans function.
### A.3 Multiple non-zero eigenvalues of zero energy
When the problem (2.4) has a multiple eigenvalue $`\lambda =\lambda _0i`$ of zero energy, the corresponding eigenvector $`(𝐮_0,𝐰_0)^T`$ satisfies the conditions of the Fredholm Alternative Theorem: $`𝐮_0,_1𝐮_0=0`$, $`𝐰_0,_0𝐰_0=0`$, such that $`h[𝐮_0,𝐰_0]=0`$ and $`l_0=𝐮_0,_0^1𝐮_0=0`$. Under parameter continuations, multiple eigenvalues are generally destroyed and new complex eigenvalues $`\lambda `$ may arise in the problem (2.4). We study this bifurcation in the case, when a multiple eigenvalue $`\lambda =\lambda _0`$ has algebraic multiplicity two and geometric multiplicity one, while $`𝜷`$.
###### Proposition A.3
Let $`ϵ`$ be the bifurcation parameter and, at $`ϵ=0`$, there exist non-zero eigenvectors $`𝐮_0,𝐮_0^{}X_c^{(u)}()`$ for $`\gamma _0`$ and $`\gamma _0<\beta _{\mathrm{min}}^2`$, such that
$`_1𝐮_0`$ $`=`$ $`\gamma _0_0^1𝐮_0,`$ (A.17)
$`_1𝐮_0^{}`$ $`=`$ $`\gamma _0_0^1𝐮_0^{}+_0^1𝐮_0,`$ (A.18)
and $`l_0=𝐮_0,_0^1𝐮_0=0`$. Assume that $`_1(ϵ)`$ and $`_0(ϵ)`$ are $`C^1`$-functions at $`ϵ=0`$, such that $`l_0^{}=𝐮_0,_0^1(0)𝐮_0^{}0`$ and $`\delta h=𝐮_0,\left(_1^{}(0)\gamma _0_0^1(0)\right)𝐮_00`$. Then, there exists $`ϵ_0>0`$ such that the linearized problem (2.4) has two complex eigenvalues $`\lambda `$ with $`\mathrm{Re}(\lambda )>0`$ in the domain:
$$𝒟_ϵ=\{ϵ:0<|ϵ|<ϵ_0,\mathrm{sign}(ϵ)=\mathrm{sign}(l_0^{}\delta h)\}.$$
* Proof. We expand solutions of (A.1) in power series of $`ϵ^{1/2}`$:
$`𝐮(x)`$ $`=`$ $`𝐮_0(x)+ϵ^{1/2}\gamma _1𝐮_0^{}(x)+ϵ𝐮_2(x)+\mathrm{O}(ϵ^{3/2}),`$ (A.19)
$`\gamma `$ $`=`$ $`\gamma _0+ϵ^{1/2}\gamma _1+ϵ\gamma _2+\mathrm{O}(ϵ^{3/2}).`$ (A.20)
The function $`𝐮_2(x)`$ solves the non-homogeneous problem in $`X_c^{(u)}()`$:
$`_1(0)𝐮_2+_1^{}(0)𝐮_0=\gamma _0_0^1(0)𝐮_2+\gamma _0_0^1(0)𝐮_0+\gamma _1^2_0^1(0)𝐮_0^{}+\gamma _2_0^1(0)𝐮_0.`$ (A.21)
Using the Fredholm Alternative Theorem, we find from (A.21) that $`\delta h=\gamma _1^2l_0^{}`$. Since $`(\gamma \gamma _0)^2=ϵ\gamma _1^2+\mathrm{O}(ϵ^{3/2})`$, the eigenvalues $`\gamma `$ bifurcate into the complex plane if $`l_0^{}0`$, $`\delta h0`$, and $`\mathrm{sign}(ϵ)=\mathrm{sign}(l_0^{}\delta h)`$.
###### Corollary A.3
Let $`l_0^{}0`$ and $`\delta h0`$. There exists $`ϵ_0>0`$ such that the problem (A.1) has two real eigenvalues $`\gamma `$ in the domain,
$$𝒟_ϵ=\{ϵ:0<|ϵ|<ϵ_0,\mathrm{sign}(ϵ)=\mathrm{sign}(l_0^{}\delta h)\}$$
with oppositely signed quadratic forms
$$𝐮,_0^1𝐮=2ϵ^{1/2}\gamma _1l_0^{}+\mathrm{O}(ϵ).$$
When $`0<\gamma _0<\beta _{\mathrm{min}}`$, multiple eigenvalue $`\lambda =\lambda _0`$ is purely imaginary, and the bifurcation of Proposition A.3 is the instability bifurcation. When $`\gamma _0<0`$, the eigenvalue $`\lambda =\lambda _0`$ is purely real and is thus already unstable. Characteristic features of bifurcation of multiple eigenvalues were analyzed in . This bifurcation is generic when purely imaginary eigenvalues of positive and negative energies $`h[𝐮,𝐰]`$ coalesce, according to Corollary A.3 . General results on the collisions of purely imaginary eigenvalues of different energies and the concept of so-called Krein signatures can be found in .
### A.4 Embedded eigenvalues
When the problem (2.4) has an embedded eigenvalue $`\lambda =\lambda _0i`$ with $`|\mathrm{Im}(\lambda _0)|>\beta _{\mathrm{min}}`$, it is generically unstable under parameter continuation . When it has a positive energy, it disappears from the continuous spectrum, while when it has a negative energy, it bifurcates as complex unstable eigenvalues with $`Re(\lambda )>0`$ . We study this bifurcation in the case, when an embedded eigenvalue has the geometric and algebraic multiplicities one, while $`𝜷`$.
###### Proposition A.4
Let $`ϵ`$ be the bifurcation parameter and, at $`ϵ=0`$, there exist a non-zero eigenvector $`𝐮_0X_c^{(u)}()`$ for $`\gamma _0`$, $`\gamma _0>\beta _{\mathrm{min}}^2`$, such that
$`_1𝐮_0=\gamma _0_0^1𝐮_0.`$ (A.22)
Assume that $`_1(ϵ)`$ and $`_0(ϵ)`$ are $`C^1`$-functions at $`ϵ=0`$, such that $`l_0=𝐮_0,_0^1(0)𝐮_0<0`$ and $`\mathrm{\Gamma }0`$ in (A.29) below. Then, there exists $`ϵ_0>0`$ such that the linearized problem (2.4) has two complex eigenvalues $`\lambda `$ with $`\mathrm{Re}(\lambda )>0`$.
* Proof. For embedded eigenvalues, we use the linearized problem in the original form (2.4). Consider perturbation series expansions near the embedded eigenvalue $`\lambda =\lambda _0`$, with $`\mathrm{Im}(\lambda _0)>\beta _{\mathrm{min}}`$:
$`𝐮(x)`$ $`=`$ $`𝐮_0(x)+ϵ𝐮_1(x)+ϵ^2𝐮_2(x)+\mathrm{O}(ϵ^3),`$ (A.23)
$`𝐰(x)`$ $`=`$ $`𝐰_0(x)+ϵ𝐰_1(x)+ϵ^2𝐰_2(x)+\mathrm{O}(ϵ^3),`$ (A.24)
and
$$\lambda =\lambda _0+ϵ\lambda _1+ϵ^2\lambda _2+\mathrm{O}(ϵ^3).$$
(A.25)
Corrections of the perturbation series (A.23) and (A.24) satisfy linear non-homogeneous equations following from the linearized problem (2.4):
$`_1(0)𝐮_1+\lambda _0𝐰_1`$ $`=`$ $`_1^{}(0)𝐮_0\lambda _1𝐰_0,`$
$`_0(0)𝐰_1\lambda _0𝐮_1`$ $`=`$ $`_0^{}(0)𝐰_0+\lambda _1𝐮_0`$ (A.26)
and
$`_1(0)𝐮_2+\lambda _0𝐰_2`$ $`=`$ $`_1^{}(0)𝐮_1{\displaystyle \frac{1}{2}}_1^{\prime \prime }(0)𝐮_0\lambda _1𝐰_1\lambda _2𝐰_0,`$
$`_0(0)𝐰_2\lambda _0𝐮_2`$ $`=`$ $`_0^{}(0)𝐰_1{\displaystyle \frac{1}{2}}_0^{\prime \prime }(0)𝐰_0+\lambda _1𝐮_1+\lambda _2𝐮_0.`$ (A.27)
Bounded solutions of the problem (A.26) exist only if the right-hand-side is orthogonal to the eigenvector $`(𝐮_0,𝐰_0)`$. The solvability condition results in the equation:
$`\lambda _1\left(𝐰_0,𝐮_0𝐮_0,𝐰_0\right)=𝐮_0,_1^{}(0)𝐮_0+𝐰_0,_0^{}(0)𝐰_0,`$
such that $`\mathrm{Re}(\lambda _1)=0`$. Since the eigenvalue $`\lambda =\lambda _0`$ belongs to the continuous spectrum of the problem (2.4), the correction terms $`(𝐮_1,𝐰_1)`$ have non-vanishing tails in the limit $`|x|\mathrm{}`$. Assuming that $`\beta _1\beta _2\mathrm{}\beta _N`$, we add Sommerfeld radiation conditions to uniquely determine the correction terms $`(𝐮_1,𝐰_1)`$:
$$\left(\begin{array}{c}𝐮_1\\ 𝐰_1\end{array}\right)\underset{j=1}{\overset{K_{\lambda _0}}{}}g_j^\pm \left(\begin{array}{cc}𝐞_j& \\ i𝐞_j& \end{array}\right)e^{ik_jx},x\pm \mathrm{},$$
(A.28)
where $`g_j^\pm `$ are some constants, $`k_j=\sqrt{(\mathrm{Im}(\lambda _0)\beta _j)/d_j}`$, and $`K_{\lambda _0}`$ is the number of branches with $`k_j`$. It follows from (A.26) that
$`\mathrm{Im}\left(𝐮_1,_1^{}(0)𝐮_0+𝐰_1,_0^{}(0)𝐰_0+\lambda _1𝐮_1,𝐰_0\lambda _1𝐰_1,𝐮_0\right)`$ (A.29)
$`=`$ $`{\displaystyle \frac{1}{2i}}\left(𝐮_1,_1𝐮_1+𝐰_1,_0𝐰_1_1𝐮_1,𝐮_1_0𝐰_1,𝐰_1\right)`$
$`=`$ $`2{\displaystyle \underset{j=1}{\overset{K_{\lambda _0}}{}}}d_jk_j\left(|g_j^+|^2+|g_j^{}|^2\right)\mathrm{\Gamma }0.`$
Again, bounded solutions of the problem (A.27) exist only if
$`\lambda _2\left(𝐰_0,𝐮_0𝐮_0,𝐰_0\right)+\lambda _1\left(𝐰_0,𝐮_1𝐮_0,𝐰_1\right)`$
$`=𝐮_0,_1^{}(0)𝐮_1+𝐰_0,_0^{}(0)𝐰_1+{\displaystyle \frac{1}{2}}𝐮_0,_1^{\prime \prime }(0)𝐮_0+{\displaystyle \frac{1}{2}}𝐰_0,_0^{\prime \prime }(0)𝐰_0,`$ (A.30)
such that
$$\mathrm{Re}(\lambda _2)=\frac{\mathrm{Im}\left(𝐮_1,_1^{}(0)𝐮_0+𝐰_1,_0^{}(0)𝐰_0\right)}{2\mathrm{I}\mathrm{m}(\lambda _0)𝐮_0,_0^1𝐮_0}=\underset{j=1}{\overset{K_{\lambda _0}}{}}\frac{d_jk_j\left(|g_j^+|^2+|g_j^{}|^2\right)}{\mathrm{Im}(\lambda _0)𝐮_0,_0^1𝐮_0}.$$
(A.31)
When $`\mathrm{\Gamma }0`$ and $`l_0=𝐮_0,_0^1(0)𝐮_0<0`$, the embedded eigenvalue $`\lambda =\lambda _0`$ becomes a complex unstable eigenvalue $`\lambda `$ with $`\mathrm{Re}(\lambda )>0`$.
###### Corollary A.4
The linearized problem (2.4) does not have complex or embedded eigenvalues $`\lambda `$ if $`l_0=𝐮_0,_0^1(0)𝐮_0>0`$ and $`\mathrm{\Gamma }0`$.
* Proof. The formal computation in (A.31) predicts that $`\mathrm{Re}(\lambda _2)<0`$, when $`\mathrm{\Gamma }0`$ and $`l_0>0`$. However, the correction terms $`(𝐮_1,𝐰_1)^T`$ in (A.28) grow exponentially in $`x`$, since $`k_j(\lambda )=\sqrt{(\mathrm{Im}(\lambda )i\mathrm{Re}(\lambda )\beta _j)/d_j}`$ implies that $`\mathrm{Im}(k_j)>0`$. The embedded eigenvalue $`\lambda =\lambda _0`$ becomes a resonant pole with $`\mathrm{Re}(\lambda )<0`$.
Characteristic features of bifurcations of embedded eigenvalues were analyzed in . This bifurcation is generic for multi-hump vector solitons at the boundaries of the existence domain $`𝜷`$ .
Summarizing, there exist four bifurcations, which may lead to unstable eigenvalues in the spectral problem (2.4): (i) $`z(_1)>1`$, (ii) $`z(𝒰)>0`$, (iii) multiple eigenvalue $`\lambda _0i`$, $`|\mathrm{Im}(\lambda _0)|<\beta _{\mathrm{min}}`$, and (iv) embedded eigenvalue $`\lambda _0i`$, $`|\mathrm{Im}(\lambda _0)|>\beta _{\mathrm{min}}`$. Let $`n_X(h)`$ be the negative index of the linearized Hamiltonian in the constrained space $`X_c^{(u)}()`$, i.e. the number of negative eigenvalues of $`_1`$ and $`_0`$ in $`X_c^{(u)}()`$. It is known from (see also ) that
$$n_X(h)=n(_1)p(𝒰)+n(_0).$$
(A.32)
The closure relation (2.9) gives then:
$$n_X(h)=N_{\mathrm{real}}+2N_{\mathrm{imag}}^{}+2N_{\mathrm{comp}}.$$
(A.33)
It is clear from (A.33) that bifurcations (i) and (ii) change the negative index $`n_X(h)`$ due to a change in $`N_{\mathrm{real}}`$ by one, while bifurcations (iii) and (iv) do not change the negative index $`n_X(h)`$ due to an exchange in $`N_{\mathrm{imag}}^{}`$ and $`N_{\mathrm{comp}}`$.
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# Multi-wavelength observations of PKS 2155-304 with H.E.S.S.
## 1 Introduction
The innermost regions of active galactic nuclei, where the largest part of their luminosity is emitted, can be probed through observations of their flux variability at different wavelengths. The physical processes in their central engines and jets are usually considered the main candidates for the origin of the observed variability. Measurements of correlated variability, spectral variations and time lags across the broad-band observations allow modelling of particle distributions and their radiation processes, as well as probing the acceleration mechanisms that are involved.
PKS 2155$``$304 is probably the most prominent and best-studied blazar-type Active Galactic Nuclei (AGN) in the Southern Hemisphere. The emission of PKS 2155$``$304, its possible variability patterns, as well as correlations across all wavebands, have been studied exhaustively over the past 20 years (see e.g. Urry et al. (1997)). Its first detection at VHE $`\gamma `$-rays by the Durham Mk VI telescopes (Chadwick et al. (1999)) classified it as a TeV blazar, like the northern hemisphere BL Lac objects Mkn 421, Mkn 501, H 1426$`+`$428, or 1ES 1959$`+`$650. Its redshift of $`z=0.117`$ makes it the second most distant confirmed TeV blazar after H 1426$`+`$428 ($`z=0.129`$). PKS 2155$``$304 was the brightest BL Lac object in the EUVE all-sky survey (Marshall, Fruscione & Carone (1995)). This source was confirmed as a high energy $`\gamma `$-ray emitter by H.E.S.S. (Aharonian et al. (2004), AH04 hereafter) at the $`45\sigma `$ significance level, when strong detections were reported for each of the dark periods of observations.
Here we report on simultaneous H.E.S.S. VHE $`\gamma `$-ray, RXTE/PCA X-ray, ROTSE optical, and NRT decimetric observations of PKS 2155$``$304 during the dark periods of October and November 2003. No simultaneous multi-wavelength campaign had before included an Atmospheric Cherenkov Telescope (ACT) that could sample the evolution of the high energy component of the spectral energy distribution (SED) of this object. We also include EGRET archival data, and other archival radio through X-ray data obtained from the NASA/IPAC Extragalactic Database (NED). Details of the observations and data reduction/analysis are given in §2. Light curves and spectra are described in §3. The attenuation of the H.E.S.S. spectrum by the EBL and an interpretation of the data using a leptonic and a hadronic model are discussed in §4.
## 2 Observations and data analysis
### 2.1 H.E.S.S.
#### 2.1.1 H.E.S.S. detector
In its first phase, the H.E.S.S. array consists of four atmospheric Cherenkov telescopes operating in stereoscopic mode. However, the data shown here were taken during the construction of the system, initially with two telescopes, with one more as of mid-September, 2003. The fourth and final telescope was added to the array in December 2003, subsequent to the data presented here. Each telescope has a tessellated $`13\mathrm{m}`$-diameter ($`107\mathrm{m}^2`$ surface area) mirror which focuses the Cherenkov light from the showers of secondary particles created by the interaction of $`\gamma `$-rays in the atmosphere onto a camera in the focal plane. This camera consists of 960 photomultipliers, each with a pixel size of $`0.16^{}`$, giving a field of view of $`5^{}`$. For the data sets presented here the $`\gamma `$-ray trigger threshold is $`100\mathrm{GeV}`$ and the spectral threshold is $`160\mathrm{GeV}`$ with an energy resolution $`15\%`$. The experiment is located in the Khomas highlands in Namibia, ($`23^{}\mathrm{S}`$, $`15^{}\mathrm{E}`$, $`1800\mathrm{m}`$ a.s.l.). A detailed description of the layout and the components of the telescope optical systems, including the segmented mirror with its support structure, the mirror facets for each telescope and the Winston cone light concentrators in front of the PMT camera can be found in Bernlöhr et al. (2003), and a description of the mirror alignment is given in Cornils et al. (2003). For details on the camera calibration see Aharonian et al. (2004). The trigger system is described in Funk, Hermann, Hinton et al. (2004). Early reports of H.E.S.S. have been given elsewhere (see e.g. Hofmann (2002)).
#### 2.1.2 Observations
In the beginning of the observation period of October 2003, a single 5 $`\sigma `$ detection by H.E.S.S. achieved in 1h, triggered an approved RXTE ToO (Target of Opportunity) on this target. The results presented here are based on observations carried out in 2003 between October 19 and November 26.
The data were taken in the Wobble mode where the source direction is positioned $`\pm 0.5^{}`$ in declination relative to the centre of the field of view of the camera during observations. This allows for both on-source observations and simultaneous estimation of the background induced by charged cosmic rays. The data reported here are selected and analyzed with the “standard analysis” described in §4 of AH04. The background is estimated here by using all events passing cuts in a ring around the source, as described in section 4.3 in AH04. The runs passing the quality selection criteria total 32.4 hours of livetime on the source. The total two-dimensional significance sky map is shown in Figure 1, along with a graph showing the $`\theta ^2`$ distribution (where $`\theta ^2`$ is the square of the angular difference between the reconstructed shower position and the source position) of the 1764 excess events observed. This yields a detection at the $`34.3\sigma `$ level, at the average rate of $`0.91\gamma /\mathrm{min}`$ and a significance of $`6.0\sigma /\sqrt{\mathrm{h}}`$.
The methods used here for reconstructing the energy of each event and for determining a spectrum are described in §6 of AH04. The measured time-average spectrum is fitted by a power law of the form $`\mathrm{d}N/\mathrm{d}E=I_0(E/1\mathrm{TeV})^\mathrm{\Gamma }`$ with $`I_0`$ the flux normalization at $`1\mathrm{TeV}`$ and $`\mathrm{\Gamma }`$ the photon index. The photon index obtained from the time-averaged spectrum is then used as a fixed parameter to estimate the integral flux above $`300\mathrm{GeV}`$ for each run. This integrated flux takes into account the effective area and threshold variations due to the source moving through the sky, giving more reliable variability information than counting rates in units of $`\gamma `$-rays/minute. Overall systematic errors are estimated to be 20% for the integral flux and $`0.1`$ for the photon index.
### 2.2 RXTE
The PCA (Jahoda et al. (1996)) units of RXTE observed PKS 2155$``$304 between October 22 and November 23 of 2003 with exposures of typically $`10\mathrm{ks}`$ in October and $`1\mathrm{ks}`$ in November. The STANDARD2 data were extracted using the ftools in the HEASOFT 5.3.1 analysis software package provided by NASA/GSFC and filtered using the RXTE Guest Observer Facility (GOF) recommended criteria. The changing Proportional Counter Units (PCU0/2/3) configuration throughout the observations was taken into account in the data reduction, and only the signals from the top layer (X1L and X1R) were used. When reducing the PCUs individually to establish the time-averaged spectrum, the average spectral fit parameters are similar within error bars but with a systematically higher $`\chi ^2`$ for spectral fits performed on PCU0 data alone. On shorter timescales this effect becomes negligible and PCU0 contributes to the statistical significance of the flux measurement. Therefore all PCUs were kept for the overall light curve which is binned in $`400\mathrm{s}`$ bins, but only PCU2 and PCU3 are used for the analysis of data segments that are simultaneous with H.E.S.S. runs and for the time-averaged spectrum.
The average spectrum used in the SED is derived by combining PCU2 and PCU3 spectra using the addspec tool weighted by the counts information delivered by fstatistic and then the corresponding response matrices were combined with addrmf. The faint-background model was used and only the 3–40 PHA channel range was kept in XSPEC v. 11.3.1, or approximately 2–$`20\mathrm{keV}`$.
To build a light curve in units of integrated flux in the 2–$`10\mathrm{keV}`$ band, spectral data were derived from $`400\mathrm{s}`$ bins to probe short timescales with adequate statistical accuracy. These segments are then fitted by a power law in XSPEC with PCU configuration-dependent response matrices generated by the ftool pcarsp v. 10.1 and a fixed column density of $`N_\mathrm{H}=1.7\times 10^{20}\mathrm{cm}^2`$ obtained from PIMMS<sup>1</sup><sup>1</sup>1See http://legacy.gsfc.nasa.gov/Tools/w3pimms.html. This yields the flux and the error (corresponding to the $`1\sigma `$ confidence level) on the flux reported in the light curves in Figure 2, in units of $`10^{11}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ in the 2–$`10\mathrm{keV}`$ band. The fits did not improve by using a broken power-law for the $`400\mathrm{s}`$ binned observations.
### 2.3 ROTSE
The ROTSE-III array is a worldwide network of four $`0.45\mathrm{m}`$ robotic, automated telescopes built for fast ($`6\mathrm{s}`$) response to GRB triggers from satellites such as HETE-2 (High Energy Transient Explorer 2) and Swift. The ROTSE-III telescopes have a wide ($`1.85^{}\times 1.85^{}`$) field of view imaged onto a Marconi $`2048\times 2048\mathrm{pixel}`$ back-illuminated thinned CCD and are operated without filters. The ROTSE-III systems are described in detail in Akerlof et al. (2003). At the time of the observations of PKS 2155$``$304 in October and November 2003, two ROTSE-III telescopes were operational in the Southern hemisphere: ROTSE-IIIa located at the Siding Spring Observatory, Australia and ROTSE-IIIc at the H.E.S.S. site. The ROTSE-IIIc telescope is located in the centre of the H.E.S.S. telescope array. A 30% share of the total observation time is available to the H.E.S.S. collaboration, which has been used to perform an automated monitoring programme of blazars, including objects that are being observed with the H.E.S.S. telescopes. Both telescopes participated in the observation campaign on PKS 2155$``$304 in October and November 2003.
The telescopes observed PKS 2155$``$304 typically 10 times per night taking sequences of 2 frames with $`60\mathrm{s}`$ exposures with a slight dithering of the pointing to reduce the impact of individual noisy pixels. The typical limiting magnitude, depending on the sky conditions, is $`18.5^{\mathrm{mag}}`$. Overall, 323 bias-subtracted and flat-fielded frames have passed visual inspection and are used to produce a light curve. A total of 6 frames were rejected due to the presence of stray light from Jupiter.
Using an overlay of 50 isolated comparison stars with similar brightness ($`12^{\mathrm{mag}}14^{\mathrm{mag}}`$) and co-located with PKS 2155$``$304 ($`<15\mathrm{arcmin}`$), a two-dimensional Gaussian is fit to the intensity distribution characterising the point-spread function by $`\sigma _{\mathrm{psf}}`$.
To estimate the local sky-background for the reference stars and the target object, an annulus with inner radius $`2\times \sigma _{\mathrm{psf}}`$ and an outer radius $`6\times \sigma _{\mathrm{psf}}`$ is chosen. Based upon a reference frame which is derived from co-adding 30 individual frames, a mask is calculated for each object excluding regions where faint objects coincide with the annulus; pixels exceeding 3 standard deviations of the local sky background are excluded. Using the local sky background, the intensity and error of each object is calculated. Using the 50 reference stars, a relative intensity and statistical error with respect to a reference frame is calculated.
The absolute flux values are obtained by calculating a relative $`R`$ magnitude by comparing the instrumental magnitude with the USNO catalogue as described in Akerlof et al. (2000). The procedure has been checked by comparing the average $`R`$ magnitude of a sample of 70 BL Lac type objects determined with ROTSE observations carried out over one year of operation with the $`V`$ magnitude listed in the $`10^{\mathrm{th}}`$ Veron Cetty & Veron catalogue of BL Lac type objects. The average $`VR`$ of 0.5 that is found is consistent with the average value for $`VR`$ obtained from cross-checking the colours with the 2MASS catalogued value for the BL Lac type objects.
Finally, the host galaxy has been resolved in optical (Falomo (1996)) and NIR (Kotilainen, Falomo, & Scarpa (1998), KFS98 hereafter) and found to be a typical giant elliptical with $`\mathrm{M}(R)=24.4`$ which translates into an apparent $`\mathrm{m}(R)=15.1`$ (here the distance moduli given by KFS98 have been used to calculate the apparent magnitude based upon the absolute magnitude quoted). The ROTSE measurements have as maximum and minimum $`\mathrm{m}(R,\mathrm{min})=13.3`$ and $`\mathrm{m}(R,\mathrm{max})=13.7`$ which corresponds to $`10\mathrm{mJy}`$ for the maximum observed flux and $`6.7\mathrm{mJy}`$ for the minimum flux taking the contribution of the host galaxy into account. These values are considerably lower than the retrieved archival data indicating that PKS 2155$``$304 was in a low state at the moment of the observations.
### 2.4 NRT
The Nançay radiotelescope is a single-dish antenna with a collecting area of $`200\times 34.56`$ m<sup>2</sup> equivalent to that of a $`94\mathrm{m}`$-diameter parabolic dish (van Driel et al. (1996)). The half-power beam width at $`11\mathrm{cm}`$ is $`1.9\mathrm{arcmin}`$ (EW) $`\times 11.5\mathrm{arcmin}`$ (NS) (at zero declination), and the system temperature is about $`45\mathrm{K}`$ in both horizontal and vertical polarizations. The point source efficiency is $`0.8\mathrm{K}\mathrm{Jy}^1`$, and the chosen filter bandwidth was 12.5 MHz for each polarization, split into two sub-bands of 6.25 MHz each. Between 4 and 14 individual 1-minute drift scans were performed for each observation, and the flux was calibrated using a calibrated noise diode emission for each drift scan. Data processing has been done with the Nançay local software NAPS and SIR.
A monitoring programme with this telescope on extragalactic sources visible by both the NRT and H.E.S.S. is in place since 2001. For the campaign described here it consisted of a measurement at $`11\mathrm{cm}`$ every two or three days. The average flux for the 8 measurements in October and November 2003 was $`0.30\pm 0.01\mathrm{Jy}`$ with possible marginal variability.
## 3 Results
### 3.1 Light curves
The October and November 2003 light curves of all the H.E.S.S., RXTE and ROTSE observations are shown in Figure 2. The H.E.S.S. light curve is binned in run-length times averaging 28 minutes each. The flux is in units of $`10^{11}\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1`$ above $`300\mathrm{GeV}`$, derived using the average photon index 3.37 obtained in section 3.3. Spectra could not be derived on a run-by-run basis due to the weak signal. As for the observations reported in AH04, the overall light curve is inconsistent with a constant flux. A $`\chi ^2`$ fit of the data to a constant yields a $`3.4\times 10^{10}`$ $`\chi ^2`$ probablility. The intra-night VHE flux on MJD52936 (Fig. 3) exhibits an increase of a factor of $`1.9\pm 0.6`$ in $`0.11\mathrm{d}`$. On MJD52932 the peak-to-peak flux shows an increase of a factor of $`2.5\pm 0.9`$ within $`0.09\mathrm{d}`$. These timescales are longer than the 30 minute doubling time reported in AH04. For these two extreme cases of VHE variability observed during this campaign only the second had a limited RXTE coverage.
The 2–$`10\mathrm{keV}`$ X-ray flux in this campaign ranges from $`F_{210\mathrm{keV}}=2.0\times 10^{11}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ to $`4.4\times 10^{11}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. The maximum is lower than the 20 November 1997 measurement of $`2.3\times 10^{10}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ (Vestrand & Sreekumar (1999)) indicating that the X-ray state seen here is not exceptionally high. The minimum seen here is consistent with historically low fluxes (Zhang, Treves et al. (2002)). The intra-night variability is also obvious here, but no flare was completely resolved.
A 60% flux variability in $`t_{\mathrm{var}}1.5\mathrm{ks}`$ on MJD52936 is the best marked transient episode in the observations reported here (bottom panel b) in Figure 3 for which the H.E.S.S. observations were made at the end of the transit inducing a large associated error on the flux estimation due to the high zenith angle of the source. This timescale is comparable to those reported by Gaidos et al. (1996) where doubling times as short as 15 min from Mkn 421 were observed in the VHE band.
This flare is the fastest rise seen in this object to date since BeppoSAX saw a $`5\times 10^4\mathrm{s}`$ rise timescale (Zhang, Treves et al. (2002)) and Kataoka et al. (2000) observed a doubling timescale of $`3\times 10^4\mathrm{s}`$ with the ASCA satellite. So far the fastest rise in this type of object was observed in Mkn 501 with a 60% increase in less than 200 s (Catanese & Sambruna (2000)) though Xue & Cui (2004) claims that this flare is likely to be an artifact.
The optical emission of PKS 2155$``$304 is dominated by the nucleus which outshines the host galaxy by a factor of $`4`$ given in KFS98. The observed variability amplitude is therefore not biased by the constant emission of the host galaxy and mainly due to the activity of the nucleus. The peak-to-peak amplitude of variability is moderate compared to the variability amplitude at shorter wavelengths and typically $`0.1^{\mathrm{mag}}`$ peak-to-peak. The object has been monitored over longer time-scales with the ROTSE-IIIc telescopes showing variations with amplitudes close to $`1^{\mathrm{mag}}`$.
### 3.2 Correlation analysis
In order to quantitatively look for correlated variability between the VHE, X-ray and optical bands, the measured fluxes are plotted against each other in Figure 4 for all the observations carried out during this campaign. For correlated VHE/X-ray variability, the RXTE analysis was slightly modified: only observation segments that happen exactly within a H.E.S.S. run are reduced and analyzed (using only PCU2, PCU3 and combinations thereof). This provides 23 simultaneous data segments for the whole campaign for which the fluxes are represented in Figure 4. There is no obvious correlation for those observations (correlation factor $`r=0.37\pm 0.13`$). For correlated optical/VHE variability, ROTSE observations that happen within a H.E.S.S. run are averaged and their errors summed quadratically. No correlation ($`r=0.24\pm 0.27)`$ is found for these observations. Also no correlation was found between the optical and X-ray band ($`r=0.02\pm 0.05)`$.
### 3.3 Spectra
The October and November 2003 H.E.S.S. data were all combined for the spectrum shown in Figure 5. The best-fitting power law ($`\chi ^2=7`$ for 8 degrees of freedom) is given by:
$$\begin{array}{c}\frac{\mathrm{d}N}{\mathrm{d}E}=(2.73\pm 0.17)\times 10^{12}\left(\frac{E}{\mathrm{TeV}}\right)^{3.37\pm 0.07\pm 0.10}\hfill \\ \hfill \mathrm{cm}^2\mathrm{s}^1\mathrm{TeV}^1\end{array}$$
(1)
which is comparable to $`\mathrm{\Gamma }=3.32\pm 0.06`$ and $`I_0=(1.96\pm 0.12)\times 10^{12}\mathrm{cm}^2\mathrm{s}^1\mathrm{TeV}^1`$ previously reported in AH04.
The result of the broken power law fit for the combined RXTE PCU2 and PCU3 spectrum is shown in Figure 6. It yields an unabsorbed flux in the 2–$`10\mathrm{keV}`$ band of $`F_{210\mathrm{keV}}=(2.66\pm 0.04)\times 10^{11}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ ($`\chi ^2=41`$, 31 degrees of freedom), a lower index of $`\mathrm{\Gamma }_L=2.81\pm 0.05`$, a break energy of $`E_b=4.9\pm 0.8\mathrm{keV}`$ and a higher index of $`\mathrm{\Gamma }_H=2.95\pm 0.04`$. A single power law fit to the same data yields $`F_{210\mathrm{keV}}=(2.69\pm 0.03)\times 10^{11}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ ($`\chi ^2=51`$, 34 degrees of freedom) and an index of $`\mathrm{\Gamma }=2.88\pm 0.13`$, a poorer fit than the broken power law, but the index still provides information that can be used for comparison with historical measurements. Indeed, the derived index is close to those measured by the BeppoSAX satellite (Giommi, Fiore, Guainazzi et al. (1998)), GINGA (Sembay et al. (1993)) and well within the range observed by EXOSAT (Treves et al. (1989)). The statistics above $`10\mathrm{keV}`$ in our RXTE observations are too poor to check the existence of a possible hard tail above $`20\mathrm{keV}`$ (Giommi, Fiore, Guainazzi et al. (1998)) which might be the signature of the onset of a high-energy component.
In order to look for flux dependent spectral variability, the RXTE data subset used in Figure 3 is divided into two energy bands, the PHA channels 0–9 (soft band) and 10–27 (hard band), corresponding to approximately 1–4 keV and 4–11 keV, respectively. A hardness ratio (HR), shown in Figure 7, is the ratio of the counting rate in the hard band over the soft band. There is a clear correlation of the HR with the rate, peaking when the rate is highest. The correlation factor between the rate and the HR is $`r=0.76\pm 0.12`$. Even though the variability timescale here is much smaller, this behavior is compatible with the hardening reported in Chiappetti et al. (1999).
### 3.4 EBL corrected spectrum
For objects at non-negligible redshifts, the large, energy-dependent opacities can cause the emitted spectrum to be greatly modified both in shape and intensity (see e.g. Stecker et al. (1992); Biller (1995); Coppi & Aharonian (1999); Vassiliev (2000)). Unfortunately, at present the knowledge of the EBL still has large uncertainties, for both direct measurements and models, as summarized in Primack et al. (2001). In order to estimate the intrinsic VHE spectrum, and thus to locate the Inverse Compton (IC) peak of the blazar’s SED, we have used three EBL models (Fig. 8) as representatives of three different flux levels for the stellar peak component. This is the EBL energy range which mostly affects the H.E.S.S. spectrum: with data up to 3–$`3.5\mathrm{TeV}`$, the peak of the $`\gamma \gamma `$ cross section is reached for soft photons with wavelengths $`5\mu \mathrm{m}`$.
The three models used here are (in order from higher to lower fluxes): the phenomenological shape used in Aharonian et al. (2003), which is based on the original Primack et al. (2001) calculation but smoothed and scaled up to match the data points below $`1\mu \mathrm{m}`$ and at 2–$`3.5\mu \mathrm{m}`$ (hereafter Phigh); the original Primack et al. (2001) calculation for a Salpeter initial mass function (hereafter Primack01); and the new 2004 calculation (Primack et al. (2004), hereafter Primack04), which takes advantage of the recent improvements in the knowledge of the cosmological parameters and of the local luminosity function of galaxies.
The opacities are calculated from the EBL SED shapes taking into account only the cosmology ($`H_0=70\mathrm{km}/\mathrm{s}/\mathrm{Mpc}`$, $`\mathrm{\Omega }_{\mathrm{Mat}}=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$). To treat all three shapes similarly, no evolution has been introduced at this point. This corresponds to a “maximum absorption” hypothesis (i.e., for increasing $`z`$, constant instead of decreasing EBL comoving energy density). But at these redshifts ($`0.1`$) and, for example, assuming the evolution given in Primack et al. (2001), the difference is still small (variation in the photon index $`\mathrm{\Delta }\mathrm{\Gamma }<0.1`$ in the range 0.3–$`1\mathrm{TeV}`$), and negligible compared to the differences between models.
The resulting absorption-corrected spectra are shown in Figure 9, together with the observed spectrum. The intrinsic spectra are all well fitted by a single power-law model, with a hard spectrum for Phigh ($`\mathrm{\Gamma }1.5`$), and soft spectra for Primack01 and Primack04 models ($`\mathrm{\Gamma }2.3`$ and $`\mathrm{\Gamma }2.8`$, respectively). This effect is directly related to the different flux levels of the stellar peak component, which imprint a different amount of softening onto the original spectrum. This direct link thus yields two simple scenarios for the location of the blazar’s high energy peak, with the dividing line represented by the EBL flux which gives $`\mathrm{\Gamma }=2.0`$ (1.3 times the Primack01 model). Models with stellar peak fluxes above this (such as Phigh, and generally all those in agreement with the direct estimates of the fluxes between 2 and $`3.5\mu \mathrm{m}`$) imply a hard intrinsic spectrum ($`\mathrm{\Gamma }<2`$), and thus an IC peak above 1–$`2\mathrm{TeV}`$. EBL models with lower fluxes (such as the Primack01 and Primack04) imply instead a soft spectrum ($`\mathrm{\Gamma }>2`$), locating the IC peak below the observed energy range ($`<200\mathrm{GeV}`$). In the following, we will discuss both scenarios for the SED modelling, using the Phigh and Primack2004 curves as the two ends of the possible range of values for the “Primack-type” shape (i.e., between the claimed EBL direct measurements at few microns and the lower limits from galaxy counts).
## 4 SED Modelling
The broadband spectral morphology of PKS 2155$``$304 is typical of the BL Lac type, with a double-humped structure in $`\nu F_\nu `$ representation, exhibiting a low-energy and a high-energy component. Its broadband emission is usually attributed to emission from a beamed relativistic jet, oriented in a direction close to the line of sight (Begelman, Blandford & Rees (1984); Blandford & Königl (1979)). The spectral energy distribution in units of power per logarithmic bandwidth $`E^2\mathrm{d}N/\mathrm{d}E`$ versus energy $`E`$ is shown in Figure 10. The EGRET measurements, between the H.E.S.S. and RXTE points, are from the third EGRET catalog (Hartman (1999)) and from a very high $`\gamma `$-ray state described in Vestrand, Stacy & Sreekumar (1995). There is a difference in spectral states, since in the former case the power law photon index is $`1.71\pm 0.24`$ whereas it is $`2.34\pm 0.26`$ in the latter which most likely consists of a mix of low and high activity state observations. The historical EGRET spectra are therefore unlikely to represent the state of PKS 2155$``$304 during the campaign presented here and are not used to put stringent constraints on the modelling. Considering the archival data and the steep X-ray spectrum in Figure 10, the peak of the low-energy component occurs in the 2–$`2000\mathrm{eV}`$ range. The archival BeppoSAX data from a high state analyzed by Chiappetti et al. (1999), and represented here above our RXTE data, show a peak at $``$ 0.1 keV. The absorbed VHE peak location is clearly below $`300\mathrm{GeV}`$, with its exact location depending on the spectrum in the EGRET range.
Whereas the current models seem to agree that the low-energy component is dominated by synchrotron radiation coming from nonthermal electrons emitted in collimated jets, the high-energy emission is assumed to be either inverse Compton scattering off the synchrotron photons (Synchrotron Self-Compton, SSC, see e.g. Maraschi et al. (1999); Bicknell & Wagner (2002)) or by external photons (see e.g. Sikora et al. (1994)). This kind of leptonic model will be discussed in section 4.2. A hadronic origin of the VHE emission using the Synchrotron-Proton Blazar (SPB) model with a dominating proton synchrotron component at high energies in a proton-electron plasma is also able to produce a double humped SED and is discussed in section 4.3. The lack of correlation between the RXTE and H.E.S.S. fluxes (and possibly also between the optical and the VHE emission) within the small variability range may indicate a different spatial origin, or a different underlying particle distribution. In the proton synchrotron model a lack of correlated variability between $`\gamma `$-ray emission and the low energy electron synchrotron component could arise if the electrons and protons are not co-accelerated.
The high-energy component above $``$ 100 GeV is attenuated by interactions with the EBL and is a lower limit for the intrinsic spectrum. The energy budget in X-rays and VHE $`\gamma `$-rays is comparable, though the maximum output at the peak energy in the high-energy component is likely to be lower than that in the lower-energy component. Interpolating between the high-state EGRET archival data and VHE data would lead to a maximum located above $`10\mathrm{GeV}`$, which is surprising since the observations reported here indicate a low state. Extrapolating the EGRET catalogue spectrum to VHE energies with a power law falls below the H.E.S.S. data and therefore requires two inflexion points in the SED. Simultaneous observations in the MeV-GeV range with the upcoming satellite GLAST will be crucial to constrain the high-energy component shape.
### 4.1 Doppler boosting and synchrotron/Compton derived parameters
The electromagnetic emission in blazars is very likely to be Doppler-boosted (or beamed) toward the observer. In the radio regime, the evidence for Doppler boosting in PKS 2155$``$304 comes from superluminal expansions observed with VLBI (Piner & Edwards (2004)). Relativistic beaming is also required in order to avoid absorption of GeV photons by X-ray photons via the $`\mathrm{e}^+/\mathrm{e}^{}`$ pair-production process (see e.g. Maraschi, Ghisellini & Celotti (1992)). It is thus possible to use the $`\gamma `$-ray variability to establish a limit for the Doppler factor $`\delta `$, with $`\delta `$ defined in the standard way as $`[\mathrm{\Gamma }(1\beta \mathrm{cos}\theta )]^1`$, where $`\mathrm{\Gamma }`$ is the bulk Lorentz factor of the plasma in the jet, $`\beta =v/c`$, and $`\theta `$ is the angle to the line of sight.
Following Dondi & Ghisellini (1995), the size of the $`\gamma `$-ray emission zone $`R`$ is derived from the time-scale of variability $`t_{\mathrm{var}}`$ (supposing the timescale of intrinsic variability is negligible compared to the light crossing time) by $`Rct_{\mathrm{var}}\delta /(1+z)`$. In this case (assuming a time scale $`t_{\mathrm{var}}2`$ ks from section 3) the size of the emission region is
$$R\delta ^15\times 10^{13}\mathrm{cm}.$$
(2)
Since for $`\gamma `$-rays in the $`300\mathrm{GeV}`$$`3\mathrm{TeV}`$ range the opacity $`\tau _{\gamma \gamma }`$ cannot significantly exceed 1, independently of the emission mechanisms, another constraint on the minimum value of the Doppler factor can be derived by estimating how much Doppler boosting is necessary for photons with observed energy $`E_\gamma `$ to escape from a source with radius $`R`$ and a flux density $`F(E_t)`$, where $`E_t=(m_ec^2)^2\delta ^2/E_\gamma (1+z)^2`$. At this point one can follow Mattox et al. (1993) and especially Eq. 3.7 in Dondi & Ghisellini (1995), writing that the opacity is
$$\tau _{\gamma \gamma }\frac{\sigma _\mathrm{T}}{5}\frac{1}{hc}\frac{d_L^2}{R}\frac{1}{\delta ^3(1+z)}F(E_t).$$
(3)
Imposing that $`\tau _{\gamma \gamma }<1`$ for photons with observed energy $`E_\gamma =1\mathrm{TeV}`$ yields a lower limit on $`\delta `$ for given $`R`$ that can be derived numerically from the observed SED. For PKS 2155$``$304 this yields
$$R^1\delta ^{6.4}5.6\times 10^{24}\mathrm{cm}^1.$$
(4)
The minimum allowable boost $`\delta `$ comes from combining Eq. 2 and Eq. 4 and yields
$$\delta 19$$
(5)
which is higher than the limit obtained in a similar way in Tavecchio et al. (1998), but at the lower end of the range usually obtained from SSC modelling as in Kataoka et al. (2000). It is useful to stress that this constraint is valid under the usual assumption that the region emitting the SED optical flux (i.e., the target photons) is cospatial with (or at least embeds) the high energy emitting region. The models used later in this paper are therefore essentially single zone models.
If the observed X-rays are synchrotron radiation from nonthermal electrons then the mean observed energy $`<E>`$ of an electron with Lorentz factor $`\gamma _\mathrm{e}`$ is given by
$$<E>\frac{\delta }{1+z}\frac{21\mathrm{}\gamma _\mathrm{e}^2qB}{15\sqrt{3}m_e}.$$
Using $`<E>=10\mathrm{keV}`$ here yields
$$B\delta \gamma _\mathrm{e}^2=1.1\times 10^{12}\mathrm{G}.$$
(6)
In the SSC scenario the same electrons can Comptonize ambient photons up to the VHE regime, thus
$$\frac{\delta }{1+z}\gamma _\mathrm{e}m_ec^23\mathrm{TeV}.$$
(7)
Combining Eq. 6 and Eq. 7 yields
$$B\delta ^11.1\times 10^{12}\mathrm{G}\times \left(\frac{3\mathrm{TeV}}{\mathrm{m}_\mathrm{e}\mathrm{c}^2}\right)^2(1+\mathrm{z})^2.$$
(8)
A numerical application yields
$$B\delta ^10.03\mathrm{G}$$
(9)
and hence the Lorentz factor is constrained to $`\gamma _\mathrm{e}1.3\times 10^6`$ which is higher than that calculated as above for Mkn 421 by Takahashi et al. (1996) who derived $`\gamma _\mathrm{e}>5\times 10^5`$ (and $`B0.2`$ G) but lower than $`\gamma _\mathrm{e}10^7`$ found in a similar way in 1ES1959$`+`$650 (Giebels et al. (2002)).
The X-ray data presented above imply that the X-ray spectrum of PKS 2155$``$304 hardens as the source brightens. This is often measured in BL Lac objects; a hardening of the spectrum when flares occur, and a blueward shift of the peak of the synchrotron emission $`\nu _{\mathrm{sync}}`$ (and presumably higher energy inverse-Compton emission) by factors that can be as large as 100 were measured in the cases of Mkn 501 (Pian et al. (1998)), 1ES 1426$`+`$428 and PKS 0548$``$322 (Costamante et al. (2001)). In the case of PKS 2005$``$489 (Perlman et al. (1999)), a more moderate shift of a factor of 3 or less of the synchrotron emission was found. The archival data suggest that $`\nu _{\mathrm{sync}}`$ lies in the UV band for PKS 2155$``$304, but no data were taken simultaneously in this campaign at that wavelength.
The lack of indication for correlated X-ray/VHE variability does not imply that PKS 2155$``$304 behaves differently from VHE sources such as Mkn 421, for which VHE/X-ray correlation has been established on a much higher variability basis (see e.g. Cui et al. (2004)) with dynamical ranges of 30 in both energy bands. Limiting those observations to the same dynamical range as observed here would not allow any claim for correlation. Future observations of PKS 2155$``$304 with a higher variability amplitude would bring more insight into this.
### 4.2 Leptonic interpretation
Interpretation with a single zone SSC model of the SED of PKS 2155$``$304 has already been proposed in the literature using two different assumptions. In Kataoka et al. (2000) the low energy tail of the SSC model is used to account for the low-energy component up to the optical in the SED. That component is decomposed into two sub-components by Chiappetti et al. (1999) where the radio to optical emission has another origin than the X-rays, which are assumed to come from the jet. These two different interpretations are used here in the context of the leptonic model described in Katarzyński et al. (2001) which has already been applied to Mkn 501 and Mkn 421 (Katarzyński et al. (2003)). To constrain this model, only the simultaneous data are used, since the archival data reported in the SED of Figure 10 are likely to not represent the state of this source (note for example the difference in optical flux and the ROTSE measurement).
When using the Primack04 absorption, the model used here can reproduce the X-ray through VHE part of the SED, but the H.E.S.S. spectrum constrains it such that the radio measurement can not be included in the synchrotron bump predicted by the single-zone model. As for Mkn 421 and Mkn 501, adding a more extended component than the VHE emitting zone can provide an explanation for this. The origin is probably the compact VLBI core which has a radio core to lobe ratio of $`1`$ (Laurent-Muehleisen et al. (1993); Piner & Edwards (2004)) and a typical size of $`10^{18}\mathrm{cm}`$, more than two orders of magnitude larger than the VHE emitting zone. This VLBI feature dominates the spectrum at low energy and is included in the SED modelling here. An uncertainty remains which is the high frequency cutoff of this VLBI component. The host galaxy contribution to the optical flux is estimated to be $`10^{11}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, deduced from the magnitudes given by KFS98 and assuming a low-redshift solar metallicity elliptical galaxy of age equal to $`13\mathrm{Gyr}`$ ($`RH=2.4`$), corresponding to a mass of $`5\times 10^{11}M_{}`$ (Fioc & Rocca-Volmerange (1997)). So even at the measured low activity state of PKS 2155$``$304 the host galaxy is not contributing much in the optical range.
The ROTSE measurement can be ascribed here to either the high-energy tail of the VLBI component or to the synchrotron part of the SSC model. Assuming a common origin for the X-rays and the optical emission, and using a variability time scale of 0.1 d to constrain the emitting zone, the model tends to predict a high IC flux, as shown in Figure 10. However, the lack of correlation between the X-rays and the optical emission in our measurements - also suggested by Dominici et al. (2004) based on less sensitive RXTE/ASM measurements - indicates that the optical emission could originate from the VLBI component, which is modelled by a slight increase in the maximal Lorentz factor of the emitting electrons. This in turn lessens the constraint on the simultaneous SSC fit of the X-ray and VHE part and allows a better fit of the VHE spectrum for smaller sizes of the emitting zone. Detailed parameters of the two hypotheses are given in Table 1.
If the absorption correction is well described by the Phigh model, the slope of H.E.S.S. data at high energy implies that the peak of the TeV emission bump is located above $`4\mathrm{TeV}`$ (or $`10^{27}\mathrm{Hz}`$). Such a high frequency peak emission imposes a strong constraint for the single-zone SSC scenario, especially since the peak of the synchrotron bump has to remain below $`1\mathrm{keV}`$ (or $`10^{17.5}\mathrm{Hz}`$) as required by the slope of the RXTE data. High values of both the jet Doppler factor and the maximal Lorentz factor of radiating particles are required to reach the necessary energy for the IC bump, that is $`\gamma _{\mathrm{max}}\delta >8\times 10^6`$. On the other hand, to keep the synchrotron peak below $`1\mathrm{keV}`$ imposes an upper limit to the magnetic field. Within these constrains, the best fit we obtain is shown in Figure 11. We can note that none of the high energy tails are well accounted for. The set of parameters for the best fit is given in Table 1. This fit marginally reproduces the observed X-ray and $`\gamma `$-ray data, but is not as satisfactory than that obtained with the Primack04 absorption correction, and in any case it is impossible to take into account the ROTSE optical point. The main changes in parameters between the two fits consist in enhancing the boosting, which then becomes quite extreme, while reducing the density and magnetic field for the Phigh absorption correction.
The constraints on $`\delta `$ derived here from either simple opacity arguments or from the one-zone model parametrisation of the SED are in the range of Doppler factors usually derived with such assumptions or models for other VHE emitters. As pointed out by Chiaberge et al. (2000) such high values are however at odds with attempts to unify the BL Lac population with the family of FR I sources (Urry & Padovani (1995)), the latter being possibly an unbeamed since off-axis viewed case of the former. The same authors suggest that models where velocity structures in the jet, such as the “spine-sheath” model (see e.g. Sol, Pelletier & Asseo (1989)) or the decelerating flow model (Georganopoulos & Kazanas (2003, 2004)) allow lower bulk Lorentz factors. Another option from Pelletier et al. (2004) is to cope with the pair creation catastrophe implied by smaller Doppler factors in their “two-flow” solution. Comparing the SED with such models, which make the BL Lac \- FR-I connection more plausible, is an interesting task but beyond the scope of this paper.
### 4.3 Hadronic models
Generally, the leptonic models constitute the preferred concept for TeV blazars, essentially because of two attractive features: (i) the capability of the (relatively) well understood shocks to accelerate electrons to TeV energies (Sikora & Madejski (2001); Pelletier (2001)) and (ii) the effective conversion of the kinetic energy of these relativistic electrons into the X-ray and VHE $`\gamma `$-ray emission components through the synchrotron and inverse Compton radiation channels. The hadronic models are generally lack these virtues. They assume that the observed $`\gamma `$-ray emission is initiated by accelerated protons interacting with ambient matter (Bednarek (1993); Dar & Laor (1997); Pohl & Schlickeiser (2000)), photon fields (PIC model, Mannheim (2000)), magnetic fields (Aharonian et al. (2002)) or both (Mücke & Protheroe (2000)).
The models of TeV blazars involving interactions of protons with photon and B-fields require particle acceleration to extreme energies exceeding $`10^{19}\mathrm{eV}`$ which is possible if the acceleration time is close to $`t_{\mathrm{acc}}=r_\mathrm{g}/c`$ ($`r_\mathrm{g}`$ is the gyro-radius). This corresponds (independent of a specific acceleration mechanism) to the maximum (theoretically possible) acceleration rates (Aharonian et al. (2002)) which can only be achieved by the conventional diffusive shock acceleration in the Bohm diffusion regime.
On the other hand, the condition of high efficiency of radiative cooling of accelerated particles requires extreme parameters characterizing the sub-parsec jets and their environments, in particular very high densities of the thermal plasma, radiation and/or B-fields. In particular, the proton-synchrotron models of TeV blazars require highly magnetized ($`B10\mathrm{G}`$) condensations of $`\gamma `$-ray emitting clouds containing Extremely High Energy (EHE) protons, where the magnetic pressure dominates over the pressure of relativistic protons (Aharonian (2000)).
Below we use the hadronic SPB model (Mücke & Protheroe (2000, 2001)) to model the average spectral energy distribution (SED) of PKS 2155$``$304 in October-November 2003. A detailed description of the model, and its implementation as a (time-independent) Monte-Carlo/numerical code, has been given in Mücke, Protheroe et al. (2003) and Reimer et al. (2004).
Considering the rather quiet activity state of PKS 2155$``$304 in Oct-Nov. 2003, we use the 3EG catalog spectrum, since it is the best determined EGRET spectrum from this source to date, as an upper limit for modelling purposes.
Flux variability provides an upper limit for the size of the emission region. To allow for a comparative study between leptonic and hadronic models we fix here the comoving emission region to $`Rct_{\mathrm{var}}\delta =5\times 10^{13}\delta \mathrm{cm}`$ deduced from the X-ray variability. We assume that the optical through X-ray emission and the $`\gamma `$-ray output stem from the same region of size $`R`$.
A reasonable model representation for the simultaneous data assuming a Primack04 model for the VHE absorption is found for the following parameters: magnetic field $`B=40\mathrm{G}`$, Doppler factor $`\delta =20`$, injection electron spectral index $`\alpha _\mathrm{e}=1.6`$, assumed to be identical to the injection proton spectral index $`\alpha _\mathrm{p}`$, maximum proton energy of order $`\gamma _{\mathrm{p},\mathrm{max}}4\times 10^9`$, e/p-ratio of 0.15 and a near-equipartition proton energy density of $`u_\mathrm{p}=27\mathrm{erg}\mathrm{cm}^3`$. The required total jet power is of the order $`L_{\mathrm{jet}}1.6\times 10^{45}\mathrm{erg}\mathrm{s}^1`$. When using the Phigh EBL model, a reasonable representation of the data may be achieved by increasing the maximum injected proton energy to $`\gamma _{\mathrm{p},\mathrm{max}}=10^{10}`$ and simultaneously increasing the e/p-ratio to 0.24, while all other parameters remain unchanged. Note that here the maximum proton gyro-radius approaches the size of the emission region. Alternatively, a doubling of the magnetic field to $`80\mathrm{G}`$ together with an increase of $`\gamma _{\mathrm{p},\mathrm{max}}`$ to $`8\times 10^9`$ and a e/p-ratio of unity (leading to $`u_B50u_\mathrm{p}`$) represents the SED-data equally well. In conclusion, none of the ”Primack-type” EBL models can explicitly be ruled out in the framework of the SPB-model by the H.E.S.S. data presented here. In all cases, proton synchrotron emission dominates the (sub-)TeV radiative output. Depending on the Doppler factor, part of the proton synchrotron radiation produced may be reprocessed to lower energies. Contributions from the muon and pion cascades are always lower than the proton synchrotron component. The low energy component is dominated by synchrotron radiation from the primary electrons, with a negligible contribution of synchrotron radiation from secondary electrons (produced by the $`\mathrm{p}`$\- and $`\mu ^\pm `$-synchrotron cascade).
On the other hand, the synchrotron radiation of secondary electrons resulting from interactions of VHE $`\gamma `$-rays with external low-energy photons with a modest $`\gamma \gamma \mathrm{e}^+/\mathrm{e}^{}`$ opacity ($`\tau _{\gamma \gamma }1`$) may lead to significant X-ray emission with a luminosity comparable to the luminosity of the primary VHE emission (Aharonian (2000)).
Models involving meson production inevitably predict neutrino emission due to the decay of charged mesons. The SPB-model for PKS 2155$``$304 explains the high energy emission dominantly as proton synchrotron radiation, making the neutrino flux completely negligible.
## 5 Conclusions
This paper reports multi-wavelength observations of the BL Lac object PKS 2155$``$304 in 2003. Although the source appeared variable in the VHE, X-rays and optical bands, the latter two indicate that PKS 2155$``$304 was in a state close to historically low measurements, even though it was easily detectable by H.E.S.S. in all nights of observations since the beginning of the detector operation (see AH04 for the observation history up to August 2003). An extreme case of VHE variability shows a peak-to-peak increase of a factor of $`2.5\pm 0.9`$ in 0.09 d. Variability on the timescale of a few ks in the 2–$`10\mathrm{keV}`$ band and of the order of $`0.1\mathrm{d}`$ for energies $`>300\mathrm{GeV}`$ were observed by RXTE and H.E.S.S. The X-ray data show a correlation of the flux with the spectrum, which becomes harder when the source is brighter. At the level of the simultaneously observed modest variability, no correlation between the VHE $`\gamma `$-rays, the X-rays and optical was seen. Observations with greater variability and better coverage are needed before it can be asserted that the VHE/X-ray pattern in PKS 2155$``$304 is different from other known VHE-emitting AGN. Since the source was in a low emission state in both the optical and X-rays compared to archival measurements, this lack of correlation has yet to be established for a higher emission state. Simultaneous observations in the X-rays/optical band and VHE $`\gamma `$-rays had never previously been performed on this scale. Its continual VHE detection makes PKS 2155$``$304 unique in the TeV BL Lac category, and probably indicates that H.E.S.S. has achieved the sensitivity level where it can detect the quiescent state of PKS 2155$``$304 at any time. A time-averaged energy spectrum is determined for the 2 observation periods and fits to a power law ($`\mathrm{\Gamma }=3.37\pm 0.07`$) in the VHE $`\gamma `$-rays, and to a broken power law ($`\mathrm{\Gamma }_L=2.81\pm 0.05`$, $`E_b=4.9\pm 0.8\mathrm{keV}`$, $`\mathrm{\Gamma }_H=2.95\pm 0.04`$) in the X-rays.
A comparison of the intrinsic spectrum with predictions from existing one-zone leptonic and one-zone hadronic models is attempted to give a plausible estimation of underlying physical parameters. The values of the parameters are in line with those inferred for other VHE-emitting blazars. In these models the VHE emission is attenuated according to two different EBL levels. This changes mainly the Doppler boosting in the leptonic model, but the high level EBL decreases the agreement significantly. In the hadronic model, the maximum injected proton energy can be changed to accomodate different EBL levels and can therefore not significantly constrain any of the EBL models used here.
###### Acknowledgements.
The authors would like to thank the anonymous referee for his correction and useful comments that improved this paper. The support of the Namibian authorities and of the University of Namibia in facilitating the construction and operation of H.E.S.S. is gratefully acknowledged, as is the support by the German Ministry for Education and Research (BMBF), the Max-Planck-Society, the French Ministry for Research, the CNRS-IN2P3 and the Astroparticle Interdisciplinary Programme of the CNRS, the U.K. Particle Physics and Astronomy Research Council (PPARC), the IPNP of the Charles University, the South African Department of Science and Technology and National Research Foundation, and by the University of Namibia. We appreciate the excellent work of the technical support staff in Berlin, Durham, Hamburg, Heidelberg, Palaiseau, Paris, Saclay, and in Namibia in the construction and operation of the equipment. The authors acknowledge the support of the ROTSE III collaboration and the sharing of observation time with the Australian ROTSE IIIa telescope operated by A. Phillips and M.C.B. Ashley from the School of Physics, Department of Astrophysics and Optics, University of New South Wales, Sydney, Australia. Special thanks also to R. Quimby from the University of Texas for providing tools for data-reduction. H. Sol and C. Boisson thank K. Katarzyński for his SSC code. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. We thank the RXTE team for their prompt response to our ToO request and the professional interactions that followed.
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# Underlying Dynamics of Typical Fluctuations of an Emerging Market Price Index: The Heston Model from Minutes to Months
## 1 Introduction
In the last decades the quantitative finance community has devoted much attention to the modeling of asset returns having as a major drive the improvement of pricing techniques by employing stochastic volatility models ameanable to analytical treatment such as Hull-White , Stein-Stein and Heston models.
Despite differences in methods and emphasis, the cross fecundation between Economics and Physics, which dates back to the early nineteenth century (see and ), has intensified recently . Following the statistical physics approach, substantial effort has been made to find microeconomic models capable of reproducing a number of recurrent features of financial time series such as: returns aggregation (probability distributions at any time scale) , volatility clustering , leverage effect (correlation between returns and volatilities) , conditional correlations and fat tails at very short time scales . A central feature of economical phenomena is the prevalence of intertwined dynamics at several time scales. In very general terms, one could divide the market dynamics onto three broad classes: the microeconomic dynamics at the time scales of books of orders and price formation, the mesoeconomic dynamics at the scales of oscillations in formed prices due to local supply and demand and the macroeconomic dynamics at the scales of aggregated economy trends.
The literature on empirical finance have emphasized the multifractal scaling and the power law tails of price fluctuations. However it has been shown that very large data sets are required in order to distinguish between a multifractal and power law tailed process and a stochastic volatility model. In this paper we, therefore, deal with the mesoeconomic dynamics as it would be described by the Heston model with stochastic volatility and exponential tails.
Recently, a semi-analytical solution for the Fokker-Planck equation describing the distribution of log-returns in the Heston model has been proposed . The authors were able to show a satisfactory agreement between return distributions of a number of developed market stock indices and the model for time scales spanning a wide interval ranging from 1 to 100 days. More recently, the same model has also been employed to describe single stocks intraday fluctuations with relative success .
In this paper we show evidence that the Heston model is also capable of describing the fluctuation dynamics of an emerging market price index, the Brazilian São Paulo Stock Exchange Index (IBOVESPA). We employ for the analysis 37 years of data since IBOVESPA inception in January, 1968 and approximately four years of intraday data as well. In this period the Brazilian economy has experienced periods of political and economical unrest, of hyperinflation, of currency crises and of major regulatory changes. These distinctive characteristics make the IBOVESPA an interesting ground for general modelling and data analysis and for testing the limits of description of the Heston model.
This paper is organized as follows. The next section surveys the Heston model, its semi-analytical solution and its connection to Ehrenfest urn models. Section 3 discusses the pre-processing necessary for isolating the fluctuations to be described by the Heston model from exogenous effects. Section 4 describes the data analysis at low (from daily fluctuations) and high (intraday fluctuations) frequencies. Conclusions and further directions are presented in Section 5.
## 2 The Heston Model
### 2.1 Semi-analytical solution
The Heston Model describes the dynamics of stock prices $`S_t`$ as a geometric Brownian motion with volatility given by a Cox-Ingersoll-Ross (or Feller) mean-reverting dynamics. In the Itô differential form the model reads:
$`dS_t`$ $`=`$ $`S_t\mu _tdt+S_t\sqrt{v_t}dW_0(t)`$ (1)
$`dv_t`$ $`=`$ $`\gamma \left[v_t\theta \right]dt+\kappa \sqrt{v_t}dW_1(t),`$
where $`v_t`$ represents the square of the volatility and $`dW_j`$ are Wiener processes with:
$`dW_j(t)`$ $`=`$ $`0,`$
$`dW_j(t)dW_k(t^{})`$ $`=`$ $`\delta (tt^{})\left[\delta _{jk}dt+(1\delta _{jk})\rho dt\right]`$ (2)
The volatility reverts towards a macroeconomic long term mean squared volatility $`\theta `$ with relaxation time given by $`\gamma ^1`$, $`\mu _t`$ represents a drift at macroeconomic scales, the coefficient $`\sqrt{v_t}`$ prevents negative volatilities and $`\kappa `$ regulates the amplitude of volatility fluctuations.
As we are mainly concerned with price fluctuations, we simplify equation (1) by introducing log-returns in a window $`t`$ as $`r(t)=\mathrm{ln}(S(t))\mathrm{ln}(S(0))`$. Using Ito’s lemma and changing variables by making $`x(t)=r(t)_0^t𝑑s\mu _s`$ we obtain a detrended version of the return dynamics that reads:
$$dx=\frac{v_t}{2}dt+\sqrt{v_t}dW_0.$$
(3)
The solution of the Fokker-Planck equation (FP) describing the unconditional distribution of log-returns was obtained by Dr$`\stackrel{˘}{a}`$gulescu and Yakovenko yielding:
$$P_t(x)=_{\mathrm{}}^+\mathrm{}\frac{dp_x}{2\pi }e^{ip_xx+\alpha \varphi _t(p_x)},$$
(4)
where
$`\alpha `$ $`=`$ $`{\displaystyle \frac{2\gamma \theta }{\kappa ^2}}`$ (5)
$`\varphi _t(p_x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }t}{2}}\mathrm{ln}\left[\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Omega }t}{2}}\right)+{\displaystyle \frac{\mathrm{\Omega }^2\mathrm{\Gamma }^2+2\gamma \mathrm{\Gamma }}{2\gamma \mathrm{\Omega }}}\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Omega }t}{2}}\right)\right]`$ (6)
$`\mathrm{\Omega }`$ $`=`$ $`\sqrt{\mathrm{\Gamma }^2+\kappa ^2\left(p_x^2ip_x\right)}`$ (7)
$`\mathrm{\Gamma }`$ $`=`$ $`\gamma +i\rho \kappa p_x.`$ (8)
This unconditional probability density has exponentially decaying tails and the following asymptotic form for short time $`t\gamma ^1`$:
$$P_t(x)=\frac{2^{1\alpha }e^{x/2}}{\mathrm{\Gamma }(\alpha )}\sqrt{\frac{\alpha }{\pi \theta t}}\left(\frac{2\alpha x^2}{\theta t}\right)^{\frac{2\alpha 1}{4}}K_{\alpha 1/2}\left(\sqrt{\frac{2\alpha x^2}{\theta t}}\right),$$
(9)
where $`K_\beta (x)`$ is the modified Bessel function of order $`\beta `$.
### 2.2 Volatility autocorrelation, volatility distribution and leverage function
The volatility can be obtained by integrating (1) and is given by:
$$v_t=\left(v_0\theta \right)e^{\gamma _1t}+\theta +\kappa _0^t𝑑W_1(u)e^{\gamma (tu)}\sqrt{v_u}.$$
(10)
A simple calculation gives the stationary autocorrelation function:
$`C(\tau \gamma ,\theta ,\kappa )\underset{t\mathrm{}}{lim}{\displaystyle \frac{v_tv_{t+\tau }v_tv_{t+\tau }}{\theta ^2}}={\displaystyle \frac{e^{\gamma \tau }}{\alpha }}.`$ (11)
The probability density for the volatility can be obtained as the stationary solution for a Fokker-Planck equation describing $`v_t`$ and reads:
$$\mathrm{\Pi }(v)=\frac{\alpha ^\alpha }{\mathrm{\Gamma }(\alpha )}\frac{v^{\alpha 1}}{\theta ^\alpha }e^{\alpha v/\theta },$$
(12)
implying that $`\alpha >1`$ is required in order to have a vanishing probability density as $`v0`$.
The leverage function describes the correlation between returns and volatilities and is given by :
$$L(\tau \gamma ,\theta ,\kappa ,\rho )\underset{t\mathrm{}}{lim}\frac{dx_t\left(dx_{t+\tau }\right)^2}{\left(dx_{t+\tau }\right)^2^2}=\rho \kappa H(\tau )G(\tau )e^{\gamma \tau },$$
(13)
where $`dx_t`$ is given by (3), $`H(\tau )`$ is the Heaviside step function and:
$$G(\tau )=\frac{v_t\mathrm{exp}\left[\frac{\kappa }{2}_t^{t+\tau }𝑑W_1(u)v_u^{\frac{1}{2}}\right]}{v_t^2}.$$
(14)
To simplify the numerical calculations we employ throughout this paper the zeroth order appoximation $`G(\tau )G(0)=\theta ^1`$. The approximation error increases with the time lag $`\tau `$ but is not critical to our conclusions.
### 2.3 Relation to Ehrenfest Urn Model
The Ehrenfest Urn (EU) as a model for the market return fluctuations has been studied in , in this section we observe that Feller processes as (1) can be produced by a sum of squared Ornstein-Uhlenbeck processes (OU), and that OU processes can be obtained as a large urn limit for the EM.
To see how those connections take shape we start by supposing that $`X_j`$ are OU processes:
$$dX_j(t)=\frac{b}{2}X_j(t)dt+\frac{a}{2}dW_j(t),$$
(15)
where $`dW_j`$ describe $`d`$ independent Wiener processes. In this case the variable $`v(t)_{j=1}^dX_j^2(t)`$ is described by a Feller process as (1) . The volatility process in (1) emerges from OU processes by applying Itô’s Lemma to get:
$$dv_t=bdt\underset{j}{\overset{d}{}}X_j^2+a\underset{j}{\overset{d}{}}X_jdW_j+\frac{a^2}{4}\underset{j}{\overset{d}{}}dW_j^2.$$
(16)
Using the definition of $`v`$ and the properties of the Wiener processes it follows that:
$$dv_t=\left[\frac{d}{4}a^2bv_t\right]dt+a\sqrt{v_t}dW.$$
(17)
The volatility process in (1) can be recovered with a few variable choices: $`a=\kappa `$, $`b=\gamma `$ and $`\theta =\frac{d}{4}\frac{\kappa ^2}{\gamma }`$. A dynamics with vanishing probability of return to the origin requires the volatility to be represented by a sum of at least two elementary OU processes, equivalently we should have $`\alpha =\frac{d}{2}1`$.
An OU process can be obtained as a limit for an Ehrenfest urn model (EM). In an EM $`N`$ numbered balls are distributed between two urns $`A`$ and $`B`$. At each time step a number $`1,\mathrm{},N`$ is chosen at random and the selected ball changes urn. We can specify $`S(n)=(s_1(n),\mathrm{},s_N(n))`$ as the system microstate at instant $`n`$, where $`s_j(n)=\{\pm 1\}`$ indicates that the ball $`j`$ is inside urn $`A`$ if $`s_j=1`$, we also can define the macrostate $`M(n)=_{j=1}^Ns_j(n)`$, which dynamics is described by a Markov chain over the space $`0,\mathrm{},N`$ with transition matrix given by $`Q(k,k+1)=1k/N`$, $`Q(k,k1)=k/N`$ and $`Q(j,k)=0`$ if $`|jk|1`$. An imbalance in the population between the two urns generates a restitution force and, consequently, a mean reverting process for $`M(n)`$. Applying the thermodynamic limit $`N\mathrm{}`$, rescaling time to $`t=n/N`$ and introducing a rescaled imbalance variable as:
$$X_t^{(N)}=\sqrt{N}\left(\frac{M(tN)}{N}\frac{1}{2}\right),$$
(18)
we recover an OU process:
$$dX_t=2X_tdt+dW_t.$$
(19)
Using this connection we could speculate on a possible microscopic model that would generate a stochastic dynamics as described by the Heston model. Supposing that market agents choose at each time step between two regimes (urns) that may represent different strategies, such as technical and fundamental trading or expectations on future bullish or bearish markets, imbalances between populations in each regime would be a source of volatility. In this picture, the condition $`d2`$ would imply that at least two independent sources of volatility would be driving the market. We shall develop this connection further elsewhere.
## 3 Data Pre-processing
### 3.1 The data
Two data sets have been used: IB1 consisting of daily data from IBOVESPA inception on January, 1968 to January, 2005 and IB2 consisting of high-frequency data sampled at 30 second intervals from March 3, 2001 to October 25, 2002 and from June 6, 2003 to August 18, 2004.
### 3.2 Inplits and Inflation
The dataset IB1 has been adjusted for eleven divisions by 10 (inplits) introduced in the period for disclosure purposes and also for inflation by the General Price Index (IGP) . In figure 1A we show the discount factor from date $`t`$ to current date $`T`$, $`I_T(t)`$, used for correcting past prices $`S_t`$ as $`S_t^T=S_tI_T(t)`$. Figure 1B shows the resulting inflation adjusted prices. The hiperinflation (average of $`25\%`$ per month) period from January 1987 to July 1994 is also indicated in both figures by dashed lines.
### 3.3 Detrending
For our analysis of price fluctuations it would be highly desirable to identify macroeconomic trends that may be present in the data. A popular detrending technique is the Hodrick-Prescott (HP) filtering which is based on decomposing the dynamics into a permanent component, or trend $`x_P(t)`$, and into a stochastic stationary component $`x(t)`$ as $`r(t)=x_P(t)+x(t)`$ by fitting a smooth function to the data. Any meaningful detrending procedure has to conserve statistical properties that define the fluctuations. However, in our experiments we have noticed that the HP filtering may introduce subdiffusive behavior at long time scales as an artifact when applied on first differences of a random walk. In this paper, in the absence of a reliable detrending procedure, we assume that the major long term trend is represented by inflation.
### 3.4 Autocorrelation of Returns
In the period span by data set IB1 the Brazilian economy (see for a brief historical account) has experienced a number of regulatory changes with consequences for price fluctuations. In it has been observed that the Hurst exponent for daily IBOVESPA returns shows an abrupt transition from an average compatible with long memory ($`H>0.5`$) to a random walk behavior ($`H=0.5`$) that coincides with major regulatory changes (known as the Collor Plan). We have confirmed the presence of memory by measuring the autocorrelation function for the entire time series (figure 2A). We have also measured the historical autocorrelation for one day lags using a 250 days moving window and confirm an abrupt behavior change coinciding with the Collor Plan in 1990 (figure 2B). As the Heston model assumes uncorrelated returns and this feature only became realistic after the Collor Plan, the following analysis is restricted to a data set consisting of daily data from January 1990 to January 2005 (IB3).
### 3.5 Extreme Events and Abnormal Runs
Abnormal runs are sequences of returns of the same sign that are incompatible with the hypotheses of uncorrelated daily returns. We follow and calculate for the data set IB3 the statistics of persistent price decrease (drawdowns) or increase (drawups). We then compare empirical distributions of runs with shuffled versions of IB3. In figure 3B we show that a seventeen business days drawdown from March 7, 1990 to March 29,1990 is statistically incompatible, within a $`98\%`$ confidence interval, with equally distributed uncorrelated time series. Observe that the largest drawup shown in figure 3C correponds to the subsequent period from March 30, 1990 to April 5, 1990 and that the Collor Plan was launched in March 1990. We therefore expunged from data set IB3 abnormally correlated runs that took place in March, 1990.
## 4 Data Analysis
### 4.1 Low Frequency
After deflating prices and expunging autocorrelated returns, four independent parameters have to be fit to the data: the long term mean volatility $`\theta `$, the relaxation time for mean reversion $`\gamma ^1`$, the amplitude of volatility fluctuations $`\kappa `$ and the correlation between price and volatility $`\rho `$. It has become apparent in that a direct least squares fit to the probability density (4) yields parameters that are not uniquely defined. The data analysis adopted in this work consists in looking for statistics predicted by the model that can be easily measured and compared. In the following subsections we describe these statistics.
#### 4.1.1 Long term mean volatility $`\theta `$
A straightforward calculation yields the second cummulant of the probability density (4) as:
$`c_2(t)`$ $`=`$ $`x(t)^2x(t)^2=\alpha {\displaystyle \frac{^2\varphi _t(k)}{k^2}}|_{k=0}`$
$`=`$ $`\theta t\left[1ϵ\right]+\theta {\displaystyle \frac{ϵ}{\gamma }}\left[1e^{\gamma t}\right],`$
where
$$ϵ=\frac{\kappa }{\gamma }\left[\rho \frac{1}{4}\frac{\kappa }{\gamma }\right].$$
(21)
As $`\kappa \gamma `$ we drop terms of $`O\left(\frac{\kappa }{\gamma }\right)`$ or superior. A non-biased estimate for the above cummulant can be calculated easily from data using:
$$\widehat{c_2}(t)=\frac{1}{N1}\underset{j=1}{\overset{N}{}}\left[x_j^{(t)}\frac{1}{N}\underset{i=1}{\overset{N}{}}x_i^{(t)}\right]^2,$$
(22)
where $`x^{(t)}`$ stands for $`t`$-days detrended log-returns. The long term mean volatility is estimated by a linear regression over the function $`\widehat{c_2}(t)`$ as shown in figure 4A.
#### 4.1.2 Relaxation time for mean reversion $`\gamma ^1`$
The relaxation time is also estimated by a linear regression over the logarithm of the empirical daily volatility autocorrelation function (11) given by:
$$\widehat{C}(\tau )=\frac{\frac{1}{N1}_{j=1}^N\left(x_j^{(1)}\right)^2\left(x_{j+\tau }^{(1)}\right)^2}{\widehat{\theta }^2}1.$$
(23)
We observe in figure 4B and 4C that IB3 data can be fit to two relaxation times: at $`\tau <20`$ we fit $`\gamma _1^1=20\pm 4`$ days and at $`\tau >20`$ we fit $`\gamma _2^1=270\pm 22`$ days.
The second relaxation time can be introduced into the original Heston model by making the long-term volatility fluctuate as a mean reverting process like:
$`d\theta `$ $`=`$ $`\gamma _2\left[\theta (t)\theta _0\right]dt+\kappa _2\sqrt{\theta (t)}dW_2(t),`$ (24)
where $`dW_2`$ is an additional independent Wiener process. The autocorrelation function for the two time scales model acquires the following form :
$`C(\tau )={\displaystyle \frac{e^{\gamma _1\tau }}{\alpha _1}}+{\displaystyle \frac{e^{\gamma _2\tau }}{\alpha _2}},`$ (25)
with $`\alpha _2=\frac{2\gamma _2\overline{\theta }}{\kappa _2^2}`$, where $`\overline{\theta }`$ stands for the average of $`\theta `$ given $`\theta _0`$.
The autocorrelation function fit in figure 4C indicates that relaxation times are of very different magnitudes and that it may be possible to solve the Fokker-Planck equation for such model approximately by adiabatic elimination. We pursue this direction further elsewhere. It is worth mentioning that another tractable alternative for introducing multiple time scales for the volatility is a superposition of OU processes as described in .
#### 4.1.3 Amplitude of volatility fluctuations $`\kappa `$
From (5) the amplitude of volatility fluctuations is given by:
$$\kappa =\sqrt{2\gamma \frac{\theta }{\alpha }}.$$
(26)
The long term volatility $`\theta `$ and the relaxation time $`\gamma ^1`$ have been estimated in the previous sections. As discussed in section 2.3, the constant $`\alpha `$ is related to the number of independent OU processes composing the stochastic volatility process as $`d=2\alpha `$. To avoid negative volatilities, $`\alpha 1`$ is required. We, therefore, assume $`\alpha =1`$ and calculate the amplitude of volatility fluctuations from (26) yielding $`\kappa =0.0032(10)`$.
#### 4.1.4 Correlation between prices and volatilities $`\rho `$
A nonzero correlation between prices and volatilities in (2.1) leads to an asymmetric probability density of returns described by (4). This asymmetry can be estimated either by directly computing the empirical distribution skewness or by calculating the empirical leverage function (13). Both procedures imply in computing highly noisy estimates for third and fourth order moments. In this section we propose estimating a confidence interval for $`\rho `$ by computing the posterior probability $`p(\rho L)`$, where $`L`$ corresponds to a data set containing the empirically measured leverage function for a given number of time lags. Bayes theorem gives the following posterior distribution:
$$p(\rho L)=\frac{1}{Z(L)}p(\rho )𝑑\gamma 𝑑\theta 𝑑\kappa p(L\rho ,\gamma ,\theta ,\kappa )p(\gamma )p(\theta )p(\kappa ),$$
(27)
where $`Z(L)`$ is a data dependent normalization constant. We assume ignorance on the parameter $`\rho `$ and fix its prior distribution to be uniform over the interval $`[1,+1]`$, so that $`p(\rho )=U\left([1,+1]\right)`$. The maximum entropy prior distributions $`p(\gamma )`$, $`p(\theta )`$ and $`p(\kappa )`$ for the remaining parameters are gaussians, since their mean and variance have been previously estimated. The model likelihood for time lags $`0<\tau <T<\gamma _1^1`$ is approximately given by:
$`p(L\rho ,\gamma ,\theta ,\kappa )`$ $`=`$ $`{\displaystyle \frac{d\sigma p(\sigma )}{\left(2\pi \sigma ^2\right)^{\frac{T}{2}}}}`$
$`\times `$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{2\sigma ^2}}{\displaystyle \underset{\tau =1}{\overset{T}{}}}\left({\displaystyle \frac{1}{M\theta ^2}}{\displaystyle \underset{t=1}{\overset{M}{}}}x_t^{(1)}\left(x_{t+\tau }^{(1)}\right)^2{\displaystyle \frac{\rho \kappa }{\theta }}e^{\gamma \tau }\right)^2\right],`$
where the first term inside the exponencial represents the empirical leverage function with $`x_t^{(1)}`$ being the $`M`$ daily returns in the data set IB3. We choose $`p(\sigma )=U\left([\sigma _{\text{min}},\sigma _{\text{max}}]\right)`$ to be uniform representing our level of ignorance on the acceptable dispersion of deviations between data and model. Having specified ignorance priors and the likelihood (4.1.4), we evaluate the posterior (27) by Monte Carlo sampling. In figure 5 we show the resulting posterior probability density and find the $`95\%`$ confidence interval to be $`\text{IC}_{\text{95\%}}(\rho )=[1.00,0.42]`$, what is strong evidence for an asymmetric probability density of returns.
#### 4.1.5 Probability density of returns
Confidence intervals for the probability density of returns can be obtained by Monte Carlo sampling a sufficiently large number of parameter sets with appropriate distributions and numerically integrating (4) for each set. Distribution features for each of the relevant parameters are summarized in the following table:
| *Parameter* | Mean | Standard Deviation |
| --- | --- | --- |
| $`\theta `$ | $`0.00101`$ $`days^1`$ | $`0.00001`$ $`days^1`$ |
| $`\gamma _1^1`$ | $`20`$ $`days`$ | $`4`$ $`days`$ |
| $`\gamma _2^1`$ | $`270`$ $`days`$ | $`22`$ $`days`$ |
| $`\kappa `$ | $`0.0032`$ $`days^1`$ | $`0.0010`$ $`days^1`$ |
| $`\rho `$ | $`\text{IC}_{\text{95\%}}=[1.00,0.42]`$ | distributed as in fig. 5. |
| $`\alpha `$ | $`1.0`$ | fixed by theoretical arguments. |
We have independently sampled gaussian distributions on the parameters $`\theta `$,$`\gamma _1`$,$`\kappa `$ and a uniform distribution $`U([1.00,0.42])`$ for the parameter $`\rho `$. In figure 6 we compare empirical probability densities with theoretical confidence intervals at $`95\%`$ finding a clear agreement at time scales ranging from 1 to 160 days.
### 4.2 High Frequency
#### 4.2.1 Autocorrelation of intraday returns
Our main aim is to describe the fluctuation dynamics at intermediate time scales of formed prices (mesosconomic time scales) by a model which assumes uncorrelated returns. The price formation process occurs at time scales from seconds to a few minutes where the placing of new orders and double auction processes take place. We propose to fix the shortest mesoeconomic time scale to be the point where the intraday autocorrelation function vanishes. In figure 7 we show that the intraday return autocorrelation function vanishes at about 20 minutes for each of the 4 years composing data set IB2. We, therefore, consider as mesoeconomic time scales over 20 minutes.
#### 4.2.2 Effective duration of a day
At first glance, it is not clear whether intraday and daily returns can be described by the same stochastic dynamics. Even less clear is whether aggregation from intraday to daily returns can be described by the same parameters. To verify this latter possibility we have to transform units by determining the effective duration in minutes of a business day $`T_{eff}`$. This effective duration must include the daily trading time at the São Paulo Stock Exchange and the impact of daily and overnight gaps over the diffusion process. The São Paulo Stock exchange opens daily for electronic trading from 10 a.m. to 5 p.m. local time and from 5:45 p.m. to 7 p.m. for after-market trading, totalizing 8h15min of trading daily.
To estimate $`T_{eff}`$ in minutes we observe that the daily return variance $`v^{(1d)}`$ is the result of the aggregation of 20 minute returns, so that, considering a diffusive process, we would have:
$$v^{(1d)}=\frac{T_{eff}}{20}v^{(20min)}.$$
(29)
It has been already observed that the volatility fluctuation dynamics are mean reverting with at least two time scales $`\gamma _1^120`$ days and $`\gamma _2^11`$ year. Considering the longest relaxation time we estimate $`T_{eff}`$ by estimating the mean daily volatility for each one of the years in IB2. In figure 8A we show linear regressions employed for estimating the mean daily variance $`v^{(1d)}`$ for each year in IB2 following the procedure described in section 4.1.1.. In figure 8B we show linear regressions employed to estimate $`v^{(20min)}`$, the effective duration confidence intervals are obtained from (29). The mean $`95\%`$ confidence interval over the 4 years analysed results in $`IC_{95\%}\left(T_{eff}\right)=[`$9h10min, 9h56min$`]`$ what is consistent with 8h15min of daily trading time plus an effective contribution of daily and overnight gaps.
#### 4.2.3 Probability density of intraday returns
For evaluating the probability density of intraday returns we have reestimated the mean volatility $`\theta `$ and the amplitude of volatility fluctuations $`\kappa `$ along the period 2001-2004 represented in the data set IB2. We then have rescaled the dimensional parameters as $`\theta ^{(ID)}=\theta /T_{eff}`$, $`\gamma ^{(ID)}=\gamma /T_{eff}`$ and $`\kappa ^{(ID)}=\kappa /T_{eff}`$. Having rescaled the distributions describing our ignorance on the appropriate returns we have employed Monte Carlo sampling to compute confidence intervals for the short time lags appoximation of the theoretical probability density described in (9). In figure 9 we compare the resulting confidence intervals and the data. We attain reasonably good fits for the longer time scales, as we approach the microeconomic time scales the theoretical description of the tails breaks with the empirical data showing fatter tails.
## 5 Conclusions and Perspectives
We have studied the Heston model with stochastic volatility and exponential tails as a model for the typical price fluctuations of the Brazilian São Paulo Stock Exchange Index (IBOVESPA). Prices have been corrected for inflation and a period spanning the last 15 years, characterized by memoryless returns, have been chosen for the analysis. We also have expunged from data a drawdown inconsistent with the supposition of independence made by the Heston model that took place in the transition between the first 22 years of long memory returns to the memoryless time series we have analysed.
The long term mean volatility $`\theta `$ has been estimated by observing the time scaling of the log-returns variance. The relaxation time for mean reversion $`\gamma ^1`$ has been estimated by observing the autocorrelation function of the log returns variance. We have verified that a modified version of the Heston model with two very different relaxation times ($`\gamma _1^120`$ days and $`\gamma _2^11`$ year) is required for describing the autocorrelation function correctly. We have used the minimum requirement for a non-vanishing volatility $`\alpha 1`$ to calculate the scale of the variance fluctuation $`\kappa `$. Finally, we employed the Bayesian statistics approach for estimating a confidence interval for the volatility-return correlation $`\rho `$, as it relies on a small data set to calculate a noisy estimate of higher order moments. The quality of the model is visually inspected by comparing the empirical probability density at time scales ranging from 1 day to 160 days, with confidence intervals obtained by a Monte Carlo simulation over the maximum ignorance parameter distributions. We have also shown that the probability density functions of log returns at intraday time scales can be described by the Heston model with the same parameters given that we introduce an effective duration for a business day that includes the effect of overnight jumps and that we consider the slow change of the volatility due to the longer relaxation time $`\gamma _2^1`$.
It is surprising and non-trivial that a single stochastic model may be capable of describing the statistical behavior of both developed and emerging markets at a wide range of time scales, despite the known instability and high susceptibility to externalities of the latter. We believe that this robust statistical behavior may point towards simple basic mechanisms acting in the market microstructure. We regard as an interesting research direction to pursue the derivation of the Heston model as a limit case for a microeconomic model like the minority game . For this matter we believe that the connection between the Heston model and Ehrenfest urns may be valuable.
Perhaps the search for underlying symmetries and simple basic mechanisms that can explain empirical observations should be regarded as the main contribution of Physics to Economics. This contribution might be particularly useful to the field of Econometrics in which a common view is that a theory built from data ‘should be evaluated in terms of the quality of the decisions that are made based on the theory’ . Clearly, these two approaches should not be considered as mutually exclusive.
We thank Victor Yakovenko and his collaborators for discussions and for providing useful MATLAB codes. We also wish to thank the São Paulo Stock Exchange (BOVESPA) for gently providing high-frequency data. This work has been partially (RV,VBPL) supported by FAPESP.
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# 𝑍₂ Topological Order and the Quantum Spin Hall Effect
## Abstract
The quantum spin Hall (QSH) phase is a time reversal invariant electronic state with a bulk electronic band gap that supports the transport of charge and spin in gapless edge states. We show that this phase is associated with a novel $`Z_2`$ topological invariant, which distinguishes it from an ordinary insulator. The $`Z_2`$ classification, which is defined for time reversal invariant Hamiltonians, is analogous to the Chern number classification of the quantum Hall effect. We establish the $`Z_2`$ order of the QSH phase in the two band model of graphene and propose a generalization of the formalism applicable to multi band and interacting systems.
The classification of electronic states according to topological invariants is a powerful tool for understanding many body phases which have bulk energy gaps. This approach was pioneered by Thouless et al.tknn (TKNN), who identified the topological invariant for the non interacting integer quantum Hall effect. The TKNN integer, $`n`$, which gives the quantized Hall conductivity for each band $`\sigma _{xy}=ne^2/h`$, is given by an integral of the Bloch wavefunctions over the magnetic Brillouin zone, and corresponds to the first Chern class of a $`U(1)`$ principal fiber bundle on a torusavron ; kohmoto . An equivalent formulation, generalizable to interacting systems, is to consider the sensitivity of the many body ground state to phase twisted periodic boundary conditionsniu1 ; arovas . This topological classification distinguishes a simple insulator from a quantum Hall state, and explains the insensitivity of the Hall conductivity to weak disorder and interactions. Nonzero TKNN integers are also intimately related to the presence of gapless edge states on the sample boundarieshatsugai .
Since the Hall conductivity violates time reversal (TR) symmetry, the TKNN integer must vanish in a TR invariant system. Nonetheless, we have recently shown that the spin orbit interaction in a single plane of graphene leads to a TR invariant quantum spin Hall (QSH) state which has a bulk energy gap, and a pair of gapless spin filtered edge states on the boundarykm . In the simplest version of our model (a $`\pi `$ electron tight binding model with mirror symmetry about the plane) the perpendicular component of the spin, $`S_z`$, is conserved. Our model then reduces to independent copies for each spin of a model introduced by Haldanehaldane , which exhibits an integer quantum Hall effect even though the average magnetic field is zero. When $`S_z`$ is conserved the distinction between graphene and a simple insulator is thus easily understood. Each spin has an independent TKNN integer $`n_{}`$, $`n_{}`$. TR symmetry requires $`n_{}+n_{}=0`$, but the difference $`n_{}n_{}`$ is nonzero and defines a quantized spin Hall conductivity.
This characterization breaks down when $`S_z`$ non conserving terms are present. Such terms will inevitably arise due to coupling to other bands, mirror symmetry breaking terms, interactions or disorder. Though these perturbations destroy the quantization of the spin Hall conductance, we argued that they do not destroy the topological order of the QSH state because Kramers’ theorem prevents TR invariant perturbations from opening a gap at the edgekm . Thus, even though the single defined TKNN number (the total Hall conductance) is zero, the QSH groundstate is distinguishable from that of a simple insulator. This suggests that there must be an additional topological classification for TR invariant systems.
In this paper we clarify the topological order of the QSH phase and introduce a $`Z_2`$ topological index that characterizes TR invariant systems. This classification is similar to the TKNN classification, and gives a simple test which can be applied to Bloch energy bands to distinguish the insulator from the QSH phase. It may also be formulated as a sensitivity to phase twisted boundary conditions. We will begin by describing our model of graphene and demonstrate that the QSH phase is robust even when $`S_z`$ is not conserved. We will then analyze the constraints of TR invariance and derive the $`Z_2`$ index.
Consider the tight binding Hamiltonian of graphene introduced in Ref. km, , which generalizes Haldane’s modelhaldane to include spin with TR invariant spin orbit interactions.
$`H=t{\displaystyle \underset{ij}{}}c_i^{}c_j+i\lambda _{SO}{\displaystyle \underset{ij}{}}\nu _{ij}c_i^{}s^zc_j`$ (1)
$`+i\lambda _R{\displaystyle \underset{ij}{}}c_i^{}(𝐬\times \widehat{𝐝}_{ij})_zc_j+\lambda _v{\displaystyle \underset{i}{}}\xi _ic_i^{}c_i.`$
The first term is a nearest neighbor hopping term on the honeycomb lattice, where we have suppressed the spin index on the electron operators. The second term is the mirror symmetric spin orbit interaction which involves spin dependent second neighbor hopping. Here $`\nu _{ij}=(2/\sqrt{3})(\widehat{𝐝}_1\times \widehat{𝐝}_2)_z=\pm 1`$, where $`\widehat{𝐝}_1`$ and $`\widehat{𝐝}_2`$ are unit vectors along the two bonds the electron traverses going from site $`j`$ to $`i`$. $`s^z`$ is a Pauli matrix describing the electron’s spin. The third term is a nearest neighbor Rashba term, which explicitly violates the $`zz`$ mirror symmetry, and will arise due to a perpendicular electric field or interaction with a substrate. The fourth term is a staggered sublattice potential ($`\xi _i=\pm 1`$), which we include to describe the transition between the QSH phase and the simple insulator. This term violates the symmetry under twofold rotations in the plane.
$`H`$ is diagonalized by writing $`\varphi _s(𝐑+\alpha 𝐝)=u_{\alpha s}(𝐤)e^{i𝐤𝐑}`$. Here $`s`$ is spin and $`𝐑`$ is a bravais lattice vector built from primitive vectors $`𝐚_{1,2}=(a/2)(\sqrt{3}\widehat{y}\pm \widehat{x})`$. $`\alpha =0,1`$ is the sublattice index with $`𝐝=a\widehat{y}/\sqrt{3}`$. For each $`𝐤`$ the Bloch wavefunction is a four component eigenvector $`|u(𝐤)`$ of the Bloch Hamiltonian matrix $`(𝐤)`$. The 16 components of $`(𝐤)`$ may be written in terms of the identity matrix, 5 Dirac matrices $`\mathrm{\Gamma }^a`$ and their 10 commutators $`\mathrm{\Gamma }^{ab}=[\mathrm{\Gamma }^a,\mathrm{\Gamma }^b]/(2i)`$murakami1 . We choose the following representation of the Dirac matrices: $`\mathrm{\Gamma }^{(1,2,3,4,5)}=(\sigma ^xI,\sigma ^zI,\sigma ^ys^x,\sigma ^ys^y,\sigma ^ys^z`$), where the Pauli matrices $`\sigma ^k`$ and $`s^k`$ represent the sublattice and spin indices. This choice organizes the matrices according to TR. The TR operator is given by $`\mathrm{\Theta }|ui(Is^y)|u^{}`$. The five Dirac matrices are even under TR, $`\mathrm{\Theta }\mathrm{\Gamma }^a\mathrm{\Theta }^1=\mathrm{\Gamma }^a`$ while the 10 commutators are odd, $`\mathrm{\Theta }\mathrm{\Gamma }^{ab}\mathrm{\Theta }^1=\mathrm{\Gamma }^{ab}`$. The Hamiltonian is thus
$$(𝐤)=\underset{a=1}{\overset{5}{}}d_a(𝐤)\mathrm{\Gamma }^a+\underset{a<b=1}{\overset{5}{}}d_{ab}(𝐤)\mathrm{\Gamma }^{ab}$$
(2)
where the $`d(𝐤)`$’s are given in Table I. Note that $`(𝐤+𝐆)=(𝐤)`$ for reciprocal lattice vectors $`𝐆`$, so $`(𝐤)`$ is defined on a torus. The TR invariance of $``$ is reflected in the symmetry (antisymmetry) of $`d_a`$ $`(d_{ab})`$ under $`𝐤𝐤`$.
For $`\lambda _R=0`$ the there is an energy gap with magnitude $`|6\sqrt{3}\lambda _{SO}2\lambda _v|`$. For $`\lambda _v>3\sqrt{3}\lambda _{SO}`$ the gap is dominated by $`\lambda _v`$, and the system is an insulator. $`3\sqrt{3}\lambda _{SO}>\lambda _v`$ describes the QSH phase. Though the Rashba term violates $`S_z`$ conservation, for $`\lambda _R<2\sqrt{3}\lambda _{SO}`$ there is a finite region of the phase diagram in Fig. 1 that is adiabatically connected to the QSH phase at $`\lambda _R=0`$. Fig. 1 shows the energy bands obtained by solving the lattice model in a zigzag strip geometrykm for representative points in the insulating and QSH phases. Both phases have a bulk energy gap and edge states, but in the QSH phase the edge states traverse the energy gap in pairs. At the transition between the two phases, the energy gap closes, allowing the edge states to “switch partners”.
The behavior of the edge states signals a clear difference between the two phases. In the QSH phase for each energy in the bulk gap there is a single time reversed pair of eigenstates on each edge. Since TR symmetry prevents the mixing of Kramers’ doublets these edge states are robust against small perturbations. The gapless states thus persist even if the spatial symmetry is further reduced (for instance by removing the $`C_3`$ rotational symmetry in (1)). Moreover, weak disorder will not lead to localization of the edge states because single particle elastic backscattering is forbiddenkm .
In the insulating state the edge states do not traverse the gap. It is possible that for certain edge potentials the edge states in Fig. 1b could dip below the band edge, reducing - or even eliminating - the edge gap. However, this is still distinct from the QSH phase because there will necessarily be an even number of Kramers pairs at each energy. This allows elastic backscattering, so that these edge states will in general be localized by weak disorder. The QSH phase is thus distinguished from the simple insulator by the number of edge state pairs modulo 2. Recently two dimensional versionszhang of the spin Hall insulator modelsmurakami2 have been introduced, which under conditions of high spatial symmetry exhibit gapless edge states. These models, however, have an even number of edge state pairs. We shall see below that they are topologically equivalent to simple insulators.
The QSH phase is not generally characterized by a quantized spin Hall conductivity. Consider the rate of spin accumulation at the opposite edges of a cylinder of circumference $`L`$, which can be computed using Laughlin’s argumentlaughlin . A weak circumferential electric field $`E`$ can be induced by adiabatically threading magnetic flux through the cylinder. When the flux increases by $`h/e`$ each momentum eigenstate shifts by one unit: $`kk+2\pi /L`$. In the insulating state (Fig. 1b) this has no effect, since the valence band is completely full. However, in the QSH state a particle-hole excitation is produced at the Fermi energy $`E_F`$. Since the particle and hole states do not have the same spin, spin accumulates at the edge. The rate of spin accumulation defines a spin Hall conductance $`dS_z/dt=G_{xy}^sE`$, where
$$G_{xy}^s=\frac{e}{h}\left(S_z_LS_z_R\right)|_{E_F}.$$
(3)
Here the expectation value of $`S_z`$ is evaluated for the left and right moving states at $`E_F`$. Since the edge states are not necessarily $`S_z`$ eigenstates this spin Hall conductance is not quantized. $`G_{xy}^s`$ is zero in the insulating phase, though, provided $`E_F`$ is in the gap at the edge. If in the insulator the edge states cross $`E_F`$, then in a clean system there could be spin accumulation at the edge (resulting from the acceleration of the edge electrons in response to $`E`$). However, if the edge states are localized then there will be no spin accumulation. Thus the nonzero spin accumulation persists only for the QSH phase, justifying the term quantum (but not quantized) spin Hall effect.
In the quantum Hall effect, the states with zero and one flux quantum threading the cylinder are distinguished by the charge polarization. The two states can not be connected by any operator that locally conserves charge. In the QSH effect there is no simple conserved quantity distinguishing the two states. However, the states are distinguishable, because the state with an edge particle-hole excitation at $`E_F`$ can not be connected to the ground state with a local TR symmetric operator. Note, however, that if a second flux is added, then there will be TR invariant interactions which do connect the state with the zero flux state. This suggests that the state with one flux added is distinguished by a $`Z_2`$ “TR polarization”.
The classification of quantum Hall states on the cylinder according to Laughlin’s argument is intimately related to the TKNN classification of the Bloch wavefunctionsniu1 . To establish the corresponding topological classification for TR invariant systems we consider TR constraints on the set of occupied Bloch wavefunctions $`|u_{i=1,2}(𝐤)`$. $`|u_i(𝐤)`$ form a rank 2 vector bundle over Brillouin zone torus. TR introduces an involution on the torus which identifies pairs of points $`𝐤`$ and $`𝐤`$. Wavefunctions at the identified points are related by $`|u_i(𝐤)=\mathrm{\Theta }|u_i(𝐤)`$, implying that the bundle is “real”. Since $`\mathrm{\Theta }^2=1`$, $`\mathrm{\Theta }`$ has period 4, so that the real bundle is “twisted”. These bundles are classified within the mathematical framework of twisted Real K theoryatiyah1 . It is found that such bundles have a $`Z\times Z_2`$ classification on a torusthanks . The first integer gives the rank of the bundle (i.e. the number of occupied bands). The $`Z_2`$ index is related to the mod 2 index of the real Dirac operatoratiyah2 . In the following we will explicitly construct this $`Z_2`$ index from the Bloch wavefunctions and show that it distinguishes the QSH phase from the simple insulator.
TR symmetry identifies two important subspaces of the space of Bloch Hamiltonians $`(𝐤)`$ and the corresponding occupied band wavefunctions $`|u_i(𝐤)`$. The “even” subspace, for which $`\mathrm{\Theta }H\mathrm{\Theta }^1=H`$ have wavefunctions with the property that $`\mathrm{\Theta }|u_i`$ is equivalent to $`|u_i`$ up to a $`U(2)`$ rotation. From Eq. 2 it is clear that in this subspace $`d_{ab}(𝐤)=0`$. TR symmetry requires that $`H(𝐤)`$ belong to the even subspace at the $`\mathrm{\Gamma }`$ point $`𝐤=0`$ as well as the three $`M`$ points shown in Fig. 2a,b. The odd subspace has wavefunctions with the property that the space spanned by $`\mathrm{\Theta }|u_i(𝐤)`$ is orthogonal to the space spanned by $`|u_i(𝐤)`$. We shall see that the $`Z_2`$ classification may be found by studying the set of $`𝐤`$ which belong to the odd subspace.
The special subspaces can be identified by considering the matrix of overlaps, $`u_i(𝐤)|\mathrm{\Theta }|u_j(𝐤)`$. From the properties of $`\mathrm{\Theta }`$ it is clear that this matrix is antisymmetric, and may be expressed in terms of a single complex number as $`ϵ_{ij}P(𝐤)`$. $`P(𝐤)`$ is in fact equal to the Pfaffian
$$P(𝐤)=\mathrm{Pf}\left[u_i(𝐤)|\mathrm{\Theta }|u_j(𝐤)\right],$$
(4)
which for a 2 by 2 antisymmetric matrix $`A_{ij}`$ simply picks out $`A_{12}`$. We shall see below that the Pfaffian is the natural generalization when there are more than two occupied bands. $`P(𝐤)`$ is not gauge invariant. Under a $`U(2)`$ transformation $`|u_i^{}=U_{ij}|u_j`$, $`P^{}=PdetU`$. Thus $`P`$ is unchanged by a $`SU(2)`$ rotation, but under a $`U(1)`$ transformation $`U=e^{i\theta }`$, $`P^{}=Pe^{2i\theta }`$. In the even subspace $`\mathrm{\Theta }|u_i`$ is equivalent to $`|u_i`$ up to a $`U(2)`$ rotation, and we have $`|P(𝐤)|=1`$. In the odd subspace $`P(𝐤)=0`$.
If no spatial symmetries constrain its form, the zeros of $`P(𝐤)`$ are found by tuning two parameters, and generically occur at points in the Brillouin zone. First order zeros occur at time reversed pairs of points $`\pm 𝐤^{}`$ with opposite “vorticity”, where the phase of $`P(𝐤)`$ advances in opposite directions around $`\pm 𝐤^{}`$. For $`\lambda _v0`$ the QSH phase is distinguished from the simple insulator by the presence of a single pair of first order zeros of $`P(𝐤)`$. The $`C_3`$ rotational symmetry of our model constrains $`𝐤^{}`$ to be at the corner of the Brillouin zone as shown in Fig 2a. If the $`C_3`$ symmetry is relaxed, $`𝐤^{}`$ can occur anywhere except the four symmetric points where $`|P(𝐤)|=1`$. The number of pairs of zeros is a $`Z_2`$ topological invariant. This can be seen by noting that two pairs $`\pm 𝐤_{1,2}^{}`$ can come together to annihilate each other when $`𝐤_1^{}=𝐤_2^{}`$. However a single pair of zeros at $`\pm 𝐤^{}`$ can not annihilate because they would have to meet at either $`\mathrm{\Gamma }`$ or $`M`$, where $`|P(𝐤)|=1`$. If TR symmetry is broken then the zeros are no longer prevented from annihilating, and the topological distinction of the QSH phase is lost.
The $`Z_2`$ index can thus be determined by counting the number of pairs of complex zeros of $`P`$. This can be accomplished by evaluating the winding of the phase of $`P(𝐤)`$ around a loop enclosing half the Brillouin zone (defined so that $`𝐤`$ and $`𝐤`$ are never both included).
$$I=\frac{1}{2\pi i}_C𝑑𝐤_𝐤\mathrm{log}(P(𝐤)+i\delta ),$$
(5)
where $`C`$ is the path shown in Fig. 2a,b.
When $`\lambda _v=0`$ (as it is in graphene) $`H`$ has a $`C_2`$ rotational symmetry, which when combined with TR constrains the form of $`(𝐤)`$, and allows $`P(𝐤)`$ to be chosen to be real. The zeros of $`P(𝐤)`$ then occur along lines, rather than at points. We find that the zeros are absent in the insulating phase, but enclose the $`M`$ point in the QSH phase as shown in Fig. 2b. In this case we find that Eq. 4 also determines the $`Z_2`$ index (given by 1/2 the number of sign changes along the path $`C`$), provided we include the convergence factor $`\delta `$. Note that though the sign of $`I`$ depends on the sign of $`\delta `$, $`I`$ mod 2 does not. We thus conclude that the QSH phase and the insulator are distinguished by the $`Z_2`$ index $`I`$.
The spin Hall insulator models studied in Refs. zhang, ; murakami2, are simple insulators with $`I=0`$. Their Hamiltonian, when expressed in the form of Eq. 2, has $`d_{ab}(𝐤)=0`$, so that $`|u_i(𝐤)`$ is in the even subspace and $`|P(𝐤)|=1`$ for all $`𝐤`$. Ref. zhang2, introduces a model which does appear to exhibit a QSH effect.
Having established the topological classification of the Bloch wavefunctions we now ask whether, in analogy with Ref. niu1, the classification can be formulated in terms of the sensitivity of the ground state wavefunction to phase twisted periodic boundary conditions. Such a formulation will address the topological stability of the many body groundstate with respect to weak disorder and electron interactions. It also provides the appropriate generalization of (4,5) for multi band Hamiltonians. Consider a $`𝐋_1\times 𝐋_2`$ sample with boundary condition $`\mathrm{\Psi }(\mathrm{},𝐫_i+𝐋_k,\mathrm{})=e^{i\alpha _k}\mathrm{\Psi }(\mathrm{},𝐫_i,\mathrm{})`$. For concreteness we consider a rectangular geometry, with $`𝐋_1=N_1(𝐚_1+𝐚_2)`$ and $`𝐋_2=N_2(𝐚_1𝐚_2)`$. For non interacting electrons, we may view the entire sample as a large unit cell with $`N_a=4N_1N_2`$ atoms imbedded in an even larger crystal. Then $`\stackrel{}{\alpha }`$ plays the role of $`𝐤`$, and the occupied single particle eigenstates $`\varphi _i(\stackrel{}{\alpha })`$ play the role of $`u_i(𝐤)`$. $`\varphi _i(\stackrel{}{\alpha })`$ form a rank $`N_a`$ bundle on the torus defined by $`\alpha _{1,2}`$. The $`Z_2`$ classification can be obtained by studying the zeros of
$$P(\stackrel{}{\alpha })=\mathrm{Pf}\left[\varphi _i(\stackrel{}{\alpha })|\mathrm{\Theta }|\varphi _j(\stackrel{}{\alpha })\right].$$
(6)
Fig. 2c compares $`|P(\stackrel{}{\alpha })|`$ in the QSH and insulating phases. for a 16 site sample with $`N_1=N_2=2`$. In the insulating phase there are no zeros. In the QSH phase the structure of the zeros in Fig. 2d is similar to Fig. 2a,b. For $`\lambda _v0`$ the first order zeros are at points, while for $`\lambda _v=0`$ they are on a loop. The zeros can not be at the four TR symmetric points. This structure persists in the QSH phase for any cell size. The $`Z_2`$ index $`I`$ can be computed by performing the integral analogous to (5) along the contour $`C`$ in Fig. 2d.
A many body formulation requires the index to be expressed in terms of the many particle groundstate $`|\mathrm{\Phi }(\stackrel{}{\alpha })`$. It is interesting to note that for non interacting electrons $`\mathrm{\Phi }(\stackrel{}{\alpha })|\mathrm{\Theta }|\mathrm{\Phi }(\stackrel{}{\alpha })=\mathrm{det}[\varphi _i(\stackrel{}{\alpha })|\mathrm{\Theta }|\varphi _j(\stackrel{}{\alpha })]=P(\stackrel{}{\alpha })^2`$. This suggests a many body generalization
$$P(\stackrel{}{\alpha })=\sqrt{\mathrm{\Phi }(\stackrel{}{\alpha })|\mathrm{\Theta }|\mathrm{\Phi }(\stackrel{}{\alpha })}.$$
(7)
We suspect that with this definition the topological structure $`P(\stackrel{}{\alpha })`$ in Fig. 2c,d will remain in the presence of weak electron interactions.
To conclude, we have introduced a $`Z_2`$ topological classifiation of TR invariant systems, analogous to the TKNN classification of quantum Hall states. This shows that the QSH phase of graphene has a topological stability that is insensitive to weak disorder and interactions.
We thank Tony Pantev for many helpful discussions. This work was supported by the NSF under MRSEC grant DMR-00-79909 and the DOE under grant DE-FG02-ER-0145118.
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# Measurement of high-𝑸^𝟐 deep inelastic scattering cross sections with longitudinally polarised positron beams at HERA
## 1 Introduction
Deep inelastic scattering (DIS) of leptons off nucleons has proved to be a key tool in the understanding of the structure of the proton and the form of the Standard Model (SM). The HERA $`ep`$ collider has made possible the exploration of DIS at high values of negative four-momentum-transfer squared, $`Q^2`$. Using data taken in the years 1994-2000 the H1 and ZEUS collaborations have reported measurements of the cross sections for charged current (CC) and neutral current (NC) DIS $`^{\mathrm{?},\mathrm{?}}`$. These measurements extend the kinematic region covered by fixed-target experiments $`^\mathrm{?}`$ to higher $`Q^2`$ and allow the HERA experiments to probe the electroweak sector of the SM.
This paper presents the measurements of the cross sections for $`e^+p`$ CC and NC DIS with longitudinally polarised positron beams. The measurements are based on the integrated luminosities collected at the luminosity weighted mean polarisations given in Table 1 with the ZEUS and H1 detectors in 2003 and 2004. During this time HERA collided protons of energy $`920\mathrm{GeV}`$ with positrons of energy $`27.6\mathrm{GeV}`$, yielding collisions at a centre-of-mass energy of $`319\mathrm{GeV}`$. The measured cross sections are compared to the SM predictions.
## 2 Polarised lepton beams
At HERA transverse polarisation of the lepton beam arises through synchrotron radiation via the Sokolov-Ternov effect $`^\mathrm{?}`$. As a part of the HERA upgrade in the year 2000 spin rotators that rotate the polarisation into the longitudinal direction were installed in the lepton beamline round the collision regions where the H1 and ZEUS detectors are located. The polarisation is continuously measured using two independent polarimeters. The TPOL $`^\mathrm{?}`$ is situated at a position of transverse lepton beam polarisation, and the LPOL $`^\mathrm{?}`$ at one of longitudinal polarisation. The luminosity weighted polarisation distribution for the data of the presented H1 measurement is shown in Fig. 1.
## 3 Kinematic variables and cross sections
Inclusive deep inelastic lepton-proton scattering can be described in terms of the kinematic variables $`x`$, $`y`$ and $`Q^2`$. The variable $`Q^2`$ is defined to be $`Q^2=q^2=(kk^{})^2`$ where $`k`$ and $`k^{}`$ are the four-momenta of the incoming and scattered lepton, respectively. Bjorken $`x`$ is defined by $`x=Q^2/2Pq`$ where $`P`$ is the four-momentum of the incoming proton. The variable $`y`$ is defined by $`Q^2=sxy`$, where $`s=4E_eE_p`$ is the square of the lepton-proton centre-of-mass energy (neglecting the masses of the incoming particles).
The electroweak Born level cross section for the CC reaction, $`e^+p\overline{\nu _e}X`$, with a longitudinally polarised positron beam (defined in Eqn. (1)), can be expressed as $`^\mathrm{?}`$
$$\frac{d^2\sigma ^{CC}(e^+p)}{dxdQ^2}=(1+𝒫)\frac{G_F^2}{4\pi x}\left(\frac{M_W^2}{M_W^2+Q^2}\right)^2\left[Y_+F_2^{CC}(x,Q^2)Y_{}xF_3^{CC}(x,Q^2)y^2F_L^{CC}(x,Q^2)\right],$$
where $`G_F`$ is the Fermi constant, $`M_W`$ is the mass of the $`W`$ boson and $`Y_\pm =1\pm (1y)^2`$. The structure functions $`F_2^{CC}`$ and $`xF_3^{CC}`$ contain sums and differences of the quark and anti-quark parton density functions (PDFs) and $`F_L^{CC}`$ is the longitudinal structure function. The longitudinal polarisation of the positron beam is defined as
$$𝒫=\frac{N_RN_L}{N_R+N_L},$$
(1)
where $`N_R`$ and $`N_L`$ are the numbers of right and left-handed positrons in the beam. Similarly the cross section for the NC reaction, $`e^+pe^+X`$, can be expressed as
$$\frac{d^2\sigma ^{NC}(e^+p)}{dxdQ^2}=\frac{2\pi \alpha ^2}{xQ^4}[H_0^++𝒫H_𝒫^+],$$
where $`\alpha `$ is the QED coupling constant and $`H_0^+`$ and $`H_𝒫^+`$ contain the unpolarised and polarised structure functions, respectively.
Charged current events are characterised by a large missing transverse momentum, $`P_{T,\mathrm{miss}}`$, which is calculated by ZEUS as
$$P_{T,\mathrm{miss}}^2=P_x^2+P_y^2=\left(\underset{i}{}E_i\mathrm{sin}\theta _i\mathrm{cos}\varphi _i\right)^2+\left(\underset{i}{}E_i\mathrm{sin}\theta _i\mathrm{sin}\varphi _i\right)^2,$$
where the sum runs over all calorimeter energy deposits $`E_i`$, and $`\theta _i`$ and $`\varphi _i`$ are the polar and azimuthal angles of the calorimeter cell as viewed from the interaction vertex. The hadronic polar angle, $`\gamma _h`$, is defined by $`\mathrm{cos}\gamma _h=(P_{T,\mathrm{miss}}^2\delta ^2)/(P_{T,\mathrm{miss}}^2+\delta ^2)`$, where $`\delta =(E_iE_i\mathrm{cos}\theta _i)=(EP_z)_i`$. H1 uses an analog definition summing over all final state particles. In the naive Quark Parton Model, $`\gamma _h`$ gives the scattering angle of the struck quark in the laboratory frame. The total transverse energy, $`E_T`$, is given by $`E_T=E_i\mathrm{sin}\theta _i`$.
Neutral current events are characterised by the presence of a high-energy isolated scattered positron in the detector. It follows from longitudinal momentum conservation that for well measured NC events $`\delta `$ peaks at twice the positron beam energy i.e. $`55\mathrm{GeV}`$.
The kinematic variables for charged current events were reconstructed from the measured $`P_{T,\mathrm{miss}}`$ and $`\delta `$ using the Jacquet-Blondel method $`^\mathrm{?}`$. For the neutral current events the ZEUS measurement uses the double-angle method $`^\mathrm{?}`$ to estimate the kinematic variables from the polar angles of the scattered positron $`\theta _e`$ and the hadronic final state $`\gamma _h`$ while H1 uses the $`e\mathrm{\Sigma }`$ method $`^\mathrm{?}`$ to estimate the kinematic quantities from the scattered positrons energy, its angle and $`\delta `$.
## 4 Monte Carlo simulation
Monte Carlo simulation (MC) was used to determine the efficiency for selecting events, the accuracy of kinematic reconstruction, to estimate the background rate and to deduce cross sections for the full kinematic region from the data. A sufficient number of events were generated to ensure that uncertainties from MC statistics were small. The MC samples were normalised to the total integrated luminosity of the data.
Neutral and charged current DIS events including radiative effects were simulated using the djangoh $`^\mathrm{?}`$ generator. The hadronic final state was simulated using the colour-dipole model of ariadne $`^\mathrm{?}`$. For the hadronisation the Lund string model of jetset 7.4 $`^\mathrm{?}`$ is used. Additional samples used are as described in $`^{\mathrm{?},\mathrm{?}}`$.
## 5 Event selection
CC events are selected by requiring $`P_{T,\mathrm{miss}}>12\mathrm{GeV}`$. In order to ensure high trigger efficiency and good kinematic resolution the analysis is restricted to the kinematic region of $`Q^2>200\mathrm{GeV}^2`$ and in $`y`$ by $`0.03<y<0.85`$ for the H1 measurement and $`y<0.9`$ for the ZEUS measurement. Non-$`ep`$ background is rejected by searching for typical beam-induced background event topologies. For the ZEUS measurement a total of 604 candidate events passed the selection criteria. The background contamination was estimated to be typically less than 1% but was as high as 5% in the lowest $`Q^2`$ bin of the negative polarisation sample. The contribution of simulated background for the H1 measurement is visible at low $`Q^2`$ in Fig. 2. It shows a comparison of data and MC distributions for the CC sample of the H1 and ZEUS measurements. The Monte Carlo gives a good description of the data.
NC events are selected by identifying the scattered electron. For this task sophisticated algorithms are used by both experiments. The main background supression is achieved by requiring a reconstructed scattered electron energy $`E_e^{}>11\mathrm{GeV}`$ for H1 and $`E_e^{}>10\mathrm{GeV}`$ for ZEUS with additional isolation criteria. The ZEUS analysis selected a total of 52004 candidate events in the kinematic region of $`Q^2>200\mathrm{GeV}^2`$ and $`y<0.95`$. The background contamination was estimated to be typically less than 1%. Figure 3 shows a comparison of data and MC expectation distributions for the NC sample. The Monte Carlo gives a generally good description of the data. Inaccuracies in the simulation of the scattered positron energy distribution are considered in the corresponding systematic uncertainty. The effects of the positron fiducial-volume cut can be seen in the distribution of the scattered positron angle.
## 6 Results
The total cross section for $`e^+p`$ CC DIS in the kinematic region $`Q^2>200\mathrm{GeV}^2`$ was measured by ZEUS at the given longitudinal positron beam polarisations to be:
$`\sigma _{\mathrm{CC}}(P=31.8\pm 0.9\%)`$ $`=`$ $`46.7\pm 2.4(\mathrm{stat}.)\pm 1.0(\mathrm{syst}.)\mathrm{pb},`$
$`\sigma _{\mathrm{CC}}(P=40.2\pm 1.1\%)`$ $`=`$ $`22.5\pm 1.6(\mathrm{stat}.)\pm 0.5(\mathrm{syst}.)\mathrm{pb}.`$
Note that the ZEUS measurement does not include the uncertainty in the luminosity of $`\pm 5\%`$ in the systematic uncertainty. The total cross section is shown as a function of the longitudinal polarisation of the positron beam in Fig. 4 including the unpolarised ZEUS measurement from the 1999-2000 data $`^\mathrm{?}`$. The data are compared to the SM prediction evaluated using the ZEUS-S PDFs $`^\mathrm{?}`$. The SM prediction describes the data well. The cross section points at polarisations of 31.8% and -40.2% respectively are 3.4 standard deviations above and 6.1 below the unpolarised measurement, respectively.
The single-differential cross-sections, $`d\sigma /dQ^2`$, $`d\sigma /dx`$ and $`d\sigma /dy`$ for charged current DIS are also shown in Fig. 4. A clear difference is observed between the measurements for positive and negative longitudinal polarisation, which is well described over the whole kinematic region by the SM evaluated using the ZEUS-S PDFs.
H1 measured the integrated polarised cross section in the kinematic range $`Q^2>400\mathrm{GeV}^2`$ and $`y<0.9`$ to be:
$`\sigma _{\mathrm{CC}}(P=+33\pm 2\%)`$ $`=`$ $`34.7\pm 1.9(\mathrm{stat}.)\pm 1.7(\mathrm{syst}.)\mathrm{pb},`$
$`\sigma _{\mathrm{CC}}(P=40.2\pm 1.5\%)`$ $`=`$ $`13.8\pm 1.0(\mathrm{stat}.)\pm 1.0(\mathrm{syst}.)\mathrm{pb}.`$
The total cross section is shown as a function of the longitudinal polarisation together with the unpolarised CC cross section measurement based on the HERA I data set $`^\mathrm{?}`$ with an integrated luminosity of $`65.2\mathrm{pb}^1`$ in Fig. 5. These are compared to the SM expectation using the H1 PDF 2000 fit $`^\mathrm{?}`$. The same figure shows the result of a linear fit to the dependence of the cross section on the longitudinal polarisation describing reasonably the data. Extrapolating the fit to the point $`P=1`$ yields:
$$\sigma _{\mathrm{CC}}(P=1)=3.7\pm 2.4(\mathrm{stat}.)\pm 2.7(\mathrm{syst}.)\mathrm{pb}.$$
This extrapolation is consistent with the SM prediction of zero cross section. The comparison of the H1 results to the ZEUS results scaled to the kinematic region corresponding to the H1 measurement in Fig. 5 shows a good agreement between the two experiments.
The cross-sections $`d\sigma /dQ^2`$ for NC DIS for positive and negative longitudinal polarisations as measured by ZEUS are shown in Fig. 6 together with the ratios of the cross sections for the positive and negative longitudinal polarisations to the unpolarised results $`^\mathrm{?}`$. Also shown is the ratio of the cross sections for the positive and negative longitudinal polarisations. Only statistical uncertainties were considered when taking ratios of the positively and negatively polarised cross sections. In taking ratios to the unpolarised cross sections the systematic uncertainties were considered uncorrelated with those of the polarised cross sections. The measurements are well described by the SM evaluated using the ZEUS-S PDFs and consistent with the expectations of the electroweak SM for polarised NC DIS, although the statistical precision of the current data set does not allow the polarisation effect to be conclusively observed.
## References
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# The Cornell High-order Adaptive Optics Survey for Brown Dwarfs in Stellar Systems–I: Observations, Data Reduction, and Detection Analyses
## 1 Introduction
The discovery of the brown dwarf Gl 229B (Nakajima et al. 1995) heralded a stream of direct detections of sub-stellar objects. Field surveys such as the Two Micron All-Sky Survey (2MASS; Skrutskie et al. 1997), the Deep Near Infrared Survey (DENIS; Epchtein et al. 1997), and the Sloan Digital Sky Survey (SDSS; Gunn and Weinberg 1995) helped raise the number of brown dwarf identifications today to close to a thousand. But despite these advances, the search for brown dwarf companions at intermediate and narrow separations (say less than a few arcseconds) to main sequence stars remains difficult. Despite a strong community effort in high-contrast imaging observations, less than a half-dozen companions have been confirmed as $`<`$100 AU (projected separation) substellar companions to main sequence stars.
To date, the most comprehensive probe of ultra-narrow separation ($``$10 AU) brown dwarf companions comes from radial velocity surveys such as Marcy & Butler (2000) : McCarthy & Zuckerman (2004), for instance, report that over 1500 F, G, K, and M stars have been observed, via this method, with sensitivities strong enough to detect brown dwarf companions between 0 and 5 AU. Using observations like these, Marcy & Butler (2000) conclude that the $``$3 AU brown dwarf companion fraction to F-M stars is less than 0.5$`\%`$. High angular resolution imaging surveys such as CHAOS, a deep adaptive optics (AO) coronagraphic search and proper motion follow-up of faint companions, are required to test if this low companion fraction extends out to intermediate distances, akin to our own outer solar sytem. Knowledge of these intermediate separation companions will help bridge the gap between radial velocity companion surveys and wide-separation companion data such as 2MASS. The CHAOS survey is not the first study to examine this search space. Recently published intermediate separation coronagraphic surveys include Liu et al. (2002), Luhman & Jayawardhana (2002), McCarthy & Zuckerman (2004), Metchev & Hillenbrand (2004), and Potter et al. (2002). Unlike these other surveys however, the CHAOS paper here represents, at the time of this writing, the only published adaptive optics survey that reports complete results for all surveyed targets, including null result observations. Reporting comprehensive results on all surveyed targets, this CHAOS paper invites a rich opportunity for statistical inquiry into brown dwarf and stellar formation theory.
In the sections below we present techniques and results for the recently completed CHAOS survey. Section 2 presents our target sample. Section 3 describes our observing techniques. In Section 4 we present the data analysis techniques we developed for this survey. In Section 5 we summarize our survey sensitivities. Section 6 describes our results. We present our conclusions in Section 7.
## 2 Target Sample
We began our candidate selection process with a careful review of the Third Catalogue of Nearby Stars (Gliese & Jahreiss 1995). Beginning with northern stars, we prioritized targets by their closeness to our solar system. Next we discarded all stars that exist in known resolvable multiple systems, as this scenario would prevent us from effectively hiding the entire parent system behind the 0.$`\mathrm{}`$9 coronagraphic mask. We double-checked for the presence of stellar companions using Hipparcos data (Perryman et al. 1997) as well as on-telescope preliminary imaging. The one known (Perryman et al. 1997) resolvable binary that we kept was Gliese 572; For this target, the secondary star’s narrow separation ($``$0.$`\mathrm{}`$4) allowed us to hide both stars behind the 0.$`\mathrm{}`$9 coronagraphic spot, to an acceptable level. We did not delete spectroscopic binaries from the target list as these unresolved targets allowed for effective coronagraphic masking; Gliese 848, 92, 567, 678, and 688 of our final target list are known spectroscopic binaries, as published in Pourbaix et al. (2004). As the next step, we removed all stars with a V magnitude fainter than $``$ 12 mags. Our previous experience using Palomar Adaptive Optics (PALAO) in 2000 indicated that stars fainter than this limit were unable to serve as effective natural guide stars. Next we searched the USNO-A2.0 Catalogue (Monet et al. 1998) for a corresponding point spread function (PSF) calibration star for each targeted star. For choosing a PSF calibration star, we required the following restrictions: 1) A separation less than a couple degrees from the target star; 2) A difference in V magnitude, relative to the target star, $``$ 1 mag; 3) An absence of any known resolvable companions. These restrictions ensured that the calibration star would deliver a measured PSF similar to the target star’s PSF. Any target star that did not have a corresponding calibration star meeting this criterion was removed from the sample. We expanded the search region further and further south from the original northern positions until the list included a total of 80 stars extending as far south as -10 degrees inclination. This final target sample included 3 A stars, 8 F stars, 13 G stars, 29 K stars, 25 M stars, and 2 stars with ambiguous spectral types. All stars possessed well-characterized proper motion values as defined by Hipparcos (Perryman et al. 1997). As nearby stars, they typically possessed high proper motion (median target proper motion $``$ 600 mas/yr) thus facilitating an efficient common proper motion follow-up strategy for candidate companions. A complete list of the target set is given in Table 1.
In the second paper of this series, we will present a thorough discussion of how selection biases in our sample may affect derived brown dwarf populations. For the time being, however, we do note that certain formation models, such as ones that support the creation of brown dwarfs within multiple systems (Clarke, Reipurth, & Delgado-Donate 2004 for example), imply that observed population levels, as derived from our mostly single-star sample, may differ significantly from statistics that include multiple systems.
## 3 Observations
### 3.1 Coronagraphic Search Observations
To conduct our survey, we used the Palomar Adaptive Optics system (PALAO; Troy et al. 2000) and accompanying PHARO science camera (Hayward et al. 2001) installed on the Palomar 200-inch Hale Telescope. PALAO provided us with the high resolution (FWHM typically $``$ 0$`\mathrm{}`$.14 in K-short) necessary for resolving close companions. The accompanying PHARO science camera (wavelength sensitivity 1-2.5 $`\mathrm{\mu m}`$ and platescale 40 mas per pixel) provided us with a coronagraphic imaging capability along with a field of view ($``$30$`\mathrm{}`$) substantially larger than any competitively sized telescope’s adaptive optics system, at the time of the survey’s commencement.
Our general observing strategy was to align the coronagraphic mask on a target star and take a series of short exposures as to not saturate many pixels in the detector. (Occasionally we saturated at the edges of the coronagraphic mask where high noise levels already prevented any meaningful companion search.) We planned our exposure time and number of exposures to allow for a maximum 8 minutes of execution time (including overheads). This helped ensure that sky conditions did not significantly change between the target exposures and following PSF calibration star exposures. When target star exposures were complete, we spent a similar amount of time taking coronagraphic images of the PSF calibration star. Immediately flanking this target pair, we took dithered images of a nearby empty sky region, using the same set-up as the target and reference star series. We repeated this process (sky, target, reference star, sky) as many times as necessary to reach our desired signal to noise. Once we completed these image sets, we inserted a neutral density filter in the optical path and conducted dithered non-coronagraph exposures of the target star. These images allowed us to characterize and record instrument and site observing conditions. Table 1 describes relevant observing information for the individual targets.
### 3.2 Common Proper Motion Observations
For candidate companions detected in the previous procedures, we checked for a physical companionship by using common proper motion observations. The nearby stars we observed tend to have high proper motions (on the order of a few hundred mas yr<sup>-1</sup>). The vast majority of false candidate companions are background stars that tend to have very small proper motions compared to the parent star. Therefore, after recording our initial measurement, we waited for a timespan long enough for the parent star to move a detectable distance, typically $``$3 sigma separation from the original position. We then repeated our observing set so that we could check to see if the candidate maintained the same position with respect to the parent star. Target stars re-observed to check for common proper motion include Gliese 740, 75, 172, 124, 69, 892, 752, 673, 41, 349, 412, 451, 390, 678, 768, 809, 49, and 688. Due to instrument scheduling constraints, Gliese 49, 41, 390, and 678 were all re-observed using the Palomar 200-inch Wide-field Infrared Camera (WIRC; Wilson et al. 2003) rather than PALAO and the accompanying PHARO science camera. Since the WIRC camera possesses no coronagraphic mode, the observations were instead conducted using standard dithered exposure sequences. WIRC, with its non-AO-corrected point spread function and lack of a coronagraphic mask, made a poorer probe of astrometry than the PHARO camera. However, the systems observed with WIRC all possessed large expected proper motions ($`>`$400 mas \[1 WIRC pixel $``$250 mas\]) and large separations ($`>`$10 arcseconds) from the parent system, making them acceptable WIRC observing targets.
## 4 Data Analysis
### 4.1 Reducing Images
We began our data reduction by median-combining each of the dithered sky sets. We then took each coronagraphed star image and subtracted the median-combined sky taken closest in time to the star image. (The typical separation in time between target and sky image was $``$5.5 minutes.) We divided each of the sky-subtracted star images by a flatfield frame that we created, using standard procedures, from twilight calibration images taken that same night. Next we median-combined each sequence of coronagraphed star frames. For this median-combination, we used the images’ residual parent star flux (that leaked from around the coronagraph) to realign any frames that may have shifted due to instrument flexure. Next we applied a bad pixel algorithm to remove suspicious pixels (defined as any pixel deviating from the surrounding 8 pixels by $``$ 5-sigma) and replace them with the median of their neighbors.
After completing this procedure for both target star and calibration star image sets, we scaled the calibration star PSF so that two 50-90 pixel annuli, one centered on the target star, the other centered on the scaled PSF, exhibited identical median values. Next we multiplied the scaled PSF by test values ranging from 0.20 to 1.76 at 0.04 intervals. For each test value we also tried shifting the scaled PSF -7 to +7 pixels, at integar steps, in each of the x and y directions. From these test combinations we selected the adjusted PSF that most closely resembled the target star, according to a least-squared fit of flux values 50 to 90 pixels from the star center. We next subtracted our adjusted PSF from the target star to arrive at a final image for the set. In the cases where we had multiple target star/calibration star observing set pairs, we co-added the final images, using the residual parent star PSF to correct any misalignments.
As our final data reduction procedure, we applied a Fourier filter to help remove non-point-like features such as unwanted internal instrument reflection and residual parent star flux. The Fourier filter application entailed our multiplying each pixel in a Fourier transformed version of the final image (where the lowest frequencies resided at the center of the array and largest frequencies resided toward the edges) by e$`^{\frac{r23}{34}}`$. r here is the separation, in units of pixels, between a given pixel and the center of the Fourier transformed array. We then applied an inverse Fourier transform to the array to produce the final filtered image. We chose the two aforementioned numerical parameters (23 and 34) of our exponential function after first running test trials on a crowded field image using a generic exponential function e$`^{\frac{rm}{\sigma }}`$. For these trial functions we set m to test values ranging from 5 to 49, at integar values; We set $`\sigma `$ to test values ranging from 1 to 39, again at integar intervals. Testing all combinations, including an equivalent gaussian version as well, we found that e$`^{\frac{r23}{34}}`$ produced the greatest signal-to-noise improvements. For sampled field stars in our tests, signal-to-noise levels improved by about 25% between non-filtered and filtered images. Along with this signal-to-noise improvement, the typical PSF FWHM decreased by about 10% as a result of the Fourier filter application.
### 4.2 Identifying Brown Dwarf Companions
Our first step in identifying brown dwarf companions was to individually inspect each final psf-subtracted and non-psf-subtracted image for any potential companions. By choosing to examine both subtracted and non-subtracted final images, we effectively recognize that the PSF-subtraction improves our ability to identify candidates close to the parent star, but, due to the introduced increased sky noise, makes it more difficult to identify candidates at larger separations from the parent star. For our identifications, the characteristic Palomar adaptive optics “waffle pattern” (see Figure 1) helped distinguish real objects from false ones. Practically, we found that this individual inspection was the most effective method of identifying candidate companions. However, for the purpose of determining quantifiable detection sensitivities, we chose to use an automated detection system as well.
Our automated algorithm operated by centering on every other pixel in the image array and creating there a 0.$`\mathrm{}`$16 diameter flux aperture and 1.$`\mathrm{}`$2-1.$`\mathrm{}`$6 diameter sky annulus. After subtracting any residual sky background, the algorithm approximated a signal-to-noise level by dividing the measured aperture flux by the combined aperture flux Poisson noise and background noise; It approximated background noise from the standard deviation of the sky annulus pixels. In the end it outputted a final array with a signal-to-noise value for each sampled pixel. For each signal-to-noise map, it also generated a map of measured background noise at each position (as estimated from the sky annuli). This outputted noise map essentially reflected the ability of the algorithm to detect (at a given thresh-hold signal-to-noise level) different brightness objects according to position on the array.
After generating maps for a given image, the program selected the signal-to-noise map pixel with the highest value, using a minimum value of five. It recorded the pixel position and then moved on to record the next highest signal to noise value greater than five. After each detection, it voided a 0.$`\mathrm{}`$4 radius around the detected candidate object. This procedure continued until there were no more positions with signal to noise values greater or equal to five. (Of course, for many images, no positions possessed signal to noise levels greater than five.) After the algorithm identified the candidate sources, we re-examined the final images to ensure that the algorithm had indeed detected a true source as opposed to a systematic effect. Again, we searched for the Palomar adaptive optics signature “waffle pattern” to ensure a true physical source. We also made comparisons to images taken at other sources to ensure that the feature was indeed unique to the target image.
We acknowledge that the use of our automated detection routine has some drawbacks. Notably, there are several instances where the algorithm over-estimates the noise level in the non-psf-subtracted and psf-subtracted images. For instance, when examining the non-psf-subtracted images at regions near the parent star, the algorithm can mistake what may be a well-ordered parent star PSF slope for a random fluctuation in background noise. In this instance, fortunately, the subsequent examination of the corresponding psf-subtracted image should ensure that the initial over-estimation in noise does not affect final results. However, such a correction will not occur when the algorithm is hunting around field stars; If a field star happens to fall in the sky annulus, the algorithm will determine that region to have excessively high background noise. Thus, only the brightest candidate objects would be detected near these field star positions. In Section 5 we discuss how we may generate limiting magnitudes and brown dwarf mass limits from these algorithm-generated noise maps.
In cases where we positively identified a potential brown dwarf companion to a parent star, we next estimated its apparent K<sub>s</sub> magnitude, using the non-coronagraphed calibration images of the parent star and published 2MASS K-magnitudes (Skrutskie et al. 1997). Resulting magnitudes are displayed in Table 2. Once we established an apparent K<sub>s</sub> magnitude, we derived a corresponding absolute K<sub>s</sub> magnitude, assuming the candidate had a distance equal to the parent system. Thanks to observational surveys such as Hipparcos (Perryman et al. 1997), all of our parent stars had well-defined parallaxes and therefore distances. With an approximate absolute K<sub>s</sub> magnitude in hand, we combined published brown dwarf observational data (Leggett et al. 2000, Leggett et al. 2002, Burgasser et al. 1999, Burgasser et al. 2000, Burgasser et al. 2002, Burgasser, McElwain, & Kirkpatrick 2003, Geballe et al. 2002, Zapatero et al. 2002, Cuby et al. 1999, Tsvetanov et al. 2000, Strauss et al. 1999, and Nakajima et al. 1995) with theoretical data from Burrows et al. (2001) to extrapolate constraints on the object’s mass. An object whose potential mass fell within acceptable brown dwarf restrictions was designated for common proper motion follow-up observations.
For our follow-up observations, we used Hipparcos published common proper motion values (Hipparcos catalogue; Perryman et al. 1997) to determine the expected movement of the parent system. Since background and field stars are unlikely to possess proper motions identical to the parent system’s, we used common proper motion as a strong support for a physical companionship. To determine the candidate companion’s relative position in different epoch images, we fit a gaussian profile to the candidate companion flux position. For the parent star, we determined position from an extrapolated gaussian profile created from the flux leaking from behind the coronagraphic mask. We could typically constrain the parent star position to within a pixel or two and the candidate position to a fraction of a pixel, depending on the signal-to-noise levels. Measuring the candidate companion’s relative position over the two epochs, we were able to distinguish physical companionships from chance alignments. We record positions in Table 2.
## 5 Survey Sensitivities
### 5.1 Determining Limiting Magnitudes
To quantify detection sensitivities from the algorithm-generated noise maps described in Section 4.2, we looked to determine the faintest detectable magnitude as a function of angular separation from each parent star. We began by sampling each of the final psf-subtracted and non-subtracted noise maps and selecting, for each pixel, the smaller of the two noise values. The resulting composite noise map array therefore reflected the best sensitivities from each of the two final images. Figure 2 displays a sample image sequence, where a psf-subtracted and non-subtracted image are combined to create a composite noise map.
Once we had generated our composite noise maps, we declared an array of sample apparent K<sub>s</sub>-magnitudes extending from 8 to 23 mags at intervals of 0.3 mags. This selection included all potential brown dwarf magnitudes that we were likely to encounter. We do note that some of the lowest luminosity brown dwarfs may have magnitudes dimmer than our 23-magnitude limit. However, since 23 magnitudes was effectively beyond even our most optimistic sensitivity estimates, we did not need to consider anything fainter than that. We next transformed the apparent magnitudes to instrument counts using the parent star calibration data described in Section 3.1.
Returning to the composite noise map, we determined the median values in a series of concentric 0$`\mathrm{}`$.20-thick rings centered on the noise map center. The median values therefore represented typical noise as a function of distance from the central star. For each noise value, we then determined the minimum apparent K<sub>s</sub>-magnitude where signal exceeded the combined Poisson noise and ring noise by a factor greater or equal to 5. In Figure 3 we plot resulting measurements for median survey sensitivities (middle curve), the best 10% of observations (lower curve), and the worst 10% of observations (top curve). Refer to Table 3 for a summary of minimum detectable magnitudes for each of the individual targets.
Another commonly used statistic for describing sensitivities for high-contrast companion surveys is the limiting differential magnitude as a function of angular separation from the parent star. In other words, how many times dimmer may a companion object be before we lose it in the parent star noise? Figure 4 plots differential magnitudes for median survey sensitivities as well as the best and worst 10% of observations.
### 5.2 Mass Sensitivities
Determining sensitivities according to companion mass is complicated by the fact that brown dwarfs of a given mass dim over time. Nonetheless, to get a general idea of detectable masses, we may assume different test ages and then use models by Burrows et al. (2001) or Baraffe et al. (2003) to transform our minimum detectable brightnesses into brown dwarf masses. Figure 5 shows a comparison of median sensitivities assuming 1 Gyr, solar age, and 10 Gyr target ages.
## 6 Results
After conducting all of our data analysis, we concluded that zero systems showed positive evidence of a brown dwarf companion. For Gliese 412 follow-up common proper motion observations, the available observing time was too short for us to positively confirm or reject common proper motion. 2MASS data tells us that, in the Gliese 412 neighborhood, the odds of our finding a field star in the PHARO field of view are about 1%. If it is a true companion, its magnitude would place it somewhere around an L9 dwarf classification. In a survey of 80 target stars, a 1% chance alignment is not particularly unusual, making a field star classification a reasonable potentiality. In the end though, these speculations cannot confirm or reject the presence of a true brown dwarf companion. At this point, we classify it as a non-brown dwarf detection until a time when we may confirm its substellar companion nature. Table 2 presents our discovered field stars meeting the automated detection routine’s sensitivity criteria.
## 7 Discussion
The observational data we have presented here clearly supports speculations of a “brown dwarf desert” at orbital separations comparable to our own outer solar system. However, we emphasize that we cannot definitively assert that a brown dwarf desert exists before applying rigorous Monte Carlo simulations that take into account any observational biases. For example, if the brown dwarf companion population were to have unusually high eccentricities, then the $``$100 AU projected separations that we believe we are investigating could in fact be representative of semi-major axes closer to 10 AU. In that case, the 100 AU (true semi-major axis) brown dwarf companion population could in fact be quite high since the members would spend the majority of their orbit outside of our field of view. To address this issue, we conducted full-scale Monte Carlo simulations that account for the effects of differing orbital parameters. In our upcoming paper (part II of this series) we discuss such population simulations at length. One important early result though of such simulations is that approximate analytical solutions presented in McCarthy & Zuckerman (2004) and Gizis et al. (2001), which assume zero inclination and zero eccentricity, suffer from sytematic observational biases that cause them to dramatically understate their uncertainties. Thus, we caution the reader against firmly asserting a brown dwarf companion desert before reading our entire upcoming analysis.
We thank Jean Mueller, Rick Burruss, Karl Dunscombe, and the Palomar Mountain crew for their support. We also thank Sarah Higdon and James Higdon for conducting observations for us. We thank our anonymous referee for their careful review and helpful suggestions. J. C. C. and S. S. E. were supported in part by NSF CAREER award AST-0328522. Facilities: Palomar Observatory(PALAO/PHARO,WIRC).
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# Preferential and Preferential-discriminative Consequence Relations 11footnote 1 This is an updated version of the paper of the same title published in The Journal of Logic and Computation. This version just contains a better presentation (so the numbering of definitions and propositions is different).
## 1 Introduction
In many situations, we are confronted with incomplete and/or inconsistent information and the classical consequence relation proves to be insufficient. Indeed, in case of inconsistent information, it leads to accept every formula as a conclusion, which amounts to loose the whole information. Therefore, we need other relations leading to non-trivial conclusions in spite of the presence of contradictions. So, several paraconsistent consequence relations have been developed. In the present paper, we will pay attention in particular to certain many-valued ones \[Bel77b, Bel77a, DdC70, CMdA00, dACM02, AA94, AA96, AA98\]. They are defined in frameworks where valuations can assign more than two different truth values to formulas. In fact, they tolerate contradictions within the conclusions, but reject the principle of explosion according to which a single contradiction entails the deduction of every formula.
In case of incomplete information, the classical consequence relation also shows its limits. Indeed, no risk is taken, the conclusions are sure, but too few. We need other relations leading to accept as conclusions formulas that are not necessarily sure, but still plausible. Eventually, some “hasty” conclusions will be rejected later, in the presence of additional information. So, a lot of plausible (generally non-monotonic) consequence relations have been developed. Choice functions are central tools to define plausible relations \[Che54, Arr59, Sen70, AM81, Leh02, Leh01, Sch92, Sch04\]. Indeed, suppose we have at our disposal a function $`\mu `$, called a choice function, which chooses in any set of valuations $`V`$, those elements which are preferred, not necessarily in the absolute sense, but when the valuations in $`V`$ are the only ones under consideration. Then, we can define a plausible consequence relation in the following natural way: a formula $`\alpha `$ follows from a set of formulas $`\mathrm{\Gamma }`$ iff every model for $`\mathrm{\Gamma }`$ chosen by $`\mu `$ is a model for $`\alpha `$.
In the present paper, we put the focus on a particular family of choice functions. Let us present it. Suppose we are given a binary preference relation $``$ on states labelled by valuations (in the style of e.g. \[KLM90, Sch04\]). This defines naturally a choice function. Indeed, choose in any set of valuations $`V`$, each element that labels a state which is $``$-preferred among those states which are labelled by the elements of $`V`$. Those choice functions which can be defined in this manner constitute the aforementioned family. The consequence relations defined by this family will be called preferential consequence relations.
For a long time, research efforts on paraconsistent relations and plausible relations were separated. However, in many applications, the information is both incomplete and inconsistent. For instance, the semantic web or big databases inevitably contain inconsistencies. This can be due to human or material imperfections as well as contradictory sources of information. On the other hand, neither the web nor big databases can contain “all” information. Indeed, there are rules of which the exceptions cannot be enumerated. Also, some information might be left voluntarily vague or in concise form. Consequently, consequence relations that are both paraconsistent and plausible are useful to reason in such applications.
Such relations first appear in e.g. \[Pri91, Bat98, KL92, AA00, KM02\]. The idea begins by taking a many-valued framework to get paraconsistency. Then, only those models that are most preferred according to some particular binary preference relation on valuations (in the style of \[Sho88, Sho87\]) are relevant for making inference, which provides plausibility. In \[AL01b, AL01a\], A. Avron and I. Lev generalized the study to families of binary preference relations which compare two valuations using, for each of them, this part of a certain set of formulas it satisfies. The present paper follows this line of research by combining many-valued frameworks and choice functions.
More explicitly, we will investigate preferential consequence relations in a general framework. According to the different assumptions which will be made about the latter, it will cover various kinds of frameworks, including e.g. the classical propositional one as well as certain many-valued ones. Moreover, in the many-valued frameworks, preferential relations are paraconsistent (in addition to be plausible). However, they do not satisfy the Disjunctive Syllogism (from $`\alpha `$ and $`\neg \alpha \beta `$ we can conclude $`\beta `$), whilst they satisfy it in classical framework.
In addition, we will investigate preferential-discriminative consequence relations. They are defined exactly as the plain version, but any conclusion such that its negation is also a conclusion is rejected. In the classical framework, they do not bring something really new. Indeed, instead of concluding everything in the face of inconsistent information, we will simply conclude nothing. On the other hand, in the many-valued frameworks, where the conclusions are non-trivial even from inconsistent information, the discriminative version will reject the contradictions among them, rendering them all the more rational.
The contribution of the present paper can now be summarized in one sentence: we characterized, in a general framework, several (sub)families of preferential(-discriminative) consequence relations. In many cases, our characterizations are purely syntactic. This has a lot of advantages, let us quote some important ones. Take some syntactic conditions that characterize a family of those consequence relations. This gives a syntactic point of view on this family defined semantically, which enables us to compare it to conditions known on the “market”, and thus to other consequence relations. This can also give rise to questions like: if we modified the conditions in such and such a natural-looking way, what would happen on the semantic side? More generally, this can open the door to questions that would not easily come to mind otherwise or to techniques of proof that could not have been employed in the semantic approach.
Several characterizations can be found in the literature for preferential relations (e.g. \[Gab85, Mak89, Mak94, KLM90, LM92, Leh02, Leh01, Sch92, Sch96, Sch00, Sch04\]). We will provide some new ones, though to do so we have been strongly inspired by techniques of K. Schlechta \[Sch04\]. In fact, our innovation is rather related to the discriminative version. To the author knowledge, the present paper is the first systematic work of characterization for preferential-discriminative consequence relations.
The rest of the paper is organized as follows. In Section 2.1, we introduce our general framework and the different assumptions which sometimes will be made about it. We will see that it covers in particular the many-valued frameworks of the well-known paraconsistent logics $`𝒪𝒰`$ and $`J_3`$. In Section 2.2, we present choice functions and some of their well-known properties. In Section 2.3, we define preferential(-discriminative) consequence relations and give examples in both the classical and the many-valued frameworks. We will also recall a characterization which involves the well-known system $`𝐏`$ of Kraus, Lehmann, and Magidor. In section 3, we provide our characterizations. Finally, we conclude in Section 4.
## 2 Background
### 2.1 Semantic structures
#### 2.1.1 Definitions and properties
We will work with general formulas, valuations, and satisfaction. A similar approach has been taken in two well-known papers \[Mak05, Leh01\].
###### Definition 1
We say that $`𝒮`$ is a semantic structure iff $`𝒮=,𝒱,`$ where $``$ is a set, $`𝒱`$ is a set, and $``$ is a relation on $`𝒱\times `$.
Intuitively, $``$ is a set of formulas, $`𝒱`$ a set of valuations for these formulas, and $``$ a satisfaction relation for these objects (i.e. $`v\alpha `$ means the formula $`\alpha `$ is satisfied in the valuation $`v`$, i.e. $`v`$ is a model for $`\alpha `$).
###### Notation 2
Let $`,𝒱,`$ be a semantic structure, $`\mathrm{\Gamma }`$, and $`V𝒱`$. Then,
$`M_\mathrm{\Gamma }:=\{v𝒱:\alpha \mathrm{\Gamma }`$, $`v\alpha \}`$,
$`T(V):=\{\alpha :VM_\alpha \}`$,
$`𝐃:=\{V𝒱:\mathrm{\Gamma },M_\mathrm{\Gamma }=V\}`$.
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, and $``$ the set of all wffs of $``$. Then,
$`T_d(V):=\{\alpha :VM_\alpha `$ and $`VM_{\neg \alpha }\}`$,
$`T_c(V):=\{\alpha :VM_\alpha `$ and $`VM_{\neg \alpha }\}`$,
$`𝐂:=\{V𝒱:\alpha `$, $`VM_\alpha `$ or $`VM_{\neg \alpha }\}`$.
Intuitively, $`M_\mathrm{\Gamma }`$ is the set of all models for $`\mathrm{\Gamma }`$ and $`T(V)`$ the set of all formulas satisfied in $`V`$. Every element of $`T(V)`$ belongs either to $`T_d(V)`$ or $`T_c(V)`$, according to whether its negation is also in $`T(V)`$. $`𝐃`$ is the set of all those sets of valuations that are definable by a set of formulas and $`𝐂`$ the set of all those sets of valuations that do not satisfy both a formula and its negation. As usual, $`M_{\mathrm{\Gamma },\alpha }`$, $`T(V,v)`$ stand for respectively $`M_{\mathrm{\Gamma }\{\alpha \}}`$, $`T(V\{v\})`$, etc.
###### Remark 3
The notations $`M_\mathrm{\Gamma }`$, $`T(V)`$, etc. should contain the semantic structure on which they are based. To increase readability, we will omit it. There will never be any ambiguity. We will omit similar things with other notations in the sequel, for the same reason.
A semantic structure defines a basic consequence relation:
###### Notation 4
We denote by $`𝒫`$ the power set operator.
Let $`,𝒱,`$ be a semantic structure.
We denote by $``$ the relation on $`𝒫()\times `$ such that $`\mathrm{\Gamma }`$, $`\alpha `$,
$$\mathrm{\Gamma }\alpha \text{iff}M_\mathrm{\Gamma }M_\alpha .$$
Let $``$ be a relation on $`𝒫()\times `$. Then,
$`(\mathrm{\Gamma }):=\{\alpha :\mathrm{\Gamma }\alpha \}`$.
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, $``$ the set of all wffs of $``$, and $`\mathrm{\Gamma }`$.
Then, we say that $`\mathrm{\Gamma }`$ is consistent iff $`\alpha `$, $`\mathrm{\Gamma }⊬\alpha `$ or $`\mathrm{\Gamma }⊬\neg \alpha `$.
The following trivial facts hold, we will use them implicitly in the sequel:
###### Remark 5
Let $`,𝒱,`$ be a semantic structure and $`\mathrm{\Gamma },\mathrm{\Delta }`$. Then:
$`M_{\mathrm{\Gamma },\mathrm{\Delta }}=M_\mathrm{\Gamma }M_\mathrm{\Delta }`$;
$`(\mathrm{\Gamma })=T(M_\mathrm{\Gamma })`$;
$`M_\mathrm{\Gamma }=M_{(\mathrm{\Gamma })}`$;
$`\mathrm{\Gamma }(\mathrm{\Delta })`$ iff $`(\mathrm{\Gamma })(\mathrm{\Delta })`$ iff $`M_\mathrm{\Delta }M_\mathrm{\Gamma }`$.
Sometimes, we will need some of the following assumptions about a semantic structure:
###### Definition 6
Suppose $`,𝒱,`$ is a semantic structure.
Then, define the following assumptions about it:
$`𝒱`$ is finite.
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, and $``$ the set of all wffs of $``$. Then, define:
$`\mathrm{\Gamma }`$, $`\alpha `$, if $`\alpha T(M_\mathrm{\Gamma })`$ and $`\neg \alpha T(M_\mathrm{\Gamma })`$, then $`M_\mathrm{\Gamma }M_\alpha M_{\neg \alpha }`$.
Suppose $``$ and $``$ are binary connectives of $``$. Then, define:
$`\alpha ,\beta `$, we have:
$`M_{\alpha \beta }=M_\alpha M_\beta `$;
$`M_{\alpha \beta }=M_\alpha M_\beta `$;
$`M_{\neg \neg \alpha }=M_\alpha `$;
$`M_{\neg (\alpha \beta )}=M_{\neg \alpha \neg \beta }`$;
$`M_{\neg (\alpha \beta )}=M_{\neg \alpha \neg \beta }`$.
Clearly, those assumptions are satisfied by classical semantic structures, i.e. structures where $``$, $`𝒱`$, and $``$ are classical. In addition, we will see, in Sections 2.1.2 and 2.1.3, that they are satisfied also by certain many-valued semantic structures.
#### 2.1.2 The semantic structure defined by $`𝒪𝒰`$
The logic $`𝒪𝒰`$ was introduced by N. Belnap in \[Bel77a, Bel77b\]. This logic is useful to deal with inconsistent information. Several presentations are possible, depending on the language under consideration. For the needs of the present paper, a classical propositional language will be sufficient. The logic has been investigated intensively in e.g. \[AA94, AA96, AA98\], where richer languages, containing an implication connective $``$ (first introduced by A. Avron \[Avr91\]), were considered.
###### Notation 7
We denote by $`𝒜`$ a set of propositional symbols (or atoms).
We denote by $`_c`$ the classical propositional language containing $`𝒜`$, the usual constants $`false`$ and $`true`$, and the usual connectives $`\neg `$, $``$, and $``$.
We denote by $`_c`$ the set of all wffs of $`_c`$.
We recall a possible meaning for the logic $`𝒪𝒰`$ (more details can be found in \[CLM99, Bel77a, Bel77b\]). Consider a system in which there are, on the one hand, sources of information and, on the other hand, a processor that listens to them. The sources provide information about the atoms only, not about the compound formulas. For each atom $`p`$, there are exactly four possibilities: either the processor is informed (by the sources, taken as a whole) that $`p`$ is true; or he is informed that $`p`$ is false; or he is informed of both; or he has no information about $`p`$.
###### Notation 8
Denote by $`0`$ and $`1`$ the classical truth values and define:
$`𝐟:=\{0\}`$; $`𝐭:=\{1\}`$; $`:=\{0,1\}`$; $`:=\mathrm{}`$.
The global information given by the sources to the processor can be modelled by a function $`s`$ from $`𝒜`$ to $`\{𝐟,𝐭,,\}`$. Intuitively, $`1s(p)`$ means the processor is informed that $`p`$ is true, whilst $`0s(p)`$ means he is informed that $`p`$ is false.
Then, the processor naturally builds information about the compound formulas from $`s`$. Before he starts to do so, the situation can be be modelled by a function $`v`$ from $`_c`$ to $`\{𝐟,𝐭,,\}`$ which agrees with $`s`$ about the atoms and which assigns $``$ to all compound formulas. Now, take $`p`$ and $`q`$ in $`𝒜`$ and suppose $`1v(p)`$ or $`1v(q)`$. Then, the processor naturally adds $`1`$ to $`v(pq)`$. Similarly, if $`0v(p)`$ and $`0v(q)`$, then he adds $`0`$ in $`v(pq)`$. Of course, such rules hold for $`\neg `$ and $``$ too.
Suppose all those rules are applied recursively to all compound formulas. Then, $`v`$ represents the “full” (or developed) information given by the sources to the processor. Now, the valuations of the logic $`𝒪𝒰`$ can be defined as exactly those functions that can be built in this manner (i.e. like $`v`$) from some information sources. More formally,
###### Definition 9
We say that $`v`$ is a four-valued valuation iff $`v`$ is a function from $`_c`$ to $`\{𝐟,𝐭,,\}`$ such that $`v(true)=𝐭`$, $`v(false)=𝐟`$ and $`\alpha ,\beta _c`$,
$`1v(\neg \alpha )`$ iff $`0v(\alpha )`$;
$`0v(\neg \alpha )`$ iff $`1v(\alpha )`$;
$`1v(\alpha \beta )`$ iff $`1v(\alpha )`$ or $`1v(\beta )`$;
$`0v(\alpha \beta )`$ iff $`0v(\alpha )`$ and $`0v(\beta )`$;
$`1v(\alpha \beta )`$ iff $`1v(\alpha )`$ and $`1v(\beta )`$;
$`0v(\alpha \beta )`$ iff $`0v(\alpha )`$ or $`0v(\beta )`$.
We denote by $`𝒱_4`$ the set of all four-valued valuations.
The definition may become more accessible if we see the four-valued valuations as those functions that satisfy Tables 1, 2, and 3 below:
| $`v(\alpha )`$ | $`v(\neg \alpha )`$ |
| --- | --- |
| $`𝐟`$ | $`𝐭`$ |
| $`𝐭`$ | $`𝐟`$ |
| $``$ | $``$ |
| $``$ | $``$ |
| Table 1. | |
| | | $`v(\beta )`$ | | | |
| --- | --- | --- | --- | --- | --- |
| | | $`𝐟`$ | $`𝐭`$ | $``$ | $``$ |
| $`v(\alpha )`$ | $`𝐟`$ | $`𝐟`$ | $`𝐭`$ | $``$ | $``$ |
| | $`𝐭`$ | $`𝐭`$ | $`𝐭`$ | $`𝐭`$ | $`𝐭`$ |
| | $``$ | $``$ | $`𝐭`$ | $``$ | $`𝐭`$ |
| | $``$ | $``$ | $`𝐭`$ | $`𝐭`$ | $``$ |
| | | $`v(\alpha \beta )`$ | | | |
| Table 2. | | | | | |
| --- | --- | --- | --- | --- | --- |
| | | $`v(\beta )`$ | | | |
| --- | --- | --- | --- | --- | --- |
| | | $`𝐟`$ | $`𝐭`$ | $``$ | $``$ |
| $`v(\alpha )`$ | $`𝐟`$ | $`𝐟`$ | $`𝐟`$ | $`𝐟`$ | $`𝐟`$ |
| | $`𝐭`$ | $`𝐟`$ | $`𝐭`$ | $``$ | $``$ |
| | $``$ | $`𝐟`$ | $``$ | $``$ | $`𝐟`$ |
| | $``$ | $`𝐟`$ | $``$ | $`𝐟`$ | $``$ |
| | | $`v(\alpha \beta )`$ | | | |
| Table 3. | | | | | |
| --- | --- | --- | --- | --- | --- |
In the logic $`𝒪𝒰`$, a formula $`\alpha `$ is considered to be satisfied iff the processor is informed that it is true (it does not matter whether he is also informed that $`\alpha `$ is false).
###### Notation 10
We denote by $`_4`$ the relation on $`𝒱_4\times _c`$ such that $`v𝒱_4`$, $`\alpha _c`$, we have
$`v_4\alpha `$ iff $`1v(\alpha )`$.
Proof systems for the consequence relation $``$ based on the semantic structure $`_c,𝒱_4,_4`$ (i.e. the semantic structure defined by $`𝒪𝒰`$) can be found in e.g. \[AA94, AA96, AA98\].
Note that the $`𝒪𝒰`$ semantic structure satisfies $`(A3)`$. In addition, if $`𝒜`$ is finite, then $`(A1)`$ is also satisfied. However, $`(A2)`$ is not satisfied by this structure. In Section 2.1.3, we turn to a many-valued semantic structure which satisfies $`(A2)`$.
#### 2.1.3 The semantic structure defined by $`J_3`$
The logic $`J_3`$ was introduced in \[DdC70\] to answer a question posed in 1948 by S. Jaśkowski, who was interested in systematizing theories capable of containing contradictions, especially if they occur in dialectical reasoning. The step from informal reasoning under contradictions and formal reasoning with databases and information was done in \[CMdA00\] (also specialized for real database models in \[dACM02\]), where another formulation of $`J_3`$ called LFI1 was introduced, and its first-order version, semantics and proof theory were studied in detail. Investigations of $`J_3`$ have also been made in e.g. \[Avr91\], where richer languages than our $`_c`$ were considered.
The valuations of the logic $`J_3`$ can be given the same meaning as those of the logic $`𝒪𝒰`$, except that the consideration is restricted to those sources which always give some information about an atom. More formally,
###### Definition 11
We say that $`v`$ is a three-valued valuation iff $`v`$ is a function from $`_c`$ to $`\{𝐟,𝐭,\}`$ such that $`v(true)=𝐭`$, $`v(false)=𝐟`$ and $`\alpha ,\beta _c`$,
$`1v(\neg \alpha )`$ iff $`0v(\alpha )`$;
$`0v(\neg \alpha )`$ iff $`1v(\alpha )`$;
$`1v(\alpha \beta )`$ iff $`1v(\alpha )`$ or $`1v(\beta )`$;
$`0v(\alpha \beta )`$ iff $`0v(\alpha )`$ and $`0v(\beta )`$;
$`1v(\alpha \beta )`$ iff $`1v(\alpha )`$ and $`1v(\beta )`$;
$`0v(\alpha \beta )`$ iff $`0v(\alpha )`$ or $`0v(\beta )`$.
We denote by $`𝒱_3`$ the set of all three-valued valuations.
As previously, the definition may become more accessible if we see the three-valued valuations as those functions that satisfy Tables 4, 5, and 6 below:
| $`v(\alpha )`$ | $`v(\neg \alpha )`$ |
| --- | --- |
| $`𝐟`$ | $`𝐭`$ |
| $`𝐭`$ | $`𝐟`$ |
| $``$ | $``$ |
| Table 4. | |
| | | $`v(\beta )`$ | | |
| --- | --- | --- | --- | --- |
| | | $`𝐟`$ | $`𝐭`$ | $``$ |
| $`v(\alpha )`$ | $`𝐟`$ | $`𝐟`$ | $`𝐭`$ | $``$ |
| | $`𝐭`$ | $`𝐭`$ | $`𝐭`$ | $`𝐭`$ |
| | $``$ | $``$ | $`𝐭`$ | $``$ |
| | | $`v(\alpha \beta )`$ | | |
| Table 5. | | | | |
| --- | --- | --- | --- | --- |
| | | $`v(\beta )`$ | | |
| --- | --- | --- | --- | --- |
| | | $`𝐟`$ | $`𝐭`$ | $``$ |
| $`v(\alpha )`$ | $`𝐟`$ | $`𝐟`$ | $`𝐟`$ | $`𝐟`$ |
| | $`𝐭`$ | $`𝐟`$ | $`𝐭`$ | $``$ |
| | $``$ | $`𝐟`$ | $``$ | $``$ |
| | | $`v(\alpha \beta )`$ | | |
| Table 6. | | | | |
| --- | --- | --- | --- | --- |
We turn to the satisfaction relation.
###### Notation 12
We denote by $`_3`$ the relation on $`𝒱_3\times _c`$ such that $`v𝒱_3`$, $`\alpha _c`$, we have
$`v_3\alpha `$ iff $`1v(\alpha )`$.
Proof systems for the consequence relation $``$ based on the semantic structure $`_c,𝒱_3,_3`$ (i.e. the semantic structure defined by $`J_3`$) have been provided in e.g. \[Avr91, DdC70\] and chapter IX of \[Eps90\]. The $`J_3`$ structure satisfies $`(A3)`$ and $`(A2)`$. In addition, if $`𝒜`$ is finite, then it satisfies $`(A1)`$ too.
### 2.2 Choice functions
#### 2.2.1 Definitions and properties
In many situations, an agent has some way to choose in any set of valuations $`V`$, those elements that are preferred (the bests, the more normal, etc.), not necessarily in the absolute sense, but when the valuations in $`V`$ are the only ones under consideration. In Social Choice, this is modelled by choice functions \[Che54, Arr59, Sen70, AM81, Leh02, Leh01\].
###### Definition 13
Let $`𝒱`$ be a set, $`𝐕𝒫(𝒱)`$, $`𝐖𝒫(𝒱)`$, and $`\mu `$ a function from $`𝐕`$ to $`𝐖`$.
We say that $`\mu `$ is a choice function iff $`V𝐕`$, $`\mu (V)V`$.
Several properties for choice functions have been put in evidence by researchers in Social Choice. Let us present two important ones (a better presentation can be found in \[Leh01\]). Suppose $`W`$ is a set of valuations, $`V`$ is a subset of $`W`$, and $`vV`$ is a preferred valuation of $`W`$. Then, a natural requirement is that $`v`$ is a preferred valuation of $`V`$. Indeed, in many situations, the larger a set is, the harder it is to be a preferred element of it, and he who can do the most can do the least. This property appears in \[Che54\] and has been given the name Coherence in \[Mou85\].
We turn to the second property. Suppose $`W`$ is a set of valuations, $`V`$ is a subset of $`W`$, and suppose all the preferred valuations of $`W`$ belong to $`V`$. Then, they are expected to include all the preferred valuations of $`V`$. The importance of this property has been put in evidence by \[Aiz85, AM81\] and has been given the name Local Monotonicity in e.g. \[Leh01\].
###### Definition 14
Let $`𝒱`$ be a set, $`𝐕𝒫(𝒱)`$, $`𝐖𝒫(𝒱)`$, and $`\mu `$ a choice function from $`𝐕`$ to $`𝐖`$.
We say that $`\mu `$ is coherent iff $`V,W𝐕`$,
$$\text{if}VW,\text{then}\mu (W)V\mu (V).$$
We say that $`\mu `$ is locally monotonic (LM) iff $`V,W𝐕`$,
$$\text{if}\mu (W)VW,\text{then}\mu (V)\mu (W).$$
In addition to their intuitive meanings, these properties are important because, as was shown by K. Schlechta in \[Sch00\], they characterize those choice functions that can be defined by a binary preference relation on states labelled by valuations (in the style of e.g. \[KLM90\]). We will take a closer look at this in Section 2.2.2.
When a semantic structure is under consideration, two new properties can be defined. Each of them conveys a simple and natural meaning.
###### Definition 15
Let $`,𝒱,`$ be a semantic structure, $`𝐕𝒫(𝒱)`$, $`𝐖𝒫(𝒱)`$, and $`\mu `$ a choice function from $`𝐕`$ to $`𝐖`$.
We say that $`\mu `$ is definability preserving (DP) iff
$$V𝐕𝐃,\mu (V)𝐃.$$
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, and $``$ the set of all wffs of $``$.
We say that $`\mu `$ is coherency preserving (CP) iff
$$V𝐕𝐂,\mu (V)𝐂.$$
Definability Preservation has been put in evidence first in \[Sch92\]. One of its advantages is that when the choice functions under consideration satisfy it, we will provide characterizations with purely syntactic conditions. To the author knowledge, the present paper is the first to introduce Coherency Preservation. An advantage of this property is that when the choice functions under consideration satisfy it, we will not need to assume $`(A2)`$ to show our characterizations (in the discriminative case).
#### 2.2.2 Preference structures
Binary preference relations on valuations have been investigated by e.g. B. Hansson to give semantics for deontic logics \[Han69\]. Y. Shoham rediscovered them to give semantics for plausible non-monotonic logics \[Sho88, Sho87\]. Then, it seems that Imielinski is one of the first persons to introduce binary preference relations on states labelled by valuations \[Imi87\]. They have been used to give more general semantics for plausible non-monotonic logics, see e.g. \[KLM90, LM92, Sch92, Sch96, Sch00, Sch04\]. Let us present them.
###### Definition 16
We say that $``$ is a preference structure on a set $`𝒱`$ iff $`=𝒮,l,`$ where $`𝒮`$ is a set, $`l`$ is a function from $`𝒮`$ to $`𝒱`$, and $``$ is a relation on $`𝒮\times 𝒮`$.
In fact, preference structures are essentially Kripke structures. The difference lies in the interpretation of $``$. In a Kripke structure, it is seen as an accessibility relation, whilst, in a preference structure, it is seen as a preference relation. We recall a possible meaning for preference structures (see e.g. \[KLM90, Sch04\] for details about meaning). Intuitively, $`𝒱`$ is a set of valuations for some language $``$ and $`𝒮`$ a set of valuations for some language $`^{}`$ richer than $``$. The elements of $`𝒮`$ are called states. $`l(s)`$ corresponds precisely to this part of $`s`$ that is about the formulas of $``$ only. We call $`l`$ a labelling function. Finally, $``$ is a preference relation, i.e. $`ss^{}`$ means $`s`$ is preferred to $`s^{}`$.
We turn to well-known properties for preference structures.
###### Definition 17
Suppose $`𝒱`$ is a set, $`=𝒮,l,`$ is a preference structure on $`𝒱`$, $`S𝒮`$, $`sS`$, $`V𝒱`$, and $`𝐕𝒫(𝒱)`$.
We say that $``$ is transitive (resp. irreflexive) iff $``$ is transitive (resp. irreflexive).
We say that $`s`$ is preferred in $`S`$ iff $`s^{}S`$, $`s^{}s`$.
$`L(V):=\{s𝒮:l(s)V\}`$ (intuitively, $`L(V)`$ contains the states labelled by the elements of $`V`$).
We say that $``$ is V-smooth (alias V-stoppered) iff $`V𝐕`$, $`sL(V)`$,
either $`s`$ is preferred in $`L(V)`$ or there exists $`s^{}`$ preferred in $`L(V)`$ such that $`s^{}s`$.
A preference structure defines naturally a choice function. The idea is to choose in any set of valuations $`V`$, each element which labels a state which is preferred among all the states labelled by the elements of $`V`$.
###### Definition 18
Suppose $`=𝒮,l,`$ is a preference structure on a set $`𝒱`$.
We denote by $`\mu _{}`$ the function from $`𝒫(𝒱)`$ to $`𝒫(𝒱)`$ such that $`V𝒱`$,
$$\mu _{}(V)=\{vV:sL(v),s\text{is preferred in}L(V)\}.$$
In \[Sch00\], Schlechta showed that Coherence and Local Monotonicity characterize those choice functions that can be defined by a preference structure. Details are given in the proposition just below. It is an immediate corollary of Proposition 2.4, Proposition 2.15, and Fact 1.3 of \[Sch00\].
###### Proposition 19
Taken from \[Sch00\].
Let $`𝒱`$ be a set, $`𝐕`$ and $`𝐖`$ subsets of $`𝒫(𝒱)`$, and $`\mu `$ a choice function from $`𝐕`$ to $`𝐖`$. Then,
$`\mu `$ is coherent iff there exists a transitive and irreflexive preference structure $``$ on $`𝒱`$ such that $`V𝐕`$, we have $`\mu (V)=\mu _{}(V)`$.
Suppose $`V,W𝐕`$, we have $`VW𝐕`$ and $`VW𝐕`$. Then,
$`\mu `$ is coherent and LM iff there exists a $`𝐕`$-smooth, transitive, and irreflexive preference structure $``$ on $`𝒱`$ such that $`V𝐕`$, we have $`\mu (V)=\mu _{}(V)`$.
In fact, in \[Sch00\], the codomain of $`\mu `$ is required to be its domain: $`𝐕`$. However, this plays no role in the proofs. Therefore, verbatim the same proofs are valid when the codomain of $`\mu `$ is an arbitrary subset $`𝐖`$ of $`𝒫(𝒱)`$. Both myself and Schlechta checked it.
### 2.3 Preferential(-discriminative) consequence relations
#### 2.3.1 Definitions
Suppose we are given a semantic structure and a choice function $`\mu `$ on the valuations. Then, it is natural to conclude a formula $`\alpha `$ from a set of formulas $`\mathrm{\Gamma }`$ iff every model for $`\mathrm{\Gamma }`$ chosen by $`\mu `$ is a model for $`\alpha `$. More formally:
###### Definition 20
Suppose $`𝒮=,𝒱,`$ is a semantic structure and $``$ a relation on $`𝒫()\times `$.
We say that $``$ is a preferential consequence relation iff there exists a coherent choice function $`\mu `$ from $`𝐃`$ to $`𝒫(𝒱)`$ such that $`\mathrm{\Gamma }`$, $`\alpha `$,
$$\mathrm{\Gamma }\alpha \text{iff}\mu (M_\mathrm{\Gamma })M_\alpha .$$
In addition, if $`\mu `$ is LM, DP, etc., then so is $``$.
These consequence relations are called “preferential” because, in the light of Proposition 19, they can be defined equivalently with preference structures, instead of coherent choice functions. They lead to “jump” to plausible conclusions which will eventually be withdrawn later, in the presence of additional information. Therefore, they are useful to deal with incomplete information. We will give an example with a classical semantic structure in Section 2.4.1.
In addition, if a many-valued semantic structure is considered, they lead to rational and non-trivial conclusions is spite of the presence of contradictions and are thus useful to treat both incomplete and inconsistent information. However, they will not satisfy the Disjunctive Syllogism. We will give an example with the $`𝒪𝒰`$ semantic structure in Section 2.4.2.
Now, we turn to a qualified version of preferential consequence. It captures the idea that the contradictions in the conclusions should be rejected.
###### Definition 21
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure, and $``$ a relation on $`𝒫()\times `$.
We say that $``$ is a preferential-discriminative consequence relation iff there is a coherent choice function $`\mu `$ from $`𝐃`$ to $`𝒫(𝒱)`$ such that $`\mathrm{\Gamma }`$, $`\alpha `$,
$$\mathrm{\Gamma }\alpha \text{iff}\mu (M_\mathrm{\Gamma })M_\alpha \text{and}\mu (M_\mathrm{\Gamma })M_{\neg \alpha }.$$
In addition, if $`\mu `$ is LM, DP, etc., then so is $``$.
If a classical semantic structure is considered, the discriminative version does not bring something really new. Indeed, the only difference will be to conclude nothing instead of everything in the face of inconsistent information. On the other hand, with a many-valued structure, the conclusions are rational even from inconsistent information. The discriminative version will then reject the contradictions in the conclusions, rendering the latter all the more rational.
In Definitions 20 and 21, the domain of the choice function is $`𝐃`$. This is natural as only the elements of $`𝐃`$ play a role in the definition of a preferential(-discriminative) consequence relation. This point of view has been adopted in e.g. \[Leh01\] (see Section 6). Now, one might want a definition with choice functions of which the domain is $`𝒫(𝒱)`$. In fact, some families of relations can be defined equivalently with $`𝐃`$ or $`𝒫(𝒱)`$. For instance, as is noted in \[Leh01\], if $`\mu `$ is a coherent choice function from $`𝐃`$ to $`𝒫(𝒱)`$, then the function $`\mu ^{}`$ from $`𝒫(𝒱)`$ to $`𝒫(𝒱)`$ defined by $`\mu ^{}(V)=V\mu (M_{T(V)})`$ is a coherent choice function which agrees with $`\mu `$ on $`𝐃`$.
Several characterizations for preferential consequence relations can be found in the literature (e.g. \[KLM90, LM92, Leh02, Leh01, Sch92, Sch96, Sch00, Sch04\]). In particular, we will recall (in Section 2.4) a characterization that involves the well-known system $`𝐏`$ of \[KLM90\].
As said previously, in the light of Proposition 19, preferential(-discriminative) consequence relations could have been introduced equivalently with preference structures. We opted for coherent choice functions for two reasons. First, they give a clearer meaning. Indeed, properties like Coherence have simple intuitive justifications, whilst preference structures contain “states”, but it is not perfectly clear what a state is in daily life. By the way, in \[KLM90\], Kraus, Lehmann, and Magidor did not consider preference structures to be ontological justifications for their interest in the formal systems investigated, but to be technical tools to study those systems and in particular settle questions of interderivability and find efficient decision procedures (see the end of Section 1.2 of \[KLM90\]).
Second, in the proofs, we will work directly with choice functions and their properties, not with preference structures. By the way, the techniques developed in the present paper (especially in the discriminative case) can certainly be adapted to new properties.
### 2.4 The system $`𝐏`$
Gabbay, Makinson, Kraus, Lehmann, and Magidor investigated extensively properties which should be satisfied by plausible non-monotonic consequence relations \[Gab85, Mak89, Mak94, KLM90, LM92\]. A certain set of properties, called the system $`𝐏`$, plays a central role in this area. It is essentially due to Kraus, Lehmann, and Magidor \[KLM90\] and has been investigated further in \[LM92\]. Let’s present it.
###### Definition 22
Suppose $``$ is a language containing the usual connectives $`\neg `$ and $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure, and $``$ a relation on $`\times `$.
Then, the system $`𝐏`$ is the set of the six following conditions: $`\alpha ,\beta ,\gamma `$,
$`\alpha \alpha `$
$`\begin{array}{c}\underset{¯}{\alpha \beta \alpha \gamma }\\ \beta \gamma \end{array}`$
$`\begin{array}{c}\underset{¯}{\alpha \beta \gamma \alpha }\\ \gamma \beta \end{array}`$
$`\begin{array}{c}\underset{¯}{\alpha \beta \gamma \alpha \beta }\\ \alpha \gamma \end{array}`$
$`\begin{array}{c}\underset{¯}{\alpha \beta \alpha \gamma }\\ \alpha \beta \gamma \end{array}`$
$`\begin{array}{c}\underset{¯}{\alpha \gamma \beta \gamma }\\ \alpha \beta \gamma \end{array}`$
Note that $`\alpha \beta `$ is a shorthand for $`\neg (\neg \alpha \neg \beta )`$. Similarly, $`\alpha \beta `$ and $`\alpha \beta `$ are shorthands. Note again that $`𝐏`$ without $`\mathrm{𝐎𝐫}`$ is called $`𝐂`$. The system $`𝐂`$ is closely related to the cumulative inference which was investigated by Makinson in \[Mak89\]. In addition, it seems to correspond to what Gabbay proposed in \[Gab85\]. Concerning the rule $`\mathrm{𝐎𝐫}`$, it corresponds to the axiom CA of conditional logic.
All the properties in $`𝐏`$ are sound if we read $`\alpha \beta `$ as “$`\beta `$ is a plausible consequence of $`\alpha `$”. In addition, $`𝐏`$ is complete in the sense that it characterizes those consequence relations that can be defined by a smooth transitive irreflexive preference structure. This is what makes $`𝐏`$ central. More formally:
###### Definition 23
Suppose $`,𝒱,`$ is a semantic structure.
Then, $`𝐃_f:=\{V𝒱:\alpha `$, $`V=M_\alpha \}`$.
Suppose $``$ is a language containing the usual connectives $`\neg `$ and $``$, and $``$ the set of all wffs of $``$.
Then define the following condition: $`v𝒱`$, $`\alpha ,\beta `$, $`\mathrm{\Gamma }`$,
$`v\neg \alpha `$ iff $`v\vDash ̸\alpha `$;
$`v\alpha \beta `$ iff $`v\alpha `$ or $`v\beta `$.
if for every finite subset $`\mathrm{\Delta }`$ of $`\mathrm{\Gamma }`$, $`M_\mathrm{\Delta }\mathrm{}`$, then $`M_\mathrm{\Gamma }\mathrm{}`$.
Note that $`(KLM2)`$ is called “assumption of compactness” in \[KLM90\].
###### Proposition 24
\[KLM90\] Suppose $``$ is a language containing the usual connectives $`\neg `$ and $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure satisfying $`(KLM0)`$$`(KLM2)`$, and $``$ a relation of $`\times `$.
Then, $``$ satisfies all the properties of $`𝐏`$ iff there exists a $`𝐃_f`$-smooth transitive irreflexive preference structure $``$ on $`𝒱`$ such that $`\alpha ,\beta `$, $`\alpha \beta `$ iff $`\mu _{}(M_\alpha )M_\beta `$.
Note that $``$ is a relation on $`\times `$, not $`𝒫()\times `$. This difference is crucial. Indeed, if we adapt the conditions of $`𝐏`$ in the obvious way to relations on $`𝒫()\times `$ and if we replace $`𝐃_f`$ by $`𝐃`$ in Proposition 24, then the latter does no longer hold. This negative result was shown by Schlechta in \[Sch92\].
Now, by Propositions 19 and 24, we immediately get the following representation theorem:
###### Proposition 25
Suppose Definition 20 (of preferential consequence relations) is adapted in the obvious way to relations on $`\times `$ (essentially, replace $`𝐃`$ by $`𝐃_f`$), $``$ is a language containing the usual connectives $`\neg `$ and $``$, $``$ the set of all wffs of $``$, $``$ a relation on $`\times `$, and $`,𝒱,`$ a semantic structure such that $`(KLM0)`$$`(KLM2)`$ hold and $`V,W𝐃_f`$, $`VW𝐃_f`$ and $`VW𝐃_f`$.
Then, LM preferential consequence relations are precisely those relations that satisfy the system $`𝐏`$.
#### 2.4.1 Example with a classical semantic structure
Let $``$ be a classical propositional language of which the atoms are $`r`$, $`q`$, and $`p`$. Intuitively, $`r`$ means Nixon is a republican, $`q`$ means Nixon is a quaker, and $`p`$ means Nixon is a pacifist. Let $``$ be the set of all wffs of $``$, $`𝒱`$ the set of all classical two-valued valuations of $``$, and $``$ the classical satisfaction relation for these objects. Then, $`𝒱`$ is the set of the 8 following valuations: $`v_0`$, $`v_1`$, $`v_2`$, $`v_3`$, $`v_4`$, $`v_5`$, $`v_6`$, and $`v_7`$, which are defined in the obvious way by the following table:
$$\begin{array}{cccc}& & & \\ & r& q& p\\ & & & \\ v_0& 0& 0& 0\\ & & & \\ v_1& 0& 0& 1\\ & & & \\ v_2& 0& 1& 0\\ & & & \\ v_3& 0& 1& 1\\ & & & \\ v_4& 1& 0& 0\\ & & & \\ v_5& 1& 0& 1\\ & & & \\ v_6& 1& 1& 0\\ & & & \\ v_7& 1& 1& 1\end{array}$$
Now, consider the class of all republicans and the class of all quakers. Consider that a republican is normal iff he is not a pacifist and that a quaker is normal iff he is a pacifist. And, consider that a valuation $`v`$ is more normal than a valuation $`w`$ from the point of view of a class $`C`$ iff
* Nixon is an individual of $`C`$ in both $`v`$ and $`w`$;
* Nixon is normal in $`v`$;
* Nixon is not normal in $`w`$.
In the following graph, there is an arrow from a valuation $`v`$ to a valuation $`w`$ iff $`v`$ is more normal than $`w`$ from the point of view of some class:
Given those considerations a natural preference structure on $`𝒱`$ is $`=𝒱,l,`$, where $`l`$ is identity and $``$ is the relation such that $`v,w𝒱`$, we have $`vw`$ iff $`(1)`$ or $`(2)`$ below holds (i.e. there is an arrow from $`v`$ to $`w`$):
* $`vr`$ and $`v\neg p`$ and $`wr`$ and $`w\vDash ̸\neg p`$;
* $`vq`$ and $`vp`$ and $`wq`$ and $`w\vDash ̸p`$.
Finally, let $``$ be the preferential consequence relation defined by the coherent choice function $`\mu _{}`$.
Then, $``$ leads us to “jump” to plausible conclusions from incomplete information and to revise previous “hasty” conclusions in the face of new and fuller information. For instance, $`r\neg p`$ and $`\{r,p\}\mid ̸\neg p`$ and $`qp`$ and $`\{q,\neg p\}\mid ̸p`$.
However, $``$ is not paraconsistent. In addition, some sets of formulas are rendered useless, because there is no preferred model for them, though there are models for them. For instance, $`\{q,r\}\alpha `$, $`\alpha `$.
#### 2.4.2 Example with the $`𝒪𝒰`$ semantic structure
Consider the $`𝒪𝒰`$ semantic structure $`_c,𝒱_4,_4`$ and suppose $`𝒜=\{r,q,p\}`$ (these objects have been defined in Section 2.1.2). In addition, make the same considerations about Nixon, the classes, normality, etc., as in Section 2.4.1, except that this time a valuation $`v`$ is considered to be more normal than a valuation $`w`$ from the point of view of a class $`C`$ iff
* in both $`v`$ and $`w`$, the processor is informed that Nixon is an individual of $`C`$;
* in $`v`$, he is informed that Nixon is normal and not informed of the contrary;
* in $`w`$, he is not informed that Nixon is normal.
See Section 2.1.2 for recalls about the sources-processor systems. Given those considerations a natural preference structure on $`𝒱_4`$ is $`=𝒱_4,l,`$, where $`l`$ is identity and $``$ is the relation such that $`v,w𝒱_4`$, we have $`vw`$ iff $`(1)`$ or $`(2)`$ below holds (i.e. $`v`$ is more normal than $`w`$ from the point of view of some class):
* $`vr`$ and $`v\neg p`$ and $`v\vDash ̸p`$ and $`wr`$ and $`w\vDash ̸\neg p`$;
* $`vq`$ and $`vp`$ and $`v\vDash ̸\neg p`$ and $`wq`$ and $`w\vDash ̸p`$.
Let $``$ be the preferential consequence relation defined by the coherent choice function $`\mu _{}`$.
Then, again we “jump” to plausible conclusions and revise previous “hasty” conclusions. For instance, $`r\neg p`$ and $`\{r,p\}\mid ̸\neg p`$ and $`qp`$ and $`\{q,\neg p\}\mid ̸p`$.
In addition, $``$ is paraconsistent. For instance, $`\{p,\neg p,q\}p`$ and $`\{p,\neg p,q\}\neg p`$ and $`\{p,\neg p,q\}q`$ and $`\{p,\neg p,q\}\mid ̸\neg q`$. And, it happens less often that a set of formulas is rendered useless because there is no preferred model for it, though there are models for it. For instance, this time, $`\{q,r\}p`$ and $`\{q,r\}\neg p`$ and $`\{q,r\}q`$ and $`\{q,r\}\mid ̸\neg q`$ and $`\{q,r\}r`$ and $`\{q,r\}\mid ̸\neg r`$.
However, $``$ does not satisfy the Disjunctive Syllogism. Indeed, for instance, $`\{\neg r,rq\}\mid ̸q`$.
## 3 Contributions
The main contributions of the present paper are summarized below. We characterized (in many cases, by purely syntactic conditions) families of preferential and preferential-discriminative consequence relations. Sometimes, we will need to make some assumptions about the semantic structure under consideration. However, no assumption will be needed for the three following families:
* the preferential consequence relations (Section 3.2);
* the DP preferential consequence relations (Section 3.1);
* the DP LM preferential consequence relations (Section 3.1).
We will assume $`(A1)`$ and $`(A3)`$ for:
* the CP preferential-discriminative consequence relations (Section 3.4);
* the CP DP preferential-discriminative consequence relations (Section 3.3);
* the CP DP LM preferential-discriminative consequence relations (Section 3.3).
And, we will need $`(A1)`$, $`(A2)`$, and $`(A3)`$ for:
* the preferential-discriminative consequence relations (Section 3.4);
* the DP preferential-discriminative consequence relations (Section 3.3);
* the DP LM preferential-discriminative consequence relations (Section 3.3).
### 3.1 The non-discriminative and definability preserving case
The characterizations in this section have already been given in Proposition 3.1 of \[Sch00\], under the assumption that a classical propositional semantic structure is considered. Using the same techniques as those of Schlechta, we show easily that his characterizations hold with any semantic structure.
###### Notation 26
Let $`,𝒱,`$ be a semantic structure and $``$ a relation on $`𝒫()\times `$.
Then, consider the following conditions: $`\mathrm{\Gamma },\mathrm{\Delta }`$,
if $`(\mathrm{\Gamma })=(\mathrm{\Delta })`$, then $`(\mathrm{\Gamma })=(\mathrm{\Delta })`$;
$`((\mathrm{\Gamma }))=(\mathrm{\Gamma })`$;
$`\mathrm{\Gamma }(\mathrm{\Gamma })`$;
$`(\mathrm{\Gamma },\mathrm{\Delta })((\mathrm{\Gamma }),\mathrm{\Delta })`$;
if $`\mathrm{\Gamma }(\mathrm{\Delta })(\mathrm{\Gamma })`$, then $`(\mathrm{\Gamma })(\mathrm{\Delta })`$.
Note that those conditions are purely syntactic when there is a proof system available for $``$ (which is the case with e.g. the classical, $`𝒪𝒰`$, and $`J_3`$ semantic structures).
###### Proposition 27
Let $`𝒮=,𝒱,`$ be a semantic structure and $``$ a relation on $`𝒫()\times `$. Then,
$``$ is a DP preferential consequence relation iff $`(`$$`0)`$, $`(`$$`1)`$, $`(`$$`2)`$, and $`(`$$`3)`$ hold;
$``$ is a DP LM preferential consequence relation iff $`(`$$`0)`$, $`(`$$`1)`$, $`(`$$`2)`$, $`(`$$`3)`$, and $`(`$$`4)`$ hold.
###### Proof
Proof of $`(0)`$. Direction: “$``$”.
By hypothesis, there exists a DP coherent choice function $`\mu `$ from $`𝐃`$ to $`𝒫(𝒱)`$ such that $`\mathrm{\Gamma }`$,
$`(\mathrm{\Gamma })=T(\mu (M_\mathrm{\Gamma }))`$. We will show:
$`(0.0)`$$``$ satisfies $`(`$$`0)`$;
$`(0.1)`$$``$ satisfies $`(`$$`1)`$;
$`(0.2)`$$``$ satisfies $`(`$$`2)`$.
Before turning to $`(`$$`3)`$, we need a preliminary result:
$`(0.3)`$$`\mathrm{\Gamma }`$, we have $`\mu (M_\mathrm{\Gamma })=M_{(\mathrm{\Gamma })}`$;
$`(0.4)`$$``$ satisfies $`(`$$`3)`$.
Direction: “$``$”.
Suppose $``$ satisfies $`(`$$`0)`$, $`(`$$`1)`$, $`(`$$`2)`$, and $`(`$$`3)`$.
Let $`\mu `$ be the function from $`𝐃`$ to $`𝒫(𝒱)`$ such that $`\mathrm{\Gamma }`$, $`\mu (M_\mathrm{\Gamma })=M_{(\mathrm{\Gamma })}`$.
Then, $`\mu `$ is well-defined.
Indeed, If $`\mathrm{\Gamma },\mathrm{\Delta }`$ and $`M_\mathrm{\Gamma }=M_\mathrm{\Delta }`$, then $`(\mathrm{\Gamma })=(\mathrm{\Delta })`$, thus, by $`(`$$`0)`$, $`(\mathrm{\Gamma })=(\mathrm{\Delta })`$.
In addition, $`\mu `$ is obviously DP. We show the following which ends the proof:
$`(0.5)`$ $`\mu `$ is a choice function;
$`(0.6)`$ $`\mu `$ is coherent;
$`(0.7)`$ $`\mathrm{\Gamma }`$, we have $`(\mathrm{\Gamma })=T(\mu (M_\mathrm{\Gamma }))`$.
Proof of $`(0.0)`$. Let $`\mathrm{\Gamma },\mathrm{\Delta }`$ and suppose $`(\mathrm{\Gamma })=(\mathrm{\Delta })`$.
Then, $`M_\mathrm{\Gamma }=M_\mathrm{\Delta }`$. Thus, $`(\mathrm{\Gamma })=T(\mu (M_\mathrm{\Gamma }))=T(\mu (M_\mathrm{\Delta }))=(\mathrm{\Delta })`$.
Proof of $`(0.1)`$. Let $`\mathrm{\Gamma }`$. Then, $`((\mathrm{\Gamma }))=(T(\mu (M_\mathrm{\Gamma })))=T(M_{T(\mu (M_\mathrm{\Gamma }))})=(\mathrm{\Gamma })`$.
Proof of $`(0.2)`$. Let $`\mathrm{\Gamma }`$. Then, $`\mathrm{\Gamma }T(M_\mathrm{\Gamma })T(\mu (M_\mathrm{\Gamma }))=(\mathrm{\Gamma })`$.
Proof of $`(0.3)`$. Let $`\mathrm{\Gamma }`$. As, $`\mu `$ is DP, $`\mu (M_\mathrm{\Gamma })𝐃`$.
Thus, $`\mathrm{\Gamma }^{}`$, $`\mu (M_\mathrm{\Gamma })=M_\mathrm{\Gamma }^{}`$. Therefore, $`\mu (M_\mathrm{\Gamma })=M_\mathrm{\Gamma }^{}=M_{T(M_\mathrm{\Gamma }^{})}=M_{T(\mu (M_\mathrm{\Gamma }))}=M_{(\mathrm{\Gamma })}`$.
Proof of $`(0.4)`$. Let $`\mathrm{\Gamma },\mathrm{\Delta }`$. As, $`M_{\mathrm{\Gamma },\mathrm{\Delta }}M_\mathrm{\Gamma }`$ and $`\mu `$ is coherent, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },\mathrm{\Delta }}\mu (M_{\mathrm{\Gamma },\mathrm{\Delta }})`$.
Therefore, $`(\mathrm{\Gamma },\mathrm{\Delta })=T(\mu (M_{\mathrm{\Gamma },\mathrm{\Delta }}))T(\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },\mathrm{\Delta }})=T(\mu (M_\mathrm{\Gamma })M_\mathrm{\Delta })`$.
Thus, by $`(0.0)`$, $`(\mathrm{\Gamma },\mathrm{\Delta })T(M_{(\mathrm{\Gamma })}M_\mathrm{\Delta })=T(M_{(\mathrm{\Gamma }),\mathrm{\Delta }})=((\mathrm{\Gamma }),\mathrm{\Delta })`$.
Proof of $`(0.5)`$. Let $`\mathrm{\Gamma }`$. Then, $`\mu (M_\mathrm{\Gamma })=M_{(\mathrm{\Gamma })}`$, which is, by $`(`$$`2)`$, a subset of $`M_\mathrm{\Gamma }`$.
Proof of $`(0.6)`$. Let $`\mathrm{\Gamma },\mathrm{\Delta }`$ and suppose $`M_\mathrm{\Gamma }M_\mathrm{\Delta }`$.
Then, $`\mu (M_\mathrm{\Delta })M_\mathrm{\Gamma }=M_{(\mathrm{\Delta })}M_\mathrm{\Gamma }=M_{(\mathrm{\Delta }),\mathrm{\Gamma }}`$.
But, by $`(`$$`3)`$, $`M_{(\mathrm{\Delta }),\mathrm{\Gamma }}M_{(\mathrm{\Delta },\mathrm{\Gamma })}=\mu (M_{\mathrm{\Delta },\mathrm{\Gamma }})=\mu (M_\mathrm{\Gamma })`$.
Proof of $`(0.7)`$. Let $`\mathrm{\Gamma }`$. Then, by $`(`$$`1)`$, $`(\mathrm{\Gamma })=((\mathrm{\Gamma }))=T(M_{(\mathrm{\Gamma })})=T(\mu (M_\mathrm{\Gamma }))`$.
Proof of $`(1)`$. Direction: “$``$”.
Verbatim the same proof as for $`(0)`$, except that in addition $`\mu `$ is LM.
We use it to show that $``$ satisfies $`(`$$`4)`$.
Let $`\mathrm{\Gamma },\mathrm{\Delta }`$ and suppose $`\mathrm{\Gamma }(\mathrm{\Delta })(\mathrm{\Gamma })`$.
Then, by $`(0.3)`$, $`\mu (M_\mathrm{\Gamma })=M_{(\mathrm{\Gamma })}M_{(\mathrm{\Delta })}=M_\mathrm{\Delta }M_\mathrm{\Gamma }`$.
Therefore, as $`\mu `$ is locally monotonic, $`\mu (M_\mathrm{\Delta })\mu (M_\mathrm{\Gamma })`$.
Thus, $`(\mathrm{\Gamma })=T(\mu (M_\mathrm{\Gamma }))T(\mu (M_\mathrm{\Delta }))=(\mathrm{\Delta })`$.
Direction: “$``$”.
Verbatim the same proof as for $`(0)`$, except that in addition $`(`$$`4)`$ is satisfied.
We use it to show that $`\mu `$ is locally monotonic.
Let $`\mathrm{\Gamma },\mathrm{\Delta }`$ and suppose $`\mu (M_\mathrm{\Gamma })M_\mathrm{\Delta }M_\mathrm{\Gamma }`$.
Then, $`M_{(\mathrm{\Gamma })}M_\mathrm{\Delta }M_\mathrm{\Gamma }`$. Therefore, $`\mathrm{\Gamma }T(M_\mathrm{\Gamma })T(M_\mathrm{\Delta })=(\mathrm{\Delta })`$.
On the other hand, $`(\mathrm{\Delta })=T(M_\mathrm{\Delta })T(M_{(\mathrm{\Gamma })})=((\mathrm{\Gamma }))`$ which is, by $`(`$$`1)`$, equal to $`(\mathrm{\Gamma })`$.
Thus, by $`(`$$`4)`$, we have $`(\mathrm{\Gamma })(\mathrm{\Delta })`$. Therefore, $`\mu (M_\mathrm{\Delta })=M_{(\mathrm{\Delta })}M_{(\mathrm{\Gamma })}=\mu (M_\mathrm{\Gamma })`$.
### 3.2 The non-discriminative and not necessarily definability preserving case
In this section, we will characterize the family of all preferential consequence relations. Unlike in Section 3.1, our conditions will not be purely syntactic (i.e. using only $``$, $``$, etc.). In fact, properties like Coherence cannot be translated in syntactic terms because the choice functions under consideration are not necessarily definability preserving. Indeed, we do no longer have at our disposal the equality: $`\mu (M_\mathrm{\Gamma })=M_{(\mathrm{\Gamma })}`$, which is of great help to perform the translation and which holds precisely because of Definability Preservation.
In Proposition 5.2.11 of \[Sch04\], K. Schlechta provided a characterization of the aforementioned family, under the assumption that a classical propositional semantic structure is considered. Note that most of his work is done in a very general, in fact algebraic, framework. Only at the end, he applied his general lemmas in a classical framework to get the characterization. The conditions he gave, as ours, are not purely syntactic (e.g. they involve the notion of model, etc.). Moreover, some limits of what can be done in this area have been put in evidence by Schlechta. Approximatively, he showed in Proposition 5.2.15 of the same book that, in an infinite classical framework, there does not exist a characterization containing only conditions which are universally quantified, of limited size, and using only simple operations (like e.g. $``$, $``$, $``$).
The purpose of the present section is to provided a new characterization, more elegant than the one of Schlechta and that hold with any semantic structure. To do so, we have been inspired by the algebraic part of the work of Schlechta (see Proposition 5.2.5 of \[Sch04\]). Technically, the idea begins by building from any function $`f`$, a coherent choice function $`\mu _f`$ such that whenever $`f`$ “covers” some coherent choice function, it necessarily covers $`\mu _f`$.
###### Definition 28
Let $`𝒱`$ be a set, $`𝐕`$ and $`𝐖`$ subsets of $`𝒫(𝒱)`$, and $`f`$ a function from $`𝐕`$ to $`𝐖`$.
We denote by $`\mu _f`$ the function from $`𝐕`$ to $`𝒫(𝒱)`$ such that $`V𝐕`$,
$$\mu _f(V)=\{vV:W𝐕,\text{if}vWV,\text{then}vf(W)\}.$$
###### Lemma 29
Let $`𝒱`$ be a set, $`𝐕`$ and $`𝐖`$ subsets of $`𝒫(𝒱)`$, and $`f`$ a function from $`𝐕`$ to $`𝐖`$.
Then, $`\mu _f`$ is a coherent choice function.
###### Proof
$`\mu _f`$ is obviously a choice function. It remains to show that it is coherent.
Suppose $`V,W𝐕`$, $`VW`$, and $`v\mu _f(W)V`$. We show $`v\mu _f(V)`$.
To do so, suppose the contrary, i.e. suppose $`v\mu _f(V)`$.
Then, as $`vV`$, we have $`Z𝐕`$, $`ZV`$, $`vZ`$, and $`vf(Z)`$.
But, $`VW`$, thus $`ZW`$. Therefore, by definition of $`\mu _f`$, $`v\mu _f(W)`$, which is impossible.
###### Lemma 30
Let $`𝒱`$ be a set, $`𝐕`$, $`𝐖`$, and $`𝐗`$ subsets of $`𝒫(𝒱)`$, $`f`$ a function from $`𝐕`$ to $`𝐖`$, and $`\mu `$ a coherent choice function from $`𝐕`$ to $`𝐗`$ such that $`V𝐕`$, $`f(V)=M_{T(\mu (V))}`$.
Then, $`V𝐕`$, $`f(V)=M_{T(\mu _f(V))}`$.
###### Proof
Let $`V𝐕`$. We show $`f(V)=M_{T(\mu _f(V))}`$.
Case 1: $`v\mu (V)`$, $`v\mu _f(V)`$.
As $`\mu (V)V`$, we have $`vV`$.
Thus, by definition of $`\mu _f`$, $`W𝐕`$, $`WV`$, $`vW`$, and $`vf(W)=M_{T(\mu (W))}\mu (W)`$.
On the other hand, as $`\mu `$ is coherent, $`\mu (V)W\mu (W)`$. Thus, $`v\mu (W)`$, which is impossible.
Case 2: $`\mu (V)\mu _f(V)`$.
Case 2.1: $`v\mu _f(V)`$, $`vf(V)`$.
Then, $`W𝐕`$, $`WV`$, $`vW`$, and $`vf(W)`$. Indeed, just take $`V`$ itself for the choice of $`W`$.
Therefore, $`v\mu _f(V)`$, which is impossible.
Case 2.2: $`\mu _f(V)f(V)`$.
Then, $`f(V)=M_{T(\mu (V))}M_{T(\mu _f(V))}M_{T(f(V))}=M_{T(M_{T(\mu (V))})}=M_{T(\mu (V))}=f(V)`$.
Now, everything is ready to show the representation result.
###### Notation 31
Let $`,𝒱,`$ be a semantic structure and $``$ a relation on $`𝒫()\times `$.
Then, consider the following condition: $`\mathrm{\Gamma }`$,
$`(\mathrm{\Gamma })=T(\{vM_\mathrm{\Gamma }:\mathrm{\Delta }`$, if $`vM_\mathrm{\Delta }M_\mathrm{\Gamma }`$, then $`vM_{(\mathrm{\Delta })}\})`$.
###### Proposition 32
Let $`,𝒱,`$ be a semantic structure and $``$ a relation on $`𝒫()\times `$.
Then, $``$ is a preferential consequence relation iff $`(`$$`5)`$ holds.
###### Proof
Direction: “$``$”.
There exists a coherent choice function $`\mu `$ from $`𝐃`$ to $`𝒫(𝒱)`$ such that $`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T(\mu (M_\mathrm{\Gamma }))`$.
Let $`f`$ be the function from $`𝐃`$ to $`𝐃`$ such that $`V𝐃`$, we have $`f(V)=M_{T(\mu (V))}`$.
By Lemma 30, $`V𝐃`$, we have $`f(V)=M_{T(\mu _f(V))}`$.
Note that $`\mathrm{\Gamma }`$, $`f(M_\mathrm{\Gamma })=M_{T(\mu (M_\mathrm{\Gamma }))}=M_{(\mathrm{\Gamma })}`$.
We show that $`(`$$`5)`$ holds. Let $`\mathrm{\Gamma }`$.
Then, $`(\mathrm{\Gamma })=T(\mu (M_\mathrm{\Gamma }))=T(M_{T(\mu (M_\mathrm{\Gamma }))})=T(f(M_\mathrm{\Gamma }))=T(M_{T(\mu _f(M_\mathrm{\Gamma }))})=T(\mu _f(M_\mathrm{\Gamma }))=T(\{vM_\mathrm{\Gamma }:W𝐃`$, if $`vWM_\mathrm{\Gamma }`$, then $`vf(W)\})=T(\{vM_\mathrm{\Gamma }:\mathrm{\Delta }`$, if $`vM_\mathrm{\Delta }M_\mathrm{\Gamma }`$, then $`vf(M_\mathrm{\Delta })\})=T(\{vM_\mathrm{\Gamma }:\mathrm{\Delta }`$, if $`vM_\mathrm{\Delta }M_\mathrm{\Gamma }`$, then $`vM_{(\mathrm{\Delta })}\})`$.
Direction: “$``$”.
Suppose $``$ satisfies $`(`$$`5)`$.
Let $`f`$ be the function from $`𝐃`$ to $`𝐃`$ such that $`\mathrm{\Gamma }`$, we have $`f(M_\mathrm{\Gamma })=M_{(\mathrm{\Gamma })}`$.
Note that $`f`$ is well-defined. Indeed, if $`\mathrm{\Gamma },\mathrm{\Delta }`$ and $`M_\mathrm{\Gamma }=M_\mathrm{\Delta }`$, then, by $`(`$$`5)`$, $`(\mathrm{\Gamma })=(\mathrm{\Delta })`$.
In addition, by $`(`$$`5)`$, we clearly have $`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T(\mu _f(M_\mathrm{\Gamma }))`$.
And finally, by Lemma 29, $`\mu _f`$ is a coherent choice function.
### 3.3 The discriminative and definability preserving case
In this section, we will characterize certain families of DP preferential-discriminative consequence relations. To do so, we will develop new techniques (especially Lemmas 39 and 40 below). We need basic notations and an inductive construction:
###### Notation 33
$``$ denotes the natural numbers including 0: $`\{0,1,2,\mathrm{},\}`$.
$`^+`$ denotes the strictly positive natural numbers: $`\{1,2,\mathrm{},\}`$.
$``$ denotes the integers.
Let $`i,j`$. Then, $`[i,j]`$ denotes the set of all $`k`$ such that $`ikj`$.
Let $``$ be a language, $``$ a binary connective of $``$, $``$ the set of all wffs of $``$, and $`\beta _1,\beta _2,\mathrm{},\beta _r`$.
Whenever we write $`\beta _1\beta _2\mathrm{}\beta _r`$, we mean $`(\mathrm{}((\beta _1\beta _2)\beta _3)\mathrm{}\beta _{r1})\beta _r`$.
###### Definition 34
Let $``$ be a language, $`\neg `$ a unary connective of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure, $``$ a relation on $`𝒫()\times `$, and $`\mathrm{\Gamma }`$. Then,
$$H_1(\mathrm{\Gamma }):=\{\neg \beta :\beta (\mathrm{\Gamma },(\mathrm{\Gamma }))(\mathrm{\Gamma })\text{and}\neg \beta (\mathrm{\Gamma },(\mathrm{\Gamma }))\}.$$
Let $`i`$ with $`i2`$. Then,
$$H_i(\mathrm{\Gamma }):=\{\neg \beta :\{\begin{array}{c}\beta (\mathrm{\Gamma },(\mathrm{\Gamma }),H_1(\mathrm{\Gamma }),\mathrm{},H_{i1}(\mathrm{\Gamma }))(\mathrm{\Gamma })\text{and}\hfill \\ \neg \beta (\mathrm{\Gamma },(\mathrm{\Gamma }),H_1(\mathrm{\Gamma }),\mathrm{},H_{i1}(\mathrm{\Gamma }))\hfill \end{array}\}.$$
$$H(\mathrm{\Gamma }):=\underset{i^+}{}H_i(\mathrm{\Gamma }).$$
###### Definition 35
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, $``$ a binary connective of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure, and $``$ a relation on $`𝒫()\times `$.
Then, consider the following conditions: $`\mathrm{\Gamma },\mathrm{\Delta }`$, $`\alpha ,\beta `$,
if $`\beta (\mathrm{\Gamma },(\mathrm{\Gamma }))(\mathrm{\Gamma })`$ and $`\neg \alpha (\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta )`$, then $`\alpha (\mathrm{\Gamma })`$;
if $`\alpha (\mathrm{\Gamma },(\mathrm{\Gamma }))(\mathrm{\Gamma })`$ and $`\beta (\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \alpha )(\mathrm{\Gamma })`$, then $`\alpha \beta (\mathrm{\Gamma })`$;
if $`\alpha (\mathrm{\Gamma })`$, then $`\neg \alpha (\mathrm{\Gamma },(\mathrm{\Gamma }))`$;
if $`\mathrm{\Delta }(\mathrm{\Gamma })`$, then $`(\mathrm{\Gamma })H(\mathrm{\Gamma })(\mathrm{\Delta },(\mathrm{\Delta }),H(\mathrm{\Delta }),\mathrm{\Gamma })`$;
if $`\mathrm{\Gamma }(\mathrm{\Delta })(\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma }))`$, then $`(\mathrm{\Gamma })H(\mathrm{\Gamma })(\mathrm{\Delta },(\mathrm{\Delta }),H(\mathrm{\Delta }))`$;
if $`\mathrm{\Gamma }`$ is consistent, then $`(\mathrm{\Gamma })`$ is consistent, $`\mathrm{\Gamma }(\mathrm{\Gamma })`$, and $`((\mathrm{\Gamma }))=(\mathrm{\Gamma })`$.
Note that those conditions are purely syntactic when there is a proof system available for $``$.
###### Proposition 36
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, $``$ and $``$ binary connectives of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure satisfying $`(A1)`$ and $`(A3)`$, and $``$ a relation on $`𝒫()\times `$. Then,
$``$ is a CP DP preferential-discriminative consequence relation iff $`(`$$`0)`$, $`(`$$`6)`$, $`(`$$`7)`$, $`(`$$`8)`$, $`(`$$`9)`$, and $`(`$$`11)`$ hold;
$``$ is a CP DP LM preferential-discriminative consequence relation iff $`(`$$`0)`$, $`(`$$`6)`$, $`(`$$`7)`$, $`(`$$`8)`$, $`(`$$`9)`$, $`(`$$`10)`$, and $`(`$$`11)`$ hold.
Suppose $`,𝒱,`$ satisfies $`(A2)`$ too. Then,
$``$ is a DP preferential-discriminative consequence relation iff $`(`$$`0)`$, $`(`$$`6)`$, $`(`$$`7)`$, $`(`$$`8)`$, and $`(`$$`9)`$ hold;
$``$ is a DP LM preferential-discriminative consequence relation iff $`(`$$`0)`$, $`(`$$`6)`$, $`(`$$`7)`$, $`(`$$`8)`$, $`(`$$`9)`$, and $`(`$$`10)`$ hold.
The proof of Proposition 36 has been relegated at the end of Section 3.3. We need first Notation 33, Definition 37 and Lemmas 38, 39, and 40 below. Here are some purely technical tools:
###### Definition 37
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, $``$ a binary connective of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure satisfying $`(A1)`$, $``$ a relation on $`𝒫()\times `$, and $`\mathrm{\Gamma }`$. Then,
$$M_\mathrm{\Gamma }^1:=\{vM_{\mathrm{\Gamma },(\mathrm{\Gamma })}:\beta T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})(\mathrm{\Gamma }),vM_{\neg \beta }\}.$$
Let $`i`$ with $`i2`$. Then,
$$M_\mathrm{\Gamma }^i:=\{vM_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{i1}:\beta T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{i1})(\mathrm{\Gamma }),vM_{\neg \beta }\}.$$
$$M_\mathrm{\Gamma }^{}:=\underset{i^+}{}M_\mathrm{\Gamma }^i$$
$$n(\mathrm{\Gamma }):=|\{i^+:M_\mathrm{\Gamma }^i\mathrm{}\}|$$
Suppose $`M_\mathrm{\Gamma }^1\mathrm{}`$. Then, we denote by $`\beta _\mathrm{\Gamma }^1`$ an element of $``$, chosen arbitrarily, such that
$`r^+`$, $`v_1,v_2,\mathrm{},v_r𝒱`$, and $`\beta _1,\beta _2,\mathrm{},\beta _r`$ with $`M_\mathrm{\Gamma }^1=\{v_1,v_2,\mathrm{},v_r\}`$,
$$\beta _\mathrm{\Gamma }^1=\beta _1\beta _2\mathrm{}\beta _r,$$
and $`j[1,r]`$, $`\beta _j(\mathrm{\Gamma })`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _j}`$, and $`v_jM_{\neg \beta _j}`$.
As $`M_\mathrm{\Gamma }^1\mathrm{}`$ and $`M_\mathrm{\Gamma }^1`$ is finite (thanks to $`(A1)`$), such an element exists.
Suppose $`i`$, $`i2`$, and $`M_\mathrm{\Gamma }^i\mathrm{}`$.
Then, we denote by $`\beta _\mathrm{\Gamma }^i`$ an element of $``$, chosen arbitrarily, such that
$`r^+`$, $`v_1,v_2,\mathrm{},v_r𝒱`$, and $`\beta _1,\beta _2,\mathrm{},\beta _r`$ with $`M_\mathrm{\Gamma }^i=\{v_1,v_2,\mathrm{},v_r\}`$,
$$\beta _\mathrm{\Gamma }^i=\beta _1\beta _2\mathrm{}\beta _r,$$
and $`j[1,r]`$, $`\beta _j(\mathrm{\Gamma })`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{i1}M_{\beta _j}`$, and $`v_jM_{\neg \beta _j}`$.
As $`M_\mathrm{\Gamma }^i\mathrm{}`$ and $`M_\mathrm{\Gamma }^i`$ is finite, such an element exists.
Suppose $`M_\mathrm{\Gamma }^{}\mathrm{}`$. Then,
$$\beta _\mathrm{\Gamma }:=\beta _\mathrm{\Gamma }^1\beta _\mathrm{\Gamma }^2\mathrm{}\beta _\mathrm{\Gamma }^{n(\mathrm{\Gamma })}$$
As $`M_\mathrm{\Gamma }^{}\mathrm{}`$, $`n(\mathrm{\Gamma })1`$. In addition, we will show in Lemma 38 below that $`n(\mathrm{\Gamma })`$ is finite and $`i^+`$ with $`in(\mathrm{\Gamma })`$, $`M_\mathrm{\Gamma }^i\mathrm{}`$. Thus, $`\beta _\mathrm{\Gamma }`$ is well-defined.
$$F(\mathrm{\Gamma }):=\{\begin{array}{cc}\{\neg \beta _\mathrm{\Gamma }\}\hfill & \text{if}M_\mathrm{\Gamma }^{}\mathrm{}\hfill \\ \mathrm{}\hfill & \text{otherwise}\hfill \end{array}$$
$$G(\mathrm{\Gamma }):=\{\alpha :\alpha (\mathrm{\Gamma }),\neg \alpha (\mathrm{\Gamma }),\text{and}T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha })(\mathrm{\Gamma })\}$$
Here are some quick results about the purely technical tools defined just above:
###### Lemma 38
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, $``$ a binary connective of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure satisfying $`(A1)`$, $``$ a relation on $`𝒫()\times `$, $`\mathrm{\Gamma }`$, and $`i,j^+`$. Then,
if $`ij`$, then $`M_\mathrm{\Gamma }^iM_\mathrm{\Gamma }^j=\mathrm{}`$;
if $`M_\mathrm{\Gamma }^i=\mathrm{}`$, then $`M_\mathrm{\Gamma }^{i+1}=\mathrm{}`$;
$`T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})(\mathrm{\Gamma })`$ iff $`M_\mathrm{\Gamma }^1=\mathrm{}`$;
if $`i2`$, then $`T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{i1})(\mathrm{\Gamma })`$ iff $`M_\mathrm{\Gamma }^i=\mathrm{}`$;
$`n(\mathrm{\Gamma })`$ is finite;
if $`in(\mathrm{\Gamma })`$, then $`M_\mathrm{\Gamma }^i\mathrm{}`$;
if $`i>n(\mathrm{\Gamma })`$, then $`M_\mathrm{\Gamma }^i=\mathrm{}`$;
if $`M_\mathrm{\Gamma }^{}\mathrm{}`$, then $`M_\mathrm{\Gamma }^{}=M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{n(\mathrm{\Gamma })}`$;
$`T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{})(\mathrm{\Gamma })`$.
###### Proof
Proofs of $`(0)`$, $`(1)`$, $`(2)`$, and $`(3)`$. Trivial.
Proof of $`(4)`$. Obvious by $`(0)`$ and $`(A1)`$.
Proof of $`(5)`$. Suppose $`i^+`$, $`M_\mathrm{\Gamma }^i=\mathrm{}`$ and $`in(\mathrm{\Gamma })`$.
Then, by $`(1)`$, $`j^+`$, $`ji`$, $`M_\mathrm{\Gamma }^j=\mathrm{}`$.
Thus, $`|\{j^+:M_\mathrm{\Gamma }^j\mathrm{}\}|i1<n(\mathrm{\Gamma })`$, which is impossible.
Proof of $`(6)`$. Suppose $`i^+`$, $`M_\mathrm{\Gamma }^i\mathrm{}`$ and $`i>n(\mathrm{\Gamma })`$.
Then, by $`(1)`$, $`j^+`$, $`ji`$, $`M_\mathrm{\Gamma }^j\mathrm{}`$.
Thus, $`|\{j^+:M_\mathrm{\Gamma }^j\mathrm{}\}|i>n(\mathrm{\Gamma })`$, which is impossible.
Proof of $`(7)`$. Obvious by $`(6)`$.
Proof of $`(8)`$. Case 1: $`M_\mathrm{\Gamma }^{}=\mathrm{}`$.
Then, $`T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{})=T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$. In addition, $`M_\mathrm{\Gamma }^1=\mathrm{}`$. Thus, by $`(2)`$, we are done.
Case 2: $`M_\mathrm{\Gamma }^{}\mathrm{}`$.
Then, by $`(7)`$, $`T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{})=T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{n(\mathrm{\Gamma })})`$.
In addition, $`n(\mathrm{\Gamma })+12`$ and, by $`(6)`$, $`M_\mathrm{\Gamma }^{n(\mathrm{\Gamma })+1}=\mathrm{}`$. Thus, by $`(3)`$, we are done.
We turn to an important lemma. Its main goal is to show that the conditions $`(`$$`6)`$, $`(`$$`7)`$, and $`(`$$`8)`$ are sufficient to establish the following important equality: $`(\mathrm{\Gamma })=T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })})`$, which provides a semantic definition of $``$ (in the discriminative manner).
###### Lemma 39
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, $``$ and $``$ binary connectives of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure satisfying $`(A1)`$ and $`(A3)`$, $``$ a relation on $`𝒫()\times `$ satisfying $`(`$$`6)`$, $`(`$$`7)`$, and $`(`$$`8)`$, and $`\mathrm{\Gamma }`$. Then,
if $`M_\mathrm{\Gamma }^{}\mathrm{}`$, then $`\beta _\mathrm{\Gamma }(\mathrm{\Gamma })`$;
if $`M_\mathrm{\Gamma }^{}\mathrm{}`$, then $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _\mathrm{\Gamma }}`$;
if $`M_\mathrm{\Gamma }^{}\mathrm{}`$, then $`M_\mathrm{\Gamma }^{}M_{\neg \beta _\mathrm{\Gamma }}=\mathrm{}`$;
if $`M_\mathrm{\Gamma }^{}\mathrm{}`$, then $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{}M_{\neg \beta _\mathrm{\Gamma }}`$;
$`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),F(\mathrm{\Gamma })}`$;
$`(\mathrm{\Gamma })=T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),F(\mathrm{\Gamma })})`$;
$`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),F(\mathrm{\Gamma })}`$;
$`(\mathrm{\Gamma })=T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })})`$.
###### Proof
Proof of $`(0)`$, $`(1)`$, and $`(2)`$. Suppose $`M_\mathrm{\Gamma }^{}\mathrm{}`$.
Then, it suffices to show by induction: $`i[1,n(\mathrm{\Gamma })]`$,
$`p_3(i)`$$`(M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^i)M_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i)}=\mathrm{}`$;
$`p_2(i)`$$`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i}`$;
$`p_1(i)`$$`\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i(\mathrm{\Gamma })`$.
As $`M_\mathrm{\Gamma }^1\mathrm{}`$, $`r^+`$, $`v_1,v_2,\mathrm{},v_r𝒱`$, and $`\beta _1,\beta _2,\mathrm{},\beta _r`$, $`M_\mathrm{\Gamma }^1=\{v_1,\mathrm{},v_r\}`$,
$`\beta _\mathrm{\Gamma }^1=\beta _1\mathrm{}\beta _r`$, and $`j[1,r]`$, $`\beta _j(\mathrm{\Gamma })`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _j}`$, and $`v_jM_{\neg \beta _j}`$.
Then, it can be shown that:
$`(0.0)`$$`p_3(1)`$ holds;
$`(0.1)`$$`p_2(1)`$ holds;
$`(0.2)`$$`p_1(1)`$ holds.
Now, let $`i[1,n(\mathrm{\Gamma })1]`$ and suppose $`p_1(i)`$, $`p_2(i)`$, and $`p_3(i)`$ hold.
As $`M_\mathrm{\Gamma }^{i+1}\mathrm{}`$, $`r^+`$, $`v_1,v_2,\mathrm{},v_r𝒱`$, and $`\beta _1,\beta _2,\mathrm{},\beta _r`$,
$`M_\mathrm{\Gamma }^{i+1}=\{v_1,\mathrm{},v_r\}`$, $`\beta _\mathrm{\Gamma }^{i+1}=\beta _1\mathrm{}\beta _r`$, and
$`j[1,r]`$, $`\beta _j(\mathrm{\Gamma })`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^iM_{\beta _j}`$, and $`v_jM_{\neg \beta _j}`$.
Then, it can be shown that:
$`(0.3)`$$`p_3(i+1)`$ holds;
$`(0.4)`$$`p_2(i+1)`$ holds.
Before turning to $`p_1(i+1)`$, we need the following:
$`(0.5)`$$`\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i\beta _1\beta _2\mathrm{}\beta _r(\mathrm{\Gamma })`$;
$`(0.6)`$$`p_1(i+1)`$ holds.
Proof of $`(0.0)`$. If $`v_jM_\mathrm{\Gamma }^1`$, then $`v_jM_{\neg \beta _j}`$. But, by $`(A3)`$, $`M_{\neg \beta _\mathrm{\Gamma }^1}M_{\neg \beta _j}`$.
Proof of $`(0.1)`$. We have $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _1}`$ which is, by $`(A3)`$, a subset of $`M_{\beta _\mathrm{\Gamma }^1}`$.
Proof of $`(0.2)`$. It suffices to show by induction: $`j[1,r]`$,
$`q(j)`$$`\beta _1\mathrm{}\beta _j(\mathrm{\Gamma })`$.
Obviously, $`q(1)`$ holds.
Let $`j[1,r1]`$. Suppose $`q(j)`$. We show $`q(j+1)`$.
By $`(A3)`$, we have $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _1\mathrm{}\beta _j}`$.
On the other hand, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg (\beta _1\mathrm{}\beta _j)}M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _{j+1}}`$.
Thus, by $`q(j)`$ and $`(`$$`7)`$ (where $`\alpha `$ is $`\beta _1\mathrm{}\beta _j`$ and $`\beta `$ is $`\beta _{j+1}`$), we get $`\beta _1\mathrm{}\beta _{j+1}(\mathrm{\Gamma })`$.
Proof of $`(0.3)`$. Let $`vM_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{i+1}`$. We show $`vM_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^{i+1})}`$.
Case 1: $`vM_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^i`$.
Then, by $`p_3(i)`$, we have $`vM_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i)}`$. But, by $`(A3)`$, $`M_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^{i+1})}M_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i)}`$.
Case 2: $`vM_\mathrm{\Gamma }^{i+1}`$.
Then, $`j[1,r]`$, $`v=v_j`$. Thus, $`vM_{\neg \beta _j}`$. But, by $`(A3)`$, $`M_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^{i+1})}M_{\neg \beta _\mathrm{\Gamma }^{i+1}}M_{\neg \beta _j}`$.
Proof of $`(0.4)`$. By $`p_2(i)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i}`$ which is, by $`(A3)`$, a subset of $`M_{\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^{i+1}}`$.
Proof of $`(0.5)`$. It suffices to show by induction $`j[1,r]`$:
$`q(j)`$$`\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i\beta _1\mathrm{}\beta _j(\mathrm{\Gamma })`$.
We will show:
$`(\mathrm{0.5.0})`$$`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i)}M_{\beta _1}`$.
Then, by $`p_1(i)`$, $`p_2(i)`$, $`(\mathrm{0.5.0})`$, and $`(`$$`7)`$ (where $`\alpha `$ is $`\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i`$ and $`\beta `$ is $`\beta _1`$), $`q(1)`$ holds.
Now, let $`j[1,r1]`$ and suppose $`q(j)`$.
Then, we will show:
$`(\mathrm{0.5.1})`$$`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i\beta _1\mathrm{}\beta _j)}M_{\beta _{j+1}}`$.
In addition, by $`p_2(i)`$ and $`(A3)`$, we get:
$`(\mathrm{0.5.2})`$$`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i\beta _1\mathrm{}\beta _j}`$.
By, $`(\mathrm{0.5.1})`$, $`(\mathrm{0.5.2})`$, $`q(j)`$, and $`(`$$`7)`$ (where $`\alpha `$ is $`\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i\beta _1\mathrm{}\beta _j`$ and $`\beta `$ is $`\beta _{j+1}`$),
we get that $`q(j+1)`$ holds.
Proof of $`(\mathrm{0.5.0})`$. Let $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i)}`$. Then, $`vM_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i)}`$.
Thus, by $`p_3(i)`$, $`vM_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^i`$. Therefore, $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^iM_{\beta _1}`$.
Proof of $`(\mathrm{0.5.1})`$. Let $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i\beta _1\mathrm{}\beta _j)}`$. Then, by $`(A3)`$, $`vM_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i)}`$.
Therefore, by $`p_3(i)`$, $`vM_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^i`$. Therefore, $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^iM_{\beta _{j+1}}`$.
Proof of $`(0.6)`$. By $`p_2(i)`$ and $`(A3)`$, we get $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i}M_{\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i\beta _1\mathrm{}\beta _r}`$.
In addition, by $`(A3)`$, we get $`M_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i\beta _1\mathrm{}\beta _r)}=M_{\neg (\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^{i+1})}`$.
Therefore, by $`(0.5)`$ and $`(`$$`6)`$ (where $`\alpha `$ is $`\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^{i+1}`$ and $`\beta `$ is $`\beta _\mathrm{\Gamma }^1\mathrm{}\beta _\mathrm{\Gamma }^i\beta _1\mathrm{}\beta _r`$),
we get that $`p_1(i+1)`$ holds.
Proof of $`(3)`$. Suppose $`M_\mathrm{\Gamma }^{}\mathrm{}`$, $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{}`$, and $`vM_{\neg \beta _\mathrm{\Gamma }}`$.
Then, by $`(0)`$, $`(1)`$, and definition of $`M_\mathrm{\Gamma }^i`$, we get $`vM_\mathrm{\Gamma }^{n(\mathrm{\Gamma })+1}`$, which is impossible by Lemma 38 $`(6)`$.
Proof of $`(4)`$. Case 1: $`M_\mathrm{\Gamma }^{}\mathrm{}`$.
By $`(3)`$, we get one direction: $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{}M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}`$.
By $`(2)`$, we get the other direction: $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{}`$.
Case 2: $`M_\mathrm{\Gamma }^{}=\mathrm{}`$.
Then, obviously, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{}=M_{\mathrm{\Gamma },(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),F(\mathrm{\Gamma })}`$.
Proof of $`(5)`$. Direction: “$``$”.
Case 1: $`M_\mathrm{\Gamma }^{}\mathrm{}`$.
Suppose the contrary of what we want to show, i.e. suppose $`\alpha (\mathrm{\Gamma })`$, $`\alpha T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }})`$.
Then, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}M_{(\mathrm{\Gamma })}M_\alpha `$. Thus, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}M_{\neg \alpha }`$.
Consequently, by $`(0)`$, $`(1)`$, and $`(`$$`6)`$, we get $`\alpha (\mathrm{\Gamma })`$, which is impossible.
Case 2: $`M_\mathrm{\Gamma }^{}=\mathrm{}`$.
Let $`\alpha (\mathrm{\Gamma })`$. Then, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{(\mathrm{\Gamma })}M_\alpha `$. In addition, by $`(`$$`8)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\neg \alpha }`$.
Consequently, $`\alpha T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})=T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),F(\mathrm{\Gamma })})`$.
Direction: “$``$”. Obvious by $`(4)`$ and Lemma 38 $`(8)`$.
Proof of $`(6)`$. Direction: “$``$”.
Case 1: $`M_\mathrm{\Gamma }^{}=\mathrm{}`$.
Case 1.1: $`H_1(\mathrm{\Gamma })\mathrm{}`$.
Then, $`\alpha `$, $`\alpha (\mathrm{\Gamma })`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\alpha `$, and $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\neg \alpha }`$. Thus, $`\alpha T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$. Therefore, by $`(5)`$, $`\alpha (\mathrm{\Gamma })`$, which is impossible.
Case 1.2: $`H_1(\mathrm{\Gamma })=\mathrm{}`$.
Clearly, $`i^+`$, if $`H_i(\mathrm{\Gamma })=\mathrm{}`$, then $`H_{i+1}(\mathrm{\Gamma })=\mathrm{}`$. Therefore, $`H(\mathrm{\Gamma })=\mathrm{}=F(\mathrm{\Gamma })`$.
Case 2: $`M_\mathrm{\Gamma }^{}\mathrm{}`$.
As, $`M_\mathrm{\Gamma }^{}M_{\mathrm{\Gamma },(\mathrm{\Gamma })}`$, we get, by $`(2)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\neg \beta _\mathrm{\Gamma }}`$.
Thus, by $`(0)`$ and $`(1)`$, we get $`\neg \beta _\mathrm{\Gamma }H_1(\mathrm{\Gamma })H(\mathrm{\Gamma })`$. Therefore, $`M_{H(\mathrm{\Gamma })}M_{F(\mathrm{\Gamma })}`$.
Direction: “$``$”.
Case 1: $`M_\mathrm{\Gamma }^{}=\mathrm{}`$.
Verbatim the proof of Case 1 of direction “$``$”.
Case 2: $`M_\mathrm{\Gamma }^{}\mathrm{}`$.
Then, the following holds:
$`(6.0)`$$`i^+`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H_1(\mathrm{\Gamma }),\mathrm{},H_i(\mathrm{\Gamma })}`$.
Now, suppose the contrary of what we want to show, i.e. suppose
$`vM_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}`$, $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}`$. Then, $`vM_{H(\mathrm{\Gamma })}`$. But, clearly, $`M_{H(\mathrm{\Gamma })}=_{i^+}M_{H_i(\mathrm{\Gamma })}`$. Therefore, $`i^+`$, $`vM_{H_i(\mathrm{\Gamma })}`$, which is impossible by $`(6.0)`$.
Proof of $`(6.0)`$. We show by induction: $`i^+`$,
$`p(i)`$$`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H_1(\mathrm{\Gamma }),\mathrm{},H_i(\mathrm{\Gamma })}`$.
We will show
$`(\mathrm{6.0.0})`$$`p(1)`$ holds.
Let $`i^+`$, suppose $`p(i)`$ holds, and suppose $`p(i+1)`$ does not hold.
Then, $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}`$, $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma }),H_1(\mathrm{\Gamma }),\mathrm{},H_{i+1}(\mathrm{\Gamma })}`$.
Thus, $`j[1,i+1]`$, $`vM_{H_j(\mathrm{\Gamma })}`$.
Case 1: $`j=1`$.
Then, $`\beta `$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\beta `$, $`\beta (\mathrm{\Gamma })`$, and $`vM_{\neg \beta }`$.
Thus $`vM_\mathrm{\Gamma }^1M_{\neg \beta _\mathrm{\Gamma }}`$, which is impossible by $`(2)`$.
Case 2: $`j2`$.
Then, $`\beta `$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H_1(\mathrm{\Gamma }),\mathrm{},H_{j1}(\mathrm{\Gamma })}M_\beta `$, $`\beta (\mathrm{\Gamma })`$, and $`vM_{\neg \beta }`$.
But, by Lemma 38 $`(7)`$, by $`(4)`$, and $`p(i)`$, we get
$`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{n(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H_1(\mathrm{\Gamma }),\mathrm{},H_i(\mathrm{\Gamma })}M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H_1(\mathrm{\Gamma }),\mathrm{},H_{j1}(\mathrm{\Gamma })}M_\beta `$.
Therefore, $`vM_\mathrm{\Gamma }^{n(\mathrm{\Gamma })+1}`$, which is impossible by Lemma 38 $`(6)`$.
Proof of $`(\mathrm{6.0.0})`$. Suppose the contrary of what we want to show, i.e.
suppose $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta _\mathrm{\Gamma }}`$, $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma }),H_1(\mathrm{\Gamma })}`$.
Then, $`vM_{H_1(\mathrm{\Gamma })}`$. Thus, $`\beta `$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\beta `$, $`\beta (\mathrm{\Gamma })`$, and $`vM_{\neg \beta }`$.
Thus $`vM_\mathrm{\Gamma }^1`$. Therefore, $`vM_\mathrm{\Gamma }^{}M_{\neg \beta _\mathrm{\Gamma }}`$, which is impossible by $`(2)`$.
Proof of $`(7)`$. Obvious by $`(5)`$ and $`(6)`$.
We turn to a second important lemma. Its main purpose is to show that any DP choice function $`\mu `$ representing (in the discriminative manner) a relation $``$ satisfies the following equality: $`\mu (M_\mathrm{\Gamma })=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}`$, which enables us to define $`\mu `$ from $``$.
###### Lemma 40
Suppose $``$ is a language, $`\neg `$ a unary connective of $``$, $``$ and $``$ binary connectives of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure satisfying $`(A1)`$ and $`(A3)`$, $`𝐕𝒫(𝒱)`$, $`\mu `$ a DP choice function from $`𝐃`$ to $`𝐕`$, $``$ the relation on $`𝒫()\times `$ such that $`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T_d(\mu (M_\mathrm{\Gamma }))`$, and $`\mathrm{\Gamma }`$. Then:
$`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma })}`$;
$``$ satisfies $`(`$$`6)`$;
$``$ satisfies $`(`$$`7)`$;
$``$ satisfies $`(`$$`8)`$;
$`M_\mathrm{\Gamma }^{}\mu (M_\mathrm{\Gamma })=\mathrm{}`$;
$`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),T_c(\mu (M_\mathrm{\Gamma }))}=\mu (M_\mathrm{\Gamma })`$;
if $`M_\mathrm{\Gamma }^{}\mathrm{}`$, then $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=\mu (M_\mathrm{\Gamma })`$.
If $`,𝒱,`$ satisfies $`(A2)`$ too, then:
if $`M_\mathrm{\Gamma }^{}=\mathrm{}`$, then $`M_{G(\mathrm{\Gamma })}=M_{T_c(\mu (M_\mathrm{\Gamma }))}`$;
if $`M_\mathrm{\Gamma }^{}=\mathrm{}`$, then $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{G(\mathrm{\Gamma })}`$;
$`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=\mu (M_\mathrm{\Gamma })`$.
If $`\mu `$ is coherency preserving, then again:
$`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=\mu (M_\mathrm{\Gamma })`$.
###### Proof
Proof of $`(0)`$. We show $`\mu (M_\mathrm{\Gamma })M_{(\mathrm{\Gamma })}`$. Let $`v\mu (M_\mathrm{\Gamma })`$ and $`\alpha (\mathrm{\Gamma })`$.
Then, $`\alpha T_d(\mu (M_\mathrm{\Gamma }))`$. Thus, $`\mu (M_\mathrm{\Gamma })M_\alpha `$. Thus, $`vM_\alpha `$ and we are done.
In addition, obviously, $`\mu (M_\mathrm{\Gamma })M_\mathrm{\Gamma }`$. Therefore, $`\mu (M_\mathrm{\Gamma })M_\mathrm{\Gamma }M_{(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma })}`$.
Proof of $`(1)`$. Let $`\alpha ,\beta `$ and suppose $`\beta (\mathrm{\Gamma },(\mathrm{\Gamma }))(\mathrm{\Gamma })`$ and $`\neg \alpha (\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta )`$.
Then, by $`(0)`$, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\beta `$. But, $`\beta (\mathrm{\Gamma })=T_d(\mu (M_\mathrm{\Gamma }))`$. Thus, $`\mu (M_\mathrm{\Gamma })M_{\neg \beta }`$.
Consequently, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \beta }M_{\neg \alpha }`$. Therefore, $`\alpha T_d(\mu (M_\mathrm{\Gamma }))=(\mathrm{\Gamma })`$.
Proof of $`(2)`$. Let $`\alpha ,\beta `$ and suppose $`\alpha (\mathrm{\Gamma },(\mathrm{\Gamma }))(\mathrm{\Gamma })`$ and $`\beta (\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \alpha )(\mathrm{\Gamma })`$.
Then, by $`(0)`$, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\alpha `$. But, $`\alpha T_d(\mu (M_\mathrm{\Gamma }))`$. Thus, $`\mu (M_\mathrm{\Gamma })M_{\neg \alpha }`$.
Thus, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\neg \alpha }M_\beta `$. But, $`\beta T_d(\mu (M_\mathrm{\Gamma }))`$. Therefore $`\mu (M_\mathrm{\Gamma })M_{\neg \beta }`$.
Thus, by $`(A3)`$, $`\mu (M_\mathrm{\Gamma })M_{\neg \alpha }M_{\neg \beta }=M_{\neg (\alpha \beta )}`$. Consequently, $`\alpha \beta T_d(\mu (M_\mathrm{\Gamma }))=(\mathrm{\Gamma })`$.
Proof of $`(3)`$. Let $`\alpha (\mathrm{\Gamma })`$. Then, $`\alpha T_d(\mu (M_\mathrm{\Gamma }))`$. Thus, $`\mu (M_\mathrm{\Gamma })M_{\neg \alpha }`$.
Thus, by $`(0)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\neg \alpha }`$.
Proof of $`(4)`$. Case 1: $`M_\mathrm{\Gamma }^{}=\mathrm{}`$. Obvious.
Case 2: $`M_\mathrm{\Gamma }^{}\mathrm{}`$.
It is sufficient to show by induction: $`i[1,n(\mathrm{\Gamma })]`$,
$`p(i)`$$`(M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^i)\mu (M_\mathrm{\Gamma })=\mathrm{}`$.
We will show:
$`(4.0)`$$`p(1)`$ holds.
Let $`i[1,n(\mathrm{\Gamma })1]`$. Suppose $`p(i)`$. We show $`p(i+1)`$.
Case 1: $`M_\mathrm{\Gamma }^{i+1}\mu (M_\mathrm{\Gamma })=\mathrm{}`$.
Then, by $`p(i)`$, we obviously get $`p(i+1)`$.
Case 2: $`vM_\mathrm{\Gamma }^{i+1}\mu (M_\mathrm{\Gamma })`$.
Then, $`\beta `$, $`\beta (\mathrm{\Gamma })`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^iM_\beta `$, and $`vM_{\neg \beta }`$.
Therefore, by $`(0)`$ and $`p(i)`$, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^iM_\beta `$. But, $`\mu (M_\mathrm{\Gamma })M_{\neg \beta }`$.
Consequently, $`\beta T_d(\mu (M_\mathrm{\Gamma }))=(\mathrm{\Gamma })`$, which is impossible.
Proof of $`(4.0)`$. Suppose the contrary of $`p(1)`$, i.e. suppose $`vM_\mathrm{\Gamma }^1\mu (M_\mathrm{\Gamma })`$.
Then, $`\beta `$, $`\beta (\mathrm{\Gamma })`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\beta `$ and $`vM_{\neg \beta }`$.
Therefore, by $`(0)`$, $`\mu (M_\mathrm{\Gamma })M_\beta `$. On the other hand, $`\mu (M_\mathrm{\Gamma })M_{\neg \beta }`$.
Therefore, $`\beta T_d(\mu (M_\mathrm{\Gamma }))=(\mathrm{\Gamma })`$, which is impossible.
Proof of $`(5)`$. As $`\mu (M_\mathrm{\Gamma })𝐃`$, $`\mathrm{\Gamma }^{}`$, $`M_\mathrm{\Gamma }^{}=\mu (M_\mathrm{\Gamma })`$.
Therefore, $`M_{T(\mu (M_\mathrm{\Gamma }))}=M_{T(M_\mathrm{\Gamma }^{})}=M_\mathrm{\Gamma }^{}=\mu (M_\mathrm{\Gamma })`$.
Thus, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),T_c(\mu (M_\mathrm{\Gamma }))}=M_{\mathrm{\Gamma },T_d(\mu (M_\mathrm{\Gamma })),T_c(\mu (M_\mathrm{\Gamma }))}=M_{\mathrm{\Gamma },T(\mu (M_\mathrm{\Gamma }))}`$. But, $`\mathrm{\Gamma }T(\mu (M_\mathrm{\Gamma }))`$.
Therefore, $`M_{\mathrm{\Gamma },T(\mu (M_\mathrm{\Gamma }))}=M_{T(\mu (M_\mathrm{\Gamma }))}=\mu (M_\mathrm{\Gamma })`$.
Proof of $`(6)`$. Suppose $`M_\mathrm{\Gamma }^{}\mathrm{}`$. Direction: “$``$”.
Case 1: $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{n(\mathrm{\Gamma })}`$, $`vM_{T_c(\mu (M_\mathrm{\Gamma }))}`$.
Then, $`\alpha T_c(\mu (M_\mathrm{\Gamma }))`$, $`vM_\alpha `$.
By Lemma 39 $`(3)`$, Lemma 38 $`(7)`$, and $`(A3)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{n(\mathrm{\Gamma })}M_{\neg \beta _\mathrm{\Gamma }}M_{\neg (\beta _\mathrm{\Gamma }\alpha )}`$.
By $`(0)`$ and Lemma 39 $`(1)`$, $`\mu (M_\mathrm{\Gamma })M_{\beta _\mathrm{\Gamma }}M_\alpha =M_{\neg \neg (\beta _\mathrm{\Gamma }\alpha )}`$.
Therefore, $`\neg (\beta _\mathrm{\Gamma }\alpha )T_d(\mu (M_\mathrm{\Gamma }))=(\mathrm{\Gamma })`$.
In addition, $`vM_\alpha M_{\neg \neg (\beta _\mathrm{\Gamma }\alpha )}`$.
Consequently, $`vM_\mathrm{\Gamma }^{n(\mathrm{\Gamma })+1}`$ (take $`\neg (\beta _\mathrm{\Gamma }\alpha )`$ for the $`\beta `$ of the definition of $`M_\mathrm{\Gamma }^i`$).
Therefore, by Lemma 38 $`(6)`$, we get a contradiction.
Case 2: $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{n(\mathrm{\Gamma })}M_{T_c(\mu (M_\mathrm{\Gamma }))}`$.
Then, by Lemma 39 $`(6)`$, Lemma 39 $`(4)`$, Lemma 38 $`(7)`$, and by $`(5)`$, we get
$`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^1\mathrm{}M_\mathrm{\Gamma }^{n(\mathrm{\Gamma })}M_{\mathrm{\Gamma },(\mathrm{\Gamma }),T_c(\mu (M_\mathrm{\Gamma }))}=\mu (M_\mathrm{\Gamma })`$.
Direction: “$``$”.
By $`(0)`$, $`(4)`$, Lemma 39 $`(4)`$, and Lemma 39 $`(6)`$, we get
$`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }^{}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),F(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}`$.
Proof of $`(7)`$. Suppose $`M_\mathrm{\Gamma }^{}=\mathrm{}`$. Direction: “$``$”.
Suppose the contrary of what we want to show, i.e. suppose $`vM_{T_c(\mu (M_\mathrm{\Gamma }))}`$, $`vM_{G(\mathrm{\Gamma })}`$.
Then, $`\alpha G(\mathrm{\Gamma })`$, $`vM_\alpha `$.
Case 1 : $`\alpha T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
As $`\alpha G(\mathrm{\Gamma })`$, $`\alpha (\mathrm{\Gamma })`$. Thus, by Lemma 39 $`(5)`$, $`\alpha T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
Therefore, $`\alpha T_c(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$. Consequently, by $`(0)`$, $`\alpha T_c(\mu (M_\mathrm{\Gamma }))`$.
Thus, $`vM_\alpha `$, which is impossible.
Case 2: $`\neg \alpha T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
As $`\alpha G(\mathrm{\Gamma })`$, $`\neg \alpha (\mathrm{\Gamma })`$. Thus, by Lemma 39 $`(5)`$, $`\neg \alpha T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
Therefore, $`\neg \alpha T_c(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$. Consequently, by $`(A3)`$, $`\alpha T_c(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
Therefore, by $`(0)`$, $`\alpha T_c(\mu (M_\mathrm{\Gamma }))`$. Thus, $`vM_\alpha `$, which is impossible.
Case 3 : $`\alpha T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$ and $`\neg \alpha T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
Then, by $`(A2)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }M_{\neg \alpha }`$. Therefore, $`\alpha T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha })`$.
But, $`\alpha G(\mathrm{\Gamma })`$. Thus, $`T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha })(\mathrm{\Gamma })`$. Thus, $`\alpha (\mathrm{\Gamma })`$. Thus, $`\alpha G(\mathrm{\Gamma })`$, impossible.
Direction: “$``$”.
Suppose the contrary of what we want to show, i.e. suppose $`vM_{G(\mathrm{\Gamma })}`$, $`vM_{T_c(\mu (M_\mathrm{\Gamma }))}`$.
Then, we will show:
$`(7.0)`$$`\alpha T_c(\mu (M_\mathrm{\Gamma }))`$, $`|M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }|<|\mu (M_\mathrm{\Gamma })|`$
But, $`\mu (M_\mathrm{\Gamma })M_\alpha `$ and, by $`(0)`$, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma })}`$. Therefore, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }`$.
Thus, $`|\mu (M_\mathrm{\Gamma })||M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }|`$, which is impossible by $`(7.0)`$.
Proof of $`(7.0)`$. We have $`\delta T_c(\mu (M_\mathrm{\Gamma }))`$, $`vM_\delta `$.
By $`(A1)`$, $`|M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\delta }|`$ is finite. To show $`(7.0)`$, it suffices to show by induction (in the decreasing direction): $`i`$ with $`i|M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\delta }|`$,
$`p(i)`$$`\alpha T_c(\mu (M_\mathrm{\Gamma }))`$, $`vM_\alpha `$ and $`|M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }||\mu (M_\mathrm{\Gamma })|i`$.
Obviously, $`p(|M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\delta }|)`$ holds (take $`\delta `$).
Let $`i`$ with $`i|M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\delta }|`$ and suppose $`p(i)`$ holds. We show $`p(i1)`$.
We have $`\alpha T_c(\mu (M_\mathrm{\Gamma }))`$, $`vM_\alpha `$ and $`|M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }||\mu (M_\mathrm{\Gamma })|i`$.
Case 1: $`T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha })(\mathrm{\Gamma })`$.
As $`\alpha T_c(\mu (M_\mathrm{\Gamma }))`$ and $`(A3)`$ holds, we get $`\neg \alpha T_c(\mu (M_\mathrm{\Gamma }))`$.
But, $`T_c(\mu (M_\mathrm{\Gamma }))T_d(\mu (M_\mathrm{\Gamma }))=\mathrm{}`$. Thus, neither $`\alpha `$ nor $`\neg \alpha `$ belongs to $`T_d(\mu (M_\mathrm{\Gamma }))=(\mathrm{\Gamma })`$.
Consequently, $`\alpha G(\mathrm{\Gamma })`$. Thus, $`vM_\alpha `$, which is impossible.
Case 2: $`\beta T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha })`$, $`\beta (\mathrm{\Gamma })`$.
By $`(0)`$, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma })}`$. On the other hand, $`\mu (M_\mathrm{\Gamma })M_\alpha `$. Thus, $`\mu (M_\mathrm{\Gamma })M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }M_\beta `$.
But, $`\beta (\mathrm{\Gamma })=T_d(\mu (M_\mathrm{\Gamma }))`$. Therefore, $`\mu (M_\mathrm{\Gamma })M_{\neg \beta }`$.
Consequently, $`\mu (M_\mathrm{\Gamma })M_\alpha M_{\neg \beta }=M_{\alpha \neg \beta }`$ and $`\mu (M_\mathrm{\Gamma })M_{\neg \alpha }M_{\neg (\alpha \neg \beta )}`$.
Therefore, $`\alpha \neg \beta T_c(\mu (M_\mathrm{\Gamma }))`$.
Moreover, $`vM_\alpha M_{\alpha \neg \beta }`$.
In addition, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha \neg \beta }M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }`$, whilst $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }M_{\neg \beta }M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha \neg \beta }`$.
Thus $`|M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha \neg \beta }||M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }|1`$. Thus, $`|M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha \neg \beta }||\mu (M_\mathrm{\Gamma })|i1`$.
Therefore, $`p(i1)`$ holds (take $`\alpha \neg \beta `$).
Proof of $`(8)`$. Suppose $`M_\mathrm{\Gamma }^{}=\mathrm{}`$.
Now, suppose the contrary of what we want to show, i.e. $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma })}`$, $`vM_{G(\mathrm{\Gamma })}`$.
Then, $`\alpha G(\mathrm{\Gamma })`$, $`vM_\alpha `$.
Case 1: $`\alpha T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
As, $`\alpha G(\mathrm{\Gamma })`$, $`\alpha (\mathrm{\Gamma })`$. Therefore, by Lemma 39 $`(5)`$, $`\alpha T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
Thus, $`\alpha T_c(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$. Therefore, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\alpha `$. Consequently, $`vM_\alpha `$, which is impossible.
Case 2: $`\neg \alpha T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
As, $`\alpha G(\mathrm{\Gamma })`$, $`\neg \alpha (\mathrm{\Gamma })`$. Therefore, by Lemma 39 $`(5)`$, $`\neg \alpha T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
Thus, $`\neg \alpha T_c(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$. Therefore, by $`(A3)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{\neg \neg \alpha }=M_\alpha `$.
Consequently, $`vM_\alpha `$, which is impossible.
Case 3: $`\alpha T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$ and $`\neg \alpha T(M_{\mathrm{\Gamma },(\mathrm{\Gamma })})`$.
Then, by $`(A2)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha }M_{\neg \alpha }`$. Thus, $`\alpha T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha })`$. But, $`\alpha G(\mathrm{\Gamma })`$. Thus, $`\alpha (\mathrm{\Gamma })`$.
Therefore, $`T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),\alpha })(\mathrm{\Gamma })`$. Consequently, $`\alpha G(\mathrm{\Gamma })`$, which is impossible.
Proof of $`(9)`$. Case 1: $`M_\mathrm{\Gamma }^{}=\mathrm{}`$.
By Lemma 39 $`(6)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),F(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma })}`$.
But, by $`(8)`$, $`(7)`$, and $`(5)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),G(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),T_c(\mu (M_\mathrm{\Gamma }))}=\mu (M_\mathrm{\Gamma })`$.
Case 2: $`M_\mathrm{\Gamma }^{}\mathrm{}`$. Obvious by $`(6)`$.
Proof of $`(10)`$.
Case 1: $`M_\mathrm{\Gamma }^{}=\mathrm{}`$.
Case 1.1: $`vM_{\mathrm{\Gamma },(\mathrm{\Gamma })}`$, $`vM_{T_c(\mu (M_\mathrm{\Gamma }))}`$.
Case 1.1.1: $`\mathrm{\Gamma }`$ is not consistent.
Then, $`\alpha T_c(\mu (M_\mathrm{\Gamma }))`$, $`vM_\alpha `$ and, as $`\mathrm{\Gamma }`$ is not consistent, $`\beta `$, $`M_\mathrm{\Gamma }M_\beta `$ and $`M_\mathrm{\Gamma }M_{\neg \beta }`$.
We have $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\mathrm{\Gamma }M_\beta M_{\beta \neg \alpha }`$.
Moreover, $`\mu (M_\mathrm{\Gamma })M_\mathrm{\Gamma }M_{\neg \beta }`$. Thus, $`\mu (M_\mathrm{\Gamma })M_{\neg \beta }M_\alpha =M_{\neg (\beta \neg \alpha )}`$.
Therefore, $`\beta \neg \alpha T_d(\mu (M_\mathrm{\Gamma }))=(\mathrm{\Gamma })`$.
In addition, $`vM_\alpha M_{\neg (\beta \neg \alpha )}`$.
Consequently, $`vM_\mathrm{\Gamma }^1`$ (take $`\beta \neg \alpha `$ for the $`\beta `$ of the definition of $`M_\mathrm{\Gamma }^1`$).
Thus, $`vM_\mathrm{\Gamma }^{}`$, which is impossible.
Case 1.1.2: $`\mathrm{\Gamma }`$ is consistent.
Thus, $`M_\mathrm{\Gamma }𝐂`$. Therefore, as $`\mu `$ is coherency preserving, $`\mu (M_\mathrm{\Gamma })𝐂`$. Thus, $`T_c(\mu (M_\mathrm{\Gamma }))=\mathrm{}`$.
Therefore, $`M_{T_c(\mu (M_\mathrm{\Gamma }))}=𝒱`$. Thus, $`vM_{T_c(\mu (M_\mathrm{\Gamma }))}`$, which is impossible.
Case 1.2: $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_{T_c(\mu (M_\mathrm{\Gamma }))}`$.
Then, by Lemma 39 $`(6)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),F(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma })}=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),T_c(\mu (M_\mathrm{\Gamma }))}`$.
Therefore, by $`(5)`$, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=\mu (M_\mathrm{\Gamma })`$.
Case 2: $`M_\mathrm{\Gamma }^{}\mathrm{}`$. Obvious by $`(6)`$.
Now comes the proof of Proposition 36 (which is stated at the beginning of Section 3.3).
###### Proof
Proof of $`(0)`$. Direction: “$``$”.
There exists a CP DP coherent choice function $`\mu `$ from $`𝐃`$ to $`𝒫(𝒱)`$ such that
$`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T_d(\mu (M_\mathrm{\Gamma }))`$.
We will show:
$`(0.0)`$$``$ satisfies $`(`$$`0)`$.
By Lemma 40 $`(1)`$, $`(2)`$, and $`(3)`$, $``$ satisfies $`(`$$`6)`$, $`(`$$`7)`$, and $`(`$$`8)`$.
By Lemma 40 $`(10)`$ and Coherence of $`\mu `$, $``$ satisfies $`(`$$`9)`$.
We will show:
$`(0.1)`$$``$ satisfies $`(`$$`11)`$.
Direction: “$``$”.
Suppose $``$ satisfies $`(`$$`0)`$, $`(`$$`6)`$, $`(`$$`7)`$, $`(`$$`8)`$, $`(`$$`9)`$, and $`(`$$`11)`$.
Then, let $`\mu `$ be the function from $`𝐃`$ to $`𝒫(𝒱)`$ such that $`\mathrm{\Gamma }`$, $`\mu (M_\mathrm{\Gamma })=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}`$.
We will show:
$`(0.2)`$$`\mu `$ is well-defined.
Clearly, $`\mu `$ is a DP choice function.
In addition, as $``$ satisfies $`(`$$`9)`$, $`\mu `$ is coherent.
We will show:
$`(0.3)`$$`\mu `$ is CP.
And finally, by Lemma 39 $`(7)`$, $`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T_d(\mu (M_\mathrm{\Gamma }))`$.
Proof of $`(0.0)`$. Let $`\mathrm{\Gamma },\mathrm{\Delta }`$ and suppose $`(\mathrm{\Gamma })=(\mathrm{\Delta })`$. Then, $`M_\mathrm{\Gamma }=M_\mathrm{\Delta }`$.
Therefore, $`(\mathrm{\Gamma })=T_d(\mu (M_\mathrm{\Gamma }))=T_d(\mu (M_\mathrm{\Delta }))=(\mathrm{\Delta })`$.
Proof of $`(0.1)`$. Let $`\mathrm{\Gamma }`$ and suppose $`\mathrm{\Gamma }`$ is consistent.
Then, $`M_\mathrm{\Gamma }𝐃𝐂`$. Thus, as $`\mu `$ is CP, $`\mu (M_\mathrm{\Gamma })𝐂`$. Therefore, $`T_d(\mu (M_\mathrm{\Gamma }))=T(\mu (M_\mathrm{\Gamma }))`$.
Consequently, $`\mathrm{\Gamma }T(M_\mathrm{\Gamma })T(\mu (M_\mathrm{\Gamma }))=T_d(\mu (M_\mathrm{\Gamma }))=(\mathrm{\Gamma })`$.
In addition, $`M_{(\mathrm{\Gamma })}=M_{T_d(\mu (M_\mathrm{\Gamma }))}=M_{T(\mu (M_\mathrm{\Gamma }))}`$. But, $`\mu (M_\mathrm{\Gamma })𝐂`$. Thus, $`M_{T(\mu (M_\mathrm{\Gamma }))}𝐂`$.
Consequently, $`(\mathrm{\Gamma })`$ is consistent.
And finally, $`(\mathrm{\Gamma })=T_d(\mu (M_\mathrm{\Gamma }))=T(\mu (M_\mathrm{\Gamma }))=T(M_{T(\mu (M_\mathrm{\Gamma }))})=T(M_{(\mathrm{\Gamma })})=((\mathrm{\Gamma }))`$.
Proof of $`(0.2)`$. Let $`\mathrm{\Gamma },\mathrm{\Delta }`$ and suppose $`M_\mathrm{\Gamma }=M_\mathrm{\Delta }`$.
Then, $`(\mathrm{\Gamma })=(\mathrm{\Delta })`$. Thus, by $`(`$$`0)`$, $`(\mathrm{\Gamma })=(\mathrm{\Delta })`$.
Consequently, $`H(\mathrm{\Gamma })=H(\mathrm{\Delta })`$. Therefore, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=M_{\mathrm{\Delta },(\mathrm{\Delta }),H(\mathrm{\Delta })}`$.
Proof of $`(0.3)`$. Suppose $`V𝐃𝐂`$. Then, $`\mathrm{\Gamma }`$, $`V=M_\mathrm{\Gamma }`$.
Case 1: $`H_1(\mathrm{\Gamma })\mathrm{}`$.
Thus, $`\beta `$, $`\beta (\mathrm{\Gamma })`$ and $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}M_\beta `$.
By $`(`$$`11)`$, $`\mathrm{\Gamma }(\mathrm{\Gamma })`$ and $`((\mathrm{\Gamma }))=(\mathrm{\Gamma })`$. Thus, $`M_{\mathrm{\Gamma },(\mathrm{\Gamma })}=M_{(\mathrm{\Gamma })}`$. Thus, $`M_{(\mathrm{\Gamma })}M_\beta `$.
Therefore, $`\beta T(M_{(\mathrm{\Gamma })})=((\mathrm{\Gamma }))=(\mathrm{\Gamma })`$, which is impossible.
Case 2: $`H_1(\mathrm{\Gamma })=\mathrm{}`$.
Then, $`H(\mathrm{\Gamma })=\mathrm{}`$. Thus, $`\mu (V)=\mu (M_\mathrm{\Gamma })=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}=M_{(\mathrm{\Gamma })}`$.
But, by $`(`$$`11)`$, $`(\mathrm{\Gamma })`$ is consistent. Therefore, $`M_{(\mathrm{\Gamma })}𝐂`$.
proof of $`(1)`$. Direction: “$``$”.
Verbatim the proof of $`(0)`$, except that in addition $`\mu `$ is LM.
Then, by Lemma 40 $`(10)`$ and LM, $``$ satisfies $`(`$$`10)`$.
Direction: “$``$”.
Verbatim the proof of $`(0)`$, except that in addition $``$ satisfies $`(`$$`10)`$.
Then, by definition of $`\mu `$ and $`(`$$`10)`$, $`\mu `$ is LM.
Proof of $`(2)`$. Direction: “$``$”.
Verbatim the proof of $`(0)`$, except that $`\mu `$ is no longer CP, whilst $`(A2)`$ now holds.
Note that, in $`(0)`$, CP was used only to show $`(`$$`11)`$ and $`(`$$`9)`$.
But, $`(`$$`11)`$ is no longer required to hold.
In addition, by Lemma 40 $`(9)`$ and Coherence of $`\mu `$, $`(`$$`9)`$ holds.
Direction: “$``$”.
Verbatim the proof of $`(0)`$, except that $`(`$$`11)`$ does no longer hold, whilst $`(A2)`$ now holds.
However, in $`(0)`$, $`(`$$`11)`$ was used only to show that $`\mu `$ is CP, which is no longer required.
Note that we do not need to use $`(A2)`$ in this direction.
Proof of $`(3)`$. Direction “$``$”.
Verbatim the proof of $`(0)`$, except that $`\mu `$ is no longer CP, whilst $`\mu `$ is now LM and $`(A2)`$ now holds.
Note that, in $`(0)`$, CP was used only to show $`(`$$`11)`$ and $`(`$$`9)`$.
But, $`(`$$`11)`$ is no longer required.
In addition, by Lemma 40 $`(9)`$ and Coherence of $`\mu `$, $`(`$$`9)`$ holds.
Similarly, by Lemma 40 $`(9)`$ and Local Monotonicity of $`\mu `$, $`(`$$`10)`$ holds.
Direction: “$``$”.
Verbatim the proof of $`(0)`$, except that $`(`$$`11)`$ does no longer hold, whilst $`(`$$`10)`$ and $`(A2)`$ now holds.
Note that, in $`(0)`$, $`(`$$`11)`$ was used only to show that $`\mu `$ is CP, which is no longer required.
Now, by definition of $`\mu `$ and by $`(`$$`10)`$, $`\mu `$ is LM.
Note that we do not need to use $`(A2)`$ in this direction.
### 3.4 The discriminative and not necessarily definability preserving case
Unlike in Section 3.3, the conditions of this section will not be purely syntactic. The translation of properties like Coherence in syntactic terms is blocked because we do no longer have the following useful equality: $`\mu (M_\mathrm{\Gamma })=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}`$, which holds when the choice functions under consideration are definability preserving (but this is not the case here). Thanks to Lemmas 29 and 30 (stated in Section 3.2), we will provide a solution with semi-syntactic conditions.
###### Notation 41
Let $``$ be a language, $`\neg `$ a unary connective of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure, and $``$ a relation on $`𝒫()\times `$.
Then, consider the following condition: $`\mathrm{\Gamma }`$,
$`(\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma }))=T(\{vM_\mathrm{\Gamma }:\mathrm{\Delta },\text{if}vM_\mathrm{\Delta }M_\mathrm{\Gamma },\text{then}vM_{(\mathrm{\Delta }),H(\mathrm{\Delta })}\})`$.
###### Proposition 42
Let $``$ be a language, $`\neg `$ a unary connective of $``$, $``$ and $``$ binary connectives of $``$, $``$ the set of all wffs of $``$, $`,𝒱,`$ a semantic structure satisfying $`(A1)`$ and $`(A3)`$, and $``$ a relation on $`𝒫()\times `$. Then,
$``$ is a CP preferential-discriminative consequence relation iff $`(`$$`0)`$, $`(`$$`6)`$, $`(`$$`7)`$, $`(`$$`8)`$, $`(`$$`11)`$ and $`(`$$`12)`$ hold.
Suppose $`,𝒱,`$ satisfies $`(A2)`$ too. Then,
$``$ is a preferential-discriminative consequence relation iff $`(`$$`0)`$, $`(`$$`6)`$, $`(`$$`7)`$, $`(`$$`8)`$, and $`(`$$`12)`$ hold.
###### Proof
Proof of $`(1)`$. Direction: “$``$”.
There exists a coherent choice function $`\mu `$ from $`𝐃`$ to $`𝒫(𝒱)`$ such that $`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T_d(\mu (M_\mathrm{\Gamma }))`$.
Then, $``$ satisfies obviously $`(`$$`0)`$.
Let $`f`$ be the function from $`𝐃`$ to $`𝐃`$ such that $`V𝐃`$, $`f(V)=M_{T(\mu (V))}`$.
Then, by Lemma 30, $`V𝐃`$, $`f(V)=M_{T(\mu _f(V))}`$.
Moreover, $`\mathrm{\Gamma }`$, $`f(M_\mathrm{\Gamma })=M_{T(\mu (M_\mathrm{\Gamma }))}M_{T(M_\mathrm{\Gamma })}=M_\mathrm{\Gamma }`$.
Therefore, $`f`$ is a choice function.
Obviously, $`f`$ is DP.
In addition, $`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T_d(\mu (M_\mathrm{\Gamma }))=T_d(M_{T(\mu (M_\mathrm{\Gamma }))})=T_d(f(M_\mathrm{\Gamma }))`$.
Consequently, by Lemma 40 $`(1)`$, $`(2)`$, and $`(3)`$, $``$ satisfies $`(`$$`6)`$, $`(`$$`7)`$, and $`(`$$`8)`$.
In addition, by Lemma 40 $`(9)`$, $`\mathrm{\Gamma }`$, $`f(M_\mathrm{\Gamma })=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}`$.
We show that $``$ satisfies $`(`$$`12)`$. Let $`\mathrm{\Gamma }`$.
Then, $`(\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma }))=T(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })})=T(f(M_\mathrm{\Gamma }))=T(M_{T(\mu _f(M_\mathrm{\Gamma }))})=T(\mu _f(M_\mathrm{\Gamma }))=`$
$`T(\{vM_\mathrm{\Gamma }:W𝐃`$, if $`vWM_\mathrm{\Gamma }`$, then $`vf(W)\})=`$
$`T(\{vM_\mathrm{\Gamma }:\mathrm{\Delta }`$, if $`vM_\mathrm{\Delta }M_\mathrm{\Gamma }`$, then $`vf(M_\mathrm{\Delta })\})=`$
$`T(\{vM_\mathrm{\Gamma }:\mathrm{\Delta }`$, if $`vM_\mathrm{\Delta }M_\mathrm{\Gamma }`$, then $`vM_{\mathrm{\Delta },(\mathrm{\Delta }),H(\mathrm{\Delta })}\})=`$
$`T(\{vM_\mathrm{\Gamma }:\mathrm{\Delta },\text{if}vM_\mathrm{\Delta }M_\mathrm{\Gamma },\text{then}vM_{(\mathrm{\Delta }),H(\mathrm{\Delta })}\})`$.
Direction: “$``$”.
Suppose $`(`$$`0)`$, $`(`$$`6)`$, $`(`$$`7)`$, $`(`$$`8)`$, and $`(`$$`12)`$ hold.
Let $`f`$ be the function from $`𝐃`$ to $`𝐃`$ such that $`\mathrm{\Gamma }`$, $`f(M_\mathrm{\Gamma })=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}`$.
By $`(`$$`0)`$, $`f`$ is well-defined.
By Lemma 39 $`(7)`$, $`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T_d(M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })})`$.
Therefore, $`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T_d(f(M_\mathrm{\Gamma }))`$.
By $`(`$$`12)`$, $`\mathrm{\Gamma }`$, $`f(M_\mathrm{\Gamma })=M_{T(\mu _f(M_\mathrm{\Gamma }))}`$.
Therefore, $`\mathrm{\Gamma }`$, $`(\mathrm{\Gamma })=T_d(f(M_\mathrm{\Gamma }))=T_d(M_{T(\mu _f(M_\mathrm{\Gamma }))})=T_d(\mu _f(M_\mathrm{\Gamma }))`$.
But, by Lemma 29, $`\mu _f`$ is a coherent choice function.
Proof of $`(0)`$. Direction: “$``$”.
Verbatim the proof of $`(1)`$, except that $`(A2)`$ does no longer hold, whilst $`\mu `$ is now CP.
Note that $`(A2)`$ was used only to apply Lemma 40 $`(9)`$ to get $`\mathrm{\Gamma }`$, $`f(M_\mathrm{\Gamma })=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}`$.
But, we will get this equality by another mean.
Indeed, if $`V𝐃𝐂`$, then, as $`\mu `$ is CP, $`\mu (V)𝐂`$, thus $`M_{T(\mu (V))}𝐂`$, thus $`f(V)𝐂`$.
Therefore $`f`$ is CP.
Consequently, by Lemma 40 $`(10)`$, we get $`\mathrm{\Gamma }`$, $`f(M_\mathrm{\Gamma })=M_{\mathrm{\Gamma },(\mathrm{\Gamma }),H(\mathrm{\Gamma })}`$.
In addition, by verbatim the proof of $`(0.1)`$ of Proposition 36, $``$ satisfies $`(`$$`11)`$.
Direction: “$``$”.
Verbatim the proof of $`(1)`$, except that $`(A2)`$ does no longer hold, whilst $``$ now satisfies $`(`$$`11)`$.
But, in this direction, $`(A2)`$ was not used in $`(0)`$.
It remains to show that $`\mu _f`$ is CP.
By verbatim the proof of $`(0.3)`$ of Proposition 36, we get that $`f`$ is CP.
Let $`V𝐃𝐂`$. Then, $`f(V)𝐂`$. Thus, $`M_{T(\mu _f(V))}𝐂`$. Thus, $`\mu _f(V)𝐂`$ and we are done.
## 4 Conclusion
We provided, in a general framework, characterizations for families of preferential(-discriminative) consequence relations. Note that we have been strongly inspired by the work of Schlechta in the non-discriminative case, whilst we developed new techniques and ideas in the discriminative case. In many cases, our conditions are purely syntactic. In fact, when the choice functions under consideration are not necessarily definability preserving, we provided solutions with semi-syntactic conditions. We managed to do so thanks to Lemmas 29 and 30. An interesting thing is that we used them both in the plain and the discriminative versions. This suggests that they can be used in yet other versions. In addition, we are quite confident that Lemmas 39 and 40 can be used to characterize other families of consequence relations defined in the discriminative manner by DP choice functions (not necessarily coherent, unlike all the families investigated here).
## 5 Acknowledgements
I owe very much to Karl Schlechta for his hints, advice, constructive criticism, and more. I acknowledge also Arnon Avron for valuable discussions.
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# Spin gases as microscopic models for non-Markovian decoherence
## Abstract
We analyze a microscopic decoherence model in which the total system is described as a spin gas. A spin gas consists of $`N`$ classically moving particles with additional, interacting quantum degrees of freedom (e.g. spins). For various multipartite entangled probe states, we analyze the decoherence induced by interactions between the probe- and environmental spins in such spin gases. We can treat mesoscopic environments ($`10^5`$ particles). We present results for a lattice gas, which could be realized by neutral atoms hopping in an optical lattice, and show the effects of non-Markovian and correlated noise, as well as finite size effects.
Decoherence is a fundamental mechanism believed to be responsible for the transition from the quantum to the classical world Zu03 . Interactions between system degrees of freedom and (uncontrollable) environment degrees of freedom lead to entanglement, manifesting itself in the decoherence of the system state. Many decoherence models have been discussed in the literature, most prominent among them oscillator bath models oscillatorbath and spin bath models Pr00 . Here, we study a physically motivated model of a mesoscopic inhomogeneous spin bath. In particular, we describe the joint state of system and environment by a spin gas Ca05 . A spin gas is a system of quantum spins with stochastic time-dependent interactions. A physical model of a spin gas is a system of $`N`$ classically moving particles with additional, internal spin degrees of freedom. Upon collision, these quantum degrees of freedom interact according to some specified Hamiltonian. Hence, in such spin gases, classical kinematics drives the evolution of the quantum state, and also the decoherence of probe systems. In general, multiple non-consecutive collisions of particles are possible. In this sense, a spin gas provides a microscopic model for non-Markovian decoherence.
In this letter, we determine an effective, time-dependent map that describes the decoherence process in a spin gas. For specific system–environment interactions, we can treat mesoscopic environments ($`10^5`$ particles) exactly and efficiently. To derive the map, we extend the description of certain states in terms of Valence Bond Solids (VBS) Ve04 to completely positive maps. We do not restrict the analysis to single-qubit probes, but also consider the effect of decoherence on various (multipartite) entangled probe states (cf. Ca04 ). Due to the stochastic nature of the interactions, our model does not display symmetries, which otherwise simplify the treatment (see e.g. Zi04 for homogeneous system–environment interactions). Throughout the paper, we concentrate on a specific realization of a spin gas, the spin lattice gas. However, our methods can be easily applied to decoherence in other spin gases, such as a Boltzmann gas Ca05 . Spin gases are not only toy models of theoretical interest, but could be experimentally realized with present-day technology and existing setups expref .
The model: The probe system $`A`$ consists of $`N_A`$ qubits prepared in some arbitrary state. The qubits of system $`A`$ interact with uncontrollable degrees of freedom of an environment $`B`$, leading to decoherence. We consider a microscopic decoherence model where the environment is described by a spin gas. The internal quantum degrees of freedom of this gas interact according to the time-dependent Hamiltonian
$`H(t)={\displaystyle \underset{k<l}{}}g(𝒓_k(t),𝒓_l(t))H_{kl}.`$ (1)
The function $`g`$ depends on the physical nature of the pairwise spin interaction described by the Hamiltonian $`H_{kl}`$. For such systems, we have shown Du05 ; Ca05 that one can efficiently compute reduced density operators of up to ten particles even for mesoscopic system sizes ($`10^5`$ particles), if all Hamiltonians $`H_{kl}`$ commute and the initial state of the system is a pure product state. Here, we extend these results to take arbitrary initial system states and mixed environmental states into account, and thereby study the decoherence of multipartite probe states in a mesoscopic environment efficiently and exactly.
For commuting Hamiltonians $`H_{kl}`$, the joint state of system and environment at a time $`t`$, $`|\mathrm{\Psi }_t=_{k,l}e^{i\phi _{kl}(t)H_{kl}}|\mathrm{\Psi }_0`$, is determined by $`N(N1)/2`$ interaction phases $`\phi _{kl}(t)_0^tg(𝒙_k(t^{}),𝒙_l(t^{}))𝑑t^{}`$. These phases are associated with the adjacency matrix $`\mathrm{\Gamma }(t)`$ of a weighted graph $`G`$. The matrix element $`\mathrm{\Gamma }_{kl}(t)=\phi _{kl}(t)`$ describes the neighborhood relation of particles $`k`$ and $`l`$, or, equivalently, the interaction history. We are interested in the state of system $`A`$, i.e. the reduced density operator $`\rho _A(t)=\mathrm{tr}_B|\mathrm{\Psi }_t\mathrm{\Psi }_t|`$. The commutation of the interaction Hamiltonians greatly simplifies the computation of this operator: To compute this operator we first take into account only interactions between particles $`kA`$ and $`lB`$, since interactions within $`B`$ do not change the state of system $`A`$. In contrast to the general case of non-commuting Hamiltonians Da04 , entanglement in the environment does not influence the decoherence properties of the system. On top of that, the interactions within system $`A`$ itself can be applied to the resulting state afterwards.
In the following, we treat the case $`H_{kl}=|11_{kl}11|`$ and initially completely polarized environment spins, $`|\mathrm{\Psi }_B=|+^N`$ where $`|\pm |0\pm |1`$. Extensions to other commuting interaction Hamiltonians and arbitrary (possibly mixed) product environmental states are straightforward. The isomorphism between completely positive maps and mixed states Ci00 —together with a generalized Valence Bond Solids (VBS) picture Ve04 ; Du05 — determines the effective map $`_t`$ that takes an initial state at $`t_0=0`$ to the state at time $`t`$, i.e. $`\rho _A(t)=_t\rho _A(0)`$. This map $`_t`$ can be equivalently described by the state $`E_t=1𝐥^A^{}_t^A|\mathrm{\Phi }\mathrm{\Phi }|`$. Here, $`|\mathrm{\Phi }=_{k=1}^{|A|}|\varphi ^+_{k^{}k}`$, $`|\varphi ^+|00+|11`$, and $`A^{}`$ is an auxiliary system with the same dimension as $`A`$. We can express the map in the Pauli basis $`\sigma _k`$, where $`\sigma _01𝐥`$. Then, $`_t\rho =\lambda _{k_1\mathrm{}k_{N_A},l_1\mathrm{}l_{N_A}}\sigma _{k_1}\mathrm{}\sigma _{k_{N_A}}\rho \sigma _{l_1}\mathrm{}\sigma _{l_{N_A}}`$, where $`\lambda _{k_1\mathrm{}k_{N_A},l_1\mathrm{}l_{N_A}}=\varphi _{k_1\mathrm{}k_{N_A}}|E_t|\varphi _{l_1\mathrm{}l_{N_A}}`$ with $`|\varphi _{k_1\mathrm{}k_{N_A}}=1𝐥^𝑨^{}(\sigma _{k_1}\mathrm{}\sigma _{k_{N_A}})^A|\mathrm{\Phi }`$. The coefficients of the map are given by the coefficients of the state $`E_t`$ written in (tensor products of) Bell bases. As in the case of states Du05 , we can separately consider maps (or equivalently the states $`E_t^{(l)}`$) resulting from the interaction of the system with a single particle $`l`$ in the environment. We find $`E_t^{(l)}=1/2(|\mathrm{\Phi }\mathrm{\Phi }|+_{k=1}^{|A|}|\chi _k\chi _k|)`$ with $`|\chi _k_{k^{}k}=1/\sqrt{2}(|00+e^{i\phi _{kl}(t)}|11)`$, where $`\phi _{kl}(t)`$ is the effective interaction phase between particles $`kA`$ and $`lB`$. The state $`E_t`$ describing the total decoherence process incorporates the influence of all particles $`lB`$. We obtain $`E_t`$ (up to normalization) by calculating the Hadamard product of all $`E_t^{(l)}`$ written in the standard basis, i.e. by component-wise multiplication. The matrix elements of $`E_t`$ expressed in the tensor Bell basis finally determine $`_t`$. We find that $`E_t`$ has non-zero components only in the subspace spanned by $`\{|\varphi _{k_1\mathrm{}k_{N_A}}\}`$ with $`k_j\{0,3\}`$. The map $``$ thus contains only tensor products of Pauli operators $`1𝐥`$ and $`\sigma _z`$. Equivalently, we can express the action of the map on an arbitrary probe input state $`\rho =_{s,s^{}}\rho _{\mathrm{𝐬𝐬}^{}}(0)|𝐬𝐬^{}|`$ by determining the evolution of the coherences $`\rho _{\mathrm{𝐬𝐬}^{}}(t)`$. With $`𝐬_A,𝐬_A^{}`$ we denote binary vectors of length $`N_A`$. We can express the coherences as $`\rho _{\mathrm{𝐬𝐬}^{}}(t)=C_{\mathrm{𝐬𝐬}^{}}(t)\rho _{\mathrm{𝐬𝐬}^{}}(0)`$ with Ca05
$`C_{\mathrm{𝐬𝐬}^{}}(t)=e^{i\frac{1}{2}_k(𝐬_A𝐬_A^{})𝚪_𝐤}{\displaystyle \underset{k=1}{\overset{N_B}{}}}\mathrm{cos}[\frac{1}{2}(𝐬_A𝐬_A^{})𝚪_𝐤].`$ (2)
The $`(𝚪_𝐤)_j=\mathrm{\Gamma }_{kj}`$ for each particle $`kB`$ are $`N_A`$-dimensional vectors. The method described above can be considered as an extension of the generalized VBS picture for states to one for completely positive maps. In this picture, we can determine the evolution of arbitrary system states in a mesoscopic spin environment, since the computational effort to calculate the maps scales only linearly with the number of particles in the environment (as opposed to exponentially for general system-environment interactions). The size of the probe system is limited to about 10 spins for numerical computations due to exponential scaling with system size $`N_A`$.
The quantum properties of the system are directly linked to the classical statistical properties of the gas through $`\mathrm{\Gamma }(t)`$. In general, it is thus necessary to know the classical $`n`$-body phase-space distributions to give a complete description of the quantum state. If we assume no control over quantum or classical degrees of freedom of the background gas we should average over all possible collision patterns at any given time Ca05 : $`\overline{C}_{s,s^{}}(t)=d\mathrm{\Gamma }p_t(\mathrm{\Gamma })C_{s,s^{}}(\mathrm{\Gamma })`$, where $`p_t(\mathrm{\Gamma })`$ is the probability that at time $`t`$ the interaction history is given by $`\mathrm{\Gamma }`$.
For some gas models and regimes (like the Boltzmann gas studied in Ca05 ), correlations play a minor role and one can find analytical expressions for single-particle phase-space distributions. In this paper, however, we study a lattice gas model that exhibits strong correlations, and produce the different random realizations of $`\mathrm{\Gamma }(t)`$ by direct simulation of the gas.
The lattice model can be possibly implemented in a quantum optical system. It has already been demonstrated that an optical lattice can be used to store ultra-cold atomic gases. The degree of control in these experiments is extraordinary, opening the door to a wide range of experiments and theoretical proposals greiner02 ; mandel03 ; paredes04 ; jaksch98 ; hofstetter02 ; duan03 ; damski03 . One can choose a parameter regime where each lattice site is occupied by at most one atom jaksch98 ; greiner02 . The internal state of the atom (e.g. two meta-stable hyperfine states) can be stored in coherent superpositions over long time-scales (few minutes). Coherent inter-atomic interactions have been achieved by cold collisions mandel03 . These correspond precisely to the Ising-type interactions chosen here. One can also find schemes expref to induce a random (incoherent) hopping of atoms from one site to its neighboring sites. Hence, we consider an $`M\times M`$ lattice containing $`N`$ particles that randomly hop from site to site with a hopping rate $`\eta `$, and have nearest-neighbor interactions with coupling constant $`g_o`$. With special relevance to possible experiments, we note that all results in this paper hold even when the environment particles themselves decohere. The only requirement is that the diagonal elements of the environment’s state —in the canonical basis— remain unchanged.
Decoherence of a single qubit: For a system consisting of a single qubit $`A=\{1\}`$ and an arbitrary environment, the time dependent map corresponding to a particular collisional history is
$`_t\rho =\lambda _{00}\rho +\lambda _{11}\sigma _z\rho \sigma _z+\lambda _{01}(1𝐥\rho \sigma _z\sigma _z\rho 1𝐥),`$ (3)
with $`\lambda _{00}=(1+r\mathrm{cos}\gamma )/2`$, $`\lambda _{11}=(1r\mathrm{cos}\gamma )/2`$, and $`\lambda _{01}=(ir\mathrm{sin}\gamma )/2`$ where $`r(t)=_{lB}\mathrm{cos}(\frac{\phi _{1l}(t)}{2})`$, $`\gamma (t)=_{lB}\frac{\phi _{1l}}{2}`$. Depending on the parameter regime, semi-quantal gases can follow various collision patterns. Accordingly, the dynamics of their quantum properties can differ considerably. If in every time step $`\delta t`$ a given particle collides with a different particle and acquires an interaction phase $`\delta _\phi `$, the dynamics will be purely Markovian. The coherence of that particle will decay exponentially fast with the number of time steps $`k=\mathrm{\Delta }t/\delta t`$, $`|\rho _{01}|=[\mathrm{cos}(\delta _\phi /2)]^k=e^{\mathrm{\Delta }t/\tau _e}`$ with $`\tau _e8\delta t/\delta _\phi ^2`$. If, on the other hand, in a small time interval $`\mathrm{\Delta }t`$ a given gas particle has collided $`k`$ times with the same particle, the coherent addition of the interaction phase leads to a Gaussian type of decay: $`|\rho _{01}|=\mathrm{cos}(k\delta _\phi /2)e^{\mathrm{\Delta }t^2/(2\tau _g^2)}`$ with $`\tau _g=2\delta t/\delta _\phi `$. The exponential and the Gaussian decay are the two extreme cases, the dynamics of the coherence $`|\rho _{01}|`$ will usually lie in between. This also holds for the coherences in a multi-qubit density matrix. A complete characterization of the system’s decoherence is obtained by averaging the maps $`_t`$ over the possible collision patterns. The resulting map has the same form as (3), but replacing the coefficients $`\lambda _{ij}`$ by their average values. An analytical expression for the time dependence of these averaged coefficients, and hence of the decoherence process, is in general hard to obtain. However, for finite lattices, one can give a precise description of the time dependence.
Decoherence of bipartite entangled states: We now study the decoherence of different, initially entangled two-qubit states. In order to explore different regimes, we imagine a scenario where the probe particles can be displaced at a constant speed $`v`$ relative to the gas. By varying the speed $`v`$ and the distance $`d`$ between the probe particles, we can highlight two effects: (i) By decreasing the probe speed, we analyze the effect of multiple interactions with the same particle in contrast to interactions with different (independent) particles. (ii) By increasing the distance, we turn from correlated to independent collisions between probe and environment particles.
Figure 1(a) shows the decay of entanglement, measured by the concurrence wootters98 , of an initial Bell state $`|\varphi ^+`$ in two extreme scenarios: (i) Fixed probe particles ($`v=0`$). (ii) Large probe speeds, $`v/a\eta `$, where $`a`$ is the inter-site spacing. A fixed value $`\phi =0.1`$ is assigned to the collisional phase every time a probe particle crosses an occupied site. These two scenarios illustrate the difference between Markovian and non-Markovian environments. A large probe speed enforces a perfect Markovian behavior which matches the analytical curve <sup>1</sup><sup>1</sup>1At every time step, each probe particle interacts with a new environment particle with probability $`\nu `$. Hence, after a number $`s`$ of time steps, the relevant coherence is given by $`|C_{00,11}|=|C_{0,1}|^2=|\nu \mathrm{exp}(i\delta _\phi /2)\mathrm{cos}(\delta _\phi /2)+(1\nu )|^{2s}`$.. Figure 1(b) shows the concurrence at a given time $`t_1`$ as a function of the distance for three different entangled states: two Bell states and a cluster state (see figure caption). For Bell states the concurrence is equal to the absolute value of their only non-zero off-diagonal element in the density matrix, and therefore Figure 1 provides direct information about the individual coherences. The figure shows the influence of correlated collisions: coherences $`\rho _{01,10}`$ (in $`|\psi ^+`$) are robust against correlated noise, and coherences $`\rho _{00,11}`$ (in $`|\varphi ^+`$) are especially fragile under correlated noise. The remaining coherences decay in the same way under correlated or uncorrelated noise (hence, the weak distance dependence of $`|G`$). From Fig. 1(b) we also see that the immediate environments of each probe become more independent as $`d`$ increases.
The different behavior under correlated and uncorrelated collisions can be readily understood. For two probe particles ($`1`$ and $`2`$) with very similar collision patterns, i.e., $`\mathrm{\Gamma }_{2j}=\mathrm{\Gamma }_{1j}+\delta _j`$ for all $`j`$, coherences associated with $`|0110|`$ will only decay by a factor $`2^{N_B}_{s_B}e^{i\delta 𝐬_B}`$, while $`|1100|`$ will be “super-damped” by $`2^{N_B}_{s_B}e^{i(2𝚪_\mathrm{𝟏}+\delta )𝐬_B}`$. Therefore, classical correlations in the collisions, e.g. induced by the geometry of the set-up, can significantly influence the entanglement properties of the system.
Decoherence of multipartite entangled states: We now apply our method to investigate the decoherence of different multipartite entangled probe states of $`N_A`$ qubits. Due to the lack of simple, computable multipartite entanglement measures, we use the negativity of bipartitions Vi02 as an indicator of multipartite entanglement in the system. That means, we consider bipartitions of the system, i.e. a partition consisting of a set of particles $`A_k`$ and its complement $`\overline{A}_k`$, and investigate the entanglement properties with respect to the $`2^{N_A1}1`$ independent bipartitions Du00 . In general, we get a broad picture of multipartite entanglement in this way. For each bipartition we can determine its negativity $`𝒩_{A_k}=(\rho ^{T_{A_k}}_11)/2`$ Vi02 . We define two multipartite entanglement measures: (i) the average negativity $`\overline{𝒩}`$, as the average over all bipartitions, $`\overline{𝒩}=1/(2^{(N_A1)}1)_{A_k}𝒩_{A_k}`$, and (ii) $`𝒩_{min}=\mathrm{min}\{𝒩_{A_k}\}`$. Zero average negativity is a necessary condition for full separability of the state, and $`𝒩_{min}=0`$ is a sufficient condition that the state is not multi-party distillably entangled.
We have examined different multipartite entangled probe states that interact with a lattice gas through a pairwise Hamiltonian $`H_{kl}=\sigma _z^{(k)}\sigma _z^{(l)}`$. These states are (linear) cluster states $`|\chi `$ Ra01 , GHZ-states and W-states $`|W_{i=1}^{N_A}|w_i`$, where $`|w_i`$ is the state corresponding to all spins in state zero except the $`i^{th}`$ that is in one. Since the noise process is basis dependent, we study variants of GHZ states corresponding to different local bases: $`|\mathrm{\Psi }|0^{N_A}+|1^{N_A}`$, $`|\mathrm{\Psi }^{}|0|+^{N_A1}+|1|^{N_A1}`$, $`|\mathrm{\Psi }^{\prime \prime }|+|0^{N_A1}+||1^{N_A1}`$. The decay of the average negativity is plotted in Fig. 2(a). To qualitatively understand the different behavior of the curves we first derive analytic results in the limit of independent environments for each particle $`kA`$. In this limit, the decoherence process can be described by a tensor product of single-qubit dephasing maps $`_t^{(k)}\rho =p_k\rho +(1p_k)\sigma _z^{(k)}\rho \sigma _z^{(k)}`$, where $`p_k=\frac{1}{2}(1+_{lB}\mathrm{cos}(2\phi _{kl}(t)))`$. Coherences decay as $`\rho _{01}^{(k)}(t)=(2p_k1)\rho _{01}(0)`$. As before, the precise time dependence of decoherence will be given by the average $`p(t)=p_k_{\mathrm{\Gamma }(t)}`$ over different realizations. We assume that this average value is the same for every probe particle $`k`$. Under the action of this dephasing map the family of GHZ states remain diagonal in the GHZ-type basis Du00 . For such states, we can obtain the spectrum of the partial transposed operators with respect to any bipartition analytically Du04 , and calculate average negativities for the three GHZ states given above. Here we give those with simple expressions: $`\overline{𝒩}=\frac{1}{2}|2p1|^{N_A}`$, and $`\overline{𝒩}^{\prime \prime }=\frac{1}{2^{N_A1}1}[𝒩_{min}^{\prime \prime }+(2^{N_A1}2)\frac{1}{2}|2p1|^{N_A1}]`$ with $`𝒩_{min}^{\prime \prime }=\frac{1}{2}\mathrm{max}\{0,|2p1|+|2p1|^{N_A1}+|2p1|^{N_A}1\}`$. For $`W`$-states, direct calculation leads to $`\overline{𝒩}^W=\frac{|2p1|^2}{N_A(2_A^N2)}_{a=1}^{N_A1}\left(\genfrac{}{}{0pt}{}{N_A}{a}\right)\sqrt{a(N_Aa)}`$.
Several observations follow from these analytic results for independent environments: (i) Standard GHZ, $`|\mathrm{\Psi }`$, and $`W`$ states remain $`N_A`$-party distillable for all times, since $`𝒩_{min}`$ only reaches zero asymptotically as $`t\mathrm{}`$. This does not hold for the $`|\mathrm{\Psi }^{\prime \prime }`$ GHZ, which has one partition that becomes disentangled at finite $`t`$, nor for $`|\mathrm{\Psi }^{}`$, for which $`\overline{𝒩}`$ vanishes at a finite time (all partitions are disentangled). In Fig. 2 we see that $`\overline{𝒩}`$ also vanishes at finite time for the cluster state. (ii) In the limit of large system sizes, the average negativity is, to a very good approximation, given by the negativity of the half–half partition (the distribution of partitions with $`k`$ particles is sharply peaked at $`k=N_A/2`$). (iii) GHZ-type states decay exponentially as we increase the system size $`N_A`$, while for $`W`$ the coherences do not vary with the system size, rendering a weak dependence of its average negativity (constant to first order). (iv) For the states studied here, if a partition is initially more entangled than another, it will remain so also at later times. This does not hold in the presence of correlated collisions, which occur when the distance between probe particles is not large with respect to the relevant times $`t`$, i.e. $`d<\sqrt{\eta t}`$.
In finite lattices and after long enough times the above description fails due to inevitable correlated collisions. However, the finite size leads to interesting effects that can also be easily understood. Figure 2(b) shows one such effect: the coherence (and hence the probe’s entanglement) lost due to the stochastic interaction with environment spins, is partially recovered after a characteristic revival time. After long times a given environment particle will have collided $`n_k`$ times with a probe particle $`k`$. After $`sM^2`$ steps, we describe the distribution of values $`n_k`$ by a Gaussian of mean value $`n_k=n=4s/M^2`$ and the variance $`\sigma ^2=(n_kn)^2=\alpha n`$, where $`\alpha `$ depends on the particular lattice model. We expect that at long times most blocking effects between environment particles will be washed out, and therefore assume that different environment particles will have independent collision distributions. Hence, the total effect of $`N_B`$ on a particular coherence will scale as $`C^{N_B}`$, where $`C`$ is the decay factor of the coherence due to a single environment particle. The value of $`C`$ is given by an average taken over a Gaussian of mean $`\phi _o`$ and width $`\sigma _\phi `$: $`C=\mathrm{cos}(2\phi )=\mathrm{cos}(2\phi _o)\mathrm{exp}(8\sigma _\phi ^2)`$. The phase $`\phi `$ is a sum of the collisional phases (with the corresponding signs) involved in the particular coherence. For example, for a standard GHZ all phases are added $`\phi =_{k=1}^{N_A}\phi _{1k}`$ leading to a mean value $`\phi _oN_An\delta _\phi `$ and a variance $`\sigma _\phi ^2\alpha ^{}N_An\delta _\phi ^2`$ <sup>2</sup><sup>2</sup>2The factor $`\alpha ^{}`$ (and $`\alpha ^{\prime \prime }`$ below) also includes the contribution of $`K_{kk^{}}=n_k^{}n_kn^2n`$, which depends on the distance between the probe particles $`k,k^{}`$.. For the $`W`$-state and a given coherence, say $`\rho _{w_1w_2}`$, the phase is $`\phi =\phi _{1k}\phi _{2k}`$, leading to a vanishing mean value and to a purely exponential decay with $`\sigma _\phi ^2\alpha ^{\prime \prime }2n\delta _\phi ^2`$ —which is independent of the system size $`N_A`$. The exponential decay of the $`W`$-state and the periodic revival of the GHZ state can be clearly identified in Fig. 2(b).
Summary: We have studied a microscopic, exact model for non-Markovian decoherence where the joint system is described by a spin gas. Using a generalized VBS picture for maps, we have determined the time-dependent maps for the decoherence process. We studied the decay of entanglement for different multipartite entangled probe states in a lattice gas with a mesoscopic number of particles. Depending on the parameters of the gas, we have shown how to reach qualitatively different regimes such as Markovian- and non-Markovian, and correlated and non-correlated decoherence processes. For finite lattices we find that, although the interactions with the environment are stochastic, entanglement in the probes can spontaneously revive at a time given by the size of the lattice and independent of the number of environment particles.
We thank J. Asbóth for reading the manuscript. This work was supported by the FWF, the European Union (IST-2001-38877,-39227, OLAQUI, SCALA), the DFG, and the ÖAW through project APART (W.D.).
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# Path-integral approach in a chiral quark-diquark model to the nucleon structure and interactions
## 1 Introduction
The origin of the short range repulsion in the nucleon-nucleon($`NN`$) interaction is a long-standing problem. The $`NN`$ interaction below the meson production threshold is phenomenologically well understood, having been fit to the phase shift analysis, but its microscopic understanding is still needed. While the long range part is well described by the meson exchange picture, there are several different approaches to the description of the short range part, including meson exchange and quark exchange . Because in the short range region the two nucleons have a substantial spatial overlap, it is commonly assumed that the internal structure of the nucleon gives a sizable contribution to the nuclear force at such short distances. To obtain a comprehensive understanding of the nuclear force from the short to long range distances, models which incorporate both the meson exchange and the internal struture of the nucleon are needed.
In a recent publication, we studied the structure of the nuclear force using the path-integral hadronization approach to a chiral quark-diquark model . This method incorporates two aspects of chiral symmetry, which naturally describes the pion exchange interaction, and of the internal structure of hadrons. This method employs an extended model of the Nambu-Jona-Lasinio type with the interactions that are not only of the quark-antiquark type but also of the quark-diquark type . The quarks and diquarks were integrated out to generate an effective Lagrangian for mesons and baryons while maintaining important symmetries, such as the gauge and chiral symmetries. The hadron structure was then described in terms of its constituents: a quark and an antiquark for mesons and a quark and a diquark for baryons. It was also pointed out that the resulting effective Lagrangian contains various interactions among hadrons, such as meson-meson, meson-baryon and baryon-baryon interactions.
In our model, all components of the $`NN`$ force are contained in an effective Lagrangian that is written in a concise form as a trace-log. Then, the expansion of the trace-log terms produces an NN force that is described as meson exchanges at long and medium ranges and quark-diquark exchanges at short ranges. The latter was then shown to contain various types of non-local interactions, including a scalar iso-scalar type and a vector iso-scalar type. We evaluated the properties of the scalar and vector type interactions; their ranges, effective masses and strengths.
In this paper, we report the results of the study of the nuclear force in the simple framework where the axial-vector diquark is neglected. In the previous work, we suggested that the neglecting the axial-vector diquark does not affect to the size of the nucleon, hence we can properly evaluate the ranges of the $`NN`$ interactions with only scalar diquarks.
We organize this paper as follows. In $`\mathrm{\S }`$ 2, we briefly give the derivation of the trace-log formula in the path-integral hadronization of the NJL model with quark-diquark correlations. In $`\mathrm{\S }`$ 3, terms containing the $`NN`$ interaction are investigated in detail, where the general structure of the $`NN`$ amplitude is presented. We present a sample numerical calculation for the case in which there is only a scalar diquark. The present study of the $`NN`$ interaction is not quantitatively complete, but it will be useful in demonstrating some important aspects of the nuclear force, in particular that the range of the short range interaction is related to the intrinsic size of the nucleon. The final section is devoted to a summary.
## 2 Effective Lagrangian for mesons and nucleons
We now briefly review the method to derive an effective Lagrangian for mesons and nucleons from a quark and diquark model of chiral symmetry following the previous work of Abu-Raddad et al . We start from the NJL Lagrangian,
$$_{\text{NJL}}=\overline{q}(i\text{/}m_0)q+\frac{G}{2}\left[(\overline{q}q)^2+(\overline{q}i\gamma _5\stackrel{}{\tau }q)^2\right].$$
(1)
Here, $`q`$ is the current quark field, $`\stackrel{}{\tau }`$ represents the isospin (flavor) Pauli matrices, $`G`$ is a dimensional coupling constant, and $`m_0`$ is the current quark mass. In this paper, we set $`m_0=0`$ i.e. we work in the chiral limit. As usual, the NJL Lagrangian is bosonized by introducing meson fields as collective auxiliary fields in the path-integral method . At an intermediate step, we find the following Lagrangian:
$`_{q\sigma \pi }^{}=\overline{q}\left(i\text{/}(\sigma +i\gamma _5\stackrel{}{\tau }\stackrel{}{\pi })\right)q{\displaystyle \frac{1}{2G}}(\sigma ^2+\stackrel{}{\pi }^{\mathrm{\hspace{0.33em}2}}).`$ (2)
Here $`\sigma `$ and $`\stackrel{}{\pi }`$ are scalar-isoscalar sigma and pseudoscalar-isovector pion fields, as generated from $`\sigma \overline{q}q`$ and $`\stackrel{}{\pi }i\overline{q}\stackrel{}{\tau }\gamma _5q`$, respectively. For our purpose, it is convenient to work in a non-linear basis rather than the linear one . This is realized through the chiral rotation from the current ($`q`$) to constituent ($`\chi `$) quark fields:
$`\chi =\xi _5q,\xi _5=\left({\displaystyle \frac{\sigma +i\gamma _5\stackrel{}{\tau }\stackrel{}{\pi }}{f}}\right)^{1/2},`$ (3)
where $`f^2=\sigma ^2+\stackrel{}{\pi }^{\mathrm{\hspace{0.33em}2}}`$. Thus, we find
$`_{\chi \sigma \pi }^{}=\overline{\chi }\left(i\text{/}f\text{/}v\text{/}a\gamma _5\right)\chi {\displaystyle \frac{1}{2G}}f^2,`$ (4)
where
$`v_\mu ={\displaystyle \frac{i}{2}}\left(_\mu \xi ^{}\xi +_\mu \xi \xi ^{}\right),a_\mu ={\displaystyle \frac{i}{2}}\left(_\mu \xi ^{}\xi _\mu \xi \xi ^{}\right)`$ (5)
are the vector and axial-vector currents written in terms of the chiral field,
$`\xi =\left({\displaystyle \frac{\sigma +i\stackrel{}{\tau }\stackrel{}{\pi }}{f}}\right)^{1/2}.`$ (6)
The Lagrangian (4) describes not only the kinetic term of the quark, but also quark-meson interactions such as the Yukawa one, Weinberg-Tomozawa one etc.
In the model we consider here, we introduce diquarks and their interaction terms with quarks. We assume local interactions between quark-diquark pairs to generate the nucleon field. As suggested by a method of constructing a local nucleon field, it is sufficient to consider two types of diquarks, a scalar, isoscalar diquark, $`D`$, and an axial-vector, isovector diquark, $`\stackrel{}{D}_\mu `$ . In this work, we consider only the scalar diquark. Hence, our microscopic Lagrangian for quarks, diquarks and mesons is given by
$`=\overline{\chi }(i\text{/}f\text{/}v\text{/}a\gamma _5)\chi {\displaystyle \frac{1}{2G}}f^2+D^{}(^2+M_S^2)D+\stackrel{~}{G}\overline{\chi }D^{}D\chi .`$ (7)
In the last term, $`\stackrel{~}{G}`$ is a coupling constant for the quark-diquark interaction.
Now, the hadronization procedure can be carried out straightforwardly by introducing the baryon fields as auxiliary fields, $`BD\chi `$, and by eliminating the quark and diquark fields in (7). The final result is written in a compact form as
$`_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{1}{2G}}f^2i\mathrm{tr}\mathrm{ln}(i\text{/}f\text{/}v\text{/}a\gamma _5){\displaystyle \frac{1}{\stackrel{~}{G}}}\overline{B}B+i\mathrm{tr}\mathrm{ln}(1𝒮).`$ (8)
Here, the trace is taken over space-time, color, flavor and Lorentz indices, and the operator $`𝒮`$ is defined by
$`𝒮=\mathrm{\Delta }^T\overline{B}SB,`$ (9)
In this equation, $`S=(i\text{/}f\text{/}v\text{/}a\gamma _5)^1,\mathrm{\Delta }=(^2M_S^2)^1`$ are the propagators of the quark and scalar diquark, respectively, and transposed diquark propagators, as denoted by the superscript $`T`$, are employed. Note that the quark propagator $`S`$ contains the interactions with pions. Through this interactions, the nucleon-pion interactions are obtained. Though the effective meson-nucleon Lagrangian (8) looks simple, it contains many important physical ingredients when the trace-log terms are expanded:
* It generates a meson Lagrangian in a chirally symmetric manner. Up to fourth order in the meson fields, it produces precisely the Lagrangian of the linear sigma model with the realization of the spontaneous breaking of chiral symmetry. Hence, the vacuum expectation value $`f`$ turns out to be the pion decay constant $`f_\pi `$.
* From the second trace-log term, a nucleon effective Lagrangian is derived. In a previous paper, the kinetic term of the nucleon was investigated, and the mass of the nucleon was computed at the one-loop level .
* In the nucleon effective Lagrangian, meson-nucleon couplings appear through the diagrams, as shown in Fig. 1. Their strengths and form factors can be computed with the use of the underlying quark-diquark dynamics. Using these vertices, meson-exchange interactions are constructed. the photon-nucleon couplings are also obtained in the same manner. The radii and magnetic moments were evaluated without the axial-vector diquark.
* There are diagrams that contain many nucleon fields. For instance, $`NN`$ interactions are expressed as one-loop diagrams, as shown in Fig. 2. This term describes the short range part of the $`NN`$ interaction. In this paper, we focus our attention mostly on the $`NN`$ interaction derived from the one-loop diagrams.
## 3 Short range interaction
We evaluate the short range interaction described by the quark-diquark loop, as shown in Fig. 2. Using the interaction vertices given in Eq. (9), it is straightforward to compute the amplitude for the quark-diquark loop:
$`_{NN}=iN_cZ^2{\displaystyle \frac{d^4k}{(2\pi )^4}}`$ (18)
$`\times `$ $`{\displaystyle \frac{\overline{B}(p_1^{})(\begin{array}{c}\text{/}\\ k\hfill \end{array}+m_q)B(p_1)\overline{B}(p_2^{})(\begin{array}{c}\text{/}\\ p\hfill \end{array}_2\begin{array}{c}\text{/}\\ p\hfill \end{array}_1^{}+\begin{array}{c}\text{/}\\ k\hfill \end{array}+m_q)B(p_2)}{[(p_1k)^2M_S^2][k^2m_q^2][(p_1^{}k)^2M_S^2][(p_2p_1^{}+k)^2m_q^2]}}.`$
where $`N_c`$ is the number of colors and $`Z`$ is the wave-function renormalization constant. In this equation, the momentum variables are assigned such that $`p_1`$ and $`p_1^{}`$ ($`p_2`$ and $`p_2^{}`$) are for the pair of contracted baryon fields, $`\overline{B}(p_1^{})\mathrm{}B(p_1)`$ ($`\overline{B}(p_2^{})\mathrm{}B(p_2)`$), and the momentum transfer is defined by $`q=p_1^{}p_1`$. Note that the momentum $`q`$ is carried by the diquark pair. The amplitude defined in this way can be interpreted as a direct term in the local potential approximation. When computing physical quantities such as phase shifts, we need to include the exchange term that is obtained by interchanging the momentum variables $`p_1^{}p_2^{}`$. Although the one-loop integral (18) converges when the scalar diquark is included, we keep the counter terms of the Pauli-Villars regularization. Because our model is a cut-off theory with a relatively small cutoff mass, $`\mathrm{\Lambda }_{PV}=0.63`$ GeV, the counter-terms play a significant role. However, because we include only the scalar diquark, we do not attempt to make a comparison at the quantitative level. Rather, in the following, we study some basic properties of the amplitude itself, mostly the interaction ranges extracted from Eq. (18).
Let us evaluate the integral in the center-of-mass system for elastic scattering:
$`p_1`$ $`=`$ $`(E_\stackrel{}{p},\stackrel{}{p}),p_2=(E_\stackrel{}{p},\stackrel{}{p}),`$
$`p_1^{}`$ $`=`$ $`(E_\stackrel{}{p}^{},\stackrel{}{p}^{}),p_2^{}=(E_\stackrel{}{p}^{},\stackrel{}{p}^{}),|\stackrel{}{p}|=|\stackrel{}{p}^{}|.`$
To proceed, we write the amplitude (18) as
$`_{NN}`$ $`=`$ $`F_S(\stackrel{}{P},\stackrel{}{q})(\overline{B}B)^2+F_V(\stackrel{}{P},\stackrel{}{q})(\overline{B}\gamma _\mu B)^2+\mathrm{},`$ (19)
where
$`\stackrel{}{P}=\stackrel{}{p}^{}+\stackrel{}{p},\stackrel{}{q}=\stackrel{}{p}^{}\stackrel{}{p}.`$ (20)
Roughly speaking, the $`q`$-dependence expresses the interaction range in the $`t`$-channel, whose Fourier transform is interpreted as the $`r`$-dependence of the local potential, while the $`P`$-dependence expresses the non-locality of the interaction.
In Eq. (19), we have defined the coefficients $`F_S`$ and $`F_V`$ as scalar and vector interactions, respectively. The ellipsis in Eq. (19) then include terms involving external momenta $`p_i`$, such as $`(\overline{B}\mathrm{\Gamma }(p_i)B)(\overline{B}B)`$ and $`(\overline{B}B)(\overline{B}\mathrm{\Gamma }(p_i)B)`$, where $`\mathrm{\Gamma }(p_i)`$ are $`4\times 4`$ matrices involving $`p_i`$. These momentum dependent terms are, however, expected to play a less important role than the dominant components of the scalar and vector terms, as in boson exchange models. Therefore, we consider here the corresponding terms of scalar and vector types. It turns out that the interaction coefficients defined in this way are attractive for $`F_S`$ and repulsive for $`F_V`$.
As anticipated, the amplitude is highly non-local, as the quark-diquark loop diagram implies. As a matter of fact, the interaction range for the diquark exchange in the $`t`$-channel is shorter than that for the quark exchange in the $`s`$-channel. This is shown explicitly in Fig. 3. Nevertheless, we proceed further and carry out an expansion in $`\stackrel{}{P}`$. We obtain
$`F_i(\stackrel{}{P},\stackrel{}{q})=F_i(\stackrel{}{P},\stackrel{}{q})|_{\stackrel{}{P}=0}+\stackrel{}{P}{\displaystyle \frac{}{\stackrel{}{P}}}F_i(\stackrel{}{P},\stackrel{}{q})|_{\stackrel{}{P}=0}+\mathrm{},`$ (21)
where $`i`$ stands for $`S`$ or $`V`$. The resulting $`\stackrel{}{q}`$ dependent function, in particular the first term, can be interpreted as the Fourier transform of a local potential as a function of the relative coordinates $`\stackrel{}{r}=\stackrel{}{x}_1\stackrel{}{x}_2`$. Then, the general structure of the $`NN`$ potential can be studied by performing the non-relativistic reduction of the amplitude, Eqs. (19) and (21). It contains central, spin-orbit and tensor components that accompany functions of non-locality $`\stackrel{}{P}`$.
We now discuss the functions $`F_S`$ and $`F_V`$ in Eq. (19) for the leading-order terms of Eq. (21). Because we cannot write the resulting $`q`$ dependent functions in an analytic form, we have carried out numerical calculations employing several different values of $`\stackrel{~}{G}`$ parameter for different binding energies and the size of the quark-diquark bound state.
At this point, we explain our parameters and the regularization scheme. We follow the scheme presented in Ref. , except for the treatment of $`\stackrel{~}{G}`$. The $`\mathrm{\Lambda }_{PV}`$ is the cutoff mass of the Pauli-Villars method used in this work to regularize divergent integrals. We note that our model is unrenormalizable. The NJL coupling constant $`G`$ and the cut-off mass $`\mathrm{\Lambda }_{PV}`$ are fixed to generate the constituent quark mass $`m_q`$ and the pion decay constant through the NJL gap equation in the meson sector . The mass of the scalar diquark $`M_S`$ is determined in the NJL model by solving the Bethe-Salpeter equation in the diquark channel . Then, we have one free parameter, the quark-diquark coupling constant $`\stackrel{~}{G}`$. The strength of $`\stackrel{~}{G}`$ controls the binding or size of the nucleon and generates the mass of the nucleon $`M_N`$. Then we perform numerical calculations by varying the coupling constant $`\stackrel{~}{G}`$. In Table 1, we list typical parameter values as used in Ref. .
First, we determine the strengths of the interactions by extracting coupling constant squares $`g_i^2`$ as
$`|F_i(0,0)|={\displaystyle \frac{g_i^2}{m_i^2}},g_i>0,`$ (22)
where the effective meson masses $`m_i`$ are evaluated using the inverse of (23) and are plotted in Fig. 6. The results are displayed in Fig. 4 as functions of the size of the nucleon, $`r^2^{1/2}`$. The coupling strengths can be compared with the empirical values $`g_S10`$ and $`g_V13`$ . The present results are strongly dependent on $`r^2`$. When only a scalar diquark is included, the scalar interaction becomes much stronger than the vector interaction. Phenomenologically, the vector (omega meson) coupling is stronger than the scalar (sigma meson) coupling.
We should comment on the effect of the axial-vector diquark. The loop integral containing the axial-vector diquarks diverges and has a large numerical value even when it is regularized. Therefore, it significantly affects the absolute values of the loop integrals as well as that of the baryon self energy, which is necessary to extract the normalization factor $`Z`$. Therefore, it is expected that the strengths are affected significantly, but the ranges and the corresponding masses, which are computed from the momentum dependence of the normalized loop integrals, are not affected much.
Let us discuss the interaction ranges. In Fig. 5, we show the $`q`$-dependence of the zeroth order coefficients in Eq. (21) for three different sizes of the nucleon. It is obvious from Fig. 5 that as the size of the nucleon becomes smaller the interaction ranges become shorter. More quantitatively, we define the interaction ranges $`R_i`$ by
$`R_i^26{\displaystyle \frac{1}{F_i(q^2)}}{\displaystyle \frac{F_i}{q^2}}|_{q^20},`$ (23)
which are related to the mass parameters of the interaction ranges as $`m_i\sqrt{6}/R_i`$. It is interesting that the range (and hence mass) parameters of the interactions are approximately proportional to the size of the nucleon $`r^2^{1/2}`$, as shown in Fig. 6. When the parameter set I listed in Table 1 is used, the mass parameters are about 650 MeV and 800 MeV for the scalar and vector interactions, respectively, which are very close to the masses of the sigma and omega mesons. If, however, we use the parameter set II for a nucleon size of about 0.5 fm, then the two masses become $`m_S800`$ MeV and $`m_V1000`$ MeV.
In the present analysis, apart from the absolute values of the interaction strengths, the interaction ranges for the scalar and vector interactions have been produced appropriately with the quark-diquark loop diagram. One question concerns the small but non-negligible difference between the ranges in the scalar and vector channels, which is consistent with empirical results. This can be roughly understood from the dimensionality of the loop integral. As seen from Eq. (18), the integrand for the vector interaction is of higher order with respect to the loop momentum than the scalar interaction. Because of this, the vector part reflects shorter distant dynamics and produces a shorter interaction range.
Although we have obtained reasonable results for the interaction ranges, the short range part of the nuclear force must be repulsive. Hence, current approach is not satisfactory. We believe that the inclusion of the axial-vector diquark will give a sizable contribution to the short range repulsion because the quark-axial-vector diquark-baryon vertex is basically spin-spin type interactions, and thus correct this shortcoming of our present approximation.
## 4 Summary
In this paper we have studied the NN interaction using a microscopic theory of quarks and diquarks. Nucleons were described as quark-diquark bound states. The quark and diquark degrees of freedom were integrated out using the path-integral method, and an effective Lagrangian was derived for mesons and baryons. The resulting trace-log formula contains various meson and nucleon interaction terms, including the $`NN`$ interaction in the short range region expressed as a quark-diquark loop. Hence, the $`NN`$ interaction can be naturally expressed as meson exchanges at long ranges and quark or diquark exchanges at short ranges.
For short range interactions, we have computed a quark-diquark loop corresponding to the direct term of the two-nucleon interaction. The resulting interaction is highly nonlocal. We have extracted scalar and vector type interactions in the local potential approximation. It turns out that the scalar term is attractive, while the vector is term repulsive.
In the present paper our numerical calculations contained only scalar diquarks, as a first step toward a full calculation. Therefore, we have concentrated our study mostly on the interaction ranges or, equivalently, the masses, because the magnitude of strengths is affected by the axial-vector diquark. On the other hand, the interaction ranges were better studied in the present study including only scalar diquarks, reflecting the size of the nucleon. Consequently, the mass parameters of the interaction ranges for the scalar and vector interactions were found to be about 650 MeV for the scalar type and about 800 MeV for the vector type, once again with the nucleon size set at 0.77 fm. Hence in our model, the scalar-isoscalar interaction emerges from the quark-diquark loop at an energy scale similar to that of the sigma and omega meson exchanges. The existence of the two components in the nuclear force, boson exchanges and quark-diquark exchanges, is a general feature when we consider a model of nucleons that is composed of a quark (and diquark) core surrounded by meson cloud.
The present result encourages us to further study baryon properties by extending the model to include the axial-vector diquark. We have already started such a study. When the axial-vector diquark is included, loop integrals diverge more strongly than in the case of the scalar diquark, due to the massive vector nature of the propagator. Although this causes the numerical study to be more complicated than in the present case, this work is in progress.
## Acknowledgments
We would like to thank Veljko Dmitrasinovic for his careful reading of this manuscript. This work was supported in part by the Sasakawa Scientific Research Grant from The Japan Science Society.
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# Non-asymptotic calibration and resolution
## 1 Introduction
We consider the problem of forecasting a new observation from the available data, which may include, e.g., all or some of the previous observations and the values of some explanatory variables. To make the process of forecasting more vivid, we imagine that the data and observations are chosen by a player called Reality and the forecasts are made by a player called Forecaster. To establish properties of forecasting algorithms, the traditional theory of machine learning makes some assumptions about the way Reality generates the observations; e.g., statistical learning theory assumes that the data and observations are generated independently from the same probability distribution. A more recent approach, prediction with expert advice (see, e.g., ), replaces the assumptions about Reality by a comparison class of prediction strategies; a typical result of this theory asserts that Forecaster can perform almost as well as the best strategies in the comparison class. This paper further explores a third possibility, suggested in , which requires neither assumptions about Reality nor a comparison class of Forecaster’s strategies. It is shown in that there exists a forecasting strategy which is automatically well calibrated; this result has been further developed in, e.g., . Almost all known calibration results, however, are asymptotic (see and for a critique of the standard asymptotic notion of calibration); a non-asymptotic result about calibration is given in , Proposition 2, but even this result involves unspecified constants and randomization. The main results of this paper (Theorems 1 and 2) establish simple explicit inequalities characterizing calibration and resolution of our deterministic forecasting algorithm.
Next we briefly describe the main features of our proof techniques and their connections with the literature. The proofs rely on the game-theoretic approach to probability suggested in . The forecasting protocol is complemented by another player, Skeptic, whose role is to gamble at the odds given by Forecaster’s probabilities. It can be said that our approach to forecasting is Skeptic-based, whereas the traditional approach is Reality-based and prediction with expert advice is Forecaster-based. The two most popular formalizations of gambling are subsequence selection rules (going back to von Mises’s collectives) and martingales (going back to Ville’s critique of von Mises’s collectives and described in detail in ). The pioneering paper on what we call the Skeptic-based approach, as well as the numerous papers developing it, used von Mises’s notion of gambling; appears to be the first paper in this direction to use Ville’s notion of gambling. Another ingredient of this paper’s approach, considering Skeptic’s continuous strategies and thus avoiding randomization by Forecaster (which was the standard feature of the previous work) goes back to and is also described in ; however, I learned it from Akimichi Takemura in June 2004 (whose observation was prompted by Glenn Shafer’s talk at the University of Tokyo).
It should be noted that, although our approach was inspired by and papers further developing , precise statements of our results and our proof techniques are completely different: they are more in the spirit of Levin’s result about the existence of neutral measures (see for details).
This version (version 4) of this technical report differs from the previous one in that it incorporates the changes made in response to the comments of the reviewers of its journal version (to be published in *Theoretical Computer Science*).
## 2 The algorithms of large numbers
In this section we describe our learning protocol and the general forecasting algorithm studied in this paper. The protocol is:
FOR $`n=1,2,\mathrm{}`$:
Reality I announces $`x_n𝐗`$.
Forecaster announces $`p_n[0,1]`$.
Reality II announces $`y_n\{0,1\}`$.
END FOR.
On each round, Reality chooses the datum $`x_n`$, then Forecaster gives his forecast $`p_n`$ for the next observation, and finally Reality discloses the actual observation $`y_n\{0,1\}`$. Reality chooses $`x_n`$ from a *data space* $`𝐗`$ and $`y_n`$ from the two-element set $`\{0,1\}`$; intuitively, Forecaster’s move $`p_n`$ is the probability he attaches to the event $`y_n=1`$. *Forecasting algorithm* is Forecaster’s strategy in this protocol. For convenience in stating the results of §6, we split Reality into two players, Reality I and Reality II.
Our learning protocol is a perfect-information protocol; in particular, Reality may take into account the forecast $`p_n`$ when deciding on her move $`y_n`$. (This feature is unusual for probability forecasting but it extends the domain of applicability of our results and we have it for free.)
Next we describe the general forecasting algorithm that we study in this paper (it was derived informally in ). A function $`𝐊:Z^2`$, where $`Z`$ is an arbitrary set and $``$ is the set of real numbers, is a *kernel on $`Z`$* if it is symmetric ($`𝐊(z,z^{})=𝐊(z^{},z)`$ for all $`z,z^{}Z`$) and positive definite ($`_{i=1}^m_{j=1}^m\lambda _i\lambda _j𝐊(z_i,z_j)0`$ for all $`(\lambda _1,\mathrm{},\lambda _m)^m`$ and all $`(z_1,\mathrm{},z_m)Z^m`$). The usual interpretation of a kernel $`𝐊(z,z^{})`$ is as a measure of similarity between $`z`$ and $`z^{}`$ (see, e.g., , §1.1). Our algorithm has one parameter, which is a kernel on the Cartesian product $`[0,1]\times 𝐗`$. The most straightforward way of constructing such kernels from kernels on $`[0,1]`$ and kernels on $`𝐗`$ is the operation of tensor product. (See, e.g., .) Let us say that a kernel $`𝐊`$ on $`[0,1]\times 𝐗`$ is *forecast-continuous* if the function $`𝐊((p,x),(p^{},x^{}))`$, where $`p,p^{}[0,1]`$ and $`x,x^{}𝐗`$, is continuous in $`(p,p^{})`$ for any fixed $`(x,x^{})𝐗^2`$.
K29 algorithm
Parameter: forecast-continuous kernel $`𝐊`$ on $`[0,1]\times 𝐗`$
FOR $`n=1,2,\mathrm{}`$:
Read $`x_n𝐗`$.
Set $`S_n(p):=_{i=1}^{n1}𝐊((p,x_n),(p_i,x_i))(y_ip_i)+\frac{1}{2}𝐊((p,x_n),(p,x_n))(12p)`$
for $`p[0,1]`$.
If $`signS_n(0)=signS_n(1)0`$, output $`p_n:=(1+signS_n(0))/2`$;
otherwise, output any root $`p`$ of $`S_n(p)=0`$ as $`p_n`$.
Read $`y_n\{0,1\}`$.
END FOR.
(Since the function $`S_n(p)`$ is continuous, the equation $`S_n(p)=0`$ indeed has a solution when $`signS_n(0)=signS_n(1)0`$ does not hold; remember that $`signS`$ is $`1`$ for $`S`$ positive, $`1`$ for $`S`$ negative, and $`0`$ for $`S=0`$.) The main term in the expression for $`S_n(p)`$ is $`_{i=1}^{n1}𝐊((p,x_n),(p_i,x_i))(y_ip_i)`$. Ignoring the other term for a moment, we can describe the intuition behind this algorithm by saying that $`p_n`$ is chosen so that $`p_i`$ are unbiased forecasts for $`y_i`$ on the rounds $`i=1,\mathrm{},n1`$ for which $`(p_i,x_i)`$ is similar to $`(p_n,x_n)`$. The term $`\frac{1}{2}𝐊((p,x_n),(p,x_n))(12p)`$, which can be rewritten as $`𝐊((p,x_n),(p,x_n))(0.5p)`$, adds an element of regularization, i.e., bias towards the “neutral” value $`p_n=0.5`$.
The K29 algorithm requires solving the equation $`S_n(p)=0`$, but this can be easily done using the bisection method or one of the numerous more sophisticated methods (see, e.g., , Chapter 9).
It is well known (see , Theorem II.3.1, for a simple proof) that there exists a function $`\mathrm{\Phi }:[0,1]\times 𝐗`$ (a *feature mapping* taking values in a Hilbert space<sup>1</sup><sup>1</sup>1Hilbert spaces in this paper are allowed to be non-separable or finite dimensional; we, however, always assume that their dimension is at least $`1`$. $``$ called the *feature space*) such that
$$𝐊(a,b)=\mathrm{\Phi }(a),\mathrm{\Phi }(b)_{},a,b[0,1]\times 𝐗$$
(1)
($`,_{}`$ standing for the inner product in $``$). It is known that, for any $`𝐊`$ and $`\mathrm{\Phi }`$ connected by (1), $`𝐊`$ is forecast-continuous if and only if $`\mathrm{\Phi }`$ is a continuous function of $`p`$ for each fixed $`x𝐗`$ (see Appendix B).
Now we can state the basic result about K29 (proved in Appendix A).
###### Theorem 1
Let $`𝐊`$ be the kernel defined by (1) for a feature mapping $`\mathrm{\Phi }:[0,1]\times 𝐗`$ continuous in its first argument. The K29 algorithm with parameter $`𝐊`$ ensures
$$\begin{array}{c}\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }(p_n,x_n)_{}^2\underset{n=1}{\overset{N}{}}p_n(1p_n)\mathrm{\Phi }(p_n,x_n)_{}^2,\hfill \\ \hfill N\{1,2,\mathrm{}\}.\end{array}$$
(2)
Let us assume, for simplicity, that
$$𝐜_𝐊:=\underset{p,x}{sup}\mathrm{\Phi }(p,x)_{}<\mathrm{}$$
(3)
(it is often a good idea to use kernels with $`\mathrm{\Phi }(p,x)_{}1`$ and, therefore, $`𝐜_𝐊=1`$). Equation (2) then implies
$$\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }(p_n,x_n)_{}\frac{𝐜_𝐊}{2}\sqrt{N},N\{1,2,\mathrm{}\}.$$
(4)
When $`\mathrm{\Phi }`$ is absent (in the sense $`\mathrm{\Phi }1`$), this shows that the forecasts $`p_n`$ are unbiased, in the sense that they are close to $`y_n`$ on average; the presence of $`\mathrm{\Phi }`$ implies, for a suitable kernel, “local unbiasedness”. This is further discussed in the first part of §5.
In the conference version of this paper we also considered the K29 algorithm, which differs from K29 in that $`S_n(p)`$ is defined as
$$S_n(p):=\underset{i=1}{\overset{n1}{}}𝐊((p,x_n),(p_i,x_i))(y_ip_i)$$
and that the requirement that $`𝐊`$ should be forecast-continuous is slightly relaxed (the joint continuity in $`(p,p^{})`$ is replaced by the separate continuity in $`p`$ and $`p^{}`$). For the K29 algorithm, the inequality (2) continues to hold if $`p_n(1p_n)`$ is removed; therefore, (4) continues to hold if the denominator $`2`$ is removed. We will sometimes use “algorithms of large numbers” as generic name for the K29 and K29 algorithms; the motivation for these names is that the main properties of these algorithms are easy corollaries of Kolmogorov’s 1929 proof of the weak law of large numbers.
## 3 Reproducing kernel Hilbert spaces
A *reproducing kernel Hilbert space* (RKHS) on a set $`Z`$ is a Hilbert space $``$ of real-valued functions on $`Z`$ such that the evaluation functional $`ff(z)`$ is continuous for each $`zZ`$. By the Riesz–Fischer theorem, for each $`zZ`$ there exists a function $`𝐊_z`$ such that
$$f(z)=𝐊_z,f_{},f.$$
(5)
The *kernel of RKHS $``$* is
$$𝐊(z,z^{}):=𝐊_z,𝐊_z^{}_{}$$
(6)
(equivalently, we could define $`𝐊(z,z^{})`$ as $`𝐊_z(z^{})`$ or as $`𝐊_z^{}(z)`$). Since (6) is a special case of (1), the function $`𝐊`$ defined by (6) is indeed a kernel on $`Z`$, as defined earlier. On the other hand, for every kernel $`𝐊`$ on $`Z`$ there exists a unique RKHS $``$ on $`Z`$ such that $`𝐊`$ is the kernel of $``$ (see, e.g., , Theorem 2).
A long list of RKHS and the corresponding kernels is given in , §7.4. Perhaps the most interesting RKHS in our current context are various Sobolev spaces $`W^{m,p}(\mathrm{\Omega })`$ ( is the standard reference for the latter). We will be interested in the especially simple space $`W^{1,2}([0,1])`$, to be defined shortly; but first let us make a brief terminological remark. The term “Sobolev space” is usually treated as the name for a topological vector space. All these spaces are normable, but different norms are not considered to lead to different Sobolev spaces as long as the topology does not change.
The *Fermi–Sobolev norm* $`f_{\mathrm{FS}}`$ of a smooth function $`f:[0,1]`$ is defined by
$$f_{\mathrm{FS}}^2:=\left(_0^1f(t)dt\right)^2+_0^1\left(f^{}(t)\right)^2dt.$$
(7)
The *Fermi–Sobolev space* on $`[0,1]`$ is the completion of the set of smooth $`f:[0,1]`$ satisfying $`f_{\mathrm{FS}}<\mathrm{}`$ with respect to the norm $`_{\mathrm{FS}}`$. It is easy to see that it is in fact an RKHS (indeed, if $`f_{\mathrm{FS}}=c<\mathrm{}`$, the mean of $`f`$ is bounded by $`c`$ in absolute value and $`\left|f(b)f(a)\right|_a^b\left|f^{}(t)\right|dtc`$ for all $`0a<b1`$). As a topological vector space, it coincides with the Sobolev space $`W^{1,2}([0,1])`$. The *Fermi–Sobolev space* on $`[0,1]^k`$ is the tensor product of $`k`$ copies of the Fermi–Sobolev space on $`[0,1]`$.
The kernel of the Fermi–Sobolev space on $`[0,1]`$ was found in (see also , §10.2); it is given by
$`𝐊(t,t^{})`$ $`=k_0(t)k_0(t^{})+k_1(t)k_1(t^{})+k_2(|tt^{}|)`$
$`=1+\left(t{\displaystyle \frac{1}{2}}\right)\left(t^{}{\displaystyle \frac{1}{2}}\right)+{\displaystyle \frac{1}{2}}\left(|tt^{}|^2|tt^{}|+{\displaystyle \frac{1}{6}}\right)`$
$`={\displaystyle \frac{1}{2}}\mathrm{min}^2(t,t^{})+{\displaystyle \frac{1}{2}}\mathrm{min}^2(1t,1t^{})+{\displaystyle \frac{5}{6}},`$ (8)
where $`k_l:=B_l/l!`$ are scaled Bernoulli polynomials $`B_l`$. We will derive the final expression for $`𝐊(t,t^{})`$ in (8) in Appendix C. For the Fermi–Sobolev space on $`[0,1]^k`$ we have
$$𝐊((t_1,\mathrm{},t_k),(t_1^{},\mathrm{},t_k^{}))=\underset{i=1}{\overset{k}{}}\left(\frac{1}{2}\mathrm{min}^2(t_i,t_i^{})+\frac{1}{2}\mathrm{min}^2(1t_i,1t_i^{})+\frac{5}{6}\right)$$
(9)
and, therefore,
$$𝐜_𝐊^2=\underset{t[0,1]}{\mathrm{max}}(\frac{1}{2}t^2+\frac{1}{2}(1t)^2+\frac{5}{6})^k=\left(\frac{4}{3}\right)^k.$$
(10)
For further information about the Fermi–Sobolev spaces, see .
## 4 The K29 algorithm in RKHS
We can now deduce the following corollary from Theorem 1.
###### Theorem 2
Let $``$ be an RKHS on $`[0,1]\times 𝐗`$ with a forecast-continuous kernel $`𝐊`$. The K29 algorithm with parameter $`𝐊`$ ensures
$$\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n,x_n)\right|f_{}\sqrt{\underset{n=1}{\overset{N}{}}p_n(1p_n)𝐊((p_n,x_n),(p_n,x_n))}$$
(11)
for all $`N`$ and all $`f`$.
* Applying K29 to the feature mapping $`(p,x)[0,1]\times 𝐗𝐊_{p,x}`$ and using (2), we obtain, for any $`f`$:
$$\begin{array}{c}\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n,x_n)\right|=\left|\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{p_n,x_n},f_{}\right|\hfill \\ \hfill =\left|\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{p_n,x_n},f_{}\right|\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{p_n,x_n}_{}f_{}\\ \hfill f_{}\sqrt{\underset{n=1}{\overset{N}{}}p_n(1p_n)𝐊((p_n,x_n),(p_n,x_n))}.\end{array}$$
When $`𝐜_𝐊`$ in (3) is finite, (11) implies
$$\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n,x_n)\right|\frac{𝐜_𝐊}{2}f_{}\sqrt{N}.$$
(11)
## 5 Informal discussion
In this section we explain why the inequalities in Theorems 1 and 2 can be interpreted as results about calibration and resolution, and then briefly discuss a puzzling aspect of the algorithms of large numbers. For concreteness, we usually talk about the K29 algorithm, but all we say can also be applied, with obvious modifications, to K29.
### Calibration, resolution, and calibration-cum-resolution
We start from the intuitive notion of calibration (for further details, see and ). The forecasts $`p_n`$, $`n=1,\mathrm{},N`$, are said to be “well calibrated” (or “unbiased in the small”, or “reliable”, or “valid”) if, for any $`p^{}[0,1]`$,
$$\frac{_{n=1,\mathrm{},N:p_np^{}}y_n}{_{n=1,\mathrm{},N:p_np^{}}1}p^{}$$
(12)
provided $`_{n=1,\mathrm{},N:p_np^{}}1`$ is not too small. The interpretation of (12) is that the forecasts should be in agreement with the observed frequencies. It will be convenient to rewrite (12) as
$$\frac{_{n=1,\mathrm{},N:p_np^{}}(y_np_n)}{_{n=1,\mathrm{},N:p_np^{}}1}0.$$
(13)
The fact that good calibration is only a necessary condition for good forecasting performance can be seen from the following standard example : if
$$(y_1,y_2,y_3,y_4,\mathrm{})=(1,0,1,0,\mathrm{}),$$
the forecasts $`p_n=1/2`$, $`n=1,2,\mathrm{}`$, are well calibrated but rather poor; it would be better to forecast with
$$(p_1,p_2,p_3,p_4,\mathrm{})=(1,0,1,0,\mathrm{}).$$
Assuming that each datum $`x_n`$ contains the information about the parity of $`n`$ (which can always be added to $`x_n`$), we can see that the problem with the forecasting strategy $`p_n1/2`$ is its lack of resolution: it does not distinguish between the data with odd and even $`n`$. In general, we would like each forecast $`p_n`$ to be as specific as possible to the current datum $`x_n`$; the resolution of a forecasting algorithm is the degree to which it achieves this goal (taking it for granted that $`x_n`$ contains all relevant information).
Analogously to (13), the forecasts $`p_n`$, $`n=1,\mathrm{},N`$, may be said to have good resolution if, for any $`x^{}𝐗`$,
$$\frac{_{n=1,\mathrm{},N:x_nx^{}}(y_np_n)}{_{n=1,\mathrm{},N:x_nx^{}}1}0$$
provided the denominator is not too small. We can also require that the forecasts $`p_n`$, $`n=1,\mathrm{},N`$, should have good “calibration-cum-resolution”: for any $`(p^{},x^{})[0,1]\times 𝐗`$,
$$\frac{_{n=1,\mathrm{},N:(p_n,x_n)(p^{},x^{})}(y_np_n)}{_{n=1,\mathrm{},N:(p_n,x_n)(p^{},x^{})}1}0$$
provided the denominator is not too small. Notice that even if forecasts have both good calibration and good resolution, they can still have poor calibration-cum-resolution.
It is easy to see that (4) implies good calibration-cum-resolution for a suitable $`\mathrm{\Phi }`$ and large $`N`$: indeed, (4) shows that the forecasts $`p_n`$ are unbiased in the neighborhood of each $`(p^{},x^{})`$ for functions $`\mathrm{\Phi }`$ that map distant $`(p,x)`$ and $`(p^{},x^{})`$ to almost orthogonal elements of the feature space (such as $`\mathrm{\Phi }`$ corresponding to the Gaussian kernel
$$𝐊((p,x),(p^{},x^{})):=\mathrm{exp}\left(\frac{(pp^{})^2+xx^{}^2}{2\sigma ^2}\right)$$
(14)
for a small “kernel width” $`\sigma >0`$).
In general, to make sense of the $``$ in the numerator and denominator of, say, (13), we replace each “crisp” point $`p^{}`$ by a “fuzzy point” $`I_p^{}:[0,1][0,1]`$; $`I_p^{}`$ is required to be continuous, and we might also want to have $`I_p^{}(p^{})=1`$ and $`I_p^{}(p)=0`$ for all $`p`$ outside a small neighborhood of $`p^{}`$. The alternative of choosing $`I_p^{}:=𝕀_{[p_{},p_+]}`$, where $`[p_{},p_+]`$ is a short interval containing $`p^{}`$ and $`𝕀_{[p_{},p_+]}`$ is its indicator function, does not work because of Oakes’s and Dawid’s examples ; $`I_p^{}`$ can, however, be arbitrarily close to $`𝕀_{[p_{},p_+]}`$.
Consider, e.g., the following approximation to the indicator function of a short interval $`[p_{},p_+]`$ containing $`p^{}`$:
$$f(p):=\{\begin{array}{cc}1\hfill & \text{if }p_{}+ϵpp_+ϵ\hfill \\ 0\hfill & \text{if }pp_{}ϵ\text{ or }pp_++ϵ\hfill \\ \frac{1}{2}+\frac{1}{2ϵ}(pp_{})\hfill & \text{if }p_{}ϵpp_{}+ϵ\hfill \\ \frac{1}{2}+\frac{1}{2ϵ}(p_+p)\hfill & \text{if }p_+ϵpp_++ϵ;\hfill \end{array}$$
(15)
we assume that $`ϵ>0`$ satisfies
$$0<p_{}ϵ<p_{}+ϵ<p_+ϵ<p_++ϵ<1.$$
It is clear that this approximation belongs to the Fermi–Sobolev space. An easy computation shows that (11) and (10) imply
$$\left|\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n)\right|\frac{1}{\sqrt{3}}\sqrt{\left(\frac{1}{ϵ}+(p_+p_{})^2\right)N}$$
(16)
for all $`N`$. We can see that (13), in the form
$$\frac{_{n=1,\mathrm{},N}f(p_n)(y_np_n)}{_{n=1,\mathrm{},N}f(p_n)}0,$$
will hold if
$$\underset{n=1}{\overset{N}{}}f(p_n)\sqrt{N}$$
(roughly, if significantly more than $`\sqrt{N}`$ forecasts fall in the neighborhood $`[p^{},p^+]`$ of $`p^{}`$).
It is clear that inequalities analogous to (16) can also be proved for “soft neighborhoods” of points $`(p^{},x^{})`$ in $`[0,1]\times 𝐗`$ (at least when $`𝐗`$ is a domain in a Euclidean space), and so Theorem 2 also implies good calibration-cum-resolution for large $`N`$. Convenient neighborhoods in $`[0,1]\times [0,1]^K`$ can be constructed as tensor products of neighborhoods (15).
Inequality (16) and analogous inequalities expressing resolution and calibration-cum-resolution are explicit in the sense that they do not involve limits, $`o`$, $`O`$, unspecified constants, etc. The price to pay is their relative complexity; therefore, we also state a simple asymptotic result about calibration-cum-resolution.
###### Corollary 1
If $`𝐗`$ is a compact metric space, some forecasting algorithm guarantees
$$\underset{N\mathrm{}}{lim}\frac{1}{N}\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n,x_n)=0$$
(17)
for all continuous functions $`f:[0,1]\times 𝐗`$.
Calibration corresponds to the case where $`f(p,x)=I_p^{}(p)`$ does not depend on $`x`$ and resolution to the case where $`f(p,x)=I_x^{}(x)`$ does not depend on $`p`$. This result was proved in in the case of calibration (there are no $`x_n`$) and Lipschitz functions $`f`$.
* Let $``$ be an RKHS on $`[0,1]\times 𝐗`$ which is *universal*, i.e., dense in the space $`C([0,1]\times 𝐗)`$, and whose kernel $`𝐊`$ is continuous and satisfies $`𝐜_𝐊<\mathrm{}`$. The notion of universality is introduced in , Definition 4, and the existence of such an $``$ is shown in , Theorem 2. For any continuous function $`f:[0,1]\times 𝐗`$ there is a $`g`$ that is $`ϵ`$-close to $`f`$ in the metric $`C([0,1]\times 𝐗)`$, and so, by (11),
$$\begin{array}{c}\underset{N\mathrm{}}{lim\; sup}\left|\frac{1}{N}\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n,x_n)\right|\underset{N\mathrm{}}{lim\; sup}\left|\frac{1}{N}\underset{n=1}{\overset{N}{}}(y_np_n)g(p_n,x_n)\right|+ϵ\hfill \\ \hfill \underset{N\mathrm{}}{lim\; sup}\frac{1}{N}\frac{𝐜_𝐊}{2}g_{}\sqrt{N}+ϵ=ϵ;\end{array}$$
since this holds for any $`ϵ>0`$, (17) also holds.
One of the algorithms achieving (17) for $`𝐗=[0,1]^k`$ is K29 applied to the Fermi–Sobolev kernel (9). It is interesting, and somewhat counterintuitive, that K29 applied to the Gaussian kernel (14) (with any $`\sigma >0`$) also achieves (17); the universality of the Gaussian kernels is proved in (Example 1).
Our discussion of calibration and resolution in this subsection has been somewhat speculative, and the reader might ask whether these two properties are really useful. This question is answered, to some degree, in , which show that probability forecasts satisfying these properties lead to good decisions (at least in the simple decision protocols considered in those papers).
### Puzzle of the iterated logarithm
Theorems 1 and 2 imply that the forecasts produced by the K29 algorithm are even closer to the actual observations on average than in the case of “genuine randomness”, where Reality produces the data and observations from a probability distribution on $`(𝐗\times \{0,1\})^{\mathrm{}}`$ and each $`p_n`$ is the conditional probability that $`y_n=1`$ given $`x_1,\mathrm{},x_n`$, $`y_1,\mathrm{},y_{n1}`$, and whatever further information may be available at this point. Indeed, let us take, for simplicity, $`\mathrm{\Phi }1`$ (and $`:=`$) in Theorem 1. According to the martingale law of the iterated logarithm (see, e.g., or Chapter 5 of ), we would expect
$$\underset{N\mathrm{}}{lim\; sup}\frac{\left|_{n=1}^N(y_np_n)\right|}{\sqrt{2A_N\mathrm{ln}\mathrm{ln}A_N}}=1,$$
(18)
where $`A_N:=_{n=1}^Np_n(1p_n)`$ is assumed to tend to $`\mathrm{}`$ as $`N\mathrm{}`$, and so expect, contrary to (4),
$$\underset{N\{1,2,\mathrm{}\}}{sup}\frac{_{n=1}^N(y_np_n)\mathrm{\Phi }(p_n,x_n)_{}}{\sqrt{N}}$$
to be infinite for $`p_n`$ not consistently very close to 0 or 1. Actually, in this case ($`\mathrm{\Phi }1`$) Forecaster can even make sure that
$$\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }(p_n,x_n)_{}=\frac{1}{2},N\{1,2,\mathrm{}\}$$
(choosing $`p_1:=1/2`$ and $`p_n:=y_{n1}`$, $`n=2,3,\mathrm{}`$).
For a general $`\mathrm{\Phi }`$, we can also expect that the probabilities $`p_n`$ contrived by the algorithms of large numbers (K29 or K29) will have better calibration and resolution than the true probabilities. There is, however, little doubt that the true probabilities are more useful than any probabilities we are able to come up with. The true probabilities are not as good at calibration and resolution, so they must be better in some other equally important respects. It remains unclear what these other respects may be, and this is what we call the puzzle of the iterated logarithm.
## 6 Optimality of the K29 algorithm
In this section we establish that the inequalities in Theorems 1 and 2 are tight, in a natural sense.
Equation (2) says that the differences $`y_np_n`$ are small on average, even when scattered in a Hilbert space by multiplying by $`\mathrm{\Phi }(p_n,x_n)`$. The next result says that it is the best Forecaster can do.
###### Theorem 3
Let $`\mathrm{\Phi }:[0,1]\times 𝐗`$, where $``$ is a Hilbert space. There is a strategy for Reality II which guarantees that
$$\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }(p_n,x_n)_{}^2\underset{n=1}{\overset{N}{}}p_n(1p_n)\mathrm{\Phi }(p_n,x_n)_{}^2$$
(19)
always holds for all $`N=1,2,\mathrm{}`$, regardless of what the other players do.
* Set
$$R_N:=\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }(p_n,x_n)_{},N=1,2,\mathrm{};$$
it is sufficient to show that on the $`N`$th round, $`N=1,2,\mathrm{}`$, Reality II can ensure that
$$R_N^2R_{N1}^2p_N(1p_N)\mathrm{\Phi }_N^2,$$
(20)
where
$$\mathrm{\Phi }_N:=\mathrm{\Phi }(p_N,x_N)_{}.$$
Fix an $`N`$. Define points $`A,C,D`$ as
$`C`$ $`:={\displaystyle \underset{n=1}{\overset{N1}{}}}(y_np_n)\mathrm{\Phi }(p_n,x_n),`$
$`A`$ $`:={\displaystyle \underset{n=1}{\overset{N1}{}}}(y_np_n)\mathrm{\Phi }(p_n,x_n)+(1p_N)\mathrm{\Phi }(p_N,x_N),`$
$`D`$ $`:={\displaystyle \underset{n=1}{\overset{N1}{}}}(y_np_n)\mathrm{\Phi }(p_n,x_n)+(p_N)\mathrm{\Phi }(p_N,x_N);`$
it is up to Reality II whether make $`R_N`$ equal to $`|OA|`$ or $`|OD|`$, where $`O`$ is the origin. Assuming, without loss of generality, that $`R_N=\mathrm{max}(|OA|,|OD|)`$, we reduce our task to showing that the maximal value of $`R_{N1}`$ for fixed $`R_N`$, $`\mathrm{\Phi }_N`$, and $`p_N`$ satisfies (20). It is geometrically obvious (see the last paragraph of this proof for a rigorous argument) that $`R_{N1}`$ attains its maximal value when $`|OA|=|OD|`$; this is illustrated in Figure 1 (remember that all four points, $`O`$, $`A`$, $`C`$, and $`D`$, lie in the same plane). Let $`B`$ be the base of the perpendicular dropped from $`O`$ onto the interval $`AD`$ and $`h:=|OB|`$. Since the triangles $`OBD`$ and $`OBC`$ are right-angled,
$`R_N^2`$ $`=h^2+\left({\displaystyle \frac{1}{2}}\mathrm{\Phi }_N\right)^2,`$
$`R_{N1}^2`$ $`=h^2+\left({\displaystyle \frac{1}{2}}\mathrm{\Phi }_Np_N\mathrm{\Phi }_N\right)^2.`$
Subtracting the second equality from the first, we obtain
$$R_N^2R_{N1}^2=\left(\frac{1}{2}\mathrm{\Phi }_N\right)^2\left(\frac{1}{2}\mathrm{\Phi }_Np_N\mathrm{\Phi }_N\right)^2=p_N(1p_N)\mathrm{\Phi }_N^2.$$
In conclusion, let us see that the maximum of $`R_{N1}`$ is indeed attained when $`|OA|=|OD|`$. Assume that $`|OA|=R_N`$, with $`|OD|`$ now allowed to be less than $`R_N`$. Because of the compactness of the disk in Figure 1 (we are only interested in two-dimensional subspaces of $``$, which are isometrically isomorphic to $`^2`$), the maximum of $`|OC|`$ is attained at some point $`C`$. Supposing $`|OD|<R_N`$, it is, however, easy to check that no $`C`$ will be a point of local maximum for $`|OC|`$; the least trivial case is perhaps where $`O`$ lies on the line $`AD`$ and $`C`$ is between $`O`$ and $`D`$.
The next result establishes the tightness of the bound in Theorem 2.
###### Theorem 4
Let $``$ be an RKHS on $`[0,1]\times 𝐗`$ with kernel $`𝐊`$. Reality II has a strategy which ensures, regardless of what the other players do, that for each $`N=1,2,\mathrm{}`$ there exists a non-zero $`f`$ such that
$$\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n,x_n)f_{}\sqrt{\underset{n=1}{\overset{N}{}}p_n(1p_n)𝐊((p_n,x_n),(p_n,x_n))}.$$
(21)
* By Theorem 3 there exists a strategy for Reality II which ensures
$$\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{p_n,x_n}_{}\sqrt{\underset{n=1}{\overset{N}{}}p_n(1p_n)𝐊((p_n,x_n),(p_n,x_n))}.$$
(22)
Taking
$$f:=\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{p_n,x_n},$$
we obtain:
$$\begin{array}{c}\underset{n=1}{\overset{N}{}}(y_np_n)f(p_n,x_n)=\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{p_n,x_n},f_{}\hfill \\ \hfill =\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{p_n,x_n},f_{}=\underset{n=1}{\overset{N}{}}(y_np_n)𝐊_{p_n,x_n}_{}f_{}\\ \hfill f_{}\sqrt{\underset{n=1}{\overset{N}{}}p_n(1p_n)𝐊((p_n,x_n),(p_n,x_n))}.\end{array}$$
If $`f0`$, our task is accomplished. Otherwise, the right-hand side of (22) will also be zero, and we can take any $`f0`$.
### Acknowledgments
I am grateful to Ilia Nouretdinov for a discussion that lead to the proof of Theorem 3 and to the anonymous reviewers of the conference and journal versions of this paper for their comments. This work was partially supported by MRC (grant S505/65) and Royal Society.
## Appendix A Proof of Theorem 1
The proof of Theorem 1 is based on the game-theoretic approach to the foundations of probability proposed in . A new player, called Skeptic, is added to the learning protocol of §2; the idea is that Skeptic is allowed to bet at the odds defined by Forecaster’s probabilities. In this proof there is no need to distinguish between Reality I and Reality II.
Binary Forecasting Game I
Players: Reality, Forecaster, Skeptic
Protocol:
$`𝒦_0:=C`$.
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Forecaster announces $`p_n[0,1]`$.
Skeptic announces $`s_n`$.
Reality announces $`y_n\{0,1\}`$.
$`𝒦_n:=𝒦_{n1}+s_n(y_np_n)`$.
END FOR.
The protocol describes not only the players’ moves but also the changes in Skeptic’s capital $`𝒦_n`$; its initial value is an arbitrary constant $`C`$.
The crucial (albeit very simple) observation is that for any continuous strategy for Skeptic there exists a strategy for Forecaster that does not allow Skeptic’s capital to grow, regardless of what Reality is doing (similar observations were made in and ). To state this observation in its strongest form, we will make Skeptic announce his strategy for each round before Forecaster’s move on that round rather than announce his full strategy at the beginning of the game. Therefore, we consider the following perfect-information game:
Binary Forecasting Game II
Players: Reality, Forecaster, Skeptic
Protocol:
$`𝒦_0:=C`$.
FOR $`n=1,2,\mathrm{}`$:
Reality announces $`x_n𝐗`$.
Skeptic announces continuous $`S_n:[0,1]`$.
Forecaster announces $`p_n[0,1]`$.
Reality announces $`y_n\{0,1\}`$.
$`𝒦_n:=𝒦_{n1}+S_n(p_n)(y_np_n)`$.
END FOR.
###### Lemma 1
Forecaster has a strategy in Binary Forecasting Game II that ensures $`𝒦_0𝒦_1𝒦_2\mathrm{}`$.
* Forecaster can use the following strategy to ensure $`𝒦_0𝒦_1\mathrm{}`$:
+ if $`S_n(0)`$ and $`S_n(1)`$ are both positive or both negative, take $`p_n:=(1+signS_n(0))/2`$;
+ otherwise, choose $`p_n`$ so that $`S_n(p_n)=0`$ (such a $`p_n`$ will exist).
A measure-theoretic version of Lemma 1 (involving randomization) was proved in , Proposition 1.
### Proof of the theorem
We start by noticing that
$$(y_np_n)^2=p_n(1p_n)+(12p_n)(y_np_n)$$
(23)
both for $`y_n=0`$ and for $`y_n=1`$. Following K29, Forecaster ensures that Skeptic will never increase his capital with the strategy
$$s_n:=\underset{i=1}{\overset{n1}{}}𝐊((p_n,x_n),(p_i,x_i))(y_ip_i)+\frac{1}{2}𝐊((p_n,x_n),(p_n,x_n))(12p_n)$$
(24)
(continuous in $`p_n`$ by our assumptions). The increase in Skeptic’s capital when he follows (24) is
$`𝒦_N𝒦_0`$ $`={\displaystyle \underset{n=1}{\overset{N}{}}}s_n(y_np_n)`$
$`={\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{n1}{}}}𝐊((p_n,x_n),(p_i,x_i))(y_np_n)(y_ip_i)`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}𝐊((p_n,x_n),(p_n,x_n))(12p_n)(y_np_n)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}𝐊((p_n,x_n),(p_i,x_i))(y_np_n)(y_ip_i)`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}𝐊((p_n,x_n),(p_n,x_n))(y_np_n)^2`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}𝐊((p_n,x_n),(p_n,x_n))(12p_n)(y_np_n)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}𝐊((p_n,x_n),(p_i,x_i))(y_np_n)(y_ip_i)`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{N}{}}}𝐊((p_n,x_n),(p_n,x_n))p_n(1p_n)`$
(we used (23) in the last equality). We can rewrite this as
$$\begin{array}{c}𝒦_N𝒦_0=\frac{1}{2}\underset{n=1}{\overset{N}{}}(y_np_n)\mathrm{\Phi }(p_n,x_n)_{}^2\frac{1}{2}\underset{n=1}{\overset{N}{}}p_n(1p_n)\mathrm{\Phi }(p_n,x_n)_{}^2,\hfill \end{array}$$
which immediately implies (2).
## Appendix B Forecast-continuity of feature mappings and kernels
In this appendix we will prove, essentially following , Lemma 3, that the forecast-continuity of a kernel $`𝐊`$ on $`[0,1]\times 𝐗`$ is equivalent to the continuity in $`p`$ of a feature mapping $`\mathrm{\Phi }(p,x)`$ satisfying (1). As a byproduct, we will also see that the forecast-continuity of a kernel $`𝐊`$ on $`[0,1]\times 𝐗`$ can be equivalently defined by requiring that
* $`𝐊((p,x),(p^{},x))`$ should be continuous in $`p`$, for all $`x𝐗`$ and all $`p^{}[0,1]`$,
* and $`𝐊((p,x),(p,x))`$ should be continuous in $`p`$, for all $`x𝐗`$.
In one direction the statement is obvious: if $`\mathrm{\Phi }(p,x)`$ is continuous in $`p`$, the continuity of the operation of taking the inner product immediately implies that $`𝐊`$ is forecast-continuous, in both senses.
Now suppose that $`𝐊`$ is forecast-continuous, as defined in the first paragraph of this appendix (this is the apparently weaker sense of forecast-continuity). To complete the proof, notice that
$`\mathrm{\Phi }(p,x)\mathrm{\Phi }(p_n,x)_{}`$
$`=\sqrt{𝐊((p,x),(p,x))2𝐊((p,x),(p_n,x))+𝐊((p_n,x),(p_n,x))}`$
$`\sqrt{𝐊((p,x),(p,x))2𝐊((p,x),(p,x))+𝐊((p,x),(p,x))}=0`$
when $`p_np`$ ($`n\mathrm{}`$).
## Appendix C Derivation of the kernel of the Fermi–Sobolev space
We first describe the standard reduction of the problem of finding the kernel of an RKHS to a variational problem. Let $`𝐊`$ be the kernel of an RKHS $``$ on $`Z`$.
Let $`cZ`$. According to (Satz III.3), the minimum of $`f_{}`$ among the functions $`f`$ satisfying $`f(c)=1`$ is attained by the function $`𝐊(,c)/𝐊(c,c)`$. Therefore, we obtain a function $`k(,c)`$ proportional to $`𝐊(,c)`$ by solving the optimization problem $`f_{}\mathrm{min}`$ under the constraint $`f(c)=1`$ (or under the constraint $`f(c)=d`$, where $`d`$ is any other constant). It remains to find the coefficient of proportionality in terms of $`k(,c)`$. If $`𝐊(,)=\alpha k(,)`$, we have:
$`𝐊(c,c)`$ $`=𝐊(,c)_{}^2;`$
$`\alpha k(c,c)`$ $`=\alpha ^2k(,c)_{}^2;`$
$`\alpha `$ $`={\displaystyle \frac{k(c,c)}{k(,c)_{}^2}}.`$
Therefore, the recipe for finding $`𝐊`$ is: for each $`cZ`$ solve the optimization problem $`f_{}\mathrm{min}`$ under the constraint $`f(c)=1`$ (the completeness of RKHS implies that the minimum is attained) and set
$$𝐊(z,c):=\frac{k(z,c)k(c,c)}{k(,c)_{}^2},$$
(25)
where $`k(,c)`$ is the solution.
Now let us apply this technique to finding the kernel corresponding to the Fermi–Sobolev space on $`[0,1]`$ with the norm given by (7). Let $`c[0,1]`$ and let $`f`$ be the solution to the optimization problem $`f_{}\mathrm{min}`$ under the constraint $`f(c)=1`$ (because of the convexity of the set $`\{f|f(c)=1\}`$, there is only one solution). First we show that the derivative $`f^{}`$ is a linear function on $`[0,c]`$ and on $`[c,1]`$, arguing indirectly. Suppose, for concreteness, that $`f^{}`$ is not linear on the interval $`(0,c)`$; in particular this interval is non-empty. There are three points $`0<t_1<t_2<t_3<c`$ such that
$$f^{}(t_2)\frac{t_3t_2}{t_3t_1}f^{}(t_1)+\frac{t_2t_1}{t_3t_1}f^{}(t_3).$$
(26)
For a small constant $`ϵ>0`$ (in particular, we assume $`2ϵ<\mathrm{min}(t_1,t_2t_1,t_3t_2,ct_3)`$), let $`g:[0,1]`$ be a smooth function such that $`_0^1g(t)dt=0`$ and:
* $`g(t)=0`$ for $`t<t_1ϵ`$;
* $`g(t)`$ is increasing for $`t_1ϵ<t<t_1+ϵ`$;
* $`g(t)=t_3t_2`$ for $`t_1+ϵ<t<t_2ϵ`$;
* $`g(t)`$ is decreasing for $`t_2ϵ<t<t_2+ϵ`$;
* $`g(t)=(t_2t_1)`$ for $`t_2+ϵ<t<t_3ϵ`$;
* $`g(t)`$ is increasing for $`t_3ϵ<t<t_3+ϵ`$;
* $`g(t)=0`$ for $`t>t_3+ϵ`$.
Since, for any $`\delta `$ (we are interested in nonzero $`\delta `$ small in absolute value),
$$f+\delta g_{FS}^2=f_{FS}^2+2\delta _0^1f^{}(t)g^{}(t)dt+\delta ^2_0^1(g^{}(t))^2dt,$$
the definition of $`f`$ implies
$$_0^1f^{}(t)g^{}(t)dt=0.$$
However, as $`ϵ0`$, the last integral tends to
$$f^{}(t_1)(t_3t_2)f^{}(t_2)(t_3t_1)+f^{}(t_3)(t_2t_1),$$
which cannot, by (26), be zero.
Once we know that $`f`$ is a quadratic polynomial to the left and to the right of $`c`$, we can easily find (this can be done conveniently using a computer algebra system) that, ignoring a multiplicative constant,
$$f(t)=3t^2+3c^26c+8=3t^2+3(1c)^2+5$$
to the left of $`c`$ and
$$f(t)=3t^2+3c^26t+8=3(1t)^2+3c^2+5$$
to the right of $`c`$. By (25), we can now find
$$𝐊(t,c)=\frac{f(t)f(c)}{f_{}^2}=f(t)/6,$$
which agrees with (8).
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# 1 Introduction
## 1 Introduction
The discovery of the pentaquark state $`\mathrm{\Theta }^+(1540)`$ has opened a new field of strong interaction and provides a new opportunity for a deeper understanding of the low energy QCD. Intense theoretical studies have been motivated to clarify the quantum numbers and to understand the under-structures of the pentaquark state $`\mathrm{\Theta }^+(1540)`$ . Although the existence of the $`\mathrm{\Theta }^+(1540)`$ is uncertain and is a subject of controversy now, the $`\mathrm{\Theta }^+(1540)`$ has already contributed to hadron spectroscopy. Unlike the chiral soliton model, the quark models take the constituent quarks or quark clusters as the elementary degrees of freedom, there exist a great number of possible quark configurations satisfy the Fermi statistics and the color singlet condition for the substructures of the pentaquark state $`\mathrm{\Theta }^+(1540)`$ if we release stringent dynamical constraints. In fact, the multiquark states are many-body problems, they are very difficult to solve. Whether or not the quarks can cluster together to form diquarks is of great importance theoretically, if we take the diquarks as the basic constituents (here the ”basic constituents” does not mean they are asymptotic states, they just exist inside the baryons or multiquark states with typical length), the problems will be greatly simplified, for example, the $`\frac{1}{2}^+`$ octet and $`\frac{3}{2}^+`$ decuplet baryons can be taken as quark-diquark bound states in the Faddeev approximation , the nonet scalar mesons below $`1GeV`$ can be taken as 4-quark states $`(qq)_{\overline{3}}(\overline{q}\overline{q})_3`$ with scalar diquarks or pseudoscalar diquarks as their basic constituents ; furthermore, we can obtain more insight into the relevant degrees of freedom and deepen our understanding about the underlying dynamics that determines the properties of the baryons and exotic multiquark states. A typical quark model for the pentaquark states is the Jaffe-Wilczek’s diquark-diquark-antiquark model . In this model, the scalar diquarks $`U^a=ϵ^{abc}d_b^T(x)C\gamma _5s_c(x)`$, $`D^a=ϵ^{abc}u_b^T(x)C\gamma _5s_c(x)`$ , $`S^a=ϵ^{abc}u_b^T(x)C\gamma _5d_c(x)`$ are taken as the basic constituents. They belong to the antitriplet $`\overline{3}`$ representation of both the color $`SU(3)_c`$ group and flavor $`SU(3)_f`$ group, in the color superconductivity theory, the attractive interactions in this channel lead to the formulation of nonzero condensates and breaking of both the color and flavor $`SU(3)`$ symmetries . The scalar diquarks correspond to the $`{}_{}{}^{1}S_{0}^{}`$ states of the diquark systems, the one-gluon exchange force and the instanton induced force can lead to significant attractions between the quarks in the $`0^+`$ channels . The pseudoscalar diquarks do not have nonrelativistic limit, can be taken as the $`{}_{}{}^{3}P_{0}^{}`$ states. As the instanton induced force results in strong attractions in the scalar diquark channel and strong repulsions in the pseudoscalar diquark channel, if the effects of the instanton are manifested, we prefer the $`S`$ type diquark to the $`P`$ type diquark in constructing interpolating currents in the QCD sum rules .
In this article, we take the point of view that the scalar diquarks are quasi-bound states of quark-quark system and study the ground state mass spectrum within the framework of the coupled Schwinger-Dyson equation (SDE) and Bethe-Salpeter equation (BSE). The coupled rainbow SDE and ladder BSE have given a lot of successful descriptions of the long distance properties of the low energy QCD and the QCD vacuum (for reviews, one can see Refs.). The SDE can naturally embody the dynamical symmetry breaking and confinement which are two crucial features of QCD, although they correspond to two very different energy scales . On the other hand, the BSE is a conventional approach in dealing with the two-body relativistic bound state problems . From the solutions of the BSE, we can obtain useful information about the under-structures of the mesons and diquarks, and obtain powerful tests for the quark theory. However, the obviously drawback may be the model dependent kernels for the gluon two-point Green’s function and the truncations for the coupled divergent SDE and BSE series in one or the other ways . Many analytical and numerical calculations indicate that the coupled rainbow SDE and ladder BSE with phenomenological potential models can give model independent results and satisfactory values . The usually used effective potential models are confining Dirac $`\delta `$ function potential, Gaussian distribution potential and flat bottom potential (FBP) . The FBP is a sum of Yukawa potentials, which not only satisfies chiral invariance and fully relativistic covariance, but also suppresses the singular point that the Yukawa potential has. It works well in understanding the dynamical chiral symmetry breaking, confinement and the QCD vacuum as well as the meson structures, such as electromagnetic form factors, radius, decay constants . In this article, we use the FBP to study the ground state mass spectrum of the scalar diquarks without fine tuning such as modifying the infrared behavior for the heavy quark systems.
The article is arranged as follows: we introduce the FBP in section II; in section III, IV and V, we solve the rainbow SDE and ladder BSE, explore the analyticity of the quark propagators, study the dynamical symmetry breaking and confinement, finally obtain the mass spectrum of the $`\pi `$, $`K`$ mesons and the scalar $`U^a`$, $`D^a`$, $`S^a`$ diquarks, and the decay constants of the $`\pi `$, $`K`$ mesons; section VI is reserved for conclusion.
## 2 Flat Bottom Potential
The present techniques in QCD calculation can not give satisfactory large $`r`$ behavior for the gluon two-point Green’s function to implement the linear potential confinement mechanism, in practical calculation, the phenomenological effective potential models always do the work. The FBP is a sum of Yukawa potentials which is an analogy to the exchange of a series of particles and ghosts with different masses (Euclidean Form),
$$G(k^2)=\underset{j=0}{\overset{n}{}}\frac{a_j}{k^2+(N+j\rho )^2},$$
(1)
where $`N`$ stands for the minimum value of the masses, $`\rho `$ is their mass difference, and $`a_j`$ is their relative coupling constant. Due to the particular condition we take for the FBP, there is no divergence in solving the SDE. In its three dimensional form, the FBP takes the following form:
$$V(r)=\underset{j=0}{\overset{n}{}}a_j\frac{\mathrm{e}^{(N+j\rho )r}}{r}.$$
(2)
In order to suppress the singular point at $`r=0`$, we take the following conditions:
$`V(0)=constant,`$
$`{\displaystyle \frac{dV(0)}{dr}}={\displaystyle \frac{d^2V(0)}{dr^2}}=\mathrm{}={\displaystyle \frac{d^nV(0)}{dr^n}}=0.`$ (3)
The $`a_j`$ can be determined by solving the equations inferred from the flat bottom condition in Eq.(3). As in previous literature , $`n`$ is set to be 9.
## 3 Schwinger-Dyson equation
The SDE can provide a natural framework for studying the nonperturbative properties of the quark and gluon Green’s functions. By studying the evolution behavior and analytic structure of the dressed quark propagators, we can obtain valuable information about the dynamical chiral symmetry breaking and confinement. In the following, we write down the rainbow SDE for the quark propagator,
$$S^1(p)=i\gamma p+\widehat{m}_{u,d,s}+4\pi \frac{d^4k}{(2\pi )^4}\gamma _\mu \frac{\lambda ^a}{2}S(k)\gamma _\nu \frac{\lambda ^a}{2}G_{\mu \nu }(kp),$$
(4)
where
$`S^1(p)`$ $`=`$ $`iA(p^2)\gamma p+B(p^2)A(p^2)[i\gamma p+m(p^2)],`$ (5)
$`G_{\mu \nu }(k)`$ $`=`$ $`(\delta _{\mu \nu }{\displaystyle \frac{k_\mu k_\nu }{k^2}})G(k^2),`$ (6)
and $`\widehat{m}_{u,d,s}`$ stands for the current quark mass that breaks chiral symmetry explicitly. For a short discussion about the full SDE for the quark propagator, one can consult Ref..
In this article, we assume that a Wick rotation to Euclidean variables is allowed, and perform a rotation analytically continuing $`p`$ and $`k`$ into the Euclidean region. Alternatively, one can derive the SDE from the Euclidean path-integral formulation of the theory, thus avoiding possible difficulties in performing the Wick rotation . As far as only numerical results are concerned, the two procedures are equal. In fact, the analytical structures of quark propagators have interesting information about confinement, we will make detailed discussion about the propagators of the $`u`$, $`d`$ and $`s`$ quarks in section V.
## 4 Bethe-Salpeter equation
The BSE is a conventional approach in dealing with the two-body relativistic bound state problems . The precise knowledge about the quark structures of the mesons and diquarks can result in better understanding of their properties. In the following, we write down the ladder BSE for the scalar diquark quasi-bound states with two quarks of different flavor ,
$`\mathrm{\Gamma }_{\overline{3}}(q,P)=`$
$`4\pi {\displaystyle \frac{d^4k}{(2\pi )^4}G_{\mu \nu }(qk)\gamma _\mu \frac{\lambda ^a}{2}S(k+\xi P)\mathrm{\Gamma }_{\overline{3}}(k,P)S^T(k+(1\xi )P)(\gamma _\nu \frac{\lambda ^a}{2})^T}.`$
Here $`T`$ denotes matrix transpose, the $`S(k)`$ is the quark propagator, $`G_{\mu \nu }(k)`$ is the gluon propagator, $`P_\mu `$ is the four-momentum of the center of mass of the scalar diquark, $`q_\mu `$ is the relative four-momentum between the two quarks, $`\gamma _\mu `$ is the bare quark-gluon vertex, and $`\mathrm{\Gamma }_{\overline{3}}(q,P)`$ is the Bethe-Salpeter amplitude of the quasi-bound state (or diquark). The $`\xi `$ is the center of mass parameter which can be chosen to vary between $`0`$ and $`1`$, for the $`S^a`$ diquark, $`\xi =\frac{1}{2}`$, for the $`U^a`$ and $`D^a`$ diquarks, as the current quark masses $`m_s>m_u`$ and $`m_s>m_d`$, $`\xi `$ is about $`\frac{1}{2}`$. We can introduce an auxiliary amplitude $`\mathrm{\Gamma }_{\overline{3}}^C(q,P)`$ to facilitate the calculation,
$$\mathrm{\Gamma }_{\overline{3}}^C(q,P)\mathrm{\Gamma }_{\overline{3}}(q,P)C,$$
(8)
here $`C=\gamma _2\gamma _4`$ is the charge conjugation matrix. The auxiliary amplitude $`\mathrm{\Gamma }_{\overline{3}}^C(q,P)`$ satisfies the following equation,
$$\mathrm{\Gamma }_{\overline{3}}^C(q,P)=\frac{8\pi }{3}\frac{d^4k}{(2\pi )^4}G_{\mu \nu }(qk)\gamma _\mu S(k+\xi P)\mathrm{\Gamma }_{\overline{3}}^C(k,P)S(k(1\xi )P)\gamma _\nu .$$
(9)
It is obviously the above equation is identical to the BSE for the pseudoscalar mesons but a reduction in the coupling strength, $`\frac{4}{3}\frac{2}{3}`$. We can introduce the Bethe-Salpeter wavefunction (BSW) $`\chi _{qq}`$ for the quasi-bound states,
$$\chi _{qq}(q,P)S(q+\xi P)\mathrm{\Gamma }_{\overline{3}}^C(p,P)S(q(1\xi )P),$$
(10)
to relate with our previously works on the pseudoscalar meson’s BSE,
$`S^1(q+\xi P)\chi _{qq}(q,P)S^1(q(1\xi )P)={\displaystyle \frac{8\pi }{3}}{\displaystyle \frac{d^4k}{(2\pi )^4}\gamma _\mu \chi _{qq}(k,P)\gamma _\nu G_{\mu \nu }(qk)},`$ (11)
$`S^1(q+\xi P)\chi _{q\overline{q}}(q,P)S^1(q(1\xi )P)={\displaystyle \frac{16\pi }{3}}{\displaystyle \frac{d^4k}{(2\pi )^4}\gamma _\mu \chi _{q\overline{q}}(k,P)\gamma _\nu G_{\mu \nu }(qk)},`$ (12)
We can perform the Wick rotation analytically and continue $`q`$ and $`k`$ into the Euclidean region <sup>2</sup><sup>2</sup>2To avoid possible difficulties in performing the Wick rotation, one can derive the BSE from the Euclidean path-integral formulation of the theory. . The BSWs of the scalar diquarks $`\chi _{qq}`$ and pseudoscalar mesons $`\chi _{q\overline{q}}`$ have the same Dirac structures, can be signed by the notation $`\chi (q,P)`$. In the lowest order approximation, the BSW $`\chi (q,P)`$ can be written as
$`\chi (q,P)=\gamma _5\left[iF_1^0(q,P)+\gamma PF_2^0(q,P)+\gamma qqPF_3^1(q,P)+i[\gamma q,\gamma P]F_4^0(q,P)\right].`$ (13)
In solving the BSEs, it is important to translate the wavefunctions $`F_i^n`$ into the same dimension,
$`F_1^0\mathrm{\Lambda }^0F_1^0,F_2^0\mathrm{\Lambda }^1F_2^0,F_3^1\mathrm{\Lambda }^3F_3^1,F_4^0\mathrm{\Lambda }^2F_4^0,qq/\mathrm{\Lambda },PP/\mathrm{\Lambda },`$
here the $`\mathrm{\Lambda }`$ is some quantity of the dimension of mass. The ladder BSEs for the scalar diquarks and pseudoscalar mesons can be projected into the following four coupled integral equations,
$`{\displaystyle \underset{j}{}}H(i,j)F_j^{0,1}(q,P)`$ $`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle d^4kK(i,j)},`$ (14)
the expressions of the $`H(i,j)`$ and $`K(i,j)`$ are cumbersome and neglected here.
We can introduce a parameter $`\lambda (P^2)`$ and solve the above equations as an eigenvalue problem. If there really exist a quasi-bound state of two-quark, the masses of the diquarks can be determined by the condition $`\lambda (P^2=M_{qq}^2)=1`$,
$`{\displaystyle \underset{j}{}}H(i,j)F_j^{0,1}(q,P)`$ $`=`$ $`\lambda (P^2){\displaystyle \underset{j}{}}{\displaystyle d^4kK(i,j)}.`$ (15)
Here we will take a short digression and give some explanations for the expressions of $`H(i,j)`$ . The $`H(i,j)`$’s are functions of the quark’s Schwinger-Dyson functions (SDF) $`A(q^2+\xi ^2P^2+2\xi qP)`$ , $`B(q^2+\xi ^2P^2+2\xi qP)`$, $`A(q^2+(1\xi )^2P^22(1\xi )qP)`$ and $`B(q^2+(1\xi )^2P^22(1\xi )qP)`$ . The relative four-momentum $`q`$ is a quantity in the Euclidean spacetime while the center of mass four-momentum $`P`$ must be continued to the Minkowski spacetime i.e. $`P^2=m_{\pi ,K,U^a,D^a,S^a}^2`$ , this results in that the $`qP`$ varies throughout a complex domain. It is inconvenient to solve the SDE with the resulting complex values of the quark momentum. We can expand the $`A`$ and $`B`$ in terms of Taylor series of $`qP`$, for example,
$`A(q^2+\xi ^2P^2+\xi qP)`$ $`=`$ $`A(q^2+\xi ^2P^2)+2\xi A(q^2+\xi ^2P^2)^{}qP+\mathrm{}.`$
The other problem is that we can not solve the SDE in the timelike region as the two-point gluon Green’s function can not be exactly inferred from the $`SU(3)`$ color gauge theory even in the low energy spacelike region. In practical calculations, we can extrapolate the values of the $`A`$ and $`B`$ from the spacelike region smoothly to the timelike region with suitable polynomial functions. To avoid possible violation with confinement in sense of the appearance of pole masses $`q^2=m^2(q^2)`$ in the timelike region, we must be care in choosing the polynomial functions .
Finally we write down the normalization condition for the BSWs of the pseudoscalar mesons,
$`N_c{\displaystyle \frac{d^4q}{(2\pi )^4}Tr\left\{\overline{\chi }\frac{S_+^1}{P_\mu }\chi (q,P)S_{}^1+\overline{\chi }S_+^1\chi (q,P)\frac{S_{}^1}{P_\mu }\right\}}=2P_\mu ,`$ (16)
here $`\overline{\chi }=\gamma _4\chi ^+\gamma _4`$, $`S_+=S(q+\xi P)`$ and $`S_{}=S(q(1\xi )P)`$. In this article, the parameters of FBP are fitted to give the correct masses and decay constants for the pseudoscalar mesons, $`\pi `$ and $`K`$, the normalization condition is needed.
## 5 Coupled rainbow SDE and ladder BSE, and the mass spectrum
In this section, we study the coupled equations of the rainbow SDE and ladder BSE for the pseudoscalar mesons ($`\pi `$ and $`K`$) and scalar diquarks ($`U^a`$, $`D^a`$ and $`S^a`$) numerically, the final results for the SDFs and BSWs can be plotted as functions of the square momentum $`q^2`$.
In order to demonstrate the confinement of quarks, we have to study the analyticity of the SDFs of the $`u`$, $`d`$ and $`s`$ quarks, and prove that there are no mass poles on the real timelike $`q^2`$ axial. In the following, we take the Fourier transform with respect to the Euclidean time T for the scalar part ($`S_s`$) of the quark propagator ,
$`S_s^{}(T)={\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{dq_4}{2\pi }}e^{iq_4T}{\displaystyle \frac{B(q^2)}{q^2A^2(q^2)+B^2(q^2)}}|_{\stackrel{}{q}=0},`$ (17)
where the 3-vector part of $`q`$ is set to zero. If S(q) has a mass pole at $`q^2=m^2(q^2)`$ in the real timelike region, the Fourier transformed $`S_s^{}(T)`$ would fall off as $`e^{mT}`$ for large T or $`\mathrm{log}S_s^{}=mT`$. In our numerical calculations, for small $`T`$, the values of $`S_s^{}`$ are positive and decrease rapidly to zero and beyond with the increase of $`T`$, which are compatible with the result (curve tendency with respect to $`T`$) from lattice simulations ; for large $`T`$, the values of $`S_s^{}`$ are negative, except occasionally a very small fraction of positive values. The negative values for $`S_s^{}`$ indicate an explicit violation of the axiom of reflection positivity , in other words, the quarks are not physical observable i.e. confinement. As colored quantity, the diquarks should also be confined and not appear as asymptotic states; the demonstration of their confinement is beyond the present work.
The $`u`$, $`d`$ and $`s`$ quarks have small current masses, the dressing or renormalization is very large and the curves of the SDFs are steep, which are corresponding to the dynamical chiral symmetry breaking phenomenon for the light quarks. At zero momentum, $`m_u(0)=m_d(0)=0.454GeV`$ and $`m_s(0)=0.684GeV`$, which are compatible with the constituent quark masses in the literature. From the solutions of BSEs for the $`\pi `$, $`K`$ mesons and $`U^a`$, $`D^a`$, $`S^a`$ diquarks as eigenvalue problems, we can obtain the masses for those pseudoscalar mesons and scalar diquarks,
$`M_\pi =135MeV,M_K=498MeV,`$
$`M_{S^a}=0.76GeV,M_{U^a}=0.98GeV,M_{D^a}=0.98MeV.`$ (18)
It is obviously
$`M_{S^a}<m_u(0)+m_d(0),M_{U^a}<m_d(0)+m_s(0),M_{D^a}<m_u(0)+m_s(0),`$ (19)
$`M_{U^a}M_{S^a}m_s(0)m_u(0)0.22GeV.`$ (20)
The attractive interaction between the quarks in the color and flavor $`SU(3)`$ $`\overline{3}`$ channel can lead to the quasi-bound states in the infrared region. The appearance of the diquarks is closely related to the dynamical symmetry breaking phenomenon, the mass splitting among the $`U^a`$, $`D^a`$ and $`S^a`$ diquarks originate from the mass splitting among $`u`$, $`d`$ and $`s`$ quarks.
The existing theoretical calculations for the masses of the scalar diquarks vary in a large range, $`M_{qq}=(0.40.7)GeV`$ with QCD sum rules ; $`M_{qq}0.5GeV`$ with random instanton liquid model ; $`M_{qq}=(0.42\pm 0.03)GeV`$ with random instanton liquid model , $`M_{qq}=0.234GeV`$ with Nambu-Jona-Lasinio Model ; $`M_{S^a}=0.74GeV,M_{U^a}=M_{D^a}=0.88GeV`$ with BSE ; $`M_{S^a}=0.82GeV,M_{U^a}=M_{D^a}=1.10GeV`$ with BSE ; $`M_{qq}=0.692GeV`$ with global color model ; $`0.14<M_{qq}<0.74GeV`$ with assumption of mass functions ; $`M_{qq}0.7GeV`$ with lattice QCD . There are large uncertainties for the masses of the quasi-bound states, $`U^a`$, $`D^a`$, and $`S^a`$. As colored quantities, the diquarks may have gauge interactions with the gluon field as fundamental scalar field with the following Lagrangian,
$`L`$ $`=`$ $`{\displaystyle \frac{1}{2}}D_\mu S_{}^{a}{}_{}{}^{+}D_\mu S^a{\displaystyle \frac{1}{2}}M_{S^a}^2S_{}^{a}{}_{}{}^{+}S^a,`$ (21)
$`D_\mu S^a`$ $`=`$ $`_\mu S^a+igf^{abc}A_\mu ^bS^c.`$ (22)
Here we use the notation $`S^a`$ to represent the scalar diquark field and $`A_\mu ^b`$ the gluon field. The direct color interactions between the scalar diquarks and the gluons may modify the mass $`M_{qq}`$ significantly.
From the plotted BSWs (see Fig.1 for the $`\pi `$ meson and Fig.2 for the $`S^a`$ diquark as examples), we can see that the BSWs of the pseudoscalar mesons and scalar diquarks have the same type (Gaussian type) momentum dependence while the quantitative values are different from each other. Just like the $`\overline{q}q`$ , $`\overline{q}Q`$ and $`\overline{Q}Q`$ pseudoscalar mesons , the gaussian type BSWs of the scalar diquarks center around zero momentum and extend to the energy scale about $`q^2=1GeV^2`$ which happen to be the energy scale for the chiral symmetry breaking, the strong interactions in the infrared region result in quasi-bound states, $`U^a`$, $`D^a`$ and $`S^a`$. The BSWs of the $`\pi `$ and $`K`$ mesons can give satisfactory values for the decay constants which are defined by
$`if_\pi P_\mu `$ $`=`$ $`0|\overline{q}\gamma _\mu \gamma _5q|\pi (P),`$ (23)
$`=`$ $`N_c{\displaystyle Tr\left[\gamma _\mu \gamma _5\chi (k,P)\right]\frac{d^4k}{(2\pi )^4}},`$
here we use $`\pi `$ to represent the pseudoscalar mesons,
$`f_\pi =127MeV;f_K=161MeV.`$ (24)
There are negative voice for the existence of the scalar diquark states in the infrared region . The coupled rainbow SDE and ladder BSE are particularly suitable for studying the flavor octet pseudoscalar mesons and vector mesons, the next-to-leading-order (NLO) contributions from the quark-gluon vertex have a significant amount of cancellation between repulsive and attractive corrections. However, for the diquarks and scalar mesons, the large repulsive corrections from the NLO contributions can significantly change the scalar meson masses and corresponding BSWs; the diquark quasi-bound states found in the ladder BSE will disappear from the spectrum. In this article, we do not take the point of view that the scalar diquarks exist in the strong interaction spectrum as asymptotic states, and our conclusion do not conflict with Ref.; the confinement precludes the observation of the free colored diquarks, the quark-quark can correlate with each other in the color and flavor $`\overline{3}`$ channels inside the baryons and multi-quark states with typical length $`l=\frac{1}{M_{qq}}`$. In fact, the NLO contributions from the quark-gluon vertex are extremely difficult to take into account if we go beyond the infrared dominated $`\delta `$ function approximation for the gluon kernel, which is obviously violate the lorentz invariance; the lorentz invariant and model-independent treatments still lack in the literatures.
In calculation, the values of current quark masses are taken as $`\widehat{m}_u=\widehat{m}_d=6MeV`$ and $`\widehat{m}_s=150MeV`$; the input parameters for the FBP are $`N=1.0\mathrm{\Lambda }`$, $`V(0)=17.0\mathrm{\Lambda }`$, $`\rho =6.0\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }=200MeV`$, which are fitted to give the correct masses of the $`\pi `$ and $`K`$ mesons. In this article, we deal with only the light flavor quarks, the FBP can give satisfactory results without fine tuning, such as modifying the infrared behavior for the heavy quarks $`c`$ and $`b`$.
## 6 Conclusion
In this article, we study the under-structures of the pseudoscalar mesons $`\pi `$ , $`K`$ and scalar diquarks $`U^a`$, $`D^a`$, $`S^a`$ in the framework of the coupled rainbow SDE and ladder BSE with the confining effective potential (FBP). After we solve the coupled rainbow SDE and ladder BSE numerically, we obtain the SDFs and BSWs of the pseudoscalar mesons $`\pi `$ , $`K`$ and scalar diquarks $`U^a`$, $`D^a`$, $`S^a`$, and the corresponding ground state mass spectrum. The $`u`$, $`d`$ and $`s`$ quarks have small current masses, the dressing or renormalization for the SDFs is very large and the curves are steep which indicate the dynamical chiral symmetry breaking phenomenon for the light quarks explicitly. The mass poles in the timelike region are absent which implement the confinement naturally. The BSWs of the pseudoscalar mesons and scalar diquarks have the same type (Gaussian type) momentum dependence while the quantitative values are different from each other. The gaussian type BSWs center around zero momentum and extend to the energy scale about $`q^2=1GeV^2`$ which happen to be the energy scale for the chiral symmetry breaking, the strong interactions in the infrared region result in bound (or quasi-bound) states. Our numerical results for the masses and decay constants of the $`\pi `$, $`K`$ mesons can reproduce the experimental values, the ground state mass spectrum of scalar diquarks are consistent with the existing theoretical calculations. The mass splitting among the $`U^a`$, $`D^a`$ and $`S^a`$ diquarks originate from the mass splitting among $`u`$, $`d`$ and $`s`$ quarks. $`M_{S^a}<m_u(0)+m_d(0)`$, $`M_{U^a}<m_d(0)+m_s(0)`$ , $`M_{D^a}<m_u(0)+m_s(0)`$ , the attractive interaction between the color and flavor $`SU(3)`$ $`\overline{3}`$ quarks can lead to the quasi-bound states in the infrared region. Once the satisfactory SDFs and BSWs of the scalar diquarks are known, we can use them to study a lot of important quantities involving the multiquark states.
## Acknowledgment
This work is supported by National Natural Science Foundation, Grant Number 10405009, and Key Program Foundation of NCEPU. The authors are indebted to Dr. J.He (IHEP), Dr. X.B.Huang (PKU) and Dr. L.Li (GSCAS) for numerous help, without them, the work would not be finished. The authors would also like to thank Prof. C. D. Roberts for providing us some important literatures.
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# Quantum state tomography with quantum shotnoise
## Abstract
We propose a scheme for a complete reconstruction of one- and two-particle orbital quantum states in mesoscopic conductors. The conductor in the transport state continuously emits orbital quantum states. The orbital states are manipulated by electronic beamsplitters and detected by measurements of average currents and zero frequency current shotnoise correlators. We show how, by a suitable complete set of measurements, the elements of the density matrices of the one- and two-particle states can be directly expressed in terms of the currents and current correlators.
According to the standard interpretation of quantum mechanics the wavefunction, or more generally the density matrix, determines the probabilities for the possible outcomes of any measurement on the quantum state. To completely characterize the wavefunction of the state is therefore of fundamental interest Pauli . It is however impossible to infer anything about an unknown state from a single measurement, a complete characterization requires an ensemble of identically prepared states and the measurement of a complete set of observables on the state Fano . A reconstruction of the quantum state wavefunction via such a series of measurements is known as Quantum State Tomography (QST) QSTrevs .
Initially, QST was performed experimentally on the discrete angular momentum state of an electron in an hydrogen atom Ashburn . During the last decade QST has been performed on e.g. the quantum state of squeezed light Smithey , the vibrational state of a molecule Dunn , the motional state of trapped ions Liebfried and of atomic wavepackets Kurtsiefer . Recently there has been an interest in QST of two-particle states in the context of quantum information processing. The entanglement of a quantum state, a potential resource for quantum information processing, is characterized by the density matrix of the state. The quantum state of polarization entangled pairs of photons has been reconstructed using QST EntQST .
To date, no QST has been performed on quantum states in solid state systems. Very recently a theoretical scheme Nori was developed for solid state two-levels systems, qubits, appropriate for e.g. the macroscopic superposition state in superconducting qubits and the spin-state of electrons in quantum dots. The set of measurements necessary to reconstruct the state involves controlled rotations and detection of the individual qubits. For coupled qubits, where entanglement between the qubits is of interest, such measurements are highly involved and have not been demonstrated.
In this paper we take a different approach and present a scheme for QST of discrete single and two-particle orbital quantum states in mesoscopic conductors. The orbital quantum states Orb ; QH1 ; QH2 are continuously emitted from the conductor during transport, making a long time measurement equivalent to an average over an ensemble of states. The orbital states can be manipulated by electronic beamsplitters, experimentally available BS1 ; BS2 , and detected by measurements of average currents and zero frequency current correlators, shotnoise Buttnoise ; Buttrev . This scheme, with all components experimentally realizable in e.g. Quantum Hall systems BS2 ; QHexp , allows for a complete characterization of the quasiparticle quantum state in mesoscopic conductors.
The key question for any QST is: what quantum states with interesting properties can be investigated with accessible experimental technics? In mesoscopic conductors, one typically measures electrical currents and current correlators. In several recent works Sukh ; Cht ; Orb ; QH1 ; QH2 , it has been shown theoretically that quantum correlations, entanglement, between two spatially separated particles can be investigated via current correlation measurements. In particular, in Refs. Orb ; QH2 it was shown how entangled orbital quasiparticle states could be generated, manipulated with experimentally available electronic BS1 ; BS2 beamsplitters and detected via current correlation measurements. Orbital one- and two-particle states investigated via current and current correlations are thus natural candidates for QST in mesoscopic conductors.
A generic setup for such orbital QST is shown in Fig. 1. A mesoscopic conductor $`S`$ acts as a source for orbital quantum states. The source is connected via four single mode leads $`A1,A2,B1`$ and $`B2`$ to two regions, $`A`$ and $`B`$, where the emitted state is manipulated and detected. The mesoscopic source has one or more reservoirs biased at $`eV`$ and an arbitrary number of reservoirs kept at ground. We note that two-particles effects are only present for two or more reservoirs biased Buttnoise . The temperature is taken to be zero. It is assumed that the scattering in the conductor is elastic, however arbitrary debasing inside the conductor can be accounted for.
The regions $`A`$ and $`B`$ each contain an electronic beamsplitter BS1 ; BS2 and an electrostatic sidegate (see e.g. QHexp ) to induce a phaseshift, $`\varphi _A`$ or $`\varphi _B`$, by modifying the length of the lead. The beamsplitters, taken to be reflectionless, are further connected to two grounded reservoirs $`+`$ and $``$ where the current is measured. The combined beamsplitter-sidegate structure can be characterized by a scattering matrix, for e.g. $`A`$ given by
$`S_A=\left(\begin{array}{cc}\sqrt{R_A}e^{i\phi _{A2}}& \sqrt{T_A}e^{i(\phi _{A3}\varphi _A)}\\ \sqrt{T_A}e^{i(\phi _{A1}+\phi _{A2})}& \sqrt{R_A}e^{i(\phi _{A1}+\phi _{A3}\varphi _A)}\end{array}\right).`$ (3)
The transmission probability $`T_A=1R_A`$ can be controlled via electrostatic gating BS1 ; BS2 ; QHexp . The phases $`\phi _{Ai},i=1,2,3`$ picked up when scattering at the beamsplitter are however assumed to be uncontrollable but fixed during the measurement.
The quantum state emitted by the mesoscopic source is in the general case a manybody state, it is a linear superposition of states with different number of particles SSQI . However, one- and two-particle observables such as current and noise are only sensitive to the one- and two-particle properties of the state. These properties are quantified by the reduced density matrix, which thus is the object of interest. Only in some special cases, typically in conductors in the tunneling limit Orb ; QH1 ; QH2 , are the emitted states true one or two-particle states. In the presence of dephasing, the emitted state is mixed. Moreover, even an emitted pure manybody state generally gives rise to a mixed reduced one- or two particle state. It is therefore appropriate to discuss the state in terms of density matrices.
To simplify the discussion we consider a spin-polarized system with scattering amplitudes independent on energy on the scale of the applied bias $`eV`$, i.e. the linear voltage regime. The emitted state then has only orbital degrees of freedom. We first consider the single-particle orbital state emitted e.g. towards $`A`$ (same considerations hold for $`B`$). Introducing operators $`b_{An}^{}`$ creating electrons in lead $`An`$, with $`n=1,2`$, propagating out from the source, the $`2\times 2`$ density matrix (not normalized) is by definition given by
$$\rho _A=\underset{n,m=1}{\overset{2}{}}\rho _{nm}b_{An}^{}|00|b_{Am}=\left(\begin{array}{cc}\rho _{11}& \rho _{12}\\ \rho _{21}& \rho _{22}\end{array}\right)$$
(4)
where we work in the basis $`\{|1_A,|2_A\}`$, with $`b_{An}^{}|0=|n_A`$, formed by the lead indices (see Fig. 1). The matrix elements $`\rho _{nm}=b_{Am}^{}b_{An}`$. The Hermitian density matrix, $`\rho _A=\rho _A^{}`$, has four independent parameters and can be written as follows
$$\rho _A=\frac{1}{2}\underset{i=0}{\overset{3}{}}c_i\sigma _i=\frac{1}{2}\left(\begin{array}{cc}c_0+c_3& c_1ic_2\\ c_1+ic_2& c_0c_3\end{array}\right),$$
(5)
where $`\{\sigma _i\}=[\mathrm{𝟏},\sigma _x,\sigma _y,\sigma _z]`$. A normalized density matrix is obtained by dividing all elements by $`c_0`$. In the same way, the two-particle density matrix is given by
$$\rho _{AB}=\underset{n,m,k,l=1}{\overset{2}{}}\rho _{nm}^{kl}b_{An}^{}b_{Bk}^{}|00|b_{Bl}b_{Am}$$
(6)
with the matrix elements $`\rho _{nm}^{kl}=b_{Am}^{}b_{Bl}^{}b_{Bk}b_{An}`$. The two-particle density matrix has 16 independent parameters and can be written
$`\rho _{AB}={\displaystyle \frac{1}{4}}{\displaystyle \underset{i,j=0}{\overset{3}{}}}c_{ij}\sigma _i\sigma _j`$ (7)
with $``$ the direct product. Expressing the real coefficients $`c_i`$ and $`c_{ij}`$ in terms of outcomes of ensemble averaged measurements thus gives a complete reconstruction of the emitted state. The accessible measurements are average current and zero frequency current correlations. Importantly, in the transport state the source continuously emits quantum states. As a consequence, the long time measurements automatically provide an ensemble average measurement. At $`A`$ the average currents at contacts $`\alpha =\pm `$ are
$`I_A^\alpha ={\displaystyle \frac{e^2V}{h}}n_A^\alpha ,n_A^\alpha =b_{A\alpha }^{}b_{A\alpha }.`$ (8)
The zero frequency correlator between currents fluctuations $`\mathrm{\Delta }I`$ in reservoirs $`A\alpha `$ and $`B\beta `$, given by $`S_{AB}^{\alpha \beta }=(1/2)𝑑t\mathrm{\Delta }I_{A\alpha }(0)\mathrm{\Delta }I_{B\beta }(t)+\mathrm{\Delta }I_{B\beta }(t)\mathrm{\Delta }I_{A\alpha }(0)`$ can be written Buttnoise
$`S_{AB}^{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{2e^3V}{h}}\left[n_A^\alpha n_B^\beta n_A^\alpha n_B^\beta \right]`$ (9)
with $`n_B^\alpha =b_{B\alpha }^{}b_{B\alpha }`$. The operators $`b_{A\alpha }`$ and $`b_{B\beta }`$ in the reservoirs at $`A`$ and $`B`$ are related to operators $`b_{An}`$ and $`b_{Bk}`$ in the leads $`An,Bk`$, with $`n,k=1,2`$, (see Fig. 1) via the scattering matrix $`S_A`$ of the beamsplitters at $`A`$ as
$`\left(\begin{array}{c}b_{A+}\\ b_A\end{array}\right)=S_A\left(\begin{array}{c}b_{A1}\\ b_{A2}\end{array}\right)`$ (14)
and similarly at $`B`$.
We start with the reconstruction of the one-particle state at $`A`$, accessible via the average current (the same procedure holds for the state at $`B`$). Here a formal approach is taken which directly can be extended to the investigation of the two-particle state. We note that the reconstruction approach is similar to QST schemes for qubits in other systems, see e.g. Refs. Kwiat2 ; Nori . There are however a number of important special features for mesoscopic systems, making a detailed investigation important. It is desirable to minimize both the type and number of experiments having to be carried out. As is clear from the following, for a complete reconstruction it is sufficient to consider only measurements of currents in one reservoir in $`A`$. Here we consider the current at $`A+`$. Using the relation between operators, Eq. (14), and Eq. (8), we have
$`I_A^+/(e^2V/h)=n_A^+=\text{tr}\left(\rho _A𝒜\right)`$ (15)
with the matrix
$`𝒜=\left(\begin{array}{cc}R_A& \sqrt{T_AR_A}e^{i(\varphi _A+\phi _A)}\\ \sqrt{T_AR_A}e^{i(\varphi _A+\phi _A)}& T_A\end{array}\right).`$ (18)
The phase $`\phi _A=\varphi _{A2}\varphi _{A3}`$ contains all the information about uncontrollable phases of the beamsplitter. From Eqs. (15) and (18) it is clear that the phase $`\phi _A`$ can be included in $`\rho _A`$ by a change of local basis $`\rho _AU_A\overline{\rho }_AU_A^{}`$ with $`U_A=\text{diag}[\text{exp}(i\phi _A/2),\text{exp}(i\phi _A/2)]`$. Below we consider the reconstruction of $`\overline{\rho }_A`$, parameterized by coefficients $`\overline{c}_i`$ \[see Eq. (5)\], thus working with $`𝒜(\phi _A=0)`$ in Eq. (18). This yields $`\rho _A`$ up to an unknown local basis rotation.
Importantly, only four settings of the beamsplitter are needed, both for the current and the current correlators, for a complete state reconstruction. The settings $`I`$ to $`IV`$ considered here are listed in the table in Fig. 2.
By constructing suitable linear combinations $`j_A(j)`$ of the observables at the different settings, in the $`\{|1_A,|2_A\}`$ basis
$`j_A(0)`$ $`=`$ $`𝒜(I)+𝒜(II)=\mathrm{𝟏},`$
$`j_A(1)`$ $`=`$ $`2𝒜(III)\left[𝒜(I)+𝒜(II)\right]=\sigma _x,`$
$`j_A(2)`$ $`=`$ $`2𝒜(IV)\left[𝒜(I)+𝒜(II)\right]=\sigma _y,`$
$`j_A(3)`$ $`=`$ $`𝒜(I)𝒜(II)=\sigma _z,`$ (19)
we obtain a complete set Fano of measurements, since the Pauli matrices $`\sigma _j`$ obey the relation $`\text{tr}(\sigma _i\sigma _j)=2\delta _{ij}`$. Here $`𝒜(I)`$ is the matrix $`𝒜`$ in Eq. (18) for the setting I etc. From Eqs. (15) and (19) and the relation $`j_A(j)=\text{tr}(\overline{\rho }_A\sigma _j=\overline{c}_j`$ we then directly obtain the coefficients $`\overline{c}_j`$, parametrizing $`\rho _A`$ in Eq. (5)
$`\overline{c}_j`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{3}{}}}Q_{jk}n_A^+(k),Q=\left(\begin{array}{cccc}1& 1& 0& 0\\ 1& 1& 2& 0\\ 1& 1& 0& 2\\ 1& 1& 0& 0\end{array}\right)`$ (24)
in terms of the measured currents for the different settings, taking the index $`\{k\}=[0,1,2,3][I,II,II,IV]`$. This completes the one-particle state reconstruction.
We then turn to the two-particle state. In Eq. (9), the quantity that is directly linked to the density matrix elements $`\rho _{nm}^{kl}`$ is the reducible correlator $`n_A^\alpha n_B^\beta `$. This correlator is directly obtained from the measured noise and the average currents. In analogy to the current, it is sufficient to consider correlations between currents in one terminal in $`A`$ and one in $`B`$. Considering here $`A+`$ and $`B+`$, one obtains from Eq. (9) and (14) the dimensionless correlator
$`{\displaystyle \frac{S_{AB}^{++}}{2e^3V/h}}+{\displaystyle \frac{I_A^+I_B^+}{(e^2V/h)^2}}=n_A^+n_B^+=\text{tr}\left(\rho _{AB}𝒜\right)`$ (25)
where the matrix $``$ is given from $`𝒜`$ in Eq. (18) by changing indices $`AB`$ in the scattering amplitudes. Similar to the one-particle state, we note from Eq. (25) that both phases $`\phi _A`$ and $`\phi _B`$ can be included in $`\rho _{AB}`$ by independent local rotations $`\rho _{AB}(U_AU_B)\overline{\rho }_{AB}(U_AU_B)^{}`$, with $`U_B=\text{diag}[\text{exp}(i\phi _B/2),\text{exp}(i\phi _B/2)]`$. Below we thus consider the reconstruction of $`\overline{\rho }_{AB}`$, parameterized as in Eq. (7) by the coefficients $`\overline{c}_{ij}`$, yielding $`\rho _{AB}`$ up to a local basis rotation entcom .
By considering the same four settings at $`B`$ as at $`A`$, we can use the linear combination operators $`j_A(j)`$ in Eq. (19) and correspondingly $`j_B(i)`$ to construct a complete set of observables, in the basis $`\{|1_A|1_B,|1_A|2_B,|2_A|1_B,|2_A|2_B\}`$,
$`j_A(j)j_B(i)=\sigma _j\sigma _i`$ (26)
since the direct products of $`\sigma `$-matrices obey $`\text{tr}\left[(\sigma _j\sigma _i)(\sigma _k\sigma _l)\right]=4\delta _{jk}\delta _{il}`$. From Eq. (25) and the relation $`j_A(j)j_B(i)=\text{tr}(\overline{\rho }_{AB}\sigma _j\sigma _i=\overline{c}_{ji}`$ we then directly obtain the coefficients $`\overline{c}_{ji}`$ as
$`\overline{c}_{ji}={\displaystyle \underset{k,l=0}{\overset{3}{}}}Q_{jk}Q_{il}n_A^+(k)n_B^+(l)`$ (27)
in terms of the measured current correlators and averaged currents. We emphasize that all elements can be determined from sixteen current correlations and eight average currents (four at $`A`$ and four at $`B`$). We also note that the reconstructed density matrix, due to nonideal measurements, might have negative eigenvalues, i.e. it might not be positive semidefinite. Schemes to correct for this for one and two-qubit states are discussed in e.g. ref. Kwiat2 .
In the context of two-particle entanglement, it is interesting to compare the QST-scheme with a Bell Inequality, recently discussed for mesoscopic system (see e.g. refs. Cht ; Orb and for a density matrix approach ref. Carlodeph ). Both schemes require the same number of current correlation measurements. The density matrix reconstructed by QST however completely determines the entanglement. In contrast, a Bell Inequality can not be used to quantify the entanglement Verstraete , there are e.g. mixed entangled states Werner that do not lead to a violation of a Bell Inequality.
It is clarifying to illustrate the above scheme with a simple example (see Fig. 2). We consider the Hanbury Brown Twiss geometry of ref. QH2 , where the number of nonzero elements of the one and two-particle density matrices are reduced due to the topological properties of the conductor. Since no scattering between the upper, $`1`$, and lower, $`2`$, leads is physically possible due to the spatial separation, the one-particle density matrix $`\overline{\rho }_A`$ only has two nonzero elements, $`\overline{\rho }_{11}`$ and $`\overline{\rho }_{22}`$. These elements are parametrized by $`\overline{c}_0`$ and $`\overline{c}_3`$, obtained by measuring $`j_A(0)`$ and $`j_A(3)`$.
The two-particle density matrix $`\overline{\rho }_{AB}`$ has four nonzero elements, $`\overline{\rho }_{11}^{22},\overline{\rho }_{22}^{11},\overline{\rho }_{21}^{12}`$ and $`\overline{\rho }_{12}^{21}`$. Using the relation between the coefficients $`\overline{c}_{ij}`$ resulting from several matrix elements being zero, $`\overline{\rho }_{AB}`$ can then be parameterized as
$$\overline{\rho }_{AB}=\frac{1}{2}\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& \overline{c}_{00}+\overline{c}_{33}& \overline{c}_{11}i\overline{c}_{21}& 0\\ 0& \overline{c}_{11}+i\overline{c}_{21}& \overline{c}_{00}\overline{c}_{33}& 0\\ 0& 0& 0& 0\end{array}\right).$$
(28)
Conseqently, only four correlations $`j_A(i)j_B(j)`$ need to be measured to completely reconstruct $`\overline{\rho }_{AB}`$, reducing the number of actual current correlations needed to be carried out to 12 for the settings considered here \[see Eq. (27)\]. It is interesting to note that in the geometry in Fig. 2, considering the tunneling limit for the beamsplitters in the source and changing to an electron-hole picture QH1 , the emitted state is a true two-particle state QH2 . Since the hole currents and current fluctuations are directly related to the electron ones, it is possible to employ our scheme to reconstruct an electron-hole state as well.
In conclusion, we have presented a scheme for orbital Quantum State Tomography, a complete reconstruction of orbital quantum states in mesoscopic conductors in the transport state. The emitted orbital states are manipulated with electronic beamsplitters and detected with currents and zero frequency current correlations. With all components experimentally available, our scheme opens up for a direct observation of quasiparticle quantum states in mesoscopic conductors.
We acknowledge discussions with E.V. Sukhorukov. This work was supported by the Swedish Research Council, the Swiss National Science Foundation and the program for Materials with Novel Electronic Properties.
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# 1 Introduction
## 1 Introduction
Two theories of gravity called the nonsymmetric gravity theory (NGT) and the metric-skew-tensor gravity (MSTG) theory have been proposed to explain the rotational velocity curves of galaxies, clusters of galaxies and cosmology without dark matter. A fitting routine for galaxy rotation curves has been used to fit a large number of galaxy rotational velocity curve data, including low surface brightness (LSB), high surface brightness (HSB) and dwarf galaxies with both photometric data and a two-parameter core model without non-baryonic dark matter . The fits to the data are remarkably good and for the photometric data only the one parameter, the mass-to-light ratio $`M/L`$, is used for the fitting, once two parameters $`M_0`$ and $`r_0`$ are universally fixed for galaxies and dwarf galaxies. The fits are close to those obtained from Milgrom’s MOND acceleration law in all cases considered . A large sample of X-ray mass profile cluster data has also been fitted .
The gravity theories require that Newton’s constant G, the coupling constant $`\gamma _c`$ that measures the strength of the coupling of the skew field to matter and the mass $`\mu `$ of the skew field, vary with distance and time, so that agreement with the solar system and the binary pulsar PSR 1913+16 data can be achieved, as well as fits to galaxy rotation curve data and galaxy cluster data. In ref. , the variation of these constants was based on a renormalization group (RG) flow description of quantum gravity theory formulated in terms of an effective classical action . Large infrared renormalization effects can cause the effective $`G`$, $`\gamma _c`$, $`\mu `$ and the cosmological constant $`\mathrm{\Lambda }`$ to run with momentum $`k`$ and a cutoff procedure leads to a space and time varying $`G`$, $`\gamma _c`$ and $`\mu `$, where $`\mu =1/r_0`$ and $`r_0`$ is the effective range of the skew symmetric field.
In the following, we shall pursue an alternative relativistic gravity theory based on scalar-tensor-vector gravity (STVG), in which $`G`$, a vector field coupling constant $`\omega `$ and the mass $`\mu `$ of the vector field are dynamical scalar fields that allow for an effective description of the variation of these “constants” with space and time. We shall not presently consider the variation of the cosmological constant $`\mathrm{\Lambda }`$ with space and time.
The gravity theory leads to the same modified acceleration law obtained from NGT and MSTG for weak gravitational fields and the same fits to galaxy rotation curve and galaxy cluster data, as well as to agreement with the solar system and pulsar PSR 1913+16 observations. An important constraint on gravity theories is the bounds obtained from weak equivalence principle tests and the existence of a “fifth” force, due to the exchange of a massive vector boson . These bounds are only useful for distances $`100A.U.`$ and they cannot rule out gravity theories that violate the weak equivalence principle or contain a fifth force at galactic and cosmological distance scales. Since the variation of $`G`$ in our modified gravity theory leads to consistency with solar system data, then we can explore the consequences of our STVG theory without violating any known local observational constraints.
An important feature of the NGT, MSTG and STVG theories is that the modified acceleration law for weak gravitational fields has a repulsive Yukawa force added to the Newtonian acceleration law. This corresponds to the exchange of a massive spin 1 boson, whose effective mass and coupling to matter can vary with distance scale. A scalar component added to the Newtonian force law would correspond to an attractive Yukawa force and the exchange of a spin 0 particle. The latter acceleration law cannot lead to a satisfactory fit to galaxy rotation curves and galaxy cluster data.
In Section 8, we investigate a cosmological solution based on a homogeneous and isotropic Friedmann-Lemaître-Robertson-Walker (FLRW) spacetime. We present a solution that can possibly fit the acoustic peaks in the CMB power spectrum by avoiding significant suppression of the baryon perturbations , and which can possibly be made to fit the recent combined satellite data for the power spectrum without non-baryonic dark matter.
All the current applications of the three gravity theories that can be directly confronted with experiment are based on weak gravitational fields. To distinguish the theories, it will be necessary to obtain experimental data for strong gravitational fields e.g. black holes. Moreover, confronting the theories with cosmological data may also allow a falsification of the gravity theories. Recently, the NGT and MSTG were studied to derive quantum fluctuations in the early universe from an inflationary-type scenario .
## 2 Action and Field Equations
Our action takes the form
$$S=S_{\mathrm{Grav}}+S_\varphi +S_S+S_M,$$
(1)
where
$$S_{\mathrm{Grav}}=\frac{1}{16\pi }d^4x\sqrt{g}\left[\frac{1}{G}(R+2\mathrm{\Lambda })\right],$$
(2)
$$S_\varphi =d^4x\sqrt{g}\left[\omega \left(\frac{1}{4}B^{\mu \nu }B_{\mu \nu }+V(\varphi )\right)\right],$$
(3)
and
$$S_S=d^4x\sqrt{g}[\frac{1}{G^3}(\frac{1}{2}g^{\mu \nu }_\mu G_\nu GV(G))$$
$$+\frac{1}{G}(\frac{1}{2}g^{\mu \nu }_\mu \omega _\nu \omega V(\omega ))+\frac{1}{\mu ^2G}(\frac{1}{2}g^{\mu \nu }_\mu \mu _\nu \mu V(\mu ))].$$
(4)
Here, we have chosen units with $`c=1`$, $`_\mu `$ denotes the covariant derivative with respect to the metric $`g_{\mu \nu }`$. We adopt the metric signature $`\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$ where $`\eta _{\mu \nu }`$ is the Minkowski spacetime metric. We have
$$R_{\mu \nu }=_\lambda \mathrm{\Gamma }_{}^{\lambda }{}_{\mu \nu }{}^{}_\nu \mathrm{\Gamma }_{}^{\lambda }{}_{\mu \lambda }{}^{}+\mathrm{\Gamma }_{}^{\lambda }{}_{\mu \nu }{}^{}\mathrm{\Gamma }_{}^{\sigma }{}_{\lambda \sigma }{}^{}\mathrm{\Gamma }_{}^{\sigma }{}_{\mu \lambda }{}^{}\mathrm{\Gamma }_{}^{\lambda }{}_{\nu \sigma }{}^{},$$
(5)
where $`\mathrm{\Gamma }_{\mu \nu }^\lambda `$ denotes the Christoffel connection:
$$\mathrm{\Gamma }_{\mu \nu }^\lambda =\frac{1}{2}g^{\lambda \sigma }(_\nu g_{\mu \sigma }+_\mu g_{\nu \sigma }_\sigma g_{\mu \nu }),$$
(6)
and $`R=g^{\mu \nu }R_{\mu \nu }`$. Moreover, $`V(\varphi )`$ denotes a potential for the vector field $`\varphi ^\mu `$, while $`V(G),V(\omega )`$ and $`V(\mu )`$ denote the three potentials associated with the three scalar fields $`G(x),\omega (x)`$ and $`\mu (x)`$, respectively. The field $`\omega (x)`$ is dimensionless and $`\mathrm{\Lambda }`$ denotes the cosmological constant. Moreover,
$$B_{\mu \nu }=_\mu \varphi _\nu _\nu \varphi _\mu .$$
(7)
The total energy-momentum tensor is given by
$$T_{\mu \nu }=T_{M\mu \nu }+T_{\varphi \mu \nu }+T_{S\mu \nu },$$
(8)
where $`T_{M\mu \nu }`$ and $`T_{\varphi \mu \nu }`$ denote the ordinary matter energy-momentum tensor and the energy-momentum tensor contribution of the $`\varphi _\mu `$ field, respectively, while $`T_{S\mu \nu }`$ denotes the scalar $`G`$, $`\omega `$ and $`\mu `$ contributions to the energy-momentum tensor. We have
$$\frac{2}{\sqrt{g}}\frac{\delta S_M}{\delta g^{\mu \nu }}=T_{M\mu \nu },\frac{2}{\sqrt{g}}\frac{\delta S_\varphi }{\delta g^{\mu \nu }}=T_{\varphi \mu \nu },\frac{2}{\sqrt{g}}\frac{\delta S_S}{\delta g^{\mu \nu }}=T_{S\mu \nu }.$$
(9)
The matter current density $`J^\mu `$ is defined in terms of the matter action $`S_M`$:
$$\frac{1}{\sqrt{g}}\frac{\delta S_M}{\delta \varphi _\mu }=J^\mu .$$
(10)
We obtain from the variation of $`g^{\mu \nu }`$, the field equations
$$G_{\mu \nu }g_{\mu \nu }\mathrm{\Lambda }+Q_{\mu \nu }=8\pi GT_{\mu \nu },$$
(11)
where $`G_{\mu \nu }=R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R`$. We have
$$Q_{\mu \nu }=G(^\alpha _\alpha \mathrm{\Theta }g_{\mu \nu }_\mu _\nu \mathrm{\Theta }),$$
(12)
where $`\mathrm{\Theta }(x)=1/G(x)`$. The quantity $`Q_{\mu \nu }`$ results from a boundary contribution arising from the presence of second derivatives of the metric tensor in $`R`$ in $`S_{\mathrm{Grav}}`$. These boundary contributions are equivalent to those that occur in Brans-Dicke gravity theory . We also have
$$T_{\varphi \mu \nu }=\omega \left[B_{\mu }^{}{}_{}{}^{\alpha }B_{\nu \alpha }g_{\mu \nu }\left(\frac{1}{4}B^{\rho \sigma }B_{\rho \sigma }+V(\varphi )\right)+2\frac{V(\varphi )}{g^{\mu \nu }}\right].$$
(13)
The $`G(x)`$ field yields the energy-momentum tensor:
$$T_{G\mu \nu }=\frac{1}{G^3}\left[_\mu G_\nu G2\frac{V(G)}{g^{\mu \nu }}g_{\mu \nu }\left(\frac{1}{2}_\alpha G^\alpha GV(G)\right)\right].$$
(14)
Similar expressions can be obtained for $`T_{\omega \mu \nu }`$ and $`T_{\mu \mu \nu }`$.
From the Bianchi identities
$$_\nu G^{\mu \nu }=0,$$
(15)
and from the field equations (11), we obtain
$$_\nu T^{\mu \nu }+\frac{1}{G}_\nu GT^{\mu \nu }\frac{1}{8\pi G}_\nu Q^{\mu \nu }=0.$$
(16)
A variation with respect to $`\varphi _\mu `$ yields the equations
$$_\nu B^{\mu \nu }+\frac{V(\varphi )}{\varphi _\mu }+\frac{1}{\omega }_\nu \omega B^{\mu \nu }=\frac{1}{\omega }J^\mu .$$
(17)
Taking the divergence of both sides with respect to $`_\mu `$, we get
$$_\mu \left(\frac{V(\varphi )}{\varphi _\mu }\right)+_\mu \left(\frac{1}{\omega }_\nu \omega B^{\mu \nu }\right)=_\mu \left(\frac{1}{\omega }J^\mu \right).$$
(18)
If we assume that the current density $`J^\mu `$ is conserved, we obtain the equation
$$_\mu \left(\frac{V(\varphi )}{\varphi _\mu }\right)+_\mu \left(\frac{1}{\omega }_\nu \omega B^{\mu \nu }\right)=_\mu \left(\frac{1}{\omega }\right)J^\mu .$$
(19)
In standard Maxwell-Proca theory, the conservation of the current $`J^\mu `$ is a separate physical assumption.
We shall choose for the potential $`V(\varphi )`$:
$$V(\varphi )=\frac{1}{2}\mu ^2\varphi ^\mu \varphi _\mu +W(\varphi ),$$
(20)
where $`W(\varphi )`$ denotes a vector field $`\varphi ^\mu `$ self-interaction contribution. We can choose as a model for the self-interaction:
$$W(\varphi )=\frac{1}{4}g(\varphi ^\mu \varphi _\mu )^2,$$
(21)
where $`g`$ is a coupling constant.
The effective gravitational “constant” $`G(x)`$ satisfies the field equations
$$_\alpha ^\alpha G+V^{}(G)+N=\frac{1}{2}G^2\left(T+\frac{\mathrm{\Lambda }}{4\pi G}\right),$$
(22)
where
$$N=\frac{3}{G}\left(\frac{1}{2}_\alpha G^\alpha GV(G)\right)+G\left(\frac{1}{2}_\alpha \omega ^\alpha \omega V(\omega )\right)$$
$$3\mathrm{\Theta }_\alpha G^\alpha G+\frac{G}{\mu ^2}\left(\frac{1}{2}_\alpha \mu ^\alpha \mu V(\mu )\right)+\frac{3G^2}{16\pi }_\alpha ^\alpha \mathrm{\Theta },$$
(23)
and $`T=g^{\mu \nu }T_{\mu \nu }`$. The scalar field $`\omega (x)`$ obeys the field equations
$$_\nu ^\nu \omega +V^{}(\omega )+F=0,$$
(24)
where
$$F=\mathrm{\Theta }_\alpha G^\alpha \omega +G\left(\frac{1}{4}B^{\mu \nu }B_{\mu \nu }+V(\varphi )\right).$$
(25)
The field $`\mu (x)`$ satisfies the equations
$$_\alpha ^\alpha \mu +V^{}(\mu )+P=0,$$
(26)
where
$$P=\left[\mathrm{\Theta }^\alpha G_\alpha \mu +\frac{2}{\mu }^\alpha \mu _\alpha \mu +\omega \mu ^2G\frac{V(\varphi )}{\mu }\right],$$
(27)
and the last term arises from the $`\mu `$ dependence of $`V(\varphi )`$ in (20).
If we adopt the condition
$$^\nu \varphi _\nu =0,$$
(28)
then (17) takes the form
$$^\nu _\nu \varphi _\mu R_{\mu }^{}{}_{}{}^{\nu }\varphi _\nu +\mu ^2\varphi _\mu \frac{W(\varphi )}{\varphi ^\mu }\frac{1}{\omega }^\nu \omega B_{\mu \nu }=\frac{1}{\omega }J_\mu .$$
(29)
The test particle action is given by
$$S_{TP}=m𝑑\tau \lambda 𝑑\tau \omega \varphi _\mu \frac{dx^\mu }{d\tau },$$
(30)
where $`\tau `$ is the proper time along the world line of the test particle and $`m`$ and $`\lambda `$ denote the test particle mass and coupling constant, respectively. The stationarity condition $`\delta S_{TP}/\delta x^\mu =0`$ yields the equations of motion for the test particle
$$m\left(\frac{d^2x^\mu }{d\tau ^2}+\mathrm{\Gamma }_{\alpha \beta }^\mu \frac{dx^\alpha }{d\tau }\frac{dx^\beta }{d\tau }\right)=f^\mu ,$$
(31)
where
$$f^\mu =\lambda \omega B_{}^{\mu }{}_{\nu }{}^{}\frac{dx^\nu }{d\tau }+\lambda ^\mu \omega \left(\varphi _\alpha \frac{dx^\alpha }{d\tau }\right)\lambda _\alpha \omega \left(\varphi ^\mu \frac{dx^\alpha }{d\tau }\right).$$
(32)
The action for the field $`B_{\mu \nu }`$ is of the Maxwell-Proca form for a massive vector field $`\varphi _\mu `$. It can be proved that this theory possesses a stable vacuum and the Hamiltonian is bounded from below. Even though the action is not gauge invariant, it can be shown that the longitudinal mode $`\varphi _0`$ (where $`\varphi _\mu =(\varphi _0,\varphi _i)(i=1,2,3)`$) does not propagate and the theory is free of ghosts. Similar arguments apply to the MSTG theory .
The Hamilton-Dirac (HD) method is a tool for investigating the constraints and the degrees of freedom of a field theory . The HD procedure checks the theory for consistency by producing the explicit constraints, and counting the number of degrees of freedom. It is a canonical initial value analysis. When the field theory is coupled dynamically to gravity, many vector field theories are ruled out, because the constraints produced by the canonical, Cauchy initial value formalism yield “derivative coupled” theories (the Christoffel connections do not cancel out). The Maxwell and Maxwell-Proca theories are prime examples of consistent vector field theories. With or without the gravitational field coupling, Maxwell’s theory has two degrees of freedom and Maxwell-Proca has three, and they are stable and satisfy a consistent Cauchy evolution analysis. Many other vector field theories are derivative coupled when the gravitational field is introduced into the action. Severe singularity problems can appear in vector-gravity coupled theories, which render them inconsistent. There are no pathological singularities in the Maxwell-Proca theory coupled to gravity, when one solves for the second time derivative in the canonical initial value formulation. In other vector theories, singularities occur that spoil the stability of the theory and rule them out as physically unviable theories.
It is possible to attribute the mass $`\mu `$ to a spontaneous symmetry breaking mechanism, but we shall not pursue this possibility at present.
## 3 Equations of Motion, Weak Fields and the Modified Gravitational Acceleration
Let us assume that we are in a distance scale regime for which the fields $`G`$, $`\omega `$ and $`\mu `$ take their approximate renormalized constant values:
$$GG_0(1+Z),\omega \omega _0A,\mu \mu _0B,$$
(33)
where $`G_0,\omega _0`$ and $`\mu _0`$ denote the “bare” values of $`G,\omega `$ and $`\mu `$, respectively, and $`Z,A`$ and $`B`$ are the associated renormalization constants.
For a static spherically symmetric field the line element is given by
$$ds^2=\gamma (r)dt^2\alpha (r)dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2).$$
(34)
The equations of motion for a test particle obtained from (31), (32) and (33) are given by
$$\frac{d^2r}{d\tau ^2}+\frac{\alpha ^{}}{2\alpha }\left(\frac{dr}{d\tau }\right)^2\frac{r}{\alpha }\left(\frac{d\theta }{d\tau }\right)^2r\left(\frac{\mathrm{sin}^2\theta }{\alpha }\right)\left(\frac{\sigma d\varphi }{d\tau }\right)^2+\frac{\gamma ^{}}{2\alpha }\left(\frac{dt}{d\tau }\right)^2$$
$$+\sigma \frac{1}{\alpha }\left(\frac{d\varphi _0}{dr}\right)\left(\frac{dt}{d\tau }\right)=0,$$
(35)
$$\frac{d^2t}{d\tau ^2}+\frac{\gamma ^{}}{\gamma }\left(\frac{dt}{d\tau }\right)\left(\frac{dr}{d\tau }\right)+\sigma \frac{1}{\gamma }\left(\frac{d\varphi _0}{dr}\right)\left(\frac{dr}{d\tau }\right)=0,$$
(36)
$$\frac{d^2\theta }{d\tau ^2}+\frac{2}{r}\left(\frac{d\theta }{d\tau }\right)\left(\frac{dr}{d\tau }\right)\mathrm{sin}\theta \mathrm{cos}\theta \left(\frac{d\varphi }{d\tau }\right)^2=0,$$
(37)
$$\frac{d^2\varphi }{d\tau ^2}+\frac{2}{r}\left(\frac{d\varphi }{d\tau }\right)\left(\frac{dr}{d\tau }\right)+2\mathrm{cot}\theta \left(\frac{d\varphi }{d\tau }\right)\left(\frac{d\theta }{d\tau }\right)=0,$$
(38)
where $`\sigma =\lambda \omega /m`$.
The orbit of the test particle can be shown to lie in a plane and by an appropriate choice of axes, we can make $`\theta =\pi /2`$. Integrating Eq.(38) gives
$$r^2\frac{d\varphi }{d\tau }=J,$$
(39)
where $`J`$ is the conserved orbital angular momentum. Integration of Eq.(36) gives
$$\frac{dt}{d\tau }=\frac{1}{\gamma }(\sigma \varphi _0+E),$$
(40)
where $`E`$ is the constant energy per unit mass.
By substituting (40) into (3) and using (39), we obtain
$$\frac{d^2r}{d\tau ^2}+\frac{\alpha ^{}}{2\alpha }\left(\frac{dr}{d\tau }\right)^2\frac{J^2}{\alpha r^3}+\frac{\gamma ^{}}{2\alpha \gamma ^2}(\sigma \varphi _0+E)^2=\sigma \frac{1}{\alpha \gamma }\left(\frac{d\varphi _0}{dr}\right)(\sigma \varphi _0+E).$$
(41)
We do not have an exact, spherically symmetric static solution to our field equations for a non-zero $`V(\varphi )`$. However, if we neglect $`V(\varphi )`$ in Eq.(17) and $`\mathrm{\Lambda }`$ in Eq.(11), then the exact static, spherically symmetric (Reissner-Nordström) solution in empty space yields the line element
$$ds^2=\left(1\frac{2GM}{r}+\frac{Q^2}{r^2}\right)dt^2\left(1\frac{2GM}{r}+\frac{Q^2}{r^2}\right)^1dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(42)
where $`M`$ is a constant of integration and
$$Q^2=4\pi G\omega ϵ^2.$$
(43)
Here, $`ϵ`$ denotes the “charge” of the spin -1 vector particle given by
$$ϵ=d^3xq,$$
(44)
where the matter current density $`J^\mu `$ is identified as $`J^\mu =(q,J^i)`$.
For large enough values of $`r`$, the solution (42) approximates the Schwarzschild metric components $`\alpha `$ and $`\gamma `$:
$$\alpha (r)\frac{1}{12GM/r},\gamma (r)1\frac{2GM}{r}.$$
(45)
It is not unreasonable to expect that a static, spherically symmetric solution of the field equations including the mass term $`\mu `$ in the field equations (17) will approximate for large values of $`r`$ the Schwarzschild metric components (45).
We assume that $`2GM/r1`$ and the slow motion approximation $`dr/dsdr/dt1`$. Then for material test particles, we obtain from (31), (32), (33), (41) and (45):
$$\frac{d^2r}{dt^2}\frac{J_N^2}{r^3}+\frac{GM}{r^2}=\sigma \frac{d\varphi _0}{dr},$$
(46)
where $`J_N`$ is the Newtonian orbital angular momentum.
For weak gravitational fields to first order, the static equations for $`\varphi _0`$ obtained from (29) are given for the source-free case by
$$\stackrel{}{}^2\varphi _0\mu ^2\varphi _0=0,$$
(47)
where $`\stackrel{}{}^2`$ is the Laplacian operator, and we have neglected any contribution from the self-interaction potential $`W(\varphi )`$. For a spherically symmetric static field $`\varphi _0`$, we obtain
$$\varphi _0^{\prime \prime }+\frac{2}{r}\varphi _0^{}\mu ^2\varphi _0=0.$$
(48)
This has the Yukawa solution
$$\varphi _0(r)=\beta \frac{\mathrm{exp}(\mu r)}{r},$$
(49)
where $`\beta `$ is a constant. We obtain from (46):
$$\frac{d^2r}{dt^2}\frac{J_N^2}{r^3}+\frac{GM}{r^2}=K\frac{\mathrm{exp}(\mu r)}{r^2}(1+\mu r),$$
(50)
where $`K=\sigma \beta `$.
We observe that the additional Yukawa force term in Eq.(50) is repulsive in accordance with the exchange of a spin 1 massive boson. We shall find that this repulsive component of the gravitational field is necessary to obtain a fit to galaxy rotation curves.
We shall write for the radial acceleration derived from (50):
$$a(r)=\frac{G_{\mathrm{}}M}{r^2}+K\frac{\mathrm{exp}(r/r_0)}{r^2}\left(1+\frac{r}{r_0}\right),$$
(51)
and $`G_{\mathrm{}}`$ is defined to be the effective gravitational constant at infinity
$$G_{\mathrm{}}=G_0\left(1+\sqrt{\frac{M_0}{M}}\right).$$
(52)
Here, $`M_0`$ denotes a parameter that vanishes when $`\omega =0`$ and $`G_0`$ is Newton’s gravitational “bare” constant. The constant $`K`$ is chosen to be
$$K=G_0\sqrt{MM_0}.$$
(53)
The choice of $`K`$, which determines the strength of the coupling of $`B_{\mu \nu }`$ to matter and the magnitude of the Yukawa force modification of weak Newtonian gravity, is based on phenomenology and is not at present derivable from the STVG action formalism.
By using (52), we can rewrite the acceleration in the form
$$a(r)=\frac{G_0M}{r^2}\left\{1+\sqrt{\frac{M_0}{M}}\left[1\mathrm{exp}(r/r_0)\left(1+\frac{r}{r_0}\right)\right]\right\}.$$
(54)
We can generalize this to the case of a mass distribution by replacing the factor $`G_0M/r^2`$ in (54) by $`G_0(r)/r^2`$. The rotational velocity of a star $`v_c`$ is obtained from $`v_c^2(r)/r=a(r)`$ and is given by
$$v_c=\sqrt{\frac{G_0(r)}{r}}\left\{1+\sqrt{\frac{M_0}{M}}\left[1\mathrm{exp}(r/r_0)\left(1+\frac{r}{r_0}\right)\right]\right\}^{1/2}.$$
(55)
The gravitational potential for a point source obtained from the modified acceleration law (54) is given by
$$\mathrm{\Phi }(r)=\frac{G_0M}{r}\left[1+\sqrt{\frac{M_0}{M}}(1\mathrm{exp}(r/r_0))\right].$$
(56)
The acceleration law (54) can be written as
$$a(r)=\frac{G(r)M}{r^2},$$
(57)
where
$$G(r)=G_0\left[1+\sqrt{\frac{M_0}{M}}\left(1\mathrm{exp}(r/r_0)\left(1+\frac{r}{r_0}\right)\right)\right]$$
(58)
is an effective expression for the variation of $`G`$ with respect to $`r`$. A good fit to a large number of galaxies has been achieved with the parameters :
$$M_0=9.60\times 10^{11}M_{},r_0=13.92\mathrm{kpc}=4.30\times 10^{22}\mathrm{cm}.$$
(59)
In the fitting of the galaxy rotation curves for both LSB and HSB galaxies, using photometric data to determine the mass distribution $`(r)`$ , only the mass-to-light ratio $`M/L`$ is employed, once the values of $`M_0`$ and $`r_0`$ are fixed universally for all LSB and HSB galaxies. Dwarf galaxies are also fitted with the parameters :
$$M_0=2.40\times 10^{11}M_{},r_0=6.96\mathrm{kpc}=2.15\times 10^{22}\mathrm{cm}.$$
(60)
By choosing values for the parameters $`G_{\mathrm{}}`$, $`(M_0)_{\mathrm{clust}}`$ and $`(r_0)_{\mathrm{clust}}`$, we are able to obtain satisfactory fits to a large sample of X-ray cluster data .
## 4 Orbital Equations of Motion
We set $`\theta =\pi /2`$ in (34), divide the resulting expression by $`d\tau ^2`$ and use Eqs.(39) and (40) to obtain
$$\left(\frac{dr}{d\tau }\right)^2+\frac{J^2}{\alpha r^2}\frac{1}{\alpha \gamma }(\sigma \varphi _0+E)^2=\frac{E}{\alpha }.$$
(61)
We have $`ds^2=Ed\tau ^2`$, so that $`ds/d\tau `$ is a constant. For material particles $`E>0`$ and for massless photons $`E=0`$.
Let us set $`u=1/r`$ and by using (39), we have $`dr/d\tau =Jdu/d\varphi `$. Substituting this into (61), we obtain
$$\left(\frac{du}{d\varphi }\right)^2=\frac{1}{\alpha \gamma J^2}(E+\sigma \varphi _0)^2\frac{1}{\alpha r^2}\frac{E}{\alpha J^2}.$$
(62)
On substituting (45) and $`dr/d\varphi =(1/u^2)du/d\varphi `$ into (62), we get after some manipulation:
$$\frac{d^2u}{d\varphi ^2}+u=\frac{EGM}{J^2}\frac{EK}{J^2}\mathrm{exp}\left(\frac{1}{r_0u}\right)\left(1+\frac{1}{r_0u}\right)+3GMu^2,$$
(63)
where $`r_0=1/\mu `$.
For material test particles $`E=1`$ and we obtain
$$\frac{d^2u}{d\varphi ^2}+u=\frac{GM}{J^2}+3GMu^2\frac{K}{J^2}\mathrm{exp}\left(\frac{1}{r_0u}\right)\left(1+\frac{1}{r_0u}\right).$$
(64)
On the other hand, for massless photons $`ds^2=0`$ and $`E=0`$ and (63) gives
$$\frac{d^2u}{d\varphi ^2}+u=3GMu^2.$$
(65)
## 5 Solar System and Binary Pulsar Observations
We obtain from Eq.(64) the orbit equation (we reinsert the speed of light c):
$$\frac{d^2u}{d\varphi ^2}+u=\frac{GM}{c^2J^2}\frac{K}{c^2J^2}\mathrm{exp}(r/r_0)\left[1+\left(\frac{r}{r_0}\right)\right]+\frac{3GM}{c^2}u^2.$$
(66)
Using the large $`r`$ weak field approximation, and the expansion
$$\mathrm{exp}(r/r_0)=1\frac{r}{r_0}+\frac{1}{2}\left(\frac{r}{r_0}\right)^2+\mathrm{}$$
(67)
we obtain the orbit equation for $`rr_0`$:
$$\frac{d^2u}{d\varphi ^2}+u=N+3\frac{GM}{c^2}u^2,$$
(68)
where
$$N=\frac{GM}{c^2J_N^2}\frac{K}{c^2J_N^2}.$$
(69)
We can solve Eq.(68) by perturbation theory and find for the perihelion advance of a planetary orbit
$$\mathrm{\Delta }\omega =\frac{6\pi }{c^2L}(GM_{}K_{}),$$
(70)
where $`J_N=(GM_{}L/c^2)^{1/2}`$, $`L=a(1e^2)`$ and $`a`$ and $`e`$ denote the semimajor axis and the eccentricity of the planetary orbit, respectively.
We now use the running of the effective gravitational coupling constant $`G=G(r)`$, determined by (58) and find that for the solar system $`rr_0`$, we have $`GG_0`$ within the experimental errors for the measurement of Newton’s constant $`G_0`$. We choose for the solar system
$$\frac{K_{}}{c^2}1.5\mathrm{km}$$
(71)
and use $`G=G_0`$ to obtain from (70) a perihelion advance of Mercury in agreement with GR. The bound (71) requires that the coupling constant $`\omega `$ varies with distance in such a way that it is sufficiently small in the solar system regime and determines a value for $`M_0`$, in Eq.(53), that is in accord with the bound (71).
For terrestrial experiments and orbits of satellites, we have also that $`GG_0`$ and for $`K_{}`$ sufficiently small, we then achieve agreement with all gravitational terrestrial experiments including Eötvös free-fall experiments and “fifth force” experiments.
For the binary pulsar PSR 1913+16 the formula (70) can be adapted to the periastron shift of a binary system. Combining this with the STVG gravitational wave radiation formula, which will approximate closely the GR formula, we can obtain agreement with the observations for the binary pulsar. The mean orbital radius for the binary pulsar is equal to the projected semi-major axis of the binary, $`r_N=7\times 10^{10}\mathrm{cm}`$, and we choose $`r_Nr_0`$. Thus, for $`G=G_0`$ within the experimental errors, we obtain agreement with the binary pulsar data for the periastron shift when
$$\frac{K_N}{c^2}4.2\mathrm{km}.$$
(72)
For a massless photon $`E=0`$ and we have
$$\frac{d^2u}{d\varphi ^2}+u=3\frac{GM}{c^2}u^2.$$
(73)
For the solar system using $`G=G_0`$ within the experimental errors gives the light deflection:
$$\mathrm{\Delta }_{}=\frac{4G_0M_{}}{c^2R_{}}$$
(74)
in agreement with GR.
## 6 Galaxy Clusters and Lensing
The bending angle of a light ray as it passes near a massive system along an approximately straight path is given to lowest order in $`v^2/c^2`$ by
$$\theta =\frac{2}{c^2}|a^{}|𝑑z,$$
(75)
where $``$ denotes the perpendicular component to the ray’s direction, and dz is the element of length along the ray and $`a`$ denotes the acceleration.
From (73), we obtain the light deflection
$$\mathrm{\Delta }=\frac{4GM}{c^2R}=\frac{4G_0\overline{M}}{c^2R},$$
(76)
where
$$\overline{M}=M\left(1+\sqrt{\frac{M_0}{M}}\right).$$
(77)
The value of $`\overline{M}`$ follows from (58) for clusters as $`rr_0`$ and
$$G(r)G_{\mathrm{}}=G_0\left(1+\sqrt{\frac{M_0}{M}}\right).$$
(78)
We choose for a cluster $`M_0=3.6\times 10^{15}M_{}`$ and a cluster mass $`M_{\mathrm{clust}}10^{14}M_{}`$, and obtain
$$\left(\sqrt{\frac{M_0}{M}}\right)_{\mathrm{clust}}6.$$
(79)
We see that $`\overline{M}7M`$ and we can explain the increase in the light bending without exotic dark matter.
From the formula Eq.(54) for $`rr_0`$ we get
$$a(r)=\frac{G_0\overline{M}}{r^2}.$$
(80)
We expect to obtain from this result a satisfactory description of lensing phenomena using Eq.(75).
An analysis of a large number of clusters shows that the MSTG and STVG theories fit well the cluster data in terms of the cluster mass, $`M_{\mathrm{clust}}`$, and an average value for the parameter $`M_0`$ .
## 7 Running of the Effective Constants $`G`$, $`\omega `$ and $`\mu `$
The scaling with distance of the effective gravitational constant $`G`$, the effective coupling constant $`\omega `$ and the effective mass $`\mu `$ is seen to play an important role in describing consistently the solar system and the galaxy and cluster dynamics, without the postulate of exotic dark matter. We have to solve the field equations (22), (24) and (26) with given potentials $`V(G),V(\omega )`$ and $`V(\mu )`$ to determine the variation of the effective constants with space and time. These equations are complicated, so we shall make simplifying approximations. In Eq.(22), we shall neglect the contributions from $`N`$ and obtain for $`T=\mathrm{\Lambda }=0`$:
$$_\nu ^\nu G+V^{}(G)=0,$$
(81)
where $`f^{}(y)=df/dy`$. The effective variation of $`G`$ with $`r`$ is determined by Eq.(58). We obtain from (81) for the static spherically symmetric equations
$$\stackrel{}{}^2G(r)V^{}(G)G^{\prime \prime }(r)+\frac{2}{r}G^{}(r)V^{}(G)=0.$$
(82)
By choosing the potential
$$V(G)=\frac{1}{2}\left(\frac{G_0}{r_0^2}\right)\left(\frac{M_0}{M}\right)\mathrm{exp}(2r/r_0)(1+\frac{2r}{r_0}\frac{r^2}{r_0^2}),$$
(83)
we obtain a solution to (82) for $`G(r)`$ given by (58). The neglect of the contributions $`N`$ can only be justified by solving the complete set of coupled equations by a perturbation calculation. We will not attempt to do this in the present work, but we plan to investigate this issue in a future publication.
We see from (58) that for $`rr_0`$ we obtain $`G(r)G_0`$. As the distance scale approaches the regime of the solar system $`r<100A.U.`$ where $`1A.U.=1.496\times 10^{13}\mathrm{cm}=4.85\times 10^9\mathrm{kpc}`$, then (54) becomes the Newtonian acceleration law:
$$a(r)=\frac{G_0M}{r^2},$$
(84)
in agreement with solar physics observations for the inner planets.
Let us make the approximation of neglecting $`F(\varphi )`$ in Eq.(24). In the static spherically symmetric case this gives
$$\omega ^{\prime \prime }(r)+\frac{2}{r}\omega ^{}(r)V^{}(\omega )=0.$$
(85)
We choose as a solution for $`\omega (r)`$:
$$\omega (r)=\omega _0\{1+\overline{\omega }[1\mathrm{exp}(\overline{\mu }r)(1+\overline{\mu }r)]\},$$
(86)
where $`\overline{\omega }`$ and $`\overline{\mu }`$ are positive constants. The potential $`V(\omega )`$ has the form
$$V(\omega )=\frac{1}{2}\omega _0^2\overline{\mu }^2\overline{\omega }^2\mathrm{exp}(2\overline{\mu }r)(1+2\overline{\mu }r\overline{\mu }^2r^2).$$
(87)
For the variation of the renormalized mass $`\mu =\mu (r)`$, we find that a satisfactory solution to Eq.(26) should correspond to a a $`\mu (r)=1/r_0(r)`$ that decreases from a value for the inner planets of the solar system, consistent with solar system observations, to a small value corresponding to $`r_0`$ for the galaxy fits, $`r_0=14`$ kpc, and to an even smaller value for the cluster data fits.
The spatial variations of $`G(r)`$, $`\omega (r)`$ and $`\mu (r)=1/r_0(r)`$ can be determined numerically from the equations (22), (24) and (26) with given potentials $`V(G),V(\omega )`$ and $`V(\mu )`$, such that for the solar system and the binary pulsar PSR 1913+16 the bounds (71) and (72) are satisfied by the solutions for $`G(r),\omega (r)`$ and $`\mu (r)`$. The spatial variations of $`G`$, the coupling constant $`\omega `$ and the range $`r_0`$ are required to guarantee consistency with solar system observations. On the other hand, their increase at galactic and cosmological distance and time scales can account for galaxy rotation curves, cluster lensing and cosmology without non-baryonic dark matter.
We have constructed a classical action for gravity that can be considered as an effective field theory description of an RG flow quantum gravity scenario as described in refs. and .
The fitting of the solar system, galaxy and the clusters of galaxies data depends on the running of the of the “constants” $`G`$, $`\omega `$ and $`r_0`$. They should increase from one distance scale to the next according to the renormalization group flow diagrams, or the solutions of the classical field equations in the present article. In a future article, the author plans to provide a more complete determination of the running of the constants. However, the present article describes the basic scenario and the ideas underlying the theory.
## 8 Cosmology
Let us now consider a cosmological solution to our STVG theory. We adopt a homogeneous and isotropic FLRW background geometry with the line element
$$ds^2=dt^2a^2(t)\left(\frac{dr^2}{1kr^2}+r^2d\mathrm{\Omega }^2\right),$$
(88)
where $`d\mathrm{\Omega }^2=d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2`$ and $`k=0,1,+1`$ for a spatially flat, open and closed universe, respectively. In this background spacetime, we have $`\varphi _0\psi 0`$, $`\varphi _i=0`$ and $`B_{\mu \nu }=0`$.
We define the energy-momentum tensor for a perfect fluid by
$$T^{\mu \nu }=(\rho +p)u^\mu u^\nu pg^{\mu \nu },$$
(89)
where $`u^\mu =dx^\mu /ds`$ is the 4-velocity of a fluid element and $`g_{\mu \nu }u^\mu u^\nu =1`$. Moreover, we have
$$\rho =\rho _M+\rho _\varphi +\rho _S,p=p_M+p_\varphi +p_S,$$
(90)
where $`\rho _i`$ and $`p_i`$ denote the components of density and pressure associated with the matter, the $`\varphi ^\mu `$ field and the scalar fields $`G`$, $`\omega `$ and $`\mu `$, respectively.
The Friedmann equations take the form
$$H^2(t)+\frac{k}{a^2(t)}=\frac{8\pi G(t)\rho (t)}{3}+\frac{\dot{a}}{a}\frac{\dot{G}}{G}+\frac{\mathrm{\Lambda }}{3},$$
(91)
$$\frac{\ddot{a}(t)}{a(t)}=\frac{4\pi G(t)}{3}(\rho (t)+3p(t))+\frac{1}{2}\left(\frac{\ddot{G}}{G}\frac{\dot{G}^2}{G^2}+\frac{2\dot{a}}{a}\frac{\dot{G}}{G}\right)+\frac{\mathrm{\Lambda }}{3},$$
(92)
where $`H(t)=\dot{a}(t)/a(t)`$.
Let us make the simplifying approximation for equations (22):
$$\ddot{𝒢}+3H\dot{𝒢}+V^{}(𝒢)=\frac{1}{2}G_0𝒢^2\left[\rho 3p+\frac{\mathrm{\Lambda }}{4\pi G_0𝒢}\right],$$
(93)
where $`𝒢(t)=G(t)/G_0`$. A solution for $`𝒢`$ in terms of a given potential $`V(𝒢)`$ and for given values of $`\rho `$ and $`p`$ can be obtained from (93).
The solution for $`𝒢`$ must satisfy a constraint at the time of big bang nucleosynthesis . The number of relativistic degrees of freedom is very sensitive to the cosmic expansion rate at 1 MeV. This can be used to constrain the time dependence of $`G`$. Recent measurements of the $`{}_{}{}^{4}He`$ mass fraction and the deuterium abundance at 1 MeV leads to the constraint $`G(t)G_0`$. We impose the conditions $`𝒢(t)1`$ as $`tt_{BBN}`$ and $`𝒢(t)1+\overline{\omega }`$ as $`tt_{SLS}`$ where $`t_{BBN}`$ and $`t_{SLS}`$ denote the times of the big bang nucleosynthesis and the surface of last scattering, respectively. A possible solution for $`𝒢`$ can take the form
$$𝒢(t)=1+\overline{\omega }\left[1\mathrm{exp}(t/T)\left(1+\frac{t}{T}\right)\right],$$
(94)
where $`\overline{\omega }`$ and $`T`$ are constants and $`\overline{\omega }`$ is a measure of the magnitude of the scalar field $`\psi `$. We have for $`t>>T`$ that $`𝒢1+\overline{\omega }`$ and for $`t<<T`$ that $`𝒢1`$. We get
$$\dot{𝒢}=\frac{\overline{\omega }t}{T^2}\mathrm{exp}(t/T),\ddot{𝒢}=\frac{\overline{\omega }}{T^2}\mathrm{exp}(t/T)\left(1\frac{t}{T}\right).$$
(95)
It follows that $`\dot{𝒢}(t)0`$ for $`t>>T`$, which allows us for a suitable choice of $`T`$ to satisfy the experimental bound from the Cassini spacecraft measurements :
$$|\dot{G}/G|10^{13}\mathrm{yr}^1.$$
(96)
A linear perturbation on the FLRW background will link the theory with observations of anisotropies in the CMB as well as galaxy clustering on large scales. The basic fields are perturbed around the background spacetime (denoted for a quantity $`Y`$ by $`\stackrel{~}{Y}`$). In the conformal metric with the time transformation $`d\eta =dt/a(t)`$:
$$ds^2=a^2(\eta )(d\eta ^2d\stackrel{}{x}^2),$$
(97)
the metric perturbations are in the conformal Newtonian gauge
$$g_{00}(\stackrel{}{x},t)=a^2(t)(1+2\mathrm{\Phi }(\stackrel{}{x},t)),g_{ij}(\stackrel{}{x},t)=a^2(t)(12\mathrm{\Phi }(\stackrel{}{x},t))\delta _{ij},$$
(98)
where $`\mathrm{\Phi }`$ is the gravitational potential. The vector field perturbations are defined by
$$\varphi _\mu (\stackrel{}{x},t)=a(t)(\stackrel{~}{\varphi }_\mu (t)+\delta \varphi _\mu (\stackrel{}{x},t)),$$
(99)
where $`\stackrel{~}{\varphi }_i(t)=0`$. Denoting by $`\chi _i`$ the scalar fields $`\chi _1=G`$, $`\chi _2=\omega `$ and $`\chi _3=\mu `$, the scalar field perturbations are
$$\chi _i(\stackrel{}{x},t)=\stackrel{~}{\chi }_i(t)+\delta \chi _i(\stackrel{}{x},t).$$
(100)
The problem that gravity theories such as MSTG and STVG face in describing cosmology with no cold dark matter (CDM) (non-baryonic dark matter) is the damping of perturbations during the recombination era. In a pure baryonic universe evolving according to Einstein’s gravitational field equations, the coupling of baryons to photons during the recombination era will suffer Silk damping, causing the collisional propagation of radiation from overdense to underdense regions . In the CDM model, the perturbations $`\delta _{CDM}`$ are undamped during recombination, because the CDM particles interact with gravity and only weakly with matter (photons). The Newtonian potential in the CDM model is approximately given by
$$k^2\mathrm{\Phi }4\pi G_0(\rho _b\delta _b+\rho _{CDM}\delta _{CDM}),$$
(101)
where $`k^2`$ denotes the square of the wave number, $`\rho _b`$ and $`\rho _{CDM}`$ denote the densities of baryons and cold dark matter, respectively, and $`\delta _i`$ denotes the perturbation density contrast for each component i of matter. If $`\rho _{CDM}`$ is sufficiently large, then $`\delta _{CDM}`$ will not be erased, whereas $`\delta _b`$ decreases during recombination . In STVG, the imperfect fluid plasma before recombination has two components: the dominant neutral vector field component that does not couple to photons and the baryon-photon component. The vector field component has zero pressure and zero shear viscosity, so the vector field perturbations are not Silk damped like the baryon-radiation perturbations, for the latter have non-vanishing pressure and shear viscosity.
In STVG, we have
$$\mathrm{\Omega }_b=\frac{8\pi G\rho _b}{3H^2},\mathrm{\Omega }_\psi =\frac{8\pi G\rho _\psi }{3H^2},\mathrm{\Omega }_S=\frac{8\pi G\rho _S}{3H^2}.$$
(102)
We also have a possible contribution from massive neutrinos
$$\mathrm{\Omega }_\nu =\frac{8\pi G\rho _\nu }{3H^2},$$
(103)
where $`\rho _\nu `$ denotes the density of neutrinos. We assume that the vector field density dominates.
The Newtonian potential in our modified gravitational model becomes
$$k^2\mathrm{\Phi }4\pi G_{\mathrm{ren}}[\rho _b\delta _b+\rho _\nu \delta _\nu +\rho _\psi \delta _\psi +\rho _S\delta _S],$$
(104)
where $`G_{\mathrm{ren}}`$ is the renormalized value of the gravitational constant. Assuming that the density $`\rho _\psi `$ is significant before and during recombination, we can consider fitting the acoustic peaks in the power spectrum in a spatially flat universe with the parameters
$$\mathrm{\Omega }=\mathrm{\Omega }_b+\mathrm{\Omega }_\nu +\mathrm{\Omega }_\psi +\mathrm{\Omega }_\nu +\mathrm{\Omega }_\mathrm{\Lambda }=1.$$
(105)
The fitting of the acoustic peaks in the CMB power spectrum does not permit a too large value of $`\mathrm{\Omega }_\mathrm{\Lambda }`$. Moreover, the neutrino contribution is constrained by the three neutrinos having a mass $`<2`$ eV. We could choose, for example, $`\mathrm{\Omega }_b=0.04,\mathrm{\Omega }_\psi =0.25,\mathrm{\Omega }_\nu =0.01,\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ as a possible choice of parameters to fit the data. There are now new data from the balloon borne Boomerang CMB observations that together with other ground based observations and WMAP data determine more accurately the height of the third acoustic peak in the angular CMB power spectrum . The ratio of the height of the first peak to the second peak determines the baryon content $`\mathrm{\Omega }_b0.04`$. The height of the third peak is determined by the amount of cold dark matter, and in the modified gravity theory by the possible amounts of scalar $`\psi `$, $`G`$, $`\omega `$ and $`\mu `$ contributions. In particular, the dominant neutral $`\psi `$ vector component perturbations will not be washed out before recombination.
Can the effects of gravitational constant renormalization together with the possible effects of the densities $`\rho _\psi `$ and $`\rho _S`$ describe a universe which can reproduce the current galaxy and CMB observations? The answer to this problem will be addressed in a future publication.
## 9 Conclusions
A modified gravity theory based an a $`D=4`$ pseudo-Riemannian metric, a spin $`1`$ vector field and a corresponding second-rank skew field $`B_{\mu \nu }`$ and dynamical scalar fields $`G`$, $`\omega `$ and $`\mu `$, yields a static spherically symmetric gravitational field with an added Yukawa potential and with an effective coupling strength and distance range. This modified acceleration law leads to remarkably good fits to a large number of galaxies and galaxy clusters . The previously published gravitational theories NGT and MSTG yielded the same modified weak gravitational field acceleration law and, therefore, the same successful fits to galaxy and cluster data. The MSTG and STVG gravity theories can both be identified generically as metric-skew-tensor gravity theories, for they both describe gravity as a metric theory with an additional degree of freedom associated with a skew field coupling to matter. For MSTG, this is a third-rank skew field $`F_{\mu \nu \lambda }`$, while for STVG the skew field is a second-rank tensor $`B_{\mu \nu }`$. However, MSTG is distinguished from STVG as being the weak field approximation to the nonsymmetric gravitational theory (NGT).
An action $`S_S`$ for the scalar fields $`G(x)`$, $`\omega (x)`$ and $`\mu (x)=1/r_0(x)`$ and the field equations resulting from a variation of the action, $`\delta S_S=0`$, can be incorporated into the NGT and MSTG theory. The dynamical solutions for the scalar fields give an effective description of the running of the constants in an RG flow quantum gravity scenario, in which strong infrared renormalization effects and increasing large scale spatial values of $`G`$ and $`\omega `$ lead, together with the modified acceleration law, to a satisfactory description of galaxy rotation curves and cluster dynamics without non-baryonic dark matter.
We have demonstrated that a cosmological model with the renormalized gravitational constant $`G_{\mathrm{ren}}`$ and contributions from the scalar fields $`G(t)`$, $`\omega (t)`$ and $`\mu (t)`$ can possibly lead to a satisfactory description of the distribution of galaxies and the CMB power spectrum in a baryon dominated universe.
The neutral vector particle $`\varphi `$ does not couple to radiation and it has zero pressure $`p`$ and zero shear viscosity. Since it dominates the period of recombination, its perturbations associated with the plasma fluid will not be washed out by Silk damping. In contrast to standard dark matter models, we should not search for new stable particles such as weakly interacting massive particles (WIMPS) or neutralinos, because the fifth force charge in STVG that is the source of the neutral vector field (skew field) is carried by the known stable baryons (and electrons and neutrinos). This new charge is the source of a fifth force skew field that modifies the gravitational field in the universe.
Acknowledgments
This work was supported by the Natural Sciences and Engineering Research Council of Canada. I thank Yujun Chen, Joel Brownstein Martin Green and Pierre Savaria for helpful discussions.
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# Search for coherent charged pion production in neutrino-carbon interactions
The K2K Collaboration
## Abstract
We report the result from a search for charged-current coherent pion production induced by muon neutrinos with a mean energy of 1.3 GeV. The data are collected with a fully active scintillator detector in the K2K long-baseline neutrino oscillation experiment. No evidence for coherent pion production is observed and an upper limit of $`0.60\times 10^2`$ is set on the cross section ratio of coherent pion production to the total charged-current interaction at 90% confidence level. This is the first experimental limit for coherent charged pion production in the energy region of a few GeV.
The charged-current (CC) coherent pion production in neutrino-nucleus scattering, $`\nu _\mu +A\mu ^{}+\pi ^++A`$, is a process in which the neutrino scatters coherently off the entire nucleus with a small energy transfer. Such a process has been measured in a number of experiments Vilain:1993pf ; Allport:1989pf ; Grabosch:1986pf ; Willocq:1993pf , providing a test of the partially conserved axial-vector current (PCAC) hypothesis Adler:1964 . The existing data agree with the Rein and Sehgal model Rein:1982pf based on the PCAC hypothesis for neutrino energies from 7 to 100 GeV, while there exists no measurement at lower energies.
The recent discovery of neutrino oscillations has renewed interest in neutrino-nucleus interactions in the sub- to few GeV region. The KEK to Kamioka (K2K) long-baseline neutrino oscillation experiment has reported Aliu:2004 a significant deficit in the forward scattering events, which limits the prediction accuracy of the neutrino energy spectrum at the far detector. CC coherent pion production is one of the candidate interactions responsible for this deficit and its study is necessary to improve the accuracy of the current and future atmospheric/accelerator-based neutrino oscillation experiments, which are expected to achieve much improved statistical precision using interactions of neutrinos in the same energy region as K2K.
This letter presents the result from a search for CC coherent pion production by neutrinos in the K2K experiment. We compare our result specifically with the Rein and Sehgal model Rein:1982pf because it is the only model that provides the kinematics of pions and is commonly used in neutrino oscillation experiments.
In the K2K experiment, protons are extracted from the KEK 12 GeV proton synchrotron and hit an aluminum target. Positively charged secondary particles, mainly pions, are focused by a magnetic horn system and decay to produce an almost pure (98%) $`\nu _\mu `$ beam with a mean energy of 1.3 GeV Ahn:2001cq . The neutrino beam energy spectrum and spatial profile are measured using a set of near neutrino detectors located 300 m downstream from the proton target. The estimated absolute flux has a large uncertainty due to difficulties in the absolute estimation of the primary proton beam intensity, the proton targeting efficiency, and hadron production cross sections. Therefore, the ratio of the CC coherent pion to the total CC cross section is measured, rather than the absolute CC coherent pion cross section. The data used for this analysis were collected with one of the near detectors, the fully active scintillator detector (SciBar), from October 2003 to February 2004, corresponding to $`1.7\times 10^{19}`$ protons on target (POT).
The SciBar detector Nitta:2004nt consists of 14,848 extruded plastic scintillator strips read out by wavelength-shifting fibers and multi-anode photomultipliers. The scintillator also acts as the neutrino interaction target; it is a fully active detector and has high efficiency for low momentum particles. Scintillator strips with dimensions of $`1.3\times 2.5\times 300\mathrm{cm}^3`$ are arranged in 64 layers. Each layer consists of two planes to measure horizontal and vertical position. The total size of the detector is $`3.0\times 3.0\times 1.7\mathrm{m}^3`$, while an inner volume of $`2.6\times 2.6\times 1.35\mathrm{m}^3`$ (9.38 tons) is used as the fiducial volume to reject incoming particles and obtain a flat efficiency for CC interactions. The minimum reconstructible track length is 8 cm. A track finding efficiency of more than 99% is achieved for single tracks with a length of more than 10 cm. The track finding efficiency for a second, shorter track is lower than that for single tracks due to overlap with the first track. This efficiency increases with the length of the second track and reaches 90% at a track length of 30 cm.
The NEUT Monte Carlo (MC) simulation program library hayato:2002 is used to simulate neutrino-nucleus interactions. The CC coherent pion production is incorporated in the simulation based on the Rein and Sehgal model Rein:1982pf , which predicts the cross section averaged over the K2K neutrino energy spectrum of $`2.85\times 10^{40}\mathrm{cm}^2`$/nucleon for carbon. The Llewellyn Smith model Llewellyn:1972 and the Rein and Sehgal model Rein:1981 are employed for quasi-elastic (QE) scattering ($`\nu _\mu +n\mu ^{}+p`$) and CC single pion (1$`\pi `$) production ($`\nu _\mu +N\mu +N+\pi `$), where N is a nucleon, respectively. The axial vector mass of the nucleon form factor is set to be 1.1 $`\mathrm{GeV}/c^2`$ for both QE and CC1$`\pi `$ interactions Bernard:2002 . For deep inelastic scattering (DIS), we use GRV94 nucleon structure functions GRV:1995 with a correction by Bodek and Yang Bodek:2002 . Nuclear effects are taken into account; for the pions originating from neutrino interactions, absorption, elastic scattering, and charge exchange inside the target nucleus are simulated. Pion cross sections are calculated using the model by Salcedo et al. Salcedo:1988 , which agrees well with past experimental data Ingram:1983 . Pion interactions outside the target nucleus are simulated based on other experimental data Carroll:1976 .
For the present analysis, the experimental signatures of CC coherent pion production are the existence of exactly two tracks, both consistent with minimum ionizing particles, and small momentum transfer defined as $`q^2(P_\mu P_\nu )^2`$, where $`P_\mu `$ and $`P_\nu `$ are the four momenta of the muon and the neutrino, respectively. According to the MC simulation, the dominant background is the CC1$`\pi `$ production, where the proton is below threshold or the neutron is invisible.
Charged current (CC) candidate events are selected by requiring that at least one reconstructed track starting in the fiducial volume is matched with a track or hits in the muon range detector (MRD) Ishii:2001sj located just behind SciBar (SciBar-MRD sample). This criterion imposes a threshold for muon momentum ($`p_\mu `$) of 450 MeV/c. According to the MC simulation, 98% of the events selected by this requirement are CC induced events, and the rest are neutral current (NC) interactions accompanied by a charged pion or proton which penetrates into the MRD. The contribution from $`\nu _e`$ is negligible ($`<0.4\%`$). The momentum of the muon is reconstructed from its range through SciBar and MRD. The resolutions for $`p_\mu `$ and the angle with respect to the neutrino beam direction ($`\theta _\mu `$) are determined to be 80 MeV/c and 1.6 degrees, respectively.
From the SciBar-MRD sample, events with two reconstructed tracks are selected. The QE candidate events are rejected by using kinematic information Aliu:2004 . Events in which the shorter track is identified as proton-like based on dE/dx information (non-QE-proton sample) are also rejected to select the non-QE-pion sample, which includes the signal candidates. The particle identification capability is verified using cosmic ray muons and the shorter tracks in the QE sample, where the latter provides a proton sample with more than 90%. The probability to mis-identify a muon track as proton-like is 1.7% with a corresponding proton selection efficiency of 90%.
The CC coherent pion candidates are extracted from the non-QE-pion sample. The background events are suppressed by requiring that the pion-like track goes forward. Even if the additional particles in the background process are not reconstructed as tracks, they can be detected as a large energy deposit or additional hits around the vertex. Figure 1(a) shows a distribution of energy deposited in the vertex strip ($`E_{\mathrm{vtx}}`$) for the non-QE-pion sample. The MC prediction for $`E_{\mathrm{vtx}}`$ is verified with the QE sample, which has no contribution from non-visible particles, as shown in Fig. 1(b). We require the events to have $`E_{\mathrm{vtx}}`$ less than 7 MeV and no additional hits around the vertex strip.
The value of $`q^2`$ reconstructed from $`p_\mu `$ and $`\theta _\mu `$ under the assumption of QE interaction is denoted $`q_{\mathrm{rec}}^2`$, and is calculated using
$`p_\nu `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{(M_p^2m_\mu ^2)+2E_\mu (M_nV)(M_nV)^2}{E_\mu +(M_nV)+p_\mu \mathrm{cos}\theta _\mu }}`$
where $`M_{p(n)}`$ is the proton (neutron) mass, $`m_\mu `$ is the muon mass and V is the nuclear potential set to 27 MeV. The $`q_{\mathrm{rec}}^2`$ for coherent pion production events, which is expected to be very small due to the small scattering angle for muons, is shifted from the true $`q^2`$ by 0.008 $`(\mathrm{GeV}/c)^2`$ with a resolution of 0.014 $`(\mathrm{GeV}/c)^2`$. Events are required to have a reconstructed $`q^2`$ of less than 0.10 $`(\mathrm{GeV}/c)^2`$.
The background contamination in the final sample is estimated by the MC simulation. In order to constrain the uncertainties, the $`q_{\mathrm{rec}}^2`$ distributions of the data in the region $`q_{\mathrm{rec}}^2>0.10(\mathrm{GeV}/c)^2`$ are fitted with MC expectations. The one track sample is used as well as two-track QE, non-QE-proton and non-QE-pion samples, and these four samples are fitted simultaneously. In the fit, the non-QE to QE relative cross section ratio, the magnitude of the nuclear effects and the momentum scale for muons are treated as free parameters. Figure 2 shows the $`q_{\mathrm{rec}}^2`$ distributions of the data with the MC simulation after the fitting. The $`\chi ^2`$ value in the regions with $`q_{\mathrm{rec}}^2>0.10(\mathrm{GeV}/c)^2`$ at the best fit is 73.2 for 82 degrees of freedom (DOF).
Figure 3 shows the $`q_{\mathrm{rec}}^2`$ distribution for the final CC coherent pion sample. The number of events in each selection step is summarized in Table 1 together with the signal efficiency and purity. In the signal region, 113 coherent pion candidates are found. The neutrino energy spectra for coherent pion events and the efficiency as a function of neutrino energy, estimated using the MC simulation, are shown in Fig. 4(a) and 4(c), respectively. The total efficiency is 21.1%. The expected number of background events in the signal region is 111.4. After subtracting the background and correcting for the efficiency, the number of coherent pion events is measured to be $`7.64\pm 50.40(\mathrm{stat}.)`$, while 470 events are expected from the MC simulation. Hence, no evidence of coherent pion production is found in the present data set.
The total number of CC interactions is estimated by using the SciBar-MRD sample. As shown in Table I, 10049 events fall into this category. Based on the MC simulation, the selection efficiency and purity for CC interactions in the sample are estimated to be 56.9% and 98.0%, respectively. The expected neutrino energy spectra and the energy dependence of the selection efficiency for CC events are shown in Fig. 4(b) and 4(d), respectively. The total number of CC events is obtained to be $`(1.73\pm 0.02(\mathrm{stat}.))\times 10^4`$. We derive the cross section ratio of CC coherent pion production to the total CC interaction to be $`(0.04\pm 0.29(\mathrm{stat}.))\times 10^2`$.
Systematic uncertainties for the cross section ratio are summarized in Table 2. The major contributions come from uncertainties of nuclear effects and the neutrino interaction models. The uncertainty due to nuclear effects is estimated by varying the cross sections of pion absorption and elastic scattering by $`\pm `$30% based on the accuracy of the reference data Ingram:1983 . The uncertainties in QE and CC1$`\pi `$ interactions are estimated by changing the axial vector mass by $`\pm `$ 0.10 $`\mathrm{GeV}/c^2`$ Bernard:2002 . For DIS, the effect of the Bodek and Yang correction is evaluated by changing the amount of correction by $`\pm `$30%. The $`q_{\mathrm{rec}}^2`$ distribution of the non-QE-proton sample (Fig. 2(c)) indicates an additional deficit of background events in the region $`q_{\mathrm{rec}}^2<0.10(\mathrm{GeV}/c)^2`$. CC1$`\pi `$ interaction dominates events in this region; its cross section has significant uncertainty due to nuclear effects. We estimate the amount of possible deficit in the same manner as described in Aliu:2004 with the one track, QE and non-QE-proton samples. We find that a 20% suppression of CC1$`\pi `$ events for $`q_{\mathrm{true}}^2<0.10(\mathrm{GeV}/c)^2`$ is allowed, which varies the cross section ratio by $`+0.14\times 10^2`$. This variation is conservatively treated as a systematic uncertainty. We also consider the uncertainties of the event selection, where the dominant error comes from track counting, detector response such as scintillator quenching, and neutrino energy spectrum shape. The total systematic uncertainty on the cross section ratio amounts to $`+0.32`$/$`0.35\times 10^2`$.
Our result is consistent with the non-existence of CC coherent pion production at K2K neutrino beam energies, and hence we set an upper limit on the cross section ratio at 90% C.L. :
$`\sigma (\text{CC coherent }\pi )/\sigma (\nu _\mu \mathrm{CC})<0.60\times 10^2.`$
For reference, the total CC cross section is calculated as $`1.07\times 10^{38}\mathrm{cm}^2/\mathrm{nucleon}`$ in the neutrino MC simulation by averaging over K2K neutrino beam energies.
The obtained upper limit is inconsistent with the model prediction by Rein and Sehgal at the level of 2.5 standard deviations. We assign a 35 % uncertainty to the theoretical prediction as described in Rein:1982pf . In addition, a finite cross section was reported by Aachen-Padova group for NC coherent pion production with 2 GeV average neutrino energy and with aluminum target Faissner:1983 . If we assume an $`A^{1/3}`$ dependence of the cross section ($`\sigma `$) and $`\sigma (\mathrm{CC})=2\sigma (\mathrm{NC})`$ according to the model of Rein and Sehgal, the discrepancy between the extrapolation from the NC measurement and the present result is as large as 3 standard deviations. There are other models predicting lower cross sections Paschos:2003hs ; Belkov:1986hn ; Kelkar:1996iv , but they do not provide the kinematics of pions and it is difficult to test them directly. Further theoretical work is necessary to construct interaction models which explain these experimental results. The non-existence of CC coherent pion production has given a solution to the low-$`q^2`$ discrepancy observed in K2K. It also reduces the uncertainty on the cross section in the relevant $`q^2`$ region, which is crucial for the future neutrino oscillation experiments.
In summary, we report on a search for CC coherent pion production by muon neutrinos with a mean energy of 1.3 GeV. The data analyzed correspond to $`1.7\times 10^{19}`$ POT recorded with the K2K-SciBar detector. No evidence of CC coherent pion production is found and an upper limit on the cross section ratio of CC coherent pion production to the total CC interaction is derived to be $`0.60\times 10^2`$ at 90% C.L. This result is the first experimental limit for CC coherent pion production by neutrinos with energies of a few GeV.
###### Acknowledgements.
We thank the KEK and ICRR directorates for their strong support and encouragement. K2K is made possible by the inventiveness and the diligent efforts of the KEK-PS machine group and beam channel group. We gratefully acknowledge the cooperation of the Kamioka Mining and Smelting Company. This work has been supported by the Ministry of Education, Culture, Sports, Science and Technology of the Government of Japan, the Japan Society for Promotion of Science, the U.S. Department of Energy, the Korea Research Foundation, the Korea Science and Engineering Foundation, NSERC Canada and Canada Foundation for Innovation, the Istituto Nazionale di Fisica Nucleare (Italy), the Spanish Ministry of Science and Technology, and Polish KBN grants: 1P03B08227 and 1P03B03826.
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# The One Page Model Checker
## 1 Introduction and Related Work
Debugging multithreaded applications is hard. Race conditions mean that failures may be nondeterministic. Deadlock can be hard to trace, because it involves behaviors from multiple concurrent threads. Tools to prove that a piece of code has no such behaviors can help find such errors, and instill confidence that the programs will work correctly when deployed.
Here we describe a method for measuring the behavior of multithreaded programs through all possible execution interleavings. Our work is straightforward and applicable to many different programming languages, although it also has some significant, fundamental limitations.
Several other techniques have been proposed which relate to model checking multithreaded applications. Visser et al built an optimized system which implements the Java VM and can prove properties about multithreaded applications. Mercer and Jones built an explicit state model checker designed around specific CPUs which verifies properties of compiled code.
In 1997, Savage et al introduced a tool which enforces a locking discipline on resources to prevent nondeterministic behavior caused by OS scheduling. While not a model checker per se, this tool also aims to help programmers reduce the uncertainty associated with multithreaded applications.
## 2 System Overview
The principle behind our system is easy to understand. Given a program with two threads, we wish to search for particular conditions like deadlock among all possible thread interleavings. For example, if the first thread executes $`print(ab)`$ and the second executes $`print(12)`$, then at least the following behaviors are possible: the first thread could run to completion, followed by the second, producing the output string “ab12”, or the second could run to completion first, producing “12ab”. If the operating system chooses to switch between the threads while they’re running, the strings “a1b2”, “a12b”, or “1a2b” might also occur. Of course, if $`print()`$ isn’t thread safe, the program might also exhibit other behaviors or crash entirely.
To test all possible interleavings, the model checker must have a way of trying different execution paths and must keep track of which paths have been explored. Our system’s simplicity and compactness comes from using the Unix-standard $`fork()`$ system call, which forks a process at the point of the function call into two independent, identical processes. $`fork()`$ can be called inside branching, looping or subroutine constructs, and both the original and newly created child processes will return to the following statement and continue as if nothing happened. This is different from most thread implementations, which must be called with a subroutine to execute in the new thread, after which the thread terminates.
$`fork()`$ allows our model checker to be implemented much like a normal recursive depth-first search. In such a search, the program’s stack is used to implicitly keep track of the current progress through the state space being searched. In place of nested function calls, our implementation creates a child process at each branching point to explore the next level of the search. This can happen surprisingly quickly, due to efficient OS techniques like copy-on-write paging which allows efficient forking; our Athlon64 3000+ running Debian GNU/Linux can perform over 7,000 $`fork()`$ operations per second.
Since our model checker operates on programs with two threads, we carefully synchronize pairs of processes to implement the possible execution orders. We call this technique “the buddy system”, and a pair of cooperating threads “buddies”. Each pair of buddies maintains a data structure in shared memory containing semaphores, execution path and other elements required for IPC and coordination.
And since $`fork()`$ creates processes, not threads, we implement threadlike behavior using shared memory. Threads do not get separate copies of program variables as processes do, so we create a structure in a shared memory segment where buddy processes keep all application variables.
### 2.1 Example
We will expand the earlier example into separate statements to show how our technique works. Thread 0 executes $`print(a);print(b)`$, while thread 1 executes $`print(1);print(2)`$. First, we instrument our code by placing calls to a function $`hook()`$ before each statement, then calling $`done()`$ at the end of execution:
```
thread0() {
hook(); print(a);
hook(); print(b);
done();
}
thread1() {
hook(); print(1);
hook(); print(2);
done();
}
```
The code is then compiled and linked with our model checking code. When it executes, $`fork()`$ is called to create a separate process for each thread, each of which executes into the first $`hook()`$. There are only two possible ways in which the threads can execute their first statements; either $`thread0`$ or $`thread1`$ goes first. Each “thread” (really a process) thus forks into two processes, resulting in two parents and two children. The pair of parents become buddies, and the pair of children become buddies. The parent processes each wait for their children to terminate, much as a recursive function calls itself and waits for the recursive call to finish. The child process for thread 1 blocks using a semaphore, waiting for its buddy in thread 0 to execute a single statement. The buddy process does this by returning from the call to $`hook()`$, which allows the first statement, $`print(a)`$, to execute. Then that process hits the next call to $`hook()`$ and signals to its buddy that it has finished executing a statement.
Now the process repeats; either thread0 can execute another statement, or thread1 can execute its first statement. Again, each buddy forks, with the parents waiting for the children to finish. The child process for thread 0 again goes first, returning to execute $`print(b)`$, then calling $`done()`$, which signals to the buddy process that thread 0 has completed execution. With no remaining alternatives, thread 1 now runs to completion, giving a resulting output of “ab12”.
Once the grandchildren of the original two threads have each terminated, the children continue running. Since the grandchildren explored the case in which thread 0 executed another statement, the children explore what happens when thread 1 runs. Thread 1 returns from its hook and executes $`print(1)`$, then signals its buddy. Once again, the children have two alternatives, so they fork another pair of grandchildren. The thread 0 grandchild executes its second statement and terminates, allowing thread 1 to complete and producing the string “a1b2”. The children again execute a statement from thread 1, $`print(2)`$, after which thread 0 runs to completion, producing “a12b”. Now the original two threads can continue, executing $`print(1)`$ from thread 1, and forking another set of children, which fork grandchildren as before, producing the strings “1ab2”, “1a2b”, and “12ab”.
## 3 Code Instrumentation, Language Independence and Statement Atomicity
In order for our system to work properly, application code must be properly instrumented, by making calls to $`hook()`$ before each program statement. These statements are assumed to execute atomically, which does not generally happen in current systems, but which can be assured using a technique we describe later in this section.
This instrumentation process is very simple, and works independently of language constructions like loops, function calls and branches. For example, we first implemented the example we gave in the last section as follows, essentially the same as we listed it before:
```
if(child) {
// thread 0
hook(); printf(‘‘a’’);
hook(); printf(‘‘b’’);
done();
} else {
// thread 1
hook(); printf(‘‘1’’);
hook(); printf(‘‘2’’);
done();
}
```
But later, we generalized it to work for arbitrary strings using a separate function and a loop:
```
void str(char *s) {
int i;
for(i=0; s[i]; i++) {
hook(); b->common.outstr[b->common.outidx++] = s[i];
}
}
main() {
...
if(child) {
// thread 0
str(‘‘ab’’);
done();
} else {
// thread 1
str(‘‘12’’);
done();
}
}
```
A naive, automated instrumentation tool might have added additional $`hook()`$ calls as follows, but that would have merely added overhead to the model checking process:
```
void str(char *s) {
int i;
hook(); for(i=0; s[i]; i++) {
hook(); b->common.outstr[b->common.outidx++] = s[i];
}
}
main() {
...
if(child) {
// thread 0
hook(); str(‘‘ab’’);
done();
} else {
// thread 1
hook(); str(‘‘12’’);
done();
}
}
```
While modifying program source code before model checking is generally deprecated, we feel our technique has several interesting features. First, adding calls to $`hook()`$ before each program statement is easy to do automatically for reasonably written source code, even without constructing a formal parser for the target language. The implementor must simply avoid placing calls where they would cause syntax errors in the program, such as in between function declarations. Redundant calls to $`hook()`$ add overhead, but don’t otherwise break our system. This makes our system straightforward to implement in a variety of languages, whereas traditional systems require significant adaptation to target languages in order to properly model their behavior.
Second, our instrumentation can be used to perform other tasks beside model checking, by changing the behavior of $`hook()`$. For instance, $`hook()`$ could be modified to implement white box testing, in which test cases are constructed which together must execute all code branches.
Third, our instrumentation can be left in place to guarantee the statement-level atomicity assumed by our system. Generally speaking, modern CPUs offer only machine instruction level atomicity – the OS may interrupt a process between any two instructions and begin execution of a different process. Model checkers like Estes work on these machine instructions directly, but their results can only be applied to that particular compilation of the application on a particular CPU. This may prove to be the only way to prove useful thread safety and liveness properties about unmodified code on a particular CPU, and would tend to suggest that dealing with multithreaded applications in their original high-level language form doesn’t even make sense. On the other hand, if calls to $`hook()`$ are left in place in distributed code, the function can be modified to essentially make each application statement into an individual critical section.
Admittedly, this adds a large amount of overhead to the code, since system calls to raise and lower a semaphore must be made for each program statement. But in modern high level scripting languages particularly, programs tend to have fewer statements, with powerful built-in commands having relatively high execution costs. Such languages may be particularly difficult to verify at the machine code level, since they run via large, complex interpreters. As a first approximation of the overhead our technique would add, we wrote a C program which forks into two processes, each of which loops 1,000,000 times. In the loop, each process grabs a semaphore, adds the current index to a variable, then releases the semaphore. The program performed the 4,000,000 semaphore operations in about 0.8 seconds on our Athlon64 3000+. We then ran a program in Perl which loops 2,000,000 times in a single process, likewise adding up the index values. It also took about 0.8 seconds, suggesting that efficient C-based hooks in the perl code to ensure statement-level atomicity might add only 50% overhead to such a program, and possibly less for a program using fewer, more costly operations than simple loops and additions.
## 4 Formal Definitions
Here we define terminology used in the algorithms described in the next sections.
* Assume there exist functions $`hook()`$, which is called before execution of each application thread program statement, and $`done()`$, which is called after the last statement in each thread. Application threads may include most usual language features, such as branching, looping and function calls (see section 3).
* We define the execution counter for a thread to be the number of times $`hook()`$ has been called since the beginning of the thread’s execution. Intuitively, this is the number of statements executed in that thread, plus one.
* If $`s_0`$ is the value of the execution counter for thread 0 and $`s_1`$ is the corresponding counter for thread 1, the pair $`s_0,s_1`$ forms the combined execution counter for the two threads.
* An execution trace is defined to be a string $`t\{0,1\}^{}`$ which represents the order in which statements from the two threads were executed to reach a particular combined execution counter. In our earlier example, the execution trace $`0011`$ corresponds to the output “ab12”.
* Let $`V_0`$ and $`V_1`$ represent the vectors of shared variable values for threads 0 and 1 having an execution trace $`t`$ and combined execution counter $`C`$. The tuple $`I=V_0,V_1,t,C`$ is a partial interleaving for the two threads (partial, since the threads may not yet have run to completion).
## 5 Search Algorithm
Our algorithm for performing a depth first search on all possible execution interleavings of two threads $`t0`$ and $`t1`$ is as follows:
* Base case: Let $`I`$ be the initial partial interleaving for threads $`t0,t1`$, representing the program state at the first call to $`hook()`$ in each thread, before any application statements have executed.
* Recursion: Given a partial interleaving $`I`$ for threads $`t0,t1`$,
+ Run any user-supplied code for checking conditions.
+ If $`done()`$ has been called in both threads, terminate the current thread and indicate successful program execution for a single complete interleaving.
+ If $`done()`$ has been called in only one thread, allow the other thread to continue to completion.
+ Otherwise, fork both threads to create children $`t0^{},t1^{}`$. Parents both wait for termination of the children. $`t0^{}`$ returns from $`hook()`$, allowing a single statement to execute while $`t1^{}`$ blocks. Then the recursive step is performed again on $`t0^{},t1^{}`$ with a new partial interleaving $`I^{}`$. When $`t0^{},t1^{}`$ have terminated, $`t1`$ returns, executing a single statement, and then the recursive step is performed again for $`t0,t1`$ with new partial interleaving $`I^{\prime \prime }`$.
Omitting $`\mathrm{\#}include`$ statements and helper functions for setting up semaphores and shared memory, our C implementation of this algorithm fits in one printed page of code (80 columns by 65 lines).
## 6 Detecting Race Conditions and Pruning the Search Space
Note that in the above example, program output is entirely dependant on the order in which the OS schedules the two threads for execution. Such nondeterministic behavior is almost never intended, and usually represents a bug in the code.
Consequently, we provide a technique for verifying that threads behave the same regardless of the order in which they are executed. This technique also makes it easy to avoid unnecessary exploration of the space of possible interleavings. Formally, we define a race condition as follows:
* Let $`I=V_0,V_1,t,C`$ and $`I^{}=V_0^{},V_1^{},t^{},C^{}`$ represent partial interleavings such that $`C=C^{}`$. That is, $`I`$ and $`I^{}`$ are partial interleavings which have reached the same execution counter in each thread but potentially through a different order of execution. $`I`$ and $`I^{}`$ form a race condition iff $`II^{}`$.
To implement this technique, we maintain a shared table keyed on the combined execution counters explored while searching the state space. Each table entry records the partial interleaving at that combined execution counter. When a particular combined execution counter is reached via a different execution trace, the current partial interleaving is compared against the stored partial interleaving. If they differ, the two partial interleavings are displayed as examples of execution paths capable of producing differing behaviors. A program with no race conditions will of course display only the single possible outcome of program execution.
The second purpose of this table is to record which combined execution counters have been reached before. Since our algorithm performs a depth first search on the possible thread interleavings, the second occurrance of a combined execution counter can only occur once all the remaining interleavings from that point on have already been explored. Consequently, if a race condition does not occur at a particular explored partial interleaving, there is no need to explore it again since the two threads are in exactly the same state as they were the last time.
Although we have not yet been able to derive a formula for the complexity of our pruning algorithm, it is clear that this pruning technique is at least an order of magnitude improvement over an exhaustive search. We ran our algorithm on pairs of threads each executing 3 to 8 statements. Values represent the total number of calls to $`hook()`$, which roughly corresponds to the number of states explored.
| Technique | 3 | 4 | 5 | 6 | 7 | 8 |
| --- | --- | --- | --- | --- | --- | --- |
| Exhaustive | 30 | 112 | 420 | 1584 | 6006 | 22876 |
| Pruning | 18 | 32 | 50 | 72 | 98 | 128 |
This addition was surprisingly easy to add to our system; it required less than a page of code in changes.
## 7 Supporting IPC
Our implementation supports multiple semaphores which may be used by the user application for interprocess communication. This complicates our system, since a thread may block until the other thread releases a particular semaphore, and complicates the actual implementation even more, due to practical issues regarding process cleanup, IPC and resource management. Here we give the algorithm for supporting an arbitrary number of application semaphores.
* Let the definition of a partial interleaving be extended to include a set of semaphores $`S=\{s_0..s_n\}`$, which may be up or down.
* Let the function $`down(i)`$ be a valid application statement (to be preceded by a call to $`hook()`$). $`down(i)`$ causes the current thread to lower $`s_i`$ if it is up, and do nothing otherwise.
* Let the function $`up(i)`$ perform the complimentary operation, with the addition that if $`s_i`$ is already up, the thread blocks until the other thread calls $`down(i)`$. If the other thread is already blocking, report deadlock and terminate both threads.
* Base case: Let $`I`$ be the initial partial interleaving for threads $`t0,t1`$, representing the program state at the first call to $`hook()`$ in each thread, before any application statements have executed.
* Recursion: Given a partial interleaving $`I`$ for threads $`t0,t1`$,
+ Run any user-supplied code for checking conditions.
+ If $`done()`$ has been called in both threads, terminate the current thread and indicate successful program execution for a single complete interleaving.
+ If $`done()`$ has been called in only one thread, allow the other thread to continue to completion. If the other thread is blocked, report an error, since the thread will block forever.
+ If a thread is blocked, allow the other thread to execute another statement.
+ Otherwise, fork both threads to create children $`t0^{},t1^{}`$. Parents both wait for termination of the children. $`t0^{}`$ returns from $`hook()`$, allowing a single statement to execute while $`t1^{}`$ blocks. Then the recursive step is performed again on $`t0^{},t1^{}`$ with a new partial interleaving $`I^{}`$. When $`t0^{},t1^{}`$ have terminated, $`t1`$ returns, executing a single statement, and then the recursive step is performed again for $`t0,t1`$ with new partial interleaving $`I^{\prime \prime }`$.
## 8 Limitations
In most thread implementations, threads share all global variables, but each has its own stack. Local variables are thus maintained independantly from other threads. Our implementation presently provides no support for such variables, and assumes that all thread state can be monitored via the shared variables and execution counter. It would be easy to create a second structure associated with each buddy process for storing variables unique to each thread, and account for that additional state information when checking for race conditions and pruning the state space.
As we described in section 3, our assumption that program statements are performed atomically is not at all guaranteed by real computers, unless our technique is employed at a machine code level. To achieve reported results in practice, statement-level atomicity would need to be enforced by the operating system, language interpreter, or by using a modified $`hook()`$ as we described.
While using $`fork()`$ to store program execution state makes our system very simple to implement, it imposes a significant amount of system overhead. The system resources for two processes, including two process table entries, are required for each level of depth in the search space, which corresponds to the number of statements executed by the combined threads. This makes even simple programs, like a pair of threads which each loops 1,000,000 times, impossible to verify with our system.
Finally, our current system is limited to programs with two threads. See section 10 for discussion on removing this limitation.
## 9 Implementation and Performance
As we showed in section 6, our pruning algorithm performs far better than an exhaustive search. Default (though modifiable) OS limitations on the number of available semaphore sets and our implementation’s inefficient use of single semaphore sets rather than multiple semaphores in multiple sets limits us to running applications which execute a total of about 22 steps. These it handles in under 0.1 seconds. The resulting maximum of 44 live processes imposes no noticeable memory consumption.
While our fundamental search algorithm can be implemented in about a page of code, our full system supporting application semaphores, race detection and pruning, and with helper functions, debugging code, and whitespace currently weighs in at 611 lines of C. The implementation requires 4 semaphores for each level of depth in the DFS, and requires $`n^3`$ storage in the DFS depth to maintain the table of partial interleavings complete with program execution paths. For our current depth limit of 22, this amounts to about 100k of memory.
## 10 Future Work
Our system is limited to programs with two threads. Since our system is modeled after the traditional recursive depth first search algorithm, there is a clear path for extending our algorithm to support any number of threads. Rather than pairs of parent processes spawning a single pair of children, each parent in a set of n parent buddies would iteratively spawn n-1 children (waiting each time for the previous child to die). The n n-member buddy sets (cliques?) would then be dispatched, exploring the paths in which each child executes its next statement. Instead of a maximum $`2k`$ live processes for a k-level DFS, up to $`nk`$ processes would exist.
Implementation, in our experience, might be time consuming, due to the inherent difficulty humans seem to have keeping track of multiprocess systems. However, with careful planning, and for programmers more experienced with multiprocess applications, this extension should not prove too difficult.
One feature that might be quite easy to add is the ability for parents to run without waiting for their children. Unchecked, this would act like a “fork bomb”, potentially swamping the system as the entire search tree unfolded at once. But with a limit on how many pairs of processes could run at once, our system would immediately be able to run on SMP systems with arbitrary numbers of processors, limited only by the size of the process table and the system memory.
At the other extreme, with only two processes running at once, it’s unnecessary to allocate each pair of buddies their own set of semaphores, as our first implementation does. With more careful use and management, we suspect that a single set would suffice, avoiding system limits on available semaphores. The state table also need not require so much storage; cryptographic hash functions can be used to reduce arbitrary amounts of data to a 128 bit digest which can be stored in place of the actual partial interleaving. An expected $`2^{64}`$ entries would have to be made before any pair of entries would have the same digest, which is far more than any PC can store today.
As we described in section 3, our system should be easy to implement in a variety of languages, possibly just by using cross-language extensions to call our original C implementation. There are also other uses for $`hook()`$ which we described but have not implemented.
## 11 Conclusion
We gave a straightforward technique for checking thread-related properties of programs with two threads. Our technique is general-purpose and language-independent, and may be particularly suited to modern high-level scripting languages due to the difficulty of machine-code level model checking for these languages and the overhead required to enforce the statement-level atomicity required by our system.
Our system performs much better than an exhaustive search of the state space, understands application semaphore use, and detects deadlock and race conditions. Furthermore, our implementation is quite compact; our exhaustive search algorithm fits in one printed page of C code, and our complete implementation fits in under ten.
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# Minkowski Vacuum Stress Tensor Fluctuations
## I Introduction
Because physically realizable states in quantum field theory are not eigenstates of the stress tensor operator, quantum stress tensor fluctuations are a universal feature of quantum fields. These fluctuations can have physical effects, including Casimir force fluctuations Barton ; Eberlein ; JR ; WKF01 , radiation pressure fluctuations WF01 , and passive fluctuations of the gravitational field F82 ; DF88 ; Kuo ; PH97 ; WF99 ; CH95 ; CCV ; Moffat ; MV99 ; HS98 ; BFP00 ; PH00 ; FW03 ; HRV04 ; BF ; BF2 ; Winitzki01 ; Vachaspati03 ; DV05 . Passive fluctuations of gravity are those driven by fluctuations of the matter field stress tensor, as opposed to the active fluctuations due to the quantum nature of gravity itself. The quantum stress tensor correlation function is singular in the limit of coincident points. However, this does not prevent us from obtaining physically meaningful results for observable quantities, such as the luminosity fluctuations of a distant source seen through the fluctuating spacetime BF . These observables are expressed as spacetime integrals of the correlation function, which can be defined by an integration by parts procedure. Alternatively, one could use other approaches, such as dimensional regularization FW04 .
In general, the stress tensor correlation function can be decomposed into three terms: a “fully normal ordered” term which is state dependent, but free of singularities, a vacuum term which is singular, but state independent, and a “cross term” which is both singular and state dependent. In many situations, one is interested in state dependent effects, so the vacuum term can be ignored. For example, radiation pressure fluctuations in a coherent state arise solely from the cross term WF01 . However, this does not mean that the vacuum term is devoid of any physical content.
The main purpose of this paper is the search for such content. Here we will be concerned with a free, massless scalar field in Minkowski spacetime, and its stress tensor correlation function in the Minkowski vacuum state. In a previous paper FR04 , we studied the subtle stress tensor correlations in non-vacuum states created by moving mirrors in two-dimensional flat spacetime. One of the key results of the present paper will be the derivation of a covariant expression for the correlation function as a sum of total derivative terms. This expression will be given in Sect. II.1 for two dimensions and in Sect. III.1 for four dimensions, with the details of the derivations presented in Appendices A and B, respectively. We will discuss spacetime averages of the energy density correlation function in Sects. II.2 and III.2, and averages along a worldline in Sects. II.3 and III.3. The results will be summarized and discussed in Sect. IV. Units in which $`\mathrm{}=c=1`$, and a spacelike metric signature will be used throughout this paper.
## II Two Dimensions
### II.1 Covariant Stress Tensor Correlation Function
We will be concerned with the stress tensor correlation function
$$C_{\mu \nu \alpha \beta }(x,x^{})=:T_{\mu \nu }(x)::T_{\alpha \beta }(x^{}):$$
(1)
for a massless, minimally coupled scalar field in two-dimensional Minkowski spacetime in the vacuum state. Here $`:T_{\mu \nu }(x):`$ is the normal ordered stress tensor operator, so $`:T_{\mu \nu }(x):=0`$. We especially seek an expression for $`C_{\mu \nu \alpha \beta }(x,x^{})`$ as a sum of terms, each of which is a total derivative of a function with at most logarithmic singularities as $`x^{}x`$. This will allow us to define integrals of the correlation function by integration by parts.
Such a form is derived in Appendix A, where it is shown that
$`C_{\mu \nu \alpha \beta }(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{384\pi ^2}}[8_\mu _\nu _\alpha _\beta f_12g_{\mu \nu }g_{\alpha \beta }\mathrm{}\mathrm{}f_2`$ (2)
$`+`$ $`(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })\mathrm{}\mathrm{}f_2+2(g_{\alpha \beta }_\mu _\nu +g_{\mu \nu }_\alpha _\beta )\mathrm{}f_2`$
$``$ $`(g_{\alpha \nu }_\mu _\beta +g_{\alpha \mu }_\nu _\beta +g_{\beta \nu }_\mu _\alpha +g_{\beta \mu }_\nu _\alpha )\mathrm{}f_2],`$
where
$$f_1=\mathrm{ln}(\mathrm{\Delta }x^2/\mathrm{}^2),$$
(3)
and
$$f_2=\mathrm{ln}^2(\mathrm{\Delta }x^2/\mathrm{}^2),$$
(4)
where $`\mathrm{}`$ is an arbitrary constant with dimensions of length. The correlation function is independent of the choice of $`\mathrm{}`$. Here $`\mathrm{}=^\mu _\mu `$ is the wave operator, and $`\mathrm{\Delta }x^2=(x^\mu x_{}^{}{}_{}{}^{\mu })(x_\mu x_{}^{}{}_{\mu }{}^{})`$. Because $`_\mu f_1=f_1/x^\mu =_\mu ^{}f_1=f_1/x_{}^{}{}_{}{}^{\mu }`$, the correlation function, Eq. (2), can be written in several equivalent forms.
The energy density correlation function becomes
$$C(x,x^{})=C_{ttt^{}t^{}}=\frac{1}{48\pi ^2}_t^4f_1=\frac{1}{48\pi ^2}_t^2_t^{}^2f_1.$$
(5)
Note that none of the $`f_2`$ terms contribute in this case. This expression allows us to compute the mean squared average energy density. Let $`g(t)`$ be a time sampling function, and $`h(x)`$ be a spatial sampling function. Then we define the averaged energy density operator as
$$\overline{\rho }=𝑑tg(t)𝑑xh(x):T_{tt}:.$$
(6)
The mean square of this operator is
$$\widehat{C}=\overline{\rho }^2=𝑑tg(t)𝑑xh(x)𝑑t^{}g(t^{})𝑑x^{}h(x^{})C(x,x^{}).$$
(7)
If we insert Eq. (5) into the above expression, and then integrate by parts, we can write
$$\widehat{C}=\frac{1}{48\pi ^2}𝑑t\ddot{g}(t)𝑑t^{}\ddot{g}(t^{})𝑑xh(x)𝑑x^{}h(x^{})f_1.$$
(8)
In the limit that the width of the spatial sampling function goes to zero, $`h(x)\delta (x)`$ and we obtain
$$\widehat{C}=\frac{1}{48\pi ^2}𝑑t\ddot{g}(t)𝑑t^{}\ddot{g}(t^{})\mathrm{ln}[(\mathrm{\Delta }t)^2/\mathrm{}^2].$$
(9)
### II.2 Averaging over Space and Time - 2D
Rather than using Eq. (9), in some cases we can also directly evaluate the integral in Eq. (7) using contour integration methods. For the explicit examples to be treated in this paper, the latter approach is more convenient. The energy density correlation function, Eq. (5), can be expressed as
$$C(x,x^{})=\frac{(\mathrm{\Delta }t^2+\mathrm{\Delta }x^2)^2+4\mathrm{\Delta }t^2\mathrm{\Delta }x^2}{4\pi ^2(\mathrm{\Delta }t^2\mathrm{\Delta }x^2)^4},$$
(10)
where $`\mathrm{\Delta }t=tt^{}`$, and $`\mathrm{\Delta }x=xx^{}`$. In this subsection we will sample this correlation function in both space and time with Lorentzian functions of width $`\alpha `$ in $`t`$ and $`t^{}`$, and $`\beta `$ in $`x`$ and $`x^{}`$. Further, let the spatial sampling functions coincide, but let the temporal ones be displaced by $`t_0`$.
Let
$$\widehat{C}(t_0)=_{\mathrm{}}^{\mathrm{}}𝑑tg_L(\alpha ,t+t_0)_{\mathrm{}}^{\mathrm{}}𝑑t^{}g_L(\alpha ,t^{})_{\mathrm{}}^{\mathrm{}}𝑑xg_L(\beta ,x)_{\mathrm{}}^{\mathrm{}}𝑑x^{}g_L(\beta ,x^{})C(x,x^{})$$
(11)
where
$$g_L(\alpha ,t)=\frac{\alpha }{\pi (t^2+\alpha ^2)},$$
(12)
and
$$_{\mathrm{}}^{\mathrm{}}𝑑tg_L(\alpha ,t)=1.$$
(13)
Now let $`ttt_0`$, so that we have
$`\widehat{C}(t_0)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑tg_L(\alpha ,t){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}g_L(\alpha ,t^{}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑xg_L(\beta ,x)`$ (14)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}dx^{}g_L(\beta ,x^{})C(tt^{}t_0,xx^{})`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau g_L(a,\tau ){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\rho g_L(b,\rho )C(\tau t_0,\rho ),`$
where $`a=2\alpha ,b=2\beta ,\tau =tt^{}`$ and $`\rho =xx^{}`$. In the last step, we used the identity
$$_{\mathrm{}}^{\mathrm{}}𝑑tg_L(\alpha ,t)_{\mathrm{}}^{\mathrm{}}𝑑t^{}g_L(\alpha ,t^{})F(tt^{})=_{\mathrm{}}^{\mathrm{}}𝑑\tau g_L(a,\tau )F(\tau ).$$
(15)
We may do the integral on $`\rho `$ first, by contour integration. The integrand has simple poles at $`\rho =\pm ib`$ and fourth order poles at $`\rho =\pm (\tau t_0)`$. We choose a contour in the upper half-plane which avoids the fourth order poles, the contour $`C_1`$ in Fig. 1. In fact, we could use other contours such as $`C_2`$ and still obtain the same answer. Even if we chose a contour which enclosed either of the fourth order poles, our answer for the real part of the integral would still be the same. This is because the contribution of either of these poles, the result of integrating around the closed circular paths, is pure imaginary. Note that the straight segments and the semicircular segments of $`C_1`$ each contain real terms which diverge as the radii of the semicircles go to zero. However, these terms cancel when the straight and semicircular contributions are added. The divergent terms on the straight segments are the boundary terms that would arise from integrating by parts along these segments only. Thus integration by parts along the straight segments and discarding the boundary terms produces the same result as integration along the complete contour.
In any case, using the residue theorem we obtain
$$_{\mathrm{}}^{\mathrm{}}𝑑\rho g_L(b,\rho )C(\tau t_0,\rho )=\frac{[(\tau t_0)^2b^2]^24b^2(\tau t_0)^2}{4\pi ^2[(\tau t_0)^2+b^2]^4}.$$
(16)
The subsequent $`\tau `$-integration was performed and yields
$$\widehat{C}(t_0)=\frac{[t_0(t_0+2a+2b)(a+b)^2][t_0(t_02a2b)(a+b)^2]}{4\pi ^2[t_{0}^{}{}_{}{}^{2}+(a+b)^2]^4}.$$
(17)
(This and several other calculations in this paper were done using the public domain algebraic manipulation program MAXIMA.) In the special case when $`t_0=0`$, we simply have
$$\widehat{C}(0)=\frac{1}{4\pi ^2(a+b)^4}.$$
(18)
Let us define
$$K(t_0,a,b)=\frac{\widehat{C}(t_0)}{\widehat{C}(0)}.$$
(19)
In general, we have that $`\widehat{C}(t_0)=\widehat{C}(t_0)`$. From Eqs. (17) and (19), we find that
$$_{\mathrm{}}^{\mathrm{}}K(t_0,a,b)𝑑t_0=0,$$
(20)
and similarly
$$_0^{\mathrm{}}K(t_0,a,b)𝑑t_0=0.$$
(21)
This result tells us that positively correlated regions ($`K>0`$), and anticorrelated regions ($`K<0`$) have equal weight.
### II.3 Sampling along a Worldline - 2D
In this subsection, we shall specialize to the case of sampling along a worldline, i.e., we will effectively set the width of the spatial sampling function to zero. Define a normal-ordered smeared stress tensor operator by
$$S(t_0)=_{\mathrm{}}^{\mathrm{}}𝑑tg(t,t_0):T_{tt}(t):,$$
(22)
where $`g(t,t_0)`$ is a sampling function whose peak is at $`t=t_0`$. Although $`S=0`$ in the vacuum state, $`S^20`$. From Eq. (9), we have that
$`S^2`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}g(t,t_0)g(t^{},t_0)C(t,t^{})`$ (23)
$`=`$ $`{\displaystyle \frac{1}{48\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t\ddot{g}(t,t_0){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}\ddot{g}(t^{},t_0)\mathrm{ln}[\mathrm{\Delta }t^2/\mathrm{}^2].`$
The case we want to consider is two regions of time-sampled energy density which are allowed to initially coincide but which are then gradually separated from one another. One sampling function has its peak at $`t^{}=0`$ and the other at $`t=t_0`$. We want to imagine sliding these regions away from one another (see Fig. 2), and examine the behavior of the vacuum correlation function as we vary $`t_0`$.
With Eq. (22), we can write
$`S(t_0)S(0)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}g(t,t_0)g(t^{},0)C(t,t^{})`$ (24)
$`=`$ $`{\displaystyle \frac{1}{48\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}\ddot{g}(t,t_0)\ddot{g}(t^{},0)\mathrm{ln}[\mathrm{\Delta }t^2/\mathrm{}^2].`$ (25)
This represents the smeared energy density correlation function for two displaced regions along a worldline. We can normalize this quantity by defining
$$K(t_0)=\frac{S(t_0)S(0)}{S^2(0)}.$$
(26)
As an example, we take the sampling function to be a Lorentzian. If we set $`b=0`$ and $`a=1`$ in Eq. (17), we find
$$K(t_0)=\frac{(t_0^22t_01)(t_0^2+2t_01)}{(t_0^2+1)^4}=\frac{(16t_0^2+t_0^4)}{(1+t_0^2)^4}.$$
(27)
The choice of $`b=0`$ corresponds to sampling in time only, with displaced sampling functions. A plot of this function appears in Fig. 3(a). The plot is somewhat deceiving because there is actually a second positive peak which, on the scale of the plot, is too small to be seen. However, it must be there since $`\widehat{C}(t_0)1/(4\pi ^2t_0^4)`$, as $`t_0\mathrm{}`$, and hence $`K(t_0)`$ has to approach $`0`$ from above for large $`t_0`$. The magnified view in Fig. 3(b) reveals the second positive peak. We can also see this by computing the extrema of $`K(t_0)`$ using
$$K^{}(t_0)=\frac{4t_0(t_0^410t_0^2+5)}{(t_0^2+1)^5}.$$
(28)
One finds that $`K^{}(t_0)=0`$ at: $`t_0=0`$ (first maximum), $`t_00.73`$ (minimum), and $`t_03.1`$ (second maximum).
As a second example, consider a compactly supported sampling function of width $`a`$ with $`g=\dot{g}=0`$ at $`t=t_0\pm a/2`$ . A simple choice of function which has this form is
$$g(t,t_0)=g(tt_0)=\frac{30}{a^5}(tt_0a/2)^2(tt_0+a/2)^2.$$
(29)
The second derivative of this function is
$$\ddot{g}(tt_0)=\frac{30[12(tt_0)^2a^2]}{a^5},$$
(30)
and
$$S^2(0)=\frac{25}{2\pi ^2a^4}.$$
(31)
Using Eqs. (25), (26), (30), and (31), one can evaluate $`K(t_0)`$, which is plotted as a function of $`t_0`$ in Fig. 4. Note that the number of maxima and minima of $`K(t_0)`$ for the compactly supported sampling function, given in Eq. (29), is the same as for the Lorentzian sampling function shown earlier. However, for the compactly supported sampling function case, the second maximum is much more pronounced. A calculation also shows that for both the Lorentzian and the compactly supported sampling functions, we have that
$$_0^{\mathrm{}}K(t_0)𝑑t_0=0.$$
(32)
We will show in Appendix C that this is true for arbitrary smooth sampling functions. In this appendix, we also prove that
$$S^2(0)>0.$$
(33)
This establishes that the behavior illustrated in Fig. 4 is independent of the details of the sampling function. The fact that $`S^2(0)>0`$ implies that nearly overlapping regions are positively correlated with one another. As $`t_0`$ increases, the correlation is replaced by anticorrelation, as shown by the negative minimum in $`K(t_0)`$. This anticorrelation implies that if we measure positive energy in a given region, there must be negative energy found nearby. Finally, when the regions are sufficiently separated, the positive correlation returns, as evidenced by the final positive peak in Fig. 4. One can understand why disjoint regions must be positively correlated from the fact that $`C(x,x^{})>0`$. When $`xx^{}`$ everywhere in the range of integration, then the integral for $`\widehat{C}`$ is well defined as an ordinary integral, and must be positive. On the other hand, when we must integrate through points where $`x=x^{}`$, then $`C(x,x^{})`$ becomes defined only as a distribution, and the integration by parts procedure can produce a negative result.
## III Four Dimensions
### III.1 Covariant Stress Tensor Correlation Function
In this section, we consider the vacuum stress tensor correlation function in four dimensions. The covariant form of this function is derived in Appendix B, with the result
$`C_{\mu \nu \alpha \beta }(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{61440\pi ^4}}[8_\mu _\nu _\alpha _\beta \mathrm{}\mathrm{}f_2+6g_{\mu \nu }g_{\alpha \beta }\mathrm{}^4f_2`$ (34)
$`+`$ $`(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })\mathrm{}^4f_26(g_{\alpha \beta }_\mu _\nu +g_{\mu \nu }_\alpha _\beta )\mathrm{}^3f_2`$
$``$ $`(g_{\alpha \nu }_\mu _\beta +g_{\alpha \mu }_\nu _\beta +g_{\beta \nu }_\mu _\alpha +g_{\beta \mu }_\nu _\alpha )\mathrm{}^3f_2].`$
Note that only the function $`f_2`$, defined in Eq. (4), appears here, in contrast to the two-dimensional result. The energy density correlation function in four dimensions is given by
$$C(x,x^{})=C_{ttt^{}t^{}}=\frac{1}{7680\pi ^4}(^2)^2\mathrm{}^2f_2=\frac{1}{7680\pi ^4}^2_{}^{}{}_{}{}^{2}\mathrm{}\mathrm{}^{}f_2,$$
(35)
where $`^2=\mathrm{}+_t^2`$ is the three-dimensional Laplacian operator. This form may be used to compute the mean squared average energy density over a spacetime region defined by a sampling function $`F(x)`$. If we define
$$\overline{\rho }=d^4xF(x):T_{tt}:,$$
(36)
then
$$\widehat{C}=\overline{\rho }^2=d^4xF(x)d^4x^{}F(x^{})C(x,x^{}).$$
(37)
After an integration by parts, this may be expressed as
$$\widehat{C}=\frac{1}{7680\pi ^4}d^4x^2\mathrm{}F(x)d^4x^{}_{}^{}{}_{}{}^{2}\mathrm{}^{}F(x^{})f_2(xx^{}).$$
(38)
At first sight, the process of obtaining finite spacetime averages of the correlation function may seem mysterious. We start with an expression for $`C_{\mu \nu \alpha \beta }(x,x^{})`$ which diverges as $`(xx^{})^8`$ as $`x^{}x`$, which seems to be a nonintegrable singularity. Yet we nonetheless obtain finite integrals of this expression. The reason that this is possible is that although $`C_{\mu \nu \alpha \beta }(x,x^{})`$ is singular as a function, it is a well-defined distribution. This is shown by the existence of the expression, Eq. (34), where $`C_{\mu \nu \alpha \beta }(x,x^{})`$ is expressed as a sum of derivatives of a function with no more than logarithmic singularities. An alternative treatment of the singularities of stress tensor correlation functions was given in Ref. FW04 . There dimensional regularization was used to render the correlation functions finite. In the limit in which $`n4`$, where $`n`$ is the spacetime dimension, time-ordered stress tensor correlation functions possess a pole term, which can be absorbed in a renormalization involving $`R^2`$ and $`R_{\mu \nu }R^{\mu \nu }`$ counterterms in the gravitational action. However, the correlation functions without time ordering, such as $`C_{\mu \nu \alpha \beta }(x,x^{})`$, have no pole term and are hence finite in dimensional regularization in the $`n4`$ limit. This is another way to understand why $`C_{\mu \nu \alpha \beta }(x,x^{})`$ is a well-defined distribution, and why the integration by parts method yields finite results.
### III.2 Averaging over Space and Time - 4D
Here we will perform a calculation analogous to that in Sect. II.2, except involving averaging over space and time in four dimensions. The energy density correlation function, Eq. (35), may be expressed as
$$C(x,x^{})=\frac{(\tau ^2+3r^2)(3\tau ^2+r^2)}{2\pi ^4(\tau ^2r^2)^6},$$
(39)
where $`\tau =tt^{}`$ and $`r=|𝐱𝐱^{}|`$. As before, we use Lorentzian sampling functions of width $`\alpha `$ in $`t`$ and in $`t^{}`$. The time-averaged correlation function is
$$\widehat{C}_T=_{\mathrm{}}^{\mathrm{}}𝑑tg_L(\alpha ,t)_{\mathrm{}}^{\mathrm{}}𝑑t^{}g_L(\alpha ,t^{})C(x,x^{})=_{\mathrm{}}^{\mathrm{}}𝑑\tau g_L(a,\tau )C(x,x^{}),$$
(40)
where $`a=2\alpha `$. The integrand in the $`\tau `$ integral has first order poles at $`\tau =\pm ia`$ and sixth order poles at $`\tau =\pm r`$. The integral may be performed by contour integration in a way analogous to the integral in Eq. (16). The result is
$$\widehat{C}_T=\frac{(3r^2a^2)(r^23a^2)}{2\pi ^4(r^2+a^2)^6}.$$
(41)
Next we wish to average $`\widehat{C}_T`$ over the spatial directions. Here it will be convenient to use a Gaussian sampling function
$$g_G(\beta ,x)=\frac{1}{\sqrt{\pi }\beta }\mathrm{e}^{x^2/\beta ^2},$$
(42)
in each of the Cartesian space coordinates, $`x,y,z,x^{},y^{},z^{}`$, and define the spacetime average as
$`\widehat{C}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑xg_G(\beta ,x){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑yg_G(\beta ,y){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑zg_G(\beta ,z)`$ (43)
$`\times `$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x^{}g_G(\beta ,x^{}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑y^{}g_G(\beta ,y^{}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑z^{}g_G(\beta ,z^{})\widehat{C}_T.`$
We may use the fact that
$$_{\mathrm{}}^{\mathrm{}}𝑑xg_G(\beta ,x)_{\mathrm{}}^{\mathrm{}}𝑑x^{}g_G(\beta ,x^{})f(xx^{})=_{\mathrm{}}^{\mathrm{}}𝑑\mathrm{\Delta }xg_G(b,\mathrm{\Delta }x)f(\mathrm{\Delta }x),$$
(44)
where $`\mathrm{\Delta }x=xx^{}`$ and $`b=\sqrt{2}\beta `$. This leads to
$`\widehat{C}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mathrm{\Delta }xg_G(b,\mathrm{\Delta }x){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mathrm{\Delta }yg_G(b,\mathrm{\Delta }y){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mathrm{\Delta }zg_G(b,\mathrm{\Delta }z)\widehat{C}_T`$ (45)
$`=`$ $`{\displaystyle \frac{4}{\sqrt{\pi }b^3}}{\displaystyle _0^{\mathrm{}}}𝑑rr^2\mathrm{e}^{r^2/\beta ^2}\widehat{C}_T,`$
where $`r^2=(\mathrm{\Delta }x)^2+(\mathrm{\Delta }y)^2+(\mathrm{\Delta }z)^2`$. If we use Eq. (41), then we can write the spacetime averaged correlation function as
$$\widehat{C}=\frac{2}{\pi ^{9/2}b^3}_0^{\mathrm{}}𝑑rr^2\frac{(3r^2a^2)(r^23a^2)}{(r^2+a^2)^6}\mathrm{e}^{r^2/\beta ^2}.$$
(46)
The integral in the above expression may evaluated in terms of the error function, erf, as
$`\widehat{C}`$ $`=`$ $`{\displaystyle \frac{1}{15\pi ^4ab^{13}}}\{\sqrt{\pi }[1\mathrm{erf}\left({\displaystyle \frac{a}{b}}\right)]\mathrm{e}^{a^2/b^2}(15b^6+90a^2b^4+60a^4b^2+8a^6)`$ (47)
$``$ $`2ab(3b^2+2a^2)(11b^2+2a^2)\}.`$
Now we wish to discuss the limits in which one sampling length scale is small compared to the other. First consider the case of a small spatial scale, $`ba`$. The exponential factor in Eq. (46) guarantees that only values of $`rb`$ contribute. Thus we can assume that $`ra`$ in the integrand and write
$$\frac{(3r^2a^2)(r^23a^2)}{(r^2+a^2)^6}\frac{3}{a^8}.$$
(48)
Then we have
$$\widehat{C}\frac{3}{2\pi ^4a^8}$$
(49)
when $`ba`$. This shows that only temporal sampling is necessary in order for $`\widehat{C}`$ to be finite. Equation (49) may also be derived from the explicit form, Eq. (47), by use of the asymptotic form of the error function for large argument.
Next we consider the opposite limit, where $`ab`$. However, $`\widehat{C}\mathrm{}`$ as $`a0`$ for fixed, nonzero $`b`$. This may be seen from the integral, Eq. (46), which becomes proportional to $`_0^{\mathrm{}}𝑑rr^6\mathrm{e}^{r^2/\beta ^2}`$ as $`a0`$. Alternatively, we can expand Eq. (47) for small $`a`$ and show that
$$\widehat{C}\frac{1}{\pi ^{7/2}ab^7},\mathrm{as}a0.$$
(50)
Thus in four dimensions, averaging over space alone is not sufficient to lead to a finite mean squared energy density. This result was obtained previously by Guth Guth and by Roura Roura .
### III.3 Sampling along a Worldline - 4D
In the previous subsection, we found that it is possible to take the limit of a vanishing spatial sampling scale, so that one is sampling along a worldline. Here we will consider that limit for displaced temporal sampling functions. First consider Lorentzian sampling functions and let
$$\widehat{C}(t_0,r)=_{\mathrm{}}^{\mathrm{}}𝑑tg_L(\alpha ,t+t_0)_{\mathrm{}}^{\mathrm{}}𝑑t^{}g_L(\alpha ,t^{})C(\tau ,r)=_{\mathrm{}}^{\mathrm{}}𝑑\tau g_L(a,\tau t_0)C(\tau ,r),$$
(51)
where $`C(\tau ,r)`$ is given by Eq. (39), and $`a=2\alpha `$. If we were to sample in space with a function whose width is small compared to $`a`$, the result is the same as setting $`r=0`$ in the above expression. More precisely, we perform the integral for nonzero $`r`$, using the same method as used to obtain Eq. (16), and then take the $`r0`$ limit. The result is
$$\widehat{C}(t_0,0)=\frac{(t_0^44at_0^36a^2t_0^2+4a^3t_0+a^4)(t_0^4+4at_0^36a^2t_0^24a^3t_0+a^4)}{\pi ^4(t_0^2+a^2)^8}.$$
(52)
This function has a form similar to that illustrated in Fig. 3, except that it has three maxima and two minima. It is somewhat difficult to graph because the relative sizes of the extrema decrease very rapidly.
In the limit that $`r=0`$, we may write the four-dimensional correlation function as
$$C(t,t^{})=\frac{3}{2\pi ^4(tt^{})^8}=\frac{1}{6720\pi ^4}_t^4_t^{}^4\mathrm{ln}[(tt^{})^2/\mathrm{}^2].$$
(53)
We can sample the energy density with arbitrary displaced sampling functions and write
$`S(t_0)S(0)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}g(tt_0)g(t^{})C(t,t^{})`$ (54)
$`=`$ $`{\displaystyle \frac{1}{6720\pi ^4}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}[_t^4g(tt_0)][_t^{}^4g(t^{})]\mathrm{ln}[\mathrm{\Delta }t^2/\mathrm{}^2].`$ (55)
It should be noted that here we did not use the form of the energy density correlation function, Eq. (35), which follows from the covariant form. Instead, we let $`r0`$, and then expressed the result in terms of time derivatives of a logarithmic function. A more rigorous approach would be to average Eq. (35) over both space and time, and then let the widths of the spatial sampling functions go to zero. However, this is difficult to do explicitly with general sampling functions. The equivalence of the two approaches needs to be studied more carefully.
Let us next consider a compactly supported sampling function given by
$$g(t)=\frac{630}{a^9}(ta/2)^4(t+a/2)^4,$$
(56)
for $`|t|a/2`$, and $`g(t)=0`$ for $`|t|a/2`$. Note that $`g(t)`$ and its first three derivatives vanish at $`t=\pm \frac{1}{2}a`$, so all surface terms vanish when we integrated by parts in Eq. (55) to obtain the second form for $`S(t_0)S(0)`$. We may again define $`K(t_0)`$ by Eq. (26) and evaluate it numerically. The result is plotted in Fig. 5.
As in two dimensions, there are regions of correlation and of anticorrelation as $`t_0`$ increases. However, the behavior in four dimensions is more complicated, with three maxima and two minima. This appears to be due to the greater number of derivatives of the sampling function in Eq. (55), as compared to Eq. (25).
In Appendix C, we show that in two and four dimensions
$$_0^{\mathrm{}}K(t_0)𝑑t_0=0,$$
(57)
and that
$$S^2(0)>0,$$
(58)
for a general $`g(t)`$. From Eq. (52), we can also explicitly verify that $`_0^{\mathrm{}}K(t_0)𝑑t_0=0`$ for the Lorentzian sampling function.
## IV Summary
In this paper, we have presented covariant expressions for the Minkowski vacuum stress tensor correlation function in two dimensions, Eq. (2), and in four dimensions, Eq. (34). These expressions are of the form of a sum of terms, each of which is a derivative of a scalar function with logarithmic singularities in the coincidence limit. These expressions allow one to write spacetime averages of the correlation function as finite integrals. We explicitly evaluated such averages of the energy density in two dimensions using Lorentzian sampling functions in both space and time. The resulting expression, Eq. (17), is symmetric in the spatial and temporal sampling widths, and is finite as either width goes to zero with the other width fixed at a nonzero value.
We next studied the correlations of the sampled 2D energy density along a worldline using displaced sampling functions. This reveals the correlation and anticorrelation of measurements of the energy density in overlapping intervals. The result is illustrated in Fig. 4 for a compactly supported sampling function. When the intervals nearly overlap, the two measurements are positively correlated, as expected. When the overlap has decreased somewhat, the two measurements become anticorrelated. This can be interpreted as telling us that if we find energy density of one sign on the first measurement, we should find the opposite sign on the next measurement. Finally, as the intervals become disjoint, the measurements are again positively correlated. Furthermore, we show that for an arbitrary sampling function, the net area under the correlation graph, e.g., the one depicted in Fig. 4, is equal to zero. It is hoped that further investigation will elucidate this interesting behavior.
The analogous calculation in four dimensions yields similar results. However, in this case there are two regions of anticorrelation and three of positive correlation. The fluctuations in the averaged energy density remain finite in the limit that the spatial width vanishes, but not in the limit that the temporal width goes to zero. Thus in four dimensions, the averaged energy density correlation function requires averaging in time to be finite.
There is a vaguely analogous result concerning quantum inequalities on the averaged expectation value of the stress tensor in an arbitrary state. There are finite lower bounds on the expectation value of the energy density averaged on a worldline in both 2D and 4D, and on the spatial average in 2D. However, the spatial average in 4D has no lower bound FHR . The search for a deeper link between quantum inequalities and the vacuum stress tensor correlation function is a topic for future research.
Another question which needs to be explored further is that of the physical effects of the passive metric fluctuations driven by vacuum stress tensor fluctuations. One approach is that adopted in Ref. BF where the Raychaudhuri equation was used as a Langevin equation to study the luminosity fluctuations and angular blurring of a distant source produced by passive metric fluctuations. The case of the Minkowski vacuum was briefly discussed in Ref. BF , where it was found that the natural quantum uncertainty in the test particles used to probe the fluctuating geometry tends to hide the effects of the metric fluctuations. However, this does not necessarily mean that these fluctuations are in principle unobservable. This is another question for further study.
###### Acknowledgements.
We would like to thank Albert Roura for valuable comments and Alan Guth for providing his unpublished notes. This work was supported in part by the National Science Foundation under Grants PHY-0139969 and PHY-0244898.
## Appendix A
In this appendix, we give the derivation of the stress tensor correlation function in two-dimensional spacetime in the Minkowski vacuum state. We first start with the form of the stress tensor for a massless, minimally coupled scalar field:
$$T_{\mu \nu }=\varphi _{,\mu }\varphi _{,\nu }\frac{1}{2}g_{\mu \nu }\varphi ^{,\rho }\varphi _{,\rho }.$$
(59)
From this expression we find the correlation function
$`C_{\mu \nu \alpha \beta }(x,x^{})`$ $`=`$ $`:T_{\mu \nu }(x)::T_{\alpha \beta }(x^{}):=:_\mu \varphi _\nu \varphi ::_\alpha ^{}\varphi _\beta ^{}\varphi :`$ (60)
$``$ $`{\displaystyle \frac{1}{2}}g_{\mu \nu }:^\rho \varphi _\rho \varphi ::_\alpha ^{}\varphi _\beta ^{}\varphi :{\displaystyle \frac{1}{2}}g_{\alpha \beta }:_\mu \varphi _\nu \varphi ::^\sigma ^{}\varphi _\sigma ^{}\varphi :`$
$``$ $`{\displaystyle \frac{1}{4}}g_{\mu \nu }g_{\alpha \beta }:^\rho \varphi _\rho \varphi ::^\sigma ^{}\varphi _\sigma ^{}\varphi :.`$
Here unprimed indices refer to the point $`x`$ and primed indices to $`x^{}`$. Next we use the identity
$$:\varphi _1\varphi _2::\varphi _3\varphi _4:=\varphi _1\varphi _3\varphi _2\varphi _4+\varphi _1\varphi _4\varphi _2\varphi _3,$$
(61)
where the $`\varphi _i`$ are quantum fields or derivatives of quantum fields. From this identity, we can show that
$$:_\mu \varphi _\nu \varphi ::_\alpha ^{}\varphi _\beta ^{}\varphi :=(_\mu _\alpha ^{}D)(_\nu _\beta ^{}D)+(_\mu _\beta ^{}D)(_\nu _\alpha ^{}D),$$
(62)
where
$$D=D(x,x^{})=\varphi (x)\varphi (x^{})$$
(63)
is the two-point function. We can express the correlation function in terms of derivatives of D as
$`C_{\mu \nu \alpha \beta }(x,x^{})`$ $`=`$ $`(_\mu _\alpha ^{}D)(_\nu _\beta ^{}D)+(_\mu _\beta ^{}D)(_\nu _\alpha ^{}D)`$ (64)
$``$ $`g_{\mu \nu }(^\rho _\alpha ^{}D)(_\rho _\beta ^{}D)g_{\alpha \beta }(_\mu ^\sigma ^{}D)(_\nu _\sigma ^{}D)`$
$``$ $`{\displaystyle \frac{1}{2}}g_{\mu \nu }g_{\alpha \beta }(^\rho ^\sigma ^{}D)(_\rho _\sigma ^{}D).`$
An equivalent expression for the case of a massive, nonminimal scalar field has been given by Martin and Verdaguer. (See Eq. 3.42 in Ref. MV99 .) The analogous expression for the electromagnetic field is given in Ref. FW03 .
Up to this point, our treatment applies to spacetimes of any dimensionality. Now we specialize to two-dimensional Minkowski spacetime. There is an infrared divergence in the two-point function for a massless scalar field in the Minkowski vacuum state in two dimensions. Thus, the field must either have a nonzero mass, or else the only physically allowed states are ones which break Lorentz invariance FV86 . Fortunately, the details of either approach have no effect on our results. If we let the scalar field have a small mass $`m`$, then the two-point function is given by
$$D=\frac{1}{4\pi }\mathrm{ln}(cm^2\mathrm{\Delta }x^2)$$
(65)
in the limit that $`m^2\mathrm{\Delta }x^2`$. Here $`c`$ is a dimensionless constant and $`\mathrm{\Delta }x^2=(x^\mu x_{}^{}{}_{}{}^{\mu })(x_\mu x_{}^{}{}_{\mu }{}^{})`$. Because the stress tensor correlation function depends only upon derivatives of $`D`$, it is independent of $`c`$ and $`m`$. The second derivative of $`D`$ is
$$_\mu _\alpha ^{}D=\frac{2\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\alpha g_{\mu \alpha }\mathrm{\Delta }x^2}{2\pi (\mathrm{\Delta }x^2)^2},$$
(66)
where $`\mathrm{\Delta }x_\mu =x_\mu x_\mu ^{}`$. We can now combine Eqs. (64) and (66) to obtain an explicit expression for the stress tensor correlation function in two dimensions:
$`C_{\mu \nu \alpha \beta }(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}[{\displaystyle \frac{8}{(\mathrm{\Delta }x^2)^4}}\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta `$ (67)
$``$ $`{\displaystyle \frac{2}{(\mathrm{\Delta }x^2)^3}}(g_{\mu \alpha }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\beta +g_{\mu \beta }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha +g_{\nu \alpha }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\beta +g_{\nu \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\alpha )`$
$`+`$ $`{\displaystyle \frac{1}{(\mathrm{\Delta }x^2)^2}}(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha }g_{\mu \nu }g_{\alpha \beta })]`$
We next wish to express $`C_{\mu \nu \alpha \beta }(x,x^{})`$ as a sum of derivatives of scalar functions. Lorentz symmetry suggests that these be functions of $`\mathrm{\Delta }x^2`$. Let $`f=f(\mathrm{\Delta }x^2)`$. Then the derivatives of $`f`$ are
$$_\mu f=2\mathrm{\Delta }x_\mu f^{},$$
(68)
$$_\mu _\nu f=2g_{\mu \nu }f^{}+4\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu f^{\prime \prime },$$
(69)
$$_\mu _\nu _\alpha f=4(g_{\mu \nu }\mathrm{\Delta }x_\alpha +g_{\mu \alpha }\mathrm{\Delta }x_\nu +g_{\nu \alpha }\mathrm{\Delta }x_\mu )f^{\prime \prime }+8\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha f^{\prime \prime \prime },$$
(70)
and
$`_\mu _\nu _\alpha _\beta f`$ $`=`$ $`4(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha }+g_{\mu \nu }g_{\alpha \beta })f^{\prime \prime }+8(g_{\mu \alpha }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\beta +g_{\mu \beta }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha `$ (71)
$`+`$ $`g_{\nu \alpha }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\beta +g_{\nu \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\alpha +g_{\mu \nu }\mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta +g_{\alpha \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu )f^{\prime \prime \prime }`$
$`+`$ $`16\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta f^{\prime \prime \prime \prime }.`$
Here primes denote derivatives of $`f`$ with respect to its argument. We will also need some expressions involving the wave operator:
$$\mathrm{}f=_\mu ^\mu f=2nf^{}+4\mathrm{\Delta }x^2f^{\prime \prime },$$
(72)
$$\mathrm{}\mathrm{}f=4n(n+2)f^{\prime \prime }+16(n+2)\mathrm{\Delta }x^2f^{\prime \prime \prime }+16(\mathrm{\Delta }x^2)^2f^{\prime \prime \prime \prime },$$
(73)
and
$`_\mu _\nu \mathrm{}f`$ $`=`$ $`4(n+2)g_{\mu \nu }f^{\prime \prime }+8[g_{\mu \nu }\mathrm{\Delta }x^2+(n+4)\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu ]f^{\prime \prime \prime }`$ (74)
$`+`$ $`16\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu \mathrm{\Delta }x^2f^{\prime \prime \prime },`$
where $`n`$ is the dimension of the spacetime.
Our goal is to express $`C_{\mu \nu \alpha \beta }(x,x^{})`$ as a sum of derivatives acting on one or more choices of $`f`$. Because $`_\mu f_1=_\mu ^{}f_1`$, we can write our results in several equivalent forms, but here and in Appendix B we will use derivatives with unprimed indices. If $`f`$ is dimensionless, then in two dimensions we will need four derivatives in each term in order that $`C_{\mu \nu \alpha \beta }(x,x^{})`$ has dimensions of $`\mathrm{length}^4`$. There are five fourth-rank tensors that we can form which have the correct dimensions and symmetry properties:
$$_\mu _\nu _\alpha _\beta f,$$
(75)
$$(g_{\alpha \beta }_\mu _\nu +g_{\mu \nu }_\alpha _\beta )\mathrm{}f,$$
(76)
$$(g_{\alpha \nu }_\mu _\beta +g_{\alpha \mu }_\nu _\beta +g_{\beta \nu }_\mu _\alpha +g_{\beta \mu }_\nu _\alpha )\mathrm{}f,$$
(77)
$$g_{\mu \nu }g_{\alpha \beta }\mathrm{}\mathrm{}f,$$
(78)
and
$$(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })\mathrm{}\mathrm{}f.$$
(79)
We would like $`f`$ to have an integrable singularity at $`\mathrm{\Delta }x^2=0`$, so a natural choice is a power of a logarithmic function. First consider
$$f_1=\mathrm{ln}(\mathrm{\Delta }x^2/\mathrm{}^2),$$
(80)
where $`\mathrm{}`$ is an arbitrary constant with dimensions of length. However, $`\mathrm{}f_1=0`$ in two dimensions, so the only nonzero tensor from the above list which can be formed from $`f_1`$ is
$`_\mu _\nu _\alpha _\beta f_1`$ $`=`$ $`{\displaystyle \frac{96}{(\mathrm{\Delta }x^2)^4}}\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta `$ (81)
$`+`$ $`{\displaystyle \frac{16}{(\mathrm{\Delta }x^2)^3}}(g_{\mu \alpha }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\beta +g_{\mu \beta }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha +g_{\nu \alpha }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\beta `$
$`+`$ $`g_{\nu \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\alpha +g_{\mu \nu }\mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta +g_{\alpha \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu )`$
$``$ $`{\displaystyle \frac{4}{(\mathrm{\Delta }x^2)^2}}(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha }+g_{\mu \nu }g_{\alpha \beta }).`$
This is not sufficient to form $`C_{\mu \nu \alpha \beta }(x,x^{})`$, so we need another choice of $`f`$, which we take to be
$$f_2=\mathrm{ln}^2(\mathrm{\Delta }x^2/\mathrm{}^2).$$
(82)
From Eqs. (72) and (73) with $`n=2`$, we find
$$\mathrm{}f_2=\frac{8}{\mathrm{\Delta }x^2}$$
(83)
and
$$\mathrm{}\mathrm{}f_2=\frac{32}{(\mathrm{\Delta }x^2)^2}.$$
(84)
This allows us to form four tensors from $`f_2`$ with the correct symmetry properties and dimension:
$$g_{\mu \nu }g_{\alpha \beta }\mathrm{}\mathrm{}f_2=32g_{\mu \nu }g_{\alpha \beta }\frac{1}{(\mathrm{\Delta }x^2)^2},$$
(85)
$$(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })\mathrm{}\mathrm{}f_2=32(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })\frac{1}{(\mathrm{\Delta }x^2)^2},$$
(86)
$$(g_{\alpha \beta }_\mu _\nu +g_{\mu \nu }_\alpha _\beta )\mathrm{}f_2=32\frac{g_{\mu \nu }g_{\alpha \beta }}{(\mathrm{\Delta }x^2)^2}+64\frac{g_{\mu \nu }\mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta +g_{\alpha \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu }{(\mathrm{\Delta }x^2)^3},$$
(87)
and
$`(g_{\alpha \nu }_\mu _\beta +g_{\alpha \mu }_\nu _\beta +g_{\beta \nu }_\mu _\alpha +g_{\beta \mu }_\nu _\alpha )\mathrm{}f_2=32{\displaystyle \frac{g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha }}{(\mathrm{\Delta }x^2)^2}}`$ (88)
$`+`$ $`{\displaystyle \frac{64}{(\mathrm{\Delta }x^2)^3}}(g_{\mu \alpha }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\beta +g_{\mu \beta }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha +g_{\nu \alpha }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\beta +g_{\nu \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\alpha ).`$
Note that $`_\mu _\nu _\alpha _\beta f_2`$ is not a suitable term because it contains logarithmic pieces that do not appear in $`C_{\mu \nu \alpha \beta }(x,x^{})`$ and which cannot be cancelled by any other terms. This leaves us with five tensors from which to form the stress tensor correlation function.
Let
$`C_{\mu \nu \alpha \beta }(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{384\pi ^2}}[c_1_\mu _\nu _\alpha _\beta f_1+c_2g_{\mu \nu }g_{\alpha \beta }\mathrm{}\mathrm{}f_2`$ (89)
$`+`$ $`c_3(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })\mathrm{}\mathrm{}f_2+c_4(g_{\alpha \beta }_\mu _\nu +g_{\mu \nu }_\alpha _\beta )\mathrm{}f_2`$
$`+`$ $`c_5(g_{\alpha \nu }_\mu _\beta +g_{\alpha \mu }_\nu _\beta +g_{\beta \nu }_\mu _\alpha +g_{\beta \mu }_\nu _\alpha )\mathrm{}f_2].`$
If we insert Eq. (81) and Eqs. (85)-(88) into this expression and compare with Eq. (67), we find five conditions on the five coefficients. The unique solution of these conditions gives
$$c_1=8,c_2=c_4=2,\mathrm{and}c_3=c_5=1.$$
(90)
As a check, the correlation function may be shown explicitly to satisfy the conservation law
$$^\mu C_{\mu \nu \alpha \beta }(x,x^{})=^\alpha ^{}C_{\mu \nu \alpha \beta }(x,x^{})=0.$$
(91)
## Appendix B
Here we repeat the derivation in the previous appendix for the case of four-dimensional Minkowski spacetime. The general form, Eq. (64), for the correlation function still holds, but the two-point function for a massless scalar field is now
$$D=\frac{1}{4\pi ^2\mathrm{\Delta }x^2}.$$
(92)
If we insert this form into Eq. (64), we find the four-dimensional analog of Eq. (67):
$`C_{\mu \nu \alpha \beta }(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^4}}[{\displaystyle \frac{32}{(\mathrm{\Delta }x^2)^6}}\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta `$ (93)
$``$ $`{\displaystyle \frac{4}{(\mathrm{\Delta }x^2)^5}}(g_{\mu \alpha }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\beta +g_{\mu \beta }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha +g_{\nu \alpha }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\beta +g_{\nu \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\alpha )`$
$``$ $`{\displaystyle \frac{8}{(\mathrm{\Delta }x^2)^5}}(g_{\mu \nu }\mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta +g_{\alpha \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu )`$
$`+`$ $`{\displaystyle \frac{1}{(\mathrm{\Delta }x^2)^4}}(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha }+4g_{\mu \nu }g_{\alpha \beta })]`$
In four dimensions, the correlation function has dimensions of $`1/\mathrm{length}^8`$. Thus any expression involving derivatives on a dimensionless function will require eight derivatives. Because there are only four free indices, there will have to be at least two wave operators. This eliminates the logarithm function $`f_1`$, Eq. (80), because in four dimensions
$$\mathrm{}\mathrm{}f_1=0.$$
(94)
However, the squared logarithm function $`f_2`$ may be used to form the following five tensors with the correct dimensions and symmetry:
$$_\mu _\nu _\alpha _\beta \mathrm{}\mathrm{}f_2,$$
(95)
$$(g_{\alpha \beta }_\mu _\nu +g_{\mu \nu }_\alpha _\beta )\mathrm{}^3f_2,$$
(96)
$$(g_{\alpha \nu }_\mu _\beta +g_{\alpha \mu }_\nu _\beta +g_{\beta \nu }_\mu _\alpha +g_{\beta \mu }_\nu _\alpha )\mathrm{}^3f_2,$$
(97)
$$g_{\mu \nu }g_{\alpha \beta }\mathrm{}^4f_2,$$
(98)
and
$$(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })\mathrm{}^4f_2.$$
(99)
We may repeatedly use Eqs. (72) and (73) with $`n=4`$ to demonstrate that, in four-dimensions,
$$\mathrm{}\mathrm{}f_2=\frac{32}{(\mathrm{\Delta }x^2)^2},$$
(100)
$$\mathrm{}^3f_2=\frac{256}{(\mathrm{\Delta }x^2)^3},$$
(101)
and
$$\mathrm{}^4f_2=\frac{6144}{(\mathrm{\Delta }x^2)^4}.$$
(102)
From these expressions, we may show that
$`_\mu _\nu _\alpha _\beta \mathrm{}\mathrm{}f_2`$ $`=`$ $`{\displaystyle \frac{61440}{(\mathrm{\Delta }x^2)^6}}\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta `$ (103)
$`+`$ $`{\displaystyle \frac{6144}{(\mathrm{\Delta }x^2)^5}}(g_{\mu \alpha }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\beta +g_{\mu \beta }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha +g_{\nu \alpha }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\beta `$
$`+`$ $`g_{\nu \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\alpha +g_{\mu \nu }\mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta +g_{\alpha \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu )`$
$``$ $`{\displaystyle \frac{768}{(\mathrm{\Delta }x^2)^4}}(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha }+g_{\mu \nu }g_{\alpha \beta }),`$
$`(g_{\alpha \nu }_\mu _\beta +g_{\alpha \mu }_\nu _\beta +g_{\beta \nu }_\mu _\alpha +g_{\beta \mu }_\nu _\alpha )\mathrm{}^3f_2={\displaystyle \frac{3072}{(\mathrm{\Delta }x^2)^4}}(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })`$ (104)
$``$ $`{\displaystyle \frac{12288}{(\mathrm{\Delta }x^2)^5}}(g_{\mu \alpha }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\beta +g_{\mu \beta }\mathrm{\Delta }x_\nu \mathrm{\Delta }x_\alpha +g_{\nu \alpha }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\beta +g_{\nu \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\alpha ),`$
and
$$(g_{\alpha \beta }_\mu _\nu +g_{\mu \nu }_\alpha _\beta )\mathrm{}^3f_2=3072\left[\frac{g_{\mu \nu }g_{\alpha \beta }}{(\mathrm{\Delta }x^2)^4}4\frac{g_{\mu \nu }\mathrm{\Delta }x_\alpha \mathrm{\Delta }x_\beta +g_{\alpha \beta }\mathrm{\Delta }x_\mu \mathrm{\Delta }x_\nu }{(\mathrm{\Delta }x^2)^5}\right].$$
(105)
We now express the correlation function as a sum of the tensors formed from $`f_2`$ as
$`C_{\mu \nu \alpha \beta }(x,x^{})`$ $`=`$ $`{\displaystyle \frac{1}{61440\pi ^4}}[c_1_\mu _\nu _\alpha _\beta \mathrm{}\mathrm{}f_2+c_2g_{\mu \nu }g_{\alpha \beta }\mathrm{}^4f_2`$ (106)
$`+`$ $`c_3(g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })\mathrm{}^4f_2+c_4(g_{\alpha \beta }_\mu _\nu +g_{\mu \nu }_\alpha _\beta )\mathrm{}^3f_2`$
$`+`$ $`c_5(g_{\alpha \nu }_\mu _\beta +g_{\alpha \mu }_\nu _\beta +g_{\beta \nu }_\mu _\alpha +g_{\beta \mu }_\nu _\alpha )\mathrm{}^3f_2].`$
If we insert the explicit forms for these tensors and compare with Eq. (93), we again find five conditions on the five coefficients, leading to the solution
$$c_1=8,c_2=c_4=6,\mathrm{and}c_3=c_5=1.$$
(107)
As required, the correlation function has a vanishing divergence on any index.
## Appendix C
In this appendix, we will prove Eqs. (32) and (33) in both two and four dimensions. We will proceed by first showing that
$$_{\mathrm{}}^{\mathrm{}}K(t_0)𝑑t_0=0.$$
(108)
Then we will prove that $`K(t_0)`$ is a symmetric function, and hence show that
$$_0^{\mathrm{}}K(t_0)𝑑t_0=0,$$
(109)
as well.
Let $`g(t)`$ be an arbitrary smooth sampling function. From Eq. (24) or (54) and Eq. (26), we have
$$K(t_0)=\frac{1}{S^2(0)}_{\mathrm{}}^{\mathrm{}}𝑑tg(tt_0)_{\mathrm{}}^{\mathrm{}}𝑑t^{}g(t^{})C(tt^{}).$$
(110)
Then
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}K(t_0)𝑑t_0`$ $`=`$ $`{\displaystyle \frac{1}{S^2(0)}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_0g(tt_0){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}g(t^{})C(tt^{})`$ (111)
$`=`$ $`{\displaystyle \frac{1}{S^2(0)}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}g(t^{}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑tC(tt^{}),`$
where we have interchanged the order of integrations, and used the fact that for $`y=tt_0`$,
$$_{\mathrm{}}^{\mathrm{}}𝑑t_0g(tt_0)=_{\mathrm{}}^{\mathrm{}}𝑑yg(y)=_{\mathrm{}}^{\mathrm{}}𝑑yg(y)=1.$$
(112)
However, if we can write $`C(tt^{})=F(tt^{})/t`$, where $`F(tt^{})0`$ as $`t\pm \mathrm{}`$, then
$$_{\mathrm{}}^{\mathrm{}}𝑑tC(tt^{})=[F(tt^{})]_{t=\mathrm{}}^{t=+\mathrm{}}=0,$$
(113)
which in turn implies that
$$_{\mathrm{}}^{\mathrm{}}K(t_0)𝑑t_0=0.$$
(114)
Recall that in two dimensions, the worldline vacuum correlation function is $`C(tt^{})=1/[4\pi ^2(tt^{})^4]`$, and in four dimensions it is $`C(tt^{})=3/[2\pi ^4(tt^{})^8]`$, so in both cases the condition Eq. (113) is satisfied. Note that in four dimensions, it is necessary to assume that we set the spatial separation $`r`$ in Eq. (39) to zero and then average over time, as discussed in Sect. III.3.
We now show that $`K(t_0)=K(t_0)`$. Let us write
$`S^2(0)K(t_0)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑tg(tt_0){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}g(t^{})C(tt^{})`$ (115)
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\overline{t}g(\overline{t}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}g(t^{})C(\overline{t}+t_0t^{})`$
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\overline{t}g(\overline{t}){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\overline{t}^{}g(\overline{t}^{}+t_0)C(\overline{t}\overline{t}^{}),`$
where we have let $`\overline{t}=tt_0`$, so $`t=\overline{t}+t_0`$, and $`\overline{t}^{}=t^{}t_0`$. If we now let $`\overline{t}^{}t,\overline{t}t^{}`$, we have
$`S^2(0)K(t_0)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑tg(t+t_0){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}g(t^{})C(t^{}t)`$ (116)
$`=`$ $`S^2(0)K(t_0),`$
where we have used the fact $`C(t^{}t)=C(tt^{})`$. Note that the symmetry of $`K(t_0)`$ depends only on that of $`C`$ and does not assume that the sampling function $`g(t)`$ is symmetric. Thus since
$$_{\mathrm{}}^{\mathrm{}}K(t_0)𝑑t_0=0,$$
(117)
and $`K(t_0)`$ is symmetric, it also follows that
$$_0^{\mathrm{}}K(t_0)𝑑t_0=0.$$
(118)
In order to determine whether a fluctuation is correlated or anti-correlated with itself, we must determine the sign of $`S^2(0)`$ in the general case. We would expect that a fluctuation should be correlated with itself, and thus that $`S^2(0)>0`$. This can be proven from the fact that $`S(0)`$, as defined by Eq. (22), is a self-adjoint operator Comment . Let $`|\psi `$ be the state under consideration, which in our case is the Minkowski vacuum. Then
$$|\mathrm{\Psi }=S(0)|\psi $$
(119)
is a well defined state vector with positive norm. Thus we have
$$\mathrm{\Psi }^2=\psi |S^{}(0)S(0)|\psi =S^2(0)>0.$$
(120)
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# Quantitative Models and Implicit Complexity
## 1 Introduction
In recent years, a large number of characterizations of complexity classes based on logics and lambda calculi have appeared. At least three different principles have been exploited, namely linear types , restricted modalities in the context of linear logic and non-size-increasing computation . Although related one to the other, these systems have been studied with different, often unrelated methodologies and few results are known about relative intentional expressive power. We believe that this area of implicit computational complexity needs unifying frameworks for the analysis of quantitative properties of computation. This would help to improve the understanding on existing systems. More importantly, unifying frameworks can be used *themselves* as a foundation for controlling the use of resources inside programming languages.
In this paper, we introduce a new semantical framework which consists of an innovative modification of realizability. The main idea underlying our proposal lies in considering bounded-time algorithms as realizers instead of taking plain Turing Machines as is usually the case in realizability constructions. Bounds are expressed abstractly as elements of a monoid. We can define a model for a given (logical or type) system by choosing a monoid flexible enough to justify all the constructs in the system. The model can then be used to study the class of representable functions.
This allows us to give new proofs of soundness (all representable functions on base types lies in certain complexity classes) for Light Affine Logic (LAL, ), Elementary Affine Logic (EAL, ), LFPL and Soft Affine Logic (SAL, ). While being the first entirely semantical proof of polytime soundness for light logics, our proof also provides a notable simplification of the original already semantical proof of polytime soundness for LFPL . A new result made possible by the semantic framework is the addition of polymorphism and a modality to LFPL.
The rest of the paper is organized as follows. In Section 2 we describe an abstract computational model that will be used in the rest of the paper. In Section 3 we introduce length spaces and show they can be used to interpret multiplicative linear logic with free weakening. Sections 4, 5 and 6 are devoted to present instances of the framework together with soundness results for elementary, soft and light affine logics. Section 7 presents a further specialization of length spaces and a new soundness theorem for LFPL based on it.
#### Related-Work
Realizability has been used in connection with resource-bounded computation in several places. The most prominent is Cook and Urquhart work , where terms of a language called $`\text{PV}^\omega `$ are used to realize formulas of bounded arithmetic. The contribution of that paper is related to ours in that realizability is used to show “polytime soundness” of a logic. There are important differences though. First, realizers in Cook and Urquhart are typed and very closely related to the logic that is being realized. Second, the language of realizers $`\text{PV}^\omega `$ only contains first order recursion and is therefore useless for systems like LFPL or LAL. In contrast, we use untyped realizers and interpret types as certain partial equivalence relations on those. This links our work to the untyped realizability model HEO (due to Kreisel ). This, in turn, has also been done by Crossley et al. . There, however, one proves externally that untyped realizers (in this case of bounded arithmetic formulas) are polytime. In our work, and this happens for the first time, the untyped realizers are used to give meaning to the logic and obtain polytime soundness as a corollary. Thus, certain resource bounds are built into the untyped realizers by their very construction. Such a thing is not at all obvious, because untyped universes of realizers tend to be Turing complete from the beginning to due definability of fixed-point combinators. We get around this problem through our notion of a resource monoid and addition of a certain time bound to Kleene applications of realizers. Indeed, we consider this as the main innovation of our paper and hope it to be useful elsewhere.
## 2 A Computational Model
In this paper, we rely on an abstract computational framework rather than a concrete one like Turing Machines. This, in particular, will simplify proofs.
Let $`L\mathrm{\Sigma }^{}`$ be the set of finite sequences over the alphabet $`\mathrm{\Sigma }`$. We assume a pairing function $`,:L\times LL`$ and a length function $`||:L`$ such that $`|x,y|=|x|+|y|+\mathrm{𝑐𝑝}`$ and $`|x|\text{length}(x)`$, where $`\text{length}(x)`$ is the number of symbols in $`x`$ and $`\mathrm{𝑐𝑝}`$ is a fixed constant. We assume a reasonable encoding of algorithms as elements of $`L`$. We write $`\{e\}(x)`$ for the (possibly undefined) application of algorithm $`eL`$ to input $`xL`$. We furthermore assume an abstract time measure $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(x))`$ such that $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(x))`$ is defined whenever $`\{e\}(x)`$ is and, moreover
* $`\{e\}(x)`$ can be evaluated on a Turing machine in time bounded by $`p(\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(x))+|e|+|x|)`$, where $`p:`$ is a fixed polynomial.
* For each Turing machine $`M`$ running in time $`f:`$, there is $`eL`$ so that $`\{e\}(\mathrm{\Phi }(x))=\mathrm{\Phi }(y)`$, (where $`y`$ is the result of running $`M`$ on input $`x`$). Furthermore, $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(\mathrm{\Phi }(x)))=O(f(|x|))`$.
* $`B=\{0,1\}^{}`$ can be embedded into $`L`$ by a map $`\mathrm{\Phi }:BL`$ such that both $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^1`$ can be computed in polynomial time.
* There are $`e_0,e_1L`$ such that for every $`xB`$, $`\{e_0\}(\mathrm{\Phi }(x))=\mathrm{\Phi }(0x)`$, $`\{e_1\}(\mathrm{\Phi }(x))=\mathrm{\Phi }(1x)`$. Moreover, $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_0\}(x))=\mathrm{𝑇𝑖𝑚𝑒}(\{e_1\}(x))=O(1)`$.
* There is $`e_{\mathrm{𝑐𝑜𝑚𝑝}}`$ (composition) such that for every $`x,y`$ it holds that $`\{e_{\mathrm{𝑐𝑜𝑚𝑝}}\}(x,y)=z`$ where $`|z|=|x|+|y|+O(1)`$ and $`\{z\}(w)=\{y\}(\{x\}(w))`$; moreover, $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑚𝑝}}\}(x,y))=O(1)`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑚𝑝}}\}(w))=\mathrm{𝑇𝑖𝑚𝑒}(\{x\}(w))+\mathrm{𝑇𝑖𝑚𝑒}(\{y\}(\{x\}(w)))+O(1)`$.
* There is $`e_{\mathrm{𝑖𝑑}}`$ (identity) such that $`\{e_{\mathrm{𝑖𝑑}}\}(x)=x`$ for every $`x`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑖𝑑}}\}(x))=O(1)`$.
* For every $`xL`$ there is $`e_{\mathrm{𝑐𝑜𝑛𝑠𝑡}}^x`$ such that $`\{e_{\mathrm{𝑐𝑜𝑛𝑠𝑡}}^x\}(y)=x`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑛𝑠𝑡}}^x\}(y))=O(1)`$.
* For every $`xL`$ there is $`e_{\mathrm{𝑡𝑒𝑛𝑠𝑐𝑜𝑛𝑠𝑡}}^x`$ such that $`\{e_{\mathrm{𝑡𝑒𝑛𝑠𝑐𝑜𝑛𝑠𝑡}}^x\}(y)=y,x`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑡𝑒𝑛𝑠𝑐𝑜𝑛𝑠𝑡}}^x\}(y))=O(1)`$.
* There is $`e_{\mathrm{𝑡ℎ𝑟𝑜𝑤𝑓𝑖𝑟𝑠𝑡}}`$ such that for every $`xL`$ $`\{e_{\mathrm{𝑡ℎ𝑟𝑜𝑤𝑓𝑖𝑟𝑠𝑡}}\}(x,y)=y`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑡ℎ𝑟𝑜𝑤𝑓𝑖𝑟𝑠𝑡}}\}(x,y))=O(1)`$.
* There is $`e_{\mathrm{𝑠𝑤𝑎𝑝}}`$ (swapping) such that $`\{e_{\mathrm{𝑠𝑤𝑎𝑝}}\}(x,y)=y,x`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑠𝑤𝑎𝑝}}\}(z))O(1)`$.
* There is $`e_{\mathrm{𝑡𝑒𝑛𝑠}}`$ (tensor) such that for every $`x`$ $`\{e_{\mathrm{𝑡𝑒𝑛𝑠}}\}(x)=y`$ where $`|y|=|x|+O(1)`$ and $`\{y\}(z,w)=\{x\}(z),w`$; moroever, $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑡𝑒𝑛𝑠}}\}(x))=O(1)`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{y\}(z,w))=\mathrm{𝑇𝑖𝑚𝑒}(\{x\}(z))+O(1)`$.
* There is $`e_{\mathrm{𝑎𝑠𝑠𝑙}}`$ (rebracketing) such that $`\{e_{\mathrm{𝑎𝑠𝑠𝑙}}\}(x,y,z)=x,y,z`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑎𝑠𝑠𝑙}}\}(x))=O(1)`$.
* There is $`e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}`$ (duplication, copying) such that $`\{e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}\}(x)=x,x`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}\}(x))=O(|x|)`$.
* There is $`e_{\mathrm{𝑒𝑣𝑎𝑙}}`$ (application) such that $`\{e_{\mathrm{𝑒𝑣𝑎𝑙}}\}(x,y)=\{x\}(y)`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑒𝑣𝑎𝑙}}\}(x,y))=\mathrm{𝑇𝑖𝑚𝑒}(\{x\}(y))+O(1)`$.
* There is $`e_{\mathrm{𝑐𝑢𝑟𝑟𝑦}}`$ (currying, “smn-theorem”) such that, for each $`x`$, $`y=\{e_{\mathrm{𝑐𝑢𝑟𝑟𝑦}}\}(x)`$ exists and satisfies $`|y|=|x|+O(1)`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑢𝑟𝑟𝑦}}\}(x))=O(1)`$; moreover, for every $`z`$, $`c_z=\{y\}(z)`$ exists and satisfies $`|c_z|=|y|+|z|+O(1)`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{y\}(z))=O(1)`$; finally, for every $`w`$, $`\{c_z\}(w)=\{x\}(z,w)`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{c_z\}(w))=\mathrm{𝑇𝑖𝑚𝑒}(\{x\}(z,w))+O(1)`$.
There are a number of ways to instantiate this framework. One noticeable and simple way consists in using call-by-value lambda calculus and is described in the following. $`\mathrm{\Sigma }`$ will be $`\{\lambda ,\mathrm{@},0,1,\}`$. To any lambda term $`M\mathrm{\Lambda }`$, we can associate a string $`M^\mathrm{\#}\mathrm{\Sigma }^{}`$ in the obvious way. For example, if $`M(\lambda x.xy)(\lambda x.\lambda y.\lambda z.x)`$, then $`M^\mathrm{\#}`$ is
$$\mathrm{@}\lambda \mathrm{@}0\lambda \lambda \lambda 10$$
In other words, free occurrences of variables are translated into $``$, while bounded occurrences of variables are translated into $`s`$, where $`s`$ is the binary representation of the deBruijn index for the occurrence. $`L`$ will just be the set of strings in $`\mathrm{\Sigma }^{}`$ corresponding to lambda terms via the mapping we just described. In the following, we will often write a lambda-term in the usual notation, but this is just syntactic sugar for the corresponding element of $`L`$. The abstract length $`|s|`$ of $`sL`$ is just $`\mathrm{𝑙𝑒𝑛𝑔𝑡ℎ}(s)`$. The map $`\mathrm{\Phi }:BL`$ is defined by induction as follows:
$`\mathrm{\Phi }(\epsilon )`$ $`=`$ $`\lambda x.\lambda y.\lambda z.z`$
$`\mathrm{\Phi }(0s)`$ $`=`$ $`\lambda x.\lambda y.\lambda z.x\mathrm{\Phi }(s)`$
$`\mathrm{\Phi }(1s)`$ $`=`$ $`\lambda x.\lambda y.\lambda z.y\mathrm{\Phi }(s)`$
Given $`M,N\mathrm{\Lambda }`$, consider the following definitions:
$`M,N`$ $``$ $`\lambda x.xMN`$
$`M_0`$ $``$ $`\lambda x.\lambda y.\lambda z.\lambda w.yx`$
$`M_1`$ $``$ $`\lambda x.\lambda y.\lambda z.\lambda w.zx`$
$`M_{\mathrm{𝑐𝑜𝑚𝑝}}`$ $``$ $`\lambda x.\lambda y.\lambda z.x(yz)`$
$`M_{\mathrm{𝑖𝑑}}`$ $``$ $`\lambda x.x`$
$`M_{\mathrm{𝑐𝑜𝑛𝑠𝑡}}^N`$ $``$ $`\lambda x.N`$
$`M_{\mathrm{𝑡𝑒𝑛𝑠𝑐𝑜𝑛𝑠𝑡}}^N`$ $``$ $`\lambda x.\lambda y.yxM`$
$`M_{\mathrm{𝑡ℎ𝑟𝑜𝑤𝑓𝑖𝑟𝑠𝑡}}`$ $``$ $`\lambda x.x(\lambda y.\lambda z.z)`$
$`M_{\mathrm{𝑠𝑤𝑎𝑝}}`$ $``$ $`\lambda x.x(\lambda y.\lambda w.\lambda z.zwy)`$
$`M_{\mathrm{𝑡𝑒𝑛𝑠}}`$ $``$ $`\lambda x.\lambda y.y(\lambda z.\lambda q.(\lambda y.\lambda w.wyq)(xz))`$
$`M_{\mathrm{𝑎𝑠𝑠𝑙}}`$ $``$ $`\lambda x.x(\lambda y.\lambda w.w(\lambda z.\lambda q.\lambda r.r(\lambda s.syz)q))`$
$`M_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}`$ $``$ $`\lambda x.\lambda y.yxx`$
$`M_{\mathrm{𝑒𝑣𝑎𝑙}}`$ $``$ $`\lambda x.x(\lambda y.\lambda w.yw)`$
$`M_{\mathrm{𝑐𝑢𝑟𝑟𝑦}}`$ $``$ $`\lambda x.\lambda y.\lambda w.x(\lambda z.zyw)`$
*Values* are abstractions and variables. We consider call-by-value reduction on lambda terms, i.e. we take $``$ as the closurure of
$$(\lambda x.M)VM\{x/V\}$$
under all applicative contexts. The application $`\{M\}(N)`$ of two lambda terms is the normal form of $`MN`$ relative to the call-by-value reduction (if one exists). We now define a (ternary) relation $`\mathrm{\Lambda }\times \times \mathrm{\Lambda }`$. In the following, we will write $`M\stackrel{n}{}N`$ standing for $`(M,n,N)`$ The precise definition of $``$ (in SOS-style) follows:
$$\begin{array}{ccccc}\begin{array}{c}\\ \stackrel{\mathit{}}{M\stackrel{0}{}M}\end{array}& & \begin{array}{c}MNn=\mathrm{max}\{1,|N||M|\}\\ \stackrel{\mathit{}}{M\stackrel{n}{}N}\end{array}& & \begin{array}{c}M\stackrel{n}{}NN\stackrel{m}{}L\\ \stackrel{\mathit{}}{M\stackrel{n+m}{}L}\end{array}\end{array}$$
It turns out that for every $`M,N`$ such that $`L`$ is the normal form of $`MN`$, there is exactly one integer $`n`$ such that $`MN\stackrel{n}{}L`$. So, defining $`\mathrm{𝑇𝑖𝑚𝑒}(\{M\}(N))`$ to be just $`n`$ is unambiguous. All the axioms listed at the beginning of this section can be proved to be satisfied by this calculus.
## 3 Length Spaces
In this section, we introduce the category of length spaces and study its properties. Lengths will not necessarily be numbers but rather elements of a commutative monoid.
A *resource monoid* is a quadruple $`M=(|M|,+,_M,𝒟_M)`$ where
* $`(|M|,+)`$ is a commutative monoid;
* $`_M`$ is a pre-order on $`|M|`$ which is compatible with $`+`$;
* $`𝒟_M:\{(\alpha ,\beta )|\alpha _M\beta \}`$ is a function such that for every $`\alpha ,\beta ,\gamma `$
$`𝒟_M(\alpha ,\beta )+𝒟_M(\beta ,\gamma )`$ $``$ $`𝒟_M(\alpha ,\gamma )`$
$`𝒟_M(\alpha ,\beta )`$ $``$ $`𝒟_M(\alpha +\gamma ,\beta +\gamma )`$
and, moreover, for every $`n`$ there is $`\alpha `$ such that $`𝒟_M(0,\alpha )n`$.
Given a resource monoid $`M=(|M|,+,_M,𝒟_M)`$, the function $`_M:|M|`$ is defined by putting $`_M(\alpha )=𝒟_M(0,\alpha )`$. We abbreviate $`\sigma +\mathrm{}+\sigma `$ ($`n`$ times) as $`n.\sigma `$.
Let us try to give some intuition about these axioms. We shall use elements of a resource monoid to bound data, algorithms, and runtimes in the following way: an element $`\phi `$ bounds an algorithm $`e`$ if $`_M(\phi )|e|`$ and, more importantly, whenever $`\alpha `$ bounds an input $`x`$ to $`e`$ then there must be a bound $`\beta _M\phi +\alpha `$ for the result $`y=\{e\}(x)`$ and, most importantly, the runtime of that computation must be bounded by $`𝒟_M(\beta ,\phi +\alpha )`$. So, in a sense, we have the option of either producing a large output fast or to take a long time for a small output. The “inverse triangular” law above ensures that the composition of two algorithms bounded by $`\phi _1`$ and $`\phi _2`$, respectively, can be bounded by $`\phi _1+\phi _2`$ or a simple modification thereof. In particular, the contribution of the unknown intermediate result in a composition cancels out using that law. Another useful intuition is that $`𝒟_M(\alpha ,\beta )`$ behaves like the difference $`\beta \alpha `$, indeed, $`(\beta \alpha )+(\gamma \beta )\gamma \alpha `$.
###### Lemma 1
If $`M`$ is a resource monoid, then $`𝒟_M`$ is antitone on its first argument and monotone on its second argument.
* If $`\alpha _M\beta `$, then
$`𝒟_M(\alpha ,\gamma )`$ $``$ $`𝒟_M(\alpha ,\beta )+𝒟_M(\beta ,\gamma )𝒟_M(\beta ,\gamma );`$
$`𝒟_M(\gamma ,\alpha )`$ $``$ $`𝒟_M(\gamma ,\alpha )+𝒟_M(\alpha ,\beta )𝒟_M(\gamma ,\beta ).`$
This concludes the proof. $`\mathrm{}`$
A *length space* on a resource monoid $`M=(|M|,+,_M,𝒟_M)`$ is a pair $`A=(|A|,_A)`$, where $`|A|`$ is a set and $`_A|M|\times L\times |A|`$ is a (infix) relation satisfying the following conditions:
* If $`\alpha ,e_Aa`$, then $`_M(\alpha )|e|`$;
* For every $`a|A|`$, there are $`\alpha ,e`$ such that $`\alpha ,e_Aa`$
* If $`\alpha ,e_Aa`$ and $`\alpha _M\beta `$, then $`\beta ,e_Aa`$;
* If $`\alpha ,e_Aa`$ and $`\alpha ,e_Ab`$, then $`a=b`$.
The last requirement implies that each element of $`|A|`$ is uniquely determined by the (nonempty) set of it realisers and in particular limits the cardinality of any length space to the number of partial equivalence relations on $`L`$.
A *morphism* from length space $`A=(|A|,_A)`$ to length space $`B=(|B|,_B)`$ (on the same resource monoid $`M=(|M|,+,_M,𝒟_M)`$) is a function $`f:|A||B|`$ such that there exist $`eL=\mathrm{\Sigma }^{}`$, $`\phi |M|`$ with $`_M(\phi )|e|`$ and whenever $`\alpha ,d_Aa`$, there must be $`\beta ,c`$ such that
* $`\beta ,c_Bf(a)`$;
* $`\beta _M\phi +\alpha `$;
* $`\{e\}(d)=c`$;
* $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))𝒟_M(\beta ,\phi +\alpha )`$
We call $`e`$ a realizer of $`f`$ and $`\phi `$ a majorizer of $`f`$. The set of all morphisms from $`A`$ to $`B`$ is denoted as $`\mathrm{𝐻𝑜𝑚}(A,B)`$. If $`f`$ is a morphism from $`A`$ to $`B`$ realized by $`e`$ and majorized by $`\phi `$, then we will write $`f:A\stackrel{e,\phi }{}B`$ or $`\phi ,e_{AB}f`$.
###### Remark 1
It is possible to alter the time bound in the definition of a morphism to $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))𝒟_M(\beta ,\phi +\alpha )_M(\alpha +\phi )`$. This allows one to accommodate linear time operations by padding the majorizer for the morphism. All the subsequent proofs go through with this alternative definition, at the expense of simplicity and ease of presentation,
Given two length spaces $`A=(|A|,_A)`$ and $`B=(|B|,_B)`$ on the same resource monoid $`M`$, we can build $`AB=(|A|\times |B|,_{AB})`$ (on $`M`$) where $`e,\alpha _{AB}(a,b)`$ iff $`_M(\alpha )|e|`$ and there are $`f,g,\beta ,\gamma `$ with
$$\begin{array}{c}f,\beta _Aa\\ g,\gamma _Bb\\ e=f,g\\ \alpha _M\beta +\gamma \end{array}$$
$`AB`$ is a well-defined length space due to the axioms on $`M`$.
Given $`A`$ and $`B`$ as above, we can build $`AB=(\mathrm{𝐻𝑜𝑚}(A,B),_{AB})`$ where $`e,\alpha _{AB}f`$ iff $`f`$ is a morphism from $`A`$ to $`B`$ realized by $`e`$ and majorized by $`\alpha `$.
Morphisms can be composed:
###### Lemma 2 (Composition)
Given length spaces $`A,B,C`$, there is a morphism
$$\mathrm{𝑐𝑜𝑚𝑝}:(BC)(AB)(AC)$$
such that $`\mathrm{𝑐𝑜𝑚𝑝}(f,g)=\lambda x.f(g(x))`$.
* Let $`f:A\stackrel{x,\phi }{}B`$ and $`g:B\stackrel{y,\psi }{}C`$. We know there are constants $`p,q,r`$ such that $`\{e_{\mathrm{𝑐𝑜𝑚𝑝}}\}(x,y)=z`$ where $`|z||x|+|y|+p`$ and $`\{z\}(w)=\{y\}(\{x\}(w))`$; moreover, $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑚𝑝}}\}(x,y))r`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑚𝑝}}\}(w))=\mathrm{𝑇𝑖𝑚𝑒}(\{x\}(w))+\mathrm{𝑇𝑖𝑚𝑒}(\{y\}(\{x\}(w)))+q`$. Now, let us now choose $`\mu `$ such that $`_M(\mu )p+q`$, We will prove that $`comp(f,g):A\stackrel{z,\phi +\psi +\mu }{}C`$. Obviously, $`_M(\phi +\psi +\mu )|z|`$. If $`\alpha ,w_Aa`$, then there must be $`\beta ,t`$ such that $`\beta ,t_Bf(a)`$ and the other conditions prescribed by the definition of a morphism hold. Moreover, there must be $`\gamma ,s`$ such that $`\gamma ,s_Cg(f(a))`$ and, again, the other conditions are satisfied. Putting them together, we get:
$$\gamma _M\beta +\psi _M\alpha +\phi +\psi _M\alpha +\phi +\psi +\mu $$
and
$`\mathrm{𝑇𝑖𝑚𝑒}(\{z\}(w))`$ $``$ $`\mathrm{𝑇𝑖𝑚𝑒}(\{x\}(w))+\mathrm{𝑇𝑖𝑚𝑒}(\{y\}(t))+q`$
$``$ $`𝒟_M(\beta ,\alpha +\phi )+𝒟_M(\gamma ,\beta +\psi )+_M(\mu )`$
$``$ $`𝒟_M(\beta +\psi ,\alpha +\phi +\psi )+𝒟_M(\gamma ,\beta +\psi )+𝒟_M(0,\mu )`$
$``$ $`𝒟_M(\gamma ,\alpha +\phi +\psi +\mu )`$
This concludes the proof, since $`comp:(BC)(AB)\stackrel{(f,g),\xi }{}AC`$ where $`\xi `$ is such that $`_M(\xi )r+|e_{\mathrm{𝑐𝑜𝑚𝑝}}|`$. $`\mathrm{}`$
Basic morphisms can be built independently on the underlying resource monoid. Noticeably, they correspond to axiom of multiplicative linear logic:
###### Lemma 3 (Basic Maps)
Given length spaces $`A,B,C`$, there are morphisms:
$`\mathrm{𝑖𝑑}`$ $`:`$ $`AA`$
$`\mathrm{𝑠𝑤𝑎𝑝}`$ $`:`$ $`ABBA`$
$`\mathrm{𝑎𝑠𝑠𝑙}`$ $`:`$ $`A(BC)(AB)C`$
$`\mathrm{𝑒𝑣𝑎𝑙}`$ $`:`$ $`A(AB)B`$
$`\mathrm{𝑐𝑢𝑟𝑟𝑦}`$ $`:`$ $`((AB)C)A(BC)`$
where
$`\mathrm{𝑖𝑑}(a)`$ $`=`$ $`a`$
$`\mathrm{𝑠𝑤𝑎𝑝}(a,b)`$ $`=`$ $`(b,a)`$
$`\mathrm{𝑎𝑠𝑠𝑙}(a,(b,c))`$ $`=`$ $`((a,b),c)`$
$`\mathrm{𝑒𝑣𝑎𝑙}(a,f)`$ $`=`$ $`f(a)`$
$`\mathrm{𝑐𝑢𝑟𝑟𝑦}(f)`$ $`=`$ $`\lambda a.\lambda b.f(a,b)`$
* We know that $`\{e_{\mathrm{𝑖𝑑}}\}(d)`$ takes constant time, say at most $`p`$. Then, let $`\phi _{\mathrm{𝑖𝑑}}M`$ be such that $`_M(\phi _{\mathrm{𝑖𝑑}})p+|e_{\mathrm{𝑖𝑑}}|`$ (this can always be done). Now, let $`\alpha ,d_Aa`$. We have that $`\alpha ,d_A\mathrm{𝑖𝑑}(a)`$, $`\alpha _M\alpha +\phi _{\mathrm{𝑖𝑑}}`$, $`\{e_{\mathrm{𝑖𝑑}}\}(d)=d`$. Moreover
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑖𝑑}}\}(d))`$ $``$ $`p_M(\phi _{\mathrm{𝑖𝑑}})=𝒟_M(0,\phi _{\mathrm{𝑖𝑑}})`$
$``$ $`𝒟_M(\alpha ,\alpha +\phi _{\mathrm{𝑖𝑑}})`$
This proves $`\mathrm{𝑖𝑑}`$ to be a morphism.
We know that $`\{e_{\mathrm{𝑠𝑤𝑎𝑝}}\}(d,c)`$ takes constant time, say at most $`p`$. Then, let $`\phi _{\mathrm{𝑠𝑤𝑎𝑝}}|M|`$ be such that $`_M(\phi _{\mathrm{𝑖𝑑}})p+|e_{\mathrm{𝑠𝑤𝑎𝑝}}|`$. Now, let $`\alpha ,e_{AB}(a,b)`$. This i that $`e=d,c`$ and $`\alpha ,c,d_{BA}(b,a)`$. We can then apply the same argument as for $`\mathrm{𝑖𝑑}`$. In particular:
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑠𝑤𝑎𝑝}}\}(e))`$ $``$ $`p_M(\phi _{\mathrm{𝑠𝑤𝑎𝑝}})=𝒟_M(0,\phi _{\mathrm{𝑠𝑤𝑎𝑝}})`$
$``$ $`𝒟_M(\alpha ,\alpha +\phi _{\mathrm{𝑠𝑤𝑎𝑝}})`$
This proves $`\mathrm{𝑠𝑤𝑎𝑝}`$ to be a morphism. We can verify $`\mathrm{𝑎𝑠𝑠𝑙}`$ to be a morphism exactly in the same way.
We know that $`\{e_{\mathrm{𝑒𝑣𝑎𝑙}}\}(d,c)=\{d\}(c)`$ and $`\{e_{\mathrm{𝑒𝑣𝑎𝑙}}\}(d,c)`$ takes constant overload time, say at most $`p`$. $`\phi _{\mathrm{𝑒𝑣𝑎𝑙}}`$ is chosen as to satisfy $`_M(\phi _{\mathrm{𝑒𝑣𝑎𝑙}})p`$. Let now $`\alpha ,e_{A(AB)}(a,f)`$. This means that $`e=d,c`$ and there are $`\beta `$ and $`\gamma `$ such that
$$\begin{array}{c}\beta ,d_Aa\\ \gamma ,c_{AB}f\\ \alpha _M\beta +\gamma \\ _M(\alpha )_M(\beta )+_M(\gamma )+\mathrm{𝑐𝑝}\end{array}$$
From $`\gamma ,c_{AB}f`$ it follows that, by the definition of a morphism, there must be $`\delta ,h`$ such that
+ $`\delta ,h_Bf(a)`$
+ $`\delta _M\beta +\gamma `$
+ $`\{c\}(d)=h`$
+ $`\mathrm{𝑇𝑖𝑚𝑒}(\{c\}(d))𝒟_M(\delta ,\beta +\gamma )`$
From $`\delta _M\beta +\gamma `$ and $`\beta +\gamma _M\alpha `$, it follows that $`\delta _M\alpha _M\alpha +\mu `$. Moreover:
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑒𝑣𝑎𝑙}}\}(d,c))`$ $``$ $`p+\mathrm{𝑇𝑖𝑚𝑒}(\{c\}(d))_M(\phi _{\mathrm{𝑒𝑣𝑎𝑙}})+𝒟_M(\delta ,\beta +\gamma )`$
$``$ $`_M(\phi _{\mathrm{𝑒𝑣𝑎𝑙}})+𝒟_M(\delta ,\beta +\gamma )+𝒟_M(\beta +\gamma ,\alpha )`$
$``$ $`𝒟_M(0,\phi _{\mathrm{𝑒𝑣𝑎𝑙}})+𝒟_M(\delta ,\alpha )`$
$``$ $`𝒟_M(\delta ,\alpha +\phi _{\mathrm{𝑒𝑣𝑎𝑙}})`$
Now, let us prove that $`\mathrm{𝑐𝑢𝑟𝑟𝑦}`$ is a morphism. First of all, we know there must be constants $`p,q,r,s,t`$ such that, for each $`e,x,y`$, there are $`d`$ and $`c_x`$ with
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑢𝑟𝑟𝑦}}\}(e))`$ $``$ $`p`$
$`d`$ $`=`$ $`\{e_{\mathrm{𝑐𝑢𝑟𝑟𝑦}}\}(e)`$
$`|d|`$ $``$ $`|e|+q`$
$`\mathrm{𝑇𝑖𝑚𝑒}(\{d\}(x))`$ $``$ $`r`$
$`c_x`$ $`=`$ $`\{d\}(x)`$
$`|c_x|`$ $``$ $`|e|+|x|+s`$
$`\mathrm{𝑇𝑖𝑚𝑒}(\{c_x\}(y))`$ $``$ $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(x,y))+t`$
$`\{e\}(x,y)`$ $`=`$ $`\{c_x\}(y)`$
Let $`\mu ,\theta ,\xi |M|`$ be such that
$`_M(\xi )`$ $``$ $`p`$
$`_M(\mu )`$ $``$ $`q`$
$`_M(\sigma )`$ $``$ $`r`$
$`_M(\theta )`$ $``$ $`s`$
$`_M(\eta )`$ $``$ $`t`$
$`_M(\chi )`$ $``$ $`\mathrm{𝑐𝑝}`$
Let now $`\gamma ,e_{ABC}f`$. We know that $`|d||e|+q`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑢𝑟𝑟𝑦}}\}(e))p`$. In order to prove that $`\mathrm{𝑐𝑢𝑟𝑟𝑦}`$ is indeed a morphism realized by $`e_{\mathrm{𝑐𝑢𝑟𝑟𝑦}}`$ and majorized by $`\mu +\xi +\sigma +\theta +\chi +\eta `$, it then suffices to prove that
$$\gamma +\mu +\sigma +\theta +\chi +\theta ,d_{ABC}\lambda a.\lambda b.f(a,b).$$
Let then $`\alpha ,x_Aa`$. There is $`c_x`$ such that $`c_x=\{d\}(x)`$, $`|c_x||e|+|x|+s`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{d\}(x))r`$. In order to prove that $`\lambda a.\lambda b.f(a,b)`$ is indeed a morphism realized by $`d`$ and majorized by $`\gamma +\mu +\sigma +\theta +\chi +\eta `$, it then suffices to prove that $`\gamma +\alpha +\mu +theta+\chi +\eta ,c_x_{BC}\lambda b.f(a,b)`$. Let then $`\beta ,y_Bb`$. There are $`\delta ,c`$ such $`\delta ,c_Cf(a,b)`$, where $`\delta \alpha +\beta +\chi +\gamma `$. Moreover, we know that
$`\mathrm{𝑇𝑖𝑚𝑒}(\{c_x\}(y))`$ $``$ $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(x,y))+t𝒟_M(\delta ,\alpha +\beta +\chi +\gamma )+t`$
$``$ $`𝒟_M(\delta ,\alpha +\beta +\gamma +\chi )+𝒟_M(0,\eta +\mu +\theta )`$
$``$ $`𝒟_M(\delta ,\alpha +\beta +\gamma +\chi +\eta +\mu +\theta )`$
This concludes the proof. $`\mathrm{}`$.
Length spaces can justify the usual rule for tensor as a map-former:
###### Lemma 4 (Tensor)
Given length spaces $`A,B,C`$, there is a morphism
$$\mathrm{𝑡𝑒𝑛𝑠}:(AB)((AC)(BC))$$
where $`\mathrm{𝑡𝑒𝑛𝑠}(f)=\lambda x.(f(\pi _1(x)),\pi _2(x))`$.
* Let $`f:A\stackrel{x,\phi }{}B`$. We know there are constants $`p,q`$ such that $`\{e_{\mathrm{𝑡𝑒𝑛𝑠}}\}(x)=y`$ where $`|y||x|+p`$ and $`\{y\}(z,w)=\{x\}(z),w`$; moroever, $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑡𝑒𝑛𝑠}}\}(x))q`$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{y\}(z,w))\mathrm{𝑇𝑖𝑚𝑒}(\{x\}(z))+r`$. Then, take $`\psi |M|`$ such that $`_M(\psi )p+r`$, put $`\sigma =\psi +\phi +\mu `$, where $`_M(\mu )\mathrm{𝑐𝑝}`$. Suppose $`\alpha ,z,w_{AC}(a,c)`$. By definition, there are $`\beta ,\gamma `$ such that
$$\begin{array}{c}\beta ,z_Aa\\ \gamma ,w_Cc\\ \alpha _M\beta +\gamma \end{array}$$
By hypothesis, there are $`\delta ,t`$ such that
$$\begin{array}{c}\delta ,t_Bf(a)\\ \delta _M\phi +\beta \\ \{e\}(z)=t\\ \mathrm{𝑇𝑖𝑚𝑒}(\{e\}(z))𝒟_M(\delta ,\phi +\beta )\end{array}$$
Then, $`\gamma +\delta +\mu ,t,w_{BC}(f(a),c)`$. Moreover,
$$\gamma +\delta +\mu _M\gamma +\phi +\beta +\mu _M\alpha +\phi +\mu _M\alpha +\sigma $$
Finally:
$`\mathrm{𝑇𝑖𝑚𝑒}(\{y\}(z,w))`$ $``$ $`\mathrm{𝑇𝑖𝑚𝑒}(\{x\}(z))+r`$
$``$ $`𝒟_M(\delta ,\phi +\beta )+_M(\psi )`$
$``$ $`𝒟_M(\delta ,\phi +\beta +\psi )`$
$``$ $`𝒟_M(\gamma +\delta +\mu ,\gamma +\phi +\beta +\mu +\psi )`$
$`=`$ $`𝒟_M(\gamma +\delta +\mu ,\gamma +\beta +\sigma )`$
$`=`$ $`𝒟_M(\gamma +\delta +\mu ,\alpha +\sigma )`$
This concludes the proof, since $`tens:(AB)\stackrel{(f,g),\xi }{}(AC)(BC)`$ where $`\xi `$ is such that $`_M(\xi )q+|e_{\mathrm{𝑡𝑒𝑛𝑠}}|`$. $`\mathrm{}`$ . $`\mathrm{}`$
Thus:
###### Lemma 5
Length spaces and their morphisms form a symmetric monoidal closed category with tensor and linear implication given as above.
A length space $`I`$ is defined by $`|I|=\{0\}`$ and $`\alpha ,e_A0`$ when $`_M(\alpha )|e|`$. For each length space $`A`$ there are isomorphisms $`AIA`$ and a unique morphism $`AI`$. The latter serves to justify full weakening.
For every resource monoid $`M`$, there is a length space $`B_M=(\{0,1\}^{},_{B_M})`$ where $`\alpha ,\mathrm{\Phi }(t)_{B_M}t`$ whenever $`_M(\alpha )|t|`$. The function $`s_0`$ (respectively, $`s_1`$) from $`\{0,1\}^{}`$ to itself which appends $`0`$ (respectively, $`1`$) to the left of its argument can be computed in constant time on the abstract computational model and, as a consequence, is a morphism from $`B_M`$ to itself.
### 3.1 Interpreting Multiplicative Affine Logic
We can now formally show that second order multiplicative affine logic (i.e. multiplicative linear logic plus full weakening) can be interpreted inside the category of length spaces on any monoid $`M`$. Doing this will simplify the analysis of richer systems presented in following sections. Formulae of (intuitionistic) multiplicative affine logic are generated by the following productions:
$$A::=\alpha |AA|AA|\alpha .A$$
where $`\alpha `$ ranges over a countable set of atoms. Rules are reported in figure 1.
A *realizability environment* is a partial function assigning length spaces (on the same resource monoid) to atoms. Realizability semantics $`A_\eta ^{}`$ of a formula $`A`$ on the realizability environment $`\eta `$ is defined by induction on $`A`$:
$`\alpha _\eta ^{}`$ $`=`$ $`\eta (\alpha )`$
$`AB_\eta ^{}`$ $`=`$ $`A_\eta ^{}B_\eta ^{}`$
$`AB_\eta ^{}`$ $`=`$ $`A_\eta ^{}B_\eta ^{}`$
$`\alpha .A_\eta ^{}`$ $`=`$ $`(|\alpha .A_\eta ^{}|,_{\alpha .A_\eta ^{}})`$
where
$`|\alpha .A_\eta ^{}|`$ $`=`$ $`{\displaystyle \underset{C𝒰}{}}|A_{\eta [\alpha C]}^{}|`$
$`\alpha ,e_{\alpha .A_\eta ^{}}a`$ $``$ $`C.\alpha ,e_{A_{\eta [\alpha C]}^{}}a`$
Here $`𝒰`$ stands for the class of all length spaces. A little care is needed when defining the product since strictly speaking it does not exist for size reasons. The standard way out is to let the product range over those length spaces whose underlying set equals the set of equivalence classes of a partial equivalence relation on $`L`$. As already mentioned, every length space is isomorphic to one such. When working with the product one has to insert these isomorphisms in appropriate places which, however, we elide to increase readability.
If $`n0`$ and $`A_1,\mathrm{},A_n`$ are formulas, the expression $`A_1\mathrm{}A_n_\eta ^{}`$ stands for $`I`$ if $`n=0`$ and $`A_1\mathrm{}A_{n1}_\eta ^{}A_n_\eta ^{}`$ if $`n1`$.
## 4 Elementary Length Spaces
In this section, we define a resource monoid $``$ such that elementary affine logic can be interpreted in the category of length spaces on $``$. We then (re)prove that functions representable in EAL are elementary time computable.
A *list* is either $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ or $`\mathrm{𝑐𝑜𝑛𝑠}(n,l)`$ where $`n`$ and $`l`$ is itself a list. The sum $`l+h`$ of two lists $`l`$ and $`h`$ is defined as follows, by induction on $`l`$:
$`\mathrm{𝑒𝑚𝑝𝑡𝑦}+h=h+\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ $`=`$ $`h`$
$`\mathrm{𝑐𝑜𝑛𝑠}(n,l)+\mathrm{𝑐𝑜𝑛𝑠}(m,h)`$ $`=`$ $`\mathrm{𝑐𝑜𝑛𝑠}(n+m,l+h)`$
For every $`e`$, binary relations $`_e`$ on lists can be defined as follows
* $`\mathrm{𝑒𝑚𝑝𝑡𝑦}_el`$;
* $`\mathrm{𝑐𝑜𝑛𝑠}(n,l)_e\mathrm{𝑐𝑜𝑛𝑠}(m,h)`$ iff there is $`d`$ such that
+ $`n3^e(m+e)d`$;
+ $`l_dh`$.
For every $`e`$ and for every lists $`l`$ and $`h`$ with $`l_eh`$, we define the natural number $`𝒟_e(l,h)`$ as follows:
$`𝒟_e(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})`$ $`=`$ $`0;`$
$`𝒟_e(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑐𝑜𝑛𝑠}(n,l))`$ $`=`$ $`3^e(n+e)+𝒟_{3^e(n+e)}(\mathrm{𝑒𝑚𝑝𝑡𝑦},l);`$
$`𝒟_e(\mathrm{𝑐𝑜𝑛𝑠}(n,l),\mathrm{𝑐𝑜𝑛𝑠}(m,h))`$ $`=`$ $`3^e(m+e)n+𝒟_{3^e(m+e)n}(l,h);`$
Given a list $`l`$, $`!l`$ stands for the list $`\mathrm{𝑐𝑜𝑛𝑠}(0,l)`$. The depth $`0ptl`$ of a list $`l`$ is defined by induction on $`l`$: $`0pt\mathrm{𝑒𝑚𝑝𝑡𝑦}=0`$ while $`0pt\mathrm{𝑐𝑜𝑛𝑠}(n,l)=0ptl+1`$. $`|l|`$ stands for the maximum integer appearing inside $`l`$, i.e. $`|\mathrm{𝑒𝑚𝑝𝑡𝑦}|=0`$ and $`|\mathrm{𝑐𝑜𝑛𝑠}(n,l)|=\mathrm{max}\{n,|l|\}`$. For every natural number $`n`$, $`[n]_{}`$ stands for $`\mathrm{𝑐𝑜𝑛𝑠}(n,\mathrm{𝑒𝑚𝑝𝑡𝑦})`$.
We can now verify that all the necessary conditions required by the definition of a resource monoid are satisfied. To do this, we need a number of preliminary results, which can all be proved by simple inductions and case-analysis:
###### Lemma 6 (Compatibility)
$`\mathrm{𝑒𝑚𝑝𝑡𝑦}_el`$ for every $`l`$. Moreover, if $`l,h,j`$ are lists and $`l_eh`$, then $`l+j_eh+j`$.
* The first claim is trivial. To prove the second, we proceed by an induction on $`j`$. If $`j=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, then $`l+j=l_eh=h+j`$. Now, suppose $`j=\mathrm{𝑐𝑜𝑛𝑠}(n,g)`$. If $`h=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, then $`l=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ and, clearly $`l+j=j_ej=h+j`$. If $`l=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, we have to prove that $`j_eh+j`$. Let $`h=\mathrm{𝑐𝑜𝑛𝑠}(m,f)`$; then
$`n`$ $``$ $`n+m3^e(n+m+e)0`$
$`g`$ $`_0`$ $`g+f`$
which means $`j_eh+j`$. Finally, suppose $`l=\mathrm{𝑐𝑜𝑛𝑠}(m,f)`$, $`h=\mathrm{𝑐𝑜𝑛𝑠}(p,r)`$. Then we know that
$`m`$ $``$ $`3^e(p+e)d`$
$`f`$ $`_d`$ $`r`$
But then, by inductive hypothesis,
$`m+n`$ $``$ $`3^e(p+e)+nd3^e(p+n+e)d`$
$`f+g`$ $`_d`$ $`r+g`$
which yields $`l+j_eh+j`$. $`\mathrm{}`$
###### Lemma 7 (Transitivity)
If $`l,h,j`$ are lists and $`l_eh`$, $`h_dj`$, then $`l_{d+e}j`$.
* We can suppose all the involved lists to be different from $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, since all the other cases are trivial. $`l=\mathrm{𝑐𝑜𝑛𝑠}(n,g)`$, $`h=\mathrm{𝑐𝑜𝑛𝑠}(m,f)`$ and $`j=\mathrm{𝑐𝑜𝑛𝑠}(p,r)`$. From the hypothesis, we have
$`n`$ $``$ $`3^e(m+e)c`$
$`m`$ $``$ $`3^d(p+d)b`$
$`g`$ $`_c`$ $`f`$
$`f`$ $`_b`$ $`r`$
But then, by inductive hypothesis, we get
$`n`$ $``$ $`3^e(m+e)c3^e(3^d(p+d)b+e)c3^e3^d(p+d+e)bc=3^{e+d}(p+d+e)(b+c)`$
$`g`$ $`_{c+b}`$ $`r`$
This means $`l_{d+e}j`$. $`\mathrm{}`$
###### Lemma 8
if $`l,h,j`$ are lists and $`l_eh`$, then $`𝒟_e(l,h)𝒟_e(l+j,h+j)`$
* We proceed by an induction on $`j`$. If $`j=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, then $`l+j=l`$ and $`h+j=h`$. Now, suppose $`j=\mathrm{𝑐𝑜𝑛𝑠}(n,g)`$. If $`h=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, then $`l=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ and, clearly $`l+j=j=h+j`$. If $`l=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, let $`h=\mathrm{𝑐𝑜𝑛𝑠}(m,f)`$; then
$`𝒟_e(l,h)`$ $`=`$ $`𝒟_e(\mathrm{𝑒𝑚𝑝𝑡𝑦},h)=3^e(m+e)+𝒟_{3^e(m+e)}(\mathrm{𝑒𝑚𝑝𝑡𝑦},f)`$
$``$ $`3^e(m+e)+3^en3^en+𝒟_{3^e(m+e)+3^en3^en}(g,g+f)`$
$``$ $`3^e(m+n+e)n+𝒟_{3^e(m+n+e)n}(g,g+f)`$
$`=`$ $`𝒟_e(j,h+j)=𝒟_e(l+j,h+j)`$
Finally, suppose $`l=\mathrm{𝑐𝑜𝑛𝑠}(m,f)`$, $`h=\mathrm{𝑐𝑜𝑛𝑠}(p,r)`$. Then we know that
$`𝒟_e(l,h)`$ $`=`$ $`3^e(m+e)p+𝒟_{3^e(m+e)p}(f,r)`$
$``$ $`3^e(m+e)p+𝒟_{3^e(m+e)p}(f+g,r+g)`$
$``$ $`3^e(m+e)+3^en(n+p)+𝒟_{3^e(m+e)+3^ennp}(f+g,r+g)`$
$`=`$ $`3^e(m+n+e)(n+p)+𝒟_{3^e(m+n+e)(n+p)}(f+g,r+g)`$
$`=`$ $`𝒟_e(l+j,h+j)`$
###### Lemma 9
If $`l,h,j`$ are lists and $`l_eh`$, $`h_dj`$, then $`𝒟_e(l,h)+𝒟_d(h,j)𝒟_{e+d}(l,j)`$.
* If either $`h=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ or $`j=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, then the thesis is trivial. So suppose $`h=\mathrm{𝑐𝑜𝑛𝑠}(n,g)`$ and $`j=\mathrm{𝑐𝑜𝑛𝑠}(m,f)`$. If $`l=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, then
$`𝒟_e(l,h)+𝒟_d(h,j)`$ $`=`$ $`3^e(n+e)+𝒟_{3^e(n+e)}(\mathrm{𝑒𝑚𝑝𝑡𝑦},g)+3^d(m+d)n+𝒟_{3^e(m+d)n}(g,f)`$
$``$ $`3^e(n+e)+3^d(m+d)n+𝒟_{3^e(n+e)+3^d(m+d)n}(\mathrm{𝑒𝑚𝑝𝑡𝑦},f)`$
$``$ $`(3^e1)n+3^ee+3^d(m+d)+𝒟_{(3^e1)n++3^ee+3^d(m+d)}(\mathrm{𝑒𝑚𝑝𝑡𝑦},f)`$
$``$ $`(3^e1)3^d(m+d)+3^ee+3^d(m+d)+𝒟_{(3^e1)3^d(m+d)+3^ee+3^d(m+d)}(\mathrm{𝑒𝑚𝑝𝑡𝑦},f)`$
$`=`$ $`3^{d+e}(m+d+e)+𝒟_{3^{d+e}(m+d+e)}(\mathrm{𝑒𝑚𝑝𝑡𝑦},f)`$
$`=`$ $`𝒟_{e+d}(l,j)`$
If $`l=\mathrm{𝑐𝑜𝑛𝑠}(p,r)`$, then
$`𝒟_e(l,h)+𝒟_d(h,j)`$ $`=`$ $`3^e(n+e)p+𝒟_{3^e(n+e)p}(r,g)+3^d(m+d)n+𝒟_{3^d(m+d)n}(g,f)`$
$``$ $`3^e(n+e)p+3^d(m+d)n+𝒟_{3^e(n+e)p+3^d(m+d)n}(r,f)`$
$``$ $`(3^e1)n+3^ee+3^d(m+d)p+𝒟_{(3^e1)n+3^ee+3^d(m+d)p}(r,f)`$
$``$ $`(3^e1)3^d(m+d)+3^ee+3^d(m+d)p+𝒟_{(3^e1)3^d(m+d)+3^ee+3^d(m+d)p}(r,f)`$
$`=`$ $`3^{d+e}(m+d+e)p+𝒟_{3^{d+e}(m+d+e)p}(r,f)`$
$`=`$ $`𝒟_{e+d}(l,j)`$
This concludes the proof. $`\mathrm{}`$
$`||`$ will denote the set of all lists, while $`_{},𝒟_{}`$ will denote $`_0`$ and $`𝒟_0`$, respectively.
###### Lemma 10
$`=(||,+,_{},𝒟_{})`$ is a resource monoid.
* $`(,+)`$ is certainly a monoid. Compatibility of $`_{}`$ follows from lemmas 6 and 7. The two required property on $`𝒟_{}`$ come directly from lemmas 8 and 9. If $`n`$, observe that $`_{}(\mathrm{𝑐𝑜𝑛𝑠}(n,\mathrm{𝑒𝑚𝑝𝑡𝑦}))=n`$. This concludes the proof. $`\mathrm{}`$
An *elementary length space* is a length space on the resource monoid $`(||,+,_{},𝒟_{})`$. Given an elementary length space $`A=(|A|,_A)`$, we can build the length space $`!A=(|A|,_{!A})`$, where $`l,e_{!A}a`$ iff $`h,e_Aa`$ and $`l_{}!h`$. The construction $`!`$ on elementary length spaces serves to capture the exponential modality of elementary affine logic. Indeed, the following two results prove the existence of morphisms and morphisms-forming rules precisely corresponding to axioms and rules from EAL.
###### Lemma 11
For every $`e`$ and for every $`l`$, $`l+l_1l`$ and $`𝒟_{e+1}(l+l,l)𝒟_e(0,l)`$.
* The inequality $`l+l_1l`$ can be proved by induction on $`l`$. The base case is trivial. If $`l=\mathrm{𝑐𝑜𝑛𝑠}(n,h)`$, then
$`n+n`$ $``$ $`3n+31=3^1(n+1)1`$
$`h+h_1`$ $`h`$
The second inequality can be proved by induction on $`l`$, too. The base case is trivial. If $`l=\mathrm{𝑐𝑜𝑛𝑠}(n,h)`$, observe that
$`𝒟_{e+1}(l+l,l)`$ $`=`$ $`3^{e+1}(n+e+1)2n+𝒟_{3^{e+1}(n+e+1)2n}(h+h,h)`$
$`𝒟_e(0,l)`$ $`=`$ $`3^e(n+e)+𝒟_{3^e(n+e)}(0,h)`$
But
$`3^{e+1}(n+e+1)2n`$ $`=`$ $`3^e(n+e+1)+2(3^e)(n+e+1)2n`$
$``$ $`3^e(n+e+1)+2n2n3^e(n+e)+1`$
This concludes the proof. $`\mathrm{}`$
###### Lemma 12 (Basic Maps)
Given elementary length spaces $`A,B`$, there are morphisms:
$`\mathrm{𝑐𝑜𝑛𝑡𝑟}`$ $`:`$ $`!A!A!A`$
$`\mathrm{𝑑𝑖𝑠𝑡𝑟}`$ $`:`$ $`!A!B!(AB)`$
where $`\mathrm{𝑐𝑜𝑛𝑡𝑟}(a)=(a,a)`$ and $`\mathrm{𝑑𝑖𝑠𝑡𝑟}(a,b)=(a,b)`$
* We know $`\{e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}\}(d)`$ takes time $`|d|+p`$, where $`p`$ is a constant. Then, let $`l,h`$ be such that $`_{}(l)p+|e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}|`$, $`_{}(h)\mathrm{𝑐𝑝}`$. Define $`l_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}`$ to be $`l+h+[1]_{}`$. Clearly, $`_{}(l_{\mathrm{𝑐𝑜𝑛𝑡𝑟}})|e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}|`$ Now, let $`j,d_{!A}a`$. This implies that $`j_{}!k`$ where $`k,d_Aa`$. Then:
$`h+!k+!k`$ $`_{}`$ $`!k+!k`$
$`_{}(h+!k+!k)`$ $``$ $`_{}(h)+_{}(!k)+_{}(!k)`$
$``$ $`\mathrm{𝑐𝑝}+_{}(!k)+_{}(!k)`$
This yields $`h+!k+!k,e_{!A!A}(a,a)`$. By lemma 11, $`h+!k+!k_{}h+!k+[1]_{}_{}h+j+[1]_{}_{}j+l_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}`$. Finally,
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}\}(d))`$ $``$ $`|d|+p_{}(k)+p𝒟_{}(!k+!k,!k+[1]_{})+_{}(l)`$
$``$ $`𝒟_{}(!k+!k,!k+[1]_{}+l)`$
$``$ $`𝒟_{}(!k+!k+h,!k+[1]_{}+l+h)`$
$`=`$ $`𝒟_{}(!k+!k+h,!k+l_{\mathrm{𝑐𝑜𝑛𝑡𝑟}})`$
This proves $`\mathrm{𝑐𝑜𝑛𝑡𝑟}`$ to be a morphism.
Let $`e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}=e_{\mathrm{𝑖𝑑}}`$. We know $`\{e_{\mathrm{𝑖𝑑}}\}(d)`$ takes constant time, say $`p`$. Then, let $`l,h`$ be such that $`_{}(l)p+|e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}|`$, $`_{}(h)\mathrm{𝑐𝑝}`$. $`l_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}`$ is then defined as $`l+!h`$. Now, let $`j,d,c_{!A!B}(a,b)`$. This means that $`j!k+!i`$, where $`k,d_Aa`$ and $`i,c_Bb`$. This in turn means that $`k+i+h,d,c_{AB}(a,b)`$ and $`!(k+i+h),d,c_{!(AB)}(a,b)`$. Moreover
$$!(k+i+h)=!k+!i+!h_{}j+!h_{}j+l_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}$$
Finally:
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}\}(d,c))`$ $``$ $`p_{}(l)`$
$``$ $`𝒟_{}(!(k+i+h),j+!h)+_{}(l)`$
$``$ $`𝒟_{}(!(k+i+h),j+!h+l)`$
$``$ $`𝒟_{}(!(k+i+h),j+l_{\mathrm{𝑑𝑖𝑠𝑡𝑟}})`$
This proves $`\mathrm{𝑑𝑖𝑠𝑡𝑟}`$ to be a morphism. $`\mathrm{}`$
###### Lemma 13 (Functoriality)
If $`f:A\stackrel{e,\phi }{}B`$, then there is $`\psi `$ such that $`f:!A\stackrel{e,\psi }{}!B`$
* Let $`\theta `$ be $`!\phi `$ and suppose $`d,l_{!A}a`$. Then $`l!h`$, where $`d,h_Aa`$. Observe that there must be $`j,c`$ such that $`c,j_Bf(a)`$, $`j_{}h+\phi `$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))𝒟_{}(j,h+\phi )`$. But then $`c,!j_{!B}f(a)`$ and, moreover
$`!j`$ $`_{}`$ $`!(h+\phi )=!h+!\phi _{}!h+\theta `$
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))`$ $``$ $`𝒟_{}(j,h+\phi )𝒟_{}(!j,!(h+\phi ))`$
$``$ $`𝒟_{}(!j,!h+!\phi ))𝒟_{}(!j,l+\theta )`$
This means that $`f:!A\stackrel{e,\theta }{}!B`$. $`\mathrm{}`$
Elementary bounds can be given on $`_{}(l)`$ depending on $`|l|`$ and $`0ptl`$:
###### Proposition 1
For every $`n`$ there is an elementary function $`p_n:`$ such that $`_{}(l)p_{0ptl}(|l|)`$.
* We prove a stronger statement by induction on $`n`$: for every $`n`$ there is an elementary function $`q_n:^2`$ such that for every $`l,e`$, $`𝒟_e(\mathrm{𝑒𝑚𝑝𝑡𝑦},l)q_{0ptl}(|l|,e)`$. First of all, we know that $`𝒟_e(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})=0`$, so $`q_0`$ is just the function which always returns $`0`$. $`q_{n+1}`$ is defined from $`q_n`$ as follows: $`q_{n+1}(x,y)=3^y(x+y)+q_n(x,3^y(x+y))`$. Indeed:
$`𝒟_e(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑐𝑜𝑛𝑠}(n,l))`$ $`=`$ $`3^e(n+e)+𝒟_{3^e(n+e)}(\mathrm{𝑒𝑚𝑝𝑡𝑦},l)`$
$``$ $`3^e(|\mathrm{𝑐𝑜𝑛𝑠}(n,l)|+e)+q_{0ptl}(|l|,3^e(n+e))`$
$``$ $`3^e(|\mathrm{𝑐𝑜𝑛𝑠}(n,l)|+e)+q_{0ptl}(|\mathrm{𝑐𝑜𝑛𝑠}(n,l)|,3^e|\mathrm{𝑐𝑜𝑛𝑠}(n,l)|+e)`$
$`=`$ $`q_{0pt\mathrm{𝑐𝑜𝑛𝑠}(n,l)}(|\mathrm{𝑐𝑜𝑛𝑠}(n,l)|,e)`$
At this point we just put $`p_n(x)=q_n(x,0)`$. $`\mathrm{}`$
We emphasize that Proposition 1 does not assert that the mapping $`(n,m)p_n(m)`$ is elementary. This, indeed, cannot be true because we know EAL to be complete for the class of elementary functions. If, however, $`A`$ is such that $`lA`$ implies $`0ptlc`$ for a fixed $`c`$, then $`(lA)p_{0ptl}(|l|)`$ is elementary and it is in this way that we will use the above proposition.
### 4.1 Interpreting Elementary Affine Logic
EAL can be obtained by endowing multiplicative affine logic with a restricted modality. The grammar of formulae is enriched with a new production $`A::=!A`$ while modal rules are reported in figure 2.
Realizability semantics is extended by $`!A_\eta ^{}=!A_\eta ^{}`$.
###### Theorem 1
Elementary length spaces form a model of EAL.
Now, consider the formula
$$\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖤𝖠𝖫}}\alpha .!(\alpha \alpha )!(\alpha \alpha )!(\alpha \alpha )$$
Binary lists can be represented as cut-free proofs with conclusion $`\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖤𝖠𝖫}}`$. Suppose you have a proof $`\pi :!^j\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖤𝖠𝖫}}!^k\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖤𝖠𝖫}}`$. From the denotation $`\pi ^{}`$ we can build a morphism $`g`$ from $`\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖤𝖠𝖫}}^{}`$ to $`B_{}`$ by internal application to $`\epsilon ,s_0,s_1`$. This map then induces a function $`f:BB`$ as follows: given $`wB`$, first compute a realizer for the closed proof corresponding to it, then apply $`g`$ to the result.
###### Remark 2
Notice that elements of $`B_{}`$ can all be majorized by lists with unit depth. Similarly, elements of $`\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖤𝖠𝖫}}^{}`$ corresponding to binary lists can be majorized by lists with bounded depth. This observation is essential to prove the following result.
###### Corollary 1 (Soundness)
Let $`\pi `$ be an EAL proof with conclusion $`!^j\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖤𝖠𝖫}}!^k\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖤𝖠𝖫}}`$ and let $`f:LL`$ be the function induced by $`\pi ^{}`$. Then $`f`$ is computable in elementary time.
The function $`f`$ in the previous result equals the function denoted by the proof $`\pi `$ in the sense of . This intuitively obvious fact can be proved straightforwardly but somewhat tediously using a logical relation or similar, see also .
## 5 Soft Length Spaces
The grammar of formulae for SAL is the same as the one of Elementary Affine Logic. Rules are reported in figure 3.
We here use a resource monoid whose underlying carrier set is $`||=||\times `$. The sum $`(l,n)+(h,m)`$ of two elements in $`||`$ is defined as $`(l+h,\mathrm{max}\{n,m\})`$. For every $`e`$, binary relations $`_e`$ on $`||`$ can be defined as follows
* $`(\mathrm{𝑒𝑚𝑝𝑡𝑦},n)_0(\mathrm{𝑒𝑚𝑝𝑡𝑦},m)`$ iff $`nm`$;
* $`(\mathrm{𝑒𝑚𝑝𝑡𝑦},n)_e(\mathrm{𝑐𝑜𝑛𝑠}(m,l),p)`$ iff there is $`d`$ such that
+ $`em+pd`$
+ $`(\mathrm{𝑒𝑚𝑝𝑡𝑦},n)_d(l,p)`$
* $`(\mathrm{𝑐𝑜𝑛𝑠}(n,l),m)_e(\mathrm{𝑐𝑜𝑛𝑠}(p,h),q)`$ iff there is $`d`$ such that
+ $`e+np+qd`$;
+ $`(l,m)_d(h,q)`$.
If $`\alpha =(l,n)||`$, then $`!\alpha `$ will be the couple $`(\mathrm{𝑐𝑜𝑛𝑠}(0,l),n)||`$. If there is $`e`$ such that $`\alpha _e\beta `$, then we will simply write $`\alpha _{}\beta `$. For every $`\alpha `$ and $`\beta `$ with $`\alpha _{}\beta `$, we define the natural number $`𝒟_{}(\alpha ,\beta )`$ as follows:
$`𝒟_{}((\mathrm{𝑒𝑚𝑝𝑡𝑦},n),(\mathrm{𝑒𝑚𝑝𝑡𝑦},m))`$ $`=`$ $`0`$
$`𝒟_{}((\mathrm{𝑒𝑚𝑝𝑡𝑦},n),(\mathrm{𝑐𝑜𝑛𝑠}(m,l),p))`$ $`=`$ $`m+p𝒟_{}((\mathrm{𝑒𝑚𝑝𝑡𝑦},n),(l,p))`$
$`𝒟_{}((\mathrm{𝑐𝑜𝑛𝑠}(n,l),m),(\mathrm{𝑐𝑜𝑛𝑠}(p,h),q))`$ $`=`$ $`pn+q𝒟_{}((l,m),(h,q))`$
Analogously, we can define $`𝒟_{}(\alpha ,\beta )`$ simply as the maximum integer $`e`$ such that $`\alpha _e\beta `$. $`|\alpha |`$ is the maximum integer appearing inside $`\alpha `$, i.e. $`|(l,n)|=\mathrm{max}\{|l|,m\}`$. The depth $`0pt\alpha `$ of $`\alpha =(l,n)`$ is $`0ptl`$.
###### Lemma 14 (Compatibility)
$`(\mathrm{𝑒𝑚𝑝𝑡𝑦},0)_0\alpha `$ for every $`\alpha `$. Moreover, if $`\alpha ,\beta ,\gamma ||`$ and $`\alpha _e\beta `$, then $`\alpha +\gamma _e\beta +\gamma `$.
* The first claim is trivial. To prove the second, we proceed by an induction on the structure of the first component of $`\gamma `$. We just consider the case where the first components of $`\alpha ,\beta ,\gamma `$ are all different from $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$. So, suppose $`\alpha =(\mathrm{𝑐𝑜𝑛𝑠}(n,l),m)`$, $`\beta =(\mathrm{𝑐𝑜𝑛𝑠}(p,h),q)`$, $`\gamma =(\mathrm{𝑐𝑜𝑛𝑠}(r,j),s)`$. By hypothesis, we get $`d`$ such that
$`e+n`$ $``$ $`p+dq`$
$`(l,m)`$ $`_d`$ $`(h,q)`$
Then, $`e+n+rp+r+dqp+r+d\mathrm{max}\{q,s\}`$ and, by induction hypothesis, $`(l+j,\mathrm{max}\{m,s\})_d(h+j,\mathrm{max}\{q,s\})`$. This implies that $`\alpha +\gamma _e\beta +\gamma `$. $`\mathrm{}`$
###### Lemma 15 (Transitivity)
If $`\alpha ,\beta ,\gamma ||`$ are lists and $`\alpha _e\beta `$, $`\beta _d\gamma `$, then $`\alpha _{d+e}\gamma `$.
* We go by induction on the structure of the first component of $`\gamma `$ and we suppose the first components of $`\alpha ,\beta ,\gamma `$ to be different from $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$. So, let $`\alpha =(\mathrm{𝑐𝑜𝑛𝑠}(n,l),m)`$, $`\beta =(\mathrm{𝑐𝑜𝑛𝑠}(p,h),q)`$ and $`\gamma =(\mathrm{𝑐𝑜𝑛𝑠}(r,j),s)`$. From the hypothesis, there are $`c,b`$ such that
$`e+n`$ $``$ $`p+cq`$
$`d+p`$ $``$ $`r+bs`$
$`(l,m)`$ $`_c`$ $`(h,q)`$
$`(h,q)`$ $`_b`$ $`(j,s)`$
But then, by inductive hypothesis, we get
$`(e+d)+n`$ $``$ $`d+p+cqr+bs+cqr+(b+c)s`$
$`(l,m)`$ $`_{c+b}`$ $`(j,s)`$
which yields $`\alpha _{d+e}\gamma `$. $`\mathrm{}`$
###### Lemma 16
if $`\alpha ,\beta ,\gamma `$ and $`\alpha _e\beta `$, then $`𝒟_{}(\alpha ,\beta )𝒟_{}(\alpha +\gamma ,\beta +\gamma )`$
* This is trivial in view of 14 and the fact that $`𝒟_{}(\alpha ,\beta )`$ is just $`\mathrm{max}\{e|\alpha _e\beta \}`$. $`\mathrm{}`$
###### Lemma 17
If $`\alpha ,\beta ,\gamma `$ and $`\alpha _e\beta `$, $`\beta _d\gamma `$, then $`𝒟_e(\alpha ,\beta )+𝒟_d(\beta ,\gamma )𝒟_{e+d}(\alpha ,\gamma )`$.
* This is trivial in view of 15 and the fact that $`𝒟_{}(\alpha ,\beta )`$ is just $`\mathrm{max}\{e|\alpha _e\beta \}`$. $`\mathrm{}`$
###### Lemma 18
$`(,+,_{},𝒟_{})`$ is a resource monoid.
* $`(||,+)`$ is certainly a commutative monoid. Compatibility of $`_{}`$ follows from lemmas 14 and 15. The two required property on $`𝒟_{}`$ come directly from lemmas 16 and 17. If $`n`$, observe that $`_{}((\mathrm{𝑐𝑜𝑛𝑠}(n,\mathrm{𝑒𝑚𝑝𝑡𝑦}),0))=n`$. This concludes the proof. $`\mathrm{}`$
A *soft length space* is a length space on the resource monoid $`(,+,_{},𝒟_{})`$.
Given a soft length space $`A=(|A|,_A)`$, we can build the length space $`!A=(|A|,_{!A})`$, where $`\alpha ,e_{!A}a`$ iff $`\beta ,e_{!A}a`$ and $`\alpha _{}!\beta `$. We write $`[n,m]_{}`$ for $`(\mathrm{𝑐𝑜𝑛𝑠}(n,\mathrm{𝑒𝑚𝑝𝑡𝑦}),m)`$.
###### Lemma 19
For every $`\alpha `$ and for every $`n,m`$ the following inequality holds:
$$n.\alpha _{n_{}(\alpha )+m}!\alpha +[m,2n]_{}$$
* Let $`\alpha =(l,p)`$. We go by induction on $`l`$. If $`l`$ is $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, then
$`n.\alpha `$ $`=`$ $`(\mathrm{𝑒𝑚𝑝𝑡𝑦},p)`$
$`!\alpha +[m,2n]_{}`$ $`=`$ $`(\mathrm{𝑐𝑜𝑛𝑠}(m,\mathrm{𝑒𝑚𝑝𝑡𝑦}),\mathrm{max}\{p,2n\})`$
$`n_{}(\alpha )+m`$ $`=`$ $`m`$
$`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ $`_0`$ $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$
This implies the thesis. Moreover, if $`l=\mathrm{𝑐𝑜𝑛𝑠}(q,h)`$, then
$`n.\alpha `$ $`=`$ $`(n.l,p)=(\mathrm{𝑐𝑜𝑛𝑠}(nq,n.h),p)`$
$`!\alpha +[m,2n]_{}`$ $`=`$ $`(\mathrm{𝑐𝑜𝑛𝑠}(m,l),\mathrm{max}\{p,2n\})`$
$`n_{}(\alpha )+m`$ $`=`$ $`n(q+p_{}(l,p))+m`$
By induction hypothesis, we get
$`(n.h,p)`$ $`_{n_{}(h,p)+q}`$ $`!(h,p)+[q,2n]_{}=(l,\mathrm{max}\{p,2n\})`$
$`(n(q+p_{}(l,p))+m)+nq`$ $`=`$ $`m+2nq+np_{}(l,p)`$
$``$ $`m+\mathrm{max}\{p,2n\}(n_{}(h,p)+q)`$
from which the desired inequality easily follows. $`\mathrm{}`$
###### Lemma 20 (Basic Maps)
Given soft length spaces $`A,B`$ and a natural number $`n1`$, there are morphisms:
$`\mathrm{𝑐𝑜𝑛𝑡𝑟}_n`$ $`:`$ $`!A\stackrel{n\text{ times}}{\stackrel{}{A\mathrm{}A}}`$
$`\mathrm{𝑑𝑖𝑠𝑡𝑟}`$ $`:`$ $`!A!B!(AB)`$
where $`\mathrm{𝑐𝑜𝑛𝑡𝑟}(a)=(\stackrel{n\text{ times}}{\stackrel{}{a,\mathrm{},a}})`$ and $`\mathrm{𝑑𝑖𝑠𝑡𝑟}(a,b)=(a,b)`$
* We define realizers $`e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n`$ for every $`n1`$ by induction on $`n`$:
$`e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^1`$ $`=`$ $`e_{\mathrm{𝑖𝑑}}`$
$`e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^{n+1}`$ $`=`$ $`(e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n)^{}e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}`$
Clearly, $`e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n`$ is a realizer for $`\mathrm{𝑐𝑜𝑛𝑡𝑟}_n`$. Moreover, $`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n\}(x))n|x|+q_n`$, where $`q_n`$ does not depend on $`x`$. Now, let $`\psi _n`$ be such that $`_{}(\psi _n)\mathrm{𝑐𝑝}n`$ and $`\phi _{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n`$ be $`[q_n,2n]_{}+\psi _n`$ for every $`n1`$. Now, let $`\alpha ,j_{!A}a`$. This implies $`\alpha _{}!(l,m)`$, where $`(l,m),j_Aa`$. Notice that
$$n.(l,m)+\psi _n,\stackrel{n\text{ times}}{\stackrel{}{j,\mathrm{},j}}_{\underset{n\text{ times}}{\underset{}{A\mathrm{}A}}}(\stackrel{n\text{ times}}{\stackrel{}{a,\mathrm{},a}})$$
By lemma 19, we finally get
$`n.(l,m)+\psi _n`$ $`_{}`$ $`!(l,m)+[q_n,2n]_{}+\psi _n`$
$`=`$ $`!(l,m)+\phi _{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n\phi _{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n+\alpha `$
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n\}(j))`$ $``$ $`n|j|+q_n`$
$``$ $`n_{}(l,m)+q_n`$
$``$ $`𝒟_{}(n.(l,m),!(l,m)+[q_n,2n]_{})`$
$``$ $`𝒟_{}(n.(l,m),(\mathrm{𝑐𝑜𝑛𝑠}(q_n,l),\mathrm{max}\{m,2n\}))`$
$``$ $`𝒟_{}(n.(l,m)+\psi _n,(\mathrm{𝑐𝑜𝑛𝑠}(q_n,l),\mathrm{max}\{m,2n\})+\psi _n)`$
$``$ $`𝒟_{}(n.(l,m)+\psi _n,[q_n,2n]_{}+\alpha +\psi _n)`$
$``$ $`𝒟_{}(n.(l,m)+\psi _n,\alpha +\phi _{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n)`$
This proves each $`e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}^n`$ to be a morphism.
Let $`e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}=e_{\mathrm{𝑖𝑑}}`$. We know $`\{e_{\mathrm{𝑖𝑑}}\}(d)`$ takes constant time, say $`p`$. Then, let $`\psi ,\mu `$ be such that $`_{}(\psi )p+|e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}|`$, $`_{}(\mu )\mathrm{𝑐𝑝}`$. $`\phi _{\mathrm{𝑑𝑖𝑠𝑡𝑟}}`$ is then defined as $`\psi +!\mu `$. Now, let $`\alpha ,d,c_{!A!B}(a,b)`$. This implies $`\alpha !\beta +!\gamma `$, where $`\beta ,d_Aa`$ and $`\gamma ,c_Bb`$. This in turn implies $`\beta +\gamma +\mu ,d,c_{AB}(a,b)`$ and $`!(\beta +\gamma +\mu ),d,c_{!(AB)}(a,b)`$. Moreover
$$!(\beta +\gamma +\mu )=!\beta +!\gamma +!\mu _{}\alpha +!\mu _{}\alpha +\phi _{\mathrm{𝑑𝑖𝑠𝑡𝑟}}$$
Finally:
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}\}(d,c))`$ $``$ $`p_{}(\psi )`$
$``$ $`𝒟_{}(!(\beta +\gamma +\mu ),\alpha +!\mu )+_{}(\psi )`$
$``$ $`𝒟_{}(!(\beta +\gamma +\mu ),\alpha +!\mu +\psi )`$
$``$ $`𝒟_{}(!(\beta +\gamma +\mu ),\alpha +\phi _{\mathrm{𝑑𝑖𝑠𝑡𝑟}})`$
This proves $`\mathrm{𝑑𝑖𝑠𝑡𝑟}`$ to be a morphism. $`\mathrm{}`$
###### Lemma 21 (Functoriality)
If $`f:A\stackrel{e,\phi }{}B`$, then there is $`\psi `$ such that $`f:!A\stackrel{e,\psi }{}!B`$
* Let $`\theta `$ be $`!\phi `$ and suppose $`\alpha ,d_{!A}a`$. Then $`\alpha !\beta `$, where $`\beta ,d_Aa`$. Observe that there must be $`\gamma ,c`$ such that $`\gamma ,c_Bf(a)`$, $`\gamma _{}\beta +\phi `$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))𝒟_{}(\gamma ,\beta +\phi )`$. But then $`!\gamma ,c_{!B}f(a)`$ and, moreover
$`!\gamma `$ $`_{}`$ $`!(\beta +\phi )=!\beta +!\phi _{}!\beta +\theta `$
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))`$ $``$ $`𝒟_{}(\gamma ,\beta +\phi )𝒟_{}(!\gamma ,!(\beta +\phi ))`$
$``$ $`𝒟_{}(!\gamma ,!\beta +!\phi )𝒟_{}(!\gamma ,\alpha +\theta )`$
This implies $`f:!A\stackrel{e,\theta }{}!B`$. $`\mathrm{}`$
###### Proposition 2
For every $`n`$ there is a polynomial $`p_n:`$ such that $`_{}(\alpha )p_{0pt\alpha }(|\alpha |)`$ for every $`\alpha ||`$.
* We go by induction on $`n`$. First of all, we know that $`𝒟_{}((\mathrm{𝑒𝑚𝑝𝑡𝑦},0),(\mathrm{𝑒𝑚𝑝𝑡𝑦},m))=0`$, so $`p_0`$ is just the function which always returns $`0`$. $`p_{n+1}`$ is defined from $`p_n`$ as follows: $`p_{n+1}(x)=x+xp_n(x)`$. Indeed:
$`𝒟_{}((\mathrm{𝑒𝑚𝑝𝑡𝑦},0),(\mathrm{𝑐𝑜𝑛𝑠}(n,l),m))`$ $`=`$ $`n+m𝒟_{}((\mathrm{𝑒𝑚𝑝𝑡𝑦},0),(l,m))`$
$``$ $`|(\mathrm{𝑐𝑜𝑛𝑠}(n,l),m)|+|(\mathrm{𝑐𝑜𝑛𝑠}(n,l),m)|p_{0pt(l,m)}(|(\mathrm{𝑐𝑜𝑛𝑠}(n,l),m)|)`$
$`=`$ $`p_{0pt(\mathrm{𝑐𝑜𝑛𝑠}(n,l),m)}((\mathrm{𝑐𝑜𝑛𝑠}(n,l),m)).`$
This concludes the proof. $`\mathrm{}`$
Again, we do not claim that $`(n,m)p_n(m)`$ is a polynomial (c.f. Remark 2).
###### Theorem 2
Soft length spaces form a model of SAL.
Binary lists can be represented in SAL as cut-free proofs with conclusion
$$\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖲𝖠𝖫}}\alpha .!(\alpha \alpha )!(\alpha \alpha )(\alpha \alpha )$$
###### Corollary 2 (Soundness)
Let $`\pi `$ be an SAL proof with conclusion $`!^j\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖲𝖠𝖫}}!^k\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖲𝖠𝖫}}`$ and let $`f:LL`$ be the function induced by $`\pi ^{}`$. Then $`f`$ is computable in polynomial time.
## 6 Light Length Spaces
The grammar of formulae for Light Affine Logic is the one from Elementary Affine Logic, enriched with a new production $`A::=\mathrm{\S }A`$. Rules are reported in figure 4.
Light length spaces are a model of Light Affine Logic. The underlying resource monoid is more complex than the ones we encountered so far. This complexity is a consequence of the strange behaviour of modality $`!`$, which is functorial but does not distribute over tensor (i.e. $`!(AB)\cong ̸!A!B`$).
A *tree* is either $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ or a triple $`\mathrm{𝑛𝑜𝑑𝑒}(n,t,T)`$ where $`n`$, $`t`$ is itself a tree and $`T`$ is a finite nonempty set of trees. $`|𝒯|`$ is the set of all trees. We write $`[n]_𝒯`$ for the tree $`\mathrm{𝑛𝑜𝑑𝑒}(n,\mathrm{𝑒𝑚𝑝𝑡𝑦},\{\mathrm{𝑒𝑚𝑝𝑡𝑦}\})`$. The sum $`t+s`$ of two trees $`t`$ and $`s`$ is defined as follows, by induction on $`n`$:
$`\mathrm{𝑒𝑚𝑝𝑡𝑦}+t`$ $`=`$ $`t+\mathrm{𝑒𝑚𝑝𝑡𝑦}=t;`$
$`\mathrm{𝑛𝑜𝑑𝑒}(n,t,T)+\mathrm{𝑛𝑜𝑑𝑒}(m,u,U)`$ $`=`$ $`\mathrm{𝑛𝑜𝑑𝑒}(n+m,t+u,TU);`$
Here, more sophisticated techniques are needed. For every $`n,e`$, binary relations $`_e^n`$ on trees can be defined as follows
* $`t_e^0u`$ for every $`t,u|𝒯|`$;
* $`\mathrm{𝑒𝑚𝑝𝑡𝑦}_e^{n+1}t`$ for every $`t|𝒯|`$;
* $`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)_e^{n+1}\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ iff there is $`d`$ such that
+ $`med`$;
+ $`t_{d^2}^n\mathrm{𝑒𝑚𝑝𝑡𝑦}`$;
+ For every $`sT`$, $`s_d^n\mathrm{𝑒𝑚𝑝𝑡𝑦}`$.
* $`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)_e^{n+1}\mathrm{𝑛𝑜𝑑𝑒}(l,u,U)`$ iff there is $`d`$ such that
+ $`ml+ed`$;
+ There is a function $`f:\{1,\mathrm{},d\}U`$ such that $`t_{d^2}^nu+_1^df(i)`$;
+ For every $`sT`$ there is $`zU`$ with $`s_d^nz`$.
For every $`e,n`$ and for every trees $`t`$ and $`u`$ with $`t_e^nu`$, we define the natural number $`𝒟_e^n(t,u)`$ as follows:
$`𝒟_e^0(t,u)`$ $`=`$ $`0`$
$`𝒟_e^{n+1}(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})`$ $`=`$ $`e+𝒟_e^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})`$
$`𝒟_e^{n+1}(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑛𝑜𝑑𝑒}(m,t,T))`$ $`=`$ $`m+e+\underset{f}{\mathrm{max}}\{𝒟_{(m+e)^2}^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},t+{\displaystyle \underset{i=1}{\overset{m+e}{}}}f(i))\}`$
$`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),\mathrm{𝑒𝑚𝑝𝑡𝑦})`$ $`=`$ $`em+𝒟_{(em)^2}^n(t,\mathrm{𝑒𝑚𝑝𝑡𝑦})`$
$`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),\mathrm{𝑛𝑜𝑑𝑒}(l,u,U))`$ $`=`$ $`l+em+\underset{f}{\mathrm{max}}\{𝒟_{(l+em)^2}^n(t,u+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i))\}`$
If $`t`$ is a tree, then $`|t|`$ is the greatest integer appearing in $`t`$, i.e. $`|\mathrm{𝑒𝑚𝑝𝑡𝑦}|=0`$ and $`|\mathrm{𝑛𝑜𝑑𝑒}(n,t,T)|=\mathrm{max}\{n,|t|,\mathrm{max}_{uT}|u|\}`$.
The depth $`0ptt`$ of a tree $`t`$ is defined as follows: $`0pt\mathrm{𝑒𝑚𝑝𝑡𝑦}=0`$ and
$$0pt\mathrm{𝑛𝑜𝑑𝑒}(n,t,T)=1+\mathrm{max}\{0ptt,\underset{uT}{\mathrm{max}}0ptu\}.$$
Given a tree $`t|𝒯|`$, we define $`!t`$ as the tree $`\mathrm{𝑛𝑜𝑑𝑒}(1,\mathrm{𝑒𝑚𝑝𝑡𝑦},\{t\})`$ and $`\mathrm{\S }t`$ as the tree $`\mathrm{𝑛𝑜𝑑𝑒}(0,t,\{\mathrm{𝑒𝑚𝑝𝑡𝑦}\})`$. In this context, a notion of isomorphism between trees is needed: we say that trees $`t`$ and $`u`$ are *isomorphic* and we write $`tu`$ iff for every $`e,n`$ and for every tree $`v`$ the following hold:
$`v_e^nt`$ $``$ $`v_e^nu`$
$`t_e^nv`$ $``$ $`u_e^nv`$
$`𝒟_e^n(v,t)`$ $`=`$ $`𝒟_e^n(v,u)`$
$`𝒟_e^n(t,v)`$ $`=`$ $`𝒟_e^n(u,v)`$
###### Lemma 22
$`\mathrm{𝑒𝑚𝑝𝑡𝑦}[0]_𝒯`$. Moreover, for every tree $`t`$, $`t+\mathrm{𝑒𝑚𝑝𝑡𝑦}t+[0]_𝒯`$.
* We have to prove that for every $`e,n`$ and for every tree $`v`$:
$`v_e^n\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ $``$ $`v_e^n[0]_𝒯`$
$`\mathrm{𝑒𝑚𝑝𝑡𝑦}_e^nv`$ $``$ $`[0]_𝒯_e^nv`$
$`𝒟_e^n(v,\mathrm{𝑒𝑚𝑝𝑡𝑦})`$ $`=`$ $`𝒟_e^n(v,[0]_𝒯)`$
$`𝒟_e^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},v)`$ $`=`$ $`𝒟_e^n([0]_𝒯,v)`$
We go by induction on $`n`$, considering the case where $`n1`$, since the base case is trivial. First of all, observe that both $`\mathrm{𝑒𝑚𝑝𝑡𝑦}_e^{n+1}t`$ and $`[0]_𝒯_e^{n+1}t`$ for every $`t`$. Moreover, $`\mathrm{𝑒𝑚𝑝𝑡𝑦}_e^{n+1}\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ and $`[0]_𝒯_e^{n+1}\mathrm{𝑒𝑚𝑝𝑡𝑦}`$. Suppose now that $`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)_e^{n+1}\mathrm{𝑒𝑚𝑝𝑡𝑦}`$. This means that there is $`d`$ such that
+ $`med`$;
+ $`t_{d^2}^n\mathrm{𝑒𝑚𝑝𝑡𝑦}`$;
+ for every $`sT`$, $`s_d^n\mathrm{𝑒𝑚𝑝𝑡𝑦}`$.
If we put $`f(i)=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ for every $`i`$, we get $`t_{d^2}^n\mathrm{𝑒𝑚𝑝𝑡𝑦}+_{i=1}^df(i)`$, which yields $`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)_e^{n+1}[0]_𝒯`$. In the same way, we can prove that if $`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)_e^{n+1}[0]_𝒯`$, then $`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)_e^{n+1}\mathrm{𝑒𝑚𝑝𝑡𝑦}`$.
We have:
$`𝒟_e^{n+1}(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})`$ $`=`$ $`e+𝒟_{e^2}^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})`$
$`𝒟_e^{n+1}(\mathrm{𝑒𝑚𝑝𝑡𝑦},[0]_𝒯)`$ $`=`$ $`e+𝒟_{e^2}^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})`$
$`𝒟_e^{n+1}([0]_𝒯,\mathrm{𝑒𝑚𝑝𝑡𝑦})`$ $`=`$ $`e+𝒟_{e^2}^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})`$
$`𝒟_e^{n+1}(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑛𝑜𝑑𝑒}(m,t,T))`$ $`=`$ $`m+e+\underset{f}{\mathrm{max}}\{𝒟_{(m+e)^2}^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},t+{\displaystyle \underset{i=1}{\overset{m+e}{}}}f(i))\}`$
$`=`$ $`𝒟_e^{n+1}([0]_𝒯,\mathrm{𝑛𝑜𝑑𝑒}(m,t,T))`$
$`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),\mathrm{𝑒𝑚𝑝𝑡𝑦})`$ $`=`$ $`em+𝒟_{(em)^2}^n(t,\mathrm{𝑒𝑚𝑝𝑡𝑦})`$
$`=`$ $`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),[0]_𝒯)`$
Moreover, observe that
$`\mathrm{𝑒𝑚𝑝𝑡𝑦}+\mathrm{𝑒𝑚𝑝𝑡𝑦}=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ $``$ $`[0]_𝒯=[0]_𝒯+\mathrm{𝑒𝑚𝑝𝑡𝑦}`$
$`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)+\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ $`=`$ $`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)+[0]_𝒯`$
This concludes the proof. $`\mathrm{}`$
###### Proposition 3 (Compatibility)
For every $`n,e`$, $`\mathrm{𝑒𝑚𝑝𝑡𝑦}_e^nt`$ for every $`t`$ and, moreover, if $`t_e^nu`$ then $`t+v_e^nu+v`$ for every $`t,u,v`$.
* $`\mathrm{𝑒𝑚𝑝𝑡𝑦}_e^nt`$ is trivial. The second statement can be proved by induction on $`n`$. The base case is trivial. In the inductive case, we can suppose all the involved trees to be different from $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$. Suppose that $`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)_e^{n+1}\mathrm{𝑛𝑜𝑑𝑒}(l,u,U)`$. We should prove $`\mathrm{𝑛𝑜𝑑𝑒}(m+k,t+v,TV)_e^{n+1}\mathrm{𝑛𝑜𝑑𝑒}(l+k,u+v,UV)`$. However,
$`m+k`$ $``$ $`(l+e)d+k=(l+k+e)d`$
$`t+v`$ $`_{d^2}^n`$ $`u+{\displaystyle \underset{i=1}{\overset{d}{}}}f(i)+v=u+v+{\displaystyle \underset{i=1}{\overset{d}{}}}f(i)`$
Moreover, for every $`zTV`$ there certanily exists $`wUV`$ such that $`z_d^nw`$. $`\mathrm{}`$
###### Proposition 4 (Transitivity)
If $`t_e^nu_d^nv`$, then $`t_{d+e}^nv`$.
* We go by induction on $`n`$. We can directly go to the inductive case, since if $`n=0`$, then the thesis is trivial. We can assume all the involved trees to be different from $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$. Let us suppose $`\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)_e^{n+1}\mathrm{𝑛𝑜𝑑𝑒}(l,u,U)`$ and $`\mathrm{𝑛𝑜𝑑𝑒}(l,u,U)_d^{n+1}\mathrm{𝑛𝑜𝑑𝑒}(k,v,V)`$ First of all, we have $`ml+ec`$ and $`lk+db`$, which yields $`mk+db+ec=k+(d+e)(b+c)`$. Moreover, by hypothesis, there are functions $`f:\{1,\mathrm{},c\}U`$ and $`g:\{1,\mathrm{},b\}V`$ such that
$`t`$ $`_{c^2}^n`$ $`u+{\displaystyle \underset{i=1}{\overset{c}{}}}f(i)`$
$`u`$ $`_{b^2}^n`$ $`v+{\displaystyle \underset{i=1}{\overset{b}{}}}g(i)`$
Therefore, by inductive hypothesis and by proposition 3:
$`t`$ $`_{c^2+b^2}^n`$ $`v+{\displaystyle \underset{i=1}{\overset{c}{}}}f(i)+{\displaystyle \underset{i=1}{\overset{b}{}}}g(i)`$
$`_{bc}^n`$ $`v+{\displaystyle \underset{i=1}{\overset{c}{}}}h(i)+{\displaystyle \underset{i=1}{\overset{b}{}}}g(i)`$
where $`h:\{1,\mathrm{},c\}V`$. We can then find a function $`k:\{1,\mathrm{},c+b\}V`$ such that
$$t_{(c+b)^2}^nv+\underset{i=1}{\overset{c+b}{}}k(i).$$
Finally, if $`zT`$ then we find $`wU`$ such that $`z_c^nw`$. We then find $`xV`$ such that $`w_b^nx`$ and so $`z_{c+b}^nx`$. $`\mathrm{}`$
###### Proposition 5
For every $`n,e`$ and for every $`t,u,v`$, $`𝒟_e^n(t,u)𝒟_e^n(t+v,u+v)`$
* We can proceed by induction on $`n`$ and, again, the case $`n=0`$ is trivial. In the inductive case, as usual, we can suppose all the involved trees to be different from $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$. We have
$`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),\mathrm{𝑛𝑜𝑑𝑒}(l,u,U))`$
$`=`$ $`l+em+\underset{f}{\mathrm{max}}\{𝒟_{(l+em)^2}^n(t,u+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i))\}`$
$`=`$ $`l+em+𝒟_{(l+em)^2}^n(t,u+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i))`$
where $`f`$ and realizes the max. By induction hypothesis,
$`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),\mathrm{𝑛𝑜𝑑𝑒}(l,u,U))`$
$``$ $`(l+k)+e(m+k)+𝒟_{((l+k)+e(m+k))^2}^n(t+v,u+v+{\displaystyle \underset{i=1}{\overset{(l+k)+e(m+k)}{}}}f(i))`$
$``$ $`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)+\mathrm{𝑛𝑜𝑑𝑒}(k,v,V),\mathrm{𝑛𝑜𝑑𝑒}(l,u,U)+\mathrm{𝑛𝑜𝑑𝑒}(k,v,V))`$
This concludes the proof. $`\mathrm{}`$
###### Proposition 6
$`𝒟_e^n(t,u)+𝒟_d^n(u,v)𝒟_{e+d}^n(t,v)`$
* We can proceed by induction on $`n`$ and, again, the case $`n=0`$ is trivial. In the inductive case, as usual, we can suppose all the involved trees to be different from $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$. Now
$`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),\mathrm{𝑛𝑜𝑑𝑒}(l,u,U))+𝒟_d^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(l,u,U),\mathrm{𝑛𝑜𝑑𝑒}(k,v,V))`$
$`=`$ $`l+em+\underset{f}{\mathrm{max}}\{𝒟_{(l+em)^2}^n(t,u+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i))\}`$
$`+k+dl+\underset{g}{\mathrm{max}}\{𝒟_{(k+dl)^2}^n(u,v+{\displaystyle \underset{i=1}{\overset{k+dl}{}}}g(i))\}`$
$`=`$ $`k+(e+d)m+𝒟_{(l+em)^2}^n(t,u+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i))`$
$`+𝒟_{(k+dl)^2}^n(u,v+{\displaystyle \underset{i=1}{\overset{k+dl}{}}}g(i))`$
$`=`$ $`k+(e+d)m+𝒟_{(l+em)^2}^n(t,u+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i))`$
$`+𝒟_{(k+dl)^2}^n(u+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i),v+{\displaystyle \underset{i=1}{\overset{k+dl}{}}}g(i)+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i))`$
$``$ $`k+(e+d)m+𝒟_{(l+em)^2+(k+dl)^2}^n(t,v+{\displaystyle \underset{i=1}{\overset{k+dl}{}}}g(i)+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i))`$
A function $`h:\{1,\mathrm{},l+em\}V`$ such that $`_{i=1}^{l+em}f(i)_{(l+em)(k+dl)}^n_{i=1}^{l+em}h(i)`$ can be easily defined, once we remember that $`\mathrm{𝑛𝑜𝑑𝑒}(l,u,U)_d^n\mathrm{𝑛𝑜𝑑𝑒}(k,v,V)`$. This yields
$`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),\mathrm{𝑛𝑜𝑑𝑒}(l,u,U))+𝒟_d^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(l,u,U),\mathrm{𝑛𝑜𝑑𝑒}(k,v,V))`$
$``$ $`k+(e+d)m+𝒟_{(l+em)^2+(k+dl)^2}^n(t,v+{\displaystyle \underset{i=1}{\overset{k+dl}{}}}g(i)+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i))`$
$`+𝒟_{(l+em)(k+dl)}^n(v+{\displaystyle \underset{i=1}{\overset{k+dl}{}}}g(i)+{\displaystyle \underset{i=1}{\overset{l+em}{}}}f(i),v+{\displaystyle \underset{i=1}{\overset{k+dl}{}}}g(i)+{\displaystyle \underset{i=1}{\overset{l+em}{}}}h(i))`$
$``$ $`k+(e+d)m+𝒟_{(k+(e+d)m)^2}^n(t,v+{\displaystyle \underset{i=1}{\overset{k+dl}{}}}g(i)+{\displaystyle \underset{i=1}{\overset{l+em}{}}}h(i))`$
$``$ $`k+(e+d)m+𝒟_{(k+(e+d)m)^2}^n(t,v+{\displaystyle \underset{i=1}{\overset{l+(d+e)m}{}}}p(i))`$
where $`p:\{1,\mathrm{},l+(d+e)m\}V`$, $`p(i)=f(i)`$ if $`il+em`$ and $`p(i)=g(i(l+em))`$ otherwise. But, then
$`𝒟_e^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),\mathrm{𝑛𝑜𝑑𝑒}(l,u,U))+𝒟_d^{n+1}(\mathrm{𝑛𝑜𝑑𝑒}(l,u,U),\mathrm{𝑛𝑜𝑑𝑒}(k,v,V))`$
$``$ $`𝒟_{e+d}^n(\mathrm{𝑛𝑜𝑑𝑒}(m,t,T),\mathrm{𝑛𝑜𝑑𝑒}(k,v,V))`$
This concludes the proof. $`\mathrm{}`$
###### Lemma 23
For every $`t,u,e`$, if $`t_e^{\mathrm{max}\{0ptt,0ptu\}}u`$, then for every $`n>\mathrm{max}\{0ptt,0ptu\}`$, $`t_e^nu`$ and $`𝒟_e^n(t,u)=𝒟_e^{\mathrm{max}\{0ptt,0ptu\}}(t,u)`$.
* A straightforward induction on $`\mathrm{max}\{0ptt,0ptu\}`$. $`\mathrm{}`$
The binary relation $`_𝒯`$ on $`|𝒯|`$ is defined by putting $`t_𝒯u`$ whenever $`0ptt0ptu`$ and $`t_0^{0ptu}u`$. $`𝒟_𝒯`$ is defined by letting $`𝒟_𝒯(t,u)=𝒟_0^{0ptu}(t,u)`$.
###### Lemma 24
$`𝒯=(|𝒯|,+,_𝒯,𝒟_𝒯)`$ is a resource monoid.
* $`(|𝒯|,+)`$ is certainly a commutative monoid. For every $`t`$, $`t_𝒯t`$, as can be proved by induction on $`t`$: $`\mathrm{𝑒𝑚𝑝𝑡𝑦}_0^0\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ by definition and, moreover, $`t=\mathrm{𝑛𝑜𝑑𝑒}(m,u,U)_0^{0ptt}t`$ because, by inductive hypothesis, $`u_0^{0ptu}u`$ which yields, by lemma 23, $`u_0^{0ptt1}u`$. In the same way, we can prove that, for every $`vU`$, $`v_0^{0ptt1}v`$. Now, suppose $`t_𝒯u`$ and $`u_𝒯v`$. This means that $`t_0^{0ptu}u`$, $`u_0^{0ptv}v`$, $`0ptt0ptu`$ and $`0ptu0ptv`$. We can then conclude that $`0ptt0ptv`$, that $`t_0^{0ptv}u`$ (by lemma 23) and $`t_0^{0ptv}v`$ (by proposition 6). This in turn yields $`t_𝒯v`$. Let us now prove compatibility: suppose $`t_𝒯u`$ and let $`v`$ be a tree. Then $`0ptt0ptu`$ and $`t_0^{0ptu}u`$. If $`0ptv0ptu`$, then $`0ptu+v=0ptu`$ and we can proceed by getting $`t+v_0^{0ptu+v}u+v`$ (by proposition 3), which means $`t+v_𝒯u+v`$. If, on the other hand, $`0ptv>0ptu`$, then we can first apply lemma 23 obtaining $`t_0^{0ptu+v}u`$ and then $`t+v_0^{0ptu+v}u+v`$ (by proposition 3). By way of lemma 23 and propositions 6 and 5 we get
$`𝒟_𝒯(t,u)+𝒟_𝒯(u,v)`$ $`=`$ $`𝒟_0^{0ptu}(t,u)+𝒟_0^{0ptv}(u,v)`$
$`=`$ $`𝒟_0^{0ptv}(t,u)+𝒟_0^{0ptv}(u,v)`$
$``$ $`𝒟_0^{0ptv}(t,v)=𝒟_𝒯(t,v)`$
$`𝒟_𝒯(t,u)`$ $`=`$ $`𝒟_0^{0ptu}(t,u)𝒟_0^{0ptu+v}(t,u)`$
$``$ $`𝒟_0^{0ptu+v}(t+v,u+v)=𝒟_𝒯(t+v,u+v)`$
This concludes the proof. $`\mathrm{}`$
A *light length space* is a length space on the resource monoid $`𝒯=(|𝒯|,+,_𝒯,𝒟_𝒯)`$. Given a light length space $`A=(|A|,_A)`$, we can define:
* The light length space $`!A=(|A|,_{!A})`$ where $`t,e_{!A}a`$ iff $`u,e_Aa`$ and $`t_𝒯!u`$.
* The light length space $`\mathrm{\S }A=(|A|,_{\mathrm{\S }A})`$ where $`t,e_{\mathrm{\S }A}a`$ iff $`u,e_Aa`$ and $`t_𝒯\mathrm{\S }u`$.
The following results states the existence of certain morphisms and will be useful when interpreting light affine logic.
###### Lemma 25 (Basic Maps)
Given light length spaces $`A,B`$, there are morphisms: $`\mathrm{𝑐𝑜𝑛𝑡𝑟}:!A!A!A`$, $`\mathrm{𝑑𝑖𝑠𝑡𝑟}:\mathrm{\S }A\mathrm{\S }B\mathrm{\S }(AB)`$ and $`\mathrm{𝑑𝑒𝑟𝑒𝑙𝑖𝑐𝑡}:!A\mathrm{\S }A`$ where $`\mathrm{𝑐𝑜𝑛𝑡𝑟}(a)=(a,a)`$ and $`\mathrm{𝑑𝑖𝑠𝑡𝑟}(a,b)=(a,b)`$ and $`\mathrm{𝑑𝑒𝑟𝑒𝑙𝑖𝑐𝑡}(a)=a`$.
* We know that $`\{e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}\}(d)`$ takes time at most $`|d|+p`$, where $`p`$ is a constant. Then, let $`t,u|𝒯|`$ be such that $`_𝒯(t)p+|e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}|`$, $`_𝒯(u)\mathrm{𝑐𝑝}`$. Define $`t_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}`$ to be $`t+u+[2]_𝒯`$. Clearly, $`_𝒯(t_{\mathrm{𝑐𝑜𝑛𝑡𝑟}})|e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}|`$. Now, let $`v,d_{!A}a`$. This means that $`v_𝒯!w`$ where $`w,d_Aa`$. Then:
$`u+!w+!w`$ $`_𝒯`$ $`!w+!w`$
$`_𝒯(u+!w+!w)`$ $``$ $`_𝒯(u)+_𝒯(!w)+_𝒯(!w)`$
$``$ $`\mathrm{𝑐𝑝}+_𝒯(!w)+_𝒯(!w)|d,d|`$
This implies $`u+!w+!w,|d,d|_{!A!A}(a,a)`$. Moreover, $`u+!w+!w=u+!w+[1]_𝒯_𝒯v+t_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}`$. Finally,
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑐𝑜𝑛𝑡𝑟}}\}(d))`$ $``$ $`|d|+p_𝒯(w)+_{trees}(t)`$
$``$ $`𝒟_𝒯(u+!w+!w,!w+t_{\mathrm{𝑐𝑜𝑛𝑡𝑟}})𝒟_𝒯(u+!w+!w,v+t_{\mathrm{𝑐𝑜𝑛𝑡𝑟}})`$
This proves $`\mathrm{𝑐𝑜𝑛𝑡𝑟}`$ to be a morphism.
Let $`e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}=e_{\mathrm{𝑖𝑑}}`$. We know that $`\{e_{\mathrm{𝑖𝑑}}\}(d)`$ takes constant time, say at most $`p`$. Then, let $`t,u|𝒯|`$ be such that $`_𝒯(t)p+|e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}|`$, $`_𝒯(u)\mathrm{𝑐𝑝}`$. $`t_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}`$ is then defined as $`t+\mathrm{\S }u`$. Now, let $`v,d,c_{\mathrm{\S }A\mathrm{\S }B}(a,b)`$. This implies that $`v\mathrm{\S }w+\mathrm{\S }x`$, where $`w,d_Aa`$ and $`x,c_Bb`$. This in turn means that $`w+x+u,d,c_{AB}(a,b)`$ and $`\mathrm{\S }(w+x+u),d,c_{AB}(a,b)`$. Moreover
$$\mathrm{\S }(w+x+u)=\mathrm{\S }w+\mathrm{\S }x+\mathrm{\S }uv+t_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}$$
Finally:
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}\}(d,c))`$ $``$ $`p_𝒯(t)`$
$``$ $`𝒟_𝒯(0,t)+𝒟_𝒯(\mathrm{\S }(w+x+u),v+\mathrm{\S }u)𝒟_𝒯(\mathrm{\S }(w+x+u),v+t_{\mathrm{𝑑𝑖𝑠𝑡𝑟}})`$
This proves $`\mathrm{𝑑𝑖𝑠𝑡𝑟}`$ to be a morphism.
Let $`e_{\mathrm{𝑑𝑒𝑟𝑒𝑙𝑖𝑐𝑡}}=e_{\mathrm{𝑖𝑑}}`$. We know that $`\{e_{\mathrm{𝑑𝑒𝑟𝑒𝑙𝑖𝑐𝑡}}\}(d)`$ takes constant time, say at most $`p`$. Then, let $`t_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}|𝒯|`$ be such that $`_𝒯(t_{\mathrm{𝑑𝑖𝑠𝑡𝑟}})p+|e_{\mathrm{𝑑𝑒𝑟𝑙𝑖𝑐𝑡}}|`$. Now, let $`v,d_{!A}a`$. This means that $`v!w`$, where $`w,d_Aa`$. This in turn means that $`\mathrm{\S }w,d_{\mathrm{\S }A}a`$. Moreover
$$\mathrm{\S }w!w!w+t_{\mathrm{𝑑𝑒𝑟𝑒𝑙𝑖𝑐𝑡}}.$$
Finally:
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e_{\mathrm{𝑑𝑖𝑠𝑡𝑟}}\}(d))`$ $``$ $`p_𝒯(t_{\mathrm{𝑑𝑒𝑟𝑒𝑙𝑖𝑐𝑡}})`$
$``$ $`𝒟_𝒯(0,t_{\mathrm{𝑑𝑒𝑟𝑒𝑙𝑖𝑐𝑡}})+𝒟_𝒯(\mathrm{\S }w,!w)`$
$``$ $`𝒟_𝒯(\mathrm{\S }w,!w+t_{\mathrm{𝑑𝑒𝑟𝑒𝑙𝑖𝑐𝑡}})`$
This proves $`\mathrm{𝑑𝑒𝑟𝑒𝑙𝑖𝑐𝑡}`$ to be a morphism. $`\mathrm{}`$
###### Lemma 26
For every $`t|𝒯|`$, there is $`u`$ such that, for every $`v`$, $`!(v+t)_𝒯!v+u`$.
* First of all we will prove the following statement by induction on $`t`$: for every $`t`$, there is an integer $`\overline{t}`$ such that for every $`u`$, $`u+t_{\overline{t}}^{\mathrm{max}\{0ptu,0ptt\}}u`$. If $`t=\mathrm{𝑒𝑚𝑝𝑡𝑦}`$, we can choose $`\overline{t}`$ to be just $`0`$, since $`u_0^nu`$ for every $`u`$. If $`t=\mathrm{𝑛𝑜𝑑𝑒}(m,v,V)`$, then we put $`\overline{t}=m+\overline{v}+_{wV}\overline{w}`$. Let $`u`$ be an arbitrary tree and let us assume, without losing generality, that $`u=\mathrm{𝑛𝑜𝑑𝑒}(l,w,W)`$. Let $`d=\overline{v}+_{wV}\overline{w}`$. We get
$`l+m`$ $``$ $`l+m+(\overline{v}+{\displaystyle \underset{wV}{}}\overline{w})(\overline{v}+{\displaystyle \underset{wV}{}}\overline{w})`$
$`=`$ $`l+\overline{t}d`$
$`v+w`$ $`_{\overline{v}}^{\mathrm{max}\{0ptv,0ptw\}}`$ $`w`$
$`_0^{\mathrm{max}\{0ptv,0ptw\}}`$ $`w+{\displaystyle \underset{i=1}{\overset{d}{}}}\mathrm{𝑒𝑚𝑝𝑡𝑦}`$
$`xV.x`$ $`_{\overline{x}}^{0ptx}`$ $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$
$`xW.x`$ $`_0^{0ptx}`$ $`x`$
Using known results, we can rewrite these inequalities as follows
$`l+m`$ $``$ $`l+\overline{t}d`$
$`v+w`$ $`_{d^2}^{\mathrm{max}\{0ptt,0ptu\}1}`$ $`w+{\displaystyle \underset{i=1}{\overset{d}{}}}\mathrm{𝑒𝑚𝑝𝑡𝑦}`$
$`xV.x`$ $`_d^{\mathrm{max}\{0ptt,0ptu\}1}`$ $`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$
$`xW.x`$ $`_d^{\mathrm{max}\{0ptt,0ptu\}1}`$ $`x`$
This yields $`u+t_{\overline{t}}^{\mathrm{max}\{0ptu,0ptt\}}t`$.
Let us now go back to the lemma we are proving. We will now prove that for every $`t`$, any term $`u=\mathrm{𝑛𝑜𝑑𝑒}(\overline{t},w,U)`$ such that $`0ptu=0ptt+1`$ satisfies the thesis. Indeed, if we put $`d=\overline{t}`$ and $`n=0ptv+t`$, we get:
$`1`$ $``$ $`\overline{t}d+1`$
$`\mathrm{𝑒𝑚𝑝𝑡𝑦}`$ $`_{d^2}^n`$ $`u`$
$`v+t`$ $`_d^n`$ $`v`$
This, in turn implies $`!(v+t)_0^{n+1}!v+u`$, which yields $`!(v+t)_𝒯!v+u`$. $`\mathrm{}`$
###### Lemma 27 (Functoriality)
If $`f:A\stackrel{e,\phi }{}B`$, then there are $`\psi ,\theta `$ such that $`f:!A\stackrel{e,\psi }{}!B`$ and $`f:\mathrm{\S }A\stackrel{e,\theta }{}\mathrm{\S }B`$.
* Let $`\xi `$ be the tree obtained from $`\phi `$ by lemma 26 and put $`\psi =\xi +\phi +[1]_𝒯`$. Suppose that $`t,d_{!A}a`$. Then $`t!u`$, where $`u,d_Aa`$. Observe that there must be $`v,c`$ such that $`v,c_Bf(a)`$, $`v_𝒯u+\phi `$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))_𝒯(u+\phi )𝒟_𝒯(v,u+\phi )`$. But then $`!v,c_{!B}f(a)`$ and moreover
$`!v`$ $`_𝒯`$ $`!(u+\phi )_𝒯!u+\xi _𝒯t+\psi `$
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))`$ $``$ $`𝒟_𝒯(v,u+\phi )𝒟_𝒯(!v,!(u+\phi )+[1]_𝒯)`$
$``$ $`𝒟_𝒯(!v,!u+\xi +[1]_𝒯)𝒟_𝒯(!v,t+\psi )`$
This means that $`f:!A\stackrel{e,\psi }{}!B`$. Now, let $`\theta `$ be $`\mathrm{\S }\phi `$ and suppose $`t,d_{\mathrm{\S }A}a`$. Then $`t\mathrm{\S }u`$, where $`u,d_Aa`$. Observe that there must be $`v,c`$ such that $`v,c_Bf(a)`$, $`v_𝒯u+\phi `$ and $`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))_𝒯(u+\phi )𝒟_𝒯(v,u+\phi )`$. But then $`\mathrm{\S }v,c_{\mathrm{\S }B}f(a)`$ and, moreover
$`\mathrm{\S }v`$ $`_𝒯`$ $`\mathrm{\S }(u+\phi )=\mathrm{\S }u+\mathrm{\S }\phi _𝒯t+\theta `$
$`\mathrm{𝑇𝑖𝑚𝑒}(\{e\}(d))`$ $``$ $`𝒟_𝒯(v,u+\phi )𝒟_𝒯(\mathrm{\S }v,\mathrm{\S }(u+\phi ))`$
$``$ $`𝒟_𝒯(\mathrm{\S }v,\mathrm{\S }u+\mathrm{\S }\phi ))𝒟_𝒯(\mathrm{\S }v,t+\theta )`$
This means that $`f:\mathrm{\S }A\stackrel{e,\theta }{}\mathrm{\S }B`$. $`\mathrm{}`$
Now, we can prove a polynomial bound on $`_T(t)`$:
###### Proposition 7
For every $`n`$ there is a polynomial $`p_n:`$ such that $`_𝒯(t)p_{0ptt}(|t|)`$.
* We prove a stronger statement by induction on $`n`$: for every $`n`$ there is a polynomial $`q_n:^2`$ such that for every $`t,e`$, $`𝒟_e^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},t)q_n(|t|,e)`$. First of all, we know that $`𝒟_e^0(\mathrm{𝑒𝑚𝑝𝑡𝑦},t)=0`$, so $`q_0`$ is just the function which always returns $`0`$. $`q_{n+1}`$ is defined from $`q_n`$ as follows: $`q_{n+1}(x,y)=x+y+q_n(x(x+y+1),(x+y)^2)`$. Indeed:
$`𝒟_e^{n+1}(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})`$ $`=`$ $`e+𝒟_e^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑒𝑚𝑝𝑡𝑦})`$
$``$ $`e+q_n(0,e)e+|\mathrm{𝑒𝑚𝑝𝑡𝑦}|`$
$`+q_n(|\mathrm{𝑒𝑚𝑝𝑡𝑦}|(|\mathrm{𝑒𝑚𝑝𝑡𝑦}|+e+1),(|\mathrm{𝑒𝑚𝑝𝑡𝑦}|+e)^2)`$
$`=`$ $`q_{n+1}(|\mathrm{𝑒𝑚𝑝𝑡𝑦}|,e)`$
$`𝒟_e^{n+1}(\mathrm{𝑒𝑚𝑝𝑡𝑦},\mathrm{𝑛𝑜𝑑𝑒}(m,t,T))`$ $`=`$ $`m+e+\underset{f}{\mathrm{max}}\{𝒟_{(m+e)^2}^n(\mathrm{𝑒𝑚𝑝𝑡𝑦},t+{\displaystyle \underset{i=1}{\overset{m+e}{}}}f(i))\}`$
$``$ $`m+e+q_n((m+e+1)(|\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)|),(m+e)^2)`$
$``$ $`|\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)|+e`$
$`+q_n((|\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)|+e+1)(|\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)|),(|\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)|+e)^2)`$
$``$ $`q_{n+1}(|\mathrm{𝑛𝑜𝑑𝑒}(m,t,T)|,e)`$
At this point, however, it suffices to put $`p_n(x)=q_n(x,0)`$. $`\mathrm{}`$
As for EALand SAL, we cannot claim $`(n,m)p_n(m)`$ to be a polynomial. However, this is not a problem since we will be able to majorize binary strings by trees with bounded depth (cf.Remark 2).
### 6.1 Interpreting Light Affine Logic
As for the $`!`$ modality, $`\mathrm{\S }A_\eta ^{}=\mathrm{\S }A_\eta ^{}`$.
###### Theorem 3
Light length spaces form a model of LAL.
Binary lists can be represented in LAL as cut-free proofs with conclusion
$$\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖫𝖠𝖫}}\alpha .!(\alpha \alpha )!(\alpha \alpha )\mathrm{\S }(\alpha \alpha )$$
###### Corollary 3 (Soundness)
Let $`\pi `$ be an LAL proof with conclusion $`\{!,\mathrm{\S }\}^j\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖫𝖠𝖫}}\{!,\mathrm{\S }\}^k\mathrm{𝐿𝑖𝑠𝑡}_{\mathrm{𝖫𝖠𝖫}}`$ and let $`f:BB`$ be the function induced by $`\pi ^{}`$. Then $`f`$ is computable in polynomial time.
## 7 Interpreting LFPL
In one of us had introduced another language, LFPL, with the property that all definable functions on natural numbers are polynomial time computable. The key difference between LFPL and other systems is that a function defined by iteration or recursion is not marked as such using modalities or similar and can therefore be used as a step function of subsequent recursive definitions.
In this section we will describe a resource monoid $``$ for LFPL, which will provide a proof of polytime soundness for that system. This is essentially the same as the proof from , but more structured and, hopefully, easier to understand.
The new approach also yields some new results, namely the justification of second-order quantification, a !-modality, and a new type of binary trees based on cartesian product which allows alternative but not simultaneous access to subtrees.
### 7.1 Overview of LFPL
LFPL is intuitionistic, affine linear logic, i.e., a linear functional language with $`,,+,\times `$. Unlike in the original presentation we also add polymorphic quantification here. In addition, LFPL has basic types for inductive datatypes, for example unary and binary natural numbers, lists, and trees. There is one more basic type, namely $`\mathrm{}`$, the resource type.
The recursive constructors for the inductive datatypes each take an additional argument of type $`\mathrm{}`$ which prevents one to invoke more constructor functions than one. Dually to the constructors one has iteration principles which make the $`\mathrm{}`$-resource available in the branches of a recursive definition. For example, the type $`T(X)`$ of $`X`$-labelled binary trees has constructors $`\mathrm{𝐥𝐞𝐚𝐟}:T(X)`$ and $`\mathrm{𝐧𝐨𝐝𝐞}:\mathrm{}XT(X)T(X)T(X)`$. The iteration principle allows one to define a function $`T(X)A`$ from closed terms $`A`$ and $`\mathrm{}XAAA`$.
In this paper we “internalise” the assumption of closedness using a $`!`$-modality.
Using this iteration principle one can encode recursive definitions by ML-style pattern matching provided recursive calls are made on structurally smaller arguments only.
Here is a fragment of an LFPL program for “treesort” written in functional notation: the additional arguments of type $`\mathrm{}`$ are supplied using @. Note that the insert function takes an extra argument of type $`\mathrm{}`$.
```
let insert x t d = match t with
Leaf -> Node(x,Leaf,Leaf)@d
| Node(y,l,r)@d’ ->
if x<=y then Node(y,insert x l d,r)@d’
else Node(y,l,insert x r d)@d’
let extract t = match t with
Leaf -> nil
| Node(x,l,r)@d ->
append (extract l) (cons(x,extract r)@d)
```
### 7.2 A Resource Monoid for LFPL
The underlying set of $``$ is the set of pairs $`(l,p)`$ where $`l`$ is a natural number and $`p`$ is a monotone polynomial in a single variable $`x`$. The addition is defined by $`(l_1,p_1)+(l_2,p_2)=(l_1+l_2,p_1+p_2)`$, accordingly, the neutral element is $`0=(0,0)`$. We have a submonoid $`_0=\{(l,p)l=0\}`$.
To define the ordering we set $`(l_1,p_1)(l_2,p_2)`$ iff $`l_1l_2`$ and $`(p_2p_1)(x)`$ is monotone and nonnegative for all $`xl_2`$. For example, we have $`(1,42x)(42,x^2)`$, but $`(1,42x)(41,x^2)`$. The distance function is defined by
$$𝒟_{}((l_1,p_1),(l_2,p_2))=(p_2p_1)(l_2)$$
We can pad elements of $``$ by adding a constant to the polynomial. The following is now obvious.
###### Lemma 28
Both $``$ and $`_0`$ are resource monoids.
A simple inspection of the proofs in Section 3.1 shows that the realisers for all maps can be chosen from $`_0`$. This is actually the case for an arbitrary submonoid of a resource monoid. We note that realisers of elements may nevertheless be drawn from all of $``$. We are thus led to the following definition.
###### Definition 1
An LFPL-space is a length space over the resource monoid $``$. A morphism from LFPL length space $`A`$ to $`B`$ is a morphism between length spaces which admits a majorizer from $`_0`$.
###### Proposition 8
LFPL length spaces with their maps form a symmetric monoidal closed category.
###### Definition 2
Let $`A`$ be an LFPL space and $`n`$. The LFPL space $`A^n`$ is defined by $`|A^n|=|A|`$ and $`\alpha ,e_{A^n}a`$ iff $`\alpha (2n1).\beta `$ for some $`\beta `$ such that $`\beta ,e_Aa`$.
So, $`A^n`$ corresponds to the subset of $`A\mathrm{}A`$ consisting of those tuples with all $`n`$ components equal to each other. The factor $`2n1`$ (“modified difference”) instead of just $`n`$ is needed in order to justify the linear time needed to compute the copying involved in the obvious morphism from $`A^{m+n}`$ to $`A^mA^n`$.
Let $`I`$ be an index set and $`A_i,B_i`$ be $`I`$-indexed families of LFPL spaces. A uniform map from $`(A_i)_i`$ to $`(B_i)_i`$ consists of a family of maps $`f_i:A_iB_i`$ such that there exist $`e,\alpha `$ with the property that $`\alpha ,ef_i`$ for all $`i`$. Recall that, in particular, the denotations of proofs with free type variables are uniform maps.
###### Proposition 9
For each $`A`$ there is a uniform (in $`m,n`$) map $`A^{m+n}A^mA^n`$. Moreover, $`A^1`$ is isomorphic to $`A`$.
The LFPL-space $`\mathrm{}`$ is defined by $`|\mathrm{}|=\{\mathrm{}\}`$ and put $`\alpha ,d_{\mathrm{}}\mathrm{}`$ if $`\alpha (1,0)`$.
For each LFPL-space $`A`$ we define LFPL-space $`!A`$ by $`|!A|=|A|`$ and $`\alpha ,t_{!A}a`$ if there exists $`\alpha ^{}=(0,p)_0`$ with $`\alpha ^{},t_Aa`$ and $`\alpha (0,(x+1)p)`$.
###### Proposition 10
There is an LFPL space $`\mathrm{}`$ and for each LFPL space $`A`$ there is an LFPL space $`!A`$ with the following properties:
* $`|!A|=|A|`$.
* If $`f:AB`$ then $`f:!A!B`$.
* $`!(AB)!A!B`$
* The obvious functions $`!A\mathrm{}^nA^n\mathrm{}^n`$ are a uniform map.
The last property means intuitively that with $`n`$ “diamonds” we can extract $`n`$ copies from an element of type $`!A`$ and get the $`n`$ “diamonds” back for later use.
* We have $`(0+1)p(0)=p(0)|t|`$. Compatibility with $``$ is obvious.
For functoriality assume that $`\varphi ,ef`$ where $`\varphi =(0,q)_0`$. We claim that $`(0,(x+1)q),ef`$ *qua* morphism from $`!A`$ to $`!B`$. Suppose that $`\alpha ,t_{!A}a`$ where $`\alpha (0,(x+1)p)`$ and $`(0,p),t_Aa`$. Since $`f`$ is a morphism, we obtain $`v,\beta `$ such that $`\beta ,v_Bf(a)`$ and $`\beta \varphi +(0,p)`$. This implies that $`\beta _0`$ as well, say, $`\beta =(0,r)`$ where $`rp+q`$. We also know that $`r(0)|v|`$ by the definition of length spaces. Now $`(0,(x+1)r),v_{!B}f(b)`$. On the other hand $`(x+1)r(x+1)(p+q)`$. The resource bounds are obvious.
Finally, consider the required morphism $`!A\mathrm{}^nA^n\mathrm{}^n`$. Clearly, it may be realised by the identity; we claim that $`0`$ can serve as a majoriser. Indeed, a majoriser of $`(a,d)|!A\mathrm{}^n|`$ is of the form $`(2n1,(x+1)p)`$ where $`(0,p)`$ majorises $`a`$ in $`A`$. Now, $`(2n1,(2n1)p)`$ is a majoriser of $`(a,d)`$ in $`A^n\mathrm{}^n`$. But $`((x+1)(2n1)p`$ is monotone and nonnegative above $`2n1`$. $`\mathrm{}`$
#### Remark
We remark at this point that we obtain an alternative resource monoid $`_S`$ for SAL whose underlying set and ordering are as in $``$, but whose addition is given by addition as $`(l_1,p_1)+(l_2,p_2)=(\mathrm{max}(l_1,l_2),p_1+p_2)`$. Length spaces over $`_S`$ with maps majorised by $`_S`$ (not $`_0`$) then also form a sound model of SAL. This points to a close relationship between LFPL and SAL and also shows a certain tradeoff between the two systems. The slightly more complex model is needed for LFPL since in LFPL the C-rule of SAL is so to say internalised in the form of the uniform map $`!A\mathrm{}^nA^n\mathrm{}^n`$. Notice that SAL’s map $`!AA^n`$ cannot be uniform. This uniformity of LFPL allows for an internal implementation of datatypes and recursion as we now show.
###### Definition 3
Let $`T_i`$ be a family of LFPL spaces such that $`|T_i|=T`$ independent of $`i`$. The LFPL space $`i.T_i`$ is defined by $`|i.T_i|=|T|`$ and $`\alpha ,e_{i.T_i}t`$ if $`\alpha ,e_{T_i}t`$ for some $`i`$.
Note that if we have a uniform family of maps $`T_iU`$ where $`U`$ does not depend on $`i`$ then we obtain a map $`i.T_iU`$ (existential elimination).
Conversely, if we have a uniform family of maps $`U_iV_{f(i)}`$ then we get a uniform family of maps $`U_ij.V_j`$ (existential introduction). We will use an informal “internal language” to denote uniform maps which when formalised would amount to an extension of LFPL with indexed type dependency in the style of Dependent ML .
### 7.3 Inductive Datatypes
In order to interpret unary natural numbers, we define $`N=n.N_n`$ where
$$N_n=\mathrm{}^nA.(AA)^nAA$$
We can internally define a successor map $`\mathrm{}N_nN_{n+1}`$ as follows: starting from $`d:\mathrm{},\stackrel{}{d}:\mathrm{}^n`$ and $`f:(AA)^nAA`$ we obtain a member of $`\mathrm{}^{n+1}`$ (from $`d`$ and $`\stackrel{}{d}`$) and we define $`f^{}:(AA)^{n+1}AA`$ as $`\lambda (u^{AA},\stackrel{}{u}^{(AA)^n}).\lambda z^A.u(f\stackrel{}{u}z)`$. From this, we obtain a map $`\mathrm{}NN`$ by existential introduction and elimination.
Of course, we also have a constant zero $`IN_0`$ yielding a map $`IN`$ by existential introduction.
Finally, we can define an iteration map
$$!(\mathrm{}AA)N_nAA$$
as follows: Given $`t:!(\mathrm{}AA)`$ and $`(\stackrel{}{d},f)N_n`$ we unpack $`t`$ using Proposition 10 to yield $`t^{}((\mathrm{}A)A)^n`$ as well as $`\stackrel{}{d}\mathrm{}^n`$. Feeding these “diamonds” one by one to the components of $`t^{}`$ we obtain $`t^{\prime \prime }(AA)^n`$. But then $`ft^{\prime \prime }`$ yields the required element of $`AA`$.
Existential elimination now yields a single map
$$!(\mathrm{}AA)NAA$$
Similarly, we can interpret binary $`X`$-labelled trees using a type family
$$T_n=\mathrm{}^n(XAAA)^nA^{n+1}A$$
and defining trees proper as $`n.T_n`$. We get maps $`\mathrm{𝐥𝐞𝐚𝐟}:T_0`$ and $`\mathrm{𝐧𝐨𝐝𝐞}:\mathrm{}XT_{n_1}T_{n_2}T_{n_1+n_2+1}`$ and an analogous iteration construct.
Finally, and this goes beyond what was already known, we can define “lazy trees” using cartesian product (also known as additive conjunction).
First, we recall from ordinary affine linear logic that an additive conjunction can be defined as
$$A\times B=C.(CA)(CB)C$$
The first projection map $`A\times BA`$ is given internally by $`\lambda (f^{CA},g^{CB},c^C).fc`$. Analogously, we have a second projection. Given maps $`f:CA`$ and $`g:CB`$ we obtain a map $`f,g:CA\times B`$ internally as $`\lambda c^C.(f,g,c)`$.
Now, following the pattern of the binary trees $`T_{m,n}`$ above, we define another family
$$T_d^\times =\mathrm{}^dA.(X(A\times A)A)^dAA$$
and $`T^\times =d.T_d^\times `$. We get maps $`\mathrm{𝐥𝐞𝐚𝐟}:\mathrm{}T_0^\times `$ and $`\mathrm{𝐧𝐨𝐝𝐞}:\mathrm{}X(T_{d_1}\times T_{d_2})T_{1+\mathrm{max}(d_1,d_2)}`$ as well as an analogous iteration construct.
We describe in detail the construction of the “node” map which is not entirely straightforward. First, we note that for any length spaces $`A,B`$ and $`m,n`$ the obvious map $`(\mathrm{}^mA)\times (\mathrm{}^nB)\mathrm{}^{\mathrm{max}(m,n)}(A\times B)`$ is a morphism. This is because a majoriser of an element of $`(\mathrm{}^mA)\times (\mathrm{}^nB)`$ must be of the form $`(k,p)`$ where $`k\mathrm{max}(m,n)`$ in view of the existence of the projection maps.
Now suppose we are given (internally) $`d:\mathrm{},x:X,\mathrm{𝑙𝑟}:T_{d_1}^\times \times T_{d_2}^\times `$. Using the just described morphism we decompose $`\mathrm{𝑙𝑟}`$ into $`\stackrel{}{d}:\mathrm{}^{\mathrm{max}(d_1,d_2)}`$ and $`\mathrm{𝑙𝑟}^{}:W_{d_1}\times W_{d_2}`$ where $`W_i=(X(A\times A)A)^iAA`$. We have stripped off the universal quantifier.
Now $`d`$ and $`\stackrel{}{d}`$ together yield an element of $`\mathrm{}^{1+\mathrm{max}(d_1,d_2)}`$. It remains to construct a member of $`W_{1+\mathrm{max}(d_1,d_2)}`$. To this end, we assume $`u:X(A\times A)A`$ and $`f:(X(A\times A)A)^{\mathrm{max}(d_1,d_2)}`$ and define the required element of $`A`$ as $`ux\mathrm{𝑙𝑟}^{}.1fa,\mathrm{𝑙𝑟}^{}.2fa`$. Here $`.1`$ and $`.2`$ denote the projections from the cartesian product. The sharing of the variables $`f`$, $`a`$, $`\mathrm{𝑙𝑟}^{}`$ is legal in the two components of a cartesian pairing, but would of course not be acceptable in a $``$ pairing. We have elided the obvious coercions from $`(\mathrm{\_})^{\mathrm{max}(d_1,d_2)}`$ to $`(\mathrm{\_})^{d_i}`$.
We remark that these cartesian trees are governed by their depth rather than their number of nodes. We also note that if $`X=I`$ we can form the function $`\lambda d^{\mathrm{}}.\lambda t^{T^\times }.\mathrm{𝐧𝐨𝐝𝐞}d()t,r:\mathrm{}T^\times T^\times `$. Iterating this map yields a function $`NT^\times `$ computing full binary trees of a given depth. Of course, on the level of the realisers, such a tree is not laid out in full as this would require exponential space, but computed lazily as subtrees are being accessed. Exploring the implications of this for programming is left to future work.
## 8 Conclusion
We have given a unified semantic framework with which to establish soundness of various systems for capturing complexity classes by logic and programming. Most notably, our framework has all of second-order multiplicative linear logic built in, so that only the connectives and modalities going beyond this need to be verified explicitly.
While resulting in a considerable simplification of previous soundness proofs, in particular for LFPL and LAL, our method has also lead to new results, in particular polymorphism and a modality for LFPL.
The method proceeds by assiging both abstract resource bounds in the form of elements from a resource monoid and resource-bounded computations to proofs (respectively, programs). In this way, our method can be seen as a combination of traditional Kleene-style realisability (which only assigns computations) and polynomial and quasi interpretation known from term rewriting (which only assigns resource bounds). An altogether new aspect is the introduction of more general notions of resource bounds than just numbers or polynomials as formalised in the concept of resource monoid. We thus believe that our methods can also be used to generalise polynomial interpretations to (linear) higher-order.
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# Strange-quark contribution to the ratio of neutral- to charged-current cross sections in neutrino-nucleus scattering
## I Introduction
In a desperate effort to salvage the conservation of energy and angular momentum, Wolfgang Pauli postulated the existence of the neutrino: a neutral and (almost) massless particle that he feared will never be detected. Seventy five years later and with three neutrino families firmly established, neutrino physics in the 21st century has taken center stage in fields as diverse as cosmology, astro, nuclear, and particle physics. With the commission of new neutrino observatories and facilities for the study of solar, atmospheric, and reactor neutrinos, conclusive evidence now exists in favor of neutrino oscillations. Although neutrino-oscillation experiments have now evolved from the discovery into the precision phase, important questions remain unanswered Kayser (2003). Depending on these answers, a radical modification to the standard model of particle physics may be required. At the very least, neutrino oscillations already demand a mild extension to the standard model: in the standard model the individual lepton numbers must be conserved.
The discovery of neutrino oscillations established two incontrovertible facts: a) that neutrinos have mass and that these masses can not all be equal and b) that the three known neutrinos $`\nu _e`$, $`\nu _\mu `$, and $`\nu _\tau `$ are linear combinations of three neutrino mass eigenstates (commonly referred as $`\nu _1`$, $`\nu _2`$, and $`\nu _3`$). Note that due to the very small neutrino masses, oscillation experiments are sensitive to only the squared mass difference ($`\mathrm{\Delta }m^2`$) of the mass eigenstates. Further, the unitary matrix linking the flavor to the mass eigenstates has a well-known counterpart in the quark sector: the CKM matrix. While precision experiments have started to determine squared mass differences and mixing angles, a strong, vibrant, and interdisciplinary program is now being established to tackle the myriad of remaining open questions Kayser (2003). Among these are: What is the mass hierarchy? Is the neutrino a Dirac or Majorana particle? Does the neutrino matrix contain a CP-violating phase that may explain our matter-dominated Universe? Is there a need for additional “sterile” neutrinos?
The apparent need for an additional sterile neutrino stems from two older experiments. The first one in 1989 measured the total and partial (into hadrons and charged leptons) widths of the $`Z^0`$ boson at the large electron-positron (LEP) collider at CERN and extracted the number of neutrinos flavors to be $`N_\nu =3.00\pm 0.08`$; note that a more recent analysis reports $`N_\nu =2.984\pm 0.008`$ Eidelman et al. (2004). The second experiment was a 1995 neutrino-oscillation experiment at the Liquid Scintillator Neutrino Detector (LSND) at the Los Alamos National Laboratory Athanassopoulos et al. (1995). This experiment reported evidence of $`\overline{\nu }_\mu \overline{\nu }_e`$ oscillation, but with a squared mass difference that is too large — and thus inconsistent — with the two independent values extracted from solar, atmospheric, and reactor experiments. Simply put, for three neutrino flavors the algebraic relation $`\mathrm{\Delta }m_{21}^2+\mathrm{\Delta }m_{32}^2=\mathrm{\Delta }m_{31}^2`$ must be satisfied, yet $`\mathrm{\Delta }m_{\mathrm{sol}}^2+\mathrm{\Delta }m_{\mathrm{atm}}^2\mathrm{\Delta }m_{\mathrm{LSND}}^2`$. The most favorable scenario that accommodates three independent $`\mathrm{\Delta }m^2`$ values is the addition of a sterile neutrino. Yet other possibilities exist: the LSND analysis may be incorrect. It is the primary goal of the Booster Neutrino Experiment (BooNE and its first phase MiniBooNE) at Fermilab to confirm the LSND result.
While MiniBooNE’s primary goal is to confirm the LSND result, this unique facility is also ideal for the study of supernova neutrinos, neutrino-nucleus scattering, and hadronic structure. An ambitious experimental program — the Fine-grained Intense Neutrino Scattering Scintillator Experiment (FINeSSE) — aims to measure the strange-quark contribution to the spin of the nucleon Potterveld et al. (2002); Bugel et al. (2003); Brice et al. (2004). FINeSSE is part of a larger program started in the late 80’s that attempts to answer a fundamental nucleon-structure question: how do the non-valence (“sea”) quarks — particularly the strange quarks — contribute to the observed properties of the nucleon? To date, the role of strange quarks in the nucleon remains a contentious issue and one that remains a subject of intense activity all over the world. In an attempt to find a satisfactory answer to this fundamental question, a number of reactions have been proposed. These include: (i) deep inelastic scattering of neutrinos on protons Bazarko et al. (1995); Goncharov et al. (2001), (ii) deep inelastic scattering of polarized charged leptons Adams et al. (1997), (iii) pseudoscalar meson scattering on a proton Dover and Fishbane (1990), and (iv) parity-violating electron scattering McKeown (1989); Beck (1989). Part of the controversy arises because these reactions do not all reach similar conclusions. For example, both reactions (i) and (ii) suggest a non-zero strangeness contribution, in contrast to reaction (iv) that indicates a strange quark contribution to the charge and magnetic moment consistent with zero. Parity-violating electron scattering in particular has received extensive experimental attention as in the SAMPLE Collaboration at the MIT-Bates accelerator Mueller et al. (1997), the HAPPEX Collaboration at the Jefferson Laboratory Aniol et al. (1999), and the A4 Collaboration at the MAMI facility in Mainz Maas (2003). However, theoretical investigations have shown that large radiative corrections and nuclear-structure effects impact negatively on the extraction of strange-quark matrix elements Musolf and Holstein (1990); Musolf et al. (1994); Horowitz and Piekarewicz (1993).
Neutrino-induced reactions provide a viable alternative to parity-violating electron scattering. While the latter is primarily sensitive to the strange electric and magnetic form factors of the nucleon, the former is particularly sensitive to the axial-vector form factor of the proton — through the combination $`(\mathrm{\Delta }u\mathrm{\Delta }d\mathrm{\Delta }s+\mathrm{\Delta }c\mathrm{\Delta }b+\mathrm{\Delta }t)`$. This sensitivity is the result of the small weak-vector charge of the proton ($`14\mathrm{sin}^2\theta _\mathrm{W}0.08`$) and the suppression of the weak anomalous magnetic moment at small $`Q^2`$. In the above expression the heavy quark flavors ($`c`$, $`b`$, and $`t`$) can be eliminated using a well-defined renormalization group procedure Bass et al. (2002, 2003). Further, the isovector combination $`(\overline{u}\gamma _\mu \gamma _5u\overline{d}\gamma _\mu \gamma _5d)`$ is constrained from neutron beta decay. This leaves the (assumed isoscalar) strange-quark contribution to the spin of the proton $`\mathrm{\Delta }s`$ to be determined from the elastic neutrino-proton reaction. Yet an absolute cross-section measurement of this reaction is an experimental challenge due to difficulties in the determination of the absolute neutrino flux. An attractive alternative has been proposed by Garvey and collaborators Garvey et al. (1992, 1993) in which the extraction of $`\mathrm{\Delta }s`$ proceeds through a measurement of the ratio of proton-to-neutron cross sections in neutral-current (NC) neutrino-nucleon scattering. This ratio is defined by the following expression:
$$R(p/n)=\frac{\sigma (\nu p\nu p)}{\sigma (\nu n\nu n)}.$$
(1)
This ratio is very sensitive to the strange-quark contribution to the spin of the nucleon as $`\mathrm{\Delta }s`$ \[or $`g_A^sG_A^{(s)}(Q^2=0)`$\] interferes with the isovector contribution ($`G_A^{(3)}`$) with one sign in the numerator and with the opposite sign in the denominator \[see Eq. 31\]. Unfortunately, $`R(p/n)`$ is difficult to measure with the desired accuracy due to experimental difficulties associated with neutron detection Brice et al. (2004). It is for this reason that FINeSSE will focus initially on the neutral- to charged-current ratio (NC/CC):
$$R(NC/CC)=\frac{\sigma (\nu p\nu p)}{\sigma (\nu n\mu ^{}p)}.$$
(2)
This ratio is “simply” determined from counting the number of events with an outgoing proton and missing mass relative to those events with an outgoing proton and a muon. Note that the CC reaction, being purely isovector, is insensitive to $`\mathrm{\Delta }s`$. As such, $`R(NC/CC)`$ is about a factor of two less sensitive to $`\mathrm{\Delta }s`$ than $`R(p/n)`$.
From a theoretical perspective extracting $`\mathrm{\Delta }s`$ from the ratio of cross sections is also attractive. As a large number of the scattering events at FINeSSE will be from nucleons bound to a Carbon nucleus, it is important to understand nuclear-structure corrections. This issue came to light in experiment E374 at the Brookhaven National Laboratory (BNL) where it was found that 80% of the events involved neutrino scattering off carbon atoms, while only 20% were from free protons. As nuclear-structure corrections in $`R(NC/CC)`$ appear to be insensitive to final-state interactions between the outgoing proton and the residual nucleus Meucci et al. (2004); Martinez et al. (2005), the ratio $`R(NC/CC)`$ may be accurately computed using a much simpler plane-wave formalism. Indeed, in Sec. II we will show how the cross section ratio $`R(NC/CC)`$ in Carbon computed in a plane-wave formalism may be expressed in a form that closely resembles the “Feynman-trace” approach used to calculate the cross section from free nucleons. We note in closing that the new generation of neutrino experiments will require a thorough understanding of neutrino-nucleus interactions since the detectors often contain complex nuclei. Experimental and theoretical work related to neutrino scattering from light and heavier nuclear targets may be found in Refs. Barish et al. (1977); Athanassopoulos et al. (1997); Nakamura et al. (2001); Auerbach et al. (2002); Mintz and Pourkaviani (1989); Kosmas and Oset (1996); Volpe et al. (2002); Maieron et al. (2003).
Our paper has been organized as follows. In Sec. II we briefly review the formalism developed in Ref. van der Ventel and Piekarewicz (2004) for the neutral-current case and point out the main modifications required to make it applicable to charged-current neutrino-nucleus scattering. Our main results — with a focus on the sensitivity of $`R(NC/CC)`$ to $`\mathrm{\Delta }s`$ — are presented in Sec. III. Finally, we summarize the main points of this work in Sec. IV.
## II Formalism
In this section the formalism for the description of charged-current neutrino-nucleus scattering is presented. As the basic outline follows closely the neutral-current formalism developed in Ref. van der Ventel and Piekarewicz (2004), we present a brief review that focusses on those modifications that arise from a finite muon mass.
### II.1 Cross section in terms of the leptonic and hadronic tensors
The lowest-order Feynman diagram for the knockout of a bound proton via the charged-current reaction $`[\nu +X(Z,A)\mu ^{}+p+X(Z1,A1)]`$ is shown in Fig. 1. Here the initial four-momentum of the (left-handed) neutrino is $`k`$ while the four-momentum and helicity of the outgoing muon are $`k^{}`$ and $`h^{}`$, respectively. The reaction proceeds via the exchange of a virtual $`W^+`$ boson with four-momentum $`q=(\omega ,𝐪)`$. The kinematical variables defining the hadronic arm are the four-momentum of the target ($`P`$) and residual nucleus ($`P^{}`$). Finally, $`p^{}`$ and $`s^{}`$ denote the four-momentum and spin component of the ejectile proton. Energy-momentum conservation demands that:
$$q=kk^{}=p^{}+P^{}P.$$
(3)
The dynamical information for this reaction is contained in the transition matrix element given by
$`i`$ $`=`$ $`\left\{\overline{\mu }(𝐤^{},h^{})\left[{\displaystyle \frac{ig}{2\sqrt{2}}}\left(\gamma ^\mu \gamma ^\mu \gamma ^5\right)\right]\nu (𝐤)\right\}i𝒟_{\mu \nu }(q)`$ (4)
$`\left\{p^{},s^{};\mathrm{\Psi }_f(P^{})\left|{\displaystyle \frac{ig}{2\sqrt{2}}}\mathrm{cos}\theta _C\widehat{J}^\nu (q)\right|\mathrm{\Psi }_i(P)\right\}.`$
In Eq. (4) the initial and final nuclear states are denoted by $`\mathrm{\Psi }_i(P)`$ and $`\mathrm{\Psi }_f(P^{})`$, respectively. Furthermore, $`g`$ is the weak coupling constant, $`\theta _C`$ is the Cabbibo angle ($`\mathrm{cos}\theta _C=0.974`$), and $`\widehat{J}^\nu (q)`$ is the weak nuclear current operator. As only low-momentum transfers ($`q_\mu q^\mu Q^2M_W^2`$) will be considered, the following approximation is valid
$$𝒟_{\mu \nu }(q)=\frac{g_{\mu \nu }+q_\mu q_\nu /M_W^2}{q^2M_W^2}\frac{g_{\mu \nu }}{M_W^2}.$$
(5)
Using this expression the transition matrix element may be written as
$$=\frac{G_F}{\sqrt{2}}\mathrm{cos}\theta _C\left[\overline{\mu }(𝐤^{},h^{})\gamma _\mu (1\gamma ^5)\nu (𝐤)\right]\left[p^{},s^{};\mathrm{\Psi }_f(P^{})\left|\widehat{J}^\mu (q)\right|\mathrm{\Psi }_i(P)\right],$$
(6)
where the Fermi coupling constant $`G_F`$ has been introduced via
$$\frac{G_F}{\sqrt{2}}=\frac{g^2}{8M_W^2}.$$
(7)
In Eq. (6) $`\mu (𝐤^{},h^{})`$ is the Dirac spinor for the outgoing muon expressed in the helicity representation. That is, (suppressing “prime” indices for clarity),
$$\mu (𝐤,h)=\sqrt{\frac{E_k+m}{2E_k}}\left(\begin{array}{c}\varphi _h(\widehat{𝐤})\\ h\frac{k}{E_k+m}\varphi _h(\widehat{𝐤})\end{array}\right),\left(E_k=\sqrt{k^2+m^2};k=|𝐤|\right),$$
(8)
where $`\varphi _{h=\pm 1}(\widehat{𝐤})`$ are two-component Pauli spinors given by
$$\varphi _{h=+1}(\widehat{𝐤})=\left(\begin{array}{c}\mathrm{cos}(\theta /2)\\ \mathrm{sin}(\theta /2)e^{i\varphi }\end{array}\right),\varphi _{h=1}(\widehat{𝐤})=\left(\begin{array}{c}\mathrm{sin}(\theta /2)e^{i\varphi }\\ \mathrm{cos}(\theta /2)\end{array}\right).$$
(9)
Here $`m`$ denotes the muon mass, and $`\theta `$ and $`\varphi `$ are the polar and azimuthal angles of the muon momentum. The neutrino spinor $`\nu (𝐤)`$ is directly obtained from the above expressions by setting the fermion mass to $`m=0`$ and the helicity to $`h=1`$. Note that left/right projection operators on plus/minus helicity states do not vanish in general due to the finite muon mass. That is,
$$𝒫_{L/R}\mu (𝐤,h=\pm 1)\frac{1}{2}(1\gamma ^5)\mu (𝐤,h=\pm 1)=𝒪(m/k)0.$$
(10)
Further, we have adopted a non-covariant normalization for the Dirac spinors of Eq. (8),
$$\mu ^{}(𝐤,h)\mu (𝐤,h^{})=\overline{\mu }(𝐤,h)\gamma ^0\mu (𝐤,h^{})=\delta _{hh^{}},$$
(11)
a choice that is in accordance with the standard normalization of the bound-state spinor Serot and Walecka (1986) and that is given by
$$𝒰_\alpha ^{}(𝐫)𝒰_\alpha (𝐫)d^3𝐫=1.$$
(12)
Following Ref. van der Ventel and Piekarewicz (2004) the differential cross section can now be written as
$$d\sigma =\frac{G_F^2\mathrm{cos}^2\theta _C}{2(2\pi )^5}d^{\mathrm{\hspace{0.17em}3}}𝐤^{}d^{\mathrm{\hspace{0.17em}3}}𝐩^{}\delta (E_k+M_AE_k^{}E_p^{}E_P^{})\mathrm{}_{\mu \nu }W^{\mu \nu },$$
(13)
where the leptonic tensor is given by
$$\mathrm{}_{\mu \nu }=\mathrm{Tr}\left[\left(\gamma _\mu \gamma _\mu \gamma _5\right)\left(\nu (𝐤)\overline{\nu }(𝐤)\right)\left(\gamma _\nu \gamma _\nu \gamma _5\right)\left(\mu (𝐤^{},h^{})\overline{\mu }(𝐤^{},h^{})\right)\right],$$
(14)
and a discussion of the hadronic tensor $`W^{\mu \nu }`$ is postponed until the next section.
We conclude this section with the evaluation of the leptonic tensor. To do so, both matrices $`\nu \overline{\nu }`$ and $`\mu \overline{\mu }`$ in Eq. (14) are first expressed in terms of Dirac matrices. For the case of the massive muon we obtain
$$\mu (𝐤^{},h^{})\overline{\mu }(𝐤^{},h^{})=\frac{\left(\text{/}k^{}+m\right)}{2E_k^{}}\left[\frac{1}{2}(1+h^{}\gamma ^5\text{/}s)\right],$$
(15)
where the four-component spin vector is given by
$$s^\mu s^\mu (𝐤^{})=\frac{1}{m}(k^{},E_k^{}\widehat{𝐤}^{}).$$
(16)
The corresponding expression for the massless left-handed neutrino may be obtained from the above equations by setting the helicity to $`h=1`$ and by taking the massless ($`m0`$) limit. Note that in the massless limit $`ms^\mu k^\mu `$. Thus we obtain,
$$\nu (𝐤)\overline{\nu }(𝐤)=\frac{\text{/}k}{2E_k}\left[\frac{1}{2}(1+\gamma ^5)\right]$$
(17)
Finally, by substituting the above expressions into the leptonic tensor of Eq. (14), which in turn we separate into ($`\mu \nu `$) symmetric and antisymmetric parts,
$$\mathrm{}^{\mu \nu }\mathrm{}_S^{\mu \nu }+\mathrm{}_A^{\mu \nu },$$
(18)
we obtain
$$\mathrm{}_S^{\mu \nu }=\frac{2}{kE_k^{}}\left(k^\mu K^\nu +K^\mu k^\nu g^{\mu \nu }kK^{}\right),$$
(19a)
$$\mathrm{}_A^{\mu \nu }=\frac{2i}{kE_k^{}}\epsilon ^{\mu \nu \alpha \beta }k_\alpha K_\beta ^{}.$$
(19b)
Note that in the above expressions the following four-vector has been introduced:
$$K^{}\frac{1}{2}(k^{}h^{}ms)\underset{m=0}{}k^{}\delta _{h^{},1},$$
(20)
where the last expression denotes the massless limit. Hence, in the $`m0`$ limit the leptonic tensor vanishes for positive-helicity ($`h^{}=+1`$) but for negative-helicity ($`h^{}=1`$) goes over to Eq. (17) of Ref. van der Ventel and Piekarewicz (2004). Finally, note that the following convention was adopted Yndurain (1999):
$$\mathrm{Tr}\left(\gamma ^5\gamma ^\mu \gamma ^\nu \gamma ^\alpha \gamma ^\beta \right)=4i\epsilon ^{\mu \nu \alpha \beta },(\epsilon ^{0123}=1,\epsilon _{0123}=+1).$$
(21)
We close this section with a comment on the conservation (or rather lack-thereof) of the leptonic tensor. While the antisymmetric component satisfies
$$q_\mu \mathrm{}_A^{\mu \nu }=\mathrm{}_A^{\mu \nu }q_\nu =0,$$
(22)
due to the antisymmetric property of the Levi-Civita tensor, this is no longer true for the symmetric part due to the finite muon mass. That is,
$$q_\mu \mathrm{}_S^{\mu \nu }0\mathrm{and}\mathrm{}_S^{\mu \nu }q_\nu 0.$$
(23)
### II.2 Differential cross section in terms of nuclear structure functions
In Eq. (13) of the previous section it was shown that the differential cross section for the CC reaction may be written as a contraction of the leptonic tensor with the hadronic tensor, where the latter is defined in terms of the expectation value of the weak nuclear operator \[see Eq. (6)\]. That is,
$$W^{\mu \nu }=\left[p^{},s^{};\mathrm{\Psi }_f(P^{})\left|\widehat{J}^\mu (q)\right|\mathrm{\Psi }_i(P)\right]\left[p^{},s^{};\mathrm{\Psi }_f(P^{})\left|\widehat{J}^\nu (q)\right|\mathrm{\Psi }_i(P)\right]^{}W_S^{\mu \nu }+W_A^{\mu \nu }.$$
(24)
Although the general form of the hadronic tensor was introduced and discussed in detail in Ref. van der Ventel and Piekarewicz (2004), some of its most salient features are underscored here for completeness. For the case of unpolarized proton emission, the hadronic tensor may be written in terms of thirteen independent structure functions,
$`W_S^{\mu \nu }`$ $`=`$ $`W_1g^{\mu \nu }+W_2q^\mu q^\nu +W_3P^\mu P^\nu +W_4p^\mu p^\nu `$ (25a)
$`+`$ $`W_5(q^\mu P^\nu +P^\mu q^\nu )+W_6(q^\mu p^\nu +p^\mu q^\nu )+W_7(P^\mu p^\nu +p^\mu P^\nu ),`$
$`W_A^{\mu \nu }`$ $`=`$ $`W_8(q^\mu P^\nu P^\mu q^\nu )+W_9(q^\mu p^\nu p^\mu q^\nu )+W_{10}(P^\mu p^\nu p^\mu P^\nu )`$ (25b)
$`+`$ $`W_{11}\epsilon ^{\mu \nu \alpha \beta }q_\alpha P_\beta +W_{12}\epsilon ^{\mu \nu \alpha \beta }q_\alpha p_\beta ^{}+W_{13}\epsilon ^{\mu \nu \alpha \beta }P_\alpha p_\beta ^{}.`$
Note that all structure functions are functions of the four Lorentz-invariant quantities, $`q^\mu q_\mu Q^2`$, $`qP`$, $`qp^{}`$, and $`Pp^{}`$. Details on the contraction between the leptonic and hadronic tensors
$`\mathrm{}_{\mu \nu }W^{\mu \nu }`$ $`=`$ $`\mathrm{}_{\mu \nu }^SW_S^{\mu \nu }+\mathrm{}_{\mu \nu }^AW_A^{\mu \nu },`$ (26)
have been reserved to Appendix A. Yet we note that the charged-current reaction is now sensitive to the three structure functions $`W_2`$, $`W_5`$ and $`W_6`$. This is in contrast to the neutral-current reaction (see Eq. (22a) of Ref. van der Ventel and Piekarewicz (2004)); the origin of this difference is the non-conservation of the symmetric part of the leptonic tensor due to the finite muon mass \[see Eq. (23)\]. However, as the antisymmetric part of the leptonic tensor is manifestly conserved, both NC and CC processes are insensitive to the $`W_8`$ and $`W_9`$ structure functions.
This concludes the model-independent description of charged-current neutrino-nucleus scattering. In summary, the cross section may be parametrized in terms of eleven nuclear-structure functions. In principle, they could be determined by a “super” Rosenbluth separation. In practice, however, this is not possible so we resort to a relativistic mean-field model to obtain explicit expressions for these quantities. This will be done in the next section.
### II.3 Model-dependent evaluation of the cross section
In the previous section a model-independent formalism was presented for charged-current neutrino-nucleus scattering. Specifically, the cross section was written in terms of a set of nuclear structure functions that parametrize our “ignorance” about the strong-interactions physics at the hadronic vertex. However, to proceed any further a number of approximations need to be made in order to obtain a numerically tractable problem.
The first “no-recoil” approximation, detailed in Eqs. (23)–(26) of Ref. van der Ventel and Piekarewicz (2004), is purely kinematical and yields the following expression for the angle-integrated differential cross section:
$$\frac{d\sigma (h^{})}{dE_p^{}}=\frac{G_F^{\mathrm{\hspace{0.17em}2}}\mathrm{cos}^2\theta _C}{2(2\pi )^4}k^{}E_k^{}p^{}E_p^{}_0^\pi \mathrm{sin}\alpha d\alpha _0^\pi \mathrm{sin}\theta d\theta _0^{2\pi }𝑑\varphi \left(\mathrm{}_{\mu \nu }W^{\mu \nu }\right).$$
(27)
Here $`\alpha `$ is the polar angle defining the direction of the outgoing proton having momentum $`p^{}|𝐩^{}|`$ and energy $`E_p^{}=\sqrt{p^2+M^2}`$. Similarly, $`\theta `$ and $`\varphi `$ define the polar and azimuthal angles of the outgoing muon with momentum $`k^{}|𝐤^{}|`$ and energy $`E_k^{}=\sqrt{k^2+m^2}`$. (For further details we refer the reader to Fig. 2 of Ref. van der Ventel and Piekarewicz (2004)). Finally, to compare the present charged-current calculation to the neutral-current one, we have integrated over the kinematic variables of the outgoing lepton.
The second approximation concerns the evaluation of the nuclear matrix element
$$J^\mu =p^{},s^{};\mathrm{\Psi }_f(P^{})\left|\widehat{J}^\mu (q)\right|\mathrm{\Psi }_i(P).$$
(28)
First, two- and many-body components of the current operator are neglected by assuming that the $`W`$-boson only couples to a single bound neutron. Second, two- and many-body rescattering processes are neglected by assuming that the detected proton is associated with the specific bound neutron to which the $`W`$-boson had coupled to. Further, as we are confident that distortion effects largely factor out from the ratio of cross sections Meucci et al. (2004); Martinez et al. (2005), final-state interactions between the outgoing proton and the residual nucleus will be neglected. Finally, the impulse approximation is invoked by assuming that the weak charged-current operator for a nucleon in the nuclear medium retains its free-space form. That is,
$$\widehat{J}_\mu \widehat{J}_\mu ^{\mathrm{CC}}\widehat{J}_{\mu 5}^{\mathrm{CC}}=F_1(Q^2)\gamma _\mu +iF_2(Q^2)\sigma _{\mu \nu }\frac{q^\nu }{2M}G_A(Q^2)\gamma _\mu \gamma _5.$$
(29)
Here $`M`$ is the nucleon mass and $`F_1`$, $`F_2`$, and $`G_A`$ are Dirac, Pauli, and axial-vector nucleon form factors. Note that the pseudoscalar form factor has been neglected, since its contribution is suppressed by the small lepton mass Kim et al. (1995). A detailed discussion of the weak charge current \[Eq. (29)\] has been reserved to Appendix B. As we have assumed that the ratio of cross sections given in Eqs. (1) and (2) are insensitive to final state interactions between the outgoing proton and the residual nucleus, both initial (bound) and final (free) nucleon propagators may be written in terms of Dirac gamma matrices — rendering the hadronic tensor analytical. Explicit expressions for both propagators and for the analytic (albeit model-dependent) hadron tensor are given in Eqs. (34–38) of Ref. van der Ventel and Piekarewicz (2004).
## III Results
In this section results are presented based on the formalism outlined in Sec. II for charged-changing neutrino scattering from <sup>12</sup>C. The angle-integrated differential cross section \[Eq. (27)\] is shown in Fig. 2 as a function of the kinetic energy $`T_p^{}`$ of the outgoing proton in the laboratory frame for three incident neutrino energies, namely, $`E_k=k=200,500`$ and $`1000`$ MeV. We display separately the contribution to the cross section from the $`1p^{3/2}`$ (solid and long-dashed–short-dashed lines) and $`1s^{1/2}`$ (dashed and dotted lines) orbitals computed in a relativistic mean-field approximation using the NL3 parameter set Lalazissis et al. (1997). Note that because of the finite muon mass, both negative ($`h^{}=1`$) and positive ($`h^{}=+1`$) helicity muons contribute to the cross section; the two smallest contributions correspond to the positive-helicity case. As the energy of the incident neutrino increases, and consequently also that of the outgoing muon, the positive-helicity contribution (which scales as $`m/E_k^{}`$) becomes less and less important until it ultimately disappears at large-enough energy. This can already be observed at $`E_k=500`$ and $`1000`$ MeV.
For the elementary process $`\nu +n\mu ^{}+p`$, the threshold laboratory energy of the incident neutrino is approximately 112 MeV. An additional kinematical constraint that is strongly affected by binding energy corrections follows from energy conservation. Using the fact that $`E_k^{}m`$ we obtain
$$E_k+(M+E_B)m+(T_p^{}+M)T_p^{}E_kE_Bm,$$
(30)
where $`E_B`$ is the (positive) binding energy of the neutron. For <sup>12</sup>C, the NL3 parameter set predicts $`E_B(1s^{1/2})53`$ MeV and $`E_B(1p^{3/2})19`$ MeV. For the particular case of a neutrino incident energy of $`E_k=200`$ MeV, the cross section displays a sharp cut-off for the knockout of the $`1s^{1/2}`$ neutron at an energy of $`T_p^{}40`$ MeV. For higher incident neutrino energies, the maximum allowed value for the kinetic energy of the outgoing proton is already sufficiently large to allow the cross section to fall off smoothly to (almost) zero. Our subsequent results will only focus on incident neutrino energies of 500 and 1000 MeV, as the ratio $`R_{NC/CC}`$ defined in Eq. (2) will be measured by the FINeSSE collaboration with neutrinos in that energy range Brice et al. (2004).
The cross section that results from adding the contributions from both muon helicities ($`h=\pm 1`$) and both neutron orbitals ($`1s^{1/2}`$ and $`1p^{3/2}`$) is depicted in Fig. 3 by the solid line. The dashed and long-dashed–short-dashed lines represent the calculation where we have summed over the two helicity values for the individual $`1p^{3/2}`$ and $`1s^{1/2}`$ orbitals, respectively. The result for the full cross section (solid line) may be compared to Fig. 8 of Ref. Alberico et al. (1997). In the kinematical region in which they can be compared, there is good agreement in both the shape and magnitude of the cross sections.
Next we investigate in Fig. 4 the contribution from the single-nucleon form factors to the differential cross section for incident neutrino energies of $`E_k=500`$ and $`1000`$ MeV. As in the previous figures, the full result is displayed by the solid line. Next in importance is the long-dashed–short-dashed line obtained by setting the weak Pauli form factor to zero ($`F_20`$). The last three lines are obtained from calculations using a single non-zero form factor. That is, the dashed line is obtained from the full calculation by setting $`F_1=F_2=0`$, the dotted line by setting $`G_A=F_2=0`$, and the dashed-dotted line by setting $`G_A=F_1=0`$. This figure clearly illustrates the relatively minor role played by the kinematically suppressed weak Pauli form factor $`F_2`$. Indeed, by itself, it yields a partial cross section that both in magnitude and in shape shows little resemblance to the full cross section. Clearly, the two dominant form factors are the weak Dirac and the axial-vector form factors, with the latter assuming the dominant role. Yet by itself, no single form factor reproduces the full cross section indicating that all interference terms, $`F_1F_2`$, $`F_1G_A`$, and $`F_2G_A`$ are important for the charged-current process. Contrast this to the neutral-current process where the Dirac form factor is strongly suppressed by the weak mixing angle ($`14\mathrm{sin}^2\theta _W0.076`$).
As mentioned earlier, systematic errors with the neutron detection make the ratio of neutral- to charged-current reactions $`R(NC/CC)`$ a more viable alternative than the proton-to-neutron neutral-current ratio $`R(p/n)`$. Thus, we now compare in Fig. 5 the cross section for the charged-current reaction with that for the neutral-current process: $`\nu +X(Z,A)\nu +p+X(Z1,A1)`$. The solid line represents the full charged-current cross section from <sup>12</sup>C, where we have summed over both bound-state orbitals and both muon helicities. The dashed line represents the corresponding neutral-current cross section with no strange-quark contribution to the spin of the proton (ı.e., $`g_A^s0`$); the long-dashed–short-dashed line is the same calculation with $`g_A^s=0.19`$. Results are shown for incident neutrino energies of 500 MeV (top graph) and 1000 MeV (bottom graph).
A comparison to Fig. 8 of Ref. Alberico et al. (1997) shows good agreement in both the shape and magnitude of the cross sections. The axial-vector form factor plays a dominant role in the neutral-current neutrino-proton reaction and makes this reaction particularly sensitive to the strange-quark contribution to the spin of the proton. Recall that the axial-vector form factor of a proton in the neutral-current case is given by van der Ventel and Piekarewicz (2004)
$$\stackrel{~}{G}_A(Q^2)=\left(g_Ag_A^s\right)G_D^A(Q^2)\underset{Q^2=\mathrm{\hspace{0.17em}0}}{}(1.26g_A^s)\underset{g_A^s=0.19}{}1.45.$$
(31)
Here $`G_D^A(Q^2)`$ is the axial-vector form factor of the nucleon (see Appendix B) and a value of $`g_A^s=0.19`$ is assumed for the strange-quark contribution to the spin of the nucleon Horowitz et al. (1993); this value seems to improve the agreement with the Brookhaven National Laboratory experiment E734 Ahrens et al. (1987). Note that this negative value of $`g_A^s`$ leads to an increase in the proton $`\stackrel{~}{G}_A`$ by about 15%.
We are now in a position to display results for the main observable of this work: the ratio of neutral- to charged-current neutrino-nucleus scattering cross sections $`R(NC/CC)`$ defined in Eq. (2). For notational simplicity let us denote the differential cross section $`d\sigma /dE`$ by $`\sigma `$, where it is implied that we have summed over the $`1s^{1/2}`$ and $`1p^{3/2}`$ orbitals of <sup>12</sup>C as well as over the two values of the helicity (when appropriate). Then, because of the dominance of the axial-vector form factor we may write the NC cross sections as
$$\frac{\sigma _{NC}(g_A^s0)}{\sigma _{NC}(g_A^s=0)}\left(1\frac{g_A^s}{g_A}\right)^2.$$
(32)
As the strange-quark contribution to the spin of the nucleon is assumed to be isoscalar, the charged-current reaction is insensitive to it. Thus,
$$\frac{R(NC/CC;g_A^s0)}{R(NC/CC;g_A^s=0)}\left(1\frac{g_A^s}{g_A}\right)^2\underset{g_A^s=0.19}{}1.32.$$
(33)
That is, assuming the dominance of the axial-vector form factor in neutral-current neutrino-proton scattering, a $`30`$% enhancement (for $`g_A^s=0.19`$) in $`R(NC/CC)`$ is expected from a non-zero strange-quark contribution to the spin of the nucleon. The $`R(NC/CC)`$ ratio is plotted in Fig. 6 as a function of the laboratory kinetic energy $`T_p^{}`$ of the outgoing proton for incident neutrino energies of 500 and 1000 MeV. The solid and long-dashed–short-dashed lines correspond to non-zero and zero values of $`g_A^s`$, respectively. Further, the dashed line has been obtained by multiplying the $`g_A^s=0`$ result by a constant enhancement factor of 1.32. The agreement between the solid and dashed lines indicates that the simple estimate given in Eq. (33) is quantitatively correct — especially at small $`T_p^{}`$ (or equivalently small momentum transfer $`q`$) where the contribution from the interference term $`\stackrel{~}{F}_2\stackrel{~}{G}_A`$ remains small. While significant, the sensitivity to $`g_A^s`$ in $`R_{NC/CC}`$ is about a factor of two less than in $`R(p/n)`$ where both proton and neutron NC cross sections are sensitive to $`g_A^s`$. We trust that after working out the systematic uncertainties in the neutron detection, this crucial experiment will also be performed.
We close this section with a brief comment on the mild oscillations displayed by the neutral- to charged-current ratio $`R_{NC/CC}`$, especially at 1000 MeV. Note that this structure is not unique to the CC cross sections (see Fig. 3) but has already been observed in the NC cross sections of Ref. van der Ventel and Piekarewicz (2004) (see Fig. 9). As neither the momentum distribution of the bound nucleons nor the nucleon form factors display any such structure, we attribute this behavior to a kinematical effect originally pointed out in Ref. Horowitz et al. (1993) (see Fig. 1) and later reproduced by us in Figs. 5-6 of Ref. van der Ventel and Piekarewicz (2004). The mild oscillations in the single-differential cross section $`d\sigma /dT_p^{}`$ is a residual effect associated with the existence of a “double-humped” structure in the double-differential cross section $`d^2\sigma /dT_p^{}d(\mathrm{cos}\alpha )`$ (here $`\alpha `$ is the polar angle of the outgoing proton with kinetic energy $`T_p^{}`$). In turn, the emergence of the double-humped structure is a purely kinematical effect that results from the inability of the reaction to produce medium-energy nucleons. That is, high-energy neutrinos are able to produce low- or high-energy nucleons but not medium energy ones. While the double-humped structure is a robust kinematical effect, the integration over $`\alpha `$ introduces model dependences that smooth out — to a greater or lesser degree — some of the structure displayed by $`d^2\sigma /dT_p^{}d(\mathrm{cos}\alpha )`$.
## IV Summary
The distribution of mass, charge, and spin in the proton are among the most fundamental properties in hadronic structure. In this context, a topic that has received tremendous attention for over fifteen years is the contribution of strange quarks to the structure of the proton. In this work we have focused on the strange-quark contribution to the spin of the proton ($`g_A^s`$). Elastic neutrino-proton scattering at low momentum transfer is particularly well suited for this study, as the axial-vector form factor of the proton — the observable that encompasses the spin structure of the proton \[see Eq. (31)\] — dominates this reaction. Indeed, the two “competing” Dirac and Pauli form factors are strongly suppressed, the former by the weak mixing angle ($`14\mathrm{sin}^2\theta _W0.076`$) and the latter by the nucleon mass ($`|𝐪|/2M`$). Yet in an effort to reduce systematic uncertainties related to the neutrino flux, a “ratio-method” has been proposed to extract $`g_A^s`$. Two ratios are particularly useful in this regard: a) proton-to-neutron yields in elastic neutrino scattering \[Eq. (1)\] and b) neutral-current-to-charged-current yields \[Eq. (2)\]. While the former shows a larger sensitivity to $`g_A^s`$, the latter is insensitive to systematic errors associated with neutron detection. As neutrino experiments involve extremely low count rates, these reactions use targets that consists of a combination of free protons and nucleons bound into nuclei. Thus, nuclear-structure effects must be considered.
In the present work we have extended the formalism developed in Ref. van der Ventel and Piekarewicz (2004) for neutral-current neutrino-nucleus scattering to the charged-current reaction. In particular, cross-section ratios have been computed within a relativistic plane wave impulse approximation. Benefiting from work done by others Meucci et al. (2004); Martinez et al. (2005), we justify the omission of final-state interactions by the suggestion that while distortion effects change the overall magnitude of the cross section, they do so without a substantial redistribution of strength. Nuclear-structure effects — which enter in our formalism exclusively in terms of the momentum distribution of the bound nucleons computed at the mean-field level — were incorporated via the accurately calibrated relativistic NL3 parameter set Lalazissis et al. (1997). The validity of the plane-wave approximation yields theoretical cross sections that may be displayed in closed, semi-analytic form. Although the structure of the weak hadronic current is the same for the neutral- and charged-current reactions, a few differences emerge. First, the finite muon mass results in muons produced with both negative and positive helicity. Further, a finite muon mass produces cross sections that display a sharp cut-off for low values of the incident neutrino energy. However, for the range of neutrino energies of interest to the FINeSSE collaboration ($`500`$$`1000`$ MeV) the positive-helicity contribution becomes negligible. Further, while the same three nucleon form factors enter the neutral- and charged-current reactions, their quantitative impact differs considerably. For example, while the Dirac form factor for the neutral-current process is strongly suppressed ($`\stackrel{~}{F}_1(Q^2=0)=0.076`$) it is large for the charged-current process ($`F_1(Q^2=0)=1`$). Hence, no single form factor dominates the charged-changing reaction. More importantly, as the strange-quark content of the nucleon is assumed to be isoscalar, the purely isovector CC reaction is insensitive to the strange-quark content of the nucleon. This renders the ratio $`R(NC/CC)`$ less sensitive to strange quark effects (by about a factor of 2) than the neutral-current ratio $`R(p/n)`$. Still, for the value of $`g_A^s=0.19`$ adopted in this work Horowitz et al. (1993), a 30% enhancemenent in $`R(NC/CC)`$ is obtained relative to a calculation with $`g_A^s=0`$. We note that our results for the charged-current cross sections were compared to similar calculations done in Ref. Alberico et al. (1997) and good agreement was found in both the shape and the magnitude of the cross section.
In summary, the sensitivity of the ratio of neutral- to charged-current cross sections to the strange-quark contribution to the spin of the nucleon $`g_A^s`$ was investigated in a relativistic plane wave impulse approximation. The enormous advantage of this formalism is that our theoretical results may be displayed in closed, semi-analytic form. The central motivation behind this work is the proposed FINeSSE program that aims to measure $`g_A^s`$ with unprecedented accuracy via the neutral- to charged-current ratio $`R(NC/CC)`$. By adopting a value of $`g_A^s=0.19`$, an increase in this ratio of approximately 30% was found relative to the $`g_A^s=0`$ result. While sensitive, this is less so than the corresponding ratio of proton-to-neutron yields $`R(p/n)`$ in neutral current neutrino-induced reactions. This measurement, however, has been hindered by difficulties associated with neutron detection. We trust that this difficulty may be overcome so that this crucial program may get off the ground.
###### Acknowledgements.
B.I.S.v.d.V gratefully acknowledges the financial support of the University of Stellenbosch and the National Research Foundation of South Africa. This material is based upon work supported by the National Research Foundation under Grant number GUN 2048567 (B.I.S.v.d.V) and by the United States Department of Energy under Grant number DE-FG05-92ER40750 (J.P.).
## Appendix A Leptonic-hadronic contraction
In Sec. II.1 it was shown that the charged-current cross section could be expressed as the contraction of a leptonic tensor $`\mathrm{}^{\mu \nu }`$ \[Eq. (14)\] with a hadronic tensor $`W^{\mu \nu }`$ \[Eq. (24)\] written in a model-independent way in terms of thirteen independent structure functions. In this appendix we carry out the contraction, which we separate into symmetric and antisymmetric parts. That is,
$$\mathrm{}^{\mu \nu }W_{\mu \nu }=\mathrm{}_{\mu \nu }^SW_S^{\mu \nu }+\mathrm{}_{\mu \nu }^AW_A^{\mu \nu }.$$
(34)
Here the symmetric part is given by
$`\left({\displaystyle \frac{4}{kE_k^{}}}\right)^1\mathrm{}_{\mu \nu }^SW_S^{\mu \nu }`$ $`=`$ $`(W_1(kK^{})+W_2f_1(q)+W_3f_1(P)+W_4f_1(p^{})`$ (35)
$`+W_5f_2(P,q)+W_6f_2(q,p^{})+W_7f_2(P,p^{})),`$
while the antisymmetric part by
$$\left(\frac{4}{kE_k^{}}\right)^1\mathrm{}_{\mu \nu }^AW_A^{\mu \nu }=i\left(W_{10}\epsilon ^{\mu \nu \alpha \beta }k_\mu K_\nu ^{}P_\alpha p^\beta +W_{11}f_3(q,P)+W_{12}f_3(q,p^{})+W_{13}f_3(P,p^{})\right).$$
(36)
Note that the following four-vector has been defined:
$$K^{}\frac{1}{2}(k^{}h^{}ms)\underset{m=0}{}k^{}\delta _{h^{},1}.$$
(37)
Further, for simplicity the following three functions have been introduced:
$`f_1(x)=2(kx)(K^{}x)x^2(kK^{}),`$ (38a)
$`f_2(x,y)=(kx)(K^{}y)+(ky)(K^{}x)(xy)(kK^{}),`$ (38b)
$`f_3(x,y)=(ky)(K^{}x)(kx)(K^{}y).`$ (38c)
From Eq. (35) we see that the three structure functions $`W_2`$, $`W_5`$ and $`W_6`$ do contribute to charged-current neutrino-nucleus scattering, in contrast to the neutral-current case (see Eq. (22a) of Ref. van der Ventel and Piekarewicz (2004)). This is due to the lack of conservation of the symmetric part of the leptonic tensor as a result of the finite muon mass \[see Eq. (19a)\]. Note, however, that in the massless limit $`f_1(q)=f_2(P,q)=f_2(q,p^{})=0`$, as required. Finally, due to the form of Eq. (19b), the charged-current process remains insensitive to the two structure functions $`W_8`$ and $`W_9`$.
## Appendix B Single nucleon form factors
In Sec. II.3 it was shown that in the impulse approximation, the single nucleon current probed in the charge-changing reaction may be written in the following standard form:
$$\widehat{J}_\mu \widehat{J}_\mu ^{\mathrm{CC}}\widehat{J}_{\mu 5}^{\mathrm{CC}}=F_1(Q^2)\gamma _\mu +iF_2(Q^2)\sigma _{\mu \nu }\frac{q^\nu }{2M}G_A(Q^2)\gamma _\mu \gamma _5,$$
(39)
where $`F_1`$, $`F_2`$, and $`G_A`$ are the Dirac, Pauli, and axial-vector form factors, respectively and the pseudoscalar form factor has been neglected. To understand the structure of the vector form factors ($`F_1`$ and $`F_2`$) we invoke the conservation of the vector current (CVC) hypothesis. To start, one parametrizes the nucleon matrix elements of the isovector electromagnetic current in the following standard form:
$`N(𝐩^{},s^{},t^{})|\widehat{J}_\mu ^{\mathrm{EM}}(T=1)|N(𝐩,s,t)=N(𝐩^{},s^{},t^{})|\overline{q}\gamma _\mu {\displaystyle \frac{\tau _3}{2}}q|N(𝐩,s,t)`$
$`\overline{U}(𝐩^{},s^{})\left[F_1^{(1)}(Q^2)\gamma _\mu +iF_2^{(1)}(Q^2)\sigma _{\mu \nu }{\displaystyle \frac{q^\nu }{2M}}\right]U(𝐩,s)\left(\tau _3\right)_{tt^{}},`$ (40)
where $`\overline{q}=(\overline{u},\overline{d})`$ is an isospin doublet of quark fields, and $`F_1^{(1)}`$ and $`F_2^{(1)}`$ are the isovector Dirac and Pauli form factors of the nucleon, respectively. In turn, these are given in terms of proton and neutron electromagnetic form factors as follows:
$$F_i^{(1)}(Q^2)=\frac{1}{2}\left(F_i^{(p)}(Q^2)F_i^{(n)}(Q^2)\right),(i=1,2).$$
(41)
The CVC hypothesis is a powerful relation that assumes that the vector part of the weak charge-changing current may be directly obtained from the isovector component of the electromagnetic current. That is,
$$\widehat{J}_\mu ^{\mathrm{EM}}(T=1)=\widehat{V}_\mu ^{(3)}=\overline{q}\gamma _\mu \frac{\tau _3}{2}q,\widehat{J}_\mu ^{\mathrm{CC}}(\pm )=\widehat{V}_\mu ^{(1)}\pm i\widehat{V}_\mu ^{(2)}=\overline{q}\gamma _\mu \left(\frac{\tau _1\pm i\tau _2}{2}\right)q.$$
(42)
Thus, a determination of the electromagnetic form factors of the nucleon — which has been done experimentally — fixes the vector part of the charge-changing currents to:
$$N(𝐩^{},s^{},t^{})|\widehat{J}_\mu ^{\mathrm{CC}}(\pm )|N(𝐩,s,t)=\overline{U}(𝐩^{},s^{})\left[F_1^{(1)}(Q^2)\gamma _\mu +iF_2^{(1)}(Q^2)\sigma _{\mu \nu }\frac{q^\nu }{2M}\right]U(𝐩,s)\left(2\tau _\pm \right)_{tt^{}}.$$
(43)
In this way the vector form factors of Eq. (39) are then simply given by
$$F_i(Q^2)=2F_i^{(1)}(Q^2)=F_i^{(p)}(Q^2)F_i^{(n)}(Q^2),(i=1,2).$$
(44)
Paraphrasing Ref. Walecka (1995): CVC implies that the vector part of the single-nucleon matrix element of the charge-changing weak current, whatever the detailed dynamic structure of the nucleon, can be obtained from elastic electron scattering through the electromagnetic interaction!
A similar procedure may be followed to determine the axial-vector form factor $`G_A`$ in terms of the isovector axial-vector current. That is,
$$\widehat{J}_{\mu 5}^{\mathrm{CC}}(\pm )=\widehat{A}_\mu ^{(1)}\pm i\widehat{A}_\mu ^{(2)}=\overline{q}\gamma _\mu \gamma _5\left(\frac{\tau _1\pm i\tau _2}{2}\right)q,$$
(45)
so that
$$N(𝐩^{},s^{},t^{})|\widehat{J}_{\mu 5}^{\mathrm{CC}}(\pm )|N(𝐩,s,t)G_A(Q^2)\overline{U}(𝐩^{},s^{})\gamma _\mu \gamma _5U(𝐩,s)\left(\tau _\pm \right)_{tt^{}}.$$
(46)
As before, the above expression neglects the contribution from the pseudoscalar form factor.
We finish this section by parameterizing the various nucleon form factors in terms of their known $`Q^2=0`$ values times form factors of a dipole form. This is identical to the procedure employed in Appendix A of Ref. van der Ventel and Piekarewicz (2004)) for the neutral-current reaction. We obtain
$`F_1^{(p)}(Q^2)=\left({\displaystyle \frac{1+\tau (1+\lambda _p)}{1+\tau }}\right)G_D^V(Q^2),F_2^{(p)}(Q^2)=\left({\displaystyle \frac{\lambda _p}{1+\tau }}\right)G_D^V(Q^2),`$ (47a)
$`F_1^{(n)}(Q^2)=\left({\displaystyle \frac{\lambda _n\tau (1\eta )}{1+\tau }}\right)G_D^V(Q^2),F_2^{(n)}(Q^2)=\left({\displaystyle \frac{\lambda _n(1+\tau \eta )}{1+\tau }}\right)G_D^V(Q^2),`$ (47b)
$`G_A(Q^2)=g_AG_D^A(Q^2),`$ (47c)
where a dipole form factor of the following form is assumed:
$`G_D^V(Q^2)=(1+Q^2/M_V^2)^2=(1+4.97\tau )^2`$ (48a)
$`G_D^A(Q^2)=(1+Q^2/M_A^2)^2=(1+3.31\tau )^2`$ (48b)
$`\eta =(1+5.6\tau )^1\tau =Q^2/(4M^2).`$ (48c)
Finally, for reference we display the value of the various nucleon form factors at $`Q^2=0`$
$`F_1^{(p)}(0)=1,F_1^{(n)}(0)=0,`$ (49a)
$`F_2^{(p)}(0)=\lambda _p=+1.79,F_2^{(n)}(0)=\lambda _n=1.91,`$ (49b)
$`G_A(0)=g_A=+1.26.`$ (49c)
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# Untitled Document
Sound-propagation gap in fluid mixtures
Supurna Sinha and M. Cristina Marchetti
Physics Department, Syracuse University, Syracuse, New York 13244
We discuss the behavior of the extended sound modes of a dense binary hard-sphere mixture. In a dense simple hard-sphere fluid the Enskog theory predicts a gap in the sound propagation at large wave vectors. In a binary mixture the gap is only present for low concentrations of one of the two species. At intermediate concentrations sound modes are always propagating. This behavior is not affected by the mass difference of the two species, but it only depends on the packing fractions. The gap is absent when the packing fractions are comparable and the mixture structurally resembles a metallic glass.
I. INTRODUCTION
Short-wavelength collective modes in dense simple hard-sphere fluids have been studied extensively by both theoretical and experimental researchers in the past few years.<sup>1-4</sup> Two interesting features of the extended hydrodynamic modes in a simple fluid are the softening of the heat mode at the wave vector $`k`$ where the static structure factor $`S(k)`$ has its first maximum and the appearance of a gap in the sound propagation in the same large-wave-vector region.<sup>l,2</sup> The softening of the heat mode corresponds to the slowing down of structural relaxation in a dense fluid and has been discussed extensively in the literature. The sound-propagation gap consists in the vanishing of the imaginary part of the frequency of the sound waves, which describes the propagation of the wave, over a finite region of $`k`$ values.
In a one-component fluid of hard spheres of size $`\sigma `$ and density $`n`$ one obtains a gap in the sound propagation for all values of the packing fraction $`\eta =(\pi /6)n\sigma ^3`$ above $`\eta 0.33.^5`$ This feature of the extended modes has been the topic of some controversy. de Schepper, Cohen, and Zuilhof,<sup>5</sup> have shown that the appearance and disappearance of the sound propagation gap can be understood as arising from the competition between generalized wave-vector-dependent reversible elastic restoring forces and dissipative forces in the fluid. Roughly speaking, a gap arises whenever the dissipation exceeds the elastic forces. They also pointed out that the behavior of the sound modes in a one-component fluid is quite insensitive to the static structure factor of the fluid and the gap is present even when setting $`S(k)=1.`$ de Schepper, Van Rijs, and Cohen<sup>6</sup> have also shown that even the long-wavelength Navier-Stokes equations, when used outside their range of validity to describe short-wavelength phenomena, predict a gap in the sound propagation. The model based on the Navier-Stokes equations suggests that there is no connection between local effects leading to wave-vector-dependent transport coefficients and susceptibilities and the trapping of sound on molecular length scales.
In a recent paper (hereafter denoted I) we used the Enskog equation as the starting point to study the short-wavelength collective excitations of a binary hard-sphere fluid mixture.<sup>7</sup> In a mixture the spatial ordering of the particles can be changed while keeping the total density of the fluid constant by changing the concentration $`x_2=n_2/(n_1+n_2)`$ of the larger species (here type-2 spheres) or the size ratio $`\alpha =\sigma _1/\sigma _2.`$ We find that sound propagation in the mixture depends on the concentration $`x_2`$: the trapping of sound modes and the corresponding propagation gap that arises in dense simple fluids only occurs in mixtures with a small concentration of small spheres in a dense fluid of large spheres $`(x_21)`$ or a small concentration of large spheres in a dense fluid of small spheres $`(x_21)`$. At intermediate values of $`x_2`$ we observe a softening of the sound propagation, but no actual gap. For either small or large values of $`x_2`$ the mixture closely resembles a simple fluid, as can be seen from the behavior of the partial structure factors $`S_{11}(k),S_{22}(k)\mathrm{and}S_{12}(k)`$ defined in I. For instance, for $`x_2=0.01`$ only the static structure factor $`S_{11}(k)`$ of the dense component is peaked \[see Fig. 1(a)\]. The other structure factors are essentially constant. Thus we expect the appearance of a sound-propagation gap as in a simple fluid. At intermediate values of the concentration $`x_2`$ the dense mixture resembles a metallic glass: there is short-range order in the spatial arrangements of both types of spheres, as indicated by the fact that all three static structure factors display considerable structure as functions of the wave vector \[see Fig. l(b)\]. The absence of a sound-propagation gap in this case is then consistent with the experimental finding that no gap occurs in the second sound of metallic glasses.<sup>8</sup> This observation seems to indicate that the local fluid structure does play a role in determining sound propagation and sound trapping at short wavelengths, in contrast to what was argued for a one-component fluid.
In this paper we investigate in more detail the extended sound modes in binary mixtures, with the objective of clarifying their dependence on concentration, size ratio and mass ratio, and the role of the fluid structure in determining the gap. As in our earlier paper, we keep the total packing fraction of the mixture $`\eta =(\pi /6)(n_1\sigma _1^3+n_2\sigma _2^3)`$ fixed at the value $`\eta =0.46`$, which is close to the packing fraction corresponding to freezing of a one-component hard-sphere fluid.
FIG.1. The partial static structure factors $`S_{11}(k),S_{22}(k)`$, and $`S_{12}(k)`$ for $`\eta =0.46`$, $`\alpha =0.7`$, $`x_2=0.01`$ (a), and $`x_2=0.5`$(b).
First we have considered the dependence of the sound modes on the ratio $`m_1/m_2`$ of the masses of the two species and the role of this ratio in determining the propagation gap. By comparing the extended hydrodynamic modes displayed in I $`(\alpha =0.7`$ and $`m_1/m_2=0.5)`$ to the modes obtained for the same values of size ratio and concentration, but equal masses, we have found that the sound propagation is essentially unchanged at all concentrations. This holds true for other values of the size ratio. In general, as long as the mass ratio is not too small $`(m_1/m_20.1)`$ the difference in mass of the two species affects only weakly the sound propagation in a mixture. This points to the fact that the size difference and the spatial ordering of the two types of spheres must play the main role in determining the absence of sound trapping at large wave vectors.
In Sec. II we discuss the dependence of the sound modes on the concentration of one of the two species. We show that in dense mixtures, as in a one-component fluid, the appearance of a sound propagation gap can be understood as the result of the competition between elastic and dissipative forces in the fluid. The absence of the sound-propagation gap in mixtures of intermediate concentrations is not determined by the details of the static structure factors, but it does depend on the fluid structure in the sense that it occurs when the packing fractions $`\eta _1=(\pi /6)n_1\sigma _1^3`$ and $`\eta _2=(\pi /6)n_2\sigma _2^3`$ of the two species are comparable. In this case all three static structure factors are peaked and the structural properties of the mixture resemble those of a metallic glass.
II. DEPENDENCE OF SOUND PROPAGATION
ON CONCENTRATION
We showed in I that at large wave vectors the extended heat and diffusion modes govern the relaxation of the densities of the two species. These are the only long-lived fluctuations in the fluid at these wave vectors. At large wave vectors the extended sound modes mainly describe the relaxation of temperature and longitudinal momentum fluctuations. Similarly, in a one-component fluid large-wave-vector density fluctuations are long-lived while temperature and momentum fluctuations relax quickly. On the basis of this observation Zuilhof and co-workerss suggested that extended sound modes in a one-component fluid can be described by a simple model obtained by neglecting the coupling to the density in the generalized hydrodynamic equations for temperature and longitudinal momentum. The resulting two coupled equations have eigenfrequencies that closely reproduce the extended sound modes at large wave vectors. Clearly, at small wave vectors these modes do not reduce to the hydrodynamic sound modes.
The same model can be used to describe the extended sound modes of a mixture, where the characteristic time scale for the relaxation of temperature and momentum fluctuations is well separated from that governing the relaxation of fluctuations in the two densities. If we neglect the coupling to fluctuations in the densities in the generalized hydrodynamic equations for temperature and longitudinal momentum, we obtain two coupled equations that are formally identical to those obtained for a one-component fluid. The eigenvalues of these equations are obtained by solving a quadratic equation, with the result
$`z_\pm (k)=`$ $``$ $`{\displaystyle \frac{1}{2}}[\mathrm{\Omega }_{ll}(k)\mathrm{\Omega }_{TT}(k)]`$ (2.1)
$`\pm `$ $`{\displaystyle \frac{1}{2}}\{[\mathrm{\Omega }_{ll}(k)\mathrm{\Omega }_{TT}(k)]^24[f(k)]^2\}^{\frac{1}{2}},`$
where $`\mathrm{\Omega }_{ll}(k)`$ and $`\mathrm{\Omega }_{TT}(k)`$ are the damping rates in the equations for longitudinal momentum and temperature fluctuations, respectively, and $`f(k)=i\mathrm{\Omega }_{Tl}(k)`$ is the elastic restoring force that couples the two equations,
$`\mathrm{\Omega }_{ll}(k)={\displaystyle \frac{2}{3\rho }}{\displaystyle \underset{a=1,2}{}}{\displaystyle \underset{b=1,2}{}}{\displaystyle \frac{2\mu _{ab}\sqrt{n_an_b}}{t_{E_{ab}}}}[1`$ $``$ $`j_0(k\sigma _{ab})`$ (2.2)
$`+`$ $`2j_2(k\sigma _{ab})]`$
$$\mathrm{\Omega }_{TT}(k)=\frac{2}{3\rho }\underset{a=1,2}{}\underset{b=1,2}{}\frac{2\mu _{ab}\sqrt{n_an_b}}{t_{E_{ab}}}[1j_0(k\sigma _{ab})]$$
(2.3)
$`f(k)=k\sqrt{2n/3\beta \rho }[1+{\displaystyle \underset{a=1,2}{}}{\displaystyle \underset{b=1,2}{}}\mathrm{\hspace{0.33em}\hspace{0.33em}2}\pi {\displaystyle \frac{n_an_b}{n}}\sigma _{ab}^3`$
$`\times \chi _{ab}{\displaystyle \frac{j_1(k\sigma _{ab})}{k\sigma _{ab}}}],`$ (2.4)
where $`n=n_1+n_2`$ and $`\rho =m_1n_1+m_2n_2`$ are the total number and mass densities, respectively. Also, $`\mu _{ab}=m_am_b/(m_a+m_b)`$ is the reduced mass, $`\chi _{ab}`$ the pair correlation function of species $`a`$ and $`b`$ at contact, $`\sigma _{ab}=(\sigma _a+\sigma _b)/2`$, and $`j_n(x)`$ is a spherical Bessel function of order $`n`$. Finally, $`t_{E_{ab}}`$ is the Enskog mean-free time, $`t_{E_{ab}}=\sqrt{2\mu _{ab}\beta }/4\sqrt{\pi n_an_b}\sigma _{ab}^2\chi _{ab})`$. At large wave vectors the two modes given in Eq. (2.1) closely resemble the extended sound modes obtained in I by solving the four coupled hydrodynamic equations. The argument of the square root in Eq. (2.1) is generally negative (and the two modes are propagating), unless the dissipative damping exceeds the elastic forces. In this case Eq. (2.1) yields two real roots, corresponding to diffusive modes. This is the region of the sound-propagation gap. The argument of the square root in Eq. (2.1) can be factorized as $`(\mathrm{\Omega }_{ll}\mathrm{\Omega }_{TT}+2f)(\mathrm{\Omega }_{ll}\mathrm{\Omega }_{TT}2f)`$. The first factor is always positive. We define the function $`\mathrm{\Delta }(k)`$ by
$`{\displaystyle \frac{\sigma _{12}}{2}}\left[{\displaystyle \frac{3\beta \rho }{2n}}\right]^{1/2}(\mathrm{\Omega }_{ll}\mathrm{\Omega }_{TT}2f)=\mathrm{\Delta }(k)k\sigma _{12},`$
with
$$\mathrm{\Delta }(k)=\frac{2\pi }{n}\underset{a=1,2}{}\underset{b=1,2}{}n_an_b\sigma _{ab}^2\sigma _{12}\chi _{ab}\left[\left[\frac{16\mu _{ab}n}{3\pi \rho }\right]^{1/2}j_2(k\sigma _{ab})j_1(k\sigma _{ab})\right].$$
(2.5)
A gap occurs when $`\mathrm{\Delta }(k)>k\sigma _{12}`$. This condition is displayed graphically in Fig. 2 at three different concentrations. The figure shows that a gap is obtained for $`x_2=0.01`$ and $`x_2=0.9`$, but not for $`x_2=0.5`$. For $`\alpha =0.7`$ and $`m_1/m_2=0.5`$ this simple model predicts no sound-propagation gap for $`0.050.7`$, in agreement with what was obtained in I solving the full coupled hydrodynamic equations.
FIG.2. The straight line is $`k\sigma _{12}`$. The curves represent the function $`\mathrm{\Delta }(k)`$ for $`\eta =0.46`$, $`\alpha =0.7`$, and $`m_1/m_2=0.5`$, and $`x_2=0.01`$ (solid line), $`x_2=0.5`$ (dashed line), and $`x_2=0.9`$ (chain dashed line), as a function of $`k\sigma _{12}`$.
In mixtures, as in a one-component fluid, the sound-propagation gap occurs as a consequence of the competition between dissipative and reversible forces. In most of the wave-vector range elastic forces exceed dissipative forces. This simply reflects the fact that sound modes propagate in both liquids and low-density gases (in the very-large-wave-vector limit the fluid resembles an ideal gas, since one is considering length scales that are too short for the interactions to be relevant). The damping term and the elastic forces depend on the packing fractions $`\eta _1`$ and $`\eta _2`$. For a fixed total packing fraction $`\eta `$ and size ratio $`(\alpha =0.7)`$, both dissipative and elastic terms decrease with increasing $`x_2`$. At all $`x_2`$ there is a region of large $`k`$ where the dissipation rates are essentially constant (and equal the sum of the Enskog times weighted by the mass fractions) and $`\mathrm{\Omega }_{ll}\mathrm{\Omega }_{TT}0`$. For $`x_21`$ and $`x_2`$ the elastic forces vanish even more rapidly than the dissipative term in this large-$`k`$ region and a gap occurs. For intermediate $`x_2`$ the elastic forces remain finite at all $`k`$ and the modes are always propagating.
In a recent paper<sup>9</sup> Campa and Cohen suggested that the static structure factors may play a role in determining sound-propagation gaps in a mixture. On the other hand, the simple model described above does not contain the static structure factors $`S_{ij}(k)`$. In fact, the appearance and disappearance of the sound-propagating gap are insensitive to the structure factors. This can also be seen by reconsidering the extended sound modes obtained in I from the solution of the four coupled hydrodynamic equation and arbitrarily setting $`S_{11}(k)=S_{22}(k)1`$ and $`S_{12}(k)=0`$ in these equations. One finds that, while the gap itself, when it occurs, is wider if the proper static structure factors are included, its concentration dependence, i.e., its appearance and disappearance, is qualitatively unchanged. The suggestion of Campa and Cohen is, however, in the right direction and it would be misleading to say that short-wavelength sound modes in a mixture are not affected by the fluid structure. The latter does enter even in the simple model yielding Eq. (2.1) through the dependence on the packing fractions of the two components of the mixture. When the packing fractions are both appreciable, there is not gap. This situation corresponds to a fluid where all three partial static structure factors are peaked (Fig. I) and the fluid structure resembles that of a metallic glass.
III. CONCLUSION
We conclude with two remarks.
1. In I we neglected the coefficient of thermal diffusion in evaluating the extended hydrodynamic modes of the mixture. This approximation was motivated by the fact that the coefficient of thermal diffusion vanishes in a first-Sonine-polynomial approximation, and in the long-wavelength limit it is always much smaller than the diffusion coefficient in mixtures of spheres of not too disparate sizes and masses.<sup>1O</sup> On the other hand, in the long-wavelength limit the coefficient of thermal diffusion contributes to the damping of the sound modes.<sup>11</sup> Since neglecting thermal diffusion is best justified in mixtures of either very low $`x_2`$ or $`x_21`$, one may ask if the absence of the sound-propagation gap at intermediate $`x_2`$ simply occurs because we have underestimated the sound damping by neglecting thermal diffusion. The model used here for the description of the extended sound modes based on the two coupled equations for temperature and longitudinal momentum fluctuations neglects all couplings to the densities and therefore does not contain the coefficient of thermal diffusion. This model still predicts the disappearance and reappearance of the sound propagation gap as a function of concentration, indicating that the dependence of the sound propagation on concentration is not an artifact of the approximation of neglecting thermal diffusion.
2. Previous work has focussed on unusual sound propagation in disparate-mass gas mixtures and binary liquid alloys at moderate density.<sup>9,12</sup> One interesting feature observed in computer simulations<sup>13</sup> is the appearance of a fast propagating sound mode above a certain nonzero value of the wave vector, signaling an effective separation of the dynamics of light and heavy particles. Campa and Cohen<sup>9</sup> have shown that the fast sound is a kinetic mode rather than a hydrodynamic mode and it yields an observable shoulder in the dynamic structure factor of gas mixtures. In moderately dense mixtures the fast sound mode is present only for mass ratios below 0.1 and no fast sound occurs for $`\eta 0.42`$. Our generalized hydrodynamic equations as derived in I cannot yield a fast sound mode since our set of independent variables only includes the conserved densities of the fluid: we only obtain the extended hydrodynamic modes and no kinetic modes. To account for the possibility of fast sound we would need to enlarge our set of independent variables to include at least the relative momentum density and the temperature difference of the two species. On the other hand, the fact that no fast sound is observed in Ref. 9 for the values of the total packing fraction and mass ratio considered here suggests that at such high densities the dynamics of light and heavy particles never decouple and the generalized hydrodynamic equations of I may indeed provide an adequate description of the dynamics at large wave vectors.<sup>14</sup> A proper test of whether our theory is relevant for real mixtures will, however, only come from a detailed comparison of the dynamic structure factor that can be obtained from the generalized hydrodynamic equations of I with neutron-scattering spectra from dense mixtures or with numerical simulations. Finally, the fact that the appearance and disappearance of the sound-propagation gap in our work does not depend on the mass ratio and takes place even for equal masses seems to indicate that the fast sound and the concentration dependence of the gap are not related.
ACKNOWLEDGMENTS
One of us (M.C.M.) thanks T. R. Kirkpatrick for stimulating discussions. This work was supported by the National Science Foundation under Contract No. DMR-87-17337.
$``$
* I. M. de Schepper and E. G. D. Cohen, Phys. Rev. A 22, 287 (1980); J. Stat. Phys. 27, 223 (1982).
* T. R. Kirkpatrick, Phys. Rev. A 32, 3120 (1985).
* J. R. D. Copley and J. M. Rowe, Phys. Rev. Lett. 32, 49 (1973).
* A. A. van Well, P. Verkerk, L. A. De Graaf, J. B. Suck, and J. R. D. Copley, Phys. Rev. A 31, 3391 (1985).
* I. M. de Schepper, E. G. D. Cohen, and M. J. Zuilhof, Phys. Lett. 103A, 120 (1984).
* I. M. de Schepper, J. C. Van Rijs, and E. G. D. Cohen, Physica 1434A, 1 (1985).
* M. C. Marchetti and S. Sinha, Phys. Rev. A 41, 3214 (1990).
* J. B. Suck, H. Rudin, H. J. Guntherodt, and H. Beck, Phys. Rev. Lett. 50, 49 (1983).
* A. Campa and E. G. D. Cohen, Phys. Rev. A 41, 5451 (1990).
* J. M. Kincaid, E. G. D. Cohen, and M. Lopez de Haro, J. Chern. Phys. 86, 963 (1987).
* J. P. Boon and S. Yip, Molecular Hydrodynamics (McGraw-Hill, New York, 1980), pp. 270 and 271.
* P. B. Lerner and I. M. Sokolov, Physica C 150, 465 (1988).
* J. Bosse, G. Jacucci, M. Ronchetti, and W. Schirmacher, Phys. Rev. Lett. 57, 3277 (1986).
* 1n a one-component hard-sphere fluid with $`\eta =0.46`$ the generalized hydrodynamic modes, provide a good representation of the dynamic structure factor up to $`k10/\sigma `$, as discussed, for instance, by E. G. D. Cohen, in Trends in Applications of Pure Mathematics to Mechanics, Vol. 249 of Lecture Notes in Physics, edited by E. Kroner and K. Krichgasser (Springer-Verlag, Berlin 1986).
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# Bright solitons in coupled defocusing NLS equation supported by coupling: Application to Bose-Einstein Condensation1footnote 11footnote 1E-mail address: [email protected] (S.K. Adhikari).
## Abstract
We demonstrate the formation of bright solitons in coupled defocusing nonlinear Schrödinger (NLS) equation supported by attractive coupling. As an application we use a time-dependent dynamical mean-field model to study the formation of stable bright solitons in two-component repulsive Bose-Einstein condensates (BECs) supported by interspecies attraction in a quasi one-dimensional geometry. When all interactions are repulsive, there cannot be bright solitons. However, bright solitons can be formed in two-component repulsive BECs for a sufficiently attractive interspecies interaction, which induces an attractive effective interaction among bosons of same type.
Recently, there have been successful observation exp1 ; exp4 and associated theoretical yyy studies of two-component Bose-Einstein condensates (BEC). There have also been experimental observation sol of bright solitons in a BEC formed due to atomic attraction and related theoretical investigations solt . In fiber optics true one-dimensional solitons are formed in the nonlinear Schrödinger (NLS) equation 1 ; agra . The solitons of BEC sol are formed in a quasi-one-dimensional geometry achieved by employing a strong transverse trap. In either case, no solitons can be formed for repulsive interactions.
In this Letter we suggest the possibility of the formation of stable bright solitons in two-component BECs in the presence of repulsive interaction among like atoms and attractive interaction among atoms of different types. We find that a sufficiently strong interspecies attraction can induce an effective attraction among like atoms responsible for the formation of bright solitons in two-component BECs. We also consider similar two-component BECs in the presence of a periodic optical-lattice potential, which are now routinely used in experiments on BECs ol . We consider the formation of these solitons using a coupled time-dependent mean-field Gross-Pitaevskii (GP) equation 11 .
Because of two components such solitons of the coupled one-dimensional NLS equation are often termed vector solitons agra in nonlinear fiber optics, where two orthogonally-polarized pulses of different widths and peak powers in general propagate undistorted. The vector solitons considered so far in nonlinear optics were always created in focusing (attractive) medium. To the best of our knowledge the present study is the first to consider the possibility of creation of vector solitons in defocusing (repulsive) medium.
The experimental study of bright solitons in quasi-one-dimensional attractive systems is quite delicate due to the possibility of collapse in such systems, although a true one-dimensional system does not exhibit collapse 11 . The two-component repulsive BECs with interspecies attraction are better suited for studying solitons as such systems may not easily collapse ska and one can have a controlled study of solitons. Experimentally, this could be realized by forming a coupled repulsive BEC in a cigar-shaped geometry and then transforming the interspecies repulsion to attraction via a Feshbach resonance fs and eventually removing the axial trap so that the BEC components attain mobility in the axial direction like soliton under the action of radial trapping alone.
Bright solitons are really eigenfunctions of the one-dimensional NLS equation. However, the experimental realization of bright solitons in trapped attractive cigar-shaped BECs has been possible under strong transverse binding which, in the case of weak or no axial binding, simulates the ideal one-dimensional situation for the formation of bright solitons. The dimensionless NLS equation in the attractive or focusing case 1
$$iu_t+u_{xx}+|u|^2u=0,$$
(1)
sustains the following bright soliton 1 :
$`u(x,t)`$ $`=`$ $`\sqrt{2w}\text{sech}[\sqrt{w}(x\delta +2vt)]`$ (2)
$`\times `$ $`\mathrm{exp}[iv(x\delta )+i(wvv)t+i\sigma ],`$
with four parameters. The parameter $`w`$ represents the amplitude as well as pulse width, $`v`$ represents velocity, the parameters $`\delta `$ and $`\sigma `$ are phase constants. The bright soliton profile is easily recognized for $`v=\delta =0`$ when Eq. (1) leads to $`|u(x,t)|=\sqrt{2w}\text{sech}(x\sqrt{w})`$. In the case of the coupled focusing NLS equations:
$`iu_t+u_{xx}+(|u|^2+|v|^2)u=0,`$ (3)
$`iv_t+v_{xx}+(|v|^2+|u|^2)v=0,`$ (4)
one could have the following bright vector solitons agra ; mana : $`u(x,t)=\mathrm{cos}(\theta )\sqrt{2w}\text{sech}(x\sqrt{w})\mathrm{exp}(it),`$ $`v(x,t)=\mathrm{sin}(\theta )\sqrt{2w}\text{sech}(x\sqrt{w})\mathrm{exp}(it),`$ where $`\theta `$ is an arbitrary angle. Equations (3) and (4) are incoherently coupled as the coupling depends only on the intensities and is therefore phase insensitive. In this Letter we consider the vector solitons in Eqs. (3) and (4) with defocusing diagonal nonlinearities.
The time-dependent Bose-Einstein condensate wave function $`\mathrm{\Psi }(𝐫,t)`$ at position $`𝐫`$ and time $`t`$ may be described by the following mean-field nonlinear GP equation 11
$`[i\mathrm{}{\displaystyle \frac{}{t}}{\displaystyle \frac{\mathrm{}^2_𝐫^2}{2m}}+V(𝐫)+gn]\mathrm{\Psi }(𝐫,t)=0,`$ (5)
with normalization $`𝑑𝐫|\mathrm{\Psi }(𝐫,t)|^2=1.`$ Here $`n|\mathrm{\Psi }(𝐫,t)|^2`$ is the boson probability density, $`g=4\pi \mathrm{}^2aN/m`$, with $`a`$ the boson-boson scattering length, $`m`$ the mass and $`N`$ the number of bosonic atoms in the condensate. The trap potential with axial symmetry may be written as $`V(𝐫)=\frac{1}{2}m\omega ^2(\rho ^2+\nu ^2z^2)`$ where $`\omega `$ and $`\nu \omega `$ are the angular frequencies in the radial ($`\rho `$) and axial ($`z`$) directions with $`\nu `$ the anisotropy parameter, which will be taken to be 0 for axially free solitons in the following.
In the presence of two types of bosons each of mass $`m,`$ Eq. (5) gets changed to the following set of coupled equations ska :
$`[i\mathrm{}{\displaystyle \frac{}{t}}{\displaystyle \frac{\mathrm{}^2_𝐫^2}{2m}}+V(𝐫)+g_{11}n_1+g_{12}n_2]\mathrm{\Psi }_1(𝐫,t)=0,`$
(6)
$`[i\mathrm{}{\displaystyle \frac{}{t}}{\displaystyle \frac{\mathrm{}^2_𝐫^2}{2m}}+V(𝐫)+g_{21}n_1+g_{22}n_2]\mathrm{\Psi }_2(𝐫,t)=0.`$
Here $`n_i|\mathrm{\Psi }_i(𝐫,t)|^2`$, $`g_{ij}=4\pi \mathrm{}^2a_{ij}N_i/m,`$ and $`i,j=1,2`$ represent the two types of bosons, and where $`N_1`$ is the number of boson 1 and $`N_2`$ that of boson 2, $`a_{ij}`$ is the scattering length for a boson of type $`i`$ and one of type $`j`$.
For the study of bright solitons we shall reduce Eqs. (Bright solitons in coupled defocusing NLS equation supported by coupling: Application to Bose-Einstein Condensation<sup>1</sup><sup>1</sup>1E-mail address: [email protected] (S.K. Adhikari).) and (Bright solitons in coupled defocusing NLS equation supported by coupling: Application to Bose-Einstein Condensation<sup>1</sup><sup>1</sup>1E-mail address: [email protected] (S.K. Adhikari).) to a minimal one-dimensional form in the cigar-shaped geometry with $`\nu =0`$. This is achieved by considering solutions of the type $`\mathrm{\Psi }_i(𝐫,t)=\varphi _i(z,t)\psi ^{(0)}(\rho )`$ where
$`|\psi ^{(0)}(\rho )|^2`$ $``$ $`{\displaystyle \frac{m\omega }{\pi \mathrm{}}}\mathrm{exp}\left({\displaystyle \frac{m\omega \rho ^2}{\mathrm{}}}\right).`$ (8)
The expression (8) corresponds to the ground state wave function in the absence of nonlinear interactions and satisfies
$`{\displaystyle \frac{\mathrm{}^2}{2m}}_\rho ^2\psi ^{(0)}(\rho )+{\displaystyle \frac{1}{2}}m\omega ^2\rho ^2\psi ^{(0)}(\rho )`$ $`=`$ $`\mathrm{}\omega \psi ^{(0)}(\rho ),`$ (9)
with normalization $`2\pi _0^{\mathrm{}}|\psi ^{(0)}(\rho )|^2\rho 𝑑\rho =1.`$ Now the dynamics is carried by $`\varphi _i(z,t)`$ and the radial dependence is frozen in the ground state $`\psi ^{(0)}(\rho )`$.
Averaging over the radial mode, i.e., multiplying Eqs. (Bright solitons in coupled defocusing NLS equation supported by coupling: Application to Bose-Einstein Condensation<sup>1</sup><sup>1</sup>1E-mail address: [email protected] (S.K. Adhikari).) and (Bright solitons in coupled defocusing NLS equation supported by coupling: Application to Bose-Einstein Condensation<sup>1</sup><sup>1</sup>1E-mail address: [email protected] (S.K. Adhikari).) by $`\psi ^{(0)}(\rho )`$ and integrating over $`\rho `$, we obtain the following one-dimensional equations abdul :
$`[i\mathrm{}{\displaystyle \frac{}{t}}+{\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \frac{^2}{z^2}}f_{11}|\varphi _1|^2f_{12}|\varphi _2|^2]\varphi _1(z,t)=0,`$ (10)
$`[i\mathrm{}{\displaystyle \frac{}{t}}+{\displaystyle \frac{\mathrm{}^2}{2m}}{\displaystyle \frac{^2}{z^2}}f_{21}|\varphi _2|^2f_{22}|\varphi _2|^2]\varphi _2(z,t)=0,`$ (11)
where
$`f_{ij}=2\pi g_{ij}{\displaystyle _0^{\mathrm{}}}|\psi ^{(0)}|^4\rho 𝑑\rho =g_{ij}{\displaystyle \frac{m\omega }{2\pi \mathrm{}}}.`$ (12)
In Eqs. (10) and (11) the normalization is given by $`_{\mathrm{}}^{\mathrm{}}|\varphi _i(z,t)|^2𝑑z=1`$.
For calculational purpose it is convenient to reduce the set (10) and (11) to dimensionless form by introducing convenient dimensionless variables. In Eqs. (10) and (11) we consider the dimensionless variables $`\tau =t\omega /2`$, $`y=z/l`$, $`\chi _i=\sqrt{l}\varphi _i`$, with $`l=\sqrt{\mathrm{}/(\omega m)}`$, so that we have the following coupled NLS equations
$`[i{\displaystyle \frac{}{\tau }}+{\displaystyle \frac{d^2}{dy^2}}n_{11}|\chi _1|^2n_{12}\left|\chi _2|^2\right]\chi _1(y,\tau )=0,`$ (13)
$`[i{\displaystyle \frac{}{\tau }}+{\displaystyle \frac{d^2}{dy^2}}n_{21}|\chi _1|^2n_{22}\left|\chi _2|^2\right]\chi _2(y,\tau )=0,`$ (14)
where $`n_{ij}=4a_{ij}N_i/l,i,j=1,2.`$ In Eqs. (13) and (14), the normalization condition is given by $`_{\mathrm{}}^{\mathrm{}}|\chi _i(y,\tau )|^2𝑑y=1.`$
Equations (13) and (14) represent the one-dimensional limit of the three-dimensional equation. For solitons we finally have to take the nonlinearity coefficients $`n_{ij}(ij)`$ to be negative corresponding to attraction. These equations have analytic solutions only under special conditions. First, when $`n_{12}=n_{21}=0`$, they become two uncoupled NLS equations which allow the following trivial soliton solutions for negative $`n_{ii}`$ 1 :
$`\chi _i(y,\tau )=\sqrt{(2w_i/|n_{ii}|)}\text{sech}(y\sqrt{w_i})\mathrm{exp}(i\tau ),`$ (15)
$`i=1,2.`$ To satisfy the normalization condition one should have $`w_i=n_{ii}^2/16`$. In Eq. (15) and below the parameter $`w`$ is to be adjusted so as to satisfy the normalization condition. When all the nonlinear interactions are attractive (negative) and $`n_{11}=n_{22}=n`$ and $`n_{12}=n_{21}=\alpha n`$ one has the solutions 1
$`\chi _i(y,\tau )=\sqrt{[2w/(n\alpha +n)]}\text{sech}(y\sqrt{w})\mathrm{exp}(i\tau ),`$ (16)
$`i=1,2`$, with $`w=(n\alpha +n)^2/16`$. Also, when $`n_{ii}=0`$, they have the following soliton solutions for negative nonlinearities
$`\chi _i(y,\tau )=\sqrt{(2w_i/|n_{ji}|)}\text{sech}(y\sqrt{w_i})\mathrm{exp}(i\tau ),`$ (17)
$`ij,`$ $`i,j=1,2`$ with $`w_i=n_{ji}^2/16`$. This case has the possibility of forming the soliton due to interspecies interaction as the intra-species interaction is zero.
It is also possible to have intra-species repulsion (positive or defocusing $`n_{ii}`$) and interspecies attraction (negative or focusing $`n_{12}`$ and $`n_{21}`$) in order to have the following soliton solutions of Eqs. (13) and (14) when $`n_{11}=n_{22}=n`$ and $`n_{12}=n_{21}=\alpha n,(\alpha >1):`$
$`\chi _i(y,\tau )=\sqrt{[2w/(n\alpha n)]}\text{sech}(y\sqrt{w})\mathrm{exp}(i\tau ),`$ (18)
$`i=1,2`$ with $`w=(n\alpha n)^2/16`$. In this case due to strong interspecies attraction solitons are formed despite of intra-species repulsion. In all these cases the functional dependence of the two analytic solutions $`\chi _i(y,\tau )`$ on $`y`$ are the same. However, there are interesting numerical solutions to Eqs. (13) and (14) where the functional dependence of $`\chi _i(y,\tau ),i=1,2,`$ could be different, which we study next. Such vector solitons will have different widths and peak powers.
We solve the coupled mean-field-hydrodynamic equations (13) and (14) for bright solitons numerically using a time-iteration method based on the Crank-Nicholson discretization scheme elaborated in Ref. sk1 . We discretize the mean-field-hydrodynamic equation using time step $`0.0005`$ and space step $`0.025`$.
We performed the time evolution of the set of equations (13) and (14) introducing harmonic oscillator potentials $`y^2`$ in these equations and starting with the eigenfunction of the linear harmonic oscillator problem with the nonlinear terms set equal to zero: $`\chi _1(y,\tau )=\chi _2(y,\tau )=\pi ^{1/4}\mathrm{exp}(y^2/2)\mathrm{exp}(i\tau )`$. During the course of time evolution the nonlinear terms are switched on very slowly and resultant solution iterated until convergence was obtained. Then the time evolution is continued and the harmonic oscillator potential terms ($`y^2`$) are slowly switched off and the resultant solution iterated 100 000 times for convergence. If converged solutions are obtained, they correspond to the required bright solutions. In the numerical investigation we take $`\omega =2\pi \times 100`$ Hz, and $`m_B`$ as the mass of <sup>87</sup>Rb. Consequently, the unit of length $`l1`$ $`\mu `$m and unit of time $`2/\omega 3`$ ms.
First we solve Eqs. (13) and (14) with $`N_1=N_2=7500`$, $`a_{11}=1`$ nm, $`a_{12}=1`$ nm and $`a_{22}=0.1`$ nm. With these parameters the nonlinearities in Eqs. (13) and (14) are $`n_{11}=30`$, $`n_{12}=30`$, $`n_{21}=30`$, and $`n_{22}=3`$.
The converged bright solitons are plotted in Fig. 1. In this case the function $`\varphi _1`$ of the first component extends over a longer region in space compared to the function $`\varphi _2`$ of the second component. It is possible to have solitons with different extensions in space by varying the parameters of the system. The scattering length can be manipulated in the bosonic systems near a Feshbach resonance fs by varying a background magnetic field. By varying the scattering length and the number of atoms we could arrive at different values of nonlinearity parameters from those in Fig. 1 and thus have different extensions of the condensates in space.
In the situation presented in Fig. 1, in Eq. (13) governing the dynamics of boson 1 the two nonlinearities are $`n_{11}=30`$ and $`n_{12}=30`$, which may superficially indicate an effective nonlinearity of 0. However, the effective nonlinearity in this equation is $`n_{11}|\chi _1|^2n_{12}|\chi _2|^2`$. Due to a more strongly bound soliton of boson 2 this effective nonlinearity could become attractive and bind the soliton of type 1.
Next we consider a two-component soliton formed in the optical-lattice potential $`V(z)=V_0\mathrm{sin}^2(4z)`$ introduced in Eqs. (13) and (14). In our numerical calculation we take $`V_0=100`$, $`n_{11}=35`$, $`n_{12}=n_{21}=30`$, and $`n_{22}=3`$. In this case in the wave function of boson 1 prominent wiggles are formed due to the optical-lattice potential. As the spacing of the optical-lattice sites are relatively large no wiggles are formed in the more localized wave function of boson 2. Equation (13) taken separately has nonlinearities $`n_{11}=35`$ and $`n_{12}=30`$ corresponding to an apparent overall repulsion. However, the strongly bound soliton of boson 2 and the interspecies attraction as well as the periodic optical-lattice potential aid in the formation of the soliton of boson 1. In this connection it should be noted that the optical-lattice potential aids in binding the soliton in the presence of nonlinear interspecies attraction. The optical-lattice potential alone cannot bind a soliton in the absence of nonlinear attraction both in one- and two-component BECs.
Finally, we consider the stability of these solitons under a small perturbation. For this purpose we consider the solitons of Eqs. (13) and (14) formed for nonlinearities $`n_{11}=30`$, $`n_{12}=n_{21}=51`$, and $`n_{22}=3`$. The wave functions $`\varphi _i(z,t)`$ of these solitons are shown in Figs. 3 (a) and (b) for $`i=1,2`$, respectively. To test their stability under small perturbation, after their formation, the nonlinearities $`n_{12}`$ and $`n_{21}`$ are suddenly changed to $`45`$ at $`t=100`$ ms so that the solitons are set into motion. Such change in the nonlinearity can be achieved by a jump in the scattering length by manipulating a background magnetic field near a Feshbach resonance fs . The solitons are found to execute stable non-periodic breathing oscillations. The stability of these solitons after the perturbation is applied is demonstrated in Figs. 3.
To further illustrate the stability of the oscillation of the solitons of Figs. 3 we plot in Fig. 4 the root mean square (rms) size $`z_{rms}`$ vs. time. Sustained oscillation for a very long time illustrates the stability of the solitons. From Fig. 4 we find that for $`t<100`$ ms the rms sizes of the two solitons are constant. However, after the application of the repulsive impulsive force at $`t=100`$ ms by reducing the attractive nonlinearity, the rms sizes suddenly jump to a larger value and execute stable oscillatory dynamics. This stable oscillation guaratees the stationary nature of the solitons under small perturbation.
The present investigation has consequences in generating bright vector solitons in directional couplers of nonlinear fiber optics agra . Vector solitons are solutions of coupled one-dimensional NLS equations with the property that the orthogonally polarized components propagate in a birefringent fiber without change in shape. In vector solitons an input pulse maintains not only its intensity profile but also its state of polarization even when it is not launched along one of the principal axes of the fiber. We have investigated the possibility that two orthogonally polarized pulses of different widths and different peak powers propagate undistorted in birefringent fibers. The present investigation suggests that it is possible to have bright vector solitons in the coupled NLS equation with diagonal defocusing (repulsive) nonlinearity and off-diagonal focusing (attractive) nonlinearity, which should be of interest in nonlinear fiber optics.
We use a coupled set of time-dependent mean-field GP equations for a two-component repulsive BEC to demonstrate the formation of bright solitons due to interspecies attraction. The interspecies attraction can neutralize the intra-species repulsion and induce an effective attraction in the mean-field GP equations responsible for the formation of bright solitons. An attractive interspecies interaction is necessary for the formation of the bright solitons as the diagonal nonlinearity $`n_{ii}`$ in the mean-field equations is taken to be repulsive (positive). In mean-field equations (13) and (14) $`n_{ii}`$s are positive (repulsive) and $`n_{ij}`$s are negative (attractive), so that the overall contribution of the nonlinear terms in these equations become attractive to support the bright solitons. We have also established the formation of these solitons in the presence of a periodic optical-lattice potential in an entirely different shape and trapping condition from a conventional soliton.
In view of the present study the appearance of bright solitons in multi-component repulsive BECs seems possible in quasi-one-dimensional geometry. Bright solitons have been created experimentally in attractive BECs in three dimensions in the presence of radial trapping only without any axial trapping sol . Also, there have been experimental studies of multicomponent BECs exp1 ; exp4 . Hence, bright solitons can be created and studied in the laboratory in the presence of radial trapping only in a two-component repulsive BEC supported by interspecies attraction and the prediction of the present study verified. We have also suggested the possibility of the formation of similar coupled solitons in fiber optics in one dimension.
Note added in proof: After the completion of this investigation we have known about another similar study sim .
###### Acknowledgements.
The work is supported in part by the CNPq of Brazil.
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# Holographic renormalization of cascading gauge theories
## 1 Introduction
The traditional definition of quantum field theories starts from a fixed point of the renormalization group (free or interacting) at high energies, which is some local conformal field theory, and defines the theory by a renormalization group flow starting from that fixed point. One of the interesting “side-effects” of the progress in our understanding of string theory in the last few years is the realization that there exist consistent quantum field theories which cannot be defined in this way. In some cases these theories may be described by a decoupling limit of some sector of string theory, and more generally they can always be defined by a background of type IIB string theory which is holographically dual to these theories, in the same sense that the $`AdS_5\times S^5`$ background of string theory is dual to the $`d=4`$ $`𝒩=4`$ supersymmetric Yang-Mills theory .
The new types of theories which were discovered seem to fall into two classes. One class of theories is “little string theories” – non-local theories which share some features with string theories (though they do not include gravity). In this paper we focus on the second class, of “cascading theories” <sup>1</sup><sup>1</sup>1An interesting relation between these two types of theories was recently discovered in .. The prototypical example of such theories (which we will focus on here, though we expect to be able to generalize our methods also to other cascading theories, including theories in other dimensions) is the theory related to fractional $`3`$-branes at a conifold singularity, first studied in (see for reviews). These theories do not have a direct definition in field theory terms (since they do not seem to have a UV fixed point), so their only known direct definition is via the holographic duality; in this paper we will attempt to understand this definition better and verify that it can be used for well-defined computations in these theories. When one introduces a finite high-energy cutoff, these theories at the cutoff scale resemble $`𝒩=1`$ supersymmetric $`SU(K)\times SU(K+P)`$ gauge theories with two bifundamental and two anti-bifundamental chiral superfields and some superpotential; when one flows down in energy from this cutoff one of the gauge theories becomes strongly coupled and the theory seems to undergo a series of Seiberg duality “cascades”, ending with a confining theory in the IR (which is in the same universality class as the $`𝒩=1`$ supersymmetric pure $`SU(P)`$ Yang-Mills theory when $`K`$ is a multiple of $`P`$). However, these gauge theories are never asymptotically free, so they cannot be used to define the theory – rather it seems that more and more degrees of freedom are needed to define the theory at higher energies, and that the ultimate definition of the cascading theories requires a theory with an infinite number of fields. It may be possible to define the cascading gauge theories by a limiting procedure, starting from a theory with a finite number of degrees of freedom which can flow to the cascade (as in ) and taking the limit in which the number of degrees of freedom goes to infinity – it would be interesting to make such a definition precise. In this paper we will not discuss the interpretation of the cascade as a gauge theory, but rather we will take the holographic dual to define the cascading theory, and attempt to see if such a definition makes sense.
In order to use a holographic dual as a definition of a field theory we need to have a prescription for the computation of all correlation functions in the field theory. For the AdS/CFT correspondence such a prescription was given in , and it can be generalized to other holographic dualities as well. In principle one should be able to perform the computations of correlation functions directly in string theory on the holographic dual background, but in practice, since this string theory is very complicated, computations can only be done in a low-energy gravitational approximation, and we will use this approximation in this paper. In this approximation the correlation functions may be defined as derivatives of the (super)gravity action on the holographic dual background with respect to sources at the boundary of space-time. When one tries to perform such computations one encounters divergences, which from the gravity point of view are IR divergences related to the infinite distance to the boundary. These divergences may be dealt with by the process of “holographic renormalization” . First, one regularizes the theory by imposing a cutoff on the radial direction (the precise meaning of this cutoff in the field theory language is not clear, but it certainly provides a good regularization). Next, one adds counter-terms to the gravitational action, which are local functions of the fields at the cutoff, in such a way that the action remains finite when the cutoff is taken to infinity. This process is very analogous to renormalization in field theory, and it seems that it should be mapped to this by the holographic duality; in both cases in a well-defined theory there is just a finite number of divergences which need to be canceled, after which correlation functions may be computed, depending on a finite number of parameters (some of which are coupling constants of the theory, while others are related to vacuum expectation values of fields or to ambiguities in the definitions of operators).
The process of holographic renormalization is by now conceptually well-understood in asymptotically anti-de Sitter spaces, where several examples have been analyzed in detail , though the general renormalization for theories with many fields in the bulk is quite complicated and has not yet been performed. In principle one would expect a similar renormalization process to apply in other holographic dualities such as those of cascading gauge theories<sup>2</sup><sup>2</sup>2 Some general properties of the stress tensor in holographic renormalization, which apply also to cascading backgrounds, were derived in . It is not obvious whether holographic renormalization should be possible also in the duals to non-local theories such as “little string theories”., and our goal in this paper is to understand how this works. One interesting question which immediately arises is the following. In standard field theories there are some correlation functions, such as the one-point function of the trace of the stress-energy tensor (related to the conformal anomaly), which are proportional to the number of degrees of freedom in the theory. As discussed above, cascading theories do not seem to have a finite number of degrees of freedom, but rather more and more degrees of freedom as one goes to shorter and shorter distance scales. So, should the correlation functions of cascading theories be finite as one takes the cutoff to infinity or should some of them diverge ?
We have attempted to perform a holographic renormalization of the cascading theories both under the assumption that all correlation functions must be finite, and under the assumption that correlation functions are allowed to diverge as the cutoff is taken to infinity, in a way which depends on the effective number of degrees of freedom. Somewhat surprisingly, we found that it is possible to renormalize the theory with finite correlation functions, but we were not able to renormalize the theory using the other assumption. Thus, we claim that cascading gauge theories should be renormalized just like standard theories, with all correlation functions finite. At first sight this seems to contradict the fact that these theories have an infinite number of high-energy degrees of freedom. We claim that this is not the case, and that these theories have an infinite number of high-energy degrees of freedom even though all correlation functions (including the conformal anomaly) are finite. We illustrate this by analyzing the thermodynamics of the cascading theories, showing that the effective number of degrees of freedom diverges at high temperatures even though all correlation functions are finite (at any fixed temperature).
At this point we should describe the precise assumptions under which we perform the holographic renormalization of the cascading gauge theories. Usually holographic renormalization is performed using a consistent truncation of the theory to a small number of fields <sup>3</sup><sup>3</sup>3As far as we know no more than two fields (coupled to gravity) were analyzed until now., for which any sources are allowed. In our case we also truncate the full spectrum of fields in the holographic dual background to a finite number of fields. One important truncation we make is that we only include fields which preserve the $`SU(2)\times SU(2)\times _{2P}\times _2`$ isometry which is present at high energies in the cascading background of , related to the global symmetry in the field theory (the metric actually has a $`U(1)`$ isometry but this is broken to a $`_{2P}`$ subgroup by the fluxes ). In infinite space the cascading theories spontaneously break the $`_{2P}`$ symmetry, as found in , so our analysis cannot be directly used to analyze the infinite volume cascading theories. However, at finite (high enough) temperature or at finite (small enough) volume (for instance on $`S^3\times `$, $`S^4`$ or $`dS_4`$) this symmetry is expected to be preserved , so our analysis may be directly used for such backgrounds. We expect that it should be straightforward (though technically difficult) to extend our analysis to include also fields which break the $`SU(2)\times SU(2)\times _{2P}\times _2`$ symmetry.
Even with the truncation to the $`(SU(2)\times SU(2)\times _{2P}\times _2)`$-invariant sector, we are left with a large number of five dimensional fields in the bulk – the metric and four scalar fields. These fields all mix together so we were not able to truncate the theory further. Moreover, some of these fields are dual to irrelevant operators <sup>4</sup><sup>4</sup>4Since the cascading theories are close to conformal field theories at high energies, with the characteristic power law behavior of conformal field theories replaced by powers multiplying logs, we will use the standard terminology of conformal field theories., so it is not known how to introduce arbitrary sources for these fields, as is usually done in the holographic renormalization process in order to systematically compute the counter-terms. Thus, we do not consider the most general sources; we allow a generic source for the five dimensional metric (this source is identified with the four dimensional metric of the space-time on which the cascading theory lives), but only constant sources for the other scalar fields (and, in particular, vanishing sources for the two scalar fields corresponding to irrelevant operators). This simplifies the analysis considerably, but there are three disadvantages. First, usually in holographic renormalization the finiteness of the action (for arbitrary sources) guarantees the finiteness of all correlation functions, but we do not allow arbitrary sources so we have to separately check that the correlation functions are finite in addition to the finiteness of the action (we check this only for one-point functions; additional counter-terms may be needed to ensure the finiteness of all correlation functions). Second, we can no longer translate the divergent terms in the action directly to counter-terms, as usually done in holographic renormalization. Therefore, we are forced to use a different procedure, of guessing the counter-terms and verifying that they lead to finite correlation functions (with a finite number of ambiguities). Again, we expect that it should be possible to generalize our analysis to include arbitrary sources (at least for all the marginal and relevant operators), though this will be technically complicated. We hope to return to this problem in the future. The counter-terms that we find in this method are far from being unique, and are certainly not the precise counter-terms that lead to finiteness of all correlation functions. However, we expect that the difference between the counter-terms we find and the correct counter-terms will not affect the unambiguous results which we obtain. Third, since we do not have arbitrary sources we cannot compute arbitrary correlation functions, but only the derivatives of the action with respect to the sources we include. Thus, our procedure allows us to compute any correlation functions of the stress-energy tensor (dual to the bulk metric), but in the scalar sector we can only compute one-point functions.
In this paper we show that, in the truncation described above, it is possible to holographically renormalize the cascading gauge theory background and to obtain finite one-point functions. We begin in section 2 by describing the background, the ansatz we use for the solutions with the sources described above, and the solutions we find. In section 3 we describe in detail the holographic renormalization process and the form of the counter-terms we find. In section 4 we discuss the example of cascading theories at finite (high) temperature, following , and compute their thermodynamic properties. We end in section 5 with our conclusions and a discussion of future directions. Various technical results are relegated to the appendix.
## 2 The action and asymptotic behavior of cascading backgrounds
In this section we construct the asymptotic (near the boundary) solutions corresponding to cascading gauge theories compactified on arbitrary manifolds, generalizing the flat space asymptotic solution found in .
### 2.1 The gravitational action and its KK reduction and truncation
We will work in the gravitational approximation to type IIB string theory, using the type IIB supergravity action. This action takes the form (in the Einstein frame)
$$\begin{array}{cc}\hfill S_{10}=\frac{1}{16\pi G_{10}}_{_{10}}& (R_{10}1\frac{1}{2}d\mathrm{\Phi }d\mathrm{\Phi }\frac{1}{2}e^\mathrm{\Phi }H_3H_3\frac{1}{2}e^\mathrm{\Phi }F_3F_3\hfill \\ & \frac{1}{4}F_5F_5\frac{1}{2}C_4H_3F_3),\hfill \end{array}$$
(2.1)
where $`_{10}`$ is the ten dimensional bulk space-time, $`G_{10}`$ is the ten dimensional gravitational constant, and we have consistently set the axion $`C_0`$ to zero (it vanishes in all the solutions we are interested in). In this action
$$F_3=dC_2,F_5=dC_4C_2H_3,$$
(2.2)
where $`C_2`$ and $`C_4`$ are the Ramond-Ramond (RR) potentials. The equations of motion following from the action (2.1) have to be supplemented by the self-duality condition
$$F_5=F_5.$$
(2.3)
It is important to remember that the self-duality condition (2.3) can not be imposed at the level of the action, as this would lead to wrong equations of motion.
Next, we perform a Kaluza-Klein (KK) reduction of this action to five dimensions, using a specific ansatz for the metric and for the various forms. This ansatz includes in particular the solution of , and it is the most general ansatz describing a deformation of this solution which preserves the $`SU(2)\times SU(2)\times _{2P}\times _2`$ symmetry of this solution<sup>5</sup><sup>5</sup>5Except for two modes of the RR fields which we consistently set to zero as mentioned in the text. (the discrete $`_2`$ symmetry acts by exchanging the two global $`SU(2)`$ factors). We take $`_{10}`$ to be a direct warped product of $`_5`$ with metric $`g_{\mu \nu }(y)`$ and the ‘squashed’ $`T^{1,1}`$ coset appearing in the solution of . So, the Einstein-frame metric ansatz is
$$ds_{10}^2=g_{\mu \nu }(y)dy^\mu dy^\nu +\mathrm{\Omega }_1^2(y)e_\psi ^2+\mathrm{\Omega }_2^2(y)\underset{a=1}{\overset{2}{}}\left(e_{\theta _a}^2+e_{\varphi _a}^2\right),$$
(2.4)
where $`y`$ denotes the coordinates of $`_5`$ (greek indices $`\mu ,\nu `$ will run from $`0`$ to $`4`$) and the one-forms $`e_\psi ,e_{\theta _a},e_{\varphi _a}`$ ($`a=1,2`$) are given by (see also ) :
$$e_\psi =\frac{1}{3}\left(d\psi +\underset{a=1}{\overset{2}{}}\mathrm{cos}\theta _ad\varphi _a\right),e_{\theta _a}=\frac{1}{\sqrt{6}}d\theta _a,e_{\varphi _a}=\frac{1}{\sqrt{6}}\mathrm{sin}\theta _ad\varphi _a.$$
(2.5)
Additionally, we assume the following ansatz for the fluxes $`H_3dB_2`$, $`F_3`$ and the dilaton $`\mathrm{\Phi }`$ :
$$F_3=Pe_\psi \left(e_{\theta _1}e_{\varphi _1}e_{\theta _2}e_{\varphi _2}\right),B_2=\stackrel{~}{k}(y)\left(e_{\theta _1}e_{\varphi _1}e_{\theta _2}e_{\varphi _2}\right),\mathrm{\Phi }=\mathrm{\Phi }(y),$$
(2.6)
where $`P`$ is an integer corresponding to the RR 3-form flux on the compact 3-cycle (and to the number of fractional branes on the conifold). Special care should be taken with the RR 5-form. From (2.2) we get the Bianchi identity
$$dF_5=F_3H_3,$$
(2.7)
which for the background fluxes (2.6) is solved by
$$F_5=dC_4\left(\stackrel{~}{K}_0+2P\stackrel{~}{k}(y)\right)e_\psi e_{\theta _1}e_{\varphi _1}e_{\theta _2}e_{\varphi _2}$$
(2.8)
with some constant $`\stackrel{~}{K}_0`$. In our ansatz the RR four-form does not depend on the compact coordinates, that is $`C_4C_4(y)`$ (note that $`C_4F_3H_30`$), and the RR five-form is proportional to the volume form of $`_5`$ (plus its dual). We define $`K(y)`$ by
$$dC_4=\frac{K(y)}{\mathrm{\Omega }_1\mathrm{\Omega }_2^4}\mathrm{vol}__5\frac{K(y)}{\mathrm{\Omega }_1\mathrm{\Omega }_2^4}\sqrt{det(g_{\mu \nu })}dy^1\mathrm{}dy^5,$$
(2.9)
and then the self-duality condition (2.3) implies
$$K(y)=\stackrel{~}{K}_0+2P\stackrel{~}{k}(y)$$
(2.10)
(again, in deriving the effective action we should keep $`C_4`$ unconstrained and impose this equation later). Altogether, from the five-dimensional perspective we allow fluctuations in the metric $`g_{\mu \nu }(y)`$, in the scalar fields $`\mathrm{\Omega }_1(y),\mathrm{\Omega }_2(y),\stackrel{~}{k}(y),\mathrm{\Phi }(y)`$ and in the four-form $`C_4(y)`$ (which is determined in terms of the others by the self-duality condition). We have set to zero various fluctuations of the form fields which are $`p`$-forms on $`_5`$, and also fluctuations of $`C_2`$ of the same form as the fluctuation of $`B_2`$ in (2.6), even though they are allowed by the symmetries. This is a consistent truncation of the full ten dimensional supergravity action.
We now perform the KK reduction of (2.1) by plugging into it the ansatz described above. Recall that
$$\mathrm{vol}_{T^{1,1}}e_\psi e_{\theta _1}e_{\varphi _1}e_{\theta _2}e_{\varphi _2}=\frac{16\pi ^3}{27}.$$
(2.11)
First, we have
$$_{_{10}}11=\mathrm{vol}_{T^{1,1}}__5\mathrm{\Omega }_1\mathrm{\Omega }_2^4\mathrm{vol}__5.$$
(2.12)
With a straightforward but somewhat tedious computation we find that in the background (2.4)
$$\begin{array}{cc}\hfill R_{10}=R_5& 2\mathrm{\Omega }_1^1g^{\lambda \nu }\left(_\lambda _\nu \mathrm{\Omega }_1\right)8\mathrm{\Omega }_2^1g^{\lambda \nu }\left(_\lambda _\nu \mathrm{\Omega }_2\right)\hfill \\ & 4g^{\lambda \nu }\left(2\mathrm{\Omega }_1^1\mathrm{\Omega }_2^1_\lambda \mathrm{\Omega }_1_\nu \mathrm{\Omega }_2+3\mathrm{\Omega }_2^2_\lambda \mathrm{\Omega }_2_\nu \mathrm{\Omega }_2\right)\hfill \\ & +24\mathrm{\Omega }_2^24\mathrm{\Omega }_1^2\mathrm{\Omega }_2^4,\hfill \end{array}$$
(2.13)
where $`R_5`$ is the five dimensional Ricci scalar of the metric
$$ds_5^2=g_{\mu \nu }(y)dy^\mu dy^\nu .$$
(2.14)
In (2.13), $`_\lambda `$ denotes the covariant derivative with respect to the metric (2.14), explicitly given by
$$\begin{array}{cc}\hfill _\lambda \mathrm{\Omega }_i& =_\lambda \mathrm{\Omega }_i,\hfill \\ \hfill _\lambda _\nu \mathrm{\Omega }_i& =_\lambda _\nu \mathrm{\Omega }_i\mathrm{\Gamma }_{\lambda \nu }^\rho _\rho \mathrm{\Omega }_i.\hfill \end{array}$$
(2.15)
Now, by plugging our ansatz into (2.1) we find that it reduces to the following effective action :
$$\begin{array}{cc}\hfill S_5=\frac{1}{16\pi G_5}__5\mathrm{vol}__5\{& \mathrm{\Omega }_1\mathrm{\Omega }_2^4\left(R_{10}\frac{1}{2}_\mu \mathrm{\Phi }^\mu \mathrm{\Phi }\right)\mathrm{\Omega }_1e^\mathrm{\Phi }\left(_\mu \stackrel{~}{k}^\mu \stackrel{~}{k}+\frac{P^2e^{2\mathrm{\Phi }}}{\mathrm{\Omega }_1^2}\right)\hfill \\ & \frac{1}{4}(\frac{(\stackrel{~}{K}_0+2P\stackrel{~}{k})^2}{\mathrm{\Omega }_1\mathrm{\Omega }_2^4}+\frac{5}{24}\mathrm{\Omega }_1\mathrm{\Omega }_2^4_{\mu _1\mathrm{}\mu _5}^{\mu _1\mathrm{}\mu _5})\}\hfill \\ \hfill +\frac{1}{16\pi G_5}P__5& d\stackrel{~}{k}C_4,\hfill \end{array}$$
(2.16)
where
$$_{\mu _1\mathrm{}\mu _5}_{[\mu _1}C_4{}_{\mu _2\mathrm{}\mu _5]}{}^{}=\frac{1}{5}\frac{K}{\mathrm{\Omega }_1\mathrm{\Omega }_2^4}\sqrt{det(g_{\mu \nu })}ϵ_{\mu _1\mathrm{}\mu _5}$$
(2.17)
($`[\mathrm{}]`$ denotes anti-symmetrization with weight one) and $`G_5`$ is the five dimensional effective gravitational constant
$$G_5\frac{G_{10}}{\mathrm{vol}_{T^{1,1}}}.$$
(2.18)
Note that our gravitational action is not the standard five dimensional action because of the factor of $`\mathrm{\Omega }_1\mathrm{\Omega }_2^4`$ in front of the five dimensional Einstein-Hilbert term.
In the five dimensional action it turns out to be possible to “integrate out” the field $`C_4`$ using the self-duality equation (2.10) and to obtain an action involving only the other fields. This leads to the action we will be using in this paper
$$\begin{array}{cc}\hfill S_5=\frac{1}{16\pi G_5}__5\mathrm{vol}__5\{& \mathrm{\Omega }_1\mathrm{\Omega }_2^4\left(R_{10}\frac{1}{2}_\mu \mathrm{\Phi }^\mu \mathrm{\Phi }\right)P^2\mathrm{\Omega }_1e^\mathrm{\Phi }\left(\frac{_\mu K^\mu K}{4P^4}+\frac{e^{2\mathrm{\Phi }}}{\mathrm{\Omega }_1^2}\right)\hfill \\ & \frac{1}{2}\frac{K^2}{\mathrm{\Omega }_1\mathrm{\Omega }_2^4}\},\hfill \end{array}$$
(2.19)
where $`R_{10}`$ is given by (2.13) and $`K(y)`$ is related to $`\stackrel{~}{k}(y)`$ by (2.10).
### 2.2 The equations of motion and the ansatz for the solution
From the effective action (2.19) we obtain the following equations of motion :
$$0=\frac{1}{\sqrt{g}}_\mu \left[\frac{e^\mathrm{\Phi }\mathrm{\Omega }_1}{2P^2}\sqrt{g}g^{\mu \nu }_\nu K\right]\frac{K}{\mathrm{\Omega }_1\mathrm{\Omega }_2^4},$$
(2.20)
$$0=\frac{1}{\sqrt{g}}_\mu \left[\mathrm{\Omega }_1\mathrm{\Omega }_2^4\sqrt{g}g^{\mu \nu }_\nu \mathrm{\Phi }\right]+\frac{\mathrm{\Omega }_1e^\mathrm{\Phi }(K)^2}{4P^2}\frac{P^2e^\mathrm{\Phi }}{\mathrm{\Omega }_1},$$
(2.21)
$$\begin{array}{cc}\hfill 0=& \mathrm{\Omega }_2^4R_512\mathrm{\Omega }_2^2(\mathrm{\Omega }_2)^2+24\mathrm{\Omega }_2^212\mathrm{\Omega }_1^28\mathrm{\Omega }_2^3\underset{5}{\text{ }\text{ }\text{ }}\mathrm{\Omega }_2\hfill \\ & \frac{1}{2}\mathrm{\Omega }_2^4(\mathrm{\Phi })^2+\frac{P^2e^\mathrm{\Phi }}{\mathrm{\Omega }_1^2}\frac{e^\mathrm{\Phi }(K)^2}{4P^2}+\frac{K^2}{2\mathrm{\Omega }_1^2\mathrm{\Omega }_2^4},\hfill \end{array}$$
(2.22)
$$\begin{array}{cc}\hfill 0=& 4\mathrm{\Omega }_1\mathrm{\Omega }_2^3R_58\mathrm{\Omega }_2^3\underset{5}{\text{ }\text{ }\text{ }}\mathrm{\Omega }_124\mathrm{\Omega }_1\mathrm{\Omega }_2^2\underset{5}{\text{ }\text{ }\text{ }}\mathrm{\Omega }_224\mathrm{\Omega }_2^2\mathrm{\Omega }_1\mathrm{\Omega }_224\mathrm{\Omega }_1\mathrm{\Omega }_2(\mathrm{\Omega }_2)^2+48\mathrm{\Omega }_1\mathrm{\Omega }_2\hfill \\ & 2\mathrm{\Omega }_1\mathrm{\Omega }_2^3(\mathrm{\Phi })^2+\frac{2K^2}{\mathrm{\Omega }_1\mathrm{\Omega }_2^5},\hfill \end{array}$$
(2.23)
$$\begin{array}{cc}\hfill \mathrm{\Omega }_1\mathrm{\Omega }_2^4R_{5\mu \nu }=& \frac{g_{\mu \nu }}{3}\left\{\frac{P^2e^\mathrm{\Phi }}{\mathrm{\Omega }_1}+\frac{K^2}{2\mathrm{\Omega }_1\mathrm{\Omega }_2^4}+\underset{5}{\text{ }\text{ }\text{ }}\left(\mathrm{\Omega }_1\mathrm{\Omega }_2^4\right)24\mathrm{\Omega }_1\mathrm{\Omega }_2^2+4\mathrm{\Omega }_1^3\right\}\hfill \\ & +_\mu _\nu \left(\mathrm{\Omega }_1\mathrm{\Omega }_2^4\right)4\mathrm{\Omega }_2^3\left(_\mu \mathrm{\Omega }_1_\nu \mathrm{\Omega }_2+_\nu \mathrm{\Omega }_1_\mu \mathrm{\Omega }_2\right)12\mathrm{\Omega }_1\mathrm{\Omega }_2^2_\mu \mathrm{\Omega }_2_\nu \mathrm{\Omega }_2\hfill \\ & +\frac{\mathrm{\Omega }_1e^\mathrm{\Phi }}{4P^2}_\mu K_\nu K+\frac{1}{2}\mathrm{\Omega }_1\mathrm{\Omega }_2^4_\mu \mathrm{\Phi }_\nu \mathrm{\Phi },\hfill \end{array}$$
(2.24)
where $`(F)^2`$ denotes $`g^{\mu \nu }_\mu F_\nu F`$ and $`\text{ }\text{ }\text{ }\text{ }_5`$ is the Laplacian in the metric (2.14).
In order to proceed further we make a convenient gauge choice for the five dimensional metric which separates the radial direction, which we will call $`\rho `$, from the four space-time dimensions of the cascading theory which we will denote by $`x^i`$ :
$$ds_5^2=h^{1/2}(x,\rho )\rho ^2\left(G_{ij}(x,\rho )dx^idx^j\right)+h^{1/2}(x,\rho )\rho ^2(d\rho )^2.$$
(2.25)
The boundary of the space will be taken to be at $`\rho 0`$. In this gauge choice some of the off-diagonal components of the metric vanish, partly fixing the diffeomorphism symmetry. We also define new scalar fields $`f_2`$ and $`f_3`$ related to the $`\mathrm{\Omega }_i`$ fields by
$$\mathrm{\Omega }_1^2=h^{1/2}f_2,\mathrm{\Omega }_2^2=h^{1/2}f_3.$$
(2.26)
The motivation for this parameterization is that in the solution of the function $`h`$ diverges logarithmically near the boundary $`\rho 0`$, but $`G_{ij}`$, $`f_2`$ and $`f_3`$ approach constant values. It is not difficult to rewrite the equations of motion (2.20)-(2.24) using the new variables $`G_{ij}`$, $`h`$, $`f_2`$, $`f_3`$, $`K`$ and $`\mathrm{\Phi }`$, and using the parameterization (2.25) of the metric, and we present the results in appendix A.1.
We now wish to find the solutions for the cascading theories on arbitrary space-time manifolds, in an expansion near the boundary at $`\rho 0`$. In the case of asymptotically anti-de Sitter spaces, fields in the bulk are dual to operators in the field theory of some dimension $`\mathrm{\Delta }`$, and they may be expanded in a power series in the radial $`\rho `$ coordinate, with a leading term of order $`\rho ^{4\mathrm{\Delta }}`$ corresponding to the source of the operator, multiplied by a power series in $`\rho ^2`$, and then a subleading term of order $`\rho ^\mathrm{\Delta }`$ corresponding to the one-point function of the operator (again multiplied by a power series in $`\rho ^2`$)<sup>6</sup><sup>6</sup>6 For integer values of $`\mathrm{\Delta }`$ there are also some logarithmic terms.. In the cascading gauge theory we expect a similar picture to arise, but with logarithmic corrections to all the terms corresponding to the logarithmic deviation from conformal invariance, and we will see below that this is indeed the case.
Let us first analyze the dimensions of the fields in the action (2.19) in the conformal case of $`P=0`$ ; note that naively the action (2.19) we wrote is singular as $`P0`$, but if we change variables from $`K(y)`$ to $`\stackrel{~}{k}(y)`$ (using (2.10)) the action becomes non-singular, so our analysis everywhere in this paper should be valid also in the $`P=0`$ case. Obviously, the metric $`G_{ij}`$ is dual to a dimension four operator which is just the stress-energy tensor. The dilaton $`\mathrm{\Phi }`$ and $`\stackrel{~}{k}`$ both correspond to dimension four scalar operators (which are the real parts of the two exactly marginal single-trace deformations of the SCFT of which preserve the $`SU(2)\times SU(2)`$ global symmetry), and one can show that combinations of the scalars $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ (or $`f_2`$ and $`f_3`$) correspond to operators of dimensions six and eight<sup>7</sup><sup>7</sup>7Note that $`h`$ does not correspond to an operator in the dual field theory, since it can be gauged away by reparametrizations of the radial coordinate.. For $`P=0`$ all these fields are decoupled (at leading order in the deformation from the solution of ), but for non-zero $`P`$ the equations of motion couple them all together, and we need to analyze all of them at the same time.
The usual procedure of holographic renormalization starts by finding the solution for arbitrary sources, continues by computing the divergences of the action as a function of the sources, and then introduces counter-terms to cancel these divergences by expressing them as a function of the local fields (the transformation from the sources to the fields is invertible). In our case we have a problem with implementing this procedure because some of the operators involved in our action are irrelevant, meaning that we cannot find the solution with arbitrary sources for these operators as a power series with bounded powers of the radial coordinate as usual. In order to find the sources we need to expand the equations of motion around some solution to linear order and look at the solution to the linearized equations of motion which is larger near the boundary. When expanding the equations of motion of appendix A.1 around the solution of , which in our parametrization is given (to leading order in $`\rho `$) by
$$\begin{array}{cc}\hfill G_{ij}(x,\rho )& =\eta _{ij},\hfill \\ \hfill \mathrm{\Phi }(x,\rho )& =\mathrm{ln}(p_0),\hfill \\ \hfill h(x,\rho )& =\frac{1}{8}P^2p_0+\frac{1}{4}K_0\frac{1}{2}P^2p_0\mathrm{ln}\rho ,\hfill \\ \hfill K(x,\rho )& =K_02P^2p_0\mathrm{ln}\rho ,\hfill \\ \hfill f_2(x,\rho )& =1,\hfill \\ \hfill f_3(x,\rho )& =1,\hfill \end{array}$$
(2.27)
(with some constants $`p_0`$ and $`K_0`$ which are the parameters of the solution), we find the following independent solutions to the linearized equations of motion which we identify with the sources for the various operators :
$$\begin{array}{cc}\hfill (i)& \delta G_{ij}=\stackrel{~}{G}_{ij}(x);\hfill \\ \hfill (ii)& \delta K=\stackrel{~}{K}(x),\delta h=\frac{1}{4}\stackrel{~}{K}(x);\hfill \\ \hfill (iii)& \delta \mathrm{\Phi }=\stackrel{~}{p}(x)/p_0,\delta h=\frac{1}{8}P^2\stackrel{~}{p}(x)\frac{1}{2}P^2\stackrel{~}{p}(x)\mathrm{ln}\rho ,\delta K=2P^2\stackrel{~}{p}(x)\mathrm{ln}\rho ;\hfill \\ \hfill (iv)& \delta f_3=\alpha _6(x)\rho ^2,\delta f_2=4\alpha _6(x)\rho ^2,\delta h=\alpha _6(x)P^2p_0\rho ^2,\delta K=2\alpha _6(x)P^2p_0\rho ^2;\hfill \\ \hfill (v)& \delta h=\alpha _8(x)\rho ^4.\hfill \end{array}$$
(2.28)
The first three sources are those of the operators which are marginal for $`P=0`$ – the stress-energy tensor and the two scalar operators of dimension four – and the last two are the sources of the two irrelevant scalar operators<sup>8</sup><sup>8</sup>8The fifth source $`\alpha _8(x)`$ naively couples only to $`h`$ which, as we mentioned, is not a physical field, but in fact by a reparametrization of the $`\rho `$ coordinate one can rewrite the corresponding solution in a way which involves $`f_2`$ and $`f_3`$. We chose to write the solution in the form above for convenience.. In all cases we wrote down only the leading $`\rho `$-dependence of the solutions – in general there are corrections to the expressions above which involve powers of $`\rho ^2\mathrm{ln}^n(\rho )`$ (for some integer $`n`$) multiplying the $`\rho `$-dependence of the terms we wrote, and which generally involve all the fields (not just the ones which are turned on at the leading order).
Finding a solution to the full non-linear equations involving all the sources above is an ill-defined question since some of these sources are irrelevant. In order to have a well-defined solution we need to set the sources $`\alpha _6(x)=\alpha _8(x)=0`$ (later we will take these sources to be infinitesimal in order to compute the correlation functions of the corresponding operators, but we cannot take them to be more than infinitesimal). Once we do this there are no negative powers of $`\rho `$ in any of the fields, so all fields have a well-defined expansion in powers of $`\rho `$ (and $`\mathrm{ln}\rho `$). We will also make another simplifying assumption and we will not introduce any sources for the other two scalar operators, leaving the corresponding fields to take (as $`\rho 0`$) the ($`x`$-independent) values they take in (2.27). It would be interesting to analyze the solutions with arbitrary sources for these fields, but we postpone this to future work. Thus, the only field that we allow an arbitrary source for is the metric, which we take to be of the form $`G_{ij}(x,\rho )=G_{ij}^{(0)}(x)+𝒪(\rho ^2\mathrm{ln}^n(\rho ))`$. This means that the solutions we construct will describe the cascading gauge theory compactified on an arbitrary manifold (with metric $`G_{ij}^{(0)}`$, since the source for the stress-energy tensor is just the metric of the cascading field theory), but without any deformations to its Lagrangian. Of course, the fact that we do not allow arbitrary sources means that even though we will be able to perform the second step of holographic renormalization – expressing the divergences of the action in terms of the sources in our equations – we will not be able to uniquely translate these divergences into functions of the local fields (since we have many fields but just one arbitrary source), and we will be forced to use other methods to determine the counter-terms. We will discuss this further in the next section.
Next, we would like to find the solution with the source described above. We do this in a perturbative expansion in $`\rho ^2`$, as usual in holographic renormalization. The difference from the usual case is that already our leading order solution contains logarithms of $`\rho `$, and when we solve the equations we find that we need even higher powers of logarithms at the higher orders in $`\rho `$. We use the following parametrization for the solution :
$$\begin{array}{cc}& G_{ij}(x,\rho )=G_{ij}^{(0)}(x)+\rho ^2\left[G_{ij}^{(2,0)}(x)+\mathrm{ln}\rho G_{ij}^{(2,1)}(x)\right]\hfill \\ & +\rho ^4\left[G_{ij}^{(4,0)}(x)+\mathrm{ln}\rho G_{ij}^{(4,1)}(x)+\mathrm{ln}^2\rho G_{ij}^{(4,2)}(x)+\mathrm{ln}^3\rho G_{ij}^{(4,3)}(x)\right]\hfill \\ & +𝒪(\rho ^6\mathrm{ln}^5\rho )\hfill \end{array}$$
(2.29)
$$\begin{array}{cc}& h(x,\rho )=\frac{1}{8}P^2p_0+\frac{1}{4}K_0\frac{1}{2}P^2p_0\mathrm{ln}\rho +\rho ^2\left[h^{(2,0)}(x)+\mathrm{ln}\rho h^{(2,1)}(x)+\mathrm{ln}^2\rho h^{(2,2)}(x)\right]\hfill \\ & +\rho ^4\left[h^{(4,0)}(x)+\mathrm{ln}\rho h^{(4,1)}(x)+\mathrm{ln}^2\rho h^{(4,2)}(x)+\mathrm{ln}^3\rho h^{(4,3)}(x)+\mathrm{ln}^4\rho h^{(4,4)}(x)\right]\hfill \\ & +𝒪(\rho ^6\mathrm{ln}^6\rho )\hfill \end{array}$$
(2.30)
$$\begin{array}{cc}& K(x,\rho )=K_02P^2p_0\mathrm{ln}\rho +\rho ^2\left[K^{(2,0)}(x)+\mathrm{ln}\rho K^{(2,1)}(x)\right]\hfill \\ & +\rho ^4\left[K^{(4,0)}(x)+\mathrm{ln}\rho K^{(4,1)}(x)+\mathrm{ln}^2\rho K^{(4,2)}(x)+\mathrm{ln}^3\rho K^{(4,3)}(x)\right]\hfill \\ & +𝒪(\rho ^6\mathrm{ln}^5\rho )\hfill \end{array}$$
(2.31)
$$\begin{array}{cc}& \mathrm{\Phi }(x,\rho )=\mathrm{ln}p_0+\rho ^2\left[p^{(2,0)}(x)+\mathrm{ln}\rho p^{(2,1)}(x)\right]\hfill \\ & +\rho ^4\left[p^{(4,0)}(x)+\mathrm{ln}\rho p^{(4,1)}(x)+\mathrm{ln}^2\rho p^{(4,2)}(x)+\mathrm{ln}^3\rho p^{(4,3)}(x)\right]\hfill \\ & +𝒪(\rho ^6\mathrm{ln}^5\rho )\hfill \end{array}$$
(2.32)
$$\begin{array}{cc}& f_2(x,\rho )=1+\rho ^2\left[a^{(2,0)}(x)+\mathrm{ln}\rho a^{(2,1)}(x)\right]\hfill \\ & +\rho ^4\left[a^{(4,0)}(x)+\mathrm{ln}\rho a^{(4,1)}(x)+\mathrm{ln}^2\rho a^{(4,2)}(x)+\mathrm{ln}^3\rho a^{(4,3)}(x)\right]\hfill \\ & +𝒪(\rho ^6\mathrm{ln}^5\rho )\hfill \end{array}$$
(2.33)
$$\begin{array}{cc}& f_3(x,\rho )=1+\rho ^2\left[b^{(2,0)}(x)+\mathrm{ln}\rho b^{(2,1)}(x)\right]\hfill \\ & +\rho ^4\left[b^{(4,0)}(x)+\mathrm{ln}\rho b^{(4,1)}(x)+\mathrm{ln}^2\rho b^{(4,2)}(x)+\mathrm{ln}^3\rho b^{(4,3)}(x)\right]\hfill \\ & +𝒪(\rho ^6\mathrm{ln}^5\rho )\hfill \end{array}$$
(2.34)
where the leading order terms are specified by the parameters $`G_{ij}^{(0)}(x)`$, $`p_0`$ and $`K_0`$ as above. Note that there is a residual reparametrization ambiguity associated with the choice of $`h`$ in (2.25). We (partially) fix $`h`$ order by order in the perturbative expansion in such a way that at each order in $`\rho `$ all fields are given by a finite order polynomial in $`\mathrm{ln}\rho `$, as indicated in (2.29)-(2.34). This still leaves some diffeomorphisms, of the form
$$\rho \widehat{\rho }=\rho \left[1+\rho ^2\left(\delta _{20}+\delta _{21}\mathrm{ln}\rho \right)+\rho ^4\left(\delta _{40}+\delta _{41}\mathrm{ln}\rho +\delta _{42}\mathrm{ln}^2\rho +\delta _{43}\mathrm{ln}^3\rho \right)+\mathrm{}\right],$$
(2.35)
unfixed, and this will result in some freedom in the solutions which we will find.
We can now find the solution by solving the equations of motion of appendix A.1 order by order in $`\rho `$. We have found the general solution with the boundary conditions described above up to fourth order in $`\rho `$, and it is explicitly given in appendix A.2. At the second order in $`\rho `$ we find that the solution depends on two arbitrary functions of $`x`$. This is related to the reparametrization freedom (2.35), which involves two arbitrary constant parameters at second order. The fact that we find two arbitrary functions rather than two arbitrary parameters is related to the fact that $`x`$-derivatives of fields show up in the equations at higher orders in $`\rho `$ than the fields themselves, so we expect the non-constant part of these functions to be determined at the next order, and indeed it is (as described in appendix A.2). At the fourth order we similarly find four arbitrary functions associated to the reparametrization freedom, and we also find additional arbitrary functions associated (as usual in holographic renormalization) to the one-point functions of the dimension four operators (the two scalar operators and the traceless part of the stress-energy tensor) which are not determined by the UV expansion near the boundary (but which must be determined by the behavior of the solution at large values of $`\rho `$, which we do not discuss here).
The solution we found is rather complicated, and its precise form is not very illuminating. It is useful to check this solution using some of the symmetries of the problem. First, the reduced type IIB action (2.19) is invariant under shifting the dilaton together with rescaling $`P`$ : $`P\alpha P`$, $`e^\mathrm{\Phi }\alpha ^2e^\mathrm{\Phi }`$. This means that our solution must also be invariant under the same transformation, and all fields except the dilaton cannot depend on $`P`$ and $`p_0`$ separately but only on the combination $`P^2p_0`$; it is easy to verify that this is indeed the case. Two additional symmetries involve reparametrizations. The reparametrization (2.35) is a symmetry of our ansatz and boundary condition, so it must take one solution to another, and we verify in appendix A.3 that this is indeed the case. We can also consider a scaling symmetry $`\rho \lambda \rho `$; this does not leave our asymptotic solution invariant, but we can make it into a symmetry if we also give an appropriate transformation to the metric $`G_{ij}^{(0)}`$, $`p_0`$ and $`K_0`$. In appendix A.3 we verify that this symmetry is also satisfied by our solution of appendix A.2. The truncated action (2.19) also has another interesting scaling property : it scales by a factor of $`\beta ^2`$ when we take
$$K\beta K,e^\mathrm{\Phi }\beta e^\mathrm{\Phi },\mathrm{\Omega }_1^4\beta \mathrm{\Omega }_1^4,\mathrm{\Omega }_2^4\beta \mathrm{\Omega }_2^4,g_{\mu \nu }\beta ^{1/2}g_{\mu \nu }.$$
(2.36)
Since this rescales the action by a constant (without acting on the coordinates, just on the fields), it is a symmetry of the equations of motion so it should also be a symmetry of our solution. In appendix A.3 we verify that the symmetry (2.36) is satisfied by our solution of appendix A.2 as well.
It is easy to verify that the solution of appendix A.2 has a good $`P0`$ limit, where it describes the asymptotic behavior of the conformal field theory of compactified on an arbitrary four-manifold. In this limit the solution simplifies considerably; no logarithms appear at second order, while at fourth order some logarithmic terms can appear in the expansions of fields which are dual to operators of dimension four. The reparametrization freedom (2.35) is reduced in this limit just to the terms with no logs, and correspondingly we have less arbitrariness in the solutions (we should set the functions $`a^{(2,1)}`$, $`a^{(4,1)}`$, $`a^{(4,2)}`$ and $`a^{(4,3)}`$ which appear in the solution to zero, and we should also set $`(a^{(4,0)}b^{(4,0)})=𝒪(P)`$).
## 3 Holographic renormalization
In this section we describe the holographic renormalization of the cascading gauge theories. We begin by regularizing the action and computing the divergences of the regularized action. We find that in addition to the familiar power-law divergences of asymptotically anti-de Sitter (AdS) geometries one encounters various logarithmic divergences. The logarithmic divergences are represented by finite order polynomials in $`\mathrm{ln}\rho `$, at least in the specific reparametrization choice<sup>9</sup><sup>9</sup>9Recall that $`h`$ is fixed order-by-order in the perturbative solution in such a way that only a finite number of powers of $`\mathrm{ln}\rho `$ appears at each order in $`\rho `$, see (2.29)-(2.34). of $`h`$ that we made in (2.25).
In the previous section we obtained the asymptotic solution of the holographic dual to the cascading gauge theory for arbitrary finite (not infinitesimal) boundary metric $`G_{ij}^{(0)}(x)`$ and for constant parameters $`\{p_0,K_0\}`$, up to order $`\rho ^4`$. As discussed above, since the sources we introduced are not arbitrary, even if we find counter-terms which give rise to a finite regularized action we are not guaranteed that all correlation functions (given by derivatives of the action with respect to arbitrary sources) will be finite. Instead, we have to directly check that we can find counter-terms that will make all the correlation functions finite. Of course, using just the asymptotic solutions that we computed above we cannot compute arbitrary correlation functions, since generic $`n`$-point functions depend on knowing the full solution and not just its asymptotic form. However, in order to compute one-point functions the asymptotic solutions are actually enough<sup>10</sup><sup>10</sup>10The one-point functions will often depend on arbitrary functions appearing in the asymptotic solutions related to expectation values as mentioned above, but no additional information about the large $`\rho `$ behavior of the solution is needed to compute the one-point functions., since they are just given by the derivatives of the action with respect to infinitesimal sources which can be translated into derivatives of the action with respect to fields near the boundary. In particular, the one-point functions corresponding to the first three sources in (2.28) are simply given by the variation of the action with respect to the parameters $`G_{ij}^{(0)}(x)`$, $`K_0`$ and $`p_0`$, respectively.
So, our procedure to determine the counter-terms will be to compute the one-point functions of the operators coupling to the sources (2.28) and to require that they are all finite. This will not determine the counter-terms uniquely, and we will see that some ambiguities will remain even in the one-point functions which we compute, but some general properties will be independent of these ambiguities and we expect them to be true for any consistent counter-terms (including the correct ones which renormalize the theory for arbitrary sources). Note that, as discussed above, some of our sources correspond to dimension four operators ($`T^{ij}`$ coupling to $`\stackrel{~}{G}_{ij}`$, $`𝒪_{p_0}`$ coupling to $`\stackrel{~}{p}`$ and $`𝒪_{K_0}`$ coupling to $`\stackrel{~}{K}`$) and we will be able to compute their one-point functions using our solution directly. On the other hand, the standard holographic operator-state correspondence (generalized to our background) implies that in order to compute the one-point correlation functions of the operators $`𝒪_6`$ and $`𝒪_8`$ coupling to the sources $`\alpha _6(x)`$ and $`\alpha _8(x)`$, respectively, one needs to know the asymptotic holographic background to order $`\rho ^6`$ and $`\rho ^8`$, respectively. Since we know the supergravity geometry only to order $`\rho ^4`$, we will not be able to compute these one-point functions, but we can still require that the contributions to them from the terms we computed should not lead to divergences. So, we can require that the renormalized one-point correlation functions of the subtracted operators<sup>11</sup><sup>11</sup>11To be defined in subsection 3.2. $`𝒪_6^s`$ and $`𝒪_8^s`$ satisfy (up to possible logarithmic corrections)
$$𝒪_6^s=𝒪(1),𝒪_8^s=𝒪(\rho ^2),\mathrm{as}\rho 0,$$
(3.1)
which is equivalent to saying that these operators do not have the leading and the first two subleading power-law divergences. We find that requiring that all dimension four one-point functions are finite and that there are no power-law divergences in $`𝒪_6^s`$ and $`𝒪_8^s`$ significantly reduces the ambiguities in the renormalized one-point functions of the stress-energy tensor.
We begin by regularizing the action in subsection 3.1, and the computation of the regularized one-point correlation functions is explained in subsection 3.2. In subsection 3.3 we discuss the local counter-terms that are needed for the renormalization of $`T_{ij}`$, $`𝒪_{p_0}`$, $`𝒪_{K_0}`$ and for the cancellation of power-law divergences in $`𝒪_6`$ and $`𝒪_8`$. In this subsection we use a particular ansatz for the counter-terms which we call the ‘minimal subtraction scheme’.
In subsection 3.4 we present results for the one-point correlation functions of the operators $`T_{ij}`$, $`𝒪_{p_0}`$, $`𝒪_{K_0}`$ in the minimal subtraction renormalization scheme. We also discuss the $`P0`$ limit of the minimal subtraction regularization scheme. Already in the minimal subtraction scheme there are some ambiguities in the results, and in more general renormalization schemes additional ambiguities appear. In subsection 3.5 we comment on the ambiguities which appear in general schemes for the renormalization.
Finally, in subsection 3.6 we discuss other possible renormalization prescriptions. Our main result in this section is that the cascading gauge theory can be renormalized with finite one-point functions, and in particular with a finite stress-energy tensor. In section 4 we show that this does not contradict the expectation<sup>12</sup><sup>12</sup>12This was originally proposed in , and further evidence was presented in . that at high temperature the number of effective degrees of freedom of the cascading gauge theory grows as $`K_{eff}^2\mathrm{ln}^2(T/\mathrm{\Lambda })`$, where $`\mathrm{\Lambda }`$ is the strong coupling scale of the cascading gauge theory. One might naively think (given the known thermodynamic properties of the cascading gauge theories ) that a different renormalization scheme would be more natural, in which the renormalized one-point functions are not finite but rather depend on the combination<sup>13</sup><sup>13</sup>13Notice that in this prescription certain $`\mathrm{ln}\rho `$ divergences are allowed, as long as they come from $`K(\rho )`$. $`K(\rho )=K_02P^2\mathrm{ln}\rho `$ (evaluated at the cutoff) rather than on $`K_0`$, but this does not appear to be possible. It would be interesting to explore this second renormalization scheme in more detail.
### 3.1 The regularized action
We write the effective action (2.19) as
$$S_5=\frac{1}{16\pi G_5}__5vol__5_{(5)}$$
(3.2)
where, when evaluated on a solution to the equations of motion,
$$\begin{array}{cc}\hfill _{(5)}=& 2\mathrm{\Omega }_1\mathrm{\Omega }_2^4\left(\underset{5}{\text{ }\text{ }\text{ }}[\mathrm{ln}\mathrm{\Omega }_2]+4([\mathrm{ln}\mathrm{\Omega }_2])^2+[\mathrm{ln}\mathrm{\Omega }_1][\mathrm{ln}\mathrm{\Omega }_2]\right)\hfill \\ & +4\mathrm{\Omega }_1(\mathrm{\Omega }_1^23\mathrm{\Omega }_2^2).\hfill \end{array}$$
(3.3)
We wish to regularize the theory by imposing a cutoff on the space $`_5`$, putting in a boundary $`_5`$ at some $`\rho =\rho _0`$. In some cases the effective action (2.19) evaluated on the equations of motion is a total derivative – for instance, this is true for the cascading gauge theory on $`\times S^3`$ or $`\times S^1\times S^2`$ – and in such cases we can rewrite (2.19) as an integral just over the boundary, but in general (for instance on $`dS_2\times S^2`$ or $`dS_4`$) this is not the case. With the ansatz (2.25) we find that we can write the action evaluated on a solution to the equations of motion as
$$\begin{array}{cc}\hfill \sqrt{g}_{(5)}=& \frac{1}{2}\left[\rho ^3\sqrt{G}f_2^{1/2}f_3^2[\mathrm{ln}h]^{}\right]^{}\hfill \\ & +\rho ^5\sqrt{G}f_2^{1/2}f_3^2\left\{\delta _0+\delta _2+\delta _4+\delta _6\right\},\hfill \end{array}$$
(3.4)
where the prime denotes a derivative with respect to $`\rho `$ and the subscript in $`\delta _i`$ indicates the power-law scaling of the terms as $`\rho 0`$, i.e., $`\delta _i\rho ^i`$. From here on derivative operators and Laplacians will be with respect to the four dimensional metric $`G_{ij}`$ rather than the five dimensional metric. We find
$$\delta _0=4f_3^2\left(f_23f_3\right),$$
(3.5)
$$\delta _2=\rho ^2\left(f_3^1f_3^{\prime \prime }3\rho ^1f_3^1f_3^{}+\frac{1}{2}\text{ }\text{ }\text{ }\text{ }h\right),$$
(3.6)
$$\begin{array}{cc}\hfill \delta _4=& \rho ^2(f_3^1f_3^{}[\mathrm{ln}\sqrt{G}]^{}+f_3^2(f_3^{})^2+2f_3^1f_3h+hf_3^1\text{ }\text{ }\text{ }f_3\hfill \\ & +\frac{1}{2}f_2^1f_3^1f_2^{}f_3^{}+\frac{1}{4}f_2^1f_2h),\hfill \end{array}$$
(3.7)
$$\begin{array}{cc}\hfill \delta _6=& \rho ^2h\left(f_3^2(f_3)^2+\frac{1}{2}f_2^1f_3^1f_2f_3\right).\hfill \end{array}$$
(3.8)
Once we introduce the cutoff as a boundary, in order to get consistent equations of motion we must introduce also a generalized Gibbons-Hawking (GH) term
$$S_{GH}=\frac{1}{8\pi G_5}__5d^4x\sqrt{det(\gamma _{\mu \nu })}\mathrm{\Omega }_1\mathrm{\Omega }_2^4\left(_\mu n^\mu +n^\mu _\mu \mathrm{ln}\left(\mathrm{\Omega }_1\mathrm{\Omega }_2^4\right)\right),$$
(3.9)
where $`n^\mu `$ is a unit space-like vector orthogonal to the four-dimensional boundary $`_5`$, and $`\gamma _{\mu \nu }`$ is the induced metric on $`_5`$
$$\gamma _{\mu \nu }g_{\mu \nu }n_\mu n_\nu .$$
(3.10)
Note that $`_\mu n^\mu `$ is nothing but the extrinsic curvature of the boundary $`_5`$, calculated with the metric (2.14), and that the whole boundary term (3.9) coincides with the Kaluza-Klein reduction of the standard Gibbons-Hawking term for the action (2.1) with a nine dimensional boundary $`_{10}`$. In the ansatz (2.25) we have
$$\begin{array}{cc}\hfill n^\mu =& \delta _\rho ^\mu \rho h^{1/4},\hfill \\ \hfill \gamma _{ij}=& \rho ^2h^{1/2}G_{ij},\hfill \\ \hfill \sqrt{det(\gamma _{\mu \nu })}=& \rho ^4h^1\sqrt{G},\hfill \end{array}$$
(3.11)
(evaluated at $`\rho =\rho _0`$) and we find
$$S_{GH}=\frac{1}{16\pi G_5}__5d^4x\sqrt{G}_{GH},$$
(3.12)
where
$$_{GH}=2\rho ^3f_2^{1/2}f_3^2\left\{\frac{1}{4}[\mathrm{ln}h]^{}4\rho ^1+[\mathrm{ln}\sqrt{G}]^{}+\frac{1}{2}f_2^1f_2^{}+2f_3^1f_3^{}\right\}.$$
(3.13)
The total regularized effective action is
$$\begin{array}{cc}\hfill S_4^\rho =& \frac{1}{16\pi G_5}__5d^4x\sqrt{G}_{(4)}^\rho ,\hfill \\ \hfill \sqrt{G}_{(4)}^\rho =& \sqrt{G}_{GH}+_{\rho _0}𝑑\rho \sqrt{g}_{(5)}.\hfill \end{array}$$
(3.14)
Generally this will diverge, and we will need to add to it some counter-term Lagrangian. We define the subtracted action to be
$$\begin{array}{cc}\hfill S_{tot}=& \frac{1}{16\pi G_5}__5d^4x\sqrt{G}\left(_{(4)}^\rho +h^1\rho ^4^{counter}\right),\hfill \end{array}$$
(3.15)
where the (local) counter-term Lagrangian must be chosen in such a way that correlation functions computed from $`S_{tot}`$ remain finite in the limit $`\rho _00`$. The renormalized action is then simply
$$S_{eff}=\underset{\rho _00}{lim}S_{tot}.$$
(3.16)
As explained in , one should distinguish between $`S_{eff}`$ and $`S_{tot}`$, as the variations required to obtain correlation functions should be performed before the limit $`\rho _00`$ is taken. This is necessary in order to implement the subtraction covariantly.
### 3.2 Regularized one-point correlation functions
As explained in the previous subsection, we can write the subtracted effective action as
$$S_{tot}=S_5+S_{GH}+S_{ct}$$
(3.17)
where $`S_5`$ is the bulk term (3.2), $`S_{GH}`$ is the generalized Gibbons-Hawking term (3.9), and we still need to determine the counter-term action
$$S_{ct}=\frac{1}{16\pi G_5}__5d^4x\sqrt{\gamma }^{counter}.$$
(3.18)
The holographic renormalization is implemented by assuming that $`^{counter}`$ is a local functional of the fields $`\{\gamma _{ij},K,\mathrm{\Phi },\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ on the regularization boundary $`_5`$. Under a generic variation of the fields in the action we have
$$\begin{array}{cc}\hfill \delta S_{tot}=& __5\sqrt{g}\left\{[\mathrm{}]_{\mu \nu }\delta g^{\mu \nu }+[\mathrm{}]\delta \mathrm{\Phi }+[\mathrm{}]\delta K+[\mathrm{}]\delta \mathrm{\Omega }_1+[\mathrm{}]\delta \mathrm{\Omega }_2\right\}\hfill \\ & +__5\sqrt{\gamma }\left\{[\mathrm{}]_{ij}\delta \gamma ^{ij}+[\mathrm{}]\delta \mathrm{\Phi }+[\mathrm{}]\delta K+[\mathrm{}]\delta \mathrm{\Omega }_1+[\mathrm{}]\delta \mathrm{\Omega }_2\right\},\hfill \end{array}$$
(3.19)
where $`[\mathrm{}]`$ in the bulk $`_5`$ integral stand for the corresponding five dimensional equations of motion (2.20)-(2.24), while the $`[\mathrm{}]`$ in the boundary $`_5`$ integral in (3.19) involve only the boundary metric $`\gamma _{ij}`$ and the boundary values of the fields $`K,\mathrm{\Phi },\mathrm{\Omega }_1,\mathrm{\Omega }_2`$. Clearly, evaluated on a solution to the bulk equations of motion, $`\delta S_{tot}`$ does not depend on $`\delta g_{55}`$, and thus we have the general expression
$$\delta S_{tot}=\delta S_{tot}[\delta \gamma _{ij},\delta K,\delta \mathrm{\Phi },\delta \mathrm{\Omega }_1,\delta \mathrm{\Omega }_2]$$
(3.20)
depending only on the values of the fields on $`_5`$.
In order to compute one-point functions we need to take derivatives of this action with respect to our sources. As mentioned above, the one-point function of the stress-energy tensor is the variation with respect to $`G_{ij}^{(0)}`$, and we can write the one-point functions of $`𝒪_{p_0}`$ and $`𝒪_{K_0}`$ as variations with respect to $`p_0`$ and $`K_0`$, respectively. Notice that $`\{\delta \gamma _{ij},\delta K,\delta \mathrm{\Phi },\delta \mathrm{\Omega }_1,\delta \mathrm{\Omega }_2\}`$ depend implicitly on the variation of the source boundary metric $`\delta G_{ij}^{(0)}`$ and on $`\delta p_0`$ , $`\delta K_0`$. Given (3.20), the subtracted one-point correlation functions of the operators dual to $`\{G_{ij}^{(0)},K_0,p_0\}`$ can be evaluated as
$$\begin{array}{cc}\hfill 𝒪_{p_0}^s\frac{16\pi G_5}{\sqrt{G}}\frac{\delta S_{tot}}{\delta p_0}=& \frac{16\pi G_5}{\sqrt{\gamma }}\frac{1}{\rho ^4h}[\frac{\delta S_{tot}}{\delta \mathrm{\Phi }}\frac{\delta \mathrm{\Phi }}{\delta p_0}+\frac{\delta S_{tot}}{\delta K}\frac{\delta K}{\delta p_0}\hfill \\ & +\frac{\delta S_{tot}}{\delta \gamma _{ij}}\frac{\delta \gamma _{ij}}{\delta p_0}+\frac{\delta S_{tot}}{\delta \mathrm{\Omega }_1}\frac{\delta \mathrm{\Omega }_1}{\delta p_0}+\frac{\delta S_{tot}}{\delta \mathrm{\Omega }_2}\frac{\delta \mathrm{\Omega }_2}{\delta p_0}],\hfill \end{array}$$
(3.21)
$$\begin{array}{cc}\hfill 𝒪_{K_0}^s\frac{16\pi G_5}{\sqrt{G}}\frac{\delta S_{tot}}{\delta K_0}=& \frac{16\pi G_5}{\sqrt{\gamma }}\frac{1}{\rho ^4h}[\frac{\delta S_{tot}}{\delta \mathrm{\Phi }}\frac{\delta \mathrm{\Phi }}{\delta K_0}+\frac{\delta S_{tot}}{\delta K}\frac{\delta K}{\delta K_0}\hfill \\ & +\frac{\delta S_{tot}}{\delta \gamma _{ij}}\frac{\delta \gamma _{ij}}{\delta K_0}+\frac{\delta S_{tot}}{\delta \mathrm{\Omega }_1}\frac{\delta \mathrm{\Omega }_1}{\delta K_0}+\frac{\delta S_{tot}}{\delta \mathrm{\Omega }_2}\frac{\delta \mathrm{\Omega }_2}{\delta K_0}],\hfill \end{array}$$
(3.22)
$$\begin{array}{cc}\hfill 𝒪_{G^{(0)}}^{ijs}\frac{16\pi G_5}{\sqrt{G}}\frac{\delta S_{tot}}{\delta G_{ij}^{(0)}}=& \frac{16\pi G_5}{\sqrt{\gamma }}\frac{1}{\rho ^4h}[\frac{\delta S_{tot}}{\delta \mathrm{\Phi }}\frac{\delta \mathrm{\Phi }}{\delta G_{ij}^{(0)}}+\frac{\delta S_{tot}}{\delta K}\frac{\delta K}{\delta G_{ij}^{(0)}}\hfill \\ & +\frac{\delta S_{tot}}{\delta \gamma _{kl}}\frac{\delta \gamma _{kl}}{\delta G_{ij}^{(0)}}+\frac{\delta S_{tot}}{\delta \mathrm{\Omega }_1}\frac{\delta \mathrm{\Omega }_1}{\delta G_{ij}^{(0)}}+\frac{\delta S_{tot}}{\delta \mathrm{\Omega }_2}\frac{\delta \mathrm{\Omega }_2}{\delta G_{ij}^{(0)}}].\hfill \end{array}$$
(3.23)
The renormalized one-point correlation functions of the corresponding operators are then simply evaluated as
$$\begin{array}{cc}\hfill 𝒪_{p_0}=& \underset{\rho _00}{lim}𝒪_{p_0}^s,\hfill \\ \hfill 𝒪_{K_0}=& \underset{\rho _00}{lim}𝒪_{K_0}^s,\hfill \\ \hfill 8\pi G_5T^{ij}=& \underset{\rho _00}{lim}𝒪_{G^{(0)}}^{ijs},\hfill \end{array}$$
(3.24)
where we have defined a standard normalization for $`T^{ij}`$. For each subtracted correlator we find it convenient to separate the regularized contribution from $`S_5+S_{GH}`$, and the counter-term contribution from $`S_{ct}`$. We can do this separation separately for every term in the equations – for example, we can write
$$\begin{array}{cc}\hfill 𝒪_\mathrm{\Phi }^s& \frac{16\pi G_5}{\sqrt{G}}\frac{\delta S_{tot}}{\delta \mathrm{\Phi }}\frac{1}{\rho ^4h}\left(𝒪_\mathrm{\Phi }^\rho +𝒪_\mathrm{\Phi }^c\right)\hfill \\ & =\frac{1}{\rho ^4h}\left(\frac{16\pi G_5}{\sqrt{\gamma }}\frac{\delta (S_5+S_{GH})}{\delta \mathrm{\Phi }}+\frac{16\pi G_5}{\sqrt{\gamma }}\frac{\delta S_{ct}}{\delta \mathrm{\Phi }}\right).\hfill \end{array}$$
(3.25)
Using (2.19), (3.9) we then find that the contributions from the original action are given by
$$𝒪_\mathrm{\Phi }^\rho =\frac{16\pi G_5}{\sqrt{\gamma }}\frac{\delta (S_5+S_{GH})}{\delta \mathrm{\Phi }}=\rho hf_2^{1/2}f_3^2\mathrm{\Phi }^{},$$
(3.26)
$$𝒪_K^\rho =\frac{16\pi G_5}{\sqrt{\gamma }}\frac{\delta (S_5+S_{GH})}{\delta K}=\frac{f_2^{1/2}\rho K^{}}{2P^2e^\mathrm{\Phi }},$$
(3.27)
$$𝒪_{\mathrm{\Omega }_1}^\rho =\frac{16\pi G_5}{\sqrt{\gamma }}\frac{\delta (S_5+S_{GH})}{\delta \mathrm{\Omega }_1}=2\rho h^{3/4}f_3^2\left([\mathrm{ln}\sqrt{G}]^{}+2f_3^1f_3^{}4\rho ^1\right),$$
(3.28)
$$\begin{array}{cc}\hfill 𝒪_{\mathrm{\Omega }_2}^\rho =\frac{16\pi G_5}{\sqrt{\gamma }}\frac{\delta (S_5+S_{GH})}{\delta \mathrm{\Omega }_2}=8\rho h^{3/4}f_2^{1/2}f_3^{3/2}(& [\mathrm{ln}\sqrt{G}]^{}+\frac{1}{2}f_2^1f_2^{}\hfill \\ & +\frac{3}{2}f_3^1f_3^{}4\rho ^1),\hfill \end{array}$$
(3.29)
and
$$\begin{array}{cc}\hfill 𝒪_{\gamma ij}^\rho =& \left(\mathrm{\Theta }_{ij}+\mathrm{\Theta }\gamma _{ij}\right)\mathrm{\Omega }_1\mathrm{\Omega }_2^4+n^\lambda _\lambda \left(\mathrm{\Omega }_1\mathrm{\Omega }_2^4\right)\gamma _{ij}\hfill \\ \hfill =& \rho ^2h^{1/2}\{f_2^{1/2}f_3^2\{\frac{1}{2}\rho G_{ij}^{}+G_{ij}(3\rho \left[\mathrm{ln}\sqrt{G}\right]^{}+\frac{3}{4}\rho [\mathrm{ln}h]^{}\hfill \\ & \rho \left[\mathrm{ln}\left(h^{5/4}f_2^{1/2}f_3^2\right)\right]^{})\}\}\hfill \end{array}$$
(3.30)
where all expressions should be evaluated at $`\rho =\rho _0`$ and
$$\mathrm{\Theta }^{ij}=\frac{1}{2}\left(^in^j+^jn^i\right),\mathrm{\Theta }=\mathrm{\Theta }^{ij}\gamma _{ij}.$$
(3.31)
Note that all the regularized correlation functions contain power-law and logarithmic divergences as the cutoff is removed, $`\rho _00`$. Thus, counter-terms must be determined to remove these divergences. Once the counter-term Lagrangian $`^{counter}`$ is specified, one can compute counter-term contributions to the subtracted operators such as $`𝒪_\mathrm{\Phi }^s`$. Then, using the asymptotic solution explicitly given in appendix A.2, the remaining variational derivatives can be evaluated to obtain the subtracted one-point functions. For example,
$$\begin{array}{cc}\hfill 𝒪_{p_0}^s=\frac{1}{\rho ^4h}[& \left(𝒪_\mathrm{\Phi }^\rho +𝒪_\mathrm{\Phi }^c\right)\frac{\delta \mathrm{\Phi }}{\delta p_0}+\left(𝒪_K^\rho +𝒪_K^c\right)\frac{\delta K}{\delta p_0}+\left(𝒪_{\mathrm{\Omega }_1}^\rho +𝒪_{\mathrm{\Omega }_1}^c\right)\frac{\delta \mathrm{\Omega }_1}{\delta p_0}\hfill \\ & +(𝒪_{\mathrm{\Omega }_2}^\rho +𝒪_{\mathrm{\Omega }_2}^c)\frac{\delta \mathrm{\Omega }_2}{\delta p_0}+(𝒪_\gamma ^{ij\rho }+𝒪_\gamma ^{ijc})\frac{\delta \gamma _{ij}}{\delta p_0}],\hfill \end{array}$$
(3.32)
with similar expressions for $`𝒪_{K_0}^s`$, $`𝒪_{G^{(0)}}^{ijs}`$.
Similarly we can analyze the subtracted one-point correlation functions of the operators $`𝒪_6`$ and $`𝒪_8`$. Given the constant infinitesimal sources $`\alpha _6`$ and $`\alpha _8`$ of (2.28) for the corresponding dual supergravity fields, we have
$$\begin{array}{cc}\hfill 𝒪_8^s=\frac{1}{\rho ^4h}[& \left(𝒪_\mathrm{\Phi }^\rho +𝒪_\mathrm{\Phi }^c\right)\frac{\delta \mathrm{\Phi }}{\delta \alpha _8}+\left(𝒪_K^\rho +𝒪_K^c\right)\frac{\delta K}{\delta \alpha _8}+\left(𝒪_{\mathrm{\Omega }_1}^\rho +𝒪_{\mathrm{\Omega }_1}^c\right)\frac{\delta \mathrm{\Omega }_1}{\delta \alpha _8}\hfill \\ & +(𝒪_{\mathrm{\Omega }_2}^\rho +𝒪_{\mathrm{\Omega }_2}^c)\frac{\delta \mathrm{\Omega }_2}{\delta \alpha _8}+(𝒪_\gamma ^{ij\rho }+𝒪_\gamma ^{ijc})\frac{\delta \gamma _{ij}}{\delta \alpha _8}],\hfill \end{array}$$
(3.33)
$$\begin{array}{cc}\hfill 𝒪_6^s=\frac{1}{\rho ^4h}[& \left(𝒪_\mathrm{\Phi }^\rho +𝒪_\mathrm{\Phi }^c\right)\frac{\delta \mathrm{\Phi }}{\delta \alpha _6}+\left(𝒪_K^\rho +𝒪_K^c\right)\frac{\delta K}{\delta \alpha _6}+\left(𝒪_{\mathrm{\Omega }_1}^\rho +𝒪_{\mathrm{\Omega }_1}^c\right)\frac{\delta \mathrm{\Omega }_1}{\delta \alpha _6}\hfill \\ & +(𝒪_{\mathrm{\Omega }_2}^\rho +𝒪_{\mathrm{\Omega }_2}^c)\frac{\delta \mathrm{\Omega }_2}{\delta \alpha _6}+(𝒪_\gamma ^{ij\rho }+𝒪_\gamma ^{ijc})\frac{\delta \gamma _{ij}}{\delta \alpha _6}],\hfill \end{array}$$
(3.34)
where the boundary field variations represent the response of the fields to turning on infinitesimal sources $`\delta \alpha _8`$ and $`\delta \alpha _6`$. Note that in equation (2.28) we only wrote down the source at leading order in $`\rho ^2`$, while naively higher orders in the source will also contribute to the one-point functions using the equations above. However, since the divergences have to cancel order by order in $`\rho ^2`$, it is easy to see that after we have canceled the divergences (in all the operators) at some order in $`\rho `$, the higher order terms in the sources (which naively could contribute to a divergence at the next order) multiply a vanishing expression, so they do not contribute. The first contribution of the higher order terms in the sources comes with one higher power of $`\rho ^2`$ than the contribution of the leading terms, but since the latter is required to be finite as $`\rho _00`$, the higher order terms never contribute in this limit.
In the next subsection we describe the construction of the counter-terms that lead to finite one-point correlation functions (3.24). As we explained earlier in the section, since we know the asymptotic solution of the dual supergravity background only to order $`\rho ^4`$, we cannot compute precisely the subtracted operators $`𝒪_6^s`$ and $`𝒪_8^s`$: at best, we expect to be able to remove only the leading ($`𝒪(\rho ^6\mathrm{ln}^\mathrm{\#}\rho )`$ and $`𝒪(\rho ^8\mathrm{ln}^\mathrm{\#}\rho )`$), next-to-leading ($`𝒪(\rho ^4\mathrm{ln}^\mathrm{\#}\rho )`$ and $`𝒪(\rho ^6\mathrm{ln}^\mathrm{\#}\rho )`$), and next-to-next-to-leading ($`𝒪(\rho ^2\mathrm{ln}^\mathrm{\#}\rho )`$ and $`𝒪(\rho ^4\mathrm{ln}^\mathrm{\#}\rho )`$) power-law divergences in their one-point correlation functions, see (3.1).
### 3.3 Local counter-terms
The counter-term Lagrangian in (3.18) must be a local functional of the fields on $`_5`$. It is useful to separate the dependence of the counter-terms on the metric $`\gamma _{ij}`$ from the dependence on the other fields. Given the structure of the divergences of the regularized correlation functions (3.26)-(3.30), it is clear that the most general form of $`^{counter}`$ must be
$$\begin{array}{c}\hfill ^{counter}=_0+_{}_\gamma +_^2_\gamma ^2+_{ic^2}_{ab\gamma }_\gamma ^{ab}+_{\text{ }\text{ }\text{ }}\underset{\gamma }{\text{ }\text{ }\text{ }}_\gamma +_{kinetic},\end{array}$$
(3.35)
where $`\{_0,_{},_^2,_{ic^2},_{\text{ }\text{ }\text{ }\text{ }}\}`$ are functions of any local fields at the boundary except for the metric $`\gamma _{ij}`$. Notice that (3.35) contains counter-terms proportional to $`\text{ }\text{ }\text{ }\text{ }_\gamma `$. Even though for constant $`K_0,p_0`$ these are total derivatives (up to order $`𝒪(\rho ^4\mathrm{ln}^\mathrm{\#}\rho )`$), so they do not contribute to the stress-energy one-point correlation function, they do contribute to the renormalization of the $`𝒪_{p_0}`$ and $`𝒪_{K_0}`$ operators. Finally, $`_{kinetic}`$, which scales as $`𝒪(\rho ^4\mathrm{ln}^\mathrm{\#}\rho )`$, contains ‘kinetic’ invariants of the boundary scalars $`_i=\{\mathrm{\Phi },K,\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$ of the type<sup>14</sup><sup>14</sup>14Counter-terms of the form $`_1_2_3`$ may also be added.
$$_1\underset{\gamma }{\text{ }\text{ }\text{ }\text{ }}_2=\rho ^2_3\underset{\gamma }{\text{ }\text{ }\text{ }\text{ }}R+\mathrm{}=\underset{\gamma }{\text{ }\text{ }\text{ }\text{ }}\left[\rho ^2_3R\right]+\mathrm{}$$
(3.36)
(where we used the form of the solution in which the leading non-constant terms in the scalar fields are proportional to the curvature $`R`$) and, thus, it is also a total derivative to order $`𝒪(\rho ^4\mathrm{ln}^\mathrm{\#}\rho )`$. Again, even though counter-terms of the type (3.36) do not contribute to the stress-energy one-point function, they are necessary for removing $`\mathrm{ln}\rho `$ divergences in the one-point functions of $`𝒪_{p_0}`$ and $`𝒪_{K_0}`$.
Let us discuss in detail the evaluation of $`𝒪_\gamma ^{ijc}`$ ($`𝒪_\mathrm{\Phi }^c`$, $`𝒪_K^c`$, $`𝒪_{\mathrm{\Omega }_1}^c`$ and $`𝒪_{\mathrm{\Omega }_2}^c`$ can be evaluated analogously). We will need the following asymptotic expansions:
$$\begin{array}{cc}\hfill _\gamma =& h^{1/2}\rho ^2\left(R_\rho +\frac{3}{2}\text{ }\text{ }\text{ }\mathrm{ln}h\right)+𝒪(\rho ^6\mathrm{ln}^\mathrm{\#}\rho ),\hfill \\ \hfill _{ij\gamma }=& R_{ij\rho }+\frac{1}{2}h^1_i_jh+\frac{1}{4}h^1G_{ij}\text{ }\text{ }\text{ }h+𝒪(\rho ^4\mathrm{ln}^\mathrm{\#}\rho ),\hfill \\ \hfill ^2=& h\rho ^4R^2+𝒪(\rho ^6\mathrm{ln}^\mathrm{\#}\rho ),\hfill \\ \hfill _{ab\gamma }_\gamma ^{ab}=& h\rho ^4R_{ab}R^{ab}+𝒪(\rho ^6\mathrm{ln}^\mathrm{\#}\rho ),\hfill \end{array}$$
(3.37)
where
$$\begin{array}{cc}\hfill R_{ij\rho }=& R_{ij}+\rho ^2(\frac{h}{4}(\text{ }\text{ }\text{ }R_{ij}2(R_{iabj}R^{ab}+R_{ik}R_j^k))\hfill \\ & +\left(\frac{_i_jR}{96}+\frac{G_{ij}^{(0)}\text{ }\text{ }\text{ }R}{192}\right)\left(2K_0P^2p_0+4P^2p_0\mathrm{ln}\rho \right)\hfill \\ & +\frac{1}{4}\left(2_i_ja^{(2,0)}+_i_ja^{(2,1)}\left(1+2\mathrm{ln}\rho \right)\right)\hfill \\ & +\frac{1}{8}G_{ij}^{(0)}(2\text{ }\text{ }\text{ }a^{(2,0)}+\text{ }\text{ }\text{ }a^{(2,1)}(1+2\mathrm{ln}\rho )))+𝒪(\rho ^4)\hfill \end{array}$$
(3.38)
and
$$R_\rho =G^{ij}R_{ij\rho }.$$
(3.39)
In (3.37) and (3.38) the differential operators on the right hand side are evaluated with the metric $`G_{ij}`$.
Given (3.35) we can compute the contribution of $`S_{ct}`$ to the subtracted operator $`𝒪_\gamma ^{ijs}`$,
$$\begin{array}{cc}\hfill 𝒪_\gamma ^{ijc}=& \frac{16\pi G_5}{\sqrt{\gamma }}\frac{\delta S_{ct}}{\delta \gamma _{ij}}\hfill \\ \hfill =& \frac{1}{2}_0\gamma ^{ij}+\left(_\gamma ^{ij}+\frac{1}{2}_\gamma \gamma ^{ij}\gamma ^{ij}\mathrm{}_\gamma +_\gamma ^i_\gamma ^j\right)_{}\hfill \\ & +_{ic^2}\left(\frac{1}{2}_{ab\gamma }_\gamma ^{ab}\gamma ^{ij}+_\gamma ^i_\gamma ^j_\gamma \underset{\gamma }{\text{ }\text{ }\text{ }}_\gamma ^{ij}\frac{1}{2}\gamma ^{ij}\underset{\gamma }{\text{ }\text{ }\text{ }}_\gamma +2_\gamma ^{aijb}_{ab\gamma }\right)\hfill \\ & +_^2\left(\frac{1}{2}_\gamma ^2\gamma ^{ij}2_\gamma _\gamma ^{ij}+2_\gamma ^i_\gamma ^j_\gamma 2\gamma ^{ij}\underset{\gamma }{\text{ }\text{ }\text{ }}_\gamma \right)+𝒪(\rho ^8\mathrm{ln}^\mathrm{\#}\rho ).\hfill \end{array}$$
(3.40)
Notice that the terms which are relevant for removing divergences in the one-point functions of $`𝒪_{p_0}`$, $`𝒪_{K_0}`$ and $`T^{ij}`$ are terms up to order $`\rho ^4`$ in $`_0`$, terms up to order $`\rho ^2`$ in $`_{}`$, and terms up to order $`\rho ^0`$ in $`_^2`$, $`_{ic^2}`$, $`_{\text{ }\text{ }\text{ }\text{ }}`$ and $`_{kinetic}`$.
Ideally we would like all the counter-terms to be functions just of the local fields at the boundary. However, already in the asymptotically AdS case it turns out that it is not possible to do this, and terms which explicitly involve $`\mathrm{ln}\rho `$ are necessary; these terms are related to the conformal anomaly. Since our theory reduces to an asymptotically AdS theory as $`P0`$, we expect that such explicit $`\rho `$-dependent terms will be required in our case as well, so we will allow them in our ansatz. We will find that terms involving up to three powers of $`\mathrm{ln}\rho `$ are required for renormalizing the cascading gauge theories.
In this subsection we will discuss a specific ansatz for the counter-terms which turns out to suffice for obtaining finite one-point functions; more general possibilities will be discussed in subsection 3.5. We begin by noting that the simple counter-term Lagrangian
$$\begin{array}{cc}\hfill _{power}^{counter}=& K2\mathrm{\Omega }_1^48\mathrm{\Omega }_2^4+_\gamma \mathrm{\Omega }_1^2\left(\frac{1}{12}K+\frac{1}{12}P^2e^\mathrm{\Phi }\frac{1}{6}\mathrm{\Omega }_1^4\right)\hfill \end{array}$$
(3.41)
removes all power law divergences in $`𝒪_{p_0}^s`$, $`𝒪_{K_0}^s`$ and $`𝒪_{G_{ij}^{(0)}}^s`$, leaving only logarithmic divergences which still need to be canceled. Thus, it is convenient to parameterize the complete counter-term Lagrangian as follows (making a specific choice for the form of $`_{kinetic}`$ that will turn out to be sufficient) :
$$\begin{array}{cc}\hfill ^{counter}=& K2\mathrm{\Omega }_1^48\mathrm{\Omega }_2^4+𝒜_4+𝒜_6+_\gamma \mathrm{\Omega }_1^2\left(\frac{1}{12}K+\frac{1}{12}P^2e^\mathrm{\Phi }\frac{1}{6}\mathrm{\Omega }_1^4+_2+_4\right)\hfill \\ & +_\gamma ^2(_^2^0+_^2^2)+_{ab\gamma }_\gamma ^{ab}(_{ic^2}^0+_{ic^2}^2)+\underset{\gamma }{\text{ }\text{ }\text{ }}_\gamma (_{\text{ }\text{ }\text{ }}^0+_{\text{ }\text{ }\text{ }}^2)\hfill \\ & +\delta _1\mathrm{ln}\rho \mathrm{\Omega }_1^6\underset{\gamma }{\text{ }\text{ }\text{ }}\mathrm{\Phi }+\delta _2\mathrm{\Omega }_1^6\underset{\gamma }{\text{ }\text{ }\text{ }}\mathrm{\Phi },\hfill \end{array}$$
(3.42)
where the subscript in $`𝒜,`$ and the superscript in $``$ indicates the scaling near the boundary of the local field configuration represented by that coefficient, for instance $`𝒜_n𝒪(\rho ^n\mathrm{ln}^\mathrm{\#}\rho )`$. Note that the counter-terms containing $`\{𝒜_6,_4,_^2^2,_{ic^2}^2,_{\text{ }\text{ }\text{ }\text{ }}^2\}`$ do not affect the one-point functions of $`𝒪_{p_0}`$, $`𝒪_{K_0}`$ and $`𝒪_{G_{ij}^{(0)}}`$, but they can contribute to the renormalization of power-law divergences in the one-point functions of the irrelevant operators $`𝒪_6`$ and $`𝒪_8`$. Counter-terms scaling as higher powers of $`\rho `$ do not contribute to any of the one-point functions we compute so we ignore them.
Since, as discussed above, we do not have a systematic way to determine the counter-terms, we will write down an ansatz for the counter-terms and check if it can lead to finite one-point functions (including (3.1)). There are two scaling symmetries which we can use to constrain the form of the counter-terms. As mentioned in the previous section, the type IIB action has a scaling symmetry of $`P\alpha P`$, $`e^\mathrm{\Phi }\alpha ^2e^\mathrm{\Phi }`$, and this will also be a property of the divergent terms in the action, so we can choose our counter-terms to depend on the dilaton and on $`P`$ only through the combination $`P^2e^\mathrm{\Phi }`$. Furthermore, both the truncated action (2.19) and the generalized Gibbons-Hawking term (3.9) have weight two under the scaling symmetry (2.36), $`S_5\beta ^2S_5`$ and $`S_{GH}\beta ^2S_{GH}`$. Thus, the divergences will have the same scaling, so in order to cancel them we need to have the same scaling also for the counter-terms that we add – this means that $`^{counter}`$ in (3.18) should scale with a factor of $`\beta `$ under this transformation, and we will use this to constrain our counter-terms.
In addition we will assume that the counter-terms contain only non-negative integer powers of $`K`$ and of $`P^2`$; again this is consistent with the structure of the divergences of the action, so it is reasonable to assume that the “correct” counter-terms (which render all correlation functions finite) should have the same property.
It is convenient to introduce the following short-hand notations
$$\begin{array}{cc}\hfill X_a& \left(1\frac{\mathrm{\Omega }_2^2}{\mathrm{\Omega }_1^2}\right),\hfill \\ \hfill X_b& K4\mathrm{\Omega }_1^4+\frac{1}{2}P^2e^\mathrm{\Phi },\hfill \end{array}$$
(3.43)
for field configurations that scale as $`\rho ^2`$. Using the asymptotic solution of section A.2 we can verify that as $`\rho 0`$
$$X_a=𝒪(\rho ^2),X_b=𝒪(\rho ^2\mathrm{ln}\rho ).$$
(3.44)
We can replace any dependence of the counter-terms on $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ by a dependence on $`X_a`$ and $`X_b`$, and it will be convenient to do this in many places because of the scaling (3.44).
Our “minimal subtraction ansatz” for the counter-terms involves choosing the kinetic terms to take the specific form they take in (3.42). In addition, when we use the counter-terms (3.41) we find that the one-point functions of $`𝒪_6^s`$ and $`𝒪_8^s`$ do not contain divergent terms with negative powers of $`h`$ that are proportional to $`R^2`$ or $`R_{ab}R^{ab}`$, and we require that all counter-terms that we add should preserve this property. This turns out to restrict $`𝒜_4`$ to be proportional to $`X_a^2\mathrm{\Omega }_1^4`$ and $`_2`$ to be proportional to $`X_a`$, and also to restrict $`𝒜_6=_4=0`$. Finally, we restrict all the counter-terms to grow no faster than $`\mathrm{ln}^3\rho `$ (multiplying the appropriate power of $`\rho `$) near the boundary, consistent with the fact that this is the scaling of the divergences in the action. Together with the scaling symmetries described above, this leads to an ansatz containing $`82`$ independent coefficients (including $`\delta _1,\delta _2`$ in (3.42)). By computing the one-point functions using the counter-terms in this ansatz we find that it is possible to get finite one-point functions for the operators $`O_{K_0},O_{p_0},T_{ij}`$ and to satisfy (3.1); these requirements give $`75`$ constraints on the $`82`$ coefficients of the ansatz, leading to a $`7`$-parameter ambiguity in the counter-terms. In the parameterization (3.42) the resulting counter-terms take the form :
$$\begin{array}{cc}\hfill \delta _1=& \frac{50}{21},\hfill \\ \hfill 𝒜_4=& \frac{18}{5}X_a^2\mathrm{\Omega }_1^4,\hfill \\ \hfill _2=& X_a\left(\frac{1}{6}K\frac{1}{30}P^2e^\mathrm{\Phi }\right),\hfill \end{array}$$
(3.45)
$$\begin{array}{cc}\hfill _^2^0=& \frac{1}{144}P^4e^{2\mathrm{\Phi }}\mathrm{ln}^3\rho \frac{1}{96}P^2e^\mathrm{\Phi }\mathrm{ln}^2\rho K\frac{1}{192}\mathrm{ln}\rho K^2\hfill \\ & +\left(\frac{1}{96}+4\kappa _1\right)P^4e^{2\mathrm{\Phi }}\mathrm{ln}^2\rho +\left(\frac{1}{96}+4\kappa _1\right)P^2e^\mathrm{\Phi }\mathrm{ln}\rho K\hfill \\ & +\left(\kappa _1+\frac{1}{1152}\right)K^2+\left(2\kappa _2\frac{43}{2304}\right)P^4e^{2\mathrm{\Phi }}\mathrm{ln}\rho \hfill \\ & +\left(\kappa _2\frac{13}{1152}\right)P^2e^\mathrm{\Phi }K+\kappa _3P^4e^{2\mathrm{\Phi }},\hfill \end{array}$$
(3.46)
$$\begin{array}{cc}\hfill _{ic^2}^0=& \frac{1}{48}P^4e^{2\mathrm{\Phi }}\mathrm{ln}^3\rho +\frac{1}{32}P^2e^\mathrm{\Phi }\mathrm{ln}^2\rho K+\frac{1}{64}\mathrm{ln}\rho K^2+\left(\frac{1}{32}12\kappa _1\right)P^4e^{2\mathrm{\Phi }}\mathrm{ln}^2\rho \hfill \\ & +\left(\frac{1}{32}12\kappa _1\right)P^2e^\mathrm{\Phi }\mathrm{ln}\rho K+\left(\frac{1}{256}3\kappa _1\right)K^2\hfill \\ & +\left(\frac{43}{768}6\kappa _2\right)P^4e^{2\mathrm{\Phi }}\mathrm{ln}\rho +\left(\frac{5}{192}3\kappa _2\right)P^2e^\mathrm{\Phi }K\hfill \\ & +\left(\frac{541}{138240}3\kappa _3\right)P^4e^{2\mathrm{\Phi }},\hfill \end{array}$$
(3.47)
$$\begin{array}{cc}\hfill _{\text{ }\text{ }\text{ }}^0=& \frac{1}{144}P^4e^{2\mathrm{\Phi }}\mathrm{ln}^3\rho +\frac{1}{96}P^2e^\mathrm{\Phi }\mathrm{ln}^2\rho K+\frac{1}{192}\mathrm{ln}\rho K^2\hfill \\ & +\left(\frac{383}{5760}+4\kappa _4\right)P^4e^{2\mathrm{\Phi }}\mathrm{ln}^2\rho +\left(\frac{2231}{40320}+4\kappa _4\right)P^2e^\mathrm{\Phi }\mathrm{ln}\rho K\hfill \\ & +\left(\kappa _4+\frac{1}{64}\right)K^2+\left(\frac{17}{26880}+2\kappa _5+\frac{7}{320}\delta _2\right)P^4e^{2\mathrm{\Phi }}\mathrm{ln}\rho \hfill \\ & +\left(\kappa _5+\frac{1}{64}\delta _2+\frac{29}{2160}\right)P^2e^\mathrm{\Phi }K+\kappa _6P^4e^{2\mathrm{\Phi }},\hfill \end{array}$$
(3.48)
$$\begin{array}{cc}\hfill _^2^2=& X_a(\frac{1}{240}P^4e^{2\mathrm{\Phi }}\mathrm{ln}^2\rho \frac{1}{240}P^2e^\mathrm{\Phi }\mathrm{ln}\rho K\frac{1}{720}K^2+P^4e^{2\mathrm{\Phi }}(\frac{1}{240}+\frac{8}{5}\kappa _1)\mathrm{ln}\rho \hfill \\ & +P^2e^\mathrm{\Phi }(\frac{4}{5}\kappa _1+\frac{1}{1152})K+P^4e^{2\mathrm{\Phi }}(\frac{2}{5}\kappa _2+\frac{43}{34560}))\hfill \\ & +X_b\left(\frac{1}{1152}K+\frac{1}{2304}P^2e^\mathrm{\Phi }\right),\hfill \end{array}$$
(3.49)
$$\begin{array}{cc}\hfill _{ic^2}^2=& X_a(\frac{1}{80}P^4e^{2\mathrm{\Phi }}\mathrm{ln}^2\rho +\frac{1}{80}P^2e^\mathrm{\Phi }\mathrm{ln}\rho K+\frac{1}{240}K^2+P^4e^{2\mathrm{\Phi }}(\frac{1}{80}\frac{24}{5}\kappa _1)\mathrm{ln}\rho \hfill \\ & \frac{12}{5}P^2e^\mathrm{\Phi }K\kappa _1\frac{6}{5}P^4e^{2\mathrm{\Phi }}\kappa _2)+X_b(\frac{1}{384}K\frac{1}{384}P^2e^\mathrm{\Phi }),\hfill \end{array}$$
(3.50)
$$\begin{array}{cc}\hfill _{\text{ }\text{ }\text{ }}^2=& X_a(\frac{1}{240}P^4e^{2\mathrm{\Phi }}\mathrm{ln}^2\rho +\frac{1}{240}P^2e^\mathrm{\Phi }\mathrm{ln}\rho K\frac{1}{480}K^2\hfill \\ & +P^4e^{2\mathrm{\Phi }}(\frac{8}{5}\kappa _4+\frac{533}{14400})\mathrm{ln}\rho +\frac{4}{5}P^2e^\mathrm{\Phi }\kappa _4K+\frac{2}{5}P^4e^{2\mathrm{\Phi }}\kappa _5)\hfill \\ & +X_b\left(\frac{25}{672}P^2e^\mathrm{\Phi }\mathrm{ln}\rho \frac{1}{768}K\left(\frac{167}{23040}+\frac{1}{64}\delta _2\right)P^2e^\mathrm{\Phi }\right).\hfill \end{array}$$
(3.51)
This result depends on seven parameters : $`\kappa _i`$ ($`i=1,\mathrm{},6`$) and $`\delta _2`$. In appendix A.4 we give a simple argument explaining why the coefficients $`\kappa _i`$ turned out to be ambiguous. In addition, we expect to find an ambiguity in the counter-terms corresponding to reparametrizing $`\rho \lambda \rho `$, because of the explicit $`\rho `$-dependence in the counter-terms. This is present already in the asymptotic AdS case, and it explains why the parameter $`\delta _2`$ turned out to be ambiguous (since this reparametrization shifts $`\delta _2`$, in addition to modifying the $`\kappa _i`$ parameters).
### 3.4 Renormalized one-point correlation functions
In this subsection we describe in detail our results in the minimal subtraction renormalization scheme; more general schemes are described in the next subsection.
The one-point function of the stress energy tensor is given by
$$\begin{array}{cc}\hfill 8\pi G_5T_{ij}=& G_{ij}^{(0)}(R_{ab}R^{ab}(\frac{1921}{276480}p_0^2P^4\frac{1}{512}K_0^2+\frac{1}{96}K_0P^2p_0)\hfill \\ & R^2\left(\frac{1}{4608}K_0^2+\frac{337}{51840}p_0^2P^4+\frac{175}{27648}K_0P^2p_0\right)\hfill \\ & +R\left(\frac{1}{16}K_0a^{(2,0)}+\frac{1}{128}P^2p_0a^{(2,0)}+\frac{5}{256}P^2p_0a^{(2,1)}\right)\hfill \\ & +\text{ }\text{ }\text{ }R(\frac{391}{82944}p_0^2P^4\frac{53}{23040}K_0^2+\frac{323}{46080}K_0P^2p_0))\hfill \\ & +R_{aijb}R^{ab}\left(\frac{17}{8640}p_0^2P^4\frac{1}{32}K_0^2+\frac{7}{192}K_0P^2p_0\right)\hfill \\ & R_i^aR_{aj}\left(\frac{1}{64}K_0^2+\frac{1}{256}p_0^2P^4+\frac{1}{64}K_0P^2p_0\right)\hfill \\ & +RR_{ij}\left(\frac{1691}{103680}p_0^2P^4\frac{1}{576}K_0^2+\frac{13}{432}K_0P^2p_0\right)\hfill \\ & R_{ij}\left(\frac{1}{16}P^2p_0a^{(2,1)}+\frac{1}{4}K_0a^{(2,0)}\right)\hfill \\ & _i_jR\left(\frac{2773}{207360}p_0^2P^4+\frac{5}{3456}K_0P^2p_0+\frac{7}{1152}K_0^2\right)\hfill \\ & +\text{ }\text{ }\text{ }R_{ij}\left(\frac{17}{17280}p_0^2P^4\frac{7}{384}K_0P^2p_0+\frac{1}{64}K_0^2\right)\hfill \\ & _i_ja^{(2,0)}\left(\frac{1}{16}P^2p_0+\frac{1}{16}K_0\right)+_i_ja^{(2,1)}\left(\frac{7}{128}P^2p_0+\frac{3}{64}K_0\right)\hfill \\ & +2G_{ij}^{(4,0)}\frac{1}{2}G_{ij}^{(0)}G_a^{(4,0)a}+\frac{3}{2}G_{ij}^{(0)}\left(b^{(4,0)}a^{(4,0)}\right)\hfill \\ & +T_{ij}^{ambiguity},\hfill \end{array}$$
(3.52)
where
$$\begin{array}{cc}\hfill T_{ij}^{ambiguity}=& (\frac{1}{2}p_0^2P^4\kappa _3+\frac{1}{2}p_0P^2\kappa _2K_0+\frac{1}{2}\kappa _1K_0^2)\times \hfill \\ & (2_i_jR+6\text{ }\text{ }\text{ }R_{ij}12R_{aijb}R^{ab}3G_{ij}^{(0)}R_{ab}R^{ab}+R^2G_{ij}^{(0)}4RR_{ij}\hfill \\ & \text{ }\text{ }\text{ }RG_{ij}^{(0)}).\hfill \end{array}$$
(3.53)
The one-point function of the trace of the stress-energy tensor in the “minimal subtraction” ansatz is unambiguously given by
$$\begin{array}{cc}\hfill 8\pi G_5T_i^i=& R_{ab}R^{ab}\left(\frac{1}{96}K_0P^2p_0+\frac{101}{4608}p_0^2P^4+\frac{1}{128}K_0^2\right)\hfill \\ & +R^2\left(\frac{11}{2304}K_0P^2p_0\frac{67}{6912}p_0^2P^4\frac{1}{384}K_0^2\right)\hfill \\ & +RP^2p_0\left(\frac{1}{32}a^{(2,0)}+\frac{1}{64}a^{(2,1)}\right)\hfill \\ & +\text{ }\text{ }\text{ }R\left(\frac{43}{2304}K_0P^2p_0+\frac{151}{11520}p_0^2P^4+\frac{1}{384}K_0^2\right)\hfill \\ & +6(b^{(4,0)}a^{(4,0)}).\hfill \end{array}$$
(3.54)
The one-point function of $`𝒪_{p_0}`$ is given by
$$\begin{array}{cc}\hfill 𝒪_{p_0}=& R_{ab}R^{ab}\left(\frac{263}{34560}p_0P^4+\frac{3}{64}K_0P^2\right)+R^2\left(\frac{79}{4608}K_0P^2\frac{83}{138240}p_0P^4\right)\hfill \\ & +RP^2\left(\frac{19}{192}a^{(2,0)}+\frac{13}{384}a^{(2,1)}\right)+\text{ }\text{ }\text{ }R\left(\frac{77}{3456}K_0P^2+\frac{33}{1280}p_0P^4\right)\hfill \\ & +\frac{1}{p_0}\left(4p^{(4,0)}3\left(b^{(4,0)}a^{(4,0)}\right)\right)+𝒪_{p_0}^{ambiguity},\hfill \end{array}$$
(3.55)
where
$$\begin{array}{cc}\hfill 𝒪_{p_0}^{ambiguity}=& (3R_{ab}R^{ab}R^2)\left(2P^4\kappa _3p_0P^2\kappa _2K_0\right)\hfill \\ & +\text{ }\text{ }\text{ }R\left(2P^4p_0\kappa _6+P^2\kappa _5K_0+\delta _2P^2\left(\frac{7}{320}K_0\frac{3}{640}P^2p_0\right)\right).\hfill \end{array}$$
(3.56)
The one-point function of $`𝒪_{K_0}`$ is given by
$$\begin{array}{cc}\hfill 𝒪_{K_0}=& R_{ab}R^{ab}\left(\frac{1}{192}K_0+\frac{7}{1152}P^2p_0\right)+R^2\left(\frac{1}{2304}K_0+\frac{5}{13824}P^2p_0\right)\hfill \\ & R\left(\frac{1}{32}a^{(2,0)}+\frac{1}{64}a^{(2,1)}\right)+\text{ }\text{ }\text{ }R\left(\frac{53}{1920}K_0+\frac{11}{34560}P^2p_0\right)\hfill \\ & +\frac{6}{P^2p_0}\left(a^{(4,0)}b^{(4,0)}\right)+𝒪_{K_0}^{ambiguity},\hfill \end{array}$$
(3.57)
where
$$\begin{array}{cc}\hfill 𝒪_{K_0}^{ambiguity}=& (3R_{ab}R^{ab}R^2)\left(2K_0\kappa _1\kappa _2p_0P^2\right)\hfill \\ & +\text{ }\text{ }\text{ }R\left(2K_0\kappa _4+p_0P^2\kappa _5\right).\hfill \end{array}$$
(3.58)
The conformal anomaly is the transformation of the action under scaling transformations (see for a review). It is generally defined as (for non-conformal theories that have some beta functions $`\beta _j`$ for the couplings to operators $`𝒪_j`$)
$$\mathrm{conformal}\mathrm{anomaly}=T_i^i\frac{1}{2}\underset{j}{}\beta _j𝒪_j.$$
(3.59)
In our case the form of the asymptotic solution (2.27) (and the analysis in the appendix of the scaling transformation (A.40)) indicates that the coupling $`K_0`$ depends on the scale, and that its derivative with respect to the logarithm of the scale is given by $`(2P^2p_0)`$. Thus, we have
$$\mathrm{conformal}\mathrm{anomaly}=T_i^i+P^2p_0𝒪_{K_0}.$$
(3.60)
Using the previous results we find that this is given by
$$\begin{array}{cc}\hfill \mathrm{conformal}\mathrm{anomaly}=& \left(3R_{ab}R^{ab}R^2\right)\left(\frac{1}{384}K_0^2\frac{1}{192}K_0P^2p_0+\frac{43}{4608}p_0^2P^4\right)\hfill \\ & +\text{ }\text{ }\text{ }R\left(\frac{1}{384}K_0^2+\frac{533}{11520}K_0P^2p_0+\frac{29}{2160}p_0^2P^4\right)\hfill \\ & +P^2p_0𝒪_{K_0}^{ambiguity}.\hfill \end{array}$$
(3.61)
Note that, as expected of the conformal anomaly, this is independent of any parameters (like ($`a^{(4,0)}b^{(4,0)}`$)) associated to the IR behavior of the theory, and that in our ansatz all the ambiguities in the conformal anomaly are related to the ambiguities in the one-point function of $`𝒪_{K_0}`$. Somewhat surprisingly, we find a finite result for the conformal anomaly, though its precise value is ambiguous because of the ambiguities in our counter-terms.
In any local field theory (see, for instance, ) the conformal anomaly is a linear combination of the Euler density, the Weyl tensor squared and $`\text{ }\text{ }\text{ }\text{ }R`$. As in other theories dual to gravitational backgrounds the conformal anomaly that we computed does not contain terms proportional to $`R_{abcd}R^{abcd}`$, so this requires the conformal anomaly to be a linear combination of a term proportional to $`(3R_{ab}R^{ab}R^2)`$ and a term proportional to $`\text{ }\text{ }\text{ }\text{ }R`$. We find that in our minimal subtraction ansatz this is indeed the case, even though it did not necessarily have to be the case because of the explicit $`\rho `$-dependence of our counter-terms.
It is straightforward to analyze the $`P0`$ limit of the minimal subtraction renormalization scheme; this provides a holographic renormalization for a truncation of the conformal field theory of . The $`P0`$ limit of the asymptotic solution was discussed in section 2.2. It is easy to verify that the one-point functions of the stress tensor and of $`𝒪_{p_0}`$ have a good $`P0`$ limit, but we should be more careful with $`𝒪_{K_0}`$. Note that in this limit equations (2.10) and (2.6) imply that we should change variables from $`K(y)`$ to the variable $`\stackrel{~}{k}(y)`$ related to the two-form, given by $`K(y)=\stackrel{~}{K}_0+P\stackrel{~}{k}(y)`$, which remains finite in the $`P0`$ limit. Let us denote by $`𝒪_{\stackrel{~}{k}}`$ the operator dual to $`\stackrel{~}{k}(y)`$. Then, it is easy to see from (3.57) that $`𝒪_{\stackrel{~}{k}}`$ is finite, provided that as $`P0`$
$$a^{(4,0)}=b^{(4,0)}+\frac{1}{6}Pp_0\stackrel{~}{k}^{(4,0)}+𝒪(P^2),$$
(3.62)
where the parameter $`\stackrel{~}{k}^{(4,0)}`$ is precisely the expectation value
$$𝒪_{\stackrel{~}{k}}=\stackrel{~}{k}^{(4,0)}.$$
(3.63)
Notice that in the $`P0`$ limit of the minimal subtraction scheme, with the scaling (3.62), the scalar one-point functions $`𝒪_{\stackrel{~}{k}}`$, $`𝒪_{p_0}`$ and the conformal anomaly $`T_i^i`$ are unambiguous. The conformal anomaly we find reproduces the known results in the AdS/CFT correspondence , up to the well-known ambiguity in the term in the conformal anomaly proportional to $`\text{ }\text{ }\text{ }\text{ }R`$. This ambiguity arises from finite (in the $`P0`$ limit) counter-terms $`_{ic^2}^0,_^2^0`$, and it is proportional to
$$T_i^i\left(_{ic^2}^0+3_^2^0\right)\text{ }\text{ }\text{ }\text{ }R.$$
(3.64)
In the minimal subtraction scheme, even though both $`_{ic^2}^0`$ and $`_^2^0`$ are ambiguous, there is no ambiguity in the combination (3.64). If we work directly in the conformal $`P=0`$ theory, the ambiguity can be reintroduced by simply shifting
$$\begin{array}{cc}& _^2^0_^2^0+\delta _^2\hfill \\ & _{ic^2}^0_{ic^2}^0+\delta _{ic^2}\hfill \end{array}$$
(3.65)
where $`\delta _{ic^2}`$ and $`\delta _^2`$ are arbitrary constants with $`\delta _{ic^2}+3\delta _^20`$. However, such a simple modification is not possible in the $`P0`$ limit of the holographic renormalization of the cascading gauge theories. The problem is that the one-point functions of the irrelevant operators (see (3.33), (3.34)) are sensitive to the renormalized $`T_i^i`$, and thus to its ambiguity, which is $`\text{ }\text{ }\text{ }\text{ }R`$. So, if a given set of counter-terms renormalizes the irrelevant operators, one would expect that a generic shift as in (3.65) would reintroduce divergences $`\text{ }\text{ }\text{ }\text{ }R`$ in these one-point functions. This is indeed what we find. We would like to emphasize that the fact that we find an unambiguous $`T_i^i`$ is only a feature of the minimal subtraction scheme – an ambiguity proportional to $`\text{ }\text{ }\text{ }\text{ }R`$ does appear in more general renormalization schemes discussed in the next subsection.
Finally, we would like to mention that the $`P0`$ limit of the counter-terms we found in our holographic renormalization correctly reproduces the unambiguous counter-terms of the conformal holographic renormalization , including also the counter-terms with explicit $`\mathrm{ln}\rho `$ dependence . To see this it is useful to note that in this limit $`\mathrm{\Omega }_1^4=\mathrm{\Omega }_2^4=h=K/4`$, so the first line of (3.42) is simply $`\frac{3}{2}K\frac{1}{8}KR`$, while the second line includes the terms $`\mathrm{ln}(\rho ^2)K\frac{1}{32}(R_{ab}R^{ab}\frac{1}{3}R^2)`$, in agreement with equation (5.42) of up to overall normalization factors in the metric and in $`G_5`$ (recall that $`K`$ is a constant in the $`P0`$ limit).
### 3.5 Ambiguities in the choice of counter-terms
In the previous subsection we presented the results of the “minimal subtraction ansatz” which leads to specific finite one-point functions; this ansatz and the resulting one-point functions have a $`7`$-parameter ambiguity. The ansatz we used is the simplest one we could find, but we do not have a good argument that it is correct (in the sense that it leads to all correlation functions being finite); in particular it is not hard to find more general choices for $`_{kinetic}`$, and to add terms with more negative powers of the $`\mathrm{\Omega }`$’s, in a way which preserves the finiteness of all one-point functions<sup>15</sup><sup>15</sup>15One rather ugly feature of our choice of $`_{kinetic}`$ is that when we include non-constant sources for the scalars, there is a term with an explicit $`\mathrm{ln}\rho `$ dependence appearing already at order $`\rho ^2`$. It is possible to choose other forms of $`_{kinetic}`$ which do not have an explicit dependence on $`\rho `$, but they do not have a good $`P0`$ limit. The correct choice should presumably be determined by introducing more general sources for the scalar operators, which we hope to do in future work..
Since with our limited choice of the sources we do not have a way to uniquely determine the counter-terms, we have studied various possibilities for the counter-terms in order to see which of our results for the one-point functions are modified by more general counter-terms and which remain true<sup>16</sup><sup>16</sup>16We have studied the most general possible counter-terms which do not contain very large explicit powers of $`\mathrm{ln}\rho `$ and do not contain large negative powers of the $`\mathrm{\Omega }`$ fields; we believe that these should be properties of the “correct” counter-terms but we do not have a rigorous argument for this.. We find that in a flat background (with $`R_{ij}=0`$) the one-point functions are completely independent of the choice of counter-terms; they are finite and given by the unambiguous results that we found in the previous subsection when $`R_{ij}=0`$. In curved space there are ambiguities in some one-point functions. Recall that already in asymptotically anti-de Sitter spaces the stress-energy tensor has a 2-parameter ambiguity related to a freedom in the definition of the stress-energy tensor in curved space. In the minimal subtraction ansatz we found one of these two ambiguities in the stress tensor (3.53) (multiplied by an arbitrary linear combination of $`P^4p_0^2`$, $`P^2p_0K_0`$ and $`K_0^2`$), and in more general renormalization schemes we find also the other ambiguity (again multiplied by an arbitrary linear combination of $`P^4p_0^2`$, $`P^2p_0K_0`$ and $`K_0^2`$), which contributes also a term proportional to $`\text{ }\text{ }\text{ }\text{ }R`$ to $`T_i^i`$. It turns out that once these ambiguities are determined (by choosing a specific definition for the stress-energy tensor in curved space) the other one-point functions are also determined, up to a freedom in shifting the one-point functions of $`𝒪_{p_0}`$ and $`𝒪_{K_0}`$ by terms proportional to $`\text{ }\text{ }\text{ }\text{ }R`$; presumably this freedom can also be interpreted as some sort of ambiguity in the definition of these operators in curved space.
Overall, with the most general ansatz for the counter-terms that we checked we find an $`11`$-parameter ambiguity in the results for the one-point functions. We expect that some of the ambiguities would remain also once all the correlation functions are renormalized, but some may disappear. One ambiguity which will always remain in any theory of this type (involving an explicit $`\mathrm{ln}\rho `$-dependence in the counter-terms) corresponds to the freedom of redefining the radial coordinate by $`\rho \lambda \rho `$. Using the most general counter-terms we find that the conformal anomaly is not necessarily given by a linear combination of $`3R_{ab}R^{ab}R^2`$ and $`\text{ }\text{ }\text{ }\text{ }R`$, as it must be in any local field theory with a gravitational dual; as mentioned above this is consistent because of the explicit $`\rho `$-dependence of our counter-terms. If we impose this as an additional constraint on the counter-terms we find only a $`9`$-parameter ambiguity in the one-point functions. Generic counter-terms do not necessarily have a good $`P0`$ limit, and this could also be imposed as an additional constraint.
In summary, in flat space we find unambiguous results for all one-point functions, while in curved space there is a finite number of ambiguity parameters (as in asymptotically AdS space). The conformal anomaly turns out to be finite but ambiguous. This is not surprising given that it is related to the number of degrees of freedom in the theory at high energies (it is independent of the IR behavior), which seems to be ill-defined in the cascading theories – the surprise is that in any specific renormalization scheme we actually find a finite answer for the conformal anomaly, as well as for all other one-point functions.
### 3.6 Other possible renormalization schemes
The most intriguing feature of our results is that we find finite results (without any $`\mathrm{ln}^\mathrm{\#}\rho _0`$ divergences) for all one-point correlation functions. This is somewhat surprising since some one-point functions, such as the conformal anomaly, are supposed to count the number of degrees of freedom in the theory, which is believed to diverge at high energies. In particular, in the conformal theory with $`P=0`$ one has (up to $`\text{ }\text{ }\text{ }\text{ }R`$ ambiguities)
$$\begin{array}{c}\hfill 8\pi G_5T_i^{iconformal}=\frac{K_0^2}{16}\left[\frac{1}{8}R_{ab}R^{ab}+\frac{1}{24}R^2\right],\end{array}$$
(3.66)
and the coefficient is proportional to the central charge (in general there are two independent coefficients in the conformal anomaly, but in theories with gravitational duals they are always equal to each other). On the other hand, the behavior of the 5-form flux in the solution of suggests that the number of degrees of freedom in the cascading theories diverges at high energies, and also thermodynamical studies of cascading gauge theories suggest that the number of effective degrees of freedom of the theory accessed at temperature $`T`$ is proportional to $`K_{eff}^2(P^2p_0\mathrm{ln}\frac{T}{\mathrm{\Lambda }})^2`$. Thus, it may be natural to guess that the conformal anomaly of the cascading gauge theory would be as in (3.66), but with a replacement
$$K_0K_{eff}(\rho _0)K_02P^2p_0\mathrm{ln}\rho _0$$
(3.67)
depending explicitly on the cutoff scale. This is quite different from the finite result that we found above (3.61). This leads to a natural question : in the context of the holographic renormalization, is it possible to ‘renormalize’ the cascading gauge theory in such a way that all one-point correlation functions contain $`K_{eff}`$ instead of $`K_0`$ ? Here, by ‘renormalize’ we mean that there are no explicit $`\mathrm{ln}^\mathrm{\#}\rho _0`$ divergences in one-point correlation functions, apart from the ones appearing implicitly in $`K_{eff}`$. With the tool-set of counter-terms as in section 3.3, requiring that all the counter-terms have a good limit as $`P0`$, and that the leading and the first two subleading power-law divergences of $`𝒪_6`$ and $`𝒪_8`$ are removed (3.1), it is possible to show that such a renormalization prescription is not possible. Specifically, one finds that the stress energy one-point function contains certain $`\mathrm{ln}^\mathrm{\#}\rho _0`$ terms (even on manifolds with $`\text{ }\text{ }\text{ }\text{ }R=\text{ }\text{ }\text{ }\text{ }R_{ij}=_i_jR=0`$) which cannot be subtracted by any counter-terms. Needless to say, it would be very interesting to explore these issues further.
## 4 Application : cascading gauge theories at finite temperature
In this section we study the high-temperature thermodynamics of cascading gauge theories, and we verify that the finite results which we found in the previous section (which are unambiguous for the thermodynamics of the theory in flat space) are consistent with the expectation that the cascading theories will have an effective “running rank” $`K_{eff}P^2p_0\mathrm{ln}\frac{T}{\mathrm{\Lambda }}`$.
The thermodynamics of cascading gauge theories was studied in . It was noted there that the $`_{2P}`$ chiral symmetry is restored in the black brane solutions which are dual to the cascading theories at high temperatures (compared to the strong coupling scale $`\mathrm{\Lambda }`$), and thus the high temperature solutions can be described using the ansatz we use in this paper. In the previous studies the stress-energy tensor was not renormalized, so the only thermodynamic property which could be extracted was the entropy density (which depends only on the horizon area). The high-temperature solutions involve a parameter $`K_{}`$ which is the value of the five-form flux at the horizon (the minimal value of the radial coordinate); the solution is constructed in a perturbation expansion in $`P^2/K_{}`$ (valid at high temperatures), and the leading term in this expansion is explicitly known . The parameter $`K_{}`$ should in principle be determined in terms of the temperature, and it appears in the computation of of the entropy density. The authors of used physically reasonable (albeit somewhat ad-hoc) arguments to argue that
$$K_{}=2P^2p_0\mathrm{ln}\frac{T}{\mathrm{\Lambda }}+\mathrm{}$$
(4.1)
where $`\mathrm{}`$ indicate sub-dominant terms in the high temperature limit $`T\mathrm{\Lambda }`$. In this section we will use the renormalized one-point functions of the stress-energy tensor to determine rigorously the relation between $`K_{}`$ and $`T`$, and thus the high-temperature thermodynamics of the cascading gauge theories.
Given the results of the previous section for the renormalized one-point functions of the stress energy tensor, one can compute the ADM mass-density (the energy density) and the pressure of the black brane solution which is holographically dual to the cascading gauge theory at finite temperature, in addition to the entropy density, to leading order in $`P^2/K_{}`$ (as in ). This allows us to explicitly verify (to leading order in $`P^2/K_{}`$) the relation <sup>17</sup><sup>17</sup>17Recall that in the absence of a chemical potential, the free energy density equals minus the pressure.
$$f=𝒫=ϵTs,$$
(4.2)
where $`f`$ and $`ϵ`$ are the free energy and the energy densities, $`s`$ is the entropy density, and $`𝒫`$ is the pressure. Additionally, requiring the first law of thermodynamics
$$dϵ=Tds$$
(4.3)
gives an equation which leads to (4.1).
The rest of this section is organized as follows. In subsection 4.1 we discuss the thermodynamics of cascading gauge theories to leading order in $`P^2/K_{}`$. In subsection 4.2 we briefly comment on the hydrodynamical properties of the cascading gauge theory plasma.
### 4.1 Thermodynamics of cascading gauge theories
Throughout this section we will use the notations and results of <sup>18</sup><sup>18</sup>18We have independently verified all the results which we actually use.. In particular, we use $`K_{}`$ to denote the 5-form flux $`K`$ evaluated at the horizon, and $`a`$ for the non-extremality parameter (which will be related to the temperature below). The ten dimensional Einstein frame metric of the non-extremal cascading solution is
$$\begin{array}{cc}\hfill ds_{10}^2=& \left(\frac{8a}{K_{}v}\right)^{1/2}e^{2P^2\eta }\left[(1v)(dx_0)^2+(dx_\alpha )^2\right]+\frac{\sqrt{K_{}}}{32}e^{2P^2(\eta 5\xi )}\frac{dv^2}{v^2(1v)}\hfill \\ & +\frac{\sqrt{K_{}}}{2}e^{2P^2(\eta \xi )}\left[e^{8P^2\omega }e_\psi ^2+e^{2P^2\omega }\left(e_{\theta _1}^2+e_{\varphi _1}^2+e_{\theta _2}^2+e_{\varphi _2}^2\right)\right],\hfill \end{array}$$
(4.4)
where $`\alpha =1,2,3`$, $`v`$ is a radial coordinate such that $`v1_{}`$ at the horizon and $`v0_+`$ at the boundary, and $`\eta ,\xi ,\omega `$ are functions of $`v`$. To leading order in $`P^2/K_{}`$ the 5-form and dilaton are given by
$$\begin{array}{cc}\hfill K=& K_{}\frac{P^2}{2}\mathrm{ln}v,\hfill \\ \hfill \mathrm{\Phi }=& \mathrm{\Phi }_{}+\frac{1}{4K_{}}\mathrm{Li}_2(1v),\hfill \end{array}$$
(4.5)
where $`\mathrm{\Phi }_{}`$ is the dilaton value at the horizon, and $`\mathrm{Li}_n(z)`$ is the polylogarithm function. As in , we choose the dilaton to vanish at the boundary, corresponding to choosing $`p_0=1`$ in the notations of the previous sections; we will write our results in this section using this choice of $`p_0`$, and the $`p_0`$-dependence can always be reinstated by recalling that factors of $`p_0`$ come together with factors of $`P^2`$. The functions $`\{\eta ,\xi ,\omega \}`$ satisfy the following ordinary differential equations :
$$\begin{array}{cc}& v(1v)\omega ^{\prime \prime }v\omega ^{}\frac{3}{4v}\omega =\frac{1}{40K_{}},\hfill \\ & v(1v)\xi ^{\prime \prime }v\xi ^{}\frac{2}{v}\xi =\frac{1}{40K_{}},\hfill \\ & v(1v)\eta ^{\prime \prime }v\eta ^{}\frac{2}{v}\eta =\frac{1}{16K_{}v}\left(2v4\mathrm{ln}v\right).\hfill \end{array}$$
(4.6)
The singularity-free solution with the correct asymptotics in the UV is uniquely determined in terms of $`\{a,K_{}\}`$. For our purposes it will be necessary to know the asymptotics of the solution near the boundary. As computed in , near the boundary $`v0`$
$$\begin{array}{cc}& \xi \frac{v}{80K_{}}+\mathrm{},\hfill \\ & \eta \frac{\mathrm{ln}v1}{8K_{}}+\mathrm{},\hfill \\ & \omega \frac{1}{30K_{}}v+\mathrm{},\hfill \\ & \mathrm{\Phi }\frac{v}{4K_{}}(\mathrm{ln}v1)+\mathrm{},\hfill \end{array}$$
(4.7)
where $`\mathrm{}`$ denote sub-dominant terms. It is straightforward to compute the Hawking temperature and the entropy density of the black brane solution (to leading order in $`P^2/K_{}`$). We find
$$\begin{array}{cc}& T=\frac{(2a)^{1/4}}{2\pi K_{}^{3/2}}\left(4K_{}P^2\right),\hfill \\ & \frac{s}{T^3}=\frac{\pi ^3K_{}}{64G_5}(K_{}+P^2).\hfill \end{array}$$
(4.8)
In order to compute the energy and the pressure of the black brane we need to relate the $`v`$ coordinate to the $`\rho `$ coordinate that we used in our solution to evaluate the stress tensor. This is done by comparing the product of warp factors in front of<sup>19</sup><sup>19</sup>19 Notice that the product of warp factors in (4.9) does not contain $`\mathrm{ln}\rho `$ factors, and thus can be consistently evaluated using the (uniformly small in $`\rho `$) $`𝒪(P^2)`$ solution (4.7) for $`\xi `$ and $`\omega `$. $`dx_0^2`$ and $`e_\psi ^2`$
$$\begin{array}{cc}\hfill h^{1/2}\rho ^2\times h^{1/2}f_2& =\left(\frac{2a}{v}\right)^{1/2}(1v)e^{2P^2(\xi 4\omega )},\hfill \\ \hfill \rho ^2\left(1+\mathrm{}\right)& =\left(\frac{2a}{v}\right)^{1/2}\left(1+\mathrm{}\right).\hfill \end{array}$$
(4.9)
Thus, to compare we need to define the radial coordinate $`\rho `$ by
$$\rho ^4\frac{v}{2a}.$$
(4.10)
Given (4.10), by translating the asymptotic form of the solution (4.7) to our ansatz (2.29)-(2.34) we find
$$\begin{array}{cc}\hfill b^{(4,0)}a^{(4,0)}=& \frac{2aP^2}{3K_{}},\hfill \\ \hfill p^{(4,0)}=& \frac{a(\mathrm{ln}(2a)1)}{2K_{}},\hfill \\ \hfill G_{00}^{(4,0)}=& \frac{P^2a}{4K_{}},\hfill \\ \hfill G_{\alpha \alpha }^{(4,0)}=& \frac{a}{4K_{}}(8K_{}+P^2).\hfill \end{array}$$
(4.11)
For the black brane geometry the renormalized one-point function of the stress tensor (3.52) is given by
$$8\pi G_5T_{ij}=\frac{1}{2}G_{ij}^{(0)}G_a^{(4,0)a}+2G_{ij}^{(4,0)}+\frac{3}{2}G_{ij}^{(0)}(b^{(4,0)}a^{(4,0)}),$$
(4.12)
where in this background $`G_{ij}^{(0)}=\eta _{ij}=\mathrm{diag}(1,1,1,1)`$ is the Minkowski metric. We are now in a position to compute the energy density $`ϵ`$ and the pressure $`𝒫`$. We find
$$\begin{array}{cc}\hfill ϵ=& \frac{a}{8\pi G_5K_{}}(P^2+3K_{}),\hfill \\ \hfill 𝒫=& \frac{a}{8\pi G_5K_{}}(K_{}P^2).\hfill \end{array}$$
(4.13)
With (4.8), (4.13) it is straightforward to verify (4.2) (to leading order in $`P^2/K_{}`$). Given (4.8) we can evaluate $`a`$ in terms of $`T`$ and $`K_{}`$. We expect that $`K_{}=K_{}(T)`$. Enforcing the first law of thermodynamics (4.3) leads to a differential equation on $`K_{}`$
$$0=2P^2T\frac{dK_{}}{dT},$$
(4.14)
which leads to (at leading order in $`P^2/K_{star}`$)
$$K_{}(T)=2P^2\mathrm{ln}\frac{T}{\mathrm{\Lambda }}$$
(4.15)
for some constant $`\mathrm{\Lambda }`$, as found from different considerations in . This allows us to write the energy density and the pressure (4.13) purely in terms of the temperature (to leading order in $`P^2/K_{}`$), and they exhibit the expected behavior of an almost-conformal theory with a number of degrees of freedom proportional to $`K_{}(T)^2`$.
The fact that we obtain a finite result for the free energy density should be useful in analyzing the deconfinement phase transition in this theory (of course, this requires going beyond the limit of $`P^2K_{}`$, but our renormalization works independently of this limit).
### 4.2 Hydrodynamics of cascading gauge theories
Small low-energy deviations from the thermodynamic equilibrium in a strongly coupled gauge theory plasma are expected to be well described by hydrodynamics. In this paper we advocated the definition of cascading gauge theories in terms of the dual string theory. Thus, the appropriate description of relaxation processes in cascading gauge theory plasma is in terms of “holographic hydrodynamics”, introduced for conformal gauge theory plasma in . The effective hydrodynamic description of relaxation of density perturbations in plasma is completely specified by two viscosity coefficients<sup>20</sup><sup>20</sup>20Not to be confused with $`\{\eta ,\xi \}`$ of the previous section., the shear viscosity $`\eta `$ and the bulk viscosity $`\xi `$, and the speed of sound waves $`c_s`$. In a conformal gauge theory plasma $`ϵ=3𝒫`$, and thus using the well-known relation
$$c_s^2=\frac{𝒫}{ϵ}$$
(4.16)
we find
$$c_s^{conformal}=\frac{1}{\sqrt{3}}.$$
(4.17)
Furthermore, conformal invariance guarantees that the bulk viscosity vanishes
$$\xi ^{conformal}=0,$$
(4.18)
while holographic hydrodynamics predicts the ratio of the shear viscosity to entropy density in the planar limit and at strong ’t Hooft coupling (namely, in the gravity approximation)<sup>21</sup><sup>21</sup>21Finite ’t Hooft coupling corrections to (4.19) were discussed in . to be
$$\frac{\eta }{s}|^{conformal}=\frac{1}{4\pi }.$$
(4.19)
It is interesting to generalize these results to non-conformal theories such as the cascading gauge theories. In particular, recall that QCD has a mass scale, and in some regimes the quark-gluon plasma of QCD could be described by strongly coupled hydrodynamics. It is thus of practical importance<sup>22</sup><sup>22</sup>22A possible application is in hydrodynamics models describing elliptic flows in heavy ion collision experiments at RHIC . to obtain hydrodynamic predictions for strongly coupled non-conformal gauge theory plasma (even though the theory we discuss here is of course very different from QCD).
Somewhat surprisingly, the ratio of shear viscosity to entropy density in all gauge theory plasmas which are dual to gravitational theories (including the cascading gauge theories), was found to be universal in the supergravity approximation,<sup>23</sup><sup>23</sup>23Of course, neither the entropy nor the shear viscosity by itself is universal. and given by (4.19). On the contrary, the speed of sound and the sound wave attenuation (which is determined in part by the bulk viscosity) are not expected to be universal. Indeed, in it was found that explicit breaking of conformal invariance in strongly coupled gauge theory plasma by fermionic (bosonic) mass terms $`m_f`$ ($`m_b`$) leads to a modified dispersion relation, with the speed of sound given in the high temperature ($`Tm_f`$, $`Tm_b`$) regime by
$$c_s=\frac{1}{\sqrt{3}}\left(1\delta _f\frac{m_f^2}{T^2}\delta _b\frac{m_b^4}{T^4}\right).$$
(4.20)
In (4.20), $`\delta _f`$ and $`\delta _b`$ are positive coefficients. The cascading gauge theory plasma differs from the non-conformal plasma discussed in in that its scale invariance is broken by a dynamically generated scale $`\mathrm{\Lambda }`$. At high temperatures the cascading gauge theory plasma is expected to resemble a conformal plasma. Indeed, using the results of the previous subsection we find that the speed of sound in this plasma is given by
$$\begin{array}{cc}\hfill c_s^2=\frac{𝒫}{ϵ}=\frac{\frac{𝒫}{T}}{\frac{ϵ}{T}}=& \frac{1}{3}\frac{4P^2}{9K_{}}+𝒪(P^4)=\frac{1}{3}\frac{2}{9\mathrm{ln}\frac{T}{\mathrm{\Lambda }}}+𝒪(P^4).\hfill \end{array}$$
(4.21)
It is amusing to note that the appearance of $`\mathrm{ln}T`$ in this correction, suggesting that the cascading gauge theory plasma has (at least some) hydrodynamic properties similar to those of weakly interacting relativistic systems (since such a correction is expected to arise in asymptotically free gauge theories at high temperature, i.e., at weak coupling).
In would be very interesting to further study the hydrodynamics of cascading gauge theory plasma, and in particular to evaluate its bulk viscosity.
## 5 Conclusions
In this paper we have performed a holographic renormalization of the cascading gauge theory of Klebanov and Tseytlin compactified on an arbitrary four-manifold, assuming that the $`_{2P}`$ global symmetry is unbroken; we have found counter-terms that can be added to the action of this theory so that the one-point functions of all operators (in the truncated action) are finite (on an arbitrary four-manifold). As discussed in detail above, the holographic renormalization in these theories is complicated by the fact that we cannot introduce arbitrary sources for the fields in the truncated action, since some of them correspond to irrelevant operators. This is an interesting problem already for asymptotically anti-de Sitter spaces, where it is also not clear how to perform a holographic renormalization for correlation functions of irrelevant operators; the difference in our case is that in the cascading theory these operators mix with the metric so we cannot consistently ignore them. Because of this complication we have performed the analysis with sources for only some of the operators, and we were not able to uniquely determine the counter-terms. However, we have proved that there exist counter-terms that make all the one-point functions finite. The choice of these counter-terms is ambiguous, and we believe that this ambiguity can be resolved (up to the usual freedom related to redefinitions of operators) by requiring that arbitrary correlation functions are finite – it would be interesting to renormalize more general correlation functions and verify that this is correct. Within our limited ansatz for the counter-terms we found that some one-point functions were uniquely determined, but others were ambiguous – we conjecture that the unambiguous one-point functions would remain the same for any consistent choice of counter-terms (including the correct one which renders all correlation functions finite), again it would be interesting to verify this <sup>24</sup><sup>24</sup>24Note that even though we discussed the ambiguity in the language of one-point functions, the ambiguity in the one-point functions that we found is independent of the state, so it may be thought of as an ambiguity in the definition of the operators themselves..
As we discussed in the introduction, our main result is that the renormalization of these theories leads to finite one-point functions despite the infinite number of high-energy degrees of freedom in these theories; we discussed in section 4 how this is consistent with the finite temperature behavior. Another possible renormalization of the cascading gauge theories involves flowing to them (at some finite scale $`\mu `$) from finite rank $`N`$ gauge theories as discussed in , assuming that the construction of is valid also at strong coupling where the gravity approximation that we have been working in is valid. It may be possible to define the confining gauge theories as a limit of the construction of in which $`N`$ and $`\mu `$ are both taken to infinity in a correlated way. Then, if one would compute correlation functions keeping the cutoff scale always above the scale $`\mu `$ one would get infinite results (say, for the conformal anomaly) because of the diverging rank of the high-energy group. However, keeping the cutoff scale above $`\mu `$ is problematic in the limit in which $`\mu `$ goes to infinity; we believe that our prescription in which the cutoff scale goes to infinity only at the end is more natural (and more analogous to the usual holographic renormalization performed in asymptotically AdS backgrounds). It would be interesting to study further various alternative renormalization schemes and to understand where they agree and where they disagree, in order to understand better how to define the cascading gauge theories. In the absence of an alternative definition for the cascading gauge theories we suggest that they can be defined in terms of their correlation functions, which one can compute using the procedure we described in this paper (at least in the gravity approximation; it should be possible to generalize this to the full string theory, but this involves understanding string theory in Ramond-Ramond backgrounds).
There are many interesting generalizations of our results. It would be interesting to analyze more general correlation functions; of course, the precise computation of higher $`n`$-point functions involves additional information beyond just the asymptotic solution that we found, but the counter-terms needed for the finiteness of these correlation functions can be found purely by using our asymptotic solutions. More precisely, $`n`$-point functions of the stress-energy tensor can be analyzed using our solution, while $`n`$-point functions of other operators require the generalization of our solution to more general sources. We expect the resulting $`2`$-point functions to agree with the results of which were computed without carefully regulating the theory. Another direction is to add additional fields to our truncated action; in particular, in order to study backgrounds like that of where the $`_{2P}`$ symmetry is spontaneously broken, we would need to add to our action fields that are charged under this symmetry. We believe that it should be possible in such backgrounds to add a finite number of fields to our effective action (2.19) and to perform the holographic renormalization as we did in this paper; it would be interesting to verify this. Finally, the asymptotic form of many other cascading backgrounds (generalizing the work of ) has recently been found , and it should be possible to generalize our results to these backgrounds.
Even without these generalizations there are several interesting applications of our results, which allow us to compute the (finite) stress-energy tensor in any solution corresponding to the compactified cascading theory which preserves the $`_{2P}`$ symmetry. One example of such a solution is the cascading theory at finite temperature, in the high temperature phase in which the $`_{2P}`$ chiral symmetry is unbroken. We have computed the thermodynamical properties of this theory at very high temperatures (compared to the strong coupling scale) in section 4 above. The solution for arbitrary temperatures is not known, but given such a solution (which can in principle be found numerically) our results allow us to precisely compute its free energy. In particular, our results would allow us to determine the temperature at which the free energy vanishes, which should be interpreted as the deconfinement temperature of the cascading gauge theory (as in ; note that in the case analyzed in a simple subtraction of the action in the two competing backgrounds was sufficient to renormalize the action, but we do not expect this to be true in more complicated backgrounds such as those of the cascading gauge theories<sup>25</sup><sup>25</sup>25It is known that background subtraction as a method for computing the free energy does not work for charged black holes in $`AdS_5`$, and for the supergravity dual to mass deformed $`𝒩=4`$ SYM theory .). Other interesting backgrounds, analyzed in , describe the confining gauge theory on $`S^3\times `$, $`S^4`$ and $`dS_4`$. Again, given any solution of this type our results allow us to precisely compute the stress-energy tensor of that solution (for example, the Casimir energy of the cascading gauge theory on $`S^3`$, which our results guarantee will be finite despite the infinite number of high-energy degrees of freedom).
## Acknowledgments
It is a pleasure to thank Vijay Balasubramanian, Marcus Berg, Nick Dorey, Eduard Gorbar, Michael Haack, Gary Horowitz, Igor Klebanov, Volodya Miransky, Rob Myers, Kostas Skenderis, Dam Son, Andrei Starinets, Matt Strassler, Marika Taylor and Arkady Tseytlin for valuable discussions. This work was supported in part by the Albert Einstein Minerva Center for Theoretical Physics. OA would like to thank the Aspen Center for Physics, the Kavli Institute for Theoretical Physics, and the Perimeter Institute for hospitality during various stages of this project. The work of OA was supported in part by the Israel-U.S. Binational Science Foundation, by the Israel Science Foundation (grant number 1399/04), by the Braun-Roger-Siegl foundation, by the European network HPRN-CT-2000-00122, by a grant from the G.I.F., the German-Israeli Foundation for Scientific Research and Development, and by Minerva. AB would like to thank the Weizmann Institute of Science, the University of Pennsylvania, the Aspen Center for Physics and the Kavli Institute for Theoretical Physics for hospitality during various stages of this project. Research at Perimeter Institute is supported in part by funds from NSERC of Canada. AB acknowledges support by an NSERC Discovery grant. AY would like to thank the Kavli Institute for Theoretical Physics for hospitality during this project. The work of AY is supported in part by a Kreitman foundation fellowship.
## Appendix A Details of the solution
### A.1 Equations of motion
In what follows we denote by a prime the derivative with respect to $`\rho `$ and by $`_i`$, $`i=0,\mathrm{},3`$, the partial derivative with respect to $`x^i`$. Also, we denote for arbitrary functions $`g_1,g_2`$ on $`_5`$
$$\begin{array}{cc}& g_1g_2G^{ij}_ig_1_jg_2,\hfill \\ & \text{ }\text{ }\text{ }g_1\frac{1}{\sqrt{G}}_i\left[\sqrt{G}G^{ij}_jg_1\right].\hfill \end{array}$$
(A.1)
Note that this is a different notation than the one we used in the beginning of section 2.2.
The equations of motion of the five dimensional supergravity (2.19) dual to the cascading gauge theory on an arbitrary curved manifold $`_5`$, in the metric parameterization (2.25), are :
$$\begin{array}{cc}\hfill 0=& \left[e^{2\mathrm{\Phi }}(G)f_2h^2\rho ^6\left(K^{}\right)^2\right]^{}+e^\mathrm{\Phi }\sqrt{G}K^{}h^1\rho ^6\{2f_2_i\left(e^\mathrm{\Phi }\sqrt{G}G^{ij}_jK\right)\hfill \\ & +e^\mathrm{\Phi }\sqrt{G}f_2K\}4P^2e^\mathrm{\Phi }(G)KK^{}h^3f_3^2\rho ^8\hfill \end{array}$$
(A.2)
$$\begin{array}{cc}\hfill 0=& \left[(G)f_2f_3^4\rho ^6\left(\mathrm{\Phi }^{}\right)^2\right]^{}+\sqrt{G}\mathrm{\Phi }^{}f_3^2\rho ^6\{2f_2_i\left(hf_3^2\sqrt{G}G^{ij}_j\mathrm{\Phi }\right)\hfill \\ & +hf_3^2\sqrt{G}f_2\mathrm{\Phi }\}+\frac{1}{2}(G)\mathrm{\Phi }^{}P^2h^1f_3^2\rho ^8\{e^\mathrm{\Phi }f_2\rho ^2((K^{})^2+h(K)^2)\hfill \\ & 4e^\mathrm{\Phi }P^4\}\hfill \end{array}$$
(A.3)
$$\begin{array}{cc}& \frac{1}{\sqrt{G}}\left[\sqrt{G}h^2\rho ^3(hf_2^2)^{}\right]^{}\frac{3}{2}f_2^1h^2\rho ^3f_2^{}(hf_2^2)^{}\rho ^3(f_2)^2\hfill \\ & +2\rho ^3f_2\text{ }\text{ }\text{ }f_2+\frac{1}{2}\rho ^3f_2h^1f_2h+f_2^2\rho ^3\text{ }\text{ }\text{ }\mathrm{ln}h\hfill \\ & =3P^2e^\mathrm{\Phi }\rho ^5h^2f_2f_3^2\rho ^5K^2h^3f_2f_3^4\hfill \\ & +\frac{1}{4}h^2\rho ^3f_2^2f_3^2P^2e^\mathrm{\Phi }\left\{(K^{})^2+h(K)^2\right\}\hfill \\ & \rho ^3h^3f_3^2\left\{(hf_2^2)^{}(hf_3^2)^{}+h(hf_2^2)(hf_3^2)\right\}+16\rho ^5h^1f_3^2f_2^3\hfill \end{array}$$
(A.4)
$$\begin{array}{cc}& \frac{1}{\sqrt{G}}\left[\sqrt{G}h^2\rho ^3(hf_3^2)^{}\right]^{}\frac{3}{2}f_3^1h^2\rho ^3f_3^{}(hf_3^2)^{}\rho ^3(f_3)^2\hfill \\ & +2\rho ^3f_3\text{ }\text{ }\text{ }f_3+\frac{1}{2}\rho ^3f_3h^1f_3h+f_3^2\rho ^3\text{ }\text{ }\text{ }\mathrm{ln}h\hfill \\ & =P^2e^\mathrm{\Phi }\rho ^5h^2f_2^1\rho ^5K^2h^3f_2^1f_3^2\frac{1}{4}h^2\rho ^3P^2e^\mathrm{\Phi }\left\{(K^{})^2+h(K)^2\right\}\hfill \\ & \frac{1}{4}\rho ^3h^3f_2^2\left\{(hf_2^2)^{}(hf_3^2)^{}+h(hf_2^2)(hf_3^2)\right\}\hfill \\ & \frac{3}{4}\rho ^3h^3f_3^2\left\{(hf_3^2)^{}(hf_3^2)^{}+h(hf_3^2)(hf_3^2)\right\}+4\rho ^5h^1(6f_32f_2)\hfill \end{array}$$
(A.5)
$$\begin{array}{cc}& R_{5ij}=G_{ij}\{\frac{1}{3}\rho ^2h^2f_2^1f_3^2P^2e^\varphi +\frac{1}{6}\rho ^2h^3f_2^1f_3^4K^2+\frac{1}{12}\rho ^3h^1f_2^2f_3^2(\hfill \\ & \frac{1}{\sqrt{G}}\left[\sqrt{G}\rho ^3f_3^6h^5[h^5f_2^2f_3^8]^{}\right]^{}+\frac{1}{\sqrt{G}}_i\left[\sqrt{G}\rho ^3f_3^6h^4G^{ij}_j[h^5f_2^2f_3^8]\right]\hfill \\ & \frac{3}{2}\rho ^3f_3^6h^5f_2^1f_2^{}[h^5f_2^2f_3^8]^{}\frac{3}{2}\rho ^3f_3^6h^4f_2^1f_2[h^5f_2^2f_3^8])\hfill \\ & 8\rho ^2h^1f_3^1+\frac{4}{3}\rho ^2h^1f_2f_3^2\}\hfill \\ & +\{_i_j[\frac{5}{4}\mathrm{ln}h+\frac{1}{2}\mathrm{ln}f_2+2\mathrm{ln}f_3]+\frac{5}{8}h^2_ih_jh\hfill \\ & +\frac{1}{8}h^1f_2^1\left(_ih_jf_2+_if_2_jh\right)+\frac{1}{2}h^1f_3^1\left(_ih_jf_3+_if_3_jh\right)\hfill \\ & G_{ij}\left(\frac{5}{16}h^2(h)^2+\frac{1}{8}h^1f_2^1hf_2+\frac{1}{2}h^1f_3^1hf_3\right)\hfill \\ & +(\frac{1}{2}h^1G_{ij}^{}\frac{1}{4}h^2G_{ij}h^{}\rho ^1h^1G_{ij})(\frac{5}{4}h^1h^{}+\frac{1}{2}f_2^1f_2^{}+2f_3^1f_3^{})\}\hfill \\ & +\{\frac{5}{16}h^2_ih_jh+\frac{1}{8}h^1f_2^1(_ih_jf_2+_if_2_jh)\hfill \\ & +\frac{1}{2}h^1f_3^1(_ih_jf_3+_if_3_jh)+\frac{1}{4}f_2^2_if_2_jf_2+f_3^2_if_3_jf_3\}\hfill \\ & +\left\{\frac{1}{4}P^2h^1f_3^2e^\mathrm{\Phi }_iK_jK+\frac{1}{2}_i\mathrm{\Phi }_j\mathrm{\Phi }\right\}\hfill \end{array}$$
(A.6)
$$\begin{array}{cc}& R_{5i\rho }=\{_i[\frac{5}{4}h^1h^{}+\frac{1}{2}f_2^1f_2^{}+2f_3^1f_3^{}]\hfill \\ & \frac{1}{2}G^{kn}G_{ni}^{}\left(\frac{5}{4}h^1_kh+\frac{1}{2}f_2^1_kf_2+2f_3^1_kf_3\right)\hfill \\ & +\left(\frac{1}{4}h^1h^{}+\rho ^1\right)\left(\frac{5}{4}h^1_ih+\frac{1}{2}f_2^1_if_2+2f_3^1_if_3\right)\hfill \\ & \frac{1}{4}h^1_ih(\frac{5}{4}h^1h^{}+\frac{1}{2}f_2^1f_2^{}+2f_3^1f_3^{})\}\hfill \\ & +\{\frac{5}{16}h^2_ihh^{}+\frac{1}{8}h^1f_2^1(_ihf_2^{}+_if_2h^{})+\frac{1}{2}h^1f_3^1(_ihf_3^{}+_if_3h^{})\hfill \\ & +\frac{1}{4}f_2^2_if_2f_2^{}+f_3^2_if_3f_3^{}\}\hfill \\ & +\left\{\frac{1}{4}h^1f_3^2P^2e^\mathrm{\Phi }_iKK^{}+\frac{1}{2}_i\mathrm{\Phi }\mathrm{\Phi }^{}\right\}\hfill \end{array}$$
(A.7)
$$\begin{array}{cc}& R_{5\rho \rho }=\{\frac{1}{3}\rho ^2h^1f_2^1f_3^2P^2e^\varphi +\frac{1}{6}\rho ^2h^2f_2^1f_3^4K^2+\frac{1}{12}\rho ^3f_2^2f_3^2(\hfill \\ & \frac{1}{\sqrt{G}}\left[\sqrt{G}\rho ^3f_3^6h^5[h^5f_2^2f_3^8]^{}\right]^{}+\frac{1}{\sqrt{G}}_i\left[\sqrt{G}\rho ^3f_3^6h^4G^{ij}_j[h^5f_2^2f_3^8]\right]\hfill \\ & \frac{3}{2}\rho ^3f_3^6h^5f_2^1f_2^{}[h^5f_2^2f_3^8]^{}\frac{3}{2}\rho ^3f_3^6h^4f_2^1f_2[h^5f_2^2f_3^8])\hfill \\ & 8\rho ^2f_3^1+\frac{4}{3}\rho ^2f_2f_3^2\}\hfill \\ & +\{[\frac{5}{4}h^1h^{}+\frac{1}{2}f_2^1f_2^{}+2f_3^1f_3^{}]^{}\hfill \\ & +\left(\frac{5}{16}h^1(h)^2+\frac{1}{8}f_2^1hf_2+\frac{1}{2}f_3^1hf_3\right)\hfill \\ & \left(\frac{5}{16}h^2(h^{})^2+\frac{1}{8}f_2^1h^1h^{}f_2^{}+\frac{1}{2}f_3^1h^1h^{}f_3^{}\right)\hfill \\ & +(\frac{5}{4}\rho ^1h^1h^{}+\frac{1}{2}f_2^1\rho ^1f_2^{}+2f_3^1\rho ^1f_3^{})\}\hfill \\ & +\left\{\left(\frac{1}{4}h^1h^{}+\frac{1}{2}f_2^1f_2^{}\right)^2+4\left(\frac{1}{4}h^1h^{}+\frac{1}{2}f_3^1f_3^{}\right)^2\right\}\hfill \\ & +\left\{\frac{1}{4}h^1f_3^2P^2e^\mathrm{\Phi }(K^{})^2+\frac{1}{2}(\mathrm{\Phi }^{})^2\right\}\hfill \end{array}$$
(A.8)
### A.2 Coefficients of the asymptotic solution
#### A.2.1 Next-to-leading order solution, $`𝒪(\rho ^2)`$
To next-to-leading order, the solution with the ansatz (2.29)-(2.34) depends on two undetermined functions, which we choose to be $`a^{(2,0)}(x)`$ and $`a^{(2,1)}(x)`$. We find that the other coefficients in the solution are then given by
$$\begin{array}{cc}\hfill G_{ij}^{(2,1)}=& \frac{1}{24}\left(12a^{(2,1)}G_{ij}^{(0)}P^2p_0\left(6R_{ij}G_{ij}^{(0)}R\right)\right)\hfill \end{array}$$
(A.9)
$$\begin{array}{cc}\hfill G_{ij}^{(2,0)}=& \frac{1}{96}\left(48a^{(2,0)}G_{ij}^{(0)}+24a^{(2,1)}G_{ij}^{(0)}+\left(2K_0+P^2p_0\right)\left(6R_{ij}G_{ij}^{(0)}R\right)\right)\hfill \end{array}$$
(A.10)
$$\begin{array}{cc}\hfill p^{(2,1)}=& 0\hfill \end{array}$$
(A.11)
$$\begin{array}{cc}\hfill p^{(2,0)}=& \frac{1}{24}P^2p_0R\hfill \end{array}$$
(A.12)
$$\begin{array}{cc}\hfill h^{(2,2)}=& P^2p_0a^{(2,1)}\hfill \end{array}$$
(A.13)
$$\begin{array}{cc}\hfill h^{(2,1)}=& P^2p_0a^{(2,0)}\frac{1}{288}\left(36\left(4K_0+P^2p_0\right)a^{(2,1)}+5P^4p_0^2R\right)\hfill \end{array}$$
(A.14)
$$\begin{array}{cc}\hfill h^{(2,0)}=\frac{1}{3456}(& 432\left(4K_0+P^2p_0\right)a^{(2,0)}\hfill \\ & +P^2p_0(216a^{(2,1)}+5(6K_0+11P^2p_0)R))\hfill \end{array}$$
(A.15)
$$\begin{array}{cc}\hfill b^{(2,1)}=& a^{(2,1)}\hfill \end{array}$$
(A.16)
$$\begin{array}{cc}\hfill b^{(2,0)}=& a^{(2,0)}\frac{1}{96}P^2p_0R\hfill \end{array}$$
(A.17)
$$\begin{array}{cc}\hfill K^{(2,1)}=& \frac{1}{24}P^2p_0\left(12a^{(2,1)}+P^2p_0R\right)\hfill \end{array}$$
(A.18)
$$\begin{array}{cc}\hfill K^{(2,0)}=& \frac{1}{96}P^2p_0\left(48a^{(2,0)}24a^{(2,1)}+\left(2K_0+7P^2p_0\right)R\right)\hfill \end{array}$$
(A.19)
where the Ricci tensor $`R_{ij}`$ and scalar $`R`$ are computed using the boundary metric $`G_{ij}^{(0)}`$.
#### A.2.2 Next-to-next-to-leading order solution, $`𝒪(\rho ^4)`$
At this order, we have several arbitrary functions, which we choose to be $`p^{(4,0)}`$, $`a^{(4,3)}`$, $`a^{(4,2)}`$, $`a^{(4,1)}`$, $`a^{(4,0)}`$, $`b^{(4,0)}`$, and $`G_{ij}^{(4,0)}\frac{1}{4}G_{ij}^{(0)}G_k^{(4,0)k}`$ (with $`G_k^{(4,0)k}G_{ij}^{(4,0)}G^{(0)ij}`$; note that $`G_k^{(4,0)k}`$ is not arbitrary but will be determined below). On the other hand, in order to find a solution we find that the Laplacians of the second order coefficients $`\{a^{(2,0)},a^{(2,1)}\}`$ are constrained to satisfy
$$\begin{array}{cc}\hfill \text{ }\text{ }\text{ }a^{(2,0)}=& \frac{1}{2160}(161P^2p_0+78K_0)\text{ }\text{ }\text{ }R,\hfill \\ \hfill \text{ }\text{ }\text{ }a^{(2,1)}=& \frac{13}{180}P^2p_0\text{ }\text{ }\text{ }R,\hfill \end{array}$$
(A.20)
where $`\text{ }\text{ }\text{ }\text{ }`$ is the Laplacian operator with the metric<sup>26</sup><sup>26</sup>26 We make a convention that $`\text{ }\text{ }\text{ }\text{ }`$ is evaluated with respect to $`G_{ij}^{(0)}(x)`$ whenever it acts on the coefficients of the perturbative solution (2.29)-(2.34), rather than on the supergravity fields. In the latter case the operator $`\text{ }\text{ }\text{ }\text{ }`$ is evaluated with respect to the full four dimensional metric $`G_{ij}(x,\rho )`$, as in the equations of motion (A.2)-(A.8). $`G_{ij}^{(0)}`$. The constant parts of $`a^{(2,0)}`$ and $`a^{(2,1)}`$ remain unfixed, as expected. The remaining functions in the solution are then given by
$$\begin{array}{cc}& G_{ij}^{(4,3)}=\frac{1}{96}P^4p_0^2\text{ }\text{ }\text{ }R_{ij}\frac{1}{288}P^4p_0^2_i_jR\frac{1}{48}P^4p_0^2R_{aijb}R^{ab}\frac{1}{144}P^4p_0^2R_{ij}R\hfill \\ & +G_{ij}^{(0)}\left(\frac{1}{576}P^4p_0^2\text{ }\text{ }\text{ }R\frac{1}{192}P^4p_0^2R_{ab}R^{ab}\frac{3}{4}a^{(4,3)}+\frac{1}{576}P^4p_0^2R^2\right)\hfill \end{array}$$
(A.21)
$$\begin{array}{cc}& G_{ij}^{(4,2)}=\frac{1}{64}P^2p_0\left(K_0+P^2p_0\right)\text{ }\text{ }\text{ }R_{ij}+\frac{1}{576}P^2p_0\left(3K_0+4P^2p_0\right)_i_jR\hfill \\ & +\frac{1}{32}P^2p_0\left(K_0+P^2p_0\right)R_{aijb}R^{ab}+\frac{1}{32}P^4p_0^2R_{ia}R_j^a\frac{1}{16}P^2p_0_i_ja^{(2,1)}\hfill \\ & +R_{ij}\left(\frac{1}{4}P^2p_0a^{(2,1)}+\frac{1}{288}P^2p_0\left(3K_02P^2p_0\right)R\right)\hfill \\ & +G_{ij}^{(0)}(\frac{1}{1152}P^2p_0(3K_0+4P^2p_0)\text{ }\text{ }\text{ }R+\frac{1}{128}P^2p_0(K_0+P^2p_0)R_{ab}R^{ab}\hfill \\ & +\frac{1}{16}\left(12a^{(4,2)}+7(a^{(2,1)})^23a^{(4,3)}\right)+\frac{1}{24}P^2p_0a^{(2,1)}R\hfill \\ & \frac{1}{2304}P^2p_0(6K_0+P^2p_0)R^2)\hfill \end{array}$$
(A.22)
$$\begin{array}{cc}& G_{ij}^{(4,1)}=\frac{1}{13824}\left(36K_{0}^{}{}_{}{}^{2}96K_0P^2p_0149P^4p_0^2\right)_i_jR\hfill \\ & +\frac{1}{512}\left(4K_{0}^{}{}_{}{}^{2}+8K_0P^2p_0+5P^4p_0^2\right)\text{ }\text{ }\text{ }R_{ij}\hfill \\ & +\frac{1}{256}\left(4K_{0}^{}{}_{}{}^{2}8K_0P^2p_05P^4p_0^2\right)R_{aijb}R^{ab}\hfill \\ & \frac{1}{64}P^2p_0\left(2K_0+P^2p_0\right)R_{ia}R_j^a\hfill \\ & +\frac{1}{64}\left(4P^2p_0_i_ja^{(2,0)}+\left(2K_0+5P^2p_0\right)_i_ja^{(2,1)}\right)\hfill \\ & +R_{ij}\left(\frac{1}{8}\left(2P^2p_0a^{(2,0)}+K_0a^{(2,1)}\right)+\frac{1}{6912}\left(36K_{0}^{}{}_{}{}^{2}+48K_0P^2p_0+79P^4p_0^2\right)R\right)\hfill \\ & +G_{ij}^{(0)}(\frac{1}{27648}(36K_{0}^{}{}_{}{}^{2}96K_0P^2p_0119P^4p_0^2)\text{ }\text{ }\text{ }R\hfill \\ & \frac{1}{3072}\left(12K_{0}^{}{}_{}{}^{2}+24K_0P^2p_0+P^4p_0^2\right)R_{ab}R^{ab}\hfill \\ & +\frac{1}{32}\left(24a^{(4,1)}4a^{(4,2)}+28a^{(2,0)}a^{(2,1)}+7(a^{(2,1)})^2+3a^{(4,3)}\right)\hfill \\ & +\frac{1}{48}(2P^2p_0a^{(2,0)}K_0a^{(2,1)})R+\frac{1}{27648}(36K_{0}^{}{}_{}{}^{2}+12K_0P^2p_085P^4p_0^2)R^2)\hfill \end{array}$$
(A.23)
$$\begin{array}{cc}& G_k^{(4,0)k}=\frac{1}{27648}P^2p_0\left(6K_0+31P^2p_0\right)\text{ }\text{ }\text{ }R\hfill \\ & +\frac{1}{3072}(12K_{0}^{}{}_{}{}^{2}+12K_0P^2p_05P^4p_0^2)R_{ab}R^{ab}+(\frac{1}{256}(448(a^{(2,0)})^264a^{(4,1)}\hfill \\ & +32a^{(4,2)}+224a^{(2,0)}a^{(2,1)}+8(a^{(2,1)})^224a^{(4,3)}768b^{(4,0)}+8K_0\text{ }\text{ }\text{ }a^{(2,0)}\hfill \\ & +8P^2p_0\text{ }\text{ }\text{ }a^{(2,0)}6K_0\text{ }\text{ }\text{ }a^{(2,1)}7P^2p_0\text{ }\text{ }\text{ }a^{(2,1)})\hfill \\ & +\frac{1}{384}\left(2\left(8K_03P^2p_0\right)a^{(2,0)}+P^2p_0a^{(2,1)}\right)R\hfill \\ & +\frac{1}{27648}\left(24K_{0}^{}{}_{}{}^{2}54K_0P^2p_037P^4p_0^2\right)R^2\hfill \end{array}$$
(A.24)
$$\begin{array}{cc}& p^{(4,3)}=\frac{1}{576}P^4p_0^2\left(\text{ }\text{ }\text{ }R+3R_{ab}R^{ab}R^2\right)\hfill \end{array}$$
(A.25)
$$\begin{array}{cc}& p^{(4,2)}=\frac{1}{768}P^2p_0\left(2K_0+11P^2p_0\right)\text{ }\text{ }\text{ }R\frac{1}{256}P^2p_0\left(2K_0+P^2p_0\right)R_{ab}R^{ab}\hfill \\ & +\frac{1}{768}P^2p_0\left(2K_0+P^2p_0\right)R^2\hfill \end{array}$$
(A.26)
$$\begin{array}{cc}& p^{(4,1)}=\frac{1}{27648}P^2p_0\left(138K_0+373P^2p_0\right)\text{ }\text{ }\text{ }R\hfill \\ & +\frac{1}{4608}P^2p_0\left(42K_019P^2p_0\right)R_{ab}R^{ab}+3\left(a^{(4,0)}b^{(4,0)}\right)\hfill \\ & +\frac{1}{384}P^2p_0\left(6a^{(2,0)}+5a^{(2,1)}\right)R+\frac{1}{27648}P^2p_0\left(102K_0+59P^2p_0\right)R^2\hfill \end{array}$$
(A.27)
$$\begin{array}{cc}& h^{(4,4)}=P^2p_0a^{(4,3)}\hfill \end{array}$$
(A.28)
$$\begin{array}{cc}& h^{(4,3)}=\frac{1}{512}P^6p_{0}^{}{}_{}{}^{3}\text{ }\text{ }\text{ }R\frac{3}{512}P^6p_{0}^{}{}_{}{}^{3}R_{ab}R^{ab}+P^2p_0a^{(4,2)}+\frac{1}{512}P^6p_{0}^{}{}_{}{}^{3}R^2\hfill \\ & +\frac{1}{16}\left(24P^2p_0(a^{(2,1)})^2\left(8K_0+3P^2p_0\right)a^{(4,3)}\right)\hfill \end{array}$$
(A.29)
$$\begin{array}{cc}& h^{(4,2)}=\frac{1}{110592}P^4p_0^2\left(324K_0+965P^2p_0\right)\text{ }\text{ }\text{ }R+\frac{25}{576}P^4p_0^2a^{(2,1)}R\hfill \\ & +\frac{1}{36864}P^4p_0^2\left(324K_0+23P^2p_0\right)R_{ab}R^{ab}+\frac{1}{12288}P^4p_0^2\left(36K_0+5P^2p_0\right)R^2\hfill \\ & +\frac{1}{64}(64P^2p_0a^{(4,1)}4(8K_0+3P^2p_0)a^{(4,2)}192P^2p_0a^{(2,0)}a^{(2,1)}\hfill \\ & +48K_0(a^{(2,1)})^2+5P^2p_0(a^{(2,1)})^23P^2p_0a^{(4,3)})\hfill \end{array}$$
(A.30)
$$\begin{array}{cc}& h^{(4,1)}=\frac{1}{442368}P^2p_0\left(408K_{0}^{}{}_{}{}^{2}+2140K_0P^2p_0+2259P^4p_0^2\right)\text{ }\text{ }\text{ }R\hfill \\ & +\frac{1}{147456}P^2p_0\left(408K_{0}^{}{}_{}{}^{2}412K_0P^2p_0+11P^4p_0^2\right)R_{ab}R^{ab}\hfill \\ & +\frac{1}{128}(192P^2p_0(a^{(2,0)})^232P^2p_0a^{(4,0)}64K_0a^{(4,1)}24P^2p_0a^{(4,1)}\hfill \\ & 4P^2p_0a^{(4,2)}+4\left(48K_0+5P^2p_0\right)a^{(2,0)}a^{(2,1)}+19P^2p_0(a^{(2,1)})^2+3P^2p_0a^{(4,3)}\hfill \\ & +160P^2p_0b^{(4,0)})\frac{1}{13824}5P^2p_0(138P^2p_0a^{(2,0)}+(60K_0+101P^2p_0)a^{(2,1)})R\hfill \\ & +\frac{1}{1327104}P^2p_0\left(1224K_{0}^{}{}_{}{}^{2}+780K_0P^2p_01057P^4p_0^2\right)R^2\hfill \end{array}$$
(A.31)
$$\begin{array}{cc}& h^{(4,0)}=\frac{1}{3538944}P^2p_0\left(480K_{0}^{}{}_{}{}^{2}+3244K_0P^2p_0+4657P^4p_0^2\right)\text{ }\text{ }\text{ }R\hfill \\ & +\frac{1}{1179648}P^2p_0\left(600K_{0}^{}{}_{}{}^{2}+908K_0P^2p_059P^4p_0^2\right)R_{ab}R^{ab}\hfill \\ & +\frac{1}{512}(8(48K_0+5P^2p_0)(a^{(2,0)})^2+8(2K_0+13P^2p_0)a^{(4,0)}8P^2p_0a^{(4,1)}\hfill \\ & +4P^2p_0a^{(4,2)}+76P^2p_0a^{(2,0)}a^{(2,1)}3P^2p_0(a^{(2,1)})^23P^2p_0a^{(4,3)}\hfill \\ & 272K_0b^{(4,0)}200P^2p_0b^{(4,0)}32P^2p_0p^{(4,0)})\hfill \\ & +\frac{1}{221184}P^2p_0\left(2\left(2706K_0+4001P^2p_0\right)a^{(2,0)}+\left(174K_0+31P^2p_0\right)a^{(2,1)}\right)R\hfill \\ & +\frac{1}{10616832}P^2p_0\left(1392K_{0}^{}{}_{}{}^{2}+220K_0P^2p_0+5231P^4p_0^2\right)R^2\hfill \end{array}$$
(A.32)
$$b^{(4,3)}=a^{(4,3)}$$
(A.33)
$$b^{(4,2)}=\frac{1}{576}P^4p_0^2\text{ }\text{ }\text{ }\text{ }R\frac{1}{192}P^4p_0^2R_{ab}R^{ab}+a^{(4,2)}+\frac{1}{576}P^4p_0^2R^2$$
(A.34)
$$\begin{array}{cc}& b^{(4,1)}=\frac{1}{6912}P^2p_0\left(12K_0+43P^2p_0\right)\text{ }\text{ }\text{ }R+\frac{1}{576}P^2p_0\left(3K_02P^2p_0\right)R_{ab}R^{ab}\hfill \\ & +a^{(4,1)}\frac{1}{192}P^2p_0a^{(2,1)}R+\frac{1}{6912}P^2p_0\left(12K_0+11P^2p_0\right)R^2\hfill \end{array}$$
(A.35)
$$K^{(4,3)}=\frac{1}{288}P^6p_{0}^{}{}_{}{}^{3}\text{ }\text{ }\text{ }\text{ }R\frac{1}{96}P^6p_{0}^{}{}_{}{}^{3}R_{ab}R^{ab}+\frac{1}{4}P^2p_0a^{(4,3)}+\frac{1}{288}P^6p_{0}^{}{}_{}{}^{3}R^2$$
(A.36)
$$\begin{array}{cc}& K^{(4,2)}=\frac{1}{192}P^4p_0^2\left(K_0+4P^2p_0\right)\text{ }\text{ }\text{ }R+\frac{1}{256}P^4p_0^2\left(4K_0+3P^2p_0\right)R_{ab}R^{ab}\hfill \\ & +\frac{1}{16}P^2p_0\left(4a^{(4,2)}3\left((a^{(2,1)})^2+a^{(4,3)}\right)\right)+\frac{1}{48}P^4p_0^2a^{(2,1)}R\hfill \\ & \frac{1}{384}P^4p_0^2\left(2K_0+P^2p_0\right)R^2\hfill \end{array}$$
(A.37)
$$\begin{array}{cc}& K^{(4,1)}=\frac{1}{13824}P^2p_0\left(18K_{0}^{}{}_{}{}^{2}+159K_0P^2p_0+308P^4p_0^2\right)\text{ }\text{ }\text{ }R\hfill \\ & \frac{1}{9216}P^2p_0\left(36K_{0}^{}{}_{}{}^{2}+156K_0P^2p_0+7P^4p_0^2\right)R_{ab}R^{ab}\hfill \\ & +\frac{1}{32}P^2p_0(96a^{(4,0)}+8a^{(4,1)}4a^{(4,2)}12a^{(2,0)}a^{(2,1)}+3(a^{(2,1)})^2+3a^{(4,3)}\hfill \\ & +96b^{(4,0)})+\frac{1}{384}P^2p_0(14P^2p_0a^{(2,0)}(4K_0+11P^2p_0)a^{(2,1)})R\hfill \\ & +\frac{1}{13824}P^2p_0\left(18K_{0}^{}{}_{}{}^{2}+69K_0P^2p_029P^4p_0^2\right)R^2\hfill \end{array}$$
(A.38)
$$\begin{array}{cc}& K^{(4,0)}=\frac{1}{221184}P^2p_0\left(132K_{0}^{}{}_{}{}^{2}+944K_0P^2p_0+1371P^4p_0^2\right)\text{ }\text{ }\text{ }R\hfill \\ & +\frac{1}{36864}P^2p_0\left(84K_{0}^{}{}_{}{}^{2}+124K_0P^2p_021P^4p_0^2\right)R_{ab}R^{ab}\hfill \\ & +\frac{1}{128}(24P^2p_0(a^{(2,0)})^2+16(6K_0+11P^2p_0)a^{(4,0)}8P^2p_0a^{(4,1)}+4P^2p_0a^{(4,2)}\hfill \\ & +12P^2p_0a^{(2,0)}a^{(2,1)}3P^2p_0(a^{(2,1)})^23P^2p_0a^{(4,3)}96K_0b^{(4,0)}144P^2p_0b^{(4,0)}\hfill \\ & 64P^2p_0p^{(4,0)})\frac{1}{3072}P^2p_0(2K_0+3P^2p_0\left)\right(22a^{(2,0)}5a^{(2,1)})R\hfill \\ & +\frac{1}{663552}P^2p_0\left(396K_{0}^{}{}_{}{}^{2}252K_0P^2p_0+713P^4p_0^2\right)R^2\hfill \end{array}$$
(A.39)
where all the differential operators are evaluated with respect to $`G_{ij}^{(0)}`$.
### A.3 Symmetries of the asymptotic solution
There are three symmetries we can use to check our general solution to (A.2)-(A.8).
First, there is a scaling symmetry
$$\begin{array}{cc}& \rho \lambda \rho ,G_{ij}\lambda ^2G_{ij},hh\hfill \\ & KK,\mathrm{\Phi }\mathrm{\Phi },f_2f_2,f_3f_3\hfill \end{array}$$
(A.40)
which translates into the following scaling symmetry of the parameters of the asymptotic solution (2.29)-(2.34)
$$p_0p_0$$
(A.41)
$$K_0K_0+2P^2p_0\mathrm{ln}\lambda $$
(A.42)
$$G_{ij}^{(0)}\lambda ^2G_{ij}^{(0)}$$
(A.43)
$$a^{(2,1)}\frac{1}{\lambda ^2}a^{(2,1)}$$
(A.44)
$$a^{(2,0)}\frac{1}{\lambda ^2}\left(a^{(2,0)}a^{(2,1)}\mathrm{ln}\lambda \right)$$
(A.45)
$$a^{(4,3)}\frac{1}{\lambda ^4}a^{(4,3)}$$
(A.46)
$$a^{(4,2)}\frac{1}{\lambda ^4}\left(a^{(4,2)}3a^{(4,3)}\mathrm{ln}\lambda \right)$$
(A.47)
$$a^{(4,1)}\frac{1}{\lambda ^4}\left(a^{(4,1)}2a^{(4,2)}\mathrm{ln}\lambda +3a^{(4,3)}\mathrm{ln}^2\lambda \right)$$
(A.48)
$$a^{(4,0)}\frac{1}{\lambda ^4}\left(a^{(4,0)}\mathrm{ln}\lambda \left(a^{(4,1)}+\mathrm{ln}\lambda \left(a^{(4,2)}+a^{(4,3)}\mathrm{ln}\lambda \right)\right)\right)$$
(A.49)
$$\begin{array}{cc}\hfill b^{(4,0)}& \frac{1}{\lambda ^4}b^{(4,0)}\frac{1}{\lambda ^4}a^{(4,1)}\mathrm{ln}\lambda \frac{\mathrm{ln}\lambda }{576\lambda ^4}P^2p_0\left(3K_02P^2p_0+3P^2p_0\mathrm{ln}\lambda \right)R_{ab}R^{ab}\hfill \\ & \frac{\mathrm{ln}\lambda }{6912\lambda ^4}P^2p_0\left(12K_0+43P^2p_0+12P^2p_0\mathrm{ln}\lambda \right)\text{ }\text{ }\text{ }R\hfill \\ & +\frac{\mathrm{ln}\lambda ^2}{\lambda ^4}\left(a^{(4,2)}a^{(4,3)}\mathrm{ln}\lambda \right)+\frac{\mathrm{ln}\lambda }{192\lambda ^4}P^2p_0a^{(2,1)}R\hfill \\ & +\frac{\mathrm{ln}\lambda }{6912\lambda ^4}P^2p_0\left(12K_011P^2p_0+12P^2p_0\mathrm{ln}\lambda \right)R^2\hfill \end{array}$$
(A.50)
$$\begin{array}{cc}\hfill G_{ij}^{(4,0)}& \frac{1}{\lambda ^2}G_{ij}^{(4,0)}+\frac{\mathrm{ln}\lambda }{\lambda ^2}(\frac{1}{16}P^2p_0_i_ja^{(2,0)}+\frac{1}{64}(2K_05P^2p_0)_i_ja^{(2,1)}\hfill \\ & +\frac{1}{512}\left(4K_{0}^{}{}_{}{}^{2}8K_0P^2p_05P^4p_0^2\right)\text{ }\text{ }\text{ }R_{ij}+\frac{1}{8}\left(2P^2p_0a^{(2,0)}K_0a^{(2,1)}\right)R_{ij}\hfill \\ & +\frac{1}{27648}\left(36K_{0}^{}{}_{}{}^{2}+96K_0P^2p_0+119P^4p_0^2\right)G_{ij}^{(0)}\text{ }\text{ }\text{ }R\hfill \\ & +\frac{1}{13824}\left(36K_{0}^{}{}_{}{}^{2}+96K_0P^2p_0+149P^4p_0^2\right)_i_jR\hfill \\ & +\frac{1}{3072}\left(12K_{0}^{}{}_{}{}^{2}+24K_0P^2p_0+P^4p_0^2\right)G_{ij}^{(0)}R_{ab}R^{ab}\hfill \\ & +\frac{1}{256}\left(4K_{0}^{}{}_{}{}^{2}+8K_0P^2p_0+5P^4p_0^2\right)R_{iabj}R^{ab}+\frac{1}{64}P^2p_0\left(2K_0+P^2p_0\right)R_{ia}R_j^a\hfill \\ & +\frac{1}{27648}\left(36K_{0}^{}{}_{}{}^{2}12K_0P^2p_0+85P^4p_0^2\right)G_{ij}^{(0)}R^2\hfill \\ & +R(\frac{1}{6912}(36K_{0}^{}{}_{}{}^{2}48K_0P^2p_079P^4p_0^2)R_{ij}\hfill \\ & +\frac{1}{48}G_{ij}^{(0)}(2P^2p_0a^{(2,0)}+K_0a^{(2,1)}))+\frac{1}{32}(24a^{(4,1)}+4a^{(4,2)}28a^{(2,0)}a^{(2,1)}\hfill \\ & 7a_{}^{(2,1)}{}_{}{}^{2}3a^{(4,3)})G_{ij}^{(0)})+\frac{\mathrm{ln}^2\lambda }{\lambda ^2}(\frac{1}{16}P^2p_0_i_ja^{(2,1)}\hfill \\ & \frac{1}{64}P^2p_0\left(K_0+P^2p_0\right)\text{ }\text{ }\text{ }R_{ij}+\frac{1}{576}P^2p_0\left(3K_0+4P^2p_0\right)\hfill \\ & +\frac{1}{1152}P^2p_0\left(3K_0+4P^2p_0\right)G_{ij}^{(0)}\text{ }\text{ }\text{ }R+\frac{1}{128}P^2p_0\left(K_0+P^2p_0\right)G_{ij}^{(0)}R_{ab}R^{ab}\hfill \\ & +\frac{1}{32}P^2p_0\left(K_0+P^2p_0\right)R_{iabj}R^{ab}+\frac{1}{32}P^4p_0^2R_{ia}R_j^a\hfill \\ & \frac{1}{2304}P^2p_0\left(6K_0+P^2p_0\right)G_{ij}^{(0)}R^2\frac{1}{4}P^2p_0a^{(2,1)}R_{ij}\hfill \\ & +R\left(\frac{1}{288}P^2p_0\left(3K_02P^2p_0\right)R_{ij}+\frac{1}{24}P^2p_0a^{(2,1)}G_{ij}^{(0)}\right)\hfill \\ & +\frac{1}{16}(12a^{(4,2)}+7a_{}^{(2,1)}{}_{}{}^{2}3a^{(4,3)})G_{ij}^{(0)})+\frac{\mathrm{ln}^3\lambda }{\lambda ^2}(\frac{1}{96}P^4p_0^2\text{ }\text{ }\text{ }R_{ij}\hfill \\ & +\frac{1}{288}P^4p_0^2_i_jR+\frac{1}{576}G_{ij}^{(0)}P^4p_0^2\text{ }\text{ }\text{ }R+\frac{1}{192}P^4p_0^2G_{ij}^{(0)}R_{ab}R^{ab}\hfill \\ & +\frac{1}{48}P^4p_0^2R_{iabj}R^{ab}+\frac{1}{144}P^4p_0^2RR_{ij}\frac{1}{576}P^4p_0^2G_{ij}^{(0)}R^2+\frac{3}{4}a^{(4,3)}G_{ij}^{(0)})\hfill \end{array}$$
(A.51)
$$\begin{array}{cc}\hfill p^{(4,0)}& \frac{\mathrm{ln}\lambda }{4608\lambda ^4}(42K_019P^2p_0+18(2K_0+P^2p_0)\mathrm{ln}\lambda \hfill \\ & +24P^2p_0\mathrm{ln}^2\lambda )P^2p_0R_{ab}R^{ab}\frac{\mathrm{ln}\lambda }{27648\lambda ^4}(138K_0+373P^2p_0\hfill \\ & +36(2K_0+11P^2p_0)\mathrm{ln}\lambda +48P^2p_0\mathrm{ln}\lambda ^2)P^2p_0\text{ }\text{ }\text{ }R+\frac{p^{(4,0)}}{\lambda ^4}\hfill \\ & +\frac{3\mathrm{ln}\lambda }{\lambda ^4}\left(b^{(4,0)}a^{(4,0)}\right)+\frac{\mathrm{ln}\lambda }{\lambda ^4}\left(\frac{1}{64}P^2p_0a^{(2,0)}\frac{5}{384}P^2p_0a^{(2,1)}\right)R\hfill \\ & +\frac{\mathrm{ln}\lambda }{27648\lambda ^4}(102K_059P^2p_0+36(2K_0+P^2p_0)\mathrm{ln}\lambda \hfill \\ & +48P^2p_0\mathrm{ln}^2\lambda )P^2p_0R^2\hfill \end{array}$$
(A.52)
The solution transforms covariantly under this transformation, and in addition, the constraints (A.20) and (A.24) are also invariant under the symmetry transformations (A.40).
Another symmetry is the residual gauge freedom associated with unfixed diffeomorphisms
$$\begin{array}{cc}\hfill \rho & \widehat{\rho }\rho \left[1+\rho ^2\left(\delta _{20}+\delta _{21}\mathrm{ln}\rho \right)+\rho ^4\left(\delta _{40}+\delta _{41}\mathrm{ln}\rho +\delta _{42}\mathrm{ln}^2\rho +\delta _{43}\mathrm{ln}^3\rho \right)\right]\hfill \\ \hfill K& K,\mathrm{\Phi }\mathrm{\Phi },hh\left(\frac{\widehat{\rho }}{\rho }\right)^4\left(\frac{\widehat{\rho }}{\rho }\right)^4,f_2f_2\left(\frac{\widehat{\rho }}{\rho }\right)^2\left(\frac{\widehat{\rho }}{\rho }\right)^2\hfill \\ \hfill f_3& f_3\left(\frac{\widehat{\rho }}{\rho }\right)^2\left(\frac{\widehat{\rho }}{\rho }\right)^2,G_{ij}G_{ij}\left(\frac{\widehat{\rho }}{\rho }\right)^4\left(\frac{\widehat{\rho }}{\rho }\right)^2\hfill \end{array}$$
(A.53)
where $`\{\delta _{20},\mathrm{},\delta _{43}\}`$ are arbitrary constants. This symmetry is realized on the parameters of the asymptotic solution as follows
$`p_0`$ $`p_0`$ (A.54)
$`K_0`$ $`K_0`$ (A.55)
$`G_{ij}^{(0)}`$ $`G_{ij}^{(0)}`$ (A.56)
$`a^{(2,0)}`$ $`4\delta _{20}+2\delta _{21}+a^{(2,0)}`$ (A.57)
$`a^{(2,1)}`$ $`4\delta _{21}+a^{(2,1)}`$ (A.58)
$`a^{(4,3)}`$ $`8\delta _{43}+a^{(4,3)}`$ (A.59)
$`a^{(4,2)}`$ $`8\delta _{21}^{}{}_{}{}^{2}+8\delta _{42}+6\delta _{43}+a^{(4,2)}+2\delta _{21}a^{(2,1)}`$ (A.60)
$`a^{(4,1)}`$ $`16\delta _{20}\delta _{21}6\delta _{21}^{}{}_{}{}^{2}+8\delta _{41}+4\delta _{42}+2\delta _{21}a^{(2,0)}+a^{(4,1)}+2\delta _{20}a^{(2,1)}+\delta _{21}a^{(2,1)}`$ (A.61)
$`a^{(4,0)}`$ $`8\delta _{20}^{}{}_{}{}^{2}6\delta _{20}\delta _{21}+\delta _{21}^{}{}_{}{}^{2}+8\delta _{40}+2\delta _{41}+2\left(\delta _{20}+\delta _{21}\right)a^{(2,0)}+a^{(4,0)}\delta _{20}a^{(2,1)}`$ (A.62)
$`b^{(4,0)}`$ $`8\delta _{20}^{}{}_{}{}^{2}6\delta _{20}\delta _{21}+\delta _{21}^{}{}_{}{}^{2}+8\delta _{40}+2\delta _{41}+2\left(\delta _{20}+\delta _{21}\right)a^{(2,0)}+b^{(4,0)}\delta _{20}a^{(2,1)}`$
$`{\displaystyle \frac{1}{48}}P^2p_0\left(\delta _{20}+\delta _{21}\right)R`$ (A.63)
$`p^{(4,0)}`$ $`p^{(4,0)}+{\displaystyle \frac{1}{12}}P^2p_0\delta _{20}R`$ (A.64)
$$\begin{array}{cc}\hfill G_{ij}^{(4,0)}& G_{ij}^{(4,0)}+\frac{1}{8}\left(P^2p_0\delta _{21}+2K_0\left(2\delta _{20}+\delta _{21}\right)\right)R_{ij}\hfill \\ & +G_{ij}^{(0)}(\frac{1}{48}(P^2p_0\delta _{21}+2K_0(2\delta _{20}+\delta _{21}))R+(2\delta _{20}+\delta _{21})a^{(2,0)}\hfill \\ & +13\delta _{20}^{}{}_{}{}^{2}+16\delta _{20}\delta _{21}+3\delta _{21}^{}{}_{}{}^{2}6\delta _{40}2\delta _{41}+\frac{1}{2}(3\delta _{20}+\delta _{21})a^{(2,1)})\hfill \end{array}$$
(A.65)
Again, in addition to the solution transforming covariantly, the constraints (A.20) and (A.24) are also invariant under the symmetry transformations (A.53).
Finally, there is a symmetry which rescales the action (and thus is a symmetry of any solution to the equations of motion) given by (2.36), which in terms of our variables is given by
$$\begin{array}{cc}& G_{ij}\beta G_{ij},h\beta h,e^\mathrm{\Phi }\beta e^\mathrm{\Phi },\hfill \\ & K\beta K,f_2f_2,f_3f_3.\hfill \end{array}$$
(A.66)
This symmetry is realized on the parameters of the asymptotic solution as follows
$$p_0\beta p_0,K_0\beta K_0,G_{ij}^{(0)}\beta G_{ij}^{(0)},G_{ij}^{(4,0)}\beta G_{ij}^{(4,0)},$$
(A.67)
with all other parameters remaining unchanged. Again, in addition to the solution transforming covariantly, the constraints (A.20) and (A.24) are also invariant under the symmetry transformations (A.66). The symmetry transformation (2.36) is very useful in constraining the possible counter-terms in the holographic renormalization of the cascading gauge theories, as discussed in section 3.3.
### A.4 Ambiguities of the minimal subtraction
In this appendix we discuss a specific simple ambiguity of the counter-term action, which is present in our minimal subtraction ansatz.
Consider a counter-term ansatz of the form
$$\delta ^{counter}=\delta _1(K,P^2e^\mathrm{\Phi })\left(_\gamma ^23_{ab\gamma }_\gamma ^{ab}\right)+\delta _2(K,P^2e^\mathrm{\Phi })\underset{\gamma }{\text{ }\text{ }\text{ }\text{ }}.$$
(A.68)
This ansatz was chosen so that it does not contribute to $`T_i^i`$, and it is easy to verify that (since $`\alpha _8`$ is a source only for $`h`$) it also does not contribute to $`𝒪_8`$.
It is straightforward to verify that (A.68) leads to
$$\delta 𝒪_{p_0}=\frac{\delta _1}{p_0}\left(R^23R_{ab}R^{ab}\right)+\frac{\delta _2}{p_0}\text{ }\text{ }\text{ }\text{ }R,$$
(A.69)
$$\delta 𝒪_{K_0}=\frac{\delta _1}{K_0}\left(R^23R_{ab}R^{ab}\right)+\frac{\delta _2}{K_0}\text{ }\text{ }\text{ }\text{ }R.$$
(A.70)
From (A.69), (A.70) it follows that $`\delta 𝒪_{p_0}`$ and $`\delta 𝒪_{K_0}`$ will be finite if the $`\delta _i`$ are finite in the limit $`\rho 0`$. Under the symmetry (2.36) we must have the scaling
$$\delta _i\beta ^2\delta _i.$$
(A.71)
The requirement of finiteness in the $`\rho 0`$ limit, along with (A.71), lead to the following choice of counter-terms :
$$\begin{array}{cc}\hfill \delta _1=& \kappa _1\left(K+2P^2e^\mathrm{\Phi }\mathrm{ln}\rho \right)^2+\kappa _2\left(K+2P^2e^\mathrm{\Phi }\mathrm{ln}\rho \right)P^2e^\mathrm{\Phi }+\kappa _3P^4e^{2\mathrm{\Phi }},\hfill \\ \hfill \delta _2=& \kappa _4\left(K+2P^2e^\mathrm{\Phi }\mathrm{ln}\rho \right)^2+\kappa _5\left(K+2P^2e^\mathrm{\Phi }\mathrm{ln}\rho \right)P^2e^\mathrm{\Phi }+\kappa _6P^4e^{2\mathrm{\Phi }}.\hfill \end{array}$$
(A.72)
The arguments above show that the only possible divergences arising from the counter-terms (A.72) are in $`𝒪_6^s`$. We find that these counter-terms indeed lead to extra divergences in $`𝒪_6^s`$, but they can be removed (preserving what has been achieved thus far) by
$$\delta ^{counter}\delta ^{counter}+\delta _{1extra}\left(_\gamma ^23_{ab\gamma }_\gamma ^{ab}\right)+\delta _{2extra}\underset{\gamma }{\text{ }\text{ }\text{ }\text{ }},$$
(A.73)
with
$$\begin{array}{cc}\hfill \delta _{1extra}=& X_a\left(\frac{4}{5}\kappa _1\left(K+2P^2e^\mathrm{\Phi }\mathrm{ln}\rho \right)P^2e^\mathrm{\Phi }+\frac{2}{5}\kappa _2P^4e^{2\mathrm{\Phi }}\right),\hfill \\ \hfill \delta _{2extra}=& X_a\left(\frac{4}{5}\kappa _4\left(K+2P^2e^\mathrm{\Phi }\mathrm{ln}\rho \right)P^2e^\mathrm{\Phi }+\frac{2}{5}\kappa _5P^4e^{2\mathrm{\Phi }}\right).\hfill \end{array}$$
(A.74)
Thus, we can always add to our action the counter-terms (A.72) and (A.73) without generating any divergences.
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# On ∗-representations of a certain class of algebras related to a graph
## Introduction
Let $`H`$ be a separable Hilbert space, and let $`A_1`$, …, $`A_n`$ be a family of self-adjoint operators in $`H`$ with fixed finite spectra, such that $`A_1+\mathrm{}+A_n=\lambda I`$ for some $`\lambda `$. Such family of operators can be treated as a representation of certain $``$-algebra with finite number of generators and polynomial relations. The corresponding algebras were studied in a number of recent papers (see \[VMS05\] and the bibliography therein). The interest to such algebras is due to their relations with deformed preprojective algebras (see, e.g., \[CBH98\]), representations of quivers (see \[KR05\]), Horn problem (see \[KPS05\]), integral operators (see \[Vas98\]) etc.
Considering a family of self-adjoint operators with a pre-defined spectra, for which $`A_1+\mathrm{}+A_n=\lambda I`$ we can assume that $`\lambda >0`$, and
$$\sigma (A_l)M_l=\{0=\alpha _0^{(l)}<\alpha _1^{(l)}<\mathrm{}<\alpha _{k_l}^{(l)}\},l=1,\mathrm{},n.$$
To study such properties of operators, it is convenient to consider them as $``$-representations of a certain $``$-algebra. Following \[VMS05\] (see also \[MSZ04\] and references therein) consider a simply-laced non-oriented graph $`\mathrm{\Gamma }`$ consisting of $`n`$ branches, such that $`l`$-th branch has $`k_l+1`$ vertices, $`l=1`$, …, $`n`$, and all branches are connected at a single root vertex. Marking the vertices of $`l`$-th branch (excluding the root vertex) by positive numbers $`(\alpha _j^{(l)})_{j=1}^{k_l}`$ increasing to the root, we get a function
$$\chi =(\alpha _1^{(1)},\mathrm{},\alpha _{k_1}^{(1)};\mathrm{};\alpha _1^{(n)},\mathrm{},\alpha _{k_n}^{(n)})$$
on the graph $`\mathrm{\Gamma }`$ defined in all veritces except the root (below this function will be called a character on $`\mathrm{\Gamma }`$). The root vertex will be marked by the number $`\lambda `$ (notice that the term character in many papers is used to denote a function $`(\chi ,\lambda )`$ on the whole graph, including the root vertex. For our needs the given notation is more convenient).
Given a graph $`\mathrm{\Gamma }`$, a character $`\chi `$ on $`\mathrm{\Gamma }`$ and a positive number $`\lambda `$, on can construct the following $``$-algebra
$$𝒜_{\mathrm{\Gamma },\chi ,\lambda }=a_l=a_l^{},l=1,\mathrm{},np_l(a_l)=0,l=1,\mathrm{},n;\underset{l=1}{\overset{n}{}}a_l=\lambda e,$$
where $`p_l(x)=x(x\alpha _1^{(l)})\mathrm{}(x\alpha _{k_l}^{(l)})`$, $`k=1`$, …, $`n`$. Then the family $`A_1`$, …, $`A_n`$ is a $``$-representation of $`𝒜_{\mathrm{\Gamma },\chi ,\lambda }`$.
Properties of the algebra $`𝒜_{\mathrm{\Gamma },\chi ,\lambda }`$, in particular, the structure of its $``$-reprsentations, crucially depend on the type of the graph $`\mathrm{\Gamma }`$: they are quite different for the cases where $`\mathrm{\Gamma }`$ is a Dynkin graph, extended Dynkin graph or none of them.
If $`\mathrm{\Gamma }`$ is a Dynkin graph then the corresponding algebra is finite-dimensional and therefore has only finite number of irreducible $``$-representations, and all of them are finite-dimensional. In this case, sets of parameters $`(\chi ,\lambda )`$, for which there exist representations, and the $``$-representations themselves are described in \[KPS05, SZ05\].
If $`\mathrm{\Gamma }`$ is an extended Dynkin graph, then the algebra is infinite-dimensional of polynomial growth \[VMS05\]. In this case, there exists a special character $`\chi _\mathrm{\Gamma }`$ on the corresponding extended Dynkin graph $`\mathrm{\Gamma }`$. For such special characters, sets $`\mathrm{\Sigma }_{\mathrm{\Gamma },\chi _\mathrm{\Gamma }}`$ of those $`\lambda `$, for which there exist representations of the corresponding algebra and the corresponding $``$-representations were were studied in \[OS99, MRS04, MSZ04\] etc. In particular, it follows that for the special characters, all irreducible $``$-representations are finite-dimensional.
In Section 3 we prove that all irreducible $``$-representations of $`𝒜_{\mathrm{\Gamma },\chi ,\lambda }`$ are finite-dimensional for any $`\chi `$ and $`\gamma `$ provided that $`\mathrm{\Gamma }`$ is an extended Dynkin graph. As a corollary we get that for $`\lambda \omega (\chi )`$ the rigidity index for irreducible representation is equal to 2.
To study the properties of $``$-representations of $`𝒜_{\mathrm{\Gamma },\chi ,\lambda }`$ in the case of arbitrary characters and more general graphs we introduce the notion of invariant functional on the set of characters. We show that such invariant functional is unique if and only if the graph is an extended Dynkin graph; for Dynkin graph there are no invariant functionals, and for other graphs there are exactly two invariant functionals (Section 2). These invariant functionals are described in terms of solutions of the equation
$$ns=\underset{l=1}{\overset{n}{}}\frac{1}{1+(s1)+\mathrm{}+(s1)^{k_l}},s1.$$
In Section 1 we show that this equation has unique solution $`s=2`$ if and only if $`\mathrm{\Gamma }`$ is an extended Dynkin graph. For ordinary Dynkin graph this equation has no solutions $`s1`$, and for other graphs it has two solutions $`1<s_1<2<s_2<n`$.
The main tool to study $``$-representations of $`𝒜_{\mathrm{\Gamma },\chi ,\lambda }`$ are reflection (Coxeter) functors introduced in \[Kru02\] (for non-involutive case, see \[GP71\]). Namely there exist two functors, $`S:\mathrm{Rep}𝒜_{\mathrm{\Gamma },\chi ,\lambda }\mathrm{Rep}𝒜_{\mathrm{\Gamma },\chi ^{},\lambda ^{}}`$ and $`T:\mathrm{Rep}𝒜_{\mathrm{\Gamma },\chi ,\lambda }\mathrm{Rep}𝒜_{\mathrm{\Gamma },\chi ^{\prime \prime },\lambda }`$, where
$`\chi ^{}`$ $`=(\alpha _{k_l}^{(1)}\alpha _{k_l1}^{(1)},\mathrm{},\alpha _{k_l}^{(1)}\alpha _0^{(1)};\mathrm{};\alpha _{k_l}^{(n)}\alpha _{k_l1}^{(n)},\mathrm{},\alpha _{k_l}^{(n)}\alpha _0^{(n)}),`$
$`\lambda ^{}`$ $`=\alpha _{k_1}^{(1)}+\mathrm{}+\alpha _{k_n}^{(n)}\lambda ,`$
$`\chi ^{\prime \prime }`$ $`=(\lambda \alpha _{k_1}^{(1)},\mathrm{},\lambda \alpha _1^{(1)};\mathrm{};\lambda \alpha _{k_n}^{(n)},\mathrm{},\lambda \alpha _1^{(n)}).`$
The action of these functors on $``$-representations gives rise to the action on pairs, $`S:(\chi ,\lambda )(\chi ^{},\lambda ^{})`$, $`T:(\chi ,\lambda )(\chi ^{\prime \prime },\lambda )`$. This action is extensively used below.
## 1. Equation which distinguishes the type of a graph
Let $`n`$ and $`k_l`$, $`l=1`$, …, $`n`$ be given natural numbers, and let $`\mathrm{\Gamma }`$ be the corresponding graph. In what follows, we will use solutions of the following equation
(1)
$$ns=\underset{l=1}{\overset{n}{}}\frac{1}{1+(s1)+\mathrm{}+(s1)^{k_l}},s1.$$
###### Theorem 1.
The equation (1) have no solutions on $`[1,\mathrm{})`$ if and only if the corresponding graph $`\mathrm{\Gamma }`$ is one of the Dynkin graphs $`A_d`$, $`d1`$, $`D_d`$, $`d4`$, $`E_6`$, $`E_7`$, or $`E_8`$. The equation (1) has a unique solution on $`[1,\mathrm{})`$ if and only if the corresponding graph $`\mathrm{\Gamma }`$ is one of the extended Dynkin graphs $`\stackrel{~}{D}_4`$, $`\stackrel{~}{E}_6`$, $`\stackrel{~}{E}_7`$, or $`\stackrel{~}{E}_8`$, this solution is $`s=2`$. In all other cases (i.e., where $`\mathrm{\Gamma }`$ is neither a Dynkin graph nor an extended Dynkin graph), the equation (1) has on $`[1,\mathrm{})`$ two solutions $`1<s_1<2<s_2<n`$.
###### Proof.
Consider auxiliary functions $`\varphi _k(x)=(1+x+\mathrm{}+x^k)^1`$, so that (1) takes the form
$$f_\mathrm{\Gamma }(s)=ns\underset{l=1}{\overset{n}{}}\varphi _{k_l}(s1)=0.$$
Calculations give the following formula
$$\varphi _k^{\prime \prime }(x)=(k+1)x^{k1}\left(kx^{k1}+2(k1)x^{k2}+\mathrm{}+(k1)2x+k\right)\varphi _k^3(x),$$
which implies that $`f_\mathrm{\Gamma }^{\prime \prime }(s)<0`$, $`s>1`$.
For each of the extended Dynkin graphs, a direct calculation shows that $`f_\mathrm{\Gamma }(2)=0`$ and $`f_\mathrm{\Gamma }^{}(2)=0`$, which gives the uniqueness of the solution for the case of extended Dynkin graphs.
Each of the Dynkin graphs $`D_4`$, $`E_6`$, $`E_7`$, $`E_8`$ is a subgraph of the corresponding extended Dynkin graph, which gives inequalities for the corresponding functions, $`f_{D_4}(s)<f_{\stackrel{~}{D}_4}(s)`$, $`f_{E_6}(s)<f_{\stackrel{~}{E}_6}(s)`$, $`f_{E_7}(s)<f_{\stackrel{~}{E}_7}(s)`$, $`f_{E_8}(s)<f_{\stackrel{~}{E}_7}(s)`$ for $`s1`$. For the graphs $`A_n`$, $`n1`$, and $`D_n`$, $`n>4`$, one checks directly the inequality $`f_\mathrm{\Gamma }(s)<0`$, $`s0`$, where $`\mathrm{\Gamma }`$ is one of these graphs.
Let $`\mathrm{\Gamma }`$ be neither Dynkin graph, nor extended Dynkin graph. Then $`\mathrm{\Gamma }`$ contains subgraph $`\mathrm{\Gamma }_0`$ which is an extended Dynkin graph, and therefore, the inequality $`f_\mathrm{\Gamma }(s)>f_{\mathrm{\Gamma }_0}(s)`$ holds for $`s1`$, in particular, $`f_\mathrm{\Gamma }(2)>0`$. To complete the proof notice that $`f_\mathrm{\Gamma }(1)=1`$ and $`f_\mathrm{\Gamma }(n)<0`$, and $`f^{\prime \prime }(s)<0`$ guarantees that there are only two solutions on the interval $`[1,\mathrm{})`$. ∎
## 2. Invariant functionals on graphs
Let $`\mathrm{\Gamma }`$ be a graph formed by $`n`$ branches connected in a single root vertex, and let $`k_l`$ be the number of vertices (excluding root) in $`l`$-th branch. Let $`\chi `$ be a character on the graph $`\mathrm{\Gamma }`$,
(2)
$$\chi =(\alpha _1^{(1)},\mathrm{},\alpha _{k_1}^{(1)};\mathrm{};\alpha _1^{(n)},\mathrm{},\alpha _{k_1}^{(n)}),$$
$$0<\alpha _1^{(l)}<\mathrm{}<\alpha _{k_l}^{(l)},l=1,\mathrm{},n.$$
Let $`\omega ()`$ be a linear functional, which takes non-negative values on characters.
###### Definition 1.
We say that $`\omega ()`$ is invariant with respect to the functor $`TS`$, if
$$TS(\chi ,\omega (\chi ))=(\stackrel{~}{\chi },\omega (\stackrel{~}{\chi }))$$
for any character $`\chi `$ on $`\mathrm{\Gamma }`$.
###### Theorem 2.
If $`\mathrm{\Gamma }`$ is a Dynkin graph (one of $`D_d`$, $`d4`$, $`E_6`$, $`E_7`$ or $`E_8`$), then there are no invariant functionals on the set of its characters.
If $`\mathrm{\Gamma }`$ is an extended Dynkin Graph (one of $`\stackrel{~}{D}_4`$, $`\stackrel{~}{E}_6`$, $`\stackrel{~}{E}_7`$ or $`\stackrel{~}{E}_8`$), then there exists unique invariant functional:
— for $`\stackrel{~}{D}_4`$, $`\omega (\alpha ;\beta ;\gamma ;\delta )=\frac{1}{2}(\alpha +\beta +\gamma +\delta )`$;
— for $`\stackrel{~}{E}_6`$, $`\omega (\alpha _1,\alpha _2;\beta _1,\beta _2;\gamma _1,\gamma _2)=\frac{1}{3}(\alpha _1+\alpha _2+\beta _1+\beta _2+\gamma _1+\gamma _2)`$;
— For $`\stackrel{~}{E}_7`$, $`\omega (\alpha _1,\alpha _2,\alpha _3;\beta _1,\beta _2,\beta _3;\gamma )=\frac{1}{4}(\alpha _1+\alpha _2+\alpha _3+\beta _1+\beta _2+\beta _3+2\gamma )`$;
— for $`\stackrel{~}{E}_8`$, $`\omega (\alpha _1,\alpha _2,\alpha _3,\alpha _4,\alpha _5;\beta _1,\beta _2;\gamma )=\frac{1}{6}(\alpha _1+\alpha _2+\alpha _3+\alpha _4+\alpha _5+2\beta _1+2\beta _2+3\gamma )`$.
In the case where $`\mathrm{\Gamma }`$ is neither Dynkin graph nor extended Dynkin graph, there exist two $`TS`$-invariant functionals. They are given by
(3)
$$\omega (\chi )=\underset{l=1}{\overset{n}{}}\underset{j=1}{\overset{k_l}{}}a_j^{(l)}\alpha _j^{(l)},a_j^{(l)}0,j=1,\mathrm{},k_l;l=1,\mathrm{},n.$$
with
(4)
$$a_j^{(l)}=\frac{(s1)^j}{1+(s1)+\mathrm{}+(s1)^{k_l}},j=1,\mathrm{},k_l;l=1,\mathrm{},n,$$
where $`s`$ is a solution of (1).
###### Proof.
Let $`\chi `$ be given by (2), then linear functional $`\omega ()`$ can be represented as (3). Then $`(TS)(\chi ,\omega (\chi ))=(\stackrel{~}{\chi },\stackrel{~}{\lambda })`$, where
$`\stackrel{~}{\chi }`$ $`=({\displaystyle \underset{l=1}{\overset{n}{}}}\alpha _{k_l}^{(l)}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \underset{j=1}{\overset{k_l}{}}}a_j^{(l)}\alpha _j^{(l)}\alpha _{k_1}^{(1)},`$
$`{\displaystyle \underset{l=1}{\overset{n}{}}}\alpha _{k_l}^{(l)}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \underset{j=1}{\overset{k_l}{}}}a_j^{(l)}\alpha _j^{(l)}\alpha _{k_1}^{(1)}+\alpha _1^{(1)},\mathrm{},`$
$`{\displaystyle \underset{l=1}{\overset{n}{}}}\alpha _{k_l}^{(l)}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \underset{j=1}{\overset{k_l}{}}}a_j^{(l)}\alpha _j^{(l)}\alpha _{k_1}^{(1)}+\alpha _{k_11}^{(1)};\mathrm{}),`$
$`\stackrel{~}{\lambda }`$ $`={\displaystyle \underset{l=1}{\overset{n}{}}}\alpha _{k_l}^{(l)}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \underset{j=1}{\overset{k_l}{}}}a_j^{(l)}\alpha _j^{(l)}.`$
Then
$`\omega (\stackrel{~}{\chi })\stackrel{~}{\lambda }`$ $`={\displaystyle \underset{l=1}{\overset{n}{}}}(\alpha _1^{(l)}((1s)a_1^{(l)}+a_2^{(l)})+\mathrm{}+\alpha _{k_l1}^{(l)}((1s)a_{k_l1}^{(l)}+a_{k_l}^{(l)})`$
$`+\alpha _{k_l}^l(ss_l1+(1s)a_{k_l}^l))`$
where $`s=_{l=1}^ns_l`$, $`s_l=_{p=1}^{k_l}a_p^{(l)}`$. Since the latter should be zero for arbitrary choice of $`\chi `$, we get the following conditions
(5)
$$a_{k_l}^{(l)}=(s1)a_{k_l1}^{(l)},\mathrm{},a_2^{(l)}=(s1)a_1^{(l)},$$
(6)
$$a_{k_l}^{(l)}=1\frac{s_l}{s1},l=1,\mathrm{},n,$$
For $`s_l`$, $`l=1`$, …, $`n`$ from (5) we have
(7)
$$s_l=a_1^{(l)}+\mathrm{}+a_{k_l}^{(l)}=a_1^{(l)}(1+(s1)+\mathrm{}+(s1)^{k_l1})$$
Substitute in (6) $`a_{k_l}^{(l)}=(s1)^{k_l1}a_1^{(l)}`$, then $`(s1)^{k_l}a_1^{(l)}=s1s_l`$, or
(8)
$$a_1^{(l)}=\frac{s1}{1+(s1)+\mathrm{}+(s1)^{k_l}}$$
Compare this with (7) we get
$$s_l=1\frac{1}{1+(s1)+\mathrm{}+(s1)^{k_l}}$$
And taking into account that $`s=_{l=1}^ns_l`$ we finally get a conditions (1) for $`s`$. Solutions of this equations are described by Theorem 1. For such $`s`$, using (5), (8), we get expression (4) for $`a_j^{(l)}`$. ∎
## 3. Irreducible representations
Let $`\mathrm{\Gamma }`$ be an extended Dynkin graph, and let $`\chi `$ be a character on it. The main result on representations of the corresponding algebra $`A_{\mathrm{\Gamma },\chi ,\lambda }`$ is the following.
###### Theorem 3.
All irreducible families of operators corresponding to extended Dynkin diagrams are finite-dimensional.
###### Proof.
Let $`\pi `$ be an irreducible representation of the algebra $`A_{\mathrm{\Gamma },\chi ,\gamma }`$, where $`\mathrm{\Gamma }`$ is an extended Dynkin graph. We consider two cases.
1. Let $`\lambda =\omega (\chi )`$. It is shown in \[Mel05, VMS05\] that the corresponding algebra is finite-dimensional over its center, and therefore, is a PI-algebra. This implies that the dimensions of all its irreducible representations are bounded.
2. Let $`\lambda <\omega (\chi )`$. We proceed as follows. We apply the $`(ST)^n`$ functors to the representation of the algebra corresponding to the pair $`(\chi ,\lambda )`$ to get representations of the algebras corresponding to other pairs $`(\chi _n,\lambda _n)`$ and show that at some step either there cannot exist representation (in this case, there are no representations for $`(\chi ,\lambda )`$), or the representation is an obvious extension of representation of a subgraph (such subgraph is a Dynkin diagram, and the corresponding algebra is finite dimensional, therefore it has a finite number of representations all of which are finite-dimensional). Then the initial representation $`\pi `$ is obtained from some finitely dimensional representation of $`A_{\mathrm{\Gamma },\chi _n,\lambda _n}`$ as a result of applying of the $`(TS)^n`$ functor, and therefore, is finite dimensional as well.
Notice first that if some of the coefficients of $`\chi `$ is greater than $`\lambda `$, then the corresponding projection is zero. Indeed, let $`P_j`$, $`j=1`$,…,$`m`$ be orthogonal projections such that $`_{j=1}^m\alpha _jP_j=\lambda I`$, where $`\alpha _j>0`$, $`j=1`$,…,$`m`$, and $`\alpha _k>\lambda `$ for some $`k`$. Then $`_{jk}\alpha _jP_kP_jP_k=(\lambda \alpha _k)P_k`$ which is possible only for $`P_k=0`$ since the left-hand side of the latter equality is non-negative, but the right-hand side is non-positive.
In the case where some of the coefficients of $`\chi `$ are equal to $`\lambda `$, the corresponding projection commutes with all other projections and therefore, is either identity (in this case all other projections are zero), or zero.
Thus, in both these cases, the representation is in fact a representation of subalgebra in $`A_{\mathrm{\Gamma },\chi ,\lambda }`$ corresponding to some subgraph.
Now let all coefficients of $`\chi _k`$, $`kn`$ be positive, and some coefficient of $`\chi _{n+1}`$ be negative or zero. Taking into account the way the functors act on characters, we easily see that this means that the corresponding coefficient of $`\chi _n`$ is grater or equal that $`\lambda _n`$.
To complete the proof, it is now sufficient to show that for any $`\lambda <\omega (\chi )`$ there exists such $`n`$, that some coefficient of $`\chi _n`$ is negative or zero.
Let $`\chi _\mathrm{\Gamma }`$ be the special character of the corresponding graph, and let $`\omega _\mathrm{\Gamma }=\omega (\chi _\mathrm{\Gamma })`$. If we norm $`\chi `$ such that $`\omega (\chi )=\omega _\mathrm{\Gamma }`$, then the character can be represented as $`\chi =\chi _\mathrm{\Gamma }+\stackrel{~}{\chi }`$, where $`\stackrel{~}{\chi }`$ is (not necessarily positive) character, such that $`\omega (\stackrel{~}{\chi })=0`$. Also represent $`\lambda =\omega _\mathrm{\Gamma }\lambda `$. That is:
— for $`\stackrel{~}{D}_4`$ $`\chi =(1+a_1;1+a_2;1+a_3;1+a_4)`$, $`\lambda =2\gamma `$, $`a_1+a_2+a_3+a_4=0`$;
— for $`\stackrel{~}{E}_6`$ $`\chi =(1+a_1,2+a_2;1+b_1,2+b_2;1+c_1,2+c_2)`$, $`\lambda =3\gamma `$, $`a_1+a_2+b_1+b_2+c_1+c_2=0`$;
— for $`\stackrel{~}{E}_7`$ $`\chi =(1+a_1,2+a_2,3+a_3;1+b_1,2+b_2,3+b_3;2+c_1)`$, $`\lambda =4\gamma `$, $`a_1+a_2+a_3+b_1+b_2+b_3+2c_1=0`$ ;
— for $`\stackrel{~}{E}_8`$ $`\chi =(1+a_1,2+a_2,3+a_3,4+a_4,5+a_5;2+b_1,4+b_2;3+c_1)`$, $`\lambda =6\gamma `$, $`a_a+a_2+a_3+a_4+a_5+2b_1+2b_2+3c_1=0`$.
Using this notation, one can directly check that for any $`k=1`$, 2, …,
— for $`\stackrel{~}{D}_4`$ $`(ST)^{2k}(\chi ,\lambda )=(\chi 2k\gamma \chi _\mathrm{\Gamma },\lambda 4k\gamma )`$;
— for $`\stackrel{~}{E}_6`$ $`(ST)^{6k}(\chi ,\lambda )=(\chi 3k\gamma \chi _\mathrm{\Gamma },\lambda 9k\gamma )`$;
— for $`\stackrel{~}{E}_7`$ $`(ST)^{12k}(\chi ,\lambda )=(\chi 4k\gamma \chi _\mathrm{\Gamma },\lambda 16k\gamma )`$;
— for $`\stackrel{~}{E}_8`$ $`(ST)^{30k}(\chi ,\lambda )=(\chi 6k\gamma \chi _\mathrm{\Gamma },\lambda 36k\gamma )`$.
Here $`\chi _\mathrm{\Gamma }`$ is the special character for the corresponding graph. The relations listed above complete the proof in the case $`\lambda <\omega (\chi )`$.
3. Let $`\lambda >\omega (\chi )`$. Apply the $`S`$ functor to the pair $`(\chi ,\lambda ))`$, then we get a pair $`(\chi ^{},\lambda ^{})`$ with $`\lambda ^{}<\omega (\chi ^{})`$ and the arguments above apply. ∎
Recall that the rigidity index \[Kat96, SV99\] of a family of operators $`A_j`$, $`j=1`$, …, $`k`$ in $`n`$-dimensional space is
$$r=n^2(2k)+\underset{j=1}{\overset{k}{}}c(A_j)$$
where $`c(A_j)`$ is the dimension of a centralizer of $`A_j`$.
###### Corollary 1.
The rigidity index is equal to 2 for all irreducible representations of $`A_{\mathrm{\Gamma },\chi ,\gamma }`$ where $`\mathrm{\Gamma }`$ is extended Dynkin graph and $`\lambda \omega (\chi )`$.
###### Proof.
Indeed, rigidity index is preserved by $`T`$ and $`S`$ functors. From the proof above it follows that any irreducible representation can be obtained from one-dimensional ones under the action of the Coxeter functors (in fact, it was shown that any irreducible representation is obtained from representation of some algebra related to ordinary Dynkin graph, but this process can be iterated for the subalgebra until we finish in one-dimensional space). Now the result follows from the directly verified fact that for one-dimensional representations $`r=2`$. ∎
The author expresses his gratitude to Professor Yu. S. Samoilenko for kind attention and fruitful discussions of the topics discussed in this paper.
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# Twisting Null Geodesic Congruences, Scri, H-Space and Spin-Angular Momentum
## 1 Introduction
When discussing or studying the general properties of a set of differential equations, it is usually best to try to adopt a coordinate system to the symmetries of the manifold in question. For example general properties of Maxwell’s equations are best studied using standard Minkowski coordinates. If however one is looking for solutions with specific properties, symmetries for example, it is clearly best to adopt coordinates closest to those of the solutions. Spherically symmetric solutions are most easily studied in spherical coordinates. In general relativity, when studying general properties of asymptotically flat space-times the most used coordinate/tetrad system, in the neighborhood of future null infinity, $`^+`$, are the Bondi coordinate/tetrad system which is associated with the asymptotic symmetry group, i.e., with the Bondi-Metzner-Sachs group. On the other hand, if particular asymptotically flat solutions are to be found, special non-Bondi coordinates and/or tetrads are often used. One sees this, for example in the study of the Robinson-Trautman metrics or the twisting type II metrics where both different coordinates and different tetrad systems are often used.
In this work we will show that for every specific asymptotically flat space-time there is a non-Bondi coordinate/tetrad system made-to-order for that solution that simplifies the discussion of the solution and brings out its properties more clearly than does the Bondi system. Often one can use the Bondi coordinates with the non-Bondi tetrad, though for certain situations, discussed later, both the non-Bondi coordinates and tetrads are preferred
To approach this issue, in Sec.II, we give a brief review and survey of the relevant properties of $`^+`$ that are needed here. This includes the possible use of two different types of coordinates, Bondi and NU coordinates, and two different types of null tetrad systems, (1) those based on a surface forming \[twist-free\] null vector field near ‘infinity’ and (2) those based on a twisting null vector field. In Sec.III, we discuss what are the different relevant functions that are defined and used on $`^+,`$ (for any given space-time), that are used in the different $`^+`$ coordinate/tetrad systems. It will be seen that depending on the different choices of coordinate/tetrad system, the descriptions of the same physical situation can be very different and that large simplifications are available by the appropriate choices. In Sec.IV we will display the Bianchi Identities for the different functions in the different coordinate and tetrad systems; all quite different from each other. In Sec.V we discuss how one tries to extract physical quantities from the variables on $`^+.`$
The major development reported in this work is given in Sec.VI. We show that for any asymptotically flat space-time there is a coordinate/tetrad (non-Bondi) system that can be adapted to the specific space-time that allows a considerable simplification in the relevant equations. Furthermore it allows us to give a definition of angular momentum (orbital and spin) that is geometric and hence independent of the choice of coordinates and thus is supertranslation invariant. The point of view that is adopted is as follows: an arbitrary asymptotically flat space-time is first completely described in a Bondi coordinate/tetrad frame by the freely chosen characteristic data, the asymptotic shear, $`\sigma ,`$ and some (irrelevant for us) initial data. {We have tacitly assumed throughout that the space-time can be described or approximated by real analytic functions.} Keeping the Bondi coordinates, we now change the tetrad, in a specific way: we go from the Bondi tetrad which has non-vanishing asymptotic shear (our free data $`\sigma `$) to a new tetrad that has (in general) twist but is asymptotically shear free. All the information that was in the shear is thus put into the geometry of the twisting shear free tetrad. The relevant function to describe the twisting tetrad is a complex (stereographic) angle at each point of $`^+`$ that is denoted by $`L[^+].`$ This process is reversible. Given $`L,`$ the $`\sigma `$ can be reconstructed. The new feature is that the shear-free twisting congruence, given by $`L,`$ is not uniquely determined by the original shear. The remaining freedom in $`L`$ is given by four complex functions of a single complex variable which can be thought of, locally, as a (complex) curve, $`\xi ^a(\tau ),`$ in C$`^4.`$ \[Specifically and more surprising, it turns out to be a curve in $`H`$ -Space.\] At this point it is a completely arbitrary curve, i.e., any choice of the curve is related to the same shear. We show that there is a natural condition that can be imposed on the curve that makes it unique (up to initial data) that leads us to define the spin angular momentum and the center of mass from the real and imaginary parts of the complex curve. These definitions are geometrically invariant under Bondi supertranslation and reduce to the standard definitions for the Kerr metric. When the asymptotic shear is pure ‘electric type’ the spin vanishes, i.e., the spin arises from the ‘magnetic’ part of the shear. Unfortunately to do this construction exactly is quite difficult and one must resort to approximations to obtain approximate explicit results. Fortunately there is a straightforward approximation method which is summarized in an appendix.
In Sec.VII we give and discuss several examples. They include the Robinson-Trautman metrics, the type II twisting metrics and space-times obtained from the class of shears referred to as the Tod-Sparling shears.
Though everything that is said and done here could have been stated in the language of the conformal compactification of space-time, we will retain the physical space-time description. From this point of view, the null boundary is reached by letting our null geodesics go to future infinity, i.e., by letting the geodesic affine parameters, $`r`$, go to infinity.
We assume that the reader is familiar with spin-coefficient formalism and the associated spin-$`s`$ harmonics though for completeness we have included in Appendix C several of the useful relations.
In addition for comparison with other works, we mention the awkward convention adopted here for historic reasons and from frequent usage, that our radial coordinate $`r`$ differs from the more conventional choice of $`r_c,`$ i.e., by $`r=\sqrt{2}r_c`$ and our $`u_B=u_c/\sqrt{2}.`$ This leads, in the Schwarzschild case, to the unconventional form
$`g_{00}du_B^2`$ $`=`$ $`2(1{\displaystyle \frac{2\sqrt{2}GM}{c^2r}})du_B^2`$ (1)
$`\psi _2^0`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}GM}{c^2}}.`$
## 2 Review of Scri<sup>+</sup>: the Future Null Boundary of Space-Time
### 2.1 Scri coordinates
Future null infinity, often referred to as $`^+,`$ is roughly speaking the set of endpoints of all future directed null geodesics. It is usually given the structure of $`S^2xR,`$ a line bundle over the sphere. This leads naturally to the global stereographic coordinates, ($`\zeta ,\overline{\zeta },),`$ for the $`S^2`$ and the $`u`$ labeling the cross sections or global slices of the bundle. From the point of view of the space-time, $`^+`$ is a null surface and the null generators are given by ($`\zeta ,\overline{\zeta },)`$ $`=constant`$.
There is a canonical slicing of $`^+`$, $`u_B=constant`$, referred to as the Bondi slicing, that it is defined from an asymptotic symmetry inherited from the interior, called the Bondi-Metzner-Sachs group. Any one Bondi slicing is related to any other by the supertranslation freedom;
$$\stackrel{~}{u}_B=u_B+\alpha (\zeta ,\overline{\zeta }).$$
(2)
with $`\alpha (\zeta `$,$`\overline{\zeta })`$ an arbitrary regular function on $`S^2.`$
All other arbitrary slicings, ($`u=\tau ),`$ often called NU slicings, are given by
$$u_B=G(\tau ,\zeta ,\overline{\zeta }).$$
(3)
### 2.2 Null Tetrads on Scri
In addition to the choice of coordinate systems there is the freedom to chose the asymptotic null tetrad system, ($`l^a,m^a,\overline{m}^a,n^a`$). Since $`^+`$ is a fixed null surface with its own generators, with null tangent vectors, \[say $`n^a],`$ it is natural to keep it fixed. The remaining tetrad freedom lies, essentially, in the choice of the other null leg, $`l^a.`$ Most often $`l^a`$ is chosen to be orthogonal to the Bondi, $`u_B=constant,`$ slices. In this case we will refer to a Bondi coordinate/tetrad system. Relative to such a system any other tetrad set, ($`l^a,m^a,n^a`$), is given by a null rotation about $`n^a,`$ i.e.,
$`l^a`$ $`=`$ $`l^a+b\overline{m}^a+\overline{b}m^a+b\overline{b}n^a,`$ (4)
$`m^a`$ $`=`$ $`m^a+bn^a,`$
$`n^a`$ $`=`$ $`n^a,`$
$`b`$ $`=`$ $`L/r+O(r^2).`$
The complex function $`L(u_B,\zeta ,\overline{\zeta },)`$, the (stereographic) angle between the null vectors, $`l^a`$ and $`l^a,`$ is at this moment completely arbitrary but later will be dynamically determined. Depending on how $`L(u_B,\zeta ,\overline{\zeta },)`$ is chosen $`l^a`$ might or might not be surface forming. We will abuse the language/notation and refer to the tetrad systems associated with $`l^a`$ as the twisting-type tetrads as distinct from a Bondi tetrad even when $`l^a`$ is surface forming.
The function $`L(u_B,\zeta ,\overline{\zeta },)`$ and its choice will later play the pivotal role in this work. It will be seen that the vector field $`l^a`$ can be constructed by an appropriately chosen $`L(u,\zeta ,\overline{\zeta },)`$ so that it is asymptotically shear-free.
## 3 Quantities Defined on $`^+`$ for Any Given Interior Space-Time
Using the spin-coefficient notation for the five complex components of the asymptotic Weyl tensor, we have, from the peeling theorem, that
$`\psi _0`$ $`=`$ $`{\displaystyle \frac{\psi _0^0}{r^5}}+0(r^6)`$
$`\psi _1`$ $`=`$ $`{\displaystyle \frac{\psi _1^0}{r^4}}+0(r^5)`$
$`\psi _2`$ $`=`$ $`{\displaystyle \frac{\psi _2^0}{r^3}}+0(r^4)`$
$`\psi _3`$ $`=`$ $`{\displaystyle \frac{\psi _3^0}{r^2}}+0(r^3)`$
$`\psi _4`$ $`=`$ $`{\displaystyle \frac{\psi _4^0}{r}}+0(r^2)`$
where the five, $`\psi _4^0,\psi _3^0,\psi _2^0,\psi _1^0,\psi _0^0,`$ are functions defined on $`^+.`$ For a given space-time their explicit expressions depend on the choices of both coordinates and tetrads. Their evolution is determined by the asymptotic Bianchi identities and the type of characteristic data that is given.
We adopt the following notation: expression written in the Bondi coordinate/tetrad system will appear without a star,‘’, e.g., $`\psi _2^0,`$ while for Bondi coordinates with a twisting tetrad a ‘’ will be used, e.g., $`\psi _2^0`$. In the case of NU coordinates and tetrad a double star will be used, e.g. $`\psi _2^0,`$etc.
Using Bondi coordinates and tetrad, the free characteristic data is given only by the (complex) asymptotic shear,
$$\sigma =\sigma (u_B,\zeta ,\overline{\zeta }),$$
while in Bondi coordinates, but with a twisting-type tetrad the free functions are
$$\sigma ^{}=\sigma ^{}(u_B,\zeta ,\overline{\zeta })\text{ \& }L(u_B,\zeta ,\overline{\zeta })$$
with $`\sigma ^{}(u_B,\zeta ,\overline{\zeta })`$ and $`L(u_B,\zeta ,\overline{\zeta })`$ carrying the same (redundant) information as did the Bondi $`\sigma (u_B,\zeta ,\overline{\zeta }).`$
In the case of NU coordinates and tetrad, the free data is given by
$$\sigma ^{}=\sigma ^{}(\tau ,\zeta ,\overline{\zeta })\text{ \& }V(\tau ,\zeta ,\overline{\zeta })\text{ }$$
with
$$V(\tau ,\zeta ,\overline{\zeta })du_B/d\tau =_\tau G(\tau ,\zeta ,\overline{\zeta }).$$
The important point is that in both cases, the Bondi coordinates with twisting tetrad and the NU coordinates and tetrad, all the information that was in $`\sigma (u,\zeta ,\overline{\zeta })`$ is transferred to the new variables, ($`\sigma ^{}`$& $`L)`$ or ($`\sigma ^{}`$& $`V)`$. Later we will show that all the information can be shifted into an appropriately chosen $`L(u_B,\zeta ,\overline{\zeta })`$ with a vanishing $`\sigma ^{}`$ and that in certain special cases (a pure ‘electric’ type $`\sigma `$) all the information can be shifted into the $`V(\tau ,\zeta ,\overline{\zeta })`$ with vanishing $`\sigma ^{}.`$
The relationship between the ($`\psi _4^0,\psi _3^0,\psi _2^0,\psi _1^0,\psi _0^0)`$ given in a Bondi tetrad and the ($`\psi _4^0,\psi _3^0,\psi _2^0,\psi _1^0,\psi _0^0)`$ of a twisting tetrad is
$`\psi _0^0`$ $`=`$ $`\psi _0^0+4L\psi _1^0+6L^2\psi _2^0+4L^3\psi _3^0+L^4\psi _4^0`$ (5)
$`\psi _1^0`$ $`=`$ $`\psi _1^0+3L\psi _2^0+3L^2\psi _3^0+L^3\psi _4^0`$ (6)
$`\psi _2^0`$ $`=`$ $`\psi _2^0+2L\psi _3^0+L^2\psi _4^0`$ (7)
$`\psi _3^0`$ $`=`$ $`\psi _3^0+L\psi _4^0`$ (8)
$`\psi _4^0`$ $`=`$ $`\psi _4^0.`$ (9)
This is used, in the next section, going from the Bondi version of the Bianchi identities to the twisting version. To go to the conventional NU tetrad a further scaling is needed.
## 4 The Asymptotic Bianchi Identities
In each of the different coordinate/tetrad cases there are relations and differential equations relating the functions defined on $`^+`$. Specifically, the $`\psi _3^0`$ & $`\psi _4^0`$, in each of the cases, are given explicitly in terms of the functions, $`\sigma `$, $`L`$ or $`V.`$ The dynamics lies in the differential equations for the remaining $`\psi ^{}s,.i.e.,`$ ($`\psi _0^0,\psi _1^0,\psi _2^0`$).
Explicitly, in a Bondi coordinate and tetrad system, they are:
$`\psi _0^0`$ $`=`$ $`\text{ }𝔡\psi _1^0+3\sigma \psi _2^0`$ (10)
$`\psi _1^0`$ $`=`$ $`\text{ }𝔡\psi _2^0+2\sigma \psi _3^0`$ (11)
$`\psi _2^0`$ $`=`$ $`\text{ }𝔡\psi _3^0+\sigma \psi _4^0`$ (12)
$`\psi _3^0`$ $`=`$ $`𝔡\overline{\sigma }^{}`$ (13)
$`\psi _4^0`$ $`=`$ $`\overline{\sigma }^{}`$ (14)
$`\psi _2^0\overline{\psi }_2^0`$ $`=`$ $`\overline{𝔡}^2\sigma \text{ }𝔡^2\overline{\sigma }+\overline{\sigma }\sigma ^{}\sigma \overline{\sigma }^{},`$ (15)
$`\mathrm{\Psi }`$ $`=`$ $`\psi _2^0+\text{ }𝔡^2\overline{\sigma }+\sigma \overline{\sigma }^{}`$ (16)
$`\mathrm{\Psi }\overline{\mathrm{\Psi }}`$ $`=`$ $`0`$ (17)
For later use we display the same set of equations in both a Bondi-coordinate/twisting Tetrad system and in NU-coordinate/tetrad system but where the simplification
$$\sigma ^{}=0\text{ and }\sigma ^{}=0$$
was used. The justification for this simplification is given in the next section.
###### Remark 1
For a Bondi Coordinate/Twisting Tetrad system, with $`\sigma ^{}=0,`$ we have, from Eq.(5), that
$`\psi _0^0`$ $`=`$ $`\text{ }𝔡\psi _1^0L\psi _1^04\dot{L}\psi _1^0`$ (18)
$`\psi _1^0`$ $`=`$ $`\text{ }𝔡\psi _2^0L\psi _2^03\dot{L}\psi _2^0`$ (19)
$`\psi _2^0`$ $`=`$ $`\text{ }𝔡\psi _3^0L\psi _3^02\dot{L}\psi _3^0`$ (20)
$`\psi _3^0`$ $`=`$ $`𝔡\overline{\sigma }^{}+L\overline{\sigma }^{}`$ (21)
$`\psi _4^0`$ $`=`$ $`\overline{\sigma }^{}`$ (22)
$`\sigma `$ $``$ $`𝔡L+{\displaystyle \frac{1}{2}}(L^2)^{}`$ (23)
$`\mathrm{\Psi }`$ $`=`$ $`\psi _2^0+2L\text{ }𝔡\overline{\sigma }^{}+L^2\overline{\sigma }^{}+\text{ }𝔡^2\overline{\sigma }+\sigma \overline{\sigma }^{}`$ (24)
$`\mathrm{\Psi }\overline{\mathrm{\Psi }}`$ $`=`$ $`0`$ (25)
###### Remark 2
Eqs. (16) and (24) are identical even though they appear different. This can be seen by using Eq.(7).
For a NU coordinate/tetrad system, with $`\sigma ^{}=0,`$ after rescaling the $`\psi `$$`s`$ with appropriate powers of $`V(\tau ,\zeta ,\overline{\zeta })`$,
$`\psi _0^0`$ $`=`$ $`\text{ }𝔡\psi _1^0+3{\displaystyle \frac{V^{}}{V}}\psi _0^0`$ (26)
$`\psi _1^0`$ $`=`$ $`\text{ }𝔡\psi _2^0+3{\displaystyle \frac{V^{}}{V}}\psi _1^0`$ (27)
$`\psi _2^0`$ $`=`$ $`\text{ }𝔡\psi _3^0+3{\displaystyle \frac{V^{}}{V}}\psi _2^0`$ (28)
$`\psi _3^0`$ $`=`$ $`\overline{𝔡}^2\text{ }𝔡\mathrm{log}V`$ (29)
$`\psi _4^0`$ $`=`$ $`\overline{𝔡}^2{\displaystyle \frac{V^{}}{V}}`$ (30)
$`\psi _2^0\overline{\psi }_2^0`$ $`=`$ $`0`$ (31)
In the first two cases, the dot indicates $`u_B`$-derivatives, while in the NU case the prime means the $`\tau `$ derivative.
We want to emphasize that we have not lost any generality when going to the Bondi/Twisting tetrad but the NU equations are valid only when the original Bondi shear $`\sigma (u_B,\zeta ,\overline{\zeta })`$ was pure ‘electric’. In other words all the information in an ‘electric’ shear can be put into the real $`V(u_B,\zeta ,\stackrel{~}{\zeta }`$). Details will be given in Sec.VI.
## 5 The Physical Quantities
Most often one extracts or attempts to extract the physical quantities by trying to identify them in the midst of this set of functions at $`^+`$ by integrals over the $`\psi ^{}`$s and complicated combinations of the $`\sigma .`$ In general this a difficult and frequently ambiguous task \- largely because of the Bondi super-translation freedom. Only in the case of the identification of four-momentum has it been done successfully and unambiguously. In Bondi coordinates and tetrad, one defines the mass aspect $`\mathrm{\Psi }(u_B,\zeta ,\overline{\zeta }),`$ which is real from Eqs.(16) and (17), by
$$\mathrm{\Psi }=\psi _2^0+\text{ }𝔡^2\overline{\sigma }+\sigma \overline{\sigma }^{}.$$
(32)
The four-momentum is extracted from the $`l=0`$ & $`1`$ coefficients of the spherical harmonics. When Eq.(12) is expressed in terms of $`\mathrm{\Psi }`$ it takes the very simple form
$$\mathrm{\Psi }^{}=\sigma ^{}\overline{\sigma }^{}$$
(33)
which, when integrated over the sphere, becomes the Bondi mass loss equation. These results can be transformed and reexpressed via the twisting tetrads or in the NU coordinate/tetrad system. Explicitly, Eq.(20) becomes identical to Eq.(33) except that now
$$\sigma ^{}\text{ }𝔡L^{}+\frac{1}{2}(L^2)^{}.$$
Unfortunately there is no agreed upon asymptotic definition of orbital or spin angular momentum.
In the following section we will address this issue from a totally different point of view. Up to now the physical quantities were assumed to be contained in the asymptotic Weyl tensor and shear and constructed by a kinematic process.
We propose that rather than the kinematic process, there is a dynamic process involving the known functions on $`^+`$ and only after it is solved can the orbital or spin angular momentum be found.
## 6 A New Idea - Dynamics on $`^+`$
We first will show that there are special choices of $`L(u_B,\zeta ,\overline{\zeta })`$ leading to the twisting tetrads for which the shear, $`\sigma ^{}=0.`$ To find this family of functions $`L(u_B,\zeta ,\overline{\zeta }),`$ a differential equation must be solved. The solution, however, is not unique. The freedom in the solution is four arbitrary complex functions of a single complex variable, i.e., it is an arbitrary complex curve in a four-dimensional parameter space. \[The parameter space is the well-studied $`H`$-space.\] Note that all the information in the Bondi characteristic data, $`\sigma (u_B,\zeta ,\overline{\zeta }),`$ will have been shifted to the $`L(u_B,\zeta ,\overline{\zeta }).`$ Given an $`L(u_B,\zeta ,\overline{\zeta })`$ with any one of these curves, we could go backwards to recover the original Bondi $`\sigma (u_B,\zeta ,\overline{\zeta }).`$ The dynamics lies in the unique determination of this curve from other considerations which are describe later. The real part of this curve determines the center of mass (or orbital momentum) while the imaginary part yields the spin-angular momentum. These quantities are invariant under supertranslations.
We begin by observing that if we start with a Bondi coordinate/tetrad system with a given shear $`\sigma (u_B,\zeta ,\overline{\zeta }),`$ then \[after an unpleasant calculation\] the shear of the new null vector $`l^a`$ after the null rotation, Eqs.(4), is related to the old one by
$$\sigma ^{}(u_B^{},\zeta ,\overline{\zeta })=\sigma (u_B,\zeta ,\overline{\zeta })\text{ }𝔡LLL^{}.$$
The function $`L(u,\zeta ,\overline{\zeta })`$ is then chosen so that the new shear vanishes, i.e., it must satisfy
$$\text{ }𝔡L+LL^{}=\sigma (u_B,\zeta ,\overline{\zeta }).$$
(34)
###### Remark 3
The special case, $`𝔡L+LL^{}=0,`$ played an important early role in the development of twistor theory, leading immediately to the relationship in flat space between shear-free null geodesic congruences and twistor theory. It leads us to conjecture that Eq.(34) might play a role in some asymptotic form of twistor theory.
Though Eq.(34) is non-linear with an arbitrary right-side and appears quite formidable, considerable understanding of it can be found by the following procedure:
Assume the existence of a complex function
$$\tau =T(u_B,\zeta ,\overline{\zeta })$$
(35)
that is invertible in the sense that it can, in principle, be written as
$$u_B=X(\tau ,\zeta ,\overline{\zeta }).$$
(36)
Then writing
$$L=\frac{\text{ }𝔡T}{T^{}},$$
(37)
and using the implicit derivatives of Eq.(36),
$`1`$ $`=`$ $`X,_\tau T^{}`$ (38)
$`0`$ $`=`$ $`𝔡_{(\tau )}X+X^{}\text{ }𝔡T,`$
we obtain
$$L=\text{ }𝔡_{(\tau )}X.$$
(39)
Prime means the $`\tau `$ derivative and $`𝔡_{(\tau )}`$ means $`𝔡`$ with $`\tau `$ held constant. Finally, with this implicit view of $`X`$ and $`T`$, Eq.(34) becomes
$$\text{ }𝔡_{(\tau )}^2X=\sigma (X,\zeta ,\overline{\zeta }).$$
(40a)
###### Remark 4
Note that this equation, derived here from a different perspective, is identical to the so-called good cut equation. In general it has a four complex dimensional solution space that has been dubbed $`H`$-space. $`H`$-space possesses a complex metric which is Ricci-flat with a self-dual Weyl tensor. It formed the basis of Penrose’s asymptotic twistor space. In its original form it had its origin in the search for complex null surfaces that were asymptotically shear-free. Here we are looking for real null geodesic congruences that are, in general, not surface forming, i.e., have twist, but are asymptotically shear-free.
###### Remark 5
Note that using the reparametrization
$$\tau \tau ^{}=F(\tau )=T^{}(u_B,\zeta ,\overline{\zeta })$$
(41)
leaves Eq.(37) and the entire construction invariant. This fact will be used later for certain simplifications.
###### Remark 6
We mention that if the shear in Eq.(40a), is pure electric \[i.e., $`\sigma (u_B,\zeta ,\overline{\zeta })`$ $`=`$ $`𝔡^2S(u_B,\zeta ,\overline{\zeta })`$ with $`S(u_B,\zeta ,\overline{\zeta })`$ a real function\] then the associated $`H`$-space is flat and has a real four-dimensional subspace that can be identified with Minkowski space. We return to this issue later.
As we just mentioned, the solutions to Eq.(40a) depend on four complex parameters, say $`z^a,`$ and can be summarized by
$$u_B=X(z^a,\zeta ,\overline{\zeta }).$$
(42)
However, since we are interested in solutions of the form, Eq.(36), i.e.,
$$u_B=X(\tau ,\zeta ,\overline{\zeta }),$$
(43)
all we must do is chose in, $`H`$-space, an arbitrary complex world-line, $`z^a=\xi ^a(\tau ),`$ and substitute it into Eq.(42) to obtain the form (43). Furthermore, it can be seen that the solution, (43), can be written as
$$u_B=X(\tau ,\zeta ,\overline{\zeta })=\xi ^a(\tau )\widehat{l}_a(\zeta ,\overline{\zeta })+X_{(l2)}(\tau ,\zeta ,\overline{\zeta })$$
with $`X_{(l2)}`$ containing spherical harmonics, $`l2`$ and
$$\widehat{l}_a(\zeta ,\overline{\zeta })=\frac{\sqrt{2}}{2}(1,\frac{\zeta +\overline{\zeta }}{1+\zeta \overline{\zeta }},i\frac{\zeta \overline{\zeta }}{1+\zeta \overline{\zeta }},\frac{1\zeta \overline{\zeta }}{1+\zeta \overline{\zeta }})$$
(44)
which has only the $`l=0`$ & $`1,`$ spherical harmonics. This result follows from the fact that $`𝔡_{(\tau )}^2`$ in the good-cut equation annihilates the $`l=0`$ & $`1`$ terms. The $`\xi ^a(\tau )`$ does get fed back into the higher harmonics via the $`\sigma (X,\zeta ,\overline{\zeta }).`$
Our solution to Eq.(34) can now be expressed implicitly by
$`L(u_B,\zeta ,\overline{\zeta })`$ $`=`$ $`𝔡_{(\tau )}X=\xi ^a(\tau )\widehat{m}_a(\zeta ,\overline{\zeta })+\text{ }𝔡_{(\tau )}X_{(l2)}(\tau ,\zeta ,\overline{\zeta })`$ (45)
$`u_B`$ $`=`$ $`X=\xi ^a(\tau )\widehat{l}_a(\zeta ,\overline{\zeta })+X_{(l2)}(\tau ,\zeta ,\overline{\zeta }).`$ (46)
###### Remark 7
We emphasize that the complex world-line, $`\xi ^a(\tau ),`$ is not in physical space but is in the parameter space, $`H`$-space. At this point we are not suggesting that there is anything profound about this observation. It is simply there and whatever meaning it might have is obscure. We however remark that this observation can be extended to the Einstein-Maxwell equations where in general the Maxwell field will have its own complex world-line. If the two independent complex world lines coincide, \[from preliminary results and special cases\] one obtains the Dirac value of the gyromagnetic ratio.
###### Remark 8
To invert the later equation, (46), to obtain $`\tau =T(u_B,\zeta ,\overline{\zeta }),`$ is, in general, virtually impossible. There however is a relatively easy method, by iteration, to get approximate inversions to any accuracy. {See Appendix A}
###### Remark 9
We have tacitly assumed that the world line is complex analytic in the complex parameter $`\tau .`$
If we expand $`\xi ^a(\tau )`$ in a Taylor series and regroup the terms with real coefficients and those with imaginary coefficients separately, we can write
$$\xi ^a(\tau )=\xi _R^a(\tau )+i\xi _I^a(\tau ).$$
This decomposition becomes important later.
###### Remark 10
The asymptotic twist, $`\mathrm{\Sigma }(u_B,\zeta ,\overline{\zeta }),`$ {a measure of how far the vector field $`l^a`$ is from being surface forming} is defined by
$$2i\mathrm{\Sigma }=\text{ }𝔡\overline{L}+L\overline{L}^{}\text{ }\overline{𝔡}L\overline{L}L^{}.$$
If $`L`$ is obtained from a shear that is pure ‘electric’, then the $`\xi ^a(\tau )`$ can be chosen so that the asymptotic twist vanishes.
It is from the fact that we can chose $`L(u,\zeta ,\overline{\zeta })`$as a solution of Eq.(34), leading to
$`\sigma ^{}=\sigma ^{}=0,`$ that is the justification for the form of Eqs.(18)-(25).
### 6.1 Determining the Complex Curve
The argument and method for the unique determination of the complex curve (and eventually the definitions of spin and center of mass motion) are not along conventional lines of thought. To try to clarify it, it is worthwhile to take a brief detour.
It is well-known that for a static charge distribution, the electric dipole moment is given by the total charge times the center of charge position given in the ‘static’ Lorentz frame. In the case of a dynamic charge it is more difficult since the center of charge depends on the choice of Lorentz frame. Nevertheless by having, in any one frame, the center of charge at the coordinate origin, the electric dipole moment will vanish. In the case of a single charged particle moving on an arbitrary real world line, (the Lienard-Wiechert Maxwell field) the electric dipole moment vanishes when the coordinate origin follows the particle’s motion or equivalently, the dipole is given by the particles displacement from the coordinate origin times the charge. The asymptotically defined dipole moment also vanishes when it is calculated (extracted) from the $`l=1`$ part of the coefficient of the $`r^3`$ term of $`\varphi _0=F_{ab}l^am^b,`$ where $`l^a`$ is a tangent vector of the light-cones that are attached to the particles world-line. Analogously, it was shown that the magnetic dipole moment of a charged particle can be ‘viewed’ as arising from a charge moving in complex Minkowski space along a complex world-line. It is given by the charge times the imaginary displacement. {We emphasize again that the complex world-lines are a bookkeeping device and no claim is made that particles are really moving in complex space. The real effect of the complex world-line picture is to create a twisting real null congruence.} If the light-cone from the complex world-line is followed to $`^+,`$ (equivalent to introducing an appropriate asymptotic twist for the null vector $`l^a`$ that is used in $`\varphi _0=F_{ab}l^am^b`$) then both the asymptotic electric and magnetic dipoles vanish. This argument is given in considerable detail in the paper.
The idea is to generalize this construction to GR. In linearized GR and in the exact case of the Kerr and charged Kerr metrics, the mass-dipole moment and the spin sit in the real and imaginary parts of the $`l=1`$ harmonic of the coefficient of the $`r^4`$ part of the Weyl component $`\psi _1^0`$ in a Bondi coordinate/tetrad system. This observation is extended to all asymptotically flat vacuum solutions by going to an asymptotically twisting tetrad, with an appropriately chosen $`H`$-space complex curve \[see below\], such that the new Weyl component, $`\psi _1^0,`$ has a vanishing $`l=1`$ harmonic in the coefficient of its $`r^4`$ part. We refer to this curve as the intrinsic complex center of charge world-line and, rough speaking, the center of mass and the spin angular momentum will be identified by its real and imaginary parts. Effectively, by chosing the world-line so that the $`l=1`$ harmonic vanishes we have shifted the ‘origin’ to the complex world-line.
To explicitly apply this idea to our discussion of twisting null congruences, we go to the set of equations, Eq.(19) and (20), namely
$`\psi _1^0`$ $`=`$ $`\text{ }𝔡\psi _2^0L\psi _2^03\dot{L}\psi _2^0`$ (47)
$`\psi _2^0`$ $`=`$ $`\text{ }𝔡\psi _3^0L\psi _3^02\dot{L}\psi _3^0`$ (48)
and require that, in (47), that the $`l=1`$ part of $`\psi _1^0`$ should vanish. This leads (via approximations) to a differential equation for the complex center-of-mass world-line. It is driven by the original Bondi shear. To see this in detail, first we recall that (48) can be rewritten as
$`\mathrm{\Psi }^{}`$ $`=`$ $`\sigma ^{}\overline{\sigma }^{}`$ (49)
$`\sigma ^{}`$ $``$ $`𝔡L^{}+{\displaystyle \frac{1}{2}}(L^2)^{}`$ (50)
$`\psi _2^0`$ $`=`$ $`\mathrm{\Psi }2L\text{ }𝔡\overline{\sigma }^{}L^2\overline{\sigma }^{}\text{ }𝔡^2\overline{\sigma }\sigma \overline{\sigma }^{}.`$ (51)
If we use the result \[see Appendix B\] that the Bondi four-momentum, $`P_a,`$ can be extracted from $`\mathrm{\Psi }`$ by
$$P_a(u_B)(M,P^i)=\frac{c^2}{8\pi G}\mathrm{\Psi }\widehat{l}_a𝑑S$$
(52)
we can invert and see that
$$\mathrm{\Psi }(u_B,\zeta ,\overline{\zeta })\chi \chi ^i\widehat{c}_i+\mathrm{}.=2\frac{G}{c^2}(\sqrt{2}M3P^i\widehat{c}_i)+\mathrm{\Psi }_{l2}\mathrm{}$$
(53)
with
$`\widehat{l}_b(\zeta ,\overline{\zeta })`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{2}}(1,C_i(\zeta ,\overline{\zeta })),`$ (54)
$`\widehat{c}_b`$ $``$ $`\widehat{l}_b\widehat{n}_b=\sqrt{2}(0,C_i(\zeta ,\overline{\zeta }))`$ (55)
$`C_a`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{2}}\widehat{c}_a,`$ (56)
$`C_i`$ a unit radial vector and $`\mathrm{\Psi }_{l2}`$ containing only harmonics $`l2.`$
The time evolution of $`P_a(u)`$ is found by integration Eq.(49), i.e., from $`\mathrm{\Psi }=\sigma ^{}\overline{\sigma }^{}𝑑u`$ or
$$P_a^{}(u_B)=\sigma ^{}\overline{\sigma }^{}\widehat{l}_a𝑑S.$$
(57)
If $`\psi _2^0,`$ from Eq.(51), is substituted into (47 ) we obtain after using Eq.(49) and regrouping, a truly ugly equation that we write as
$$\psi _1^0\text{ }𝔡^3\overline{\sigma }+\text{ }𝔡\mathrm{\Psi }+3L^{}\mathrm{\Psi }=(NL)$$
(58)
with the non-linear terms, (NL), expressed by
$`(NL)`$ $``$ $`𝔡[\overline{\sigma }^{}\sigma ]+2\sigma \text{ }𝔡\overline{\sigma }^{}+3L\text{ }𝔡^2\overline{\sigma }+L^3\overline{\sigma }^{\mathrm{}}+3L\sigma \overline{\sigma }^{}`$
$`+3L^2\text{ }𝔡\overline{\sigma }^{}+3\dot{L}[2L\text{ }𝔡\overline{\sigma }^{}+L^2\overline{\sigma }^{}+\text{ }𝔡^2\overline{\sigma }+\sigma \overline{\sigma }^{}].`$
The grouping has the non-linear higher order terms put into (NL) while the ‘controllable’ (lower order) terms are on the left side. It is from this equation and (57) we find the equations of motion. We first require, as we said earlier, that $`\psi _1^0`$ (and of course $`\psi _1^0`$) have no $`l=1`$ harmonics. From the fact that $`𝔡^3\overline{\sigma },`$ from its spin $`s=2`$ structure, has no $`l=1`$ terms we see that Eq.(58) reduces to the condition that
$$(\text{ }𝔡\mathrm{\Psi }+3L^{}\mathrm{\Psi })_{l=1}=(NL)_{l=1}.$$
(60)
It is this equation that we must analyze, or more accurately approximate, by looking only at the leading spherical harmonic terms. \[See the following subsection for an alternative derivation of Eq.(60).\]
From Eqs.(53 ) and (45), considering only the lowest order $`l=1`$ terms we have
$`\mathrm{\Psi }(u_B,\zeta ,\overline{\zeta })`$ $`=`$ $`2{\displaystyle \frac{G}{c^2}}(\sqrt{2}M3P^ic_i)`$
$`L(u_B,\zeta ,\overline{\zeta })`$ $`=`$ $`\xi ^a(\tau )\widehat{m}_a(\zeta ,\overline{\zeta })`$
from which it follows that
$`𝔡\mathrm{\Psi }`$ $`=`$ $`12{\displaystyle \frac{G}{c^2}}P^a\widehat{m}_a`$ (61)
$`3L^{}\mathrm{\Psi }`$ $`=`$ $`6{\displaystyle \frac{G}{c^2}}\xi ^a(\tau )\widehat{m}_a\{\sqrt{2}M3P^ic_i\}`$ (62)
$`=`$ $`\sqrt{2}6{\displaystyle \frac{G}{c^2}}M\xi ^a(\tau )\widehat{m}_a+18{\displaystyle \frac{G}{c^2}}\xi ^aP^b(\tau )\widehat{m}_a\widehat{c}_b`$ (63)
$`=`$ $`\sqrt{2}6{\displaystyle \frac{G}{c^2}}M\xi ^a(\tau )\widehat{m}_a+9i{\displaystyle \frac{G}{c^2}}\xi ^aP^b(\tau )ϵ_{abcd}\widehat{m}^ct^d+Y\{l2\}`$ (64)
where we have used the Clebsch-Gordon expansion \[see Appendix C\] of the product $`\widehat{m}_ac_b`$. From Eq.(60), we then have
$$12\frac{G}{c^2}P^a\widehat{m}_a\sqrt{2}6\frac{G}{c^2}M\xi ^a(\tau )\widehat{m}_a+9i\frac{G}{c^2}\xi ^aP^b(\tau )ϵ_{abcd}\eta ^{cf}\widehat{m}_ft^d=(NL)_{l=1}GN^a\widehat{m}_a$$
or, using the rescaled $`u_B`$ to the conventional Bondi $`u_c`$ coordinate1, by $`u_B=\frac{u_c}{\sqrt{2}}`$.
$$P^e=M\xi ^e(\tau )i\frac{3\sqrt{2}}{4}\xi ^aP^b(\tau )ϵ_{abcd}\eta ^{ce}t^d+D^e.$$
(65)
We have the three-momentum expressed in terms of the mass $`M`$ and the complex velocity, $`\xi ^a.`$ Remembering, from Eq.(25) that $`\mathrm{\Psi }`$ is real we see, from Eq.(53), that $`M`$ and $`P^c`$ are real. By defining the real and imaginary parts of $`\xi ^a,`$
$$\xi ^a=\xi _R^a+i\xi _I^a,$$
(66)
and taking the real and imaginary parts of (65), (ignoring $`D^f`$ in this approximation)
$`P^e`$ $`=`$ $`M\xi _R^e(\tau )+{\displaystyle \frac{3\sqrt{2}}{4}}\xi _I^aP^b(\tau )ϵ_{abcd}\eta ^{ce}t^d`$ (67)
$`M\xi _I^e(\tau )`$ $`=`$ $`{\displaystyle \frac{3}{4}}\xi _R^aP^b(\tau )ϵ_{abcd}\eta ^{ce}t^d.`$ (68)
If we identify the spin-angular momentum (from the analogy with the Kerr metric) by
$$S^b=M\xi _I^b$$
and $`\xi _R^b`$ $`v^b`$ with a real velocity vector and linearize, we have
$`P^c`$ $`=`$ $`Mv^c`$
$`S^b`$ $`=`$ $`{\displaystyle \frac{M^{}}{M}}S^b`$
which, with
$$P_a^{}(u_B)=(Mv^c)^{}=\sigma ^{}\overline{\sigma }^{}\widehat{l}_a𝑑S,$$
yields the equations of motion for the center of mass and spin-angular momentum - the equations of motion of a spinning particle with mass and momentum loss via radiation $``$ all derived from the Bianchi identities and driven by the original Bondi shear.
###### Remark 11
Note that from the linearization of Eqs.(67) and (68), the coupling between the spin and the velocity vector has disappeared. The details of the coupling will be investigated in a future note.
There is an important point to be made here. We know that the complex curve we wish to determine is given by the expression
$$z^a=\xi ^b(\tau ),$$
with the complex paramter $`\tau ,`$ while in the above expressions for the determination of the curve we have treated $`\xi ^b`$ to be a complex funtion of the real parameter $`u`$. This is justified by realizing that we have been making severe approximations, throwing out higher non-linear terms, not using more and higher Clebsch-Gordon expansions and from the linearization of the inversion of $`\tau =T(u,\zeta ,\overline{\zeta }),`$ Appendix A. In the lowest order of the inversion we do have $`\tau =u.`$
The basic idea we are espousing is that there is, in $`H`$-Space, a unique complex curve that can be called the intrinsic complex center of mass world-line. It is determined from a geometric structure on the physical space-time, the null direction field, independent of the choice of coordinates on Scri. In order to obtain - as an approximation - what one would normally call the orbital and spin angular momentum in linear theory one must use Eq.(6) and extract by an integration, from the Bondi $`\psi _1^0,`$ its $`l=1`$ coefficient. The values of this ‘normally’ defined angular momentum will, in general, depended on the choice of the scri coordinates or the Bondi frame.
###### Remark 12
If we linearize (decouple) Eq.(65) we have
$$P^e=M\xi ^e(u)$$
so that if we treat $`M`$ and $`P^e`$ as very slowly varying we can integrate it as
$$\xi ^e=\xi _R^a+i\xi _I^a=\frac{P^e}{M}u$$
or, since $`M`$ and $`P^e`$ are real,
$`\xi _R^a={\displaystyle \frac{P^e}{M}}u`$
$`\xi _I^a=const=\text{initial conditions,}`$
thus obtaining the usual motion for the center of mass.
### 6.2 An Alternative Route to the Complex Curve
There is an alternative method of finding the curve that in some sense is more direct than the one just given.
We begin by inverting Eqs.(5) - (9) and then writing
$$\psi _1^0=\psi _1^03L\psi _2^0+3L^2\psi _3^0L^3\psi _4^0$$
and its $`u_B`$-derivative
$$\psi _1^0=(\psi _1^03L\psi _2^0+3L^2\psi _3^0L^3\psi _4^0)^{}$$
(69)
We then impose our basic requirement for the determination of the curve, namely that the lowest harmonic, i.e., the $`l=1`$ term in $`\psi _1^0`$and $`\psi _1^0`$ should vanish. Eq.(69) can then be written as
$$\{\psi _1^0\}_{l=1}=\{(3L\psi _2^03L^2\psi _3^0+L^3\psi _4^0)^{}\}_{l=1}+[\text{Terms]}_{l2}.$$
(70)
Eq.(70), using the Bianchi identities for $`\psi _1^0`$ and $`\psi _1^0`$ and expressions $`\psi _3^0`$and $`\psi _4^0,`$ i.e., Eqs.(11) - (14), again leads to the complex curve by the following argument.
Using
$$\mathrm{\Psi }=\psi _2^0+\text{ }𝔡^2\overline{\sigma }+\sigma \overline{\sigma }^{}$$
(71)
instead of $`\psi _2^0`$ and the evolution equation for $`\mathrm{\Psi },`$ Eq.(33), (really a Bianchi identity)
$$\mathrm{\Psi }^{}=\sigma ^{}\overline{\sigma }^{}$$
(72)
Eq.(70) can be rewritten as
$$\{3L^{}\text{ }\mathrm{\Psi }\psi _1^0\}_{l=1}=NL_{l=1}+[\text{ Terms]}_{l2},$$
(73)
while the Bianchi identity for $`\psi _1^0^{}`$ becomes
$$\psi _1^0=\text{ }𝔡\mathrm{\Psi }+NL+[\text{Terms]}_{l2}.$$
(74)
In both cases, for our approximation, we have simply grouped the higher harmonic terms and the nasty non-linear terms into $`[`$Terms\]<sub>l⩾2</sub> and $`NL`$.
By substituting the $`\psi _1^0`$ from Eq.(74) into (73) we obtain
$$(\text{ }𝔡\mathrm{\Psi }+3L^{}\mathrm{\Psi })_{l=1}=(NL)_{l=1},$$
(75)
which is identical to Eq.(60). This relationship, with Eqs.(45) and (46),
$`L(u_B,\zeta ,\overline{\zeta })`$ $`=`$ $`\xi ^a(\tau )\widehat{m}_a(\zeta ,\overline{\zeta })+\text{ }𝔡_{(\tau )}X_{(l2)}(\tau ,\zeta ,\overline{\zeta })`$ (76)
$`u_B`$ $`=`$ $`\xi ^a(\tau )\widehat{l}_a(\zeta ,\overline{\zeta })+X_{(l2)}(\tau ,\zeta ,\overline{\zeta })`$ (77)
leads, via the argument in the previous section, to the equations for the complex curve.
## 7 Examples
### 7.1 The Robinson-Trautman metrics
We begin with the Bianchi Identities Eqs.(26) - (31), which are essentially the same as Eqs.(18)-(25), differing only by the coordinate transformation $`\tau u_B,`$ given by $`u_B=X(\tau ,\zeta ,\overline{\zeta }),`$ from one set to the other. The tetrad system, as well as the stereographic angle $`L,`$ are the same with both sets. In other words, geometrically, they represent the same situation. The $`L`$ or the $`l^a`$ describe a surface forming congruence. We imposed, in the last section, the condition that the $`l=1`$ harmonics of the tetrad component, $`\psi _1^0,`$ vanish. In this example we simply strengthen this $`condition`$ by requiring that
$$\psi _1^0=0.$$
This condition, from Eq.(27) forces $`\psi _2^0`$ to be a function only of $`\tau ,`$ i.e., $`\psi _2^0=\psi _2^0(\tau )=2\sqrt{2}GM(\tau )/c^2,`$ so that Eq.(28) becomes
$$\psi _2^0=\overline{𝔡}𝔡\overline{𝔡}𝔡\mathrm{log}V+3\frac{V^{}}{V}\psi _2^0.$$
(78)
(The strange numerical factor $`2\sqrt{2}`$ arises from our use of the radial coordinate $`r`$ that differs from the conventional one by the factor$`\sqrt{2}.`$)
By the reparametrization freedom, Eq.(41), a change in the $`\tau `$ coordinate can be made, i.e., $`\tau =F(\tau ^{}),`$ so that the $`\psi _2^0`$ is constant, thus yielding
$$\frac{V^{}}{V}\psi _2^0=\overline{𝔡}𝔡\overline{𝔡}𝔡\mathrm{log}V$$
(79)
the well-known Robinson-Trautman equation describing an asymptotic Robinson-Trautman space-time. If we further assumed that $`\psi _0=0`$ this would lead to the Robinson-Trautman metric.
Note that since
$`u`$ $`=`$ $`X=\xi ^a(\tau )\widehat{l}_a(\zeta ,\overline{\zeta })+X_{(l2)}(\tau ,\zeta ,\overline{\zeta })`$
$`V`$ $`=`$ $`X^{}=\xi ^a(\tau )\widehat{l}_a(\zeta ,\overline{\zeta })+X_{(l2)}^{}(\tau ,\zeta ,\overline{\zeta })`$
we see that Eq.(79), is a differential equation for $`V`$ and more specifically the coefficients of the $`l=0`$ $`\&`$ $`1`$ harmonics have the form of Newton’s second law of motion
$$M\xi ^{a\prime \prime }=F^a$$
so that the Robinson-Trautman space-time describes a gravitational rocket ship.
If we do have a solution to the Robinson-Trautman equation, i.e., $`V=V(\tau ,\zeta ,\overline{\zeta }),`$ the associated Bondi shear can be easily constructed from the implicit relations
$`u_B`$ $`=`$ $`X(\tau ,\zeta ,\overline{\zeta })={\displaystyle V(\tau ,\zeta ,\overline{\zeta })𝑑\tau }`$ (80)
$`\sigma (u_B,\zeta ,\overline{\zeta })`$ $`=`$ $`𝔡_{(\tau )}^2X(\tau ,\zeta ,\overline{\zeta }).`$
### 7.2 Twisting Type II Metrics
We can describe a more general situation by going to Eqs.(18)-(25) and again require that the Weyl components
$`\psi _1^0`$ $`=`$ $`0`$
$`\psi _0^0`$ $`=`$ $`0.`$
This leads to the equations
$`𝔡\psi _2^0+L\psi _2^0+3L^{}\psi _2^0`$ $`=`$ $`0`$ (81)
$`𝔡\psi _3^0+L\psi _3^0+2L^{}\psi _3^0`$ $`=`$ $`\psi _2^0,`$ (82)
with
$`\psi _3^0`$ $`=`$ $`𝔡\overline{\sigma }^{}+L\overline{\sigma }^{},`$ (83)
$`\sigma `$ $``$ $`𝔡L+{\displaystyle \frac{1}{2}}(L^2)^{},`$
$`\mathrm{\Psi }`$ $`=`$ $`\psi _2^0+2L\text{ }𝔡\overline{\sigma }^{}+L^2\overline{\sigma }^{}+\text{ }𝔡^2\overline{\sigma }+\sigma \overline{\sigma }^{},`$
$`\mathrm{\Psi }\overline{\mathrm{\Psi }}`$ $`=`$ $`0,`$
that are precisely the equations for the twisting algebraically special type II metrics. A detailed analysis of these equations and their generalization to the type II Einstein-Maxwell equations is in preparation.
### 7.3 The Sparling-Tod Shear
George Sparling and Paul Tod found a class of Bondi shears such that the good-cut equation (40a)
$$\text{ }𝔡^2X=\sigma (X,\zeta ,\overline{\zeta })$$
(84)
can be integrated exactly. An example of their class is
$$\sigma (u,\zeta ,\overline{\zeta })=\frac{(q_am^a)^2}{u^3}$$
(85)
with $`q_a`$ the constant vector
$$q_a=(0,q_i).$$
(86)
The solution to (84) can be given by
$$u^2=X^2(z^a,\zeta ,\overline{\zeta })=Z^2+S^2$$
(87)
with
$`Z`$ $`=`$ $`z^al_a(\zeta ,\overline{\zeta }),`$ (88)
$`S`$ $`=`$ $`s^al_a(\zeta ,\overline{\zeta }).`$
The four $`z^a=(t,x,y,z)`$ are the $`H`$-Space coordinates (four constants of integration) and the $`s^a(z^b)`$ are four functions of $`z^a`$ determined by the algebraic equation
$$s^az^b(l_am_bl_bm_a)=q_am^a.$$
(89)
Using the identity (see Appendix C for the notation and conventions)
$$(\widehat{l}_a\widehat{m}_b\widehat{m}_a\widehat{l}_b)=\frac{1}{2}(\widehat{t}_a\widehat{m}_b\widehat{m}_a\widehat{t}_b)\frac{1}{2}iϵ_{abcd}\widehat{m}^ct^d$$
(90)
and the ansatz
$`s^a`$ $`=`$ $`S_0\text{ }t^a+S_1\text{ }q^a+i\text{ }S_2\sqrt{2}\text{ }\epsilon ^a{}_{bc}{}^{}q_{}^{b}z^c`$ (91)
$`\sqrt{2}\text{ }\epsilon ^a_{bc}`$ $`=`$ $`\text{ }\epsilon ^a{}_{bcd}{}^{}t_{}^{d}`$ (92)
one find that the ($`S_0,S_1,S_2`$) are uniquely determined so that
$$s^a=\frac{1}{z_ez^e}\{z_bq^bt^az^bt_bq^a+i\text{ }\epsilon ^a{}_{bcd}{}^{}t_{}^{d}q^bz^c\}.$$
(93)
Note that though $`s^a,`$ in Eq.(89) was defined up to a term proportional to $`z^a`$, the ansatz determined it uniquely.
Eqs.(93), (87) and (LABEL:SZ) determine the solution $`u_B=X(z^b,\zeta ,\overline{\zeta })`$ so when an $`H`$-Space world-line, $`z^b=\xi ^b(\tau ),`$ is chosen we have the null direction field
$`L(u_B,\zeta ,\overline{\zeta })`$ $`=`$ $`𝔡_{(\tau )}X(\xi ^a(\tau ),\zeta ,\overline{\zeta })={\displaystyle \frac{Z\text{ }𝔡Z+S\text{ }𝔡S}{\sqrt{Z^2+S^2}}}`$ (94)
$`u_B`$ $`=`$ $`\sqrt{Z^2+S^2}.`$
### 7.4 Arbitrary Construction
For further general examples, we can reverse the process of first giving the shear and then finding the $`L(u_B,\zeta ,\overline{\zeta })`$. Instead, we can choose an arbitrary, $`s=0,`$ function with the form
$$u_B=X(\tau ,\zeta ,\overline{\zeta })=\xi ^a(\tau )\widehat{l}_a(\zeta ,\overline{\zeta })+X_{(l2)}(\tau ,\zeta ,\overline{\zeta })$$
and then define
$$L(u_B,\zeta ,\overline{\zeta })=\text{ }𝔡X=\xi ^a(\tau )\widehat{m}_a(\zeta ,\overline{\zeta })+\text{ }𝔡X_{(l2)}(\tau ,\zeta ,\overline{\zeta }).$$
The associated Bondi shear is then given, implicitly, by
$`\sigma (u_B,\zeta ,\overline{\zeta })`$ $`=`$ $`𝔡^2X_{(l2)}(\tau ,\zeta ,\overline{\zeta })`$
$`u`$ $`=`$ $`\xi ^a(\tau )\widehat{l}_a(\zeta ,\overline{\zeta })+X_{(l2)}(\tau ,\zeta ,\overline{\zeta }).`$
## 8 Conclusion
In this work we have shown that for all asymptotically flat space-times there is a hidden structure that must be extracted dynamically from the known asymptotic Weyl tensor and its related characteristic data, the Bondi shear, $`\sigma `$. This hidden structure is a specific field of asymptotically shear-free null directions - or an asymptotically shear-free null geodesic congruence. Within the information for the description of this direction field is a complex world-line that is defined in the $`H`$-space associated with the given asymptotically flat space-time. One can try to give meaning to this world-line by defining it to be the complex center of mass of the interior gravitating system. Though there is no proof or even overwhelming evidence that this assignment is completely reasonable, we can ask the question: what physical justification can be given for this assignment. On the surface it certainly is strange - where does a complex world-line in $`H`$-space enter in any direct physical observation.
We would like to try to give a justification by three different types of argument.
1. Examples: In several known cases, e.g., the Kerr-metric, the charged-Kerr metric, the Lienard-Wiechert-Maxwell fields, the Kerr-Maxwell field, all have a center of mass (or center of charge in the Maxwell case) that is identical to the more general construction described here. In addition, the construction given here coincides with the definitions from linearized theory.
2. Though it appears to us to be surprising and remarkable, one does get, in a first approximation, the conventional equations of motion of a complicated spinning gravitating source directly from the asymptotic information - with no details needed about the interior source. These equations contain the radiation reaction from the gravitational radiation that arises from the Bondi mass loss.
3. It is possible to ask the question: can one give, at least in principle, an observational means of “observing” this complex motion? The answer, we believe, is yes - though to really do so is impossible. It involves having a huge number of observers surrounding the gravitating source and looking at all the null rays reaching them, then picking out shear-free null direction fields and by angular integration finding the complex world-line. This type of argument must be tightened and made more precise.
Work has begun on the application of these ideas to asymptotically flat Einstein-Maxwell fields. We remark that it appears (in a preliminary stage) that there are two complex world-lines, a complex center of mass obtained from the asymptotic Weyl tensor and a complex center of charge obtained form the asymptotic Maxwell field. Their imaginary parts give respectively the spin and magnetic dipole moments. If the two world-lines coincide then the Dirac value of gyromagnetic ratio follows.
Many years ago, by a rather strange and not at all understood complex coordinate transformation applied to the Reissner-Nordstrom Einstein-Maxwell field, the charged Kerr Einstein-Maxwell field was first found. The work described here appears to clarify and even explain this mysterious transformation.
Nevertheless, there is still much to be understood about the present material, i.e., the asymptotic shear-free congruences and the complex H–space curves. Is it just an coincidence that it appears to be so connected with physical issues or does something deeper lie behind it?
## 9 Acknowledgments
This material is based upon work (partially) supported by the National Science Foundation under Grant No. PHY-0244513. Any opinions, findings, and conclusions or recommendations expressed in this material are those of the authors and do not necessarily reflect the views of the National Science. E.T.N. thanks the NSF for this support. C.N.K. would like to thank CONICET for support. G.S.O. acknowledges the financial support from CONACyT through Grand No.44515-F, VIEP-BUAP through Grant No. II161-04/EXA and Sistema Nacional de Investigadores (SNI-M$`\backslash `$’exico).
## 10 Appendix A
In order to obtain $`\tau =T(u_B,\zeta ,\overline{\zeta }),`$ the inversion of
$$u_B=X(\tau ,\zeta ,\overline{\zeta })=z^a(\tau )\widehat{l}_a(\zeta ,\overline{\zeta })+X_{(l2)}(\tau ,\zeta ,\overline{\zeta }),$$
we first note that $`\tau `$ can be replaced by any function of $`\tau ,`$ i.e., the Eq.(37) is invariant under $`\widehat{\tau }=F(\tau ).`$ See Eq.(41). Thus by taking
$$z^0(\tau )=\tau $$
we have
$$u_B=\tau +z^i(\tau )\widehat{l}_i(\zeta ,\overline{\zeta })+X_{(l2)}(\tau ,\zeta ,\overline{\zeta })$$
or
$$\tau =u_Bz^i(\tau )\widehat{l}_i(\zeta ,\overline{\zeta })+X_{(l2)}(\tau ,\zeta ,\overline{\zeta }).$$
This can be solved by iteration to any order:
$`\tau _0`$ $`=`$ $`u_B`$
$`\tau _1`$ $`=`$ $`u_Bz^i(u_B)\widehat{l}_i(\zeta ,\overline{\zeta })+X_{(l2)}(u_B,\zeta ,\overline{\zeta })`$
$`\tau _n`$ $`=`$ $`u_Bz^i(\tau _{n1})\widehat{l}_i(\zeta ,\overline{\zeta })+X_{(l2)}(\tau _{n1},\zeta ,\overline{\zeta }).`$
## 11 Appendix B
Normalization and conventions: Our goal here is to express the function $`\mathrm{\Psi }`$ in terms of the Bondi mass/momenta and also show how to extract them from the $`\mathrm{\Psi }.`$
We start, in a given frame in Minkowski space, with $`C_a`$ a unit radial space-like vector and $`T_a`$ a unit time-like vector given by
$`C^a`$ $`=`$ $`(0,\mathrm{cos}\varphi \mathrm{sin}\theta ,\mathrm{sin}\varphi \mathrm{sin}\theta ,\mathrm{cos}\theta )=C_a,`$ (95)
$`T_a`$ $`=`$ $`(1,0,0,0).`$ (96)
with the null vectors
$`\widehat{l}_a`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(1,\mathrm{cos}\varphi \mathrm{sin}\theta ,\mathrm{sin}\varphi \mathrm{sin}\theta ,\mathrm{cos}\theta )={\displaystyle \frac{1}{\sqrt{2}}}(T_a+C_a)`$ (97)
$`\widehat{n}_a`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(1,\mathrm{cos}\varphi \mathrm{sin}\theta ,\mathrm{sin}\varphi \mathrm{sin}\theta ,\mathrm{cos}\theta )={\displaystyle \frac{1}{\sqrt{2}}}(T_aC_a)`$ (98)
and
$$\widehat{c}_a=\widehat{l}_a\widehat{n}_a=\sqrt{2}(0,\mathrm{cos}\varphi \mathrm{sin}\theta ,\mathrm{sin}\varphi \mathrm{sin}\theta ,\mathrm{cos}\theta )=\sqrt{2}C_a$$
We let $`a=(0,1,2,3)=(0,i)`$ and $`\widehat{l}_a=(l_0,l_i)`$ and $`\widehat{c}_i=\sqrt{2}C_i`$, etc.
We then define the $`l=0`$ and $`l=1`$ parts of the mass aspect $`\mathrm{\Psi }`$ by
$`\mathrm{\Psi }`$ $``$ $`\psi _2^0+2L\text{ }𝔡\overline{\sigma }^{}+L^2\overline{\sigma }^{}+\text{ }𝔡^2\overline{\sigma }+\sigma \overline{\sigma }^{}`$ (99)
$`\mathrm{\Psi }`$ $`=`$ $`\chi \chi ^i\widehat{c}_i+\mathrm{}.`$ (100)
Therefore
$`\mathrm{\Psi }\widehat{l}_a`$ $`=`$ $`\chi \widehat{l}_a\chi ^i\widehat{c}_i\widehat{l}_a+\mathrm{}`$
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}[\chi T_a\sqrt{2}\chi ^iC_iC_a+\chi C_a\sqrt{2}\chi ^iC_iT_a].`$
From the integral identities
$`{\displaystyle 𝑑S}`$ $`=`$ $`4\pi `$
$`{\displaystyle C_i𝑑S}`$ $`=`$ $`0`$
$`{\displaystyle C_iC_j𝑑S}`$ $`=`$ $`{\displaystyle \frac{4\pi }{3}}\delta _{ij}`$
we have that
$`{\displaystyle \mathrm{\Psi }\widehat{l}_a𝑑S}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle 𝑑S(\chi T_a\sqrt{2}\chi ^iC_iC_a)}`$ (102)
$`=`$ $`{\displaystyle \frac{4\pi }{\sqrt{2}}}[\chi T_a\chi ^i{\displaystyle \frac{\sqrt{2}}{3}}]={\displaystyle \frac{4\pi }{\sqrt{2}}}(\chi ,\chi _i{\displaystyle \frac{\sqrt{2}}{3}}).`$
with $`\chi _i=\chi ^i.`$
On the other hand, from reference () and a change of notation\[\*coordinatechange\], we have that the 4-momentum for an asymptotically flat space-time is defined by
$`P_a`$ $`=`$ $`{\displaystyle \frac{c^2}{8\pi G}}{\displaystyle \mathrm{\Psi }\widehat{l}_a𝑑S},`$ (103)
$`\mathrm{\Psi }`$ $`=`$ $`2\sqrt{2}{\displaystyle \frac{G}{c^2}}(M+\mathrm{}..`$ (104)
then from Eqs.(103),(102) and (11),
$`P_a`$ $`=`$ $`{\displaystyle \frac{c^2}{8\pi G}}[{\displaystyle \frac{4\pi }{\sqrt{2}}}(\chi ,\chi _i{\displaystyle \frac{\sqrt{2}}{3}})]`$ (105)
$`(P_0,P_i)`$ $`=`$ $`(M,P_i)={\displaystyle \frac{c^2}{2\sqrt{2}G}}(\chi ,\chi _i{\displaystyle \frac{\sqrt{2}}{3}}).`$
We find that:
$`\chi `$ $`=`$ $`2\sqrt{2}{\displaystyle \frac{G}{c^2}}M`$ (106)
$`\chi _i`$ $`=`$ $`6GP_i=6{\displaystyle \frac{G}{c^2}}P^i`$ (107)
and therefore we have
$$\mathrm{\Psi }=\chi \chi ^i\widehat{c}_i+\mathrm{}.=2\frac{G}{c^2}(\sqrt{2}M3P^i\widehat{c}_i)+\mathrm{}.$$
(108)
## 12 Appendix C
Relations between the following quantities have been used through this work:
$`\widehat{l}^a`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{2}}(1,{\displaystyle \frac{\zeta +\overline{\zeta }}{1+\zeta \overline{\zeta }}},i{\displaystyle \frac{\zeta \overline{\zeta }}{1+\zeta \overline{\zeta }}},{\displaystyle \frac{1+\zeta \overline{\zeta }}{1+\zeta \overline{\zeta }}}),`$ (109)
$`\widehat{m}^a`$ $`=`$ $`𝔡\widehat{l}^a={\displaystyle \frac{\sqrt{2}}{2}}(0,{\displaystyle \frac{1\overline{\zeta }^2}{1+\zeta \overline{\zeta }}},{\displaystyle \frac{i(1+\overline{\zeta }^2)}{1+\zeta \overline{\zeta }}},{\displaystyle \frac{2\overline{\zeta }}{1+\zeta \overline{\zeta }}}),`$
$`\widehat{\overline{m}}^a`$ $`=`$ $`\overline{𝔡}\widehat{l}^a={\displaystyle \frac{\sqrt{2}}{2}}(0,{\displaystyle \frac{1\zeta ^2}{1+\zeta \overline{\zeta }}},{\displaystyle \frac{i(1+\zeta ^2)}{1+\zeta \overline{\zeta }}},{\displaystyle \frac{2\zeta }{1+\zeta \overline{\zeta }}}),`$
$`\widehat{t}^a`$ $`=`$ $`\sqrt{2}(1,0,0,0),`$
$`\widehat{n}^a`$ $`=`$ $`\widehat{t}^a\widehat{l}^a={\displaystyle \frac{\sqrt{2}}{2}}(1,{\displaystyle \frac{\zeta +\overline{\zeta }}{1+\zeta \overline{\zeta }}},i{\displaystyle \frac{\zeta \overline{\zeta }}{1+\zeta \overline{\zeta }}},{\displaystyle \frac{1\zeta \overline{\zeta }}{1+\zeta \overline{\zeta }}}),`$
$`\widehat{c}^a`$ $`=`$ $`\widehat{l}^a\widehat{n}^a=\sqrt{2}(0,{\displaystyle \frac{\zeta +\overline{\zeta }}{1+\zeta \overline{\zeta }}},i{\displaystyle \frac{\zeta \overline{\zeta }}{1+\zeta \overline{\zeta }}},{\displaystyle \frac{1+\zeta \overline{\zeta }}{1+\zeta \overline{\zeta }}}),`$
with the products
$`(\widehat{m}_a\widehat{\overline{m}}_b\widehat{\overline{m}}_a\widehat{m}_b)`$ $`=`$ $`{\displaystyle \frac{i}{2}}ϵ_{abcd}\widehat{c}^ct^d,`$
$`(\widehat{m}_a\widehat{c}_b\widehat{c}_a\widehat{m}_b)`$ $`=`$ $`iϵ_{abcd}\widehat{m}^ct^d.`$
$`\widehat{m}_a\widehat{c}_b`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\widehat{m}_a\widehat{c}_b\widehat{m}_b\widehat{c}_a)+{\displaystyle \frac{1}{2}}(\widehat{m}_a\widehat{c}_b+\widehat{m}_b\widehat{c}_a)`$
$`=`$ $`i{\displaystyle \frac{1}{2}}ϵ_{abcd}\widehat{m}^ct^d+(l=2)Terms`$
and
$$ϵ_{abc}=\frac{1}{\sqrt{2}}ϵ_{abcd}t^d.$$
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# Local conversion of GHZ states to approximate W states
## Abstract
Genuine 3-qubit entanglement comes in two different inconvertible types represented by the GHZ state and the W state. We describe a specific method based on local positive operator valued measures (POVMs) and classical communication that can convert the ideal $`N`$-qubit GHZ state to a state arbitrarily close to the ideal $`N`$-qubit W state. We then experimentally implement this scheme in the 3-qubit case and characterize the input and the final state using 3-photon quantum state tomography.
Entanglement is at the heart of quantum mysteries and the power of quantum information. While 2-qubit entanglement is fairly well understood with good entanglement measures entmeasures , understanding multi-qubit entanglement remains a considerable challenge. With 3-qubits it is no longer enough simply to ask *if* the qubits are entangled or not; one must also ask *how* the qubits are entangled. The two most common examples of 3-qubit entangled states are the $`|\mathrm{GHZ}=1/\sqrt{2}\left(|HHH+|VVV\right)`$ ghz and $`|\mathrm{W}=1/\sqrt{3}\left(|HVV+|VHV+|VVH\right)`$ ZHG ; cirac , where $`|H`$ and $`|V`$ represent horizontally- and vertically- polarized photon states. An important characteristic of 3-particle GHZ states is that loss of any one of the qubits leaves the other two in a mixed state with only classical correlations. It is well-known that if two qubits are maximally entangled, then neither can be entangled to a third; the W state is the 3-qubit state in which each pair of qubits have the same and maximum amount of bipartite entanglement. This feature makes the entanglement of the W state maximally symmetrically robust against loss of any single qubit. It has been shown that the states $`|\mathrm{GHZ}`$ and $`|\mathrm{W}`$ represent two distinct classes of 3-qubit entanglement that cannot be interconverted under any local operations and classical communication (LOCC)cirac ; acin . Experimental realizations of GHZ states and more recently W states have been performed in optical and trapped ion experiments dik ; GHZexp ; 4ghz ; Wexp ; blatt .
While conversion of a GHZ state to an exact W state is not possible via LOCC, a specific scheme based on partial quantum measurement (positive operator valued measures or POVMs) and classical communication can, however, convert a GHZ state to an *arbitrarily good approximation* to a W state with a tradeoff between the fidelity of the final state and the probability of success. In any real experiment, where there is inevitable noise in state production and measurement, arbitrarily good approximations are indistinguishable from directly-prepared W states. In the present work, we first discuss a POVM scheme, inspired by the procrustean method procrustean , for converting the 3-qubit state $`|\mathrm{GHZ}`$ into an approximate $`|\mathrm{W}`$. Then we generalize the scheme to convert between $`N`$-qubit analogues of the GHZ and W states. Finally, we experimentally apply the POVM scheme to a 3-photon GHZ state and characterize the change from the input to the output using quantum state tomography tomography .
In the diagonal basis, where $`|D=1/\sqrt{2}\left(|H+|V\right)`$ and $`|A=1/\sqrt{2}\left(|H|V\right)`$, we can rewrite $`|\mathrm{GHZ}=1/2\left(|DDD+|DAA+|ADA+|AAD\right)`$. It is apparent in this basis that $`|\mathrm{GHZ}`$ is a superposition of an unwanted term, $`|DDD,`$ and a W-state. We define a local POVM with elements $`\epsilon _1=|AA\left|+a^2\right|DD|`$ and $`\epsilon _2=(1a^2)|DD|`$, where $`a`$ is a real number between 0 (perfect measurement) and 1 (no measurement). This POVM is applied to each photon in the state and if all of the parties find element $`\epsilon _1`$, then the new state is $`|\psi =𝒩\left[a^3/2|DDD+a\sqrt{3}/2|\mathrm{W}^{}\right],`$with the normalization constant, $`𝒩=2/\sqrt{a^6+3a^2}.`$ The state $`|\mathrm{W}^{}=1/\sqrt{3}\left(|DAA+|ADA+|AAD\right)`$ is simply related to $`|\mathrm{W}`$ by three single-qubit rotations. It is clear from this state, that the unwanted term $`|DDD`$ is reduced relative to the term $`|\mathrm{W}^{}`$. It is also clear that one achieves a pure $`|\mathrm{W}^{}`$ only in the limit as $`a0`$ where the probability of success also goes to zero. Nevertheless, the fidelity of this state with the desired $`|\mathrm{W}^{}`$ is $`_\mathrm{W}^{}=\left|\psi |\mathrm{W}^{}\right|^2=3/(a^4+3)`$, which rapidly rises from 3/4 to 1 as $`a`$ decreases, i.e., as the strength of the measurement increases. Conversely, the fidelity with the GHZ state $`_{\mathrm{GHZ}}=\left|\psi |\mathrm{GHZ}\right|^2=(a^4+6a^2+9)/(4a^4+12)`$ drops from 1 to 3/4 as $`a`$ decreases.
Our method can be generalized to convert an $`N`$-qubit GHZ state, $`|N+,`$ where $`|N\pm =1/\sqrt{2}\left(|H^N\pm |V^N\right)`$, into an arbitrarily good approximation to the $`N`$-qubit W state $`|\mathrm{W}_N^{}=1/\sqrt{N}\left(|DA\mathrm{}A+|AD\mathrm{}A+\mathrm{}+|AA\mathrm{}D\right)`$ generalW . We use the fact that the GHZ states $`|N\pm `$ satisfy the following relation:
$$|N\pm =\frac{1}{\sqrt{2}}\left[|\left(N\text{-}M\right)+|M\pm +|\left(N\text{-}M\right)|M\right],$$
(1)
where $`M<N`$. Notice that this factorization preserves the evenness or oddness in the number of negative signs. Through repeated application of these two rules, one can factor $`|N+`$ in terms of only single-qubit states $`|D|1+`$ and $`|A|1`$. This reexpresses the GHZ state as an equally weighted superposition of *all* $`2^{N1}`$ terms with an *even number* of $`|A`$s. When $`N`$ is odd, the GHZ state can be directly rewritten as,
$$|N+=\frac{1}{\sqrt{2^{N1}}}\left[\sqrt{N}|\mathrm{W}_N^{}+\sqrt{2^{N1}N}|\varphi \right],$$
(2)
where the state $`|\varphi `$ is a superposition of all those terms containing an odd number and at least 3 $`|D`$s. When $`N`$ is even, application of a local transformation $`|D|A,`$ $`|A|D`$ to any qubit allows the GHZ state to be written in the form of Eq. 2. Applying the same local POVM as in the 3-qubit case on each of the $`N`$ qubits, and given that each POVM returns element $`\epsilon _1`$ the unwanted amplitudes by *at least* a factor of $`a^3`$ while reducing the desired amplitude by only a single factor of $`a`$. In general, the fidelity of the resultant state by this prescription with $`|\mathrm{W}_N^{}`$ is given by
$$_\mathrm{W}^{}^N=\frac{2a^2N}{\left(1+a^2\right)^N\left(1a^2\right)^N}$$
regardless of whether $`N`$ is even or odd.
The details on our experimental method for creating 3-photon GHZ states can be found in ghztomo . Ultraviolet laser pulses from a frequency-doubled Ti:Sapphire laser make two passes through a type-II phase-matched $`\beta `$-barium borate (BBO) crystal aligned, with walk-off compensation to produce 2-photon pairs each in the Bell state $`|\varphi ^+`$ kwiatdownconv . These 2 independent photon pairs can be further entangled when combined at the polarizing beamsplitter (PBS1) and the four photons take four separate outputs A and B. Recall that a PBS works by reflecting horizontally-polarized light $`|H`$ and transmitting vertically-polarized light $`|V`$. Thus, two photons that were incident from different sides can only pass into different output modes when their polarizations were both $`|H`$ or both $`|V`$. In this sense, the PBS acts as a quantum parity check paritycheck . Given that the parity check succeeds on the two photons from the independent pairs, our state is transformed from the product state $`|\varphi ^+_{12}|\varphi ^+_{34}`$ to the 4-photon GHZ state $`|4+=1/\sqrt{2}\left(|HHHH_{AB14}+|VVVV_{AB14}\right)`$ 4ghz . We project photon 4 onto the state $`|D`$ and when this projection succeeds leaves $`|\mathrm{GHZ}=1/\sqrt{2}\left(|HHH_{AB1}+|VVV_{AB1}\right)`$.
A tomographically complete set of measurements for a 3-photon polarization state requires 64 polarization measurements. We use the 64 combinations of the single-photon projections $`|H,`$ $`|V,`$ $`|D,`$ and $`|R`$ on each of the 3 photons. These projections are implemented using a quarter-wave plate and polarizer for each of photons A and B, and a half- or quarter-wave plate and PBS2 for photon 1. Successful projections are signalled by four-photon coincidence measurements, 3 photons for the state and 1 trigger photon, using single-photon counting APDs and coincidence logic. The most-likely physical density matrix for our 3-qubits is extracted using maximum-likelihood reconstruction maxlike ; james .
We begin with the GHZ state that was characterized previously via 3-photon quantum state tomography ghztomo . We rewrite the density matrix in the $`|D/A`$ basis; this gives the density matrix shown in Fig. 2b (real part) and 2c (imaginary part). A comparison to the ideal GHZ written in the same basis is shown in Fig. 2a (real part only, the imaginary part is all zero). The colour plots display the absolute value of each element and show that the two matrices have the same structure.
Our POVM was implemented using three partial polarizers. Instead of orienting the polarizers in the $`|D/A`$ basis, we rotated the polarization of each photon by $`45^{}`$ using half-wave plates. These rotations were accomplished by using the existing half-wave plate in mode 1, and by adding two additional half-wave plates in modes A and B. These extra rotations allowed the partial polarizers to operate in the $`|H/V`$ basis, and therefore our POVM also operates in the $`|H/V`$ basis. Each polarizer comprised of two uncoated glass microscope slides such that the angle of incidence for the input light was at 56, near Brewster’s angle (Figure 1). The configuration of the plates are such that the beam experienced minimal additional transverse shift and maintained high coupling efficiency into single-mode fibres. Such partial polarizers have been used to study hidden nonlocality and entanglement concentration of maximally-entangled mixed states partpols . We placed each such element so that vertically-polarized light was $`P`$-polarized and horizontally-polarized light was $`S`$-polarized; we measured 88% transmission for the vertically-polarized light and 33% for the horizontally-polarized light. The transmission of the vertical light is thus only 38% of that for the horizontal and we can can describe the experiment using the POVM elements $`\epsilon _1`$ and $`\epsilon _2`$ with $`a^238\%`$. With this attenuation value, and beginning with the ideal GHZ state, the fidelity of the state with the ideal W state given 3 POVM outcomes $`\epsilon _1`$ is expected to increase from 75% to 95%.
We used the same 64 tomographic measurement settings for the W state as for the GHZ state. Data for each setting was accumulated for 1800 seconds and yielded a maximum of 120 four-fold coincidence counts (for the $`|VVV`$ projection). To account for laser power drift, which was small but not insignificant, we divided the four-folds by the square of the singles at the trigger detector. Background four-folds from a two-fold coincidence count and an uncorrelated accidental were estimated for each measurement setting and subtracted from the measured coincidences. Using the maximum likelihood reconstruction, our most likely density matrix is shown in Fig. 3b (real part) and 3c (imaginary part). The ideal W state density matrix consists of only 9 real elements – 3 diagonals of 1/3 height corresponding to $`|HVV`$, $`|VHV`$, and $`|VVH`$ and 6 maximal positive coherences between them. It is clear from the data that the dominant elements in the density matrix are those same 9 elements. In Fig. 3a., we show the effect of the POVM with our experimentally measured attenuation on the ideal GHZ state. The diagonal elements are attenuated much more strongly, by $`a^2`$, while the coherences remain maximal and are thus reduced only by $`a`$. Note that one specific unwanted term contained in the diagonal element for $`|VVV`$ is much more significant after the application of the POVM; this noise contribution, in the ideal case, is untouched by the POVM, and in the real case, least reduced.
We characterize the changes in our states using fidelity. The fidelity of a density matrix, $`\rho `$, with a pure quantum state, $`|\psi `$, is given by $`=\psi \left|\rho \right|\psi `$. We calculate this fidelity with a general GHZ (W) state, $`|\mathrm{GHZ}_G`$ $`(|\mathrm{W}_G)`$, which is related to $`|\mathrm{GHZ}`$ $`(|\mathrm{W})`$ by 3 local unitary rotations. The initial state has a fidelity of $`_{\mathrm{GHZ}_G}=\left(79.4\pm 1.6\right)\%`$ and $`_{\mathrm{W}_G}=\left(60.5\pm 1.9\right)\%`$ as compared with the ideal $`100\%`$ and $`75\%`$. After successful application of three local POVMs, $`_{\mathrm{GHZ}_G}=\left(59.8\pm 2.5\right)\%`$ and $`_{\mathrm{W}_G}=\left(68.4\pm 2.4\right)\%`$. Uncertainties in quantities extracted from these density matrices were calculated using a Monte Carlo routine and assumed Poissonian errors. A theoretical calculation based on our measured initial state and measured $`a`$ has yields a final state with $`_{\mathrm{W}_G}=75\%.`$ Thus much of the difference with the expected fidelity in the ideal case is a result of the quality of the initial state. Nevertheless, the overlap with a W state has been significantly improved while the overlap with a GHZ state has been strongly reduced.
We have described a method for converting $`N`$-qubit GHZ states to arbitrarily good approximations to $`N`$-qubit W states based on generalized quantum measurements (POVMs). We have implemented this scheme for the 3-qubit case and characterized the input and output states using multiphoton quantum state tomography. We have quantitatively shown that the transformation induced by the partial polarizers results in a decrease in the overlap of the state with a GHZ state while increasing the overlap with the desired W state. Multiparticle entanglement is essential to the success of quantum information processing. The theory and experimental work presented here extends our abilities to manipulate and understand the relationship between different types of complex entangled states.
The authors thank Antonio Aćin, Časlav Brukner, Klaus Hornberger, Morgan Mitchell, and Andrew White for valuable discussions. This work was supported by ARC Seibersdorf Research GmbH, the Austrian Science Foundation (FWF), project number SFB 015 P06, NSERC, and the European Commission, contract number IST-2001-38864 (RAMBOQ).
Figure 1. Experimental setup for production of GHZ state and its conversion to an approximate W state. A double-pass pulsed parametric down-conversion crystal (BBO) is used to create two pairs of polarization-entangled photons both in the state $`|\varphi ^+`$. Extra crystals (COMP) compensate for walkoff effects. A polarizing beamsplitter (PBS) performs a parity check on two of the photons, one from each pair; given that the parity check succeeds, signalled by the two photons taking different output modes, the emerging photons are in a four photon GHZ state. Projection of photon 4 onto the state $`|D`$ leaves the remaining 3 photons in the desired 3-photon GHZ state. Each photon in the GHZ state was rotated locally – photons A and B were rotated using additional half-wave plates (HWP) and photon 1 was rotated using the existing half-wave plate, rotating the polarization only $`45^{}`$ instead of $`90^{}`$. Our 3-qubit local POVM is performed using 3 partial polarizers (PP). Each partial polarizer consists of two microscope slides (MS) mounted such that the light is incident at $`56^{},`$ Brewster’s angle for $`n=1.5`$. This configuration had a measured transmission of 88% for p-polarization and 33% for s-polarization. Each POVM was oriented to reduce the horizontal component of the light relative to the vertical component. The three-photon polarization measurements for tomography were taken using quarter-wave plates (QWP) and rotatable polarizers (POL) for photons A and B, and a half- or quarter-wave plate and a fixed polarizer (actually a second PBS) for photon A before photon-counting detectors (DET and TRIG). A fixed QWP is mode A was used to compensate birefringence in the PBS.
Figure 2. Density matrix of an ideal and experimentally measured 3-photon GHZ state. The reconstructed density matrix of an a) ideal GHZ state (real part only) and our measured GHZ state, real part b) & imaginary part c), from reference ghztomo . The state is displayed in the $`D/A`$-basis, where $`D=\frac{1}{\sqrt{2}}\left(|H+|V\right)`$ and $`A=\frac{1}{\sqrt{2}}\left(|H|V\right)`$ are defined as in the text. The false colour plots display the absolute value of each component of the matrices and are meant to show the structure of the matrix - namely that our GHZ state is characterized in this basis by 4 diagonal elements of equal height with maximal positive coherences. The experimentally measured density matrix has a fidelity of 77% with the ideal GHZ state and 79% with any state related to the ideal GHZ state via local unitary transformations.
Figure 3. Density matrix of approximate W output state after the POVM procedure. The reconstructed density matrix of our output state, b) real part & c) imaginary part, after the three local POVM operations. The application of the POVMs have suppressed several of the matrix components such that the final state contains only 9 major elements. These are the same 9 elements for the ideal W state. The operation has increased the fidelity of our state with a W state from 61% to 68% while at the same time reducing the fidelity of our state with a GHZ state from 79% to 60%. For comparison, we show a) the action of the POVM operation on the ideal GHZ state. Although this state still contains large coherences with the HHH component, this state has 95% fidelity with the W state.
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