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# On the origin of reentrance in 2D Josephson Junction Arrays
## I Introduction
According to the current paradigm, paramagnetic Meissner effect (PME) can be related to the presence of $`\pi `$-junctions , either resulting from the presence of magnetic impurities in the junction or from unconventional pairing symmetry . Other possible explanations of this phenomenon are based on flux trapping and flux compression effects including also an important role of the surface of the sample . Besides, in the experiments with unshunted 2D-JJA, we have previously reported that PME manifests itself through a dynamic reentrance (DR) of the AC magnetic susceptibility as a function of temperature. These results have been further corroborated by Nielsen et al. and De Leo et al. who argued that PME can be simply related to magnetic screening in multiply connected superconductors. So, the main question is: which parameters are directly responsible for the presence (or absence) of DR in artificially prepared arrays?
Previously (also within the single plaquette approximation), Barbara et al. have briefly discussed the effects of varying $`\beta _L`$ on the observed dynamic reentrance with the main emphasis on the behavior of 2D-JJA samples with high (and fixed) values of $`\beta _C`$. However, to our knowledge, up to date no systematic study (either experimental or theoretical) has been done on how the $`\beta _C`$ value itself affects the reentrance behavior. In the present work, by a comparative study of the magnetic properties of shunted and unshunted 2D-JJA, we propose an answer to this open question. Namely, by using experimental and theoretical results, we will demonstrate that only arrays with sufficiently large value of the Stewart-McCumber parameter $`\beta _C`$ will exhibit the dynamic reentrance behavior (and hence PME).
## II Experimental Results
To measure the complex AC susceptibility in our arrays we used a high-sensitive home-made susceptometer based on the so-called screening method in the reflection configuration . The experimental system was calibrated by using a high-quality niobium thin film. Previously , we have shown that the calibrated output complex voltage of the experimental setup corresponds to the complex AC susceptibility.
To experimentally investigate the origin of the reentrance, we have measured $`\chi ^{}(T)`$ for three sets of shunted and unshunted samples obtained from different makers (Westinghouse and Hypress) under the same conditions of the amplitude of the excitation field $`h_{ac}`$ ($`1mOe<h_{ac}<10Oe`$), external magnetic field $`H_{dc}`$ ($`0<H_{dc}<500Oe`$) parallel to the plane of the sample, and frequency of AC field $`\omega =2\pi f`$ (fixed at $`f=20kHz`$). Unshunted 2D-JJAs are formed by loops of niobium islands linked through $`NbAlO_xNb`$ Josephson junctions while shunted 2D-JJAs have a molybdenum shunt resistor (with $`R_{sh}2.2\mathrm{\Omega }`$) short-circuiting each junction (see Fig.1). Both shunted and unshunted samples have rectangular geometry and consist of $`100\times 150`$ tunnel junctions. The unit cell for both types of arrays has square geometry with lattice spacing $`a46\mu m`$ and a single junction area of $`5\times 5\mu m^2`$. The critical current density for the junctions forming the arrays is about $`600A/cm^2`$ at $`4.2K`$. Besides, for the unshunted samples $`\beta _C(4.2K)30`$ and $`\beta _L(4.2K)30`$, while for shunted samples $`\beta _C(4.2K)1`$ and $`\beta _L(4.2K)30`$ where $`\beta _L(T)=\frac{2\pi LI_C(T)}{\mathrm{\Phi }_0}`$ and $`\beta _C(T)=\frac{2\pi C_JR_J^2I_C(T)}{\mathrm{\Phi }_0}`$. Here, $`C_J0.58pF`$ is the capacitance, $`R_J10.4\mathrm{\Omega }`$ the quasi-particle resistance (of unshunted array), and $`I_C(4.2K)150\mu A`$ critical current of the Josephson junction. $`\mathrm{\Phi }_0`$ is the quantum of magnetic flux. The parameter $`\beta _L`$ is proportional to the number of flux quanta that can be screened by the maximum critical current in the junctions, while the Stewart-McCumber parameter $`\beta _C`$ basically reflects the quality of the junctions in arrays.
It is well established that both magnetic and transport properties of any superconducting material can be described via a two-component response , the intragranular (associated with the grains exhibiting bulk superconducting properties) and intergranular (associated with weak-link structure) contributions . Likewise, artificially prepared JJAs (consisting of superconducting islands, arranged in a symmetrical periodic lattice and coupled by Josephson junctions) will produce a similar response .
Since our shunted and unshunted samples have the same value of $`\beta _L`$ and different values of $`\beta _C`$, it is possible to verify the dependence of the reentrance effect on the value of the Stewart-McCumber parameter. For the unshunted 2D-JJA (Fig. 2a) we have found that for an AC field lower than $`50mOe`$ (when the array is in the Meissner-like state) the behavior of $`\chi ^{}(T)`$ is quite similar to homogeneous superconducting samples, while for $`h_{ac}>50mOe`$ (when the array is in the mixed-like state with practically homogeneous flux distribution) these samples exhibit a clear reentrant behavior of susceptibility . At the same time, the identical experiments performed on the shunted samples produced no evidence of any reentrance for all values of $`h_{ac}`$ (see Fig. 2b). It is important to point out that the analysis of the experimentally obtained imaginary component of susceptibility $`\chi ^{\prime \prime }(T)`$ shows that for the highest AC magnetic field amplitudes (of about $`200mOe`$) dissipation remains small. Namely, for typical values of the AC amplitude, $`h_{ac}=100mOe`$ (which corresponds to about $`10`$ vortices per unit cell) the imaginary component is about $`15`$ times smaller than its real counterpart. Hence contribution from the dissipation of vortices to the observed phenomena can be safely neglected.
To further study this unexpected behavior we have also performed experiments where we measure $`\chi ^{}(T)`$ for different values of $`H_{dc}`$ keeping the value of $`h_{ac}`$ constant. The influence of DC fields on reentrance in unshunted samples is shown in Fig. 3. On the other hand, the shunted samples still show no signs of reentrance, following a familiar pattern of field-induced gradual diminishing of superconducting phase (very similar to a zero DC field flat-like behavior seen in Fig.2b).
To understand the influence of DC field on reentrance observed in unshunted arrays, it is important to emphasize that for our sample geometry this parallel field suppresses the critical current $`I_C`$ of each junction without introducing any detectable flux into the plaquettes of the array. Thus, a parallel DC magnetic field allows us to vary $`I_C`$ independently from temperature and/or applied perpendicular AC field. The measurements show (see Fig.3) that the position of the reentrance is tuned by $`H_{dc}`$. We also observe that the value of temperature $`T_{min}`$ (at which $`\chi ^{}(T)`$ has a minimum) first shifts towards lower temperatures as we raise $`H_{dc}`$ (for small DC fields) and then bounces back (for higher values of $`H_{dc}`$). This non-monotonic behavior is consistent with the weakening of $`I_C(T)`$ and corresponds to Fraunhofer-like dependence of the Josephson junction critical current on DC magnetic field applied in the plane of the junction. We measured $`I_C`$ from transport current-voltage characteristics, at different values of $`H_{dc}`$ at $`T=4.2K`$ and found that $`\chi ^{}(T=4.2K)`$, obtained from the isotherm $`T=4.2K`$ (similar to that given in Fig.3), shows the same Fraunhofer-like dependence on $`H_{dc}`$ as the critical current $`I_C(H_{dc})`$ of the junctions forming the array (see Fig.4). This gives further proof that only the junction critical current is varied in this experiment. This also indicates that the screening currents at low temperature (i.e., in the reentrant region) are proportional to the critical currents of the junctions. In addition, this shows an alternative way to obtain $`I_C(H_{dc})`$ dependence in big arrays. And finally, a sharp Fraunhofer-like pattern observed in both arrays clearly reflects a rather strong coherence (with negligible distribution of critical currents and sizes of the individual junctions) which is based on highly correlated response of all single junctions forming the arrays, thus proving their high quality. Such a unique behavior of Josephson junctions in our samples provides a necessary justification for suggested theoretical interpretation of the obtained experimental results. Namely, based on the above-mentioned properties of our arrays, we have found that practically all the experimental results can be explained by analyzing the dynamics of just a single unit cell in the array.
## III Theoretical Interpretation and Numerical Simulations
To understand the different behavior of the AC susceptibility observed in shunted and unshunted 2D-JJAs, in principle one would need to analyze in detail the flux dynamics in these arrays. However, as we have previously reported , because of the well-defined periodic structure of our arrays (with no visible distribution of junction sizes and critical currents), it is reasonable to expect that the experimental results obtained from the magnetic properties of our 2D-JJAs can be quite satisfactory explained by analyzing the dynamics of a single unit cell (plaquette) of the array. An excellent agreement between a single-loop approximation and the observed behavior (seen through the data fits) justifies a posteriori our assumption. It is important to mention that the idea to use a single unit cell to qualitatively understand PME was first suggested by Auletta et al. . They simulated the field-cooled DC magnetic susceptibility of a single-junction loop and found a paramagnetic signal at low values of external magnetic field.
In our calculations and numerical simulations, the unit cell is a loop containing four identical Josephson junctions and the measurements correspond to the zero-field cooling (ZFC) AC magnetic susceptibility. We consider the junctions of the single unit cell as having capacitance $`C_J`$, quasi-particle resistance $`R_J`$ and critical current $`I_C`$. We have used this simple four-junctions model to study the magnetic behavior of our 2D-JJA by calculating the AC complex magnetic susceptibility $`\chi =\chi ^{}+i\chi ^{\prime \prime }`$ as a function of $`T`$, $`\beta _C`$ and $`\beta _L`$. Specifically, shunted samples are identified through low values of the McCumber parameter ($`\beta _C1`$) while high values ($`\beta _C1`$) indicate an unshunted 2D-JJA.
If we apply an AC external field $`B_{ac}(t)=\mu _0h_{ac}\mathrm{cos}\omega t`$ normally to the 2D-JJA and a DC field $`B_{dc}=\mu _0H_{dc}`$ parallel to the array, then the total magnetic flux $`\mathrm{\Phi }(t)`$ threading the four-junction superconducting loop is given by $`\mathrm{\Phi }(t)=\mathrm{\Phi }_{ext}(t)+LI(t)`$ where $`L`$ is the loop inductance, $`\mathrm{\Phi }_{ext}(t)=SB_{ac}(t)+ldB_{dc}`$ is the flux related to the applied magnetic field (with $`l\times d`$ being the size of the single junction area, and $`Sa^2`$ being the projected area of the loop), and the circulating current in the loop reads
$$I(t)=I_C(T)\mathrm{sin}\varphi _i(t)+\frac{\mathrm{\Phi }_0}{2\pi R_J}\frac{d\varphi _i}{dt}+\frac{C_J\mathrm{\Phi }_0}{2\pi }\frac{d^2\varphi _i}{dt^2}$$
(1)
Here $`\varphi _i(t)`$ is the gauge-invariant superconducting phase difference across the $`i`$th junction, and $`\mathrm{\Phi }_0`$ is the magnetic flux quantum.
Since the inductance of each loop is $`L=\mu _0a64pH`$ and the critical current of each junction is $`I_C150\mu A`$, for the mixed-state region (above $`50mOe`$) we can safely neglect the self-field effects because in this region $`LI(t)`$ is always smaller than $`\mathrm{\Phi }_{ext}(t)`$. Besides, since the length $`l`$ and the width $`w`$ of each junction in our array is smaller than the Josephson penetration depth $`\lambda _J=\sqrt{\mathrm{\Phi }_0/2\pi \mu _0dj_{c0}}`$ (where $`j_{c0}`$ is the critical current density of the junction, and $`d=2\lambda _L+\xi `$ is the size of the contact area with $`\lambda _L(T)`$ being the London penetration depth of the junction and $`\xi `$ an insulator thickness), namely $`lw5\mu m`$ and $`\lambda _J20\mu m`$ (using $`j_{c0}600A/cm^2`$ and $`\lambda _L39nm`$ for $`Nb`$ at $`T=4.2K`$), we can adopt the small-junction approximation for the gauge-invariant superconducting phase difference across the $`i`$th junction (for simplicity we assume as usual that $`\varphi _1=\varphi _2=\varphi _3=\varphi _4\varphi _i`$)
$$\varphi _i(t)=\varphi _0(H_{dc})+\frac{2\pi B_{ac}(t)S}{\mathrm{\Phi }_0}$$
(2)
where $`\varphi _0(H_{dc})=\varphi _0(0)+2\pi \mu _0H_{dc}dl/\mathrm{\Phi }_0`$ with $`\varphi _0(0)`$ being the initial phase difference.
To properly treat the magnetic properties of the system, let us introduce the following Hamiltonian
$$(t)=J\underset{i=1}{\overset{4}{}}[1\mathrm{cos}\varphi _i(t)]+\frac{1}{2}LI^2(t)$$
(3)
which describes the tunneling (first term) and inductive (second term) contributions to the total energy of a single plaquette. Here, $`J(T)=(\mathrm{\Phi }_0/2\pi )I_C(T)`$ is the Josephson coupling energy.
The real part of the complex AC susceptibility is defined as
$$\chi ^{}(T,h_{ac},H_{dc})=\frac{M}{h_{ac}}$$
(4)
where
$$M(T,h_{ac},H_{dc})=\frac{1}{V}\frac{}{h_{ac}}$$
(5)
is the net magnetization of the plaquette. Here $`V`$ is the sample’s volume, and $`<\mathrm{}>`$ denotes the time averaging over the period $`2\pi /\omega `$, namely
$$<A>=\frac{1}{2\pi }_0^{2\pi }d(\omega t)A(t)$$
(6)
Taking into account the well-known analytical approximation of the BCS gap parameter (valid for all temperatures), $`\mathrm{\Delta }(T)=\mathrm{\Delta }(0)\mathrm{tanh}\left(2.2\sqrt{\frac{T_cT}{T}}\right)`$ for the explicit temperature dependence of the Josephson critical current
$$I_C(T)=I_C(0)\left[\frac{\mathrm{\Delta }(T)}{\mathrm{\Delta }(0)}\right]\mathrm{tanh}\left[\frac{\mathrm{\Delta }(T)}{2k_BT}\right]$$
(7)
we successfully fitted all our data using the following set of parameters: $`\varphi _0(0)=\frac{\pi }{2}`$ (which corresponds to the maximum Josephson current within a plaquette), $`\beta _L(0)=32`$, $`\beta _C(0)=32`$ (for unshunted array) and $`\beta _C(0)=1.2`$ (for shunted array). The corresponding fits are shown by solid lines in Figs.2 and 3 for the experimental values of AC and DC field amplitudes.
In the mixed-state region and for low enough frequencies (this assumption is well-satisfied because in our case $`\omega \omega _{LR}`$ and $`\omega \omega _{LC}`$ where $`\omega _{LR}=R/L`$ and $`\omega _{LC}=1/\sqrt{LC}`$ are the two characteristic frequencies of the problem) from Eqs.(3)-(6) we obtain the following approximate analytical expression for the susceptibility of the plaquette
$`\chi ^{}(T,h_{ac},H_{dc})`$ $``$ $`\chi _0(T)[\beta _L(T)f_1(b)\mathrm{cos}\left({\displaystyle \frac{2H_{dc}}{H_0}}\right)`$ (8)
$`+`$ $`f_2(b)\mathrm{sin}\left({\displaystyle \frac{H_{dc}}{H_0}}\right)\beta _C^1(T)]`$ (9)
where $`\chi _0(T)=\pi S^2I_C(T)/V\mathrm{\Phi }_0`$, $`H_0=\mathrm{\Phi }_0/(2\pi \mu _0dl)10Oe`$, $`f_1(b)=J_0(2b)J_2(2b)`$, and $`f_2(b)=J_0(b)bJ_1(b)3J_2(b)+bJ_3(b)`$ with $`b=2\pi S\mu _0h_{ac}/\mathrm{\Phi }_0`$ and $`J_n(x)`$ being the Bessel function of the $`n`$th order.
Notice also that the analysis of Eq.(8) reproduces the observed Fraunhofer-like behavior of the susceptibility in applied DC field (see Fig.4) and the above-mentioned fine tuning of the reentrance effect (see also Ref.13). Indeed, according to Eq.(8) (and in agreement with the observations), for small DC fields the minimum temperature $`T_{min}`$ (indicating the beginning of the reentrant transition) varies with $`H_{dc}`$ as follows, $`1T_{min}/T_CH_{dc}/H_0`$.
To further test our interpretation and verify the influence of the parameter $`\beta _C`$ on the reentrance, we have also performed extensive numerical simulations of the four-junction model previously described but without a simplifying assumption about the explicit form of the phase difference based on Eq.(2). More precisely, we obtained the temperature behavior of the susceptibility by solving the set of equations responsible for the flux dynamics within a single plaquette and based on Eq.(1) for the total current $`I(t)`$, the equation for the total flux $`\mathrm{\Phi }(t)=\mathrm{\Phi }_{ext}(t)+LI(t)`$ and the flux quantization condition for four junctions, namely $`\varphi _i(t)=\frac{\pi }{2}\left(n+\frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}\right)`$ where $`n`$ is an integer. Both Euler and fourth-order Runge-Kutta integration methods provided the same numerical results. In Fig.5 we show the real component of the simulated susceptibility $`\chi (T)`$ corresponding to the fixed value of $`\beta _C(T=4.2K)=1`$ (shunted samples) and different values of $`\beta _L(T=4.2K)=1`$, $`10`$, $`15`$, $`20`$, $`30`$, $`40`$, $`50`$, $`60`$, $`90`$, $`150`$ and $`200`$. As expected, for this low value of $`\beta _C`$ reentrance is not observed for any values of $`\beta _L`$. On the other hand, Fig.6 shows the real component of the simulated $`\chi (T)`$ but now using fixed value of $`\beta _L(T=4.2K)=30`$ and different values of $`\beta _C(T=4.2K)=1`$, $`2`$, $`5`$, $`10`$, $`20`$, $`30`$ and $`100`$. This figure clearly shows that reentrance appears for values of $`\beta _C>20`$. In both cases we used $`h_{ac}=70mOe`$. We have also simulated the curve for shunted ($`\beta _L=30`$, $`\beta _C=1`$) and unshunted ($`\beta _L=30`$, $`\beta _C=30`$) samples for different values of $`h_{ac}`$ (see Fig.7). In this case the values of the parameters $`\beta _L`$ and $`\beta _C`$ were chosen from our real 2D-JJA samples. Again, our simulations confirm that dynamic reentrance does not occur for low values of $`\beta _C`$, independently of the values of $`\beta _L`$ and $`h_{ac}`$.
The following comment is in order regarding some irregularities visibly seen in Figs.(5)-(7) around the transition regions from non-reentrant to reentrant behavior. It is important to emphasize that the above irregularities are just artifacts of the numerical simulations due to the conventional slow-converging real-time reiteration procedures . They neither correspond to any experimentally observed behavior (within the accuracy of the measurements technique and data acquisition), nor they reflect any irregular features of the considered here theoretical model (which predicts a smooth temperature dependence seen through the data fits). As usual, to avoid this kind of artificial (non-physical) discontinuity, more powerful computers are needed.
Based on the above extensive numerical simulations, a resulting phase diagram $`\beta _C\beta _L`$ (taken for $`T=1K`$, $`h_{ac}=70mOe`$, and $`H_{dc}=0`$) is depicted in Fig.8 which clearly demarcates the border between the reentrant (white area) and non-reentrant (shaded area) behavior in the arrays for different values of $`\beta _L(T)`$ and $`\beta _C(T)`$ parameters at given temperature. In other words, if $`\beta _L`$ and $`\beta _C`$ parameters of any realistic array have the values inside the white area, this array will exhibit a reentrant behavior. In addition, this diagram shows that one can prepare a reentrance exhibiting array by changing one of the parameters (usually, it is much easier to change $`\beta _C`$ by tuning the shunt resistance rather than the geometry related inductance parameter $`\beta _L`$).
It is instructive to mention that a hyperbolic-like character of $`\beta _L`$ vs $`\beta _C`$ law (seen in Fig.8) is virtually present in the approximate analytical expression for the susceptibility of the plaquette given by Eq.(8) (notice however that this expression can not be used to produce any quantitative prediction because the neglected in Eq.(8) frequency-related terms depend on $`\beta _L`$ and $`\beta _C`$ parameters as well). A qualitative behavior of the envelope of the phase diagram (depicted in Fig.8) with DC magnetic field $`H_{dc}`$ (for $`T=1K`$ and $`h_{ac}=70mOe`$), obtained using Eq.(8), is shown in Fig.9.
And finally, to understand how small values of $`\beta _C`$ parameter affect the flux dynamics in shunted arrays, we have analyzed the $`\mathrm{\Phi }_{tot}`$($`\mathrm{\Phi }_{ext}`$) diagram. Similarly to those results previously obtained from unshunted samples , for a shunted sample at fixed temperature this curve is also very hysteretic (see Fig.10). In both cases, $`\mathrm{\Phi }_{tot}`$ vs. $`\mathrm{\Phi }_{ext}`$ shows multiple branches intersecting the line $`\mathrm{\Phi }_{tot}=0`$ which corresponds to diamagnetic states. For all the other branches, the intersection with the line $`\mathrm{\Phi }_{tot}=\mathrm{\Phi }_{ext}`$ corresponds to the boundary between diamagnetic states (negative values of $`\chi ^{}`$) and paramagnetic states (positive values of $`\chi ^{}`$). As we have reported before , for unshunted 2D-JJA at temperatures below $`7.6K`$ the appearance of the first and third branches adds a paramagnetic contribution to the average value of $`\chi ^{}`$. When $`\beta _C`$ is small (shunted arrays), the analysis of these curves shows that there is no reentrance at low temperatures because in this case the second branch appears to be energetically stable, giving an extra diamagnetic contribution which overwhelms the paramagnetic contribution from subsequent branches. In other words, for low enough values of $`\beta _C`$ (when the samples are ZFC and then measured at small values of the magnetic field), most of the loops will be in the diamagnetic states, and no paramagnetic response is registered. As a result, the flux quanta cannot get trapped into the loops even by the following field-cooling process in small values of the magnetic field. In this case the superconducting phases and the junctions will have the same diamagnetic response and the resulting measured value of the magnetic susceptibility will be negative (i.e., diamagnetic) as well. On the other hand, when $`\beta _C`$ is large enough (unshunted arrays), the second branch becomes energetically unstable, and the average response of the sample at low temperatures is paramagnetic (Cf. Fig.7 from Ref. ).
In conclusion, our experimental and theoretical results have demonstrated that the reentrance phenomenon (and concomitant PME) in artificially prepared Josephson Junction Arrays is related to the damping effects associated with the Stewart-McCumber parameter $`\beta _C`$. Namely, reentrant behavior of AC susceptibility takes place in the underdamped (unshunted) array (with large enough value of $`\beta _C`$) and totally disappears in overdamped (shunted) arrays.
###### Acknowledgements.
We thank P. Barbara, C.J. Lobb, R.S. Newrock, and A. Sanchez for useful discussions. S.S. and F.M.A.M. gratefully acknowledge financial support from Brazilian Agency FAPESP under grant 03/00296-5.
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# 1 Introduction.
## 1 Introduction.
The mean problem we have to face in performing the quantization of the gravitational field is the resulting non evolutionary dynamics, following by the operator translation of the theory constraints. In our points of view, the problem is relied to an ambiguous definition of the label time, in the splitting procedure extended to the quantum regime , , .
Indeed, as shown in , , , , , , , a dualism exists in quantum gravity, between including a reference fluid into the dynamics and explicitly breaking down time displacement invariance of the theory. According to this paradigm, we get a self consistent canonical quantization of gravity with the properly evolutionary features, only when a physical reference frame is included into to the system evolution and make possible a physical (3+1)-slicing of the spacetime. In this sense the analysis here presented completes the proof of the dualism inferred in ;
In fact, we start by assigning a time variable, suitable to construct the T-product functions and show how, in the classical limit of the resulting evolutionary quantum gravity, an Eckart fluid appears.
This issue provides a physical meaning to the slicing procedure, in view of the time-like character of the fluid 4-velocity field. Thus, we see how a correct implementation of the classical constraints toward the quantum dynamics, provides a clock for theory and restore the time evolution of the state functional.
In section 2 we present the procedure by which a generating functional can be built, in the case of the scalar field, starting from a field theory approach. The result is achieved using both the Lagrangian and the Hamiltonian formalism, after extending the generating functional definition to the phase-space.
In section 3, we resume the main steps by which the canonical quantization of gravity is performed, i.e. the (3+1)-slicing procedure and its consequences on the dynamics, like the frozen formalism, are discussed.
Section 4 deals with the gravitational Hamilton equations, as rewritten in a new set of variables, the Misner-like ones; they are introduced in order to make possible the evaluation of a gravitational generating functional $`Z`$ referred to the phase-space.
In the section 5, the procedure illustrated for the scalar field is implemented to the gravitational Hamilton equations and the generating functional is explicitly built up.
Through the generating functional formalism, we relax, in the quantum regime, the primary constraints of the gravitational theory, handling on the functional $`𝒞`$, taken on the Lagrangian multipliers $`\xi `$ and $`\xi ^i`$. In fact, the assumption of a physical time at the ground of the generating functional formalism implies that the corresponding diffeomorphisms invariance is broken, in the quantum gravity approach. Our point of view is that the problem of the stationary evolution for the gravitational quantum field, stands in a wrong implementation, to the quantum sector, of the classical primary constraints (as well as the implementation of the secondary ones). By the constraints relaxation performed in section 6, we obtain, in the limit of small $`\mathrm{}`$, a Schrödinger equation describing the evolution of the quantum gravitational field. We provide the analysis resulting in this equation in section 7.
At the end, in section 8, we talk about the classical limit of the founded dynamical equation. This limit accounts for the presence, in the system, of a fluid, having the Eckart form. The presence of this fluid provides the label time $`t`$ with a physical meaning, that is it assures us about the existence of a non arbitrary time for the considered gravitational system. In section 9 brief concluding remarks follow.
## 2 Generating functional in the scalar field case.
As first step, we present the method of searching for the generating functional, used in the scalar field case ; we show it in both the Lagrangian and Hamiltonian formalism. Then, in section 5, we will follow this method in order to find the gravitational generating functional expression. We start to consider a self-interacting scalar field $`\phi (x)`$, described by the following action
$$S=\left[\frac{1}{2}_\nu \phi (y)^\nu \phi (y)\frac{1}{2}\mu ^2\phi ^2(y)\frac{1}{4}g_0\phi ^4(y)\right]𝑑y^4,$$
where $`\mu `$ and $`g_o`$ denote assigned constants; the associated classical equation of motion read as
$$\frac{\delta S}{\delta \phi \left(x\right)}=\left(\mathrm{}_x+\mu ^2\right)\phi \left(x\right)+g_0\phi ^3\left(x\right)=0$$
(1)
where $`x=(\stackrel{}{x},t)`$ denotes event coordinates.
Now we consider $`\phi `$ as a quantum scalar field, satisfying the canonical (equal times) commutation relations,
$`[\widehat{\phi }(x),\widehat{\phi }\left(y\right)]_{x^0=y^0}`$ $`=\mathrm{\hspace{0.33em}0}`$ (2)
$`[\widehat{\pi }(x),\widehat{\pi }\left(y\right)]_{x^0=y^0}`$ $`=\mathrm{\hspace{0.33em}0}`$
$`[\widehat{\pi }\left(x\right),\widehat{\phi }\left(y\right)]_{x^0=y^0}`$ $`=i\mathrm{}\delta ^3\left(\stackrel{}{x}\stackrel{}{y}\right),`$
being $`\pi (x)\frac{\phi }{t}`$ the conjugate momentum to $`\phi (x)`$.
We define the *two-points T-product* as
$`0\left|T\left(\phi \left(x\right)\phi \left(y\right)\right)\right|0`$ $`=`$ $`+\theta \left(x^0y^0\right)0\left|\left(\phi \left(x\right)\phi \left(y\right)\right)\right|0`$
$`+\theta \left(y^0x^0\right)0\left|\left(\phi \left(y\right)\phi \left(x\right)\right)\right|0;`$
above, we assumed the existence of a vacuum state for the theory, denoted by $`|0`$. (For a field theory, based on the ground state expectation values, see )
Let us search for restrictions on the T-products general form, implied by (1), via the commutation relations (2). Starting with the expectation value of the classical equation, taken on the ground state,
$`0\left|T\left(\left[\mathrm{}_x\phi \left(x\right)+\mu ^2\phi \left(x\right)+g_0\phi ^3\left(x\right)\right]\phi \left(y\right)\right)\right|0=0,`$
we arrive to an equation for the two-points T-products
$$\left[\mathrm{}_x+\mu ^2\right]0\left|T\left(\phi \left(x\right)\phi \left(y\right)\right)\right|0+g_00\left|T\left(\phi ^3\left(x\right)\phi \left(y\right)\right)\right|0=i\mathrm{}\delta ^4\left(xy\right)$$
(For a T-products theory see ).
This equation can be generalized in view of n-points T-products as follows:
$`\left[\mathrm{}_x+\mu ^2\right]0\left|T\left(\phi \left(x\right)\phi \left(x_1\right)\mathrm{}\varphi \left(x_n\right)\right)\right|0+g_00\left|T\left(\phi ^3\left(x\right)\phi \left(x_1\right)\mathrm{}\phi \left(x_n\right)\right)\right|0`$
$`=i\mathrm{}{\displaystyle \underset{i=1}{\overset{n}{}}}\delta ^4\left(xx_i\right)0\left|T\left(\phi \left(x_1\right)\mathrm{}\phi \left(x_{i1}\right)\phi \left(x_{i+1}\right)\mathrm{}\phi \left(x_n\right)\right)\right|0`$ (3)
Introducing the generating functional $`Z(J)`$ (, , ) is useful to rewrite the whole set of the equation (3), namely for any value of n, in a compact form.
We assign the source current $`J(x)`$ and let us construct the quantity
$$Z(J)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{i^n}{n!}𝑑x_1^4\mathrm{}𝑑x_n^40\left|T\left(\phi \left(x_1\right)\mathrm{}\phi \left(x_n\right)\right)\right|0J\left(x_1\right)\mathrm{}J\left(x_n\right),$$
which satisfies the normalization condition $`Z(0)=1`$;
in a shorter form, $`Z`$ can be restated in terms of the following vacuum expectation value
$$Z\left(J\right)=0\left|T\left(e^{i{\scriptscriptstyle J\phi }}\right)\right|0,$$
where we have:
$$iJ\phi =iJ\left(x\right)\phi \left(x\right)𝑑x^4.$$
We show explicitly the link between the T-product functions and the functional derivatives of $`Z`$ as
$$0\left|T\left(\phi \left(x_1\right)\mathrm{}\phi \left(x_n\right)\right)\right|0=\frac{1}{i^n}\frac{\delta ^nZ\left(J\right)}{\delta J\left(x_n\right)\mathrm{}\delta J\left(x_1\right)}|_{J=0};$$
so we can get the equation for any T-product by performing the derivatives of the following fundamental expression
$$\left[\mathrm{}_x+\mu ^2\right]\frac{1}{i}\frac{\delta Z\left(J\right)}{\delta J\left(x\right)}+g_0\frac{1}{i^3}\frac{\delta ^3Z\left(J\right)}{\delta ^3J\left(x\right)}=\mathrm{}J\left(x\right)Z\left(J\right).$$
(4)
To solve (4), we introduce $`\stackrel{~}{Z}(\phi )`$, the functional Fourier transformation of $`Z(J)`$,
$$Z\left(J\right)=\underset{x}{}d\phi (x)\stackrel{~}{Z}(\phi )e^{i{\scriptscriptstyle \phi (x)J(x)𝑑x^4}}$$
In the next, we set the Lebesgue measure as $`\delta \phi =_xd\phi \left(x\right)`$.
If we substitute, into (4), the Fourier expansion of $`Z`$, we get the equation
$$\delta \phi e^{i{\scriptscriptstyle \phi J}}\left\{\stackrel{~}{Z}\left(\phi \right)\left[\left(\mathrm{}_x+\mu ^2\right)\phi \left(x\right)+g_0\phi ^3\left(x\right)\right]+\frac{\mathrm{}}{i}\frac{\delta \stackrel{~}{Z}\left(\phi \right)}{\delta \phi \left(x\right)}\right\}=0.$$
With some algebra and using (1), we arrive to the expression
$`{\displaystyle \frac{\delta }{\delta \phi \left(x\right)}}\left[ln\stackrel{~}{Z}(\phi ){\displaystyle \frac{i}{\mathrm{}}}S(\phi )\right]=0,`$
by which we finally recognize the functional $`\stackrel{~}{Z}`$ in the form
$$\stackrel{~}{Z}(J)=Ce^{\frac{i}{\mathrm{}}S\left(\phi \right)},$$
so that the generating functional follows:
$$Z\left(J\right)=C\delta \phi (x)e^{\frac{i}{\mathrm{}}S\left(\phi \right)+i{\scriptscriptstyle J(x)\phi (x)d^4x}},$$
being $`C`$ the normalization constant.
In view of applying this procedure to the gravitational case, we have to translate it, into the Hamiltonian formalism, because of some difficulties that Lagrangian formalism can not overcome. So, we get the Hamiltonian of the system out of the scalar field action,
$$H=\left[\frac{1}{2}\pi ^2\frac{1}{2}^j\phi _j\phi +\frac{1}{2}\mu ^2\phi ^2+\frac{1}{4}g_0\phi ^4\right]𝑑y^3,$$
via the above conjugate momentum of $`\phi (x)`$.
The Hamiltonian equations we find, look as follows:
$`\dot{\pi }`$ $`=`$ $`{\displaystyle \frac{\delta H}{\delta \phi }}=\left(_j^2\phi \mu ^2\phi \right)g_0\phi ^3`$ (5)
$`\dot{\phi }`$ $`=`$ $`+{\displaystyle \frac{\delta H}{\delta \pi }}=\pi .`$ (6)
Like before, we transform each equation of motion in a T-products one, using again the commutation relations (2)
$$\{\begin{array}{ccc}_t0\left|T\left[\pi (x)\phi (y)\right]\right|0\hfill & +\hfill & _j^20\left|T\left[\phi (x)\phi (y)\right]\right|0+\mu ^20\left|T\left[\phi (x)\phi (y)\right]\right|0\hfill \\ & +\hfill & g_00\left|T\left[\phi ^3(x)\phi (y)\right]\right|0=i\mathrm{}\delta ^4(xy)\hfill \\ _t0\left|T\left[\phi (x)\pi (y)\right]\right|0\hfill & \hfill & 0\left|T\left[\pi (x)\pi (y)\right]\right|0=+i\mathrm{}\delta ^4(xy).\hfill \end{array}$$
Let us provide the definition of the generating functional, extended to the phase-space
$$Z(J,W)=0\left|T\left(e^{i{\scriptscriptstyle (J\phi +W\pi )}}\right)\right|0$$
and its Fourier Transform
$$Z(J,W)=\delta \phi \delta \pi \stackrel{~}{Z}(\phi ,\pi )e^{i{\scriptscriptstyle \left[J(x)\phi (x)+W(x)\pi (x)\right]𝑑x^4}}$$
(7)
where $`W`$ is the current source associated to the conjugate momentum.
Thus, the equations of motion become functional ones in $`Z(J,W)`$
$$_t\frac{1}{i}\frac{\delta Z}{\delta W(x)}+\left[_j^2+\mu ^2\right]\frac{1}{i}\frac{\delta Z}{\delta J(x)}+g_0\frac{1}{i^3}\frac{\delta ^3Z}{\delta J(x)^3}=\mathrm{}Z(\phi )J(x)$$
(8)
$$_t\frac{1}{i}\frac{\delta Z}{\delta J(x)}+\frac{1}{i}\frac{\delta Z}{\delta W(x)}+=\mathrm{}Z(\phi )W(x).$$
(9)
Hence, the first of these equations can be rewritten in terms of $`\stackrel{~}{Z}(J,W)`$ and takes the form
$`{\displaystyle \delta \phi \delta \pi \stackrel{~}{Z}(\phi ,\pi )e^{i{\scriptscriptstyle (J\phi +W\pi )}}}`$ $`\left[_t\pi (x)\left(_j+\mu ^2\right)\phi (x)+g_0\phi ^3(x)\right]`$
$`={\displaystyle \frac{\mathrm{}}{i}}{\displaystyle \delta \phi \delta \pi \left(\frac{\delta \stackrel{~}{Z}(\phi ,\pi )}{\delta \phi \left(x\right)}\right)e^{i{\scriptscriptstyle (J\phi +W\pi )}}}.`$
At the end, with simple algebra, we get
$$\frac{\delta }{\delta \phi \left(x\right)}\left[ln\stackrel{~}{Z}(\phi ,\pi )\frac{i}{\mathrm{}}S(\phi ,\pi )\right]=0,$$
which admits the solution
$$\stackrel{~}{Z}(\phi ,\pi )=C(\pi )e^{\frac{i}{\mathrm{}}S(\phi ,\pi )};$$
(10)
here, $`C(\pi )`$ denotes an integration $`\phi `$-constant functional. In analogy from (9), we get also
$$\stackrel{~}{Z}(\phi ,\pi )=D(\phi )e^{\frac{i}{\mathrm{}}S(\phi ,\pi )}$$
(11)
being $`D(\phi )`$ an integration $`\pi `$-constant functional. Compatibility of (10), (11) requires that the following expression for $`Z(J,W)`$ holds
$$Z(J,W)=N\delta \pi \delta \phi e^{\frac{i}{\mathrm{}}S(\phi ,\pi )+i{\scriptscriptstyle (J\phi +W\pi )}}$$
where $`N`$ is pure normalization constant. The above is just the procedure which will allow us to get.
We conclude this section by remarking that, in the scalar field case, the Hamiltonian formulation of the generating functional is equivalent to the Lagrangian one, as soon as we address the identifications $`\pi =_t\phi `$ and $`W0`$.
## 3 Canonical quantization of gravity in the frozen formalism.
Before we begin to search for the gravitational generating functional, we resume the necessary steps allowing us to arrive to the canonical quantum equation for the gravitational field, that is the Wheeler-DeWitt equation , , ;
We start to consider a pseudo-Riemannian 4-dimensional manifold $`M^4`$, on which a metric tensor $`h_{\mu \nu }(y^\rho )`$ is defined. The signature we choose is $`(+++)`$. (The Greek indexes stand for 0,1,2,3, as the Latin ones stand for 1,2,3).
We apply a (3+1)-slicing assigning the parametric equation of a space-like hypersurfaces family $`_t^3:y^\rho =y^\rho (t,x^k)`$ on the considered manifold.
This definition provides a local reference basis $`\{e_i^\mu ,n^\mu \}`$ in $`M^4`$, made up by the tangent vectors $`e_i^\mu _iy^\mu `$ and the normal vector $`n^\mu (y^\rho )`$ to $`_t^3`$. Projecting the *deformation vector* $`N^\mu `$ on this basis we get the *lapse function* $`N(t,x^k)`$ and the *shift vector* $`N^i(t,x^k)`$ :
$$N^\mu _ty^\mu =Nn^\mu +N^ie_i^\mu .$$
(12)
In terms of this quantity we find the metric tensor expression in the $`(t,x^k)`$-reference system
$$h_{\mu \nu }=\left(\begin{array}{cc}N_iN^iN^2& N_i\\ N_i& g_{ij}\end{array}\right).$$
(13)
The term $`g_{ij}h_{\mu \nu }e_i^\mu e_j^\nu `$ is the so-called induced metric tensor, i.e. the projection of the metric tensor on the space-like hypersurface, which becomes our dynamical variable.
Now the splitted Hilbert-Einstein action (the A.D.M. action) looks like as the following one , , , , , , :
$$S_{HE}=\underset{\mathrm{\Sigma }^3\times \mathrm{}}{}dtd^3xN\sqrt{g}\left(R+k_{ij}k^{ij}k^2\right),$$
(14)
where $`k_{ij}`$ denotes the extrinsic curvature of the tre-hypersurfaces $`\mathrm{\Sigma }^3`$, while $`R`$ is the curvature scalar associated to the tre-metric.
From the Lagrangian contained in the previous expression, we find the classical momenta associated to $`g_{ij}`$ and to $`N`$ and $`N^i`$, in view of rewriting the action in the Hamiltonian formalism. After the Legendre transformation, the action rewrite as
$$S_{HE}=_{\mathrm{\Sigma }^3\times \mathrm{}}\left\{\pi ^{ij}_tg_{ij}NN^i_i\right\}d^3x𝑑t$$
(15)
where $``$ and $`_i`$ are called super-Hamiltonian and super-momentum respectively. The momenta $`P_N`$ and $`P_{N^i}`$, associated to the lapse function and to the shift vector, vanish identically, as
$$P_N=\frac{\delta L}{\delta \left(_tN\right)}=0,P_{N^i}=\frac{\delta L}{\delta (_tN^i)}=0$$
(16)
and these equations constitute the primary constraints for the gravity dynamics. The secondary constraints follow directly from the first ones, via the Hamilton equations (see also , , ):
$`(g_{ij},\pi ^{ij})`$ $`=`$ $`G_{ijkl}\pi ^{ij}\pi ^{kl}\sqrt{h}{}_{}{}^{3}R=0,`$ (17)
$`_i(g_{ij},\pi ^{ij})`$ $`=`$ $`2{}_{}{}^{3}_{j}^{}\pi _i^j=0,`$ (18)
where $`G_{ijkl}`$ is the super-metric and stands as
$$G_{ijkl}\frac{1}{2\sqrt{h}}(h_{ik}h_{jl}+h_{il}h_{jk}h_{ij}h_k).$$
The canonical quantization procedure is obtained by translating secondary constraints into operator ones, on the wave functional $`\psi `$, which describes the physical states, i.e.
$$\widehat{}(\widehat{g}_{ij},\widehat{\pi }^{ij})\psi =0,$$
$$\widehat{}_i(\widehat{g}_{ij},\widehat{\pi }^{ij})\psi =0,$$
where we introduced the canonical quantum operators $`\widehat{g}_{ij},\widehat{\pi }^{ij}`$, for which the commutation relations
$`[\widehat{g}_{ab}(\stackrel{}{x},t),\widehat{g}_{cd}(\stackrel{}{x}^{},t)]`$ $`=`$ $`0,`$
$`[\widehat{\pi }^{ab}(\stackrel{}{x},t),\widehat{\pi }^{cd}(\stackrel{}{x}^{},t)]`$ $`=`$ $`0,`$
$`[\widehat{g}_{ab}(\stackrel{}{x},t),\widehat{\pi }^{cd}(\stackrel{}{x}^{},t)]`$ $`=`$ $`i\mathrm{}\delta _{ab}^{cd}\delta (\stackrel{}{x}\stackrel{}{x}^{})`$ (19)
have to be valid; this last request implies the following representation for the tre-metric and momenta:
$$\widehat{g}_{ij}(x)\psi =g_{ij}(x)\psi ,$$
$$\widehat{\pi }^{kl}(x)\psi =i\mathrm{}\frac{\delta \psi }{\delta g_{kl}(x)}.$$
At the end, choosing the appropriate normal ordering , we can rewrite the operator constraints as
$$\widehat{}(x)\psi =[\mathrm{}^2:G_{ijkl}(x)\frac{\delta ^2\psi }{\delta g_{ij}(x)\delta g_{kl}(x)}:\sqrt{g}R]\psi =0,$$
$$\widehat{}_i(x)\psi =:2i\mathrm{}g_{ik}_j\frac{\delta \psi }{\delta g_{jk}(x)}:=0.$$
The first of the above functional equations is known as the Wheeler-DeWitt one and provides the quantum dynamics of the system, while the tree super-momentum equations state the theory invariance under tre-diffeomorphisms. This last invariance provides the wave functional depending only on tre-geometries $`\{g_{ij}\}`$, i.e. $`\psi =\psi (\{g_{ab}\})`$ , , .
As we can see, the Wheeler-DeWitt equation is a Schörodinger equation for a non-dynamical system, characterized by a wave functional satisfying $`_t\psi =0`$; In detail, the $`N`$-independence of $`\psi `$ is responsible for the time non-evolution of the wave functional.
The point of view addressed in this paper is that such stationary evolution for the gravitational field is caused by a too straightforward quantum implementation of the theory constraints, in the canonical form. In particular, we concentrate our attention on the quantum formulation of the primary constraints; in the canonical Wheeler-DeWitt approach, the translation of these constraints, $`P_N=0`$ and $`P_{N^i}=0`$, leads to the independence of the wave functional on the lapse function and on the shift vector, i.e.
$`\widehat{P}_N\psi =i\mathrm{}{\displaystyle \frac{\delta }{\delta N}}\psi =0,`$
$`\widehat{P}_{N^i}\psi =i\mathrm{}{\displaystyle \frac{\delta }{\delta N^i}}\psi =0,`$
an issue consistent with the non-evolutionary character of the WDE along the slicing and of its invariance under the tre-diffeomorphisms.
In the next section, in the spirit of the Gupta-Bleuler approach , we will propose a relaxation of these constraints, allowing an evolution for the wave functional $`\psi (g_{ij},t)`$, which depends so on time and on the tre-metric field.
## 4 The gravitational theory in the Misner-like variables.
We start by introducing the Misner-like variables $`(\alpha (x),u(x))`$ , via the relations
$`g_{ij}(x)e^{\alpha (x)}\left(e^{u(x)}\right)_{ij}`$ (20)
$`g^{ij}(x)e^{\alpha (x)}\left(e^{u(x)}\right)^{ij},`$ (21)
where $`u(x)`$ is a traceless matrix and, in order to safe the tensorial language, we assign a different character to the quantity $`e^{u(x)}`$ and $`u(x)`$ itself, i.e. covariant and mist one respectively.
The aim of the below analysis is to rewrite the gravitational Hamiltonian $`H(g_{ij},\pi ^{ij})`$ in the Misner-like variables, so then we can get the equations of motion in $`(\alpha (x),u(x))`$ and use them to find the generating functional for gravity, like shown in the previous section.
The relation between the old momenta $`\pi ^{ij}`$, and the new ones is provided by requiring the canonical nature of the transformation, i.e.
$$\pi ^{ij}_tg_{ij}=p_\alpha _t\alpha +p_i^k_tu_k^i.$$
(22)
Substituting the relation (21) into (22) we can identify the new momenta as
$$\{\begin{array}{cc}p_\alpha =𝚷\pi ^{ij}g_{ij}\hfill & \\ p_i^k=\pi _i^k\frac{1}{3}𝚷\delta _i^k,\hfill & \end{array}$$
where $`𝚷`$ denotes the quantity $`\pi ^{ij}g_{ij}`$. Starting from the expression of the super-Hamiltonian $``$ and super-momentum $`_i`$, we rewrite them in terms of the new variables as
$$\{\begin{array}{cc}=\hfill & \frac{1}{2}e^{(3/2)\alpha }\left[\mathrm{\hspace{0.17em}2}p_i^jp_j^i\frac{1}{3}p_\alpha ^2\right]e^{(3/2)\alpha }R(u,\alpha )\hfill \\ _i=\hfill & 2_jp_i^j\frac{2}{3}_ip_\alpha +p_\alpha _i\alpha +2U_{ni}^jp_j^n,\hfill \end{array}$$
where we denote with $`U_{dq}^n(u,u)`$ the Christoffel terms constructed only by derivatives of $`u`$-variable,
$$\{\begin{array}{ccc}U_{dq}^n(u,u)\hfill & \hfill & \frac{1}{4}[+_du_q^n+_qu_d^n+_qu_l^m(e^u)^{nl}(e^u)_{md}+(_du_l^m)(e^u)^{nl}(e^u)_{mq}\hfill \\ & & (_lu_d^m)(e^u)^{nl}(e^u)_{mq}(_lu_q^m)(e^u)^{nl}(e^u)_{md}],\hfill \end{array}$$
and the full Christoffel symbol rewrites as
$$\mathrm{\Gamma }_{dq}^nU_{dq}^n+\frac{1}{2}\left[(_d\alpha )\delta _q^n+(_q\alpha )\delta _d^n(_m\alpha )(e^u)^{nm}(e^u)_{dq}\right].$$
Starting from the expression of $``$ and $`_i`$, as written in the new variables, we can get the Hamilton equations for the gravitational system, in the Misner-like representation, which stand as
$`\dot{u}_a^b`$ $`=`$ $`2Ne^{(3/2)\alpha }p_a^b2\left[U_{bn}^aN^n+_aN^b\right]`$
$`\dot{p}_a^b`$ $`=`$ $`\mathrm{\Theta }_a^b+2\mathrm{\Omega }_{jna}^{ib}N^np_i^j`$
$`\dot{\alpha }`$ $`=`$ $`+{\displaystyle \frac{1}{3}}Ne^{(3/2)\alpha }p_\alpha \left[(_l\alpha )N^l+{\displaystyle \frac{2}{3}}_jN^j\right]`$
$`\dot{p_\alpha }`$ $`=`$ $`Ne^{(3/2)\alpha }\left\{{\displaystyle \frac{3}{4}}e^{3\alpha }\left[+2p_i^jp_j^i{\displaystyle \frac{1}{3}}p_\alpha ^2\right]+{\displaystyle \frac{3}{2}}R(u,\alpha )\right\}\mathrm{\Theta }`$
$`_l\left(N^lp_\alpha \right)`$
being
$$d^3xNe^{(3/2)\alpha }\left(\frac{\delta R(u,\alpha )}{\delta \alpha }\right)\delta \alpha d^3x\mathrm{\Theta }\delta \alpha $$
$$d^3xNe^{(3/2)\alpha }\left(\frac{\delta R(u,\alpha )}{\delta u_b^a}\right)\delta u_b^ad^3x\mathrm{\Theta }_a^b\delta u_b^a$$
and
$`{\displaystyle d^3x\left(\frac{\delta U_{jn}^i}{\delta u_b^a}\right)N^n\delta u_b^a}{\displaystyle d^3x\mathrm{\Omega }_{jna}^{ib}\delta u_b^a}.`$
The introduction of these new variables, together with the Hamiltonian approach, is necessary to overcome some difficulties in applying the method, shown in section 2, to the gravitational case; such difficulties can be summarized by the following point:
* the presence of the inverse metric tensor in the equations of motion, which we can not easily rewrite as a function of the direct metric $`g_{ij}`$; now, both the inverse and direct metric are function of the dynamical variables, $`(\alpha ,u)`$, and so no difficulty survives.
* the impossibility to take all the temporal derivatives out of the field T-products, because of their non linearity; using an Hamiltonian approach, quadratic terms containing time derivatives of the tre-metric tensor are removed, because we rewrite them via the conjugate momenta. In this sense, the phase-space representation overcomes the gravity non-linearity problem in building up the generating functional (the non-linearity in the spatial derivatives plays no role here).
## 5 Gravitational generating Functional.
Using the Misner-like variables and the Hamiltonian formulation, the Hilbert-Einstein action looks like
$`S={\displaystyle _{M^4}}d^3x𝑑t`$ $`[`$ $`p_n^j_tu_j^n+p_\alpha _t\alpha +P_N_tN+P_{N^i}_tN^i`$
$`\xi P_N\xi ^iP_{N^i}NN^i_i],`$
where the Lagrange multipliers $`\xi (x),\xi ^i(x)`$ are included in order to obtain the primary constraints ($`P_NP_{N^i}=0`$) as Hamilton equations, and being the integration domain $`M^4=\mathrm{\Sigma }^3\times `$ ($`\mathrm{\Sigma }^3`$ is taken compact and boundaryless).
The whole set of Hamilton equations is, therefore,
$$\begin{array}{cccc}\text{I)}\hfill & \hfill +\dot{u}_j^n\frac{\delta H}{\delta p_n^j}& =& 0\hfill \\ \text{II)}\hfill & \hfill \dot{p}_j^n\frac{\delta H}{\delta u_n^j}& =& 0\hfill \\ \text{III)}\hfill & \hfill +\dot{\alpha }\frac{\delta H}{\delta p_\alpha }& =& 0\hfill \\ \text{IV)}\hfill & \hfill \dot{p}_\alpha \frac{\delta H}{\delta \alpha }& =& 0\hfill \\ \text{V)}\hfill & \hfill \dot{P}_N& =& 0\hfill \\ \text{VI)}\hfill & \hfill \dot{P}_{N^i}_i& =& 0\hfill \\ \text{VII)}\hfill & \hfill +\dot{N}\xi & =& 0\hfill \\ \text{VIII)}\hfill & \hfill +\dot{N^i}\xi ^i& =& 0\hfill \\ \text{IX)}\hfill & \hfill P_N& =& 0\hfill \\ \text{X)}\hfill & \hfill P_{N^i}& =& 0.\hfill \end{array}$$
In order to extend the analysis of section 2 to the gravitational case, we assume that the following statements are valid for the quantum theory:
* We can define a ground (vacuum) state $`|0`$, for the gravitational theory, by which we can built the T-product of the operator field and of its momentum. For the same reason, we need the theory is provided with a definite time variables, which has to have the meaning of physical time; we identify it with the label time.
* By a second quantization procedure $`g_{ij}(x)`$ and $`\pi ^{ij}(x)`$ become operators which satisfy the canonical commutation relations (19). This hypothesis has to hold for any representation of the tre-metric and of the conjugate momentum, like the considered Misner-like one.
Now, we introduce the generating functional $`Z`$ , ,
$`Z(J^a,W_a)0\left|T\left(e^{i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}}\right)\right|0,`$ (23)
where we adopt a notation of the form
$$Z(J,W)=Z(J_i^j,J_\alpha ,J_N,J_{N^i},W_j^i,W_\alpha ,W_N,W_{N^i}),$$
$`{\displaystyle \left(g_aJ^a+P^aW_a\right)}`$ $`=`$ $`{\displaystyle }[u_j^n(x)J_n^j(x)+p_n^j(x)W_j^n(x)+\alpha (x)J_\alpha (x)`$
$`+p_\alpha (x)W_\alpha (x)+N(x)J_N(x)+N^i(x)J_{N^i}(x)`$
$`+P_N(x)W_N(x)+P_{N^i}(x)W_{N^i}(x)]d^3xdt.`$
As in the scalar case, even here $`J^a`$ and $`W_a`$ have to be interpreted as quantum source currents, because we will see they give a $`\mathrm{}`$ contribution in the action.
Now, we integrate the equations I)-VIII) to determine the generating functional, in the gravitational theory, via an algorithm which we discuss here, in detail, only for equation I).
From the first equation, we build the corresponding T-products dynamics
$$_t0\left|T\left(u_j^n\left(x\right)p_n^j\left(y\right)\right)\right|00\left|T\left(\frac{\delta H}{\delta p_n^j\left(x\right)}p_n^j\left(y\right)\right)\right|0=i\mathrm{}\delta ^4\left(xy\right),$$
which can be rewritten formally in terms of $`Z`$, as
$$_t\frac{1}{i}\frac{\delta Z}{\delta J_n^j(x)}H(i\frac{\delta }{\delta J^a},i\frac{\delta }{\delta W_a})Z=\mathrm{}W_j^n(x)Z.$$
(24)
Substituting, in the previous equation, the following Fourier Transform
$$Z(J,W)=D(g_a,P^a)\stackrel{~}{Z}(g_a,P^a)e^{i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}},$$
(25)
where the Lebesgue measure $`D(g_a,P^a)`$ is
$$D(g_a,P^a)=\delta u_j^n\delta \alpha \delta N\delta N^i\delta p_n^j\delta P_\alpha \delta P_N\delta P_{N^i},$$
we get
$$D(g_a,P^a)e^{i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}}\stackrel{~}{Z}\left(\dot{u}_j^n\frac{\delta H}{\delta p_n^j}\right)=D(g_a,P^a)e^{i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}}\frac{\mathrm{}}{i}\frac{\delta \stackrel{~}{Z}}{\delta p_n^j}.$$
This functional equation admits the solution
$$\stackrel{~}{Z}(g_a,P^a)=C\left(p_n^j\right)e^{\frac{i}{\mathrm{}}S}.$$
$`C\left(p_n^j\right)`$ denotes a functional depending on all the variables (including $`\xi `$ and $`\xi ^i`$) but $`p_n^j`$.
Inferring the above procedure for all the other equations, but the last two, we arrive to a final expression for $`Z`$, i.e.
$$Z(J^a,W_a)=D(g_a,P^a)\delta \xi \delta \xi ^i𝒞(\xi ,\xi ^i)e^{\frac{i}{\mathrm{}}S+i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}},$$
being $``$ a normalization constant and $`𝒞(\xi ,\xi ^i)`$ an arbitrary functional on the Lagrangian multipliers, allowed by the theory; we remark that the integration over such multipliers is required by the $`Z`$ independence on them (the T-products calculated by $`Z`$ are observable of the theory and, therefore, can not depend on Lagrangian multipliers), and it is possible in view of the linearity, characterizing the whole system of equation I)-VIII).
The implementation of the primary constraints IX)-X), via the above algorithm, would imply the $`Z`$ independence on $`W_N`$ and $`W_{N^i}`$, which would lead again to the Wheeler-DeWitt dynamics; however, our aim is to investigate the quantum dynamics resulting from different implementation of this primary constraints. As we shall see in the next section, the role played here by equations IX)-X), in the quantum sector, is resumed by the form of the functional $`𝒞(\xi ,\xi ^i)`$. (In fact, the existence of this arbitrariness in the theory is a direct consequence of disregarding the primary constraints).
## 6 The Gupta-Bleuler approach.
We have ended the previous section, saying the form we choose for the functional $`𝒞(\xi ,\xi ^i)`$ reflects the way by which we perform the implementation of the primary constraints in the quantum theory.
We start by observing how the frozen formalism requirement, that the operators $`\widehat{P}_N`$ and $`\widehat{P}_{N^i}`$ annihilate the Schörodinger wave functional, is equivalent, in the Heisenberg approach, to have vanishing T-products of these operators;
indeed, the vanishing behavior of T-products, of any number of points, implies that the generating functional does not depend on the currents $`W_N`$ and $`W_{N^i}`$, i.e.
$$i\frac{\delta Z}{\delta W_N}=0,$$
$$i\frac{\delta Z}{\delta W_{N^i}}=0.$$
Below, we show that the quantum implementation of the primary constraints is summarized by the functional form taken by $`𝒞(\xi ,\xi ^i)`$, and, therefore, this quantity is the fundamental degree of freedom, in our approach.
We start by observing that the resulting generating functional
$`Z(J,W)=`$ $``$ $`{\displaystyle D(g_a,P^a)\delta \xi \delta \xi ^i𝒞(\xi ,\xi ^i)}`$
$`e^{\frac{i}{\mathrm{}}S+{\scriptscriptstyle \left[u_j^nJ_n^j+p_n^jW_j^n+\alpha J_\alpha +p_\alpha W_\alpha +NJ_N+N^iJ_{N^i}\right]d^3x𝑑t}},`$
contains integrals of the form
$$\delta \xi (\stackrel{}{x},t)\delta \xi ^i(\stackrel{}{x},t)𝒦(\xi )𝒦_i(\xi ^i)e^{\frac{i}{\mathrm{}}{\scriptscriptstyle d^3x𝑑t\xi (\stackrel{}{x},t)P_N(\stackrel{}{x},t)}}e^{\frac{i}{\mathrm{}}{\scriptscriptstyle d^3x𝑑t\xi ^i(\stackrel{}{x},t)P_{N^i}(\stackrel{}{x},t)}},$$
where we recasted the functional dependence on $`\xi `$ and on $`\xi ^i`$, by introducing the functionals $`𝒦(\xi )`$ and $`𝒦_i(\xi ^i)`$, in place of $`𝒞(\xi ,\xi ^i)`$. Then, we consider the following two different cases:
* As soon as we require the constance of the functionals $`𝒦`$ and $`𝒦_𝒾`$, the above integral become a delta-functional on the variable $`P_N`$ and $`P_{N^i}`$, apart from a normalization constant. Thus we get for $`Z`$
$$Z(J,W)=D(g_a,P^a)\delta \left(P_N\right)\delta (P_{N^i})e^{\frac{i}{\mathrm{}}S_0+i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}}.$$
(26)
being
$$S_0=d^3x𝑑t\left[p_i^j_tu_j^i+p_\alpha _t\alpha +P_N_tN+P_{N^i}_tN^iNN^i_i\right].$$
Evaluating the delta-functionals, we obtain $`Z`$, calculated for the usual gravitational action $`S_{HE}`$, (15) ,corresponding to vanishing momenta $`P_N`$ and $`P_{N^i}`$, and to a stationary evolution for the field, in the Schrödinger formulation, i.e.
$$Z(J,W)=D(g_a,P^a)e^{\frac{i}{\mathrm{}}S_{HE}+i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}}.$$
Because of the constance of $`𝒦(\xi )`$ as a consequence of the delta term, we also get the following result for the T-product function; below we treat separately the quantities $`\xi `$ and $`\xi ^i`$ since they enter in an equivalent way in the problem:
$`0\left|T\left(\widehat{P}_N(x_1)\mathrm{}\widehat{P}_N(x_n)\right)\right|0={\displaystyle \frac{1}{i^n}}{\displaystyle \frac{\delta ^nZ(J,W)}{\delta W_N(x_n)\mathrm{}\delta W_N(x_1)}}|_{J=0}`$
$`={\displaystyle D(g_a,P^a)\delta \left(P_N\right)\delta \left(P_{N^i}\right)P_N(x_1)\mathrm{}P_N(x_n)e^{\frac{i}{\mathrm{}}S_0}}=0.`$
Obviously, the same arguments are valid for T-product functions containing $`P_{N^i}`$, if we take $`𝒦_i(\xi ^i)`$ as a constant.
We can conclude, that having a constant $`𝒞`$-term is equivalent to implement the primary constraints in their strong form and, therefore, to restate the WDE dynamics , , .
* Now, we will see that a different choice for $`𝒦(\xi )`$ and $`𝒦_i(\xi ^i)`$, in detail a Gaussian form, is equivalent to interpret the primary constraints in a relaxed way.
We start again from the expression (26) and we perform a constraint weakening by applying the standard procedure of the generating functional approach .
More precisely, in order to weak the constraint associated to the momentum $`P_N`$, we turn the $`P_N`$-delta, centered on the null function, into a $`\lambda `$-centered one. Furthermore, we introduce a Gaussian weight in the $`\lambda `$-parameter function, the corresponding normalization associated to this term being included into the new constant $`^{}`$. We show this method only for $`P_N`$, but we can extend it to all the other constraints associated to $`P_{N^i}`$. This scheme yields the following generating functional:
$`Z(J,W)`$ $`=`$ $`^{}{\displaystyle D(g_a,P^a)\delta \left(P_N\lambda \right)\delta \left(P_{N^i}\right)e^{\frac{i}{\mathrm{}}S_0+i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}}}`$
$`\left({\displaystyle \delta \lambda e^{\frac{1}{2\sigma }{\scriptscriptstyle d^4x\lambda ^2(x)}}}\right).`$
Rewriting delta-functionals as integrals, we get
$`Z(J,W)`$ $`=`$ $`^{}{\displaystyle D(g_a,P^a)\delta \xi \delta \xi ^ie^{\frac{i}{\mathrm{}}{\scriptscriptstyle d^4x\xi (P_N\lambda )}}e^{\frac{i}{\mathrm{}}{\scriptscriptstyle d^4x\xi ^iP_{N^i}}}}`$
$`e^{\frac{i}{\mathrm{}}S_0+i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}}\left({\displaystyle \delta \lambda e^{\frac{1}{2\sigma }{\scriptscriptstyle d^4x\lambda ^2(x)}}}\right).`$
Hence, we can identify the $`𝒦(\xi )`$ expression as:
$$𝒦(\xi )=\delta \lambda e^{{\scriptscriptstyle d^4x{\scriptscriptstyle \frac{i}{\mathrm{}}}\xi \lambda }}e^{\frac{1}{2\sigma }\lambda ^2};$$
by simple algebra, making the following substitution
$$\overline{\lambda }\left(\frac{1}{\sqrt{2\sigma }}\lambda +\xi \frac{i}{\mathrm{}}\sqrt{\frac{\sigma }{2}}\right)$$
and performing the $`\overline{\lambda }`$-integral, which result to be independent on $`\xi `$, finally we get
$$𝒦(\xi )=𝒜e^{{\scriptscriptstyle d^4x\left({\scriptscriptstyle \frac{\sigma }{2\mathrm{}^2}}\right)\xi ^2}},$$
being $`𝒜`$ a normalization constant.
From this formula, we can infer how, having a $`P_N`$-delta not centered in zero, so that $`P_N`$ is not vanishing, leads to a Gaussian form for the functional $`𝒦(\xi )`$. In the limit in which the dispersion $`\sqrt{\sigma }`$ verge to zero, we obtain again the constance of $`𝒦(\xi )`$ and so the classic constraints. We can get the same result for $`𝒦_i(\xi ^i)`$.
As a consequence of this relaxation now we should get non-vanishing T-products of operators $`\widehat{P}_N`$ and $`\widehat{P}_{N^i}`$; we can show this, starting from a generating functional obtained by assigning to $`𝒦(\xi )`$ and $`𝒦_i(\xi ^i)`$ a Gaussian form, as above, i.e.
$$\{\begin{array}{ccc}𝒦(\xi )\hfill & \hfill & e^{\frac{1}{2A}{\scriptscriptstyle d^4x\xi ^2(x)}}\hfill \\ 𝒦_i(\xi ^i)\hfill & \hfill & e^{\left(\frac{1}{2A}^3\right)\mathrm{\Sigma }_i{\scriptscriptstyle d^4x(\xi ^i(x))^2}},\hfill \end{array}$$
where $`A\frac{\mathrm{}^2}{\sigma }`$.
Dealing again with $`\xi `$-component only, we have
$$Z(J,W)=D(g_a,P^a)\delta \xi \delta \xi ^i𝒦_i(\xi ^i)\left(e^{\frac{1}{2A}{\scriptscriptstyle d^3x𝑑t\xi ^2(x)}}\right)e^{\frac{i}{\mathrm{}}S+i{\scriptscriptstyle \left(g_aJ^a+P^aW_a\right)}},$$
and, through this generating functional, we can write the T-products as
$`T_{1\mathrm{}n}`$ $``$ $`0\left|T\left(\widehat{P}_N(x_1)\mathrm{}\widehat{P}_N(x_n)\right)\right|0`$
$`=`$ $`^{}{\displaystyle D(g_a,P^a)\delta \xi \delta \xi ^i\left[P_N(x_1)\mathrm{}P_N(x_n)\right]e^{\frac{1}{2A}\left({\scriptscriptstyle d^3x𝑑t\xi ^2(x)}\right)}e^{\frac{i}{\mathrm{}}S}}.`$
Because of $`\xi `$-term in the action, we can express every $`P_N`$ as a functional derivative in $`\xi `$, acting on the action exponential, so that the functions $`T_{1\mathrm{}n}`$ become
$$T_{1\mathrm{}n}=\frac{^{}}{A^n}D(g_a,P^a)\delta \xi \delta \xi ^i\left[\xi (x_1)\mathrm{}\xi (x_n)\right]e^{\frac{1}{2A}\left({\scriptscriptstyle d^3x𝑑t\xi ^2(x)}\right)}e^{\frac{i}{\mathrm{}}S};$$
in this integral, we point out the $`\xi `$-term, i.e.
$$E(P_N)\delta \xi e^{\frac{i}{\mathrm{}}\left({\scriptscriptstyle d^3x𝑑t\xi (x)P_N(x)}\right)}e^{\frac{1}{2A}\left({\scriptscriptstyle d^4x\xi ^2(x)}\right)}\left[\xi (x_1)\mathrm{}\xi (x_n)\right].$$
If we turn it into
$$E(P_N)=e^{\frac{A}{2\mathrm{}^2}P_N^2}\delta \xi \left[\xi (x_1)\mathrm{}\xi (x_n)\right]e^{\frac{1}{2A}{\scriptscriptstyle d^3x𝑑t\left(\xi +{\scriptscriptstyle \frac{iA}{\mathrm{}}}P_N\right)^2}},$$
we can see that it is not vanishing, being a Gaussian integral not null-centered, though around an imaginary value.
As a consequence, the whole expression of the T-products $`T_{1\mathrm{}n}`$ no longer vanish and this is the expected quantum relaxation of the primary constraints.
## 7 Evolutionary quantum gravity.
The next step, in our work, is to show the dynamical consequences, in the Schrödinger approach, of the generating functional resulting from the Gaussian choice for $`𝒞(\xi ,\xi ^i)`$.
In this representation, having non-vanishing T-products $`T_{1\mathrm{}n}`$ of operators $`\widehat{P}_N`$ and $`\widehat{P}_{N^i}`$, corresponds to require the following condition on $`\psi `$:
$$\widehat{P}_N\psi 0,\widehat{P}_{N^i}\psi 0.$$
As in the canonical formulation, we have $`\widehat{P}_N=i\mathrm{}\frac{\delta }{\delta N}`$ and $`\widehat{P}_{N^i}=i\mathrm{}\frac{\delta }{\delta N^i}`$, so that now the wave functional $`\psi `$ has to depend on the the slicing structure, that is on $`N`$ and $`N^i`$. Moreover, via the Hamilton equations (extended to the operator form), the violation of the secondary constraints follows directly, i.e.
$$\widehat{}\psi 0,\widehat{}_i\psi 0,$$
as well as the explicit breakdown of the quantum theory invariance, under 4-diffeomorphisms.
To continue our analysis, we point out that:
* The dependence of $`\psi `$ on the variables $`N`$ and $`N^i`$ can be seen as a dependence on the label time $`t`$. Reminding the definition of the deformation vector (12), we can refer to the lapse function and shift vector dependence, as to that one on $`y^\mu `$, the field which fixes the splitting; thus, we can write the temporal derivative as
$$_t\psi d^3x\frac{\delta \psi }{\delta y^\mu }_ty^\mu .$$
As we have seen, because of the relaxation, we get
$$\frac{\delta \psi }{\delta y^\mu }0_t\psi 0$$
and so the $`\psi `$ temporal dependence outcomes.
* In the Schrödinger representation, we can express the wave functional $`\psi `$ as a path integral, which is also a transition amplitude between two states, i.e. the propagator , .
(For an extension of path integral approach to the General Relativity see also ).
The metric configurations $`g_a,g_0`$ (corresponding to different times, $`t`$ and $`t_0`$, and to associated hypersurfaces $`\mathrm{\Sigma }^3`$ and $`\mathrm{\Sigma }_0^3`$) characterize the final and the initial state for the transition.
In agreement to the analysis of the previous section, we assign the following path integral, which correspnds to the adopted scheme of constraints relaxation:
$$g_b,t|g_a,t_0=D(g_a,P^a)\delta \xi \delta \xi ^i𝒞(\xi ,\xi ^i)e^{\frac{i}{\mathrm{}}S},$$
(27)
where $`𝒞(\xi ,\xi ^i)`$ has to be thought Gaussian (indeed, the below analysis holds even for a non-Gaussian form of $`𝒞`$, but it is constant)
This integral does not contain boundary terms, because the hypersurfaces are compact and boundaryless.
We remind that $`g_a`$ stands for the metric configuration, in the time t, given by the variables $`(u,\alpha ,N,N^i)`$, as well as $`P^a`$ stands for $`(p,P_\alpha ,P_N,P_{N^i})`$.
Starting from (27), we rewrite the corresponding gravitational action as
$$S=d^3xP^aD(g_a)\left(H_0+H_{(\xi ,\xi ^i)}\right)𝑑t$$
Here, we have pointed out the term $`H_{\xi ,\xi ^i}`$, having the form
$$H_{(\xi ,\xi ^i)}d^3x\left(\xi P_N+\xi ^iP_{N^i}\right);$$
which is reabsorbed by the integration on $`\xi `$ e $`\xi ^i`$, as follows:
$$g_a,t|g_0,t_0=D(g_a,P^a)(P_N,P_{N^i})e^{\frac{i}{\mathrm{}}S_0};$$
$`(P_N,P_{N^i})`$ is a functional depending on the momenta $`P_N`$ and $`P_{N^i}`$, resulting from the above integration.
Now, we can obtain the dynamical equation, taking
$`i\mathrm{}_t\psi =i\mathrm{}_tg_a,t|g_0,t_0`$ $`=`$ $`{\displaystyle D(g_a,P^a)(P_N,P_{N^i})(_tS_0)e^{\frac{i}{\mathrm{}}S_0}}`$
$`=`$ $`{\displaystyle D(g_a,P^a)(P_N,P_{N^i})H(g_a,P^a)e^{\frac{i}{\mathrm{}}S_0}}`$
$`=`$ $`{\displaystyle D(g_a,P^a)(P_N,P_{N^i})}:\widehat{H}:e^{\frac{i}{\mathrm{}}S_0},`$
where explicitly $`\psi \psi (g_a,t)`$.
Above, in the last passage, we have changed the super-Hamiltonian and the super-momentum into operators, disregarding higher power in $`\mathrm{}`$, as follows:
$`(g_a,{\displaystyle \frac{\delta S_0}{\delta g_a}})e^{\frac{i}{\mathrm{}}S_0}`$ $`=`$ $`:\widehat{}(u,\alpha ,i\mathrm{}{\displaystyle \frac{\delta }{\delta u}},i\mathrm{}{\displaystyle \frac{\delta }{\delta \alpha }}):e^{\frac{i}{\mathrm{}}:S_0:},`$
$`_i(g_a,{\displaystyle \frac{\delta S_0}{\delta g_a}})e^{\frac{i}{\mathrm{}}S_0}`$ $`=`$ $`:\widehat{}_i(u,\alpha ,i\mathrm{}{\displaystyle \frac{\delta }{\delta u}},i\mathrm{}{\displaystyle \frac{\delta }{\delta \alpha }}):e^{\frac{i}{\mathrm{}}:S_0:},`$
where we used the Hamilton-Jacobi relation between variables and conjugate momenta $`P^a=\frac{\delta S_0}{\delta g_a}`$.
At the end, if we remark that the functional derivatives contained in $`:\widehat{H}:`$ commute with $``$ and with the integration measure, we get, in the limit of small $`\mathrm{}`$, the Schrödinger equation for the gravitational quantum field
$$i\mathrm{}_t\psi =:\widehat{H}:\psi .$$
(28)
In the derivation of this evolutionary approach to quantum gravity, we need not a specific form for $`𝒞`$, but requiring it is not a constant functional. Indeed, it is possible to show that having a non-constant $`𝒞`$ corresponds to a non-vanishing vacuum expectation value of the momenta $`\widehat{P}_N`$ and $`\widehat{P}_{N^i}`$, i.e. $`0\left|P_N\right|00`$ and $`0\left|P_{N^i}\right|00`$.
This issue indicates that the existence of an evolutionary quantum gravity does not require that the primary constraints are violated in a strong form (i.e. all the T-products for $`\widehat{P}_N`$ and $`\widehat{P}_{N^i}`$ vanish).
By other words, to get a Schrödinger dynamics, it is sufficient to break the primary constraints (and therefore the secondary ones too) only in terms of their expectation values (see also ).
## 8 Classical limit of the theory and the Eckart fluid.
In this section, we want to show how the dynamical equation (28), founded for the gravitational field through the relaxation, implies the presence of a particular fluid in the system, when taking the classical limit. Such a fluid results to have the Eckart fluid characteristics . (For other approaches, which correlates the presence of a fluid in the dynamics to the appearance of a time in quantum gravity, see , , , , , , , , , ).
Let us consider the following development for the wave functional $`\psi `$, over the super-Hamiltonian and super-momentum eigenfunctions
$$\psi (u_j^n,\alpha ,N,N^i,t)=D\omega Dk_i\chi (\omega ,k_i)e^{\frac{i}{\mathrm{}}(tt_0){\scriptscriptstyle d^3x\left[N\omega +N^ik_i\right]}};$$
(29)
$`\omega \omega (\stackrel{}{x})`$ and $`k_ik_i(\stackrel{}{x})`$ are the eigenvalues functions, and
$$\chi (\omega ,k_i)=\chi (u_b^a,\alpha ,\omega ,k_i)$$
is the corresponding eigenfunctional, i.e. we get
$`\widehat{}\chi (\omega ,k_i)`$ $`=`$ $`\omega \chi (\omega ,k_i),`$ (30)
$`\widehat{}_i\chi (\omega ,k_i)`$ $`=`$ $`k_i\chi (\omega ,k_i).`$ (31)
Now, we mean to take the classical limit of these eigenvalues problems, via a WKB approach, which provide us the following form for $`\chi (\omega ,k_i)`$, to be inserted into (30) and (31):
$$\chi (\omega ,k_i)=e^{\frac{i}{\mathrm{}}𝒫}$$
being $``$ and $`𝒫`$ the modulus and the phase of $`\chi `$, respectively.
In the limit $`\mathrm{}0`$, the substitution result yields the gravitational Hamilton-Jacobi equation and its super-momentum equivalent, containing two additional terms, which correspond to the classical limit of $`\omega `$ and $`k_i`$ i.e. $`\overline{\omega }`$ and $`\overline{k}_i`$ respectively:
$`(g_a,{\displaystyle \frac{\delta 𝒫}{\delta g_a}})=\overline{\omega },`$ (32)
$`_i(g_a,{\displaystyle \frac{\delta 𝒫}{\delta g_a}})=\overline{k}_i.`$ (33)
We recast this Hamilton-Jacobi system of equations into the Einstein one.
Using well-known results , , it is possible to show that the following relations hold:
1. $`G_{\mu \nu }n^\mu n^\nu =XT_{\mu \nu }n^\mu n^\nu =\frac{}{2\sqrt{g}}=\frac{\overline{\omega }}{2\sqrt{g}}`$
2. $`G_{\mu \nu }_iy^\mu n^\nu =XT_{\mu \nu }_iy^\mu n^\nu =\frac{_i}{2\sqrt{g}}=\frac{\overline{k}_i}{2\sqrt{g}}`$
3. $`G_{\mu \nu }_iy^\mu _jy^\nu =XT_{\mu \nu }_iy^\mu _jy^\nu =0`$
where $`X`$ denotes the Einstein constant.
Taking into account these three relations, we arrive to construct the energy-momentum tensor, associated to the classical limit of the eigenvalues problems (32) and (33),
$$T_{\mu \nu }=ϵn_\mu n_\nu +2s_{(\mu }n_{\nu )};$$
the identifications
$`ϵ`$ $`=`$ $`{\displaystyle \frac{\overline{\omega }}{2X\sqrt{g}}},`$
$`s_\mu `$ $`=`$ $`{\displaystyle \frac{\overline{k}^i}{2X\sqrt{g}}}_iy^\rho h_{\rho \mu },`$
take place.
In particular the above third relation assure us the tensor $`T_{\mu \nu }`$ does not contain any pure spatial component.
The so founded energy-momentum tensor $`T_{\mu \nu }`$ has the expression of that one associated to an Eckart fluid, as soon as we identify $`ϵ`$ as the fluid energy density, $`n^\mu `$ as its 4-velocity and $`s_\mu `$ as its thermal conduction vector .
This fluid is co-moving with the spatial hypersurface and contains a time-like vector corresponding to its 4-velocity, so that its presence gives a physical meaning to the label time $`t`$, used in the splitting procedure. We can infer that, in quantum gravity, the existence of a time and the presence of such a fluid in the gravitational dynamics are both aspects of the same physical entity.
## 9 Concluding remarks.
The main issue of our analysis has to be regarded the construction of a phase-space generating functional, for an evolutionary quantum gravity. We have shown how the quantum implementation of the primary constraints can be controlled through the form taken by a free functional of the associated Lagrangian multipliers. When this functional is a constant one, the implementation of the constraints takes place in its strong form (the operators $`\widehat{P}_N`$ and $`\widehat{P}_{N^i}`$ vanish identically), so that the standard Wheeler-DeWitt approach is reproduced. But, if we address Gaussian functional form on $`\xi `$ and $`\xi ^i`$, then, the resulting dynamics corresponds to the relaxation of the primary constraints, i.e. the T-products associated to $`\widehat{P}_N`$ and $`\widehat{P}_{N^i}`$ are different from zero, for any number of points.
The physical meaning of the time variable, we have introduced in taking the T-products (namely the label time), is recognized as soon as we extend the relaxation scheme of the primary constraints to the path-integral formulation of the Shrödinger picture.
In fact, in this picture, it is enough to require that the free functional, on $`\xi `$ and $`\xi ^i`$, has a non-constant form, to get an evolutionary quantum dynamics, in the limit of small $`\mathrm{}`$, as described by the Schrödinger equation. Having a non-constant functional ensures only that $`0\left|\widehat{P}_N\right|00`$ and $`0\left|\widehat{P}_{N^i}\right|00`$, but it does not imply the vanishing behavior of any T-product function. We can infer that, in the limit of small $`\mathrm{}`$, the appearance of a time in quantum gravity is associated to weaker relaxation of the primary constraints. Indeed, the physical meaning of the label time outcomes when taking the classical limit of the Schrödinger equation.
Summarizing, we started from the classical Hamiltonian dynamics and assumed the label time suitable for constructing T-products functions. The classical limit for the theory provided the appearance of an Eckart fluid within the Einstein equations, as required by the selfconsistence of the spacetime splitting and by the existence of a real time variable. Thus, we have shown that if we have a physical time in quantum level, it results into a reference fluid in the classical behavior.
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# Noncommutative Maxwell-Chern-Simons theory in three dimensions and its dual
## I Introduction
The study of non-commutative (NC) space-time and its implication in physics has a long history snyder . The renewal of interest is due to the recent developments in NC geometry cons and string theory sw . The Moyal product defined in NC space-time, which replaces the ordinary product in commutative space-time, introduces highly non-local and non-linear interaction terms which are not present in ordinary theories. Because of this, NC theories have many novel features which are not shared by their commutative counterparts rev . The UV/IR mixing, which is characteristic of NC theories, usually breaks down renormalizability min ; bala . Supersymmetry is essential to recover renormalizability riv and NC supersymmetric field theories alb as well as supersymmetric quantum mechanical models kh have been constructed and investigated. Fermionic field theories on NC manifolds have been studied and shown to be free of the fermion doubling problem bal . Due to the interesting new properties they share it is of utmost importance to investigate the NC generalizations of well established notions of commutative gauge theories.
Different descriptions of the same theory in commutative space-time have been useful in several branches of physics because they usually lead to the concept of duality duality . Many interesting studies aiming to extend the duality in commutative space-time to the NC setting have been performed gop ; trg ; sg ; botta ; cw ; davi . The generalization of the well known equivalence between the Maxwell-Chern-Simons (MCS) theory and the self-dual (SD) model mcssd to NC space-time has also been investigated. The master action technique, which was used to establish the equivalence between these models in commutative space-time, has been adopted in sg and botta and these authors have reached different conclusions regarding the equivalence in the NC setting. In botta , after eliminating some of the fields from the master action, the perturbative solution to the field equations were used and it was argued that the NCMCS theory is equivalent to the NCSD model when the Chern-Simons (CS) term has a cubic contribution like in the non-Abelian case. In sg , however, which also used the master action method, it was argued that the NCMCS theory constructed by applying the inverse Seiberg-Witten (SW) map sw , is equivalent to a theory where the cubic interaction of the vector field is absent in the CS term. A different approach to study the equivalence has been adopted in cw . Using an iterative embedding method cwjrf for the NCSD model, a dual equivalent theory was constructed to all orders in the NC parameter. This dual model differs from NCMCS theory in the coefficient of the cubic interaction of the CS term and this breaks gauge invariance. In davi , the SW mapped NCMCS theory was argued to be equivalent to a theory where the effect of noncommutativity appears through a non-covariant term. This term vanishes in the commutative limit and the SD model is then recovered. It is then imperative, using alternative approaches, to reexamine the relation between NCMCS theory and NCSD model since the previous studies are inconclusive. Also, this result has interesting implications for deriving the bosonization rules for the NC massive Thirring model sg ; botta .
In this paper we use a procedure which was applied to get a dual description of the sigma model bp and was also used recently to show the equivalence between massive Abelian gauge theories in $`3+1`$ dimensions hm . We first apply the procedure to the partition function of the SW mapped NCMCS theory to order $`\theta `$ and derive the dual theory also to order $`\theta `$. We then argue that this result can be extended to all orders in $`\theta `$. From the dual theory constructed, we show that the equivalence between the MCS theory and the SD model do not get generalized to the NC setting. In our way to derive the SW map for the NCMCS theory we found that the presence of a massive coupling constant turns the map ambiguous. An infinite number of terms can be present in the map but we choose the minimal set required by the map.
## II Ambiguity in the Seiberg-Witten Map
The SW map is obtained by requiring that an ordinary gauge transformation on $`A_\mu `$ with parameter $`\lambda `$ is equivalent to a NC gauge transformation on $`\widehat{A}_\mu `$ with gauge parameter $`\widehat{\lambda }`$ so that ordinary gauge fields that are gauge equivalent are mapped into NC gauge fields that are also equivalent. In four dimension, where it was originally derived, the SW map for the Abelian gauge theory to first order in $`\theta `$ is given by
$`\widehat{A}_\mu `$ $`=`$ $`A_\mu {\displaystyle \frac{1}{2}}\theta ^{\alpha \beta }A_\alpha (2_\beta A_\mu _\mu A_\beta ),`$ (1)
$`\widehat{\lambda }`$ $`=`$ $`\lambda +{\displaystyle \frac{1}{2}}\theta ^{\alpha \beta }_\alpha \lambda A_\beta .`$ (2)
The NC action, when expanded to first order in $`\theta `$,
$$\widehat{S}=\frac{1}{4}d^4x\widehat{f}^{\mu \nu }(\widehat{f}_{\mu \nu }+2\theta ^{\alpha \beta }_\alpha \widehat{A}_\mu _\beta \widehat{A}_\nu ),$$
(3)
with $`\widehat{f}_{\mu \nu }=_\mu \widehat{A}_\nu _\nu \widehat{A}_\mu `$, gives rise to the SW action
$$S_{SW}=\frac{1}{4}d^4x\left[f^2+2\theta ^{\alpha \beta }(f^{\mu \nu }f_{\mu \alpha }f_{\nu \beta }\frac{1}{4}f_{\alpha \beta }f^2)\right].$$
(4)
The question we are interested in is the freedom allowed by the SW map. Due to its nature we can add to the map (1) any gauge invariant term built with $`\theta `$ and derivatives of the gauge field with the right dimension and the new map will still be a SW map. The question is then how the SW action will be affected. To answer this question let us note that by adding to the map (1) a term like
$$\delta \widehat{A}_\mu =\theta ^{\alpha \beta }T_{\mu \alpha \beta },$$
(5)
we get a contribution to the action (4) like
$$\delta \widehat{S}=d^4x\theta ^{\alpha \beta }f^{\mu \nu }_\mu T_{\nu \alpha \beta }.$$
(6)
Then if this integral vanishes we will not get any new contribution to the SW action. Since in four dimensions the gauge field has dimension one the only gauge invariant terms we can add to the SW map have $`T_{\mu \alpha \beta }`$ of the form $`_\mu f_{\alpha \beta }`$, $`_\alpha f_{\mu \beta }`$ and $`^\rho f_{\rho \beta }\eta _{\alpha \mu }`$. The first term is a gauge transformation to order $`\theta `$ taik and gives no contribution to the SW action. The second one is proportional to the first after applying the Bianchi identity. Finally, the third term gives no contribution to the action since the integral in (6) vanishes. Then the SW map to order $`\theta `$ is essentially unique in four dimensions. However, as we shall see, in three dimensions the situation is completely different.
In three dimensions the NCMCS theory is described by the Lagrangian
$$\widehat{}_{NCMCS}=\frac{1}{4g^2}\widehat{F}_{\mu \nu }\widehat{F}^{\mu \nu }+\frac{\mu }{2}ϵ_{\mu \nu \lambda }\widehat{A}^\mu (\widehat{F}^{\nu \lambda }+\frac{2i}{3}\widehat{A}^\nu \widehat{A}^\lambda ),$$
(7)
where $`\widehat{F}_{\mu \nu }=_\mu \widehat{A}_\nu _\nu \widehat{A}_\nu i[\widehat{A}_\mu ,\widehat{A}_\nu ]_{}`$ while the NCSD model with a compensating Stückelberg field has a Lagrangian given by
$$\widehat{}_{NCSD}=\frac{g^2}{2}(\widehat{f}_\mu \widehat{b}_\mu )(\widehat{f}^\mu \widehat{b}^\mu )\frac{1}{2k}ϵ_{\mu \nu \lambda }\widehat{f}^\mu (^\nu \widehat{f}^\lambda \frac{2i}{3}\widehat{f}^\nu \widehat{f}^\lambda ),$$
(8)
where $`\widehat{b}_\mu =i\widehat{𝒰}^1_\mu \widehat{𝒰},\widehat{𝒰}U(1)`$. The NCMCS theory is invariant under the $`U(1)`$ gauge transformation
$$\widehat{A}_\mu \widehat{U}^1\widehat{A}_\mu \widehat{U}+i\widehat{U}^1_\mu \widehat{U},$$
(9)
while the NC Stückelberg-SD Lagrangian is invariant under
$`\widehat{f}_\mu `$ $``$ $`\widehat{U}^1\widehat{f}_\mu \widehat{U}+i\widehat{U}^1_\mu \widehat{U},`$
$`\widehat{𝒰}`$ $``$ $`\widehat{𝒰}\widehat{U}.`$ (10)
We should remark that for the pure NCCS theory the SW map has the form (1) if the CS coefficient $`\mu `$ is chosen to be dimensionless so that the gauge field has dimension one. The pure NCCS theory has the remarkable property that the SW action has no dependence whatsoever in $`\theta `$ gran .
In the NCMCS theory and NCSD model the situation is rather different since one of the couplings must be dimensionfull and this choice determines the gauge field dimensionality. If we make the usual choice for the gauge field dimensionality to be one then $`g^2`$ in the NCMCS theory has dimension one and we can use the SW map (1) to obtain
$$_{SW}=\frac{1}{4g^2}\left[F_{\mu \nu }F^{\mu \nu }+2\theta ^{\alpha \beta }F_{\alpha \mu }F_{\beta \nu }F^{\mu \nu }\frac{1}{2}\theta ^{\alpha \beta }F_{\alpha \beta }F_{\mu \nu }F^{\mu \nu }\right]+\frac{\mu }{4}ϵ^{\mu \nu \lambda }A_\mu F_{\nu \lambda }.$$
(11)
The fact that $`g^2`$ has dimension one means now that the SW map (1) has an arbitrariness since we can add an infinite number of gauge invariant terms, all linear in $`\theta `$, but with different powers of derivatives of $`F_{\mu \nu }`$. These arbitrary terms in the SW map have the form $`g^6\theta ^{\alpha \beta }T_{\mu \alpha \beta }`$ where the $`g^6`$ factor was chosen so that $`T_{\mu \alpha \beta }`$ is a dimensionless function of $`F_{\mu \nu }`$ and its derivatives times an appropriate power of $`g`$. We should then ask whether such terms contribute to the SW action (11). We find that their contribution has the form
$$d^3xF^{\mu \nu }(_\mu T_{\nu \alpha \beta }\mu g^2ϵ_{\mu \nu \rho }T_{}^{\rho }{}_{\alpha \beta }{}^{}).$$
(12)
Let us now examine the first terms in the expansion of $`T_{\mu \alpha \beta }`$ in powers of $`1/g`$. The leading terms are
$$\frac{1}{g^4}ϵ_{\alpha \beta \rho }F_{}^{\rho }{}_{\mu }{}^{},\frac{1}{g^4}ϵ_{\mu [\alpha }^{}{}_{}{}^{\rho }F_{\beta ]\rho }.$$
(13)
The first term can be removed by a gauge transformation and a rigid translation while for the second term (12) vanishes so both can be disregarded. The next terms have the form
$$\frac{1}{g^6}_\mu F_{\alpha \beta },\frac{1}{g^6}_{[\alpha }F_{\beta ]\mu },$$
(14)
and again the first term can be removed by a gauge transformation while the second is proportional to the first after using the Bianchi identity. Higher order terms, however, can contribute. For instance, to order $`1/g^8`$ we find that $`ϵ_{\mu \alpha \beta }F^2`$ gives a non trivial contribution since (12) does not vanish. Its contribution to the SW action (11) is
$$\frac{1}{g^4}\theta ^{\alpha \beta }ϵ_{\alpha \beta \mu }F^2_\nu F^{\mu \nu }\frac{2\mu }{g^2}\theta ^{\alpha \beta }F_{\alpha \beta }F^{\mu \nu }F_{\mu \nu }.$$
(15)
Notice that we get a contribution of order $`1/g^2`$ and the coefficient of such a contribution could be chosen to cancel the corresponding term in (11).
The ambiguity found here is not of the same sort as that found by successive applications of the SW map taik . Here it arises because the model has a dimensionfull coupling constant. If we require the SW map to be universal in the sense that it applies to any gauge theory then such terms are not present. We will take this point of view from now on.
In nel the SW map for the NC Stückelberg-Proca theory has been obtained by requiring that in the unitary gauge it gives the Proca theory. Using the same criterion, the SW map for the NC Stückelberg-SD model is found to be
$`\widehat{f}_\mu `$ $`=`$ $`f_\mu {\displaystyle \frac{1}{2}}\theta ^{\alpha \beta }b_\alpha (2_\beta f_\mu _\mu b_\beta ),`$
$`\widehat{b}_\mu `$ $`=`$ $`b_\mu +{\displaystyle \frac{1}{2}}\theta ^{\alpha \beta }_\alpha b_\mu b_\beta ,`$ (16)
while the gauge parameter transforms as
$$\widehat{\alpha }=\alpha \frac{1}{2}\theta ^{\alpha \beta }b_\alpha _\beta \alpha .$$
(17)
Applying the map to (8) we obtain the SW mapped action
$`L_{SWSD}`$ $`=`$ $`{\displaystyle d^3x\frac{g^2}{2}\left[(f_\mu b_\mu )(f^\mu b^\mu )+\theta ^{\alpha \beta }(f_\mu b_\mu )(2b_\alpha _\beta f_\mu b_\alpha _\mu b_\beta +_\alpha b_\mu b_\beta )\right]}`$ (18)
$``$ $`{\displaystyle \frac{1}{4k}}{\displaystyle d^3xϵ_{\mu \nu \lambda }f^{\mu \nu }f^\lambda }\theta ^{\alpha \beta }ϵ_{\mu \nu \lambda }\left[f^{\mu \nu }b_\alpha (2_\beta f^\lambda ^\mu b_\beta )+{\displaystyle \frac{4}{3}}f^\mu _\alpha f^\nu _\beta f^\lambda \right].`$
## III Equivalence of the MCS theory and the SD model
In order to make the procedure of deriving the dual theory in NC space-time more transparent and also to set up our notation, we present a brief derivation of the well known equivalence between the MCS theory and the SD model in commutative space-time. The MCS theory described by the Lagrangian
$$_{MCS}=\frac{1}{4g^2}F_{\mu \nu }F^{\mu \nu }+\frac{\mu }{2}ϵ_{\mu \nu \lambda }A^\mu ^\nu A^\lambda ,$$
(19)
is invariant under the $`U(1)`$ gauge transformation $`A_\mu A_\mu +_\mu \alpha `$ while the SD model, whose Lagrangian is
$$_{SD}=\frac{g^2}{2}f_\mu f^\mu \frac{1}{2k}ϵ_{\mu \nu \lambda }f^\mu ^\nu f^\lambda ,$$
(20)
has no such an invariance since the $`f_\mu f^\mu `$ term breaks the symmetry. Their equivalence has been analyzed using a phase space path integral approach rab and it was shown that the SD model is equivalent to a gauge fixed version of MCS theory. Also, this equivalence has been studied within the generalized canonical framework of Batalin and Fradkin in rab1 . It was shown that the gauge invariant formulation obtained by the Hamiltonian embedding of SD model is equivalent to the $`U(1)`$ invariant MCS theory, clarifying the equivalence between both theories in spite of the fact that they have different gauge structures. The procedure employed here also sheds light into this issue as we shall see.
The MCS theory is also invariant under a global shift of the vector field $`A_\mu A_\mu +\xi _\mu `$ apart from the $`U(1)`$ gauge invariance. We first elevate this global shift symmetry to a local one by gauging it by an appropriate antisymmetric gauge field $`G_{\mu \nu }`$ which transforms as $`G_{\mu \nu }G_{\mu \nu }+_\mu \xi _\nu _\nu \xi _\mu `$. To have the same physical content as our starting MCS theory we then constrain this gauge field to be non-propagating. This is done by introducing a Lagrange multiplier $`\mathrm{\Phi }`$ which imposes the dual field strength of this gauge field to be flat. The result is
$``$ $`=`$ $`{\displaystyle \frac{1}{4g^2}}(F_{\mu \nu }G_{\mu \nu })(F^{\mu \nu }G^{\mu \nu })+{\displaystyle \frac{\mu }{4}}ϵ_{\mu \nu \lambda }P^\mu (F^{\nu \lambda }G^{\nu \lambda }){\displaystyle \frac{\mu }{8}}ϵ_{\mu \nu \lambda }P^\mu ^\nu P^\lambda `$ (21)
$`+`$ $`{\displaystyle \frac{1}{4}}ϵ_{\mu \nu \lambda }G^{\mu \nu }^\lambda \mathrm{\Phi }+{\displaystyle \frac{1}{4}}ϵ_{\mu \nu \lambda }J^\mu (F^{\nu \lambda }G^{\nu \lambda }),`$
where we have introduced an auxiliary field $`P_\mu `$ to linearize the CS term. This field has a $`U(1)`$ gauge invariance $`P_\mu P_\mu +_\mu \chi `$ when the multiplier field transforms as $`\mathrm{\Phi }\mathrm{\Phi }+\mu \chi `$ and $`A_\mu A_\mu `$. The last term in the Lagrangian is a source $`J^\mu `$ coupling to the local shift invariant combination of $`A_\mu `$ and $`G_{\mu \nu }`$. The MCS theory is recovered from the above Lagrangian by eliminating the $`\mathrm{\Phi }`$ field using its equation of motion.
To show the equivalence to the SD model we start from the partition function
$$Z=D\mathrm{\Phi }DP_\mu DA_\mu DG_{\mu \nu }e^{i{\scriptscriptstyle d^3x}}.$$
(22)
Integrations over $`G_{\mu \nu }`$ and $`A_\mu `$ are Gaussian and can be done trivially leading to
$$Z_{dual}=D\mathrm{\Phi }DP_\mu e^{i{\scriptscriptstyle d^3x_{eff}}}.$$
(23)
After the redefinitions $`\mu P_\mu =f_\mu `$ and $`\mathrm{\Phi }=\mathrm{\Lambda }`$, we get the effective Lagrangian
$$_{eff}=\frac{g^2}{8}(f_\mu _\mu \mathrm{\Lambda })(f^\mu ^\mu \mathrm{\Lambda })\frac{1}{8\mu }ϵ_{\mu \nu \lambda }f^\mu ^\nu f^\lambda +\frac{g^2}{8}J_\mu J^\mu +\frac{g^2}{4}(f^\mu ^\mu \mathrm{\Lambda })J_\mu .$$
(24)
This theory is invariant under the $`U(1)`$ gauge transformation $`f_\mu f_\mu +_\mu \alpha `$ when the Stückelberg field transforms as $`\mathrm{\Lambda }\mathrm{\Lambda }+\alpha `$. We also note that the MCS coupling constant $`g^2`$ and the Chern-Simons parameter $`\mu `$ have both appeared as inverse couplings when compared with (20). We can now fix the gauge invariance in (24), for instance by choosing the unitary gauge $`\mathrm{\Lambda }=0`$, to recover the self-dual model given in (20). We thus conclude that the $`U(1)`$ invariant MCS theory is dual to the $`U(1)`$ invariant Stückelberg formulation of self-dual model.
From the partition functions (22) and (23) we derive the mapping between the n-point correlators for these theories. For the 2-point function, we get
$$ϵ_{\mu \nu \lambda }F^{\nu \lambda }(x)ϵ_{\alpha \beta \rho }F^{\beta \rho }(y)g^4(f_\mu _\mu \mathrm{\Lambda })(x)(f_\alpha _\alpha \mathrm{\Lambda })(y)+g^2g_{\mu \alpha }\delta (xy),$$
(25)
leading the identification (up to non-propagating contact terms) between the gauge invariant combinations
$$ϵ_{\mu \nu \lambda }F^{\nu \lambda }g^2(f_\mu _\mu \mathrm{\Lambda }).$$
(26)
This equivalence between SD model and MCS theory has been extended to include interaction with matter cwjrf . It has been shown that the SD model minimally coupled to charged dynamical fermionic and bosonic matter fields is equivalent to a MCS theory non-minimally coupled to matter. In the weak coupling limit, it was shown in nb that the non-Abelian MCS theory is equivalent to non-Abelian SD model and recently it was shown that, perturbatively, this equivalence exists in all regimes of the coupling constant botta1 .
After re-expressing the NCMCS theory (7) in terms of $`A_\mu `$ and $`\theta ^{\alpha \beta }`$ using the SW map (1) we apply the above procedure to construct the corresponding dual theory. Then by comparing this dual theory with SW mapped NC Stückelberg-SD model, we study the status of their equivalence. We take up this in the next section.
## IV Seiberg-Witten mapped Maxwell-Chern-Simons theory and duality
By applying the SW map (1) to the NCMCS Lagrangian (7) we get to order $`\theta `$
$`_{SW}`$ $`=`$ $`{\displaystyle \frac{1}{4g^2}}\left[F_{\mu \nu }F^{\mu \nu }+2\theta ^{\alpha \beta }F_{\alpha \mu }F_{\beta \nu }F^{\mu \nu }{\displaystyle \frac{1}{2}}\theta ^{\alpha \beta }F_{\alpha \beta }F_{\mu \nu }F^{\mu \nu }\right]`$ (27)
$`+`$ $`{\displaystyle \frac{\mu }{4}}ϵ_{\mu \nu \lambda }P^\mu F^{\nu \lambda }{\displaystyle \frac{\mu }{8}}ϵ_{\mu \nu \lambda }P^\mu ^\nu P^\lambda ,`$
where an auxiliary field $`P_\mu `$ was introduced to linearize the CS term. We have also used the fact that the NCCS term gets mapped to the usual commutative CS term by the SW map gran . After rewriting the above Lagrangian using auxiliary fields $`B_{\mu \nu }`$ and $`C_{\mu \nu }`$ as
$`_{SW}`$ $`=`$ $`{\displaystyle \frac{1}{4g^2}}C_{\mu \nu }B^{\mu \nu }{\displaystyle \frac{\mu }{8}}ϵ_{\mu \nu \lambda }P^\mu ^\nu P^\lambda +{\displaystyle \frac{\mu }{4}}ϵ_{\mu \nu \lambda }P^\mu F^{\nu \lambda }`$ (28)
$``$ $`{\displaystyle \frac{1}{4g^2}}\left[F_{\mu \nu }F^{\mu \nu }+2\theta ^{\alpha \beta }C_{\alpha \mu }C_{\beta \nu }F^{\mu \nu }{\displaystyle \frac{1}{2}}\theta ^{\alpha \beta }C_{\alpha \beta }C_{\mu \nu }F^{\mu \nu }B_{\mu \nu }F^{\mu \nu }\right],`$
we can now gauge the shift invariance of $`A_\mu `$ field as in the commutative case. Due to the introduction of $`B_{\mu \nu }`$ and $`C_{\mu \nu }`$ we see that $`G_{\mu \nu }`$ will appear quadratically and this will simplify the calculation considerably. So, we introduce a gauge field $`G_{\mu \nu }`$ to promote the global shift invariance of $`A_\mu `$ to a local one. We then get
$`_{SW}`$ $`=`$ $`{\displaystyle \frac{1}{4g^2}}C_{\mu \nu }B^{\mu \nu }{\displaystyle \frac{\mu }{24}}ϵ_{\mu \nu \lambda }P^\mu ^\nu P^\lambda +{\displaystyle \frac{\mu }{4}}ϵ_{\mu \nu \lambda }P^\mu (F^{\nu \lambda }G^{\nu \lambda }){\displaystyle \frac{1}{4}}ϵ_{\mu \nu \lambda }G^{\mu \nu }^\lambda \mathrm{\Phi }`$ (29)
$``$ $`{\displaystyle \frac{1}{4g^2}}\left[(F_{\mu \nu }G_{\mu \nu })+(2\theta ^{\alpha \beta }C_{\alpha \mu }C_{\beta \nu }{\displaystyle \frac{1}{2}}\theta ^{\alpha \beta }C_{\alpha \beta }C_{\mu \nu })B_{\mu \nu }\right](F^{\mu \nu }G^{\mu \nu }).`$
Starting with the partition function
$$Z=DP_\mu D\mathrm{\Phi }DC_{\mu \nu }DB_{\mu \nu }DA_\mu DG_{\mu \nu }e^{i{\scriptscriptstyle d^x_{SW}}},$$
(30)
we can integrate over $`G_{\mu \nu }`$, $`A_\mu `$ and $`B_{\mu \nu }`$ to get the partition function corresponding to the effective Lagrangian
$`_{eff}`$ $`=`$ $`{\displaystyle \frac{\mu }{8}}ϵ_{\mu \nu \lambda }P^\mu ^\nu P^\lambda {\displaystyle \frac{1}{4g^2}}C_{\mu \nu }C^{\mu \nu }+{\displaystyle \frac{1}{4}}ϵ_{\mu \nu \lambda }C^{\mu \nu }(\mu P^\lambda ^\lambda \mathrm{\Phi })`$ (31)
$``$ $`{\displaystyle \frac{1}{4g^2}}C^{\mu \nu }\left[2\theta ^{\alpha \beta }C_{\alpha \mu }C_{\beta \nu }{\displaystyle \frac{1}{2}}\theta ^{\alpha \beta }C_{\alpha \beta }C_{\mu \nu }\right].`$
We have neglected higher order terms in $`\theta `$ in performing the Gaussian integrals. It is easy to see that in the commutative limit we get (20) when $`C_{\mu \nu }`$ is eliminated by using its field equation and setting $`\mathrm{\Phi }=0`$.
In the NC case $`C_{\mu \nu }`$ can be eliminated perturbatively in $`\theta `$. We then get
$`_{dual}`$ $`=`$ $`{\displaystyle \frac{g^2}{8}}(f_\mu _\mu \mathrm{\Lambda })(f^\mu ^\mu \mathrm{\Lambda })+{\displaystyle \frac{g^4}{32}}\theta ^{\alpha \beta }ϵ_{\alpha \beta \lambda }(f^\lambda ^\lambda \mathrm{\Lambda })(f^\mu ^\mu \mathrm{\Lambda })(f_\mu _\mu \mathrm{\Lambda })`$ (32)
$``$ $`{\displaystyle \frac{1}{8\mu }}ϵ_{\mu \nu \lambda }f^\mu ^\nu f^\lambda ,`$
where we have identified $`\mu P_\mu =f_\mu `$ and $`\mathrm{\Phi }=\mathrm{\Lambda }`$. As in the commutative case the strong coupling limit of the original theory gets mapped into the weak coupling limit of the dual. It is easy to see that in the limit of vanishing $`\theta `$ the above Lagrangian (in the unitary gauge where $`\mathrm{\Lambda }=0`$) correctly reproduces the SD Lagrangian (20).
It is interesting to note that the explicit form of the order $`\theta `$ term in the $`C_{\mu \nu }`$ field equation is not needed at all to find the above Lagrangian. This happens because there are nice cancellations and it is easy to be convinced that to obtain the dual Lagrangian to $`n`$-th order in $`\theta `$ we need the perturbative solution for $`C_{\mu \nu }`$ only to order $`(n1)`$.
We can couple a source term $`ϵ_{\mu \nu \lambda }F^{\mu \nu }J^\lambda `$ to the Lagrangian (27) and this leads to the map between the 2-point functions
$`ϵ_{\mu \nu \lambda }F^{\nu \lambda }(x)ϵ_{\alpha \beta \rho }F^{\beta \rho }(y)g^4\stackrel{~}{f}_\mu (x)\stackrel{~}{f}(y)+g^2g_{\mu \alpha }\delta (xy)`$
$`+{\displaystyle \frac{g^8}{64}}\overline{\theta }_\mu \stackrel{~}{f}^\nu \stackrel{~}{f}_\nu +2\overline{\theta }_\nu \stackrel{~}{f}^\nu \stackrel{~}{f}_\mu \overline{\theta }_\alpha \stackrel{~}{f}^\beta \stackrel{~}{f}_\beta +2\overline{\theta }_\beta \stackrel{~}{f}^\beta \stackrel{~}{f}_\alpha +g^4(\overline{\theta }_\mu \stackrel{~}{f}_\alpha +\overline{\theta }_\alpha \stackrel{~}{f}_\mu +\overline{\theta }_\beta \stackrel{~}{f}_\beta g_{\mu \alpha }),`$ (33)
where $`\overline{\theta }_\mu =ϵ_{\mu \nu \lambda }\theta ^{\nu \lambda }`$ and $`\stackrel{~}{f}_\mu =f_\mu _\mu \mathrm{\Lambda }`$. In the limit $`\theta 0`$ we recover the map obtained in (25).
Here we note that all the $`\theta `$ dependence of the SW mapped NCMCS theory comes from the Maxwell term alone as the NCCS term gets mapped to usual commutative CS term. Since it is possible to express the SW mapped Maxwell action to all orders in $`\theta `$ in terms of the commutative field strength $`F_{\mu \nu }`$ and $`\theta `$ alone gop (an exact closed form for the SW mapped Maxwell action is given in hsy ), it is easy to convince from (27) and (29) that the procedure adopted here can be used to construct the dual theory to all orders in $`\theta `$ using a perturbative solution for the $`C_{\mu \nu }`$ field equation.
One important point to note is that the theory described by the Lagrangian (32), which is equivalent to the SW mapped NCMCS theory, is not the same as the SW mapped action for NCSD model (18). This clearly shows that the SW mapped theories are not equivalent.
## V Conclusion
In this paper we have constructed and studied the dual description of the NCMCS theory and investigated the status of the equivalence between this theory and SD model. We have derived the dual theory starting from the SW mapped NCMCS Lagrangian which is given in terms of commutative fields and the NC parameter. The equivalence was obtained at the level of partition functions and it allowed us to get the mapping between the n-point correlators of both theories. We have shown that the dual theory does not coincide with the SW mapped NC Stückelberg-SD theory. However, in the commutative limit, we recover the well known equivalence between them. We have also shown that the two-point correlators map reduces to the one obtained in the commutative case in this limit. We have argued that this result can be extended to all orders in $`\theta `$ due to the structure of the SW mapped NCMCS Lagrangian. We have also verified that even after accounting for the ambiguous terms in the SW map, the dual theory and SW mapped NC Stückelberg-SD model are not equivalent.
Hence, we have shown that the equivalence between the MCS theory and the SD model in commutative space-time does not survive in the NC case. In this respect we are in agreement with the results obtained earlier in sg and cw where it was argued that these NC theories are not equivalent. But unlike the NCCS term used in sg , we have used the standard NC $`U(1)`$ invariant CS term with a cubic interaction as in cw and botta . The non-equivalence between the NCSD model and NCMCS theory shown here will come as an obstacle in generalizing the bosonization of the commutative Thirring model to NC space-time as was pointed out in sg ; cw .
ACKNOWLEDGMENTS:
EH thanks FAPESP for support through grant 03/09044-9. The work of VOR was partially supported by CNPq, FAPESP and PRONEX under contract CNPq 66.2002/1998-99.
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# An invariant approach to dynamical fuzzy spaces with a three-index variable
## 1 Introduction and a general proposal
Several thought experiments in semi-classical quantum gravity and string theory - suggest that the classical notion of space-time in general relativity should be replaced with a quantum one in quantum treatment of space-time dynamics. An interesting candidate is given by the non-commutative geometry or the fuzzy space -. The central principle of general relativity is the invariance under the general coordinate transformation. It would be obviously interesting if a fuzzy space can be obtained from formulation invariant under a fuzzy analog of the general coordinate transformation. This kind of formulation would lead to fuzzy general relativity, which may describe the gravitational dynamics and the evolution of the universe in terms of dynamical fuzzy spaces<sup>1</sup><sup>1</sup>1The present author considered evolving fuzzy spaces in -..
A fuzzy space is characterized by the algebraic relations $`f_af_b=C_{ab}^{}{}_{}{}^{c}f_c`$ among the functions $`f_a`$ on the fuzzy space. Therefore it would be a natural assumption that a dynamical fuzzy space can be described by treating the three-index variable $`C_{ab}^{}{}_{}{}^{c}`$ dynamically. It will be argued in the following section that a fuzzy analog of the general coordinate transformation is given by the general linear transformation on $`f_a`$. Thus the central proposal of this paper is that dynamical fuzzy spaces are described by equations for $`C_{ab}^{}{}_{}{}^{c}`$ invariant under the general linear transformation.
In the following section, a fuzzy analog of the general coordinate transformation will be discussed and the general form of the models will be proposed. In Section 3, it will be shown that the solutions to the equations can be constructed from the invariant tensors of Lie groups. In Section 4, the Lie group will be specified to SO(3), and some series of the solutions will be explicitly constructed. In Section 5, I will discuss the construction of a scalar field theory on a fuzzy two-sphere by using the result in the preceding section. The final section will be devoted to discussions and comments.
## 2 A fuzzy general coordinate transformation and the general form of the models
Let me first review the general coordinate transformation in the usual commutative space $`R^d`$. A basis of the continuous functions on $`R^d`$ can be given by
$$\{1,x^i,x^ix^j,x^ix^jx^k,\mathrm{}\},$$
(1)
where $`x^i(i=1,\mathrm{},d)`$ are the coordinates of $`R^d`$. Let $`f_a`$ with an index $`a`$ denote these independent functions in the set. A continuous function on $`R^d`$ is given by a linear combination of $`f_a`$.
Let me consider a general coordinate transformation,
$$x_{}^{}{}_{}{}^{i}=f^i(x).$$
(2)
It is a natural restriction that the coordinate transformation is not singular, and is invertible. Since the right-hand side is a continuous function, the general coordinate transformation can be represented by a linear transformation,
$$x_{}^{}{}_{}{}^{i}=M_{i}^{}{}_{}{}^{a}f_a,$$
(3)
where $`M_{i}^{}{}_{}{}^{a}`$ are real.
Let me now consider a fuzzy space with a finite number of independent functions $`f_a(a=1,\mathrm{},n)`$ on the fuzzy space. This is a fuzzy space corresponding to a compact space in usual continuous theory. I assume that all the $`f_a`$ denote real functions on a fuzzy space, and that the variable $`C_{ab}^{}{}_{}{}^{c}`$, which determines the algebraic relations $`f_af_b=C_{ab}^{}{}_{}{}^{c}f_c`$, is real. I do not assume the associativity or the commutativity of the algebra, so that $`C_{ab}^{}{}_{}{}^{c}`$ has no other constraints.
An important assumption of this paper is to interpret the transformation rule (3) as a partial appearance of a more general linear transformation,
$$f_a^{}=M_{a}^{}{}_{}{}^{b}f_b,$$
(4)
where $`M_{a}^{}{}_{}{}^{b}`$ can take any real values provided that the matrix $`M_{a}^{}{}_{}{}^{b}`$ is invertible, which comes from the assumed invertibility of the coordinate transformation (2). Therefore the fuzzy general coordinate transformation of this paper is given by the GL($`n`$,$`R`$) transformation on $`f_a`$. Under the transformation (4), the three-index variable $`C_{ab}^{}{}_{}{}^{c}`$ transforms in the way,
$$C_{ab}^{}{}_{}{}^{c}=M_{a}^{}{}_{}{}^{a^{}}M_{b}^{}{}_{}{}^{b^{}}C_{a^{}b^{}}^{}{}_{}{}^{c^{}}(M^1)_{c^{}}^{}{}_{}{}^{c}.$$
(5)
I impose that the equations of motion for $`C_{ab}^{}{}_{}{}^{c}`$ must be invariant under this GL($`n`$,$`R`$) transformation.
For convenience, let me introduce a graphical expression. The three-index variable $`C_{ab}^{}{}_{}{}^{c}`$ can be graphically represented as in Fig. 1.
An example of an invariant equation of motion is given by
$$C_{ia}^{}{}_{}{}^{j}C_{jb}^{}{}_{}{}^{k}C_{kc}^{}{}_{}{}^{i}+C_{ai}^{}{}_{}{}^{j}C_{cj}^{}{}_{}{}^{k}C_{bk}^{}{}_{}{}^{i}=0,$$
(6)
which is graphically represented in Fig. 2. It is clear that such invariant equations of motion can be made in infinitely many ways.
In the next, let me discuss the construction of an action. From the transformation property (5), the lower and upper indices must be contracted to make an invariant under the GL($`n`$,$`R`$) transformation. Since an action must be invariant and $`C_{ab}^{}{}_{}{}^{c}`$ has more lower indices than upper, it is necessary to introduce an additional variable which has more upper indices to construct an invariant action. A way to achieve this is to define such a variable from the $`C_{ab}^{}{}_{}{}^{c}`$ itself. For example, one can define
$$(C_{ai}^{}{}_{}{}^{j}C_{jb}^{}{}_{}{}^{i})^1,$$
(7)
which has two upper indices.
This procedure, however, will not work when the matrix in the parentheses in (7) is not invertible, and will also let an action be a complicated function of $`C_{ab}^{}{}_{}{}^{c}`$. Therefore, rather than defining a quantity like (7) from the beginning, I simply introduce a new variable with upper indices $`g^{ab}`$ and assume it be determined from the equations of motion derived from an action. I impose its reality and the symmetry of the two indices,
$$g^{ab}=g^{ba}.$$
(8)
The grahical representation of $`g^{ab}`$ is given in Fig. 3.
Thus an invariant action is a function of the two dynamical real variables $`g^{ab}`$ and $`C_{ab}^{}{}_{}{}^{c}`$,
$$S(g^{ab},C_{ab}^{}{}_{}{}^{c}),$$
(9)
where all the indices are contracted. The graphical representation of an invariant action is a closed diagram with oriented lines connecting blobs and three-vertices.
## 3 The classical solutions to the equations of motion
The quantum mechanical treatment of the model presented in the previous section would be obviously very interesting, but in this paper I restrict myself to the classical solutions to the equations of motion derived from the invariant action (9). Let me suppose that the equations of motion for $`C_{ab}^{}{}_{}{}^{c}`$,
$$\frac{S}{C_{ab}^{}{}_{}{}^{c}}=0,$$
(10)
are satisfied at $`C_{ab}^{}{}_{}{}^{c}=C_{ab}^{0}{}_{}{}^{c},g^{ab}=g_0^{ab}`$. For simplicity, let me assume $`g_0^{ab}`$ is invertible as a matrix. At the general values of $`C_{ab}^{}{}_{}{}^{c}`$ and $`g^{ab}`$, the invariance of the action under the GL($`n`$,$`R`$) transformation implies
$$\frac{S}{g^{ab}}\delta g^{ab}+\frac{S}{C_{ab}^{}{}_{}{}^{c}}\delta C_{ab}^{}{}_{}{}^{c}=0,$$
(11)
where $`\delta g^{ab}`$ and $`\delta C_{ab}^{}{}_{}{}^{c}`$ are the infinitesimally small GL($`n`$,$`R`$) transformation of $`g^{ab}`$ and $`C_{ab}^{}{}_{}{}^{c}`$,
$`\delta g^{ab}`$ $`=`$ $`\delta M_{i}^{}{}_{}{}^{a}g^{ib}\delta M_{i}^{}{}_{}{}^{b}g^{ai},`$ (12)
$`\delta C_{ab}^{}{}_{}{}^{c}`$ $`=`$ $`\delta M_{a}^{}{}_{}{}^{i}C_{ib}^{}{}_{}{}^{c}+\delta M_{b}^{}{}_{}{}^{i}C_{ai}^{}{}_{}{}^{c}\delta M_{i}^{}{}_{}{}^{c}C_{ab}^{}{}_{}{}^{i},`$ (13)
where $`\delta M_{i}^{}{}_{}{}^{a}`$ can take any infinitesimally small real values. At the solution to (10), $`C_{ab}^{}{}_{}{}^{c}=C_{ab}^{0}{}_{}{}^{c},g^{ab}=g_0^{ab}`$, (11) becomes
$$\delta M_{i}^{}{}_{}{}^{a}g_0^{ib}\frac{S}{g^{ab}}|_{\begin{array}{c}g=g_0\\ C=C^0\end{array}}=0,$$
(14)
where I have used (8), (10) and (12). Therefore, since $`\delta M_{i}^{}{}_{}{}^{a}`$ is arbitrary and $`g_0^{ab}`$ is invertible,
$$\frac{S}{g^{ab}}|_{\begin{array}{c}g=g_0\\ C=C^0\end{array}}=0$$
(15)
are satisfied. This means that the equations of motion for $`g^{ab}`$ are always simultaneously satisfied when the equations of motion for $`C_{ab}^{}{}_{}{}^{c}`$ are satisfied. Therefore it is enough to consider only the equations of motion (10) to find the classical solutions, provided $`g^{ab}`$ is invertible.
Now let me consider a Lie group which has a representation of dimension $`n`$. The representation can be either reducible or irreducible. Let me embed the representation into a classical solution: The lower index of a classical solution is assumed to be transformed in the representation, while the upper one in the inverse representation. In fact, the following discussions do not depend on whether the representation is real or complex, provided that the invariant tensors considered below are real. Let $`g_0^{ab}`$ be a real symmetric rank-two invariant tensor under the inverse representation. I assume the tensor $`g_0^{ab}`$ is invertible as a matrix. Let me introduce $`I_{ab}^{\alpha }{}_{}{}^{c}(\alpha =1,2,\mathrm{},N)`$, which span all the real tensors invariant under the same transformation property as $`C_{ab}^{0}{}_{}{}^{c}`$. Let me assume $`C_{ab}^{0}{}_{}{}^{c}`$ is given by a linear combination of these invariant tensors,
$$C_{ab}^{0}{}_{}{}^{c}=A_\alpha I_{ab}^{\alpha }{}_{}{}^{c},$$
(16)
where $`A_\alpha (\alpha =1,2,\mathrm{},N)`$ are real coefficients. Since the action is obviously invariant under the transformation of the Lie group, the left-hand side of (10) becomes a real invariant tensor when $`g^{ab}`$ and $`C_{ab}^{}{}_{}{}^{c}`$ are substituted with the invariant tensors above. Therefore, after exchanging the upper and lower indices by using $`g_0^{ab}`$ and its inverse $`g_{ab}^0`$, the left-hand side of (10) can be expressed as a linear combination of $`I_{ab}^{\alpha }{}_{}{}^{c}`$,
$$g_{ai}^0g_{bj}^0g_0^{ck}\frac{S}{C_{ij}^{}{}_{}{}^{k}}|_{\begin{array}{c}g=g_0\\ C=C^0\end{array}}=B_\alpha (A_1,A_2,\mathrm{},A_N)I_{ab}^{\alpha }{}_{}{}^{c},$$
(17)
where $`B_\alpha (A_1,A_2,\mathrm{},A_N)`$ are some functions of $`A_1,A_2,\mathrm{},A_N`$, which are determined from the specific form of the left-hand side of (10). Therefore the equations of motion (10) are reduced to the following simultaneous equations for $`A_1,A_2,\mathrm{},A_N`$,
$$B_\alpha (A_1,A_2,\mathrm{},A_N)=0,(\alpha =1,\mathrm{},N).$$
(18)
These equations are much easier to solve than (10). Since the numbers of the variables and the equations are the same, the simultaneous equations (18) will generally have some number of solutions. The solutions are generally complex, but real solutions can be actually found in some interesting cases.
## 4 Explicit solutions for SO(3)
The Lie algebra of SO(3) is given by
$$[J_z,J_\pm ]=\pm J_\pm ,[J_+,J_{}]=2J_z,$$
(19)
and an irreducible representation labeled with an integer spin $`j`$ can be explicitly given by
$`J_z|j,m`$ $`=`$ $`m|j,m,`$ (20)
$`J_\pm |j,m`$ $`=`$ $`\sqrt{j(j+1)m(m\pm 1)}|j,m\pm 1.`$ (21)
It is convenient to associate the index of the model with a pair $`(j,m)`$, where $`j`$ is the spin, and $`m`$ the $`J_z`$ eigenvalue $`(m=j,j+1,\mathrm{},j)`$. The above representation of SO(3) is generally complex but the invariant tensors can be taken real,
$`g_0^{(j_1,m_1)(j_2,m_2)}`$ $`=`$ $`\delta _{j_1j_2}\delta _{m_1m_2}(1)^{m_1},`$ (22)
$`I_{(j_1,m_1)(j_2,m_2)(j_3,m_3)}`$ $`=`$ $`\left(\begin{array}{ccc}j_1& j_2& j_3\\ m_1& m_2& m_3\end{array}\right),`$ (25)
where the right-hand side in the second line is the $`3j`$-symbol . These invariant tensors are essentially unique when the spins of the representations are given.
Let me consider a representation given by the direct sum of a number of the irreducible representations. Here I consider the case that each irreducible representation appears at most once in the direct sum. Then the ansatz (16) becomes
$$C_{(j_1,m_1)(j_2,m_2)(j_3,m_3)}^0=A(j_1,j_2,j_3)I_{(j_1,m_1)(j_2,m_2)(j_3,m_3)},$$
(26)
where the spins $`j_i`$ must be contained in the representation. Since the 3$`j`$-symbol has the cyclic symmetry $`1231`$, this is also imposed on $`A(j_1,j_2,j_3)`$. It is assumed that $`A(j_1,j_2,j_3)`$ vanishes if $`j_i`$ do not satisfy the triangle inequalities because $`I_{(j_1,m_1)(j_2,m_2)(j_3,m_3)}`$ vanishes identically.
The 3$`j`$-symbol satisfies the identity,
$`{\displaystyle \underset{M_i}{}}(1)^{_{i=1}^3l_i+M_i}\left(\begin{array}{ccc}l_1& l_2& j_3\\ M_1& M_2& m_3\end{array}\right)\left(\begin{array}{ccc}l_2& l_3& j_1\\ M_2& M_3& m_1\end{array}\right)\left(\begin{array}{ccc}l_3& l_1& j_2\\ M_3& M_1& m_2\end{array}\right)=`$ (33)
$`\left(\begin{array}{ccc}j_1& j_2& j_3\\ m_1& m_2& m_3\end{array}\right)\left\{\begin{array}{ccc}j_1& j_2& j_3\\ l_1& l_2& l_3\end{array}\right\},`$ (38)
where $`\{:::\}`$ denotes the $`6j`$-symbol. This identity can be used to compute
$$D_{abc}^0=C_{ail}^0C_{bki^{}}^0C_{cl^{}k^{}}^0g_0^{ii^{}}g_0^{ll^{}}g_0^{kk^{}}$$
(39)
in Fig. 4, where the roman indices abbreviate the pairs $`(j,m)`$. The result is
$`D_{(j_1,m_1)(j_2,m_2)(j_3,m_3)}^0=(1)^{_{i=1}^3j_i}\left(\begin{array}{ccc}j_1& j_2& j_3\\ m_1& m_2& m_3\end{array}\right)`$ (42)
$`\times {\displaystyle \underset{l_i}{}}(1)^{_{i=1}^3l_i}\left\{\begin{array}{ccc}j_1& j_2& j_3\\ l_1& l_2& l_3\end{array}\right\}A(j_1,l_3,l_2)A(j_2,l_1,l_3)A(j_3,l_2,l_1).`$ (45)
In the above derivation I have also used the property that the 3$`j`$-symbol in (22) changes its sign by $`(1)^{_{i=1}^3j_i}`$, when a pair of the rows are interchanged. Let me consider the equations of motion,
$$C_{abc}^0gD_{abc}^0=0,$$
(46)
where $`g`$ is a real coupling constant. The graphical representation is given in Fig. 5. The action which leads to the equations of motion (46) can be constructed in several ways. This will not be discussed here, because the classical solutions are the main interest in this paper. From the result (42), the equations of motion are reduced to
$$A(j_1,j_2,j_3)g\underset{l_i}{}(1)^{_{i=1}^3l_i+j_i}\left\{\begin{array}{ccc}j_1& j_2& j_3\\ l_1& l_2& l_3\end{array}\right\}A(j_1,l_3,l_2)A(j_2,l_1,l_3)A(j_3,l_2,l_1)=0.$$
(47)
A series of the solutions to (47) parameterized by a spin parameter $`L`$ can be constructed in the following way. The 6$`j`$-symbol satisfies the identity,
$`{\displaystyle \underset{l_i}{}}(1)^{_{i=1}^3l_i}{\displaystyle \underset{i=1}{\overset{3}{}}}(2l_i+1)\left\{\begin{array}{ccc}j_1& j_2& j_3\\ l_1& l_2& l_3\end{array}\right\}\left\{\begin{array}{ccc}j_1& l_3& l_2\\ b& a_2& a_3\end{array}\right\}\left\{\begin{array}{ccc}j_2& l_1& l_3\\ b& a_3& a_1\end{array}\right\}\left\{\begin{array}{ccc}j_3& l_2& l_1\\ b& a_1& a_2\end{array}\right\}`$ (48)
$`=(2b+1)(1)^{b+_{i=1}^3a_i+j_i}\left\{\begin{array}{ccc}j_1& j_2& j_3\\ a_1& a_2& a_3\end{array}\right\}.`$ (49)
Therefore,
$$A(j_1,j_2,j_3)=\frac{1}{\sqrt{(2L+1)g}}\underset{i=1}{\overset{3}{}}\sqrt{2j_i+1}\left\{\begin{array}{ccc}j_1& j_2& j_3\\ L& L& L\end{array}\right\}$$
(50)
is the solution to (47) for any given spin $`L`$. Note that $`j_i2L`$ for $`A(j_1,j_2,j_3)`$ to be non-zero, because the six arguments of the 6$`j`$-symbol must be the edge lengths of a tetrahedron. Therefore the solution (50) can be embedded in a finite-dimensional model where the index of the model runs over $`(j,m)(j=0,1,\mathrm{},2L;m=j,j+1,\mathrm{},j)`$.
I will now discuss another series of the solutions which will be used in the construction of a scalar field theory on a fuzzy two-sphere. The only non-vanishing components of $`C_{ab}^{0}{}_{}{}^{c}`$ are assumed to be
$$C_{(1,m_1)(L,m_2)}^{0}{}_{}{}^{(L,m_3)}=A\left(\begin{array}{ccc}1& L& L\\ m_1& m_2& m_3\end{array}\right)g_0^{(L,m_3)(L,m_3)},$$
(51)
where $`L`$ is a given spin and $`A`$ is a real coefficient, which will be determined from the equations of motion. The rank-two symmetric tensor $`g_0^{ab}`$ is as before. The index of the model runs over $`(1,m)(m=1,0,1),(L,m)(m=L,L+1,\mathrm{},L)`$. The cyclic symmetry for the indices of $`C_{ab}^{0}{}_{}{}^{c}`$ is not imposed this time, and $`C_{ab}^{0}{}_{}{}^{c}`$ are non-zero only when its first index has a spin one and the others a spin $`L`$. The equations of motion considered are
$$C_{ab}^{}{}_{}{}^{c}gC_{ai}^{}{}_{}{}^{j}C_{ki^{}}^{}{}_{}{}^{b^{}}C_{k^{}c^{}}^{}{}_{}{}^{j^{}}g^{ii^{}}g_{jj^{}}g^{kk^{}}g_{bb^{}}g^{cc^{}}=0,$$
(52)
where $`g`$ is the coupling constant and $`g_{ab}`$ is the inverse of $`g^{ab}`$. Note that the invertibility of $`g^{ab}`$ must be assumed from the beginning in defining this model. The graphical representation is given in Fig. 6.
An action which leads to the equations of motion (52) can be constructed, for example, by summing the squares of the equations of motion with appropriate contractions of the indices by $`g_{ab},g^{ab}`$. Using the identity (33), it can be shown that (51) is actually a solution to (52) if
$$A=\frac{1}{\sqrt{g\left\{\begin{array}{ccc}1& L& L\\ 1& L& L\end{array}\right\}}}=\sqrt{\frac{L(L+1)(2L+1)}{g(L^2+L1)}}.$$
(53)
## 5 A scalar field theory on a fuzzy two-sphere
In the continuum case, the spherical harmonic functions give the independent functions on a two-sphere. They are labeled with the SO(3) indices $`(j,m)`$, where the SO(3) is the rotational symmetry in the three-dimensional Euclidean space embedding the two-sphere. In the discussions in Section 2, $`f_a`$ denotes the functions on a fuzzy space. Therefore it would seem reasonable to consider the solutions (26), (50) to represent the fuzzy two-sphere, because the ranges of the indices match between the spherical harmonics and the model. But I have not succeeded in this direction. The main reason for this failure is that it is difficult to obtain the spectra of the Laplacian $`j(j+1)`$ from the solutions (26), (50). For example, using the properties of the $`3j`$\- and $`6j`$-symbols, one finds
$`C_{akl}^0C_{bl^{}k^{}}^0g_0^{kk^{}}g_0^{ll^{}}={\displaystyle \frac{1}{g}}g_{ab}^0,`$ (54)
where the spins of $`a,b2L`$. There is no dependence on $`j`$ in (54). I have computed some other invariants, but have not found the appropriate dependence on $`j`$.
A way to construct successfully a scalar field theory on a fuzzy two-sphere can be obtained from the solutions (51), (53). This construction is almost similar to the original one . In the paper , the kinetic term of a scalar field theory is obtained from the quadratic Casimir. The generator $`J_i`$ of SO(3) in the spin $`L`$ representation has two spin $`L`$ indices and a spin 1. This index structure is the same as $`C_{(1,m_1)(L,m_2)}^{0}{}_{}{}^{(L,m_3)}`$, and the kinetic term can be constructed in a similar way. An extra issue in the present model is that the indices contain $`(1,m)`$ as well as $`(L,m)`$. It must be checked that the scalar field components with these extra indices do not destroy the wanted spectra.
Let me consider the following four-index invariant,
$$K_{b^{}a}^{a^{}b}=C_{ia}^{0}{}_{}{}^{a^{}}C_{i^{}b^{}}^{0}{}_{}{}^{b}g_0^{ii^{}},$$
(55)
shown in Fig. 7.
This invariant can be regarded as an operator on $`\varphi _b^a`$,
$$K_{b^{}a}^{a^{}b}\varphi _b^a,$$
(56)
where $`\varphi _b^a`$ will be identified as a scalar field on a fuzzy two-sphere.
Since the only non-vanishing components of $`C_{ab}^{0}{}_{}{}^{c}`$ are (51), the operator $`K_{b^{}a}^{a^{}b}`$ is non-trivial only on $`\varphi _{(L,m_2)}^{(L,m_1)}`$, while it vanishes on $`\varphi _{(1,m_2)}^{(L,m_1)}`$, $`\varphi _{(L,m_2)}^{(1,m_1)}`$, $`\varphi _{(1,m_2)}^{(1,m_1)}`$. From the composition rule of two spins, the scalar field $`\varphi _{(L,m_2)}^{(L,m_1)}`$ can be decomposed into the components with total spins $`J=0,1,\mathrm{},2L`$. The scalar field with a total spin $`J`$ has the components,
$$(\varphi _{J,m})_{(L,m_2)}^{(L,m_1)}=g_0^{(L,m_1)(L,m_1)}\left(\begin{array}{ccc}L& L& J\\ m_1& m_2& m\end{array}\right).$$
(57)
By using the following identities of the $`3j`$\- and $`6j`$-symbols,
$`{\displaystyle \underset{h,m_h}{}}(1)^{m_h}(2h+1)\left\{\begin{array}{ccc}j_1& J_1& h\\ J_2& j_2& f\end{array}\right\}\left(\begin{array}{ccc}j_1& J_1& h\\ m_1& M_1& m_h\end{array}\right)\left(\begin{array}{ccc}j_2& J_2& h\\ m_2& M_2& m_h\end{array}\right)`$ (58)
$`=(1)^{j_2+J_1}{\displaystyle \underset{m_f}{}}(1)^{m_f}\left(\begin{array}{ccc}j_1& j_2& f\\ m_1& m_2& m_f\end{array}\right)\left(\begin{array}{ccc}J_1& J_2& f\\ M_1& M_2& m_f\end{array}\right),`$ (59)
and
$$\underset{m_1,m_2}{}\left(\begin{array}{ccc}j_1& j_2& j_3\\ m_1& m_2& m_3\end{array}\right)\left(\begin{array}{ccc}j_1& j_2& j_3^{}\\ m_1& m_2& m_3^{}\end{array}\right)=\frac{1}{2j_3+1}\delta _{j_3j_3^{}}\delta _{m_3m_3^{}},$$
(60)
one can show that $`\varphi _{J,m}`$ is an eigen vector of the operator $`K_{b^{}a}^{a^{}b}`$,
$$K_{b^{}a}^{a^{}b}(\varphi _{J,m})_b^a=(1)^JA^2\left\{\begin{array}{ccc}L& L& J\\ L& L& 1\end{array}\right\}(\varphi _{J,m})_b^{}^a^{}=\frac{J(J+1)2L(L+1)}{2g(L^2+L1)}(\varphi _{J,m})_b^{}^a^{}.$$
(61)
Let me define the kinetic operator $`\stackrel{~}{K}`$ by
$$\stackrel{~}{K}_{b^{}a}^{a^{}b}=K_{b^{}a}^{a^{}b}+\frac{L(L+1)}{g(L^2+L1)}I_{b^{}a}^{a^{}b},$$
(62)
where $`I_{b^{}a}^{a^{}b}`$ denotes the identity operator, whose graphical representation is shown in Fig. 8.
The spectra of the operator $`\stackrel{~}{K}`$ in the $`\varphi _{(L,m_1)}^{(L,m_2)}`$ sector are given by $`J(J+1)/2g(L^2+L1)(J=0,1,\mathrm{},2L)`$. This is the spectra of a massless free scalar field on a fuzzy two-sphere of a size $`\sqrt{2g(L^2+L1)}`$. The fields in the other sectors $`\varphi _{(1,m_2)}^{(L,m_1)}`$, $`\varphi _{(L,m_2)}^{(1,m_1)}`$, $`\varphi _{(1,m_2)}^{(1,m_1)}`$ obtain a mass of approximately $`\sqrt{1/g}`$ on account of the identity operator in (62). These fields components can be physically decoupled by regarding $`g`$ to be small enough<sup>2</sup><sup>2</sup>2Since the model has no scales, only a relative scale is physically relevant..
## 6 Discussions and comments
In this paper, I studied models with a three-index variable and discussed whether they can be used as models for dynamical fuzzy spaces. The classical solutions of the models contain many solutions with Lie group symmetries. I considered the solutions with SO(3) symmetry as specific examples, and have constructed a scalar field theory on a fuzzy two-sphere. The construction in Section 4 and 5 would be straightforwardly applied to the general SO($`n`$) group to obtain the fuzzy complex quadratics discussed in .
The Lie group symmetries can be different from the orthogonal group, provided that the invariant tensors can be taken real. Moreover even in the SO(3) case, the model contains other solutions than what were used in the construction of a fuzzy-two sphere. It would be obviously interesting to find the interpretation of the other solutions as fuzzy spaces.
In the construction of a scalar field theory on a fuzzy two-sphere in Section 5, the scalar field was introduced as an additional degree of freedom, and the construction was rather adhoc. It is clear that the dynamics of $`C_{ab}^{}{}_{}{}^{c}`$ and $`g^{ab}`$ is more interesting. One could analyze their quadratic fluctuations around the classical solutions. Some of the fluctuations would be identified as scalar fields, and higher-spin fields would be also found. One might suspect that the models contain too many fields, but this might turn out to be nice, since it was recently argued that higher-spin fields must appear in gauge theory on non-associative fuzzy spaces .
Actions were not explicitly given, because the main interest of this paper was the classical solution. On the other hand, actions will be needed to perform the analysis in the preceding paragraph and also to study the quantum properties of the models. The quantum process could describe the transitions between distinct fuzzy spaces with different symmetries, and is worth to study. Considering a real action with complex variables would be also interesting, since the discussions on real variables in this paper can be essentially applied also to the complex case and the solutions can be constructed more freely.
The argument about the fuzzy general coordinate transformation has remained inconsistent. In the discussions about the fuzzy general coordinate transformation in Section 2, $`f_a`$ denote the functions on a fuzzy space. However, in the successful construction of the scalar field theory in Section 5, the scalar field has two indices and cannot be identified with $`f_a`$. Another formulation of field theory consistent with the argument or another interpretation of the fuzzy general coordinate transformation seems to be required.
It is shown in Section 3 that a lot of solutions can be constructed from the invariant tensors of Lie groups. A question is how many of the solutions have the Lie group symmetries. If most of the solutions do, it would be interesting to use the models as fuzzy higher dimensions - to explain the origin of the symmetries in our world. Incorporation of fermionic degrees of freedom and supersymmetry will be also interesting as phenomenology.
###### Acknowledgments.
The author was supported by the Grant-in-Aid for Scientific Research No.13135213 and No.16540244 from the Ministry of Education, Science, Sports and Culture of Japan.
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# 1 Introduction
## 1 Introduction
The Wilkinson Microwave Anisotropy Probe (WMAP) has placed stringent bounds on the amount of cold dark matter (CDM) in the universe. The amount of CDM deduced from WMAP data is given by
$`\mathrm{\Omega }_{CDM}h^2=0.1126_{0.009}^{+0.008},`$ (1)
where $`\mathrm{\Omega }_{CDM}=\rho _{CDM}/\rho _c`$, and where $`\rho _{CDM}`$ is the matter density of cold dark matter and $`\rho _c`$ is the critical matter density needed to close the universe, and $`h`$ is the Hubble parameter measured in units of 100km/s/Mpc. It is reasonable to assume that similar amounts of dark matter exist in our Milky Way and in the terrestrial neighborhood, and there are many ongoing experiments for the detection of such dark matter in the laboratory. On the theoretical side the WMAP data on cold dark matter puts stringent constraints on unified models of fundamental interactions since such models are called upon to predict or at least accommodate the WMAP data on CDM. As is well known, supergravity unified models with R-parity conservation allow for the possibility that the lightest neutralino may be the lightest supersymmetric particle (LSP) which could serve as a dark matter candidate <sup>2</sup><sup>2</sup>2There is also revived interest in the possibility that the LSP in SUGRA models could be the gravitino. For an update see Ref.. A hallmark of many unified models is Yukawa unification. In this paper we carry out a detailed investigation of the possibility of accommodating the WMAP cold dark matter data in the neutralino LSP scenario but under the constraints of Yukawa unification and including the effect of CP phases<sup>3</sup><sup>3</sup>3For an analysis of dark matter with CP phases but without inclusion of the Yukawa unification constraints see Refs.. An analysis of dark matter with quasi Yukawa unification was given in Refs... Another important constraint is the FCNC constraint given by $`bs+\gamma `$ which is discussed in some detail in this paper and included in the analysis.
Since the main focus of the analysis is the Yukawa unification constraint on dark matter in SUGRA models<sup>4</sup><sup>4</sup>4For recent works on dark matter analyses in SUGRA models see Ref., we briefly discuss some broad features of this constraint with details to follow later. In the supersymmetric framework the unification of the Yukawa couplings of the third generation, as predicted in several grand unification models, is rather sensitive to the parameters of the SUSY models. Thus, the compatibility of $`b\tau `$ unification at the grand unification scale with the observed $`b`$ and $`\tau `$ masses depends sensitively on the sign of $`\mu `$<sup>5</sup><sup>5</sup>5We use the sign convention on $`\mu `$ as in Ref.. (where $`\mu `$ is the Higgs mixing parameter) as well as on the details of the sparticle spectrum . Moreover, for most of the available parameter space b-$`\tau `$ unification is in conflict with other experimental constraints such as the FCNC process $`bs\gamma `$. The more stringent $`bt\tau `$ unification is predicted in the minimal $`SO(10)`$ models where the quarks and leptons, residing in the 16-plet spinor representation of $`SO(10)`$, gain masses via coupling with a 10-plet tensor representation of $`SO(10)`$<sup>6</sup><sup>6</sup>6More Higgs multiplets are needed to break the gauge symmetry correctly down to the standard model gauge symmetry, but typically these additional Higgs fields do not have couplings to quarks and leptons. . Finally, we mention that there can be GUT scale threshold corrections to the Yukawa unification. However, typically they are expected to be small .
Let us now be more specific and review the situation of Yukawa coupling unification in the mSUGRA case with no phases. With universal Yukawa couplings at the grand unification scale, the masses of the bottom and the top quark are naturally higher than the $`\tau `$ lepton mass. This phenomenon arises because of the color interactions which causes the Yukawa couplings of the quarks to increase as one goes down to lower energy scales. Thus, the running quark masses end up larger than the running charged lepton masses. To convert the running mass to the pole mass, one needs to include the supersymmetric as well as the standard model (SM) threshold corrections. In particular, it is well known that the supersymmetric threshold correction to the bottom quark mass, $`\mathrm{\Delta }m_b`$, can be very large. The value of $`\mathrm{\Delta }m_b`$ is enhanced for large values of $`\mathrm{tan}\beta `$, where $`\mathrm{tan}\beta `$ is the ratio $`H_u/H_d`$ and where $`H_u`$ gives mass to the up quark and $`H_d`$ gives mass to the down quark and the lepton. In the mSUGRA case with no CP phases, $`\mathrm{\Delta }m_b`$ takes the sign of $`\mu `$ (unless the trilinear terms are very large). A negative SUSY threshold correction to $`m_b`$ is required in models with $`b\tau `$ unification, in order to obtain a b-quark mass in the allowed range. Therefore, b-$`\tau `$ unification points toward a negative value of the $`\mu `$-parameter. However, a negative $`\mu `$ parameter makes the SUSY contribution to BR($`bs\gamma `$) positive and hence it adds to the SM contribution and the charged Higgs contribution. As a result, a heavy spectrum is required in order not to exceed the upper bound for this branching ratio.
The SUSY contribution to the muon anomalous magnetic moment also takes the sign of $`\mu `$ in mSUGRA , and more generally this contribution depends on CP phases . However, experimentally the situation is less clear regarding the implications of the $`g_\mu 2`$ data. Thus, while the BNL experiment has significantly improved the accuracy of the $`g_\mu 2`$ measurement , ambiguities in the hadronic error, which is needed to compute the deviation of the observed value from the Standard Model prediction, still persist. Currently, the largest source of error in the computation of the Standard Model prediction is the O($`\alpha ^2`$) hadronic vacuum polarization correction. The most recent evaluation of this correction are done by (i) Davier et.al. using the $`\tau `$ decay data, and by (ii) Hagiwara et.al. using the low energy data from $`e^+e^{}`$ hadrons. Assuming that the entire difference $`\mathrm{\Delta }a_\mu `$ (where $`a_\mu `$ is defined so that the effective operator is $`a_\mu (e/2m_\mu )\overline{\mu }\sigma _{\alpha \beta }\mu F^{\alpha \beta }`$) between experiment and theory comes from supersymmetry, one finds that the supersymmetric contribution for the case of Davier et.al. is $`(7.6\pm 9.0)\times 10^{10}`$, while for the case of Hagiwara et.al., the value is $`(23.9\pm 10.0)\times 10^{10}`$. In this analysis we adopt solution (i). In fact, in most of the parameter space we explore the sparticle spectrum is rather heavy, and the SUSY contribution $`\mathrm{\Delta }a_\mu `$ small, and thus the $`a_\mu `$ prediction is essentially the same as the Standard Model prediction which is consistent with the current data. Solution (ii) puts a more stringent constraint. However, the constraint can be softened if the universality condition on the soft terms is removed .
As indicated earlier this paper is devoted mostly to an analysis of the WMAP data with the $`b\tau `$ unification constraint. However, we will also briefly discuss b-$`\tau `$-t unification. As is well known such a unification requires large $`\mathrm{tan}\beta `$ and for this reason much of the parameter space is excluded since it does not correctly break the electro-weak symmetry. Several studies has been done, and models such as, e.g., the D-term splitting in SO(10) or non-universal Higgs masses can indeed give rise to a viable t-b-$`\tau `$ unification . However, typically these solutions require a very heavy SUSY spectrum. Thus, the predicted dark matter abundance of neutralinos, in models with R-parity conservation, will be too high and thus will over-close the universe. As mentioned earlier, models with quasi unification have also been investigated .
In the present work we first analyze the relic density within mSUGRA and show that there exist regions of the parameter space where the WMAP relic density constraint, the Yukawa unification constraint, and the BR($`bs\gamma `$) constraint can all be simultaneously satisfied. We then extend the mSUGRA parameter space retaining universality on the magnitude of the soft parameters but allowing non-universality for the phases in some sectors. The SUSY contribution to $`\mathrm{\Delta }m_b`$ is phase–dependent and this allows one to determine the phases in some cases such as to obtain $`m_b(M_Z)`$ in the experimental range and thus achieve b-$`\tau `$ unification. Indeed, one finds that with the inclusion of phases $`b\tau `$ unification is achievable in a large area of the parameter-space. Further, it is possible to find arrangement of phases such that the prediction of the electric dipole moments (EDMs) is in agreement with the experimental bounds. The case of full Yukawa unification is, however, still (almost) incompatible with the experimental value for BR($`bs\gamma `$) . However, it is worth keeping in mind that small flavor mixings in the sfermion mass matrices can substantially change the predictions for BR($`bs\gamma `$) while leaving other predictions essentially untouched. Thus, in principle, non-zero CP-phases could also allow viable b-$`\tau `$-t unification modulo mixings in the squark flavor sector. However, we do not pursue this line of investigation in the work here.
The outline of the rest of the paper is as follows: In Sec.(2) we give a discussion of the parameter space of the model and the details of the procedure of the calculations. In Sec.(3) we discuss the calculation of the BR($`bs\gamma `$) and resolve some of the ambiguities present in the literature in the large $`\mathrm{tan}\beta `$ enhanced contributions, by carrying out an independent analysis of the parameters $`ϵ_b^{^{}}(t),ϵ_t^{^{}}(s),ϵ_{bb}`$, which codify these contributions. In Sec.(4) we carry out an analysis of the relic density with the $`b\tau `$ unification constraint, and with the BR($`bs\gamma `$) constraint both within mSUGRA and in extended SUGRA models with phases. In Sec.(5) we give an analysis of the relic density for the case of the full Yukawa unification. In Sec.(6) we discuss the consistency of the analysis with large phases with the EDM constraints and give examples of models with large phases consistent with the WMAP data, b-$`\tau `$ unification and with the EDM constraints. Conclusions are given in Sec.(7).
## 2 Constraints on SUGRA models with CP-phases
Within mSUGRA there are only two physical phases, which can be chosen as $`\theta _\mu `$, the phase of the Higgs mixing parameter $`\mu `$, and $`\alpha _0`$ the phase of the universal trilinear term $`A_0`$. These phases are severely constrained by the non-observation of the electric dipole moments (EDM). The present upper bounds for the EDM of the electron, of the neutron, and of the mercury <sup>199</sup>Hg atom are
$$|d_e|<4.23\times 10^{27}\mathrm{e}\mathrm{cm},|d_n|<6.5\times 10^{26}\mathrm{e}\mathrm{cm},C_{\mathrm{Hg}}<3.0\times 10^{26}\mathrm{cm},$$
(2)
where $`C_{\mathrm{Hg}}`$ is defined as in Ref.. Large phases can be accommodated in several scenarios such as models with heavy sfermions , models with the cancellation mechanism , models with phases only in the third generation , or models with a non-trivial soft flavor structure . Here, we use the cancellation mechanism <sup>7</sup><sup>7</sup>7For a more complete list of references and for a discussion of the effects of CP phases on low energy processes see Ref.. which becomes possible if the SUGRA parameter space is extended to allow for different gaugino phases. The model we consider is thus described by the following parameters
$$m_0,m_{1/2},\mathrm{tan}\beta ,|A_0|,\theta _\mu ,\alpha _0,\xi _1,\xi _2,\xi _3,$$
(3)
where, $`\xi _i`$ is the phase of the gaugino mass $`M_i,i=1,2,3`$. The value of $`|\mu |`$ is determined by imposing electroweak symmetry breaking (EWSB).
In the analysis we use a top-down approach, and thus impose Yukawa unification at the GUT scale, $`M_{\mathrm{GUT}}`$. For $`b\tau `$ unification we have two independent Yukawa couplings at the grand unification scale, i.e., one common $`h_{\mathrm{uni}}`$ for the b and the $`\tau `$, and one for the top-quark. We use these to fit the experimental value of the $`\tau `$ and the top masses. Unless another value is specified, we fix the top mass at 178 GeV, which is its current experimental central value . The value of $`\alpha _s`$ is fixed to be 0.1185. For the $`\tau `$ mass at the electroweak scale $`M_Z`$, we use 1.7463 GeV, which takes into account the Standard Model radiative correction. Naturally, we also take into account the SUSY correction, as derived in , when calculating $`m_\tau `$. In the case of the full Yukawa unification we impose $`h_b=h_t=h_\tau =h_{\mathrm{uni}}`$ at $`M_{\mathrm{GUT}}`$. Therefore, the value of $`\mathrm{tan}\beta `$ is fixed, since the two parameters $`h_{\mathrm{uni}}`$ and $`\mathrm{tan}\beta `$ are varied so as to obtain agreement with experimental values of $`m_\tau `$ and $`m_{\mathrm{top}}`$. As the b-quark couples to the same Higgs doublet ($`H_d`$) as the $`\tau `$ lepton, its mass is fixed by $`h_{\mathrm{uni}}`$. Therefore, $`m_b(M_Z)`$ is a prediction of our model and we require its value to be within the 2$`\sigma `$ range,
$$2.69\mathrm{GeV}<m_b(M_Z)<3.10\mathrm{GeV}$$
(4)
as described in . In addition to the above, the other important constraints of the analysis are the relic density and the BR($`bs\gamma `$) constraint (see section 3).
The procedure for the calculation of the particle and sparticle masses is as follows; After choosing a given set of the parameters in Eq.(3), we run the renormalization group equations (RGEs) down to the SUSY scale, defined as the average of the two stop masses. At the SUSY scale the scalar potential is minimized and $`|\mu |`$ is calculated along with the SUSY threshold corrections to, e.g., the b quark and the $`\tau `$ lepton masses and the couplings are corrected accordingly. Hereafter, the sparticles are decoupled and the SM RGEs are used to run down to $`M_Z`$. At the electro-weak scale we check if the gauge couplings, the Weinberg angle, the top quark and the $`\tau `$ lepton masses are in agreement with their experimental values. If not, the RGEs are run iteratively until convergence is achieved. In the analysis we use the two-loop SUSY renormalization group equations except for the trilinear terms, the gaugino and sfermion masses, which are calculated at the one-loop level. The SUSY renormalization group equation will also be influenced by the CP-phases. However, it is easy to see that neither the phase of the $`\mu `$-term nor the phases of the gaugino masses will run. But, the phases of the trilinear terms run, and in general there will be three different phases at the low energy scale namely, $`\alpha _t`$, $`\alpha _b`$ and $`\alpha _\tau `$. $`\alpha _t`$ is important as it affects $`\mathrm{\Delta }m_b`$ as well as BR($`bs\gamma `$) . However, its value is almost fixed by the gluino phase. As shown in Ref., the approximate relation $`A_{\mathrm{top}}M_3`$ holds at low energy.
The regions of the mSUGRA parameter space that allows for acceptable relic abundance can be classified as: (i) the $`\chi \stackrel{~}{\tau }`$ coannihilation region, (ii) the resonance region, and (iii) the Hyperbolic Branch/Focus Point (HB/FP) region. In a previous work , we pointed out the strong variation of $`\mathrm{\Delta }m_b`$ with CP phases. In that work we focussed on the effects induced by the SUSY corrections on the spectrum and their consequences for the neutralino relic density. It was shown that the CP-phases have a very large impact on the value of the CP-odd Higgs mass $`M_A`$, which in turn affects the predicted dark matter abundance in the so-called resonance region. The analysis of Ref. used a bottom-up approach by fixing the value of $`m_b(M_Z)`$ to its central value. In this work we use a top-down approach and large effects of the CP phases are not seen. In fact, the predicted neutralino relic abundance, turns out almost independent of the phases in the resonance region. In the stau coannihilation region there is also very little dependence on the CP-phases, except for the trilinear phase. As we show below, the HB/FP region cannot give rise to Yukawa unification within our model. In the calculation of the relic density we take into account the CP even-CP odd Higgs mixing. In the MSSM, after spontaneous breaking of the electro-weak symmetry one has at the tree-level two CP even Higgs ($`h^0,H^0`$) and one CP odd Higgs ($`A`$). In the presence of CP violating phases these mix, producing mass eigenstates ($`H_1^0,H_2^0,H_3^0`$), which are no longer eigen-functions of CP <sup>8</sup><sup>8</sup>8For further details regarding the implications of these CP even-CP odd Higgs mixings on neutralino dark matter analysis see Ref..
The most important supersymmetric threshold correction is the one to the bottom mass. At the loop level the effective b quark coupling with the Higgs is given by
$$_{bbH^0}=(h_b+\delta h_b)\overline{b}_Rb_LH_1^0+\mathrm{\Delta }h_b\overline{b}_Rb_LH_2^0+H.c.$$
(5)
The correction to the b quark mass is then given directly in terms of $`\mathrm{\Delta }h_b`$ and $`\delta h_b`$ by
$$\mathrm{\Delta }m_b=[Re(\frac{\mathrm{\Delta }h_b}{h_b})\mathrm{tan}\beta +Re(\frac{\delta h_b}{h_b})].$$
(6)
We use the full analysis of $`\mathrm{\Delta }m_b`$ derived in . The largest contributions to $`\mathrm{\Delta }m_b`$ are the gluino and the chargino exchange contributions. The gluino exchange contribution is proportional to $`M_3\mu `$, and will therefore depend on the phase combination $`\theta _\mu +\xi _3`$. The chargino exchange contribution is usually smaller, except for very large values of $`|A_t|`$, since it is proportional to $`A_t\mu `$. Its dominant phase dependence is given by $`\theta _\mu +\alpha _t`$, and it has the opposite sign of the gluino contribution in a large region of the parameter space. When evaluating $`h_b`$ at $`M_{\mathrm{SUSY}}`$, we take into account threshold corrections using the relation <sup>9</sup><sup>9</sup>9This relation resums the SUSY self-energy leading order logarithmic corrections .
$$h_b^{\mathrm{SM}}=h_b^{\mathrm{SUSY}}(1+\mathrm{\Delta }m_b).$$
(7)
The SM Yukawa coupling is evolved down to the electroweak scale, and the bottom quark mass
$$m_b(M_Z)=h_b^{\mathrm{SM}}\frac{v}{\sqrt{2}}\mathrm{cos}\beta ,$$
(8)
is calculated and compared with experiment. Similar expressions hold for the $`\tau `$ lepton with b replaced by $`\tau `$. For the top quark at the Z scale one has
$`m_t(M_Z)={\displaystyle \frac{v}{\sqrt{2}}}\mathrm{sin}\beta h_t^{\mathrm{SUSY}}(1+\mathrm{\Delta }m_t)`$ (9)
where
$$\mathrm{\Delta }m_t=[Re(\frac{\mathrm{\Delta }h_t}{h_t})\mathrm{cot}\beta +Re(\frac{\delta h_t}{h_t})].$$
(10)
A full analysis of $`\mathrm{\Delta }m_t`$ is given in Ref. . However, in the region of interest which corresponds to large $`\mathrm{tan}\beta `$ the correction to the top quark Yukawa is essentially negligible.
## 3 BR($`bs\gamma `$) with CP–phases
The present average for the BR($`bs\gamma `$) derived from the available experimental data is found to be,
$$\mathrm{BR}(\mathrm{b}\mathrm{s}\gamma )=(3.54_{0.28}^{+0.30})\times 10^4,$$
(11)
by the Heavy Flavor Averaging Group . The error includes an uncertainty due to the decay spectrum as well as the statistical error. The theoretical SM prediction is ,
$$\mathrm{BR}(\mathrm{b}\mathrm{s}\gamma )=(3.70\pm 0.30)\times 10^4.$$
(12)
The above result uses the $`\overline{MS}`$ running charm mass instead of the pole mass. It was claimed in Ref. that this consideration reduces the NNLO uncertainty in the SM calculation. However, other analyses question the theoretical precision of Eq.(12) predicting a lower central value for the SM result. In any case, the result of Eq.(12) appears to be a good benchmark value for the SM prediction to work with.
The dominant SUSY contributions from the charged Higgs exchange include the $`\mathrm{tan}\beta `$ enhanced NLO corrections, which contribute to the Wilson coefficients $`C_7`$ and $`C_8`$ (these are coefficients of the operators $`O_7=\frac{e}{16\pi ^2}m_b(\overline{s}_L\sigma _{\mu \nu }b_R)F_{\mu \nu }`$ and $`O_8=\frac{g_s}{16\pi ^2}m_b(\overline{s}_L\sigma _{\mu \nu }T^ab_R)G_{\mu \nu }^a`$). These contributions can be codified in $`ϵ_b^{^{}}(t)`$, $`ϵ_t^{^{}}(b)`$ and $`ϵ_{bb}`$ which enter in the Lagrangian for effective interaction involving the charged Goldstone boson and the charged Higgs boson as follows
$``$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{2}M_W}}G^+\{{\displaystyle \underset{d}{}}m_tV_{td}\overline{t}_Rd_L{\displaystyle \underset{u}{}}m_bV_{ub}{\displaystyle \frac{1+ϵ_b^{}(u)\mathrm{tan}\beta }{1+ϵ_{bb}^{}\mathrm{tan}\beta }}\overline{u}_Lb_R\}`$
$`+`$ $`{\displaystyle \frac{g}{\sqrt{2}M_W}}H^+\{{\displaystyle \underset{d}{}}m_tV_{td}\overline{t}_Rd_L{\displaystyle \frac{1+ϵ_t^{}(d)\mathrm{tan}\beta }{\mathrm{tan}\beta }}+{\displaystyle \underset{u}{}}m_bV_{ub}\overline{u}_Lb_R{\displaystyle \frac{\mathrm{tan}\beta }{1+ϵ_{bb}^{}\mathrm{tan}\beta }}\}+H.c.`$
where $`V_{ij}`$ is the CKM mixing matrix. Evaluation of $`ϵ_b^{^{}}(t)`$, $`ϵ_t^{^{}}(b)`$, and $`ϵ_{bb}`$ exist in the literature , but there is some ambiguity concerning the signs of some of the terms among the above groups. To resolve this we carry out an independent analysis of these quantities for the same loop diagrams as in the previous works, including also their dependence on CP phases, which was taken into account only in one analysis previously. Our analysis is derived using the work of Ref.. We find
$`ϵ_b^{^{}}(t)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \frac{2\alpha _s}{3\pi }}e^{i\xi _3}D_{b2j}^{}D_{t1i}[{\displaystyle \frac{m_t}{m_b}}\mathrm{cot}\beta A_tD_{b1j}D_{t2i}^{}+\mu D_{b2j}D_{t1i}^{}+m_t\mathrm{cot}\beta D_{b2j}D_{t2i}^{}`$
$`+{\displaystyle \frac{m_t^2}{m_b}}\mathrm{cot}\beta D_{b1j}D_{t1i}^{}{\displaystyle \frac{m_W^2}{m_b}}\mathrm{sin}\beta \mathrm{cos}\beta D_{b1j}D_{t1i}^{}]{\displaystyle \frac{1}{|m_{\stackrel{~}{g}}|}}H({\displaystyle \frac{m_{\stackrel{~}{t}_i^2}}{|m_{\stackrel{~}{g}}|^2}},{\displaystyle \frac{m_{\stackrel{~}{b}_j^2}}{|m_{\stackrel{~}{g}}|^2}})`$
$`+2{\displaystyle \underset{k=1}{\overset{4}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}[{\displaystyle \frac{m_t}{m_b}}\mathrm{cot}\beta A_tD_{b1j}D_{t2i}^{}+\mu D_{b2j}D_{t1i}^{}+m_t\mathrm{cot}\beta D_{b2j}D_{t2i}^{}`$
$`+{\displaystyle \frac{m_t^2}{m_b}}\mathrm{cot}\beta D_{b1j}D_{t1i}^{}{\displaystyle \frac{m_W^2}{m_b}}\mathrm{sin}\beta \mathrm{cos}\beta D_{b1j}D_{t1i}^{}]`$
$`\times (\alpha _{bk}^{}D_{b1j}^{}\gamma _{bk}^{}D_{b2j}^{})(\beta _{tk}D_{t1i}+\alpha _{tk}^{}D_{t2i}){\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \frac{1}{m_{\chi _k^0}}}H({\displaystyle \frac{m_{\stackrel{~}{t}_i^2}}{m_{\chi _k^0}^2}},{\displaystyle \frac{m_{\stackrel{~}{b}_j^2}}{m_{\chi _k^0}^2}})`$
In the above $`D_q`$ is the matrix that diagonalizes the squark mass<sup>2</sup> matrix $`M_{\stackrel{~}{q}}^2`$, i.e.,
$`D_q^{}M_{\stackrel{~}{q}}^2D_q=\mathrm{diag}(M_{\stackrel{~}{q}_1}^2,M_{\stackrel{~}{q}_2}^2)`$ (15)
and $`H(a,b)`$ is defined by
$`H(a,b)={\displaystyle \frac{a}{(1a)(ab)}}\mathrm{ln}a+{\displaystyle \frac{b}{(1b)(ba)}}\mathrm{ln}b`$ (16)
where $`\alpha _{bk},\beta _{bk},\gamma _{bk}`$ for the b quark and the corresponding coefficients for the t quark are as defined in Ref.. Similarly for $`ϵ_t^{^{}}(s)`$ we find
$`ϵ_t^{^{}}(s)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \frac{2\alpha _s}{3\pi }}e^{i\xi _3}D_{s1i}^{}D_{t2j}[{\displaystyle \frac{m_s}{m_t}}\mathrm{tan}\beta A_s^{}D_{s2i}D_{t1j}^{}+\mu ^{}D_{s1i}D_{t2j}^{}+m_s\mathrm{tan}\beta D_{s2i}D_{t2j}^{}`$
$`+{\displaystyle \frac{m_s^2}{m_t}}\mathrm{tan}\beta D_{s1i}D_{t1j}^{}{\displaystyle \frac{m_W^2}{m_t}}\mathrm{sin}\beta \mathrm{cos}\beta D_{s1i}D_{t1j}^{}]{\displaystyle \frac{1}{|m_{\stackrel{~}{g}}|}}H({\displaystyle \frac{m_{\stackrel{~}{s}_i}^2}{|m_{\stackrel{~}{g}}|^2}},{\displaystyle \frac{m_{\stackrel{~}{t}_j}^2}{|m_{\stackrel{~}{g}}|^2}})`$
$`2{\displaystyle \underset{k=1}{\overset{4}{}}}{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}[{\displaystyle \frac{m_s}{m_t}}\mathrm{tan}\beta A_s^{}D_{s2i}D_{t1j}^{}+\mu ^{}D_{s1i}D_{t2j}^{}+m_s\mathrm{tan}\beta D_{s2i}D_{t2j}^{}`$
$`+{\displaystyle \frac{m_s^2}{m_t}}\mathrm{tan}\beta D_{s1i}D_{t1j}^{}{\displaystyle \frac{m_W^2}{m_t}}\mathrm{sin}\beta \mathrm{cos}\beta D_{s1i}D_{t1j}^{}]`$
$`\times (\beta _{sk}^{}D_{s1i}^{}+\alpha _{sk}D_{s2i}^{})(\alpha _{tk}D_{t1j}\gamma _{tk}D_{t2j}){\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \frac{1}{m_{\chi _k^0}}}H({\displaystyle \frac{m_{\stackrel{~}{s}_i}^2}{m_{\chi _k^0}^2}},{\displaystyle \frac{m_{\stackrel{~}{t}_j}^2}{m_{\chi _k^0}^2}})`$
Finally, our analysis of $`ϵ_{bb}`$ gives
$`ϵ_{bb}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \frac{2\alpha _s}{3\pi }}e^{i\xi _3}D_{b1i}^{}D_{b2j}[{\displaystyle \frac{M_Zm_W}{m_b}}{\displaystyle \frac{\mathrm{cos}\beta }{\mathrm{cos}\theta _W}}\{({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{3}}\mathrm{sin}^2\theta _W)D_{b1i}D_{b1j}^{}`$ (19)
$`{\displaystyle \frac{1}{3}}\mathrm{sin}^2\theta _WD_{b2i}D_{b2j}^{}\}\mathrm{sin}\beta +\mu ^{}D_{b1i}D_{b2j}^{}]{\displaystyle \frac{1}{|m_{\stackrel{~}{g}}|}}H({\displaystyle \frac{m_{\stackrel{~}{b}_i}^2}{|m_{\stackrel{~}{g}}|^2}},{\displaystyle \frac{m_{\stackrel{~}{b}_j}^2}{|m_{\stackrel{~}{g}}|^2}})`$
$`{\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \underset{k=1}{\overset{2}{}}}g^2[{\displaystyle \frac{M_Zm_W}{m_b}}{\displaystyle \frac{\mathrm{cos}\beta }{\mathrm{cos}\theta _W}}\{({\displaystyle \frac{1}{2}}{\displaystyle \frac{2}{3}}\mathrm{sin}^2\theta _W)D_{t1i}D_{t1j}^{}+{\displaystyle \frac{2}{3}}\mathrm{sin}^2\theta _WD_{t2i}D_{t2j}^{}\}\mathrm{sin}\beta `$
$`{\displaystyle \frac{m_t^2}{m_b}}\mathrm{cot}\beta \{D_{t1i}D_{t1j}^{}+D_{t2i}D_{t2j}^{}\}{\displaystyle \frac{m_t}{m_b}}\mathrm{cot}\beta A_t^{}D_{t2i}D_{t1j}^{}]`$
$`\times (V_{k1}^{}D_{t1i}^{}K_tV_{k2}^{}D_{t2i}^{})(K_bU_{k2}^{}D_{t1j}){\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \frac{1}{|m_{\stackrel{~}{\chi }_k^+}|}}H({\displaystyle \frac{m_{\stackrel{~}{t}_i}^2}{|m_{\stackrel{~}{\chi }_k^+}|^2}},{\displaystyle \frac{m_{\stackrel{~}{t}_j}^2}{|m_{\stackrel{~}{\chi }_k^+}|^2}})`$
The form factor $`H(a,b)`$ in the above equation can have $`a=b`$ and in this case it reads
$`H(a,a)={\displaystyle \frac{1}{(a1)^2}}[1a+\mathrm{ln}a]`$ (20)
Before proceeding further we give a brief comparison of these results with the results of the previous works. The analysis of $`ϵ_b^{^{}}(t)`$ may be compared to $`ϵ_{tb}`$ of Ref. in the limit of large $`\mathrm{tan}\beta `$ and small squark mixings. In this case the limit of the first two lines in Eq.(3) agrees with the result of Ref.. However, the limit of the last three lines of Eq.(3) have an opposite sign to that of Ref.. Here our analysis is in agreement with the result of Ref..
Next we give a computation of $`ϵ_t^{^{}}(s)`$. Approximating Eq.(3) we find
$`ϵ_t^{^{}}(s)={\displaystyle \underset{i=1}{\overset{2}{}}}{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \frac{2\alpha _s}{3\pi }}e^{i\xi _3}\mu ^{}|D_{s1i}|^2|D_{t2j}|^2{\displaystyle \frac{1}{|m_{\stackrel{~}{g}}|}}H({\displaystyle \frac{m_{\stackrel{~}{s}_i^2}}{|m_{\stackrel{~}{g}}|^2}},{\displaystyle \frac{m_{\stackrel{~}{t}_j^2}}{|m_{\stackrel{~}{g}}|^2}})`$
$`{\displaystyle \frac{h_s^2}{16\pi ^2}}{\displaystyle \frac{A_s^{}}{m_{\chi _k^0}}}X_{3k}X_{4k}|D_{s2i}|^2|D_{t1j}|^2H({\displaystyle \frac{m_{\stackrel{~}{s}_i^2}}{|m_{\stackrel{~}{\chi }_k^+}|^2}},{\displaystyle \frac{m_{\stackrel{~}{t}_j}^2}{|m_{\stackrel{~}{\chi }_k^+}|^2}})`$ (21)
The analysis of Ref. computed only the first line of Eq.(21) and for this part we agree with their work when we take the large $`\mathrm{tan}\beta `$ limit and the limit of small mixing angles of our result. The work Ref. gives results corresponding to Eq.(21). However, here we find that we have a disagreement with the sign of the second part of their Eq.(16).
Our analysis of $`ϵ_{bb}`$ given by our Eq.(19) agrees with the analysis of Ref. in the limit of large $`\mathrm{tan}\beta `$ and in the limit of small squark mixings and here there is a general agreement (taking account of typo corrections) among various groups in the limit of no CP phases. Our analysis like that of Ref. takes into account the full dependence on CP phases. In the numerical analysis to be presented below, we have used the code provided by micrOMEGAs in the CP conserving case. This code agrees with the codes used by other groups . In the CP violating case we have combined the codes of Refs. with our own codes of the SUSY contributions with CP phases.
For the uncertainty in BR($`bs\gamma `$) we use a linear combination of the errors on Eqs. (11) and (12). At the 2-$`\sigma `$ level one has
$$2.3\times 10^4<\mathrm{BR}(\mathrm{b}\mathrm{s}\gamma )<4.7\times 10^4.$$
(22)
The numerical analysis given below is controlled essentially by the upper bound in Eq.(22). In order to obtain the correct value of $`m_b(m_Z)`$ one needs the phase combination $`\theta _\mu +\xi _3`$ to be close to $`\pi `$<sup>10</sup><sup>10</sup>10The phase combination is drawn to smaller values for large $`\mathrm{tan}\beta `$ and for the full Yukawa unification it ends up close to $`\pi /2`$.. In this case the chargino contribution to the BR($`bs\gamma `$) is positive and therefore the lower bound is not reached for the values of the SUSY parameters in our study. We turn now to the details of the numerical analysis.
## 4 WMAP Dark Matter Constraints and $`b\tau `$ Yukawa unification
It is useful to first summarize our results in the mSUGRA case, where all CP phases are either zero or $`\pi `$. The relation $`h_b=h_\tau `$ can be satisfied for a wide range of soft masses in the MSSM with real universal soft terms. To discuss the dependence of $`m_b(M_Z)`$ on $`\mathrm{tan}\beta `$ we consider two representative set of soft parameters: (i) $`m_{\frac{1}{2}}=800`$ GeV, $`A_0=0`$, $`m_0=300`$ GeV, and (ii) $`m_{\frac{1}{2}}=800`$ GeV, $`A_0=0`$, $`m_0=600`$ GeV. In Fig. (1) we study the $`\mu >0`$ and $`\mu <0`$ cases for each set. The lines corresponding to case (i) are interrupted when $`m_{\stackrel{~}{\tau }}<m_\chi `$, while for case (ii) large values of $`\mathrm{tan}\beta `$ are incompatible with EWSB. We include a reference line ignoring the SUSY threshold corrections (i.e., $`\mathrm{\Delta }m_b=0`$). Fig. (1) exhibits the well known phenomenon, that $`\mathrm{\Delta }m_b`$ is positive for $`\mu `$ positive and therefore the theoretical prediction for the b quark pole mass is too high, lying outside the experimental range. Thus b-$`\tau `$ unification does not occur in this case. When $`\mu <0`$, on the other hand, $`\mathrm{\Delta }m_b`$ is negative and the theoretical prediction for the b quark mass can lie within the experimental range for values of $`\mathrm{tan}\beta `$ between roughly 25 and 45. A similar analysis of $`m_b(M_Z)`$ but as a function of $`m_0`$ is given in Fig.(2). Here we consider only the $`\mu <0`$ case and find that the theoretical prediction of $`m_b(M_Z)`$ can lie within the corridor allowed by experiment for a range of $`m_0`$ values. However, one finds that for very high values of $`m_0`$, i.e., for values above 5 TeV and beyond, a region which includes the Hyperbolic Branch (HB)/Focus Point (FP) region, Yukawa unification is not achieved with universal soft parameters.
We extend the analysis now to include the relic density constraints. Fig.(3) shows an area plot in the $`m_0m_{\frac{1}{2}}`$ plane and all the three interesting dark matter regions in mSUGRA can be seen. In Fig.(3) the coannihilation area and the resonance area overlap. The HB/FP area, which is the region adjacent to the area with no EWSB ($`\mu ^2<0`$), is incompatible with any kind of Yukawa unification within the framework of universality of soft parameters at the GUT scale as already seen in the analysis of Fig.(2).
The HB/FP region moves to lower values of $`m_0`$ as $`tan\beta `$ and $`m_t`$ decrease and this variation, especially with $`m_t`$, can be very large (we present our analysis in Figs.(2,3) using $`m_t=176`$ GeV in order that the HB/FP area appear below 20 TeV). With $`m_t`$ at its lower bound the HB/FP region appears at values of $`m_0`$ of 5-6 TeV, but here $`m_b`$ is already too high. With large values of $`tan\beta `$, $`m_A^2`$ becomes negative before $`\mu `$ becomes small and therefore there is no inversion of the gaugino/Higgsino components in the composition of $`\chi ^0`$. Correspondingly, the HB/FB is not reached. This is the case for the line of $`\mathrm{tan}\beta =48`$ in Fig.2. Therefore, overlapping of the HB/FP region and the allowed $`m_b`$ area is not possible. Moreover, phases cannot improve the situation as $`\mathrm{\Delta }m_b`$ is very small in the HB/FP (below 5%), and thus cannot lower the value of $`m_b`$ sufficiently.
In Fig. (4) we further analyze the mSUGRA case with area plots in the $`m_0m_{\frac{1}{2}}`$ plane for four values of $`\mathrm{tan}\beta `$: 30, 35, 40, 45. Fig. (4) shows that the relic abundance, $`b\tau `$ unification, and $`bs\gamma `$ constraints can be simultaneously satisfied for a narrow range of parameters for values of $`\mathrm{tan}\beta `$ in the range 30–45. The BR($`bs\gamma `$) constraint is a major restriction, since both SUSY and Higgs contributions to the branching ratio add to the one from the Standard Model, and thus one needs a relatively heavy spectrum such that the BR($`bs\gamma `$) prediction remains below the experimental upper bound. The b-$`\tau `$ unification constraint and the WMAP constraint further reduce the parameter space. Even so, one finds that there exist regions of the parameter space for all the four cases in Fig. (4) consistent with the WMAP data under the $`b\tau `$ unification and $`bs\gamma `$ constraints.
### 4.1 Effects of CP phases
To determine the impact of phases on the above picture we choose two representative points from Fig. 4:
$$a)tan\beta =30,m_0=290\mathrm{GeV},\mathrm{M}_{1/2}=800\mathrm{GeV},\mathrm{A}_0=0\mathrm{GeV}$$
(23)
$$b)tan\beta =40,m_0=710\mathrm{GeV},\mathrm{M}_{1/2}=800\mathrm{GeV},\mathrm{A}_0=0\mathrm{GeV}$$
(24)
These points are chosen because in the absence of phases the WMAP relic density constraints are satisfied by different mechanisms for these two cases. Thus, for the point in Eq.(23) the WMAP constraint is satisfied due to $`\chi \stackrel{~}{\tau }`$ coannihilations. In contrast, for the point in Eq.(24) the WMAP constraint is satisfied due to a resonance in the Higgs mediated annihilation of $`\chi \chi `$. To determine the effect of phases we study the most relevant phases for the processes that we consider. The phases $`\xi _1`$ and $`\xi _2`$ have little impact on $`bs\gamma `$ and $`\mathrm{\Delta }m_b`$. For simplicity we set them to zero in this section. We have already discussed the phase combinations that play an important role in the analysis of $`\mathrm{\Delta }m_b`$. For the analysis of $`bs\gamma `$ we find that the same phase combinations, i.e, $`Arg(\mu A_t)`$ and $`Arg(\mu M_3)`$ are the important ones.
We now discuss the specifics of the point in Eq. (23), which as already stated is in the $`\chi \stackrel{~}{\tau }`$ coannihilation region. In Fig.(5) we analyze the BR($`bs\gamma `$) and the $`b\tau `$ unification constraints in the $`\theta _\mu \xi _3`$ plane, and as is seen the point satisfies the $`bs\gamma `$ as well as $`m_b(M_Z)`$ constraint in the mSUGRA case with a negative $`\mu `$. The figure illustrates that the inclusion of phases changes the value of $`m_b(M_Z)`$ and BR($`bs\gamma `$) drastically. However, the two above mentioned constraints have a tendency to conflict with each other, even with the inclusion of CP phases. Nevertheless, we find that there exists a substantial overlap of the areas allowed by the bounds on $`m_b(M_Z)`$ and BR($`bs\gamma `$) . At the same time the predicted $`\mathrm{\Omega }_{CDM}h^2`$ remains inside the WMAP bounds because the phases do not affect significantly the ratio $`m_\chi /m_{\stackrel{~}{\tau }}`$ and hence the relic density prediction remains dominated by coannihilations. It is also instructive to study the effects of $`\alpha _0`$. In Fig. (6) we analyze the dependence on $`|A_0|`$ and $`\alpha _0`$ for the point $`\xi _3=0.3`$ rad and $`\theta _\mu =2.4`$ rad of Fig.(5). In the analysis of Fig. (6) the ratio $`m_{\stackrel{~}{\tau }}/m_\chi `$ does not exceed 1.08, and thus we remain in the coannihilation region allowing for the satisfaction of the relic density constraints. Furthermore, it is also possible to satisfy BR($`bs\gamma `$) and $`m_b`$ bounds simultaneously for a wide range of $`|A_0|`$ and $`\alpha _0`$.
Next we analyze the implication of phases for the point in Eq. (24). As already indicated this point is within the resonance region in the mSUGRA case. However, the point produces a value for the BR($`bs\gamma `$) outside the experimental bounds as may be seen from Fig. (4). The effect of varying $`\xi _3`$ and $`\theta _\mu `$ is analyzed in Fig. (7). Here one finds a substantial overlap of the areas allowed by the bounds on $`m_b(M_Z)`$ and BR($`bs\gamma `$) while the relic density prediction remains within the WMAP bounds. As already stated, this analysis is substantially different from the one given in Ref. at $`\mathrm{tan}\beta =40`$. There $`m_b(M_Z)`$ was fixed and the dependence of $`\mathrm{\Delta }m_b`$ on the phases has a big effect on the resonant channels. For the present case, the bottom Yukawa has only a small fluctuation due to the unification condition at the GUT scale. Thus its effect on the Higgs mass parameters through the RGE’s is not as large as the one found in the analysis of Ref. . Thus in the analysis of Ref. no unification condition was assumed for the Yukawa couplings, and the only requirement on them was to predict fixed values for the fermion masses. In the case of $`m_b`$, the effects induced by the phases via $`\mathrm{\Delta }m_b`$ were compensated by variations on $`h_b`$ so as to obtain a fixed $`m_b(M_Z)`$. Since such adjustments of $`h_b`$ induced large changes on the Higgs parameters, the relic density was very sensitive to the phases. In the present case, $`h_b`$ is approximately fixed by the condition $`h_b=h_\tau `$ at the GUT scale (where $`h_\tau `$ is determined by $`\mathrm{tan}\beta `$ and $`m_\tau `$).
The fluctuation of $`\mathrm{\Delta }m_b`$ with the phases enters in the prediction of $`m_b(M_Z)`$ which is allowed to vary in its experimental range. Thus in contrast to the analysis of Ref. $`h_b`$ is not adjusted as the phases vary in the present analysis. The value of $`h_b`$ is approximately fixed by the condition $`h_b=h_\tau `$ at the GUT scale. Consequently, the phases do not have a big effect on the $`\mathrm{\Omega }h^2`$ prediction in the present scenario. For example, in Fig. (7), $`\mathrm{\Omega }h^2`$ varies only in the approximate range 0.10 – 0.13. The effects of variations with $`\alpha _0`$ for the point Eq. (24) are analyzed in Fig. (8). Specifically, Fig. (8) gives an analysis of the neutralino relic density in the $`|A|_0\alpha _0`$ plane for the input of Eq. (24) along with $`\theta _\mu =.5,\xi _3=1.7`$. One finds a considerable structure here exhibiting the important effects of $`\alpha _0`$ in this case. The relic density remains within the WMAP bounds, in the dark hatched area, while the area above the dashed line has a relic density below the lower bound.
## 5 WMAP dark matter constraint and full Yukawa unification
In the above we discussed the satisfaction of the WMAP relic density constraints consistent with the BR($`bs+\gamma `$) and $`b\tau `$ unification constraints within mSUGRA and its extensions including phases. It was seen that the loop corrections to the b quark mass (and to the $`\tau `$ lepton mass) play an important role in accomplishing $`b\tau `$ Yukawa unification at the GUT scale consistent with the experimental values for the $`\tau `$ lepton and the b quark masses. The values of $`\mathrm{tan}\beta `$ used in the analysis above were fairly large, lying in the range up to $`4045`$. When $`\mathrm{tan}\beta `$ exceeds these values the possibility that full Yukawa unification for the third generation holds becomes feasible. Here we investigate this possibility in further detail to determine if WMAP relic density and the BR($`bs+\gamma `$) constraints can also be simultaneously satisfied. In the analysis we will allow for the dependence on CP phases.
In Fig. (9) we present an analysis of full Yukawa unification and we also display the constraints of relic density and of BR($`bs+\gamma `$). We impose full Yukawa unification at the GUT scale, the value of $`tan\beta `$ is therefore fixed by the experimental $`\tau `$ and top masses. As before, $`m_b(M_Z)`$ is a prediction. Typically, there are two main constraints on $`m_0`$ and $`m_{1/2}`$ for a given $`A_0`$. These are the condition of radiative EWSB (or almost equivalently $`m_A<120`$ GeV) and the condition that the LSP be neutral. The constraints on $`m_0`$ and $`m_{1/2}`$ such that both conditions are met were described in an early paper and also emphasized in Ref. which gave the relation
$$m_A^2=\alpha m_{1/2}^2\beta m_0^2\mathrm{constant}$$
(25)
where the coefficients $`\alpha `$ and $`\beta `$ are positive and $`0.1`$, and the constant is $`M_Z`$. Thus for fixed $`m_A`$ one has a hyperbolic branch. Furthermore, the requirement that LSP be neutral, i.e. $`m_\chi <m_{\stackrel{~}{\tau }}`$, makes another cut in the allowed area. While Fig. (9) exhibits a narrow area where the WMAP relic density constraint is satisfied, one finds that $`m_b(M_Z)`$ is outside the experimental bounds (the line $`m_b(M_Z)=2.50`$ GeV is presented as a reference line in Fig. (9)).
In the whole figure the value of $`m_b(M_Z)`$ lies below the lower experimental bound. The region satisfying the BR($`bs\gamma `$) bounds is also exhibited in Fig. (9). The charged Higgs contribution is enhanced in this case due to the low values of its mass, $`m_{H^+}`$ (lines corresponding to the values $`m_{H^+}=300`$ and $`500`$ GeV are given as reference). Since the SUSY contribution is also positive the value of BR($`bs+\gamma `$) lies below its upper experimental limit only for the small region found at $`m_{1/2}2900`$ GeV. For the SM contribution we followed the considerations of by using the $`\overline{MS}`$ running charm mass, so that $`\frac{m_c}{m_b}=0.29`$. In this case the central value of Eq. 12 for the SM prediction is obtained. However, as argued in Ref. the theoretical SM prediction is possibly lower. Thus as an illustration we also give an analysis using the pole mass ratio $`\frac{m_c}{m_b}=0.29`$ which leads to a SM prediction of $`3.33\times 10^4`$.
We investigate now the implications of extending the parameter space by CP phases for a selected point in the coannihilation region. Fig. (10) shows the $`\theta _\mu \xi _3`$ plane for $`m_0=880`$ GeV, $`M_{1/2}`$=1500 GeV and $`A_0=0`$. The value of $`tan\beta `$ lies in the range 51 – 54.5. The prediction for the relic density remains in the WMAP range, since the neutralino remains in the coannihilation area. The regions where $`m_b(M_Z)`$ and BR($`bs\gamma `$) lie inside the experimental bounds are shown. There is only a rather tiny region, roughly at $`\theta _\mu =\pi /2`$ and $`\xi _3=0`$, where the Yukawa unification constraints and the BR($`bs\gamma `$) are simultaneously satisfied. This area is affected by the uncertainty in the determination of the SM value for BR($`bs\gamma `$). The area is significantly enlarged when the ratio $`\frac{m_c}{m_b}=0.29`$ is used in the BR($`bs\gamma `$) computation.
## 6 Consistency with the EDM constraints.
With inclusion of phases, one has to account for the satisfaction of the EDM constraints. In the following we demonstrate that, there exist regions in the parameter space, where the WMAP, the $`b\tau `$ unification, the BR($`bs\gamma `$) as well as the EDM constraints are all satisfied when the phases are large. In Table 1 we define two points, one for $`\mathrm{tan}\beta =40`$ and another for $`\mathrm{tan}\beta =45`$ where all constraints are satisfied as shown in Table 2.
A more detailed exhibition of the value of $`m_b(M_Z)`$ and BR($`bs+\gamma `$) as $`\theta _\mu `$ and $`\xi _3`$ varies is given in Figs. (11) and (12) while $`\xi _1`$, $`\xi _2`$ are set at the values given in Table 1. It was shown in Ref., that if the EDM constraints are satisfied for a given point, there exists a scaling region, where $`m_0,m_{1/2},A_0`$ scale by a common factor $`\lambda `$, in which the EDM constraints also are satisfied, for a reasonable range of $`\lambda `$ around one. The allowed range of $`\lambda `$ depends on other dynamical parameters. For Point (i) $`m_0`$ and $`m_{1/2}`$ are related by $`m_0=0.832m_{1/2}`$ while for Point (ii) this relation becomes $`m_0=1.80m_{1/2}`$. For point (ii) the EDM constrains are satisfied down to $`m_{1/2}750`$ GeV. In Fig. (13) we display BR($`bs\gamma `$) and the neutralino relic density for the case of two points in Table 1. The analysis shows that $`m_b(M_Z)`$ remains inside its experimental range for the range of parameters shown in the figure. The qualitative behavior of the relic density in both cases can be understood by comparison with the corresponding cases in Fig. 4. For $`\mathrm{tan}\beta =40`$ the line $`m_0=0.832m_{1/2}`$ has a sizeable overlap with the WMAP area, whereas for $`\mathrm{tan}\beta =45`$ the line $`m_0=1.832m_{1/2}`$ intersects the WMAP area. The values of $`m_b(M_Z)`$ ranges from 2.80 to 2.86 GeV for the case (i) and from 2.84 to 2.96 GeV for case (ii).
## 7 Conclusions
The main focus of this work is an analysis of the neutralino relic density consistent with the WMAP data under the constraint of $`b\tau `$ Yukawa unification, and the constraint of $`bs+\gamma `$ branching ratio. In the analysis of the $`bs+\gamma `$ branching ratio. we have included the $`\mathrm{tan}\beta `$ enhanced NLO corrections which contribute to the Wilson co-efficients. These enhancements are codified via the epsilon terms defined in Eq.(13). There is ambiguity in the sign of some of the terms among the various groups. To resolve this we carried out an independent calculation of these quantities as discussed in Sec.3. The analysis is carried out within SUGRA unified models where universality on the magnitudes of soft parameters at the GUT scale is assumed, but we allow for CP violating phases and specifically allow non-universality of the phases in the gaugino mass sector. First we give an analysis for the case when all the soft parameters are real. This is the mSUGRA case, and here we find that for values of $`\mathrm{tan}\beta `$ in the range 27-48, one obtains an amount of dark matter consistent with WMAP as well as consistency with $`b\tau `$ unification and with the $`bs+\gamma `$ constraint. An interesting phenomenon that arises is the following: There are three regions in the $`m_0m_{\frac{1}{2}}`$ parameter space where relic density and the $`bs+\gamma `$ constraint can be satisfied in general. These consist of the coannihilation region, the resonance region, and the HB/FP region. Of these only the first two can satisfy the Yukawa unification constraint. Thus the constraint of Yukawa unification narrows the available parameter space by eliminating the HB/FP region. We then extend this analysis to include phases and show that new regions of the parameter space allow for consistency with the WMAP data and other constraints extending the allowed region of the parameter space. In b-$`\tau `$ unification case we find explicit phase arrangements such that the EDM bounds are satisfied, $`m_b(M_Z)`$ and the rate for BR($`bs\gamma `$) lie within their experimental ranges, and the prediction of the neutralino relic density lies within the WMAP bounds. We have also given an analysis of the full $`b\tau t`$ Yukawa unification constraint with inclusion of CP phases. We find a small area where $`m_b`$ is predicted inside the experimental range and the BR($`bs\gamma `$) bound is satisfied. Furthermore, the relic density of neutralinos lies within the WMAP bounds due to $`\chi \stackrel{~}{\tau }`$ coannihilations. However, this area is rather small and moreover we could not find phase arrangements satisfying the EDM constraints. It is conjectured that inclusion of additional non-universalities could rectify the situation.
Acknowledgments
MEG acknowledges support from the ’Consejería de Educación de la Junta de Andalucía’, the Spanish DGICYT under contract BFM2003-01266 and European Network for Theoretical Astroparticle Physics (ENTApP), member of ILIAS, EC contract number RII-CT-2004-506222. The research of TI and PN was supported in part by NSF grant PHY-0139967. PN also acknowledges support from the Alexander von Humboldt Foundation and thanks the Max Planck Institute, Munich for hospitality extended him. SS is supported by Fundacão de Amparo à Pesquisa do Estado de São Paulo (FAPESP).
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# Extended Jaynes-Cummings models and (quasi)-exact solvability
## 1 Introduction
Several new theoretical aspects of quantum mechanics have been developped in the last years. In series of papers (see e.g. and for a recent review) it is shown that the traditional self adjointness requirement of the Hamiltonian operator is not a necessary condition to guarantee a real spectrum and that the weaker condition of PT-invariance of the Hamiltonian is sufficient for that purpose. An alternative possibility for an operator to admit a real spectrum is also developed in . It is the notion of pseudo-hermiticity. Following the ideas of , we remind here that a Hamiltonian is called $`\eta `$ pseudo-hermitian if it satisfies the relation $`\eta H\eta ^1=H^{}`$, where $`\eta `$ denotes a linear hermitian operator. It is this new notion (i.e pseudo-hermiticity property) of non hermitian Hamiltonians which explains the reality of their energy spectrum. This important property has further been considered in Refs..
Another direction of development of quantum mechanics is the notion of quasi exact solvability . It provides techniques to construct linear operators preserving a finite dimensional subspace $`𝒱`$ of the Hilbert space. Accordingly, the so called Quasi Exactly Solvable operators, once restricted on $`𝒱`$ can be diagonalized by means of algebraic methods. The QES property is strongly connected to finite dimensional representation of Lie or graded Lie algebras . Amongst many models used to describe quantum properties of physical systems, the Jaynes-Cummings model play an important role . It describes, in a simple way the interaction of photons with a spin-1/2 particle. From the mathematical point of view, the Jaynes-Cummings model is described by a self-adjoint operator and it is completely solvable in a sense that the entire spectrum can be computed algebraically.
The purpose of this paper is to consider operators generalizing the Jaynes-Cummings Hamiltonians which are neither self-adjoint nor PT-invariant but which are pseudo-hermitian with respect to two different operators. In particular, from the original Jaynes-Cummings model (JCM in the following), we construct an extended one by adding a polynomial of the form $`P(a^{},a)`$ ($`a^{},a`$ are the usual creation and annihilation operators) of degree $`d2`$ in the diagonal part of the hamiltonian. Some particular choices of $`P`$ are constructed in such a way that the resulting operator becomes QES. The non-diagonal interaction part is also modified in such a way that (i) multiple photon exchanges are allowed and (ii) the full operator can be hermitian or pseudo-hermitian.
Here is the plan of the paper. In section 2, we revisit the Hamiltonian considered in Ref. and express it in terms of differential operator of a real variable $`x`$. This reveals its exact solvability in terms of differential operators acting on sets of polynomials of appropriate degrees in $`x`$. In Sect. 3 we propose a family of operators which generalize the original JC Hamiltonian in several respects. The (pseudo)-hermiticity of these operators are analysed and the spectra and the eigenvectors are computed in details for a few of them. The differences in the spectrum corresponding to Hermitian and pseudo-Hermitian are pointed out. In particular, the energy eigenvalues are entirely real in spite of the fact that they are associated to a non hermitian and non $`PT`$-invariant Hamiltonian. The reality of those eigenvalues is ensured by the pseudo-hermiticity of the Hamiltonians. The section 4 is devoted to QES extensions of the JCM. These are constructed in such a way that, both, one-photon and two-photons exchange terms coexist in the non-diagonal interacting terms. By construction, these new models preserve finite dimensional vector spaces of the Hilbert spaces ,the algebraic part of the spectrum is computed in Sect 5. Further properties of these new types of QES operators, say $`H_T`$, can be discussed. Namely, following the ideas of we show in Sect. 6 that the solutions of the spectral equation $`H_T\psi =E\psi `$ for generic values of $`E`$ lead to new types of recurence relations. The relations between $`H_T`$ and specific graded algebras are pointed out in Sect 7. Finally, the section 8 is kept for concluding remarks.
## 2 Exactly solvable pseudo-hermitian Hamiltonian
In this section we consider the Hamiltonian describing a system of a spin-$`\frac{1}{2}`$ particle in the external magnetic field, $`\stackrel{}{B}`$ which couples to a harmonic oscillator through some nonhermitian interaction
$$H=\mu \stackrel{}{\sigma }\stackrel{}{B}+\mathrm{}\omega a^{}a+\rho (\sigma _+a\sigma _{}a^{}).$$
(1)
Here $`\stackrel{}{\sigma }`$ denotes Pauli matrices, $`\rho `$ is some arbitrary real parameter and $`\sigma _\pm \frac{1}{2}[\sigma _x\pm i\sigma _y]`$. $`\sigma _+`$ and $`\sigma _{}`$ can be expressed in matrix form
$$\sigma _+=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),\sigma _{}=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right).$$
(2)
Our purpose is to relate the Hamiltonian above to an appropriate differential operator preserving a family of spaces of polynomials in the variable $`x`$, following the ideas of exactly and quasi-exactly solvable operators . With this aim, we use the usual creation and annihilation operators respectively $`a^{}`$ and $`a`$ which are defined as follows
$$a^{}=\frac{p+im\omega x}{\sqrt{2m\omega \mathrm{}}},a=\frac{pim\omega x}{\sqrt{2m\omega \mathrm{}}},$$
(3)
where $`p=i\frac{d}{dx}`$. The external magnetic field is chosen in $`z`$-direction (i.e $`\stackrel{}{B}=B_0\stackrel{}{z}`$) in order to reduce the Hamiltonian defined in Eq.(1) and it has the form
$$H=\frac{ϵ}{2}\sigma _z+\mathrm{}\omega a^{}a+\rho (\sigma _+a\sigma _{}a^{}),$$
(4)
where $`ϵ=2\mu B_0`$. As $`\sigma _\pm ^{}=\sigma _{}`$, it is pointed out that this Hamiltonian is not hermitian
$`H^{}`$ $`={\displaystyle \frac{ϵ}{2}}\sigma _z+\mathrm{}\omega a^{}a\rho (\sigma _+a\sigma _{}a^{}),`$ (5)
$`H.`$
Thus as,
$`PTH(PT)^1`$ $`={\displaystyle \frac{ϵ}{2}}\sigma _z+\mathrm{}\omega a^{}a+\rho (\sigma _+a^{}\sigma _{}a),`$ (6)
$`H,`$
one can see that the Hamiltonian (1) is not PT symmetric i.e $`HH^{PT}`$ .
The next step is to write $`H`$ in terms of differential operators(i.e $`p=i\frac{d}{dx}`$) and of variable $`x`$. The purpose of these transformations is to reveal the exact solvability of the operator $`H`$ by using the quasi-exactly solvable (QES) technique as has been considered in Ref.. Replacing the operators $`a^{}`$ and $`a`$ by their expressions(as given in Eq.(3))in the Eq.(4), the Hamiltonian of the model is written now as follows
$$H=\frac{ϵ}{2}\sigma _z+\frac{p^2m\omega +m^2\omega ^2x^2}{2m}+\rho \frac{[\sigma _+(pim\omega x)\sigma _{}(p+im\omega x)]}{\sqrt{2m\omega \mathrm{}}}$$
(7)
In order to reveal the solvability of the above operator $`H`$, we first perform the standard (often called ” gauge”) transformation
$$\stackrel{~}{H}=R^1HR,R=exp(\frac{m\omega x^2}{2}).$$
(8)
After some algebra, the new Hamiltonian $`\stackrel{~}{H}`$ is obtained and is given by
$`\stackrel{~}{H}`$ $`={\displaystyle \frac{ϵ}{2}}\sigma _z{\displaystyle \frac{1}{2m}}{\displaystyle \frac{d^2}{dx^2}}+\omega x{\displaystyle \frac{d}{dx}}+\rho {\displaystyle \frac{[\sigma _+p\sigma _{}(p+2im\omega x)]}{\sqrt{2m\omega \mathrm{}}}}`$ (9)
$`={\displaystyle \frac{ϵ}{2}}\sigma _z+{\displaystyle \frac{p^2}{2m}}+\omega x{\displaystyle \frac{d}{dx}}+\rho {\displaystyle \frac{[\sigma _+p\sigma _{}(p+2im\omega x)]}{\sqrt{2m\omega \mathrm{}}}}`$
Replacing $`\sigma _z`$, $`\sigma _+`$ and $`\sigma _{}`$ by their matrix form, the final form of the Hamiltonian $`\stackrel{~}{H}`$ reads
$`\stackrel{~}{H}`$ $`={\displaystyle \frac{ϵ}{2}}\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)+\left(\begin{array}{cc}\frac{p^2}{2m}+\omega x\frac{d}{dx}& 0\\ 0& \frac{p^2}{2m}+\omega x\frac{d}{dx}\end{array}\right)\rho \left(\begin{array}{cc}0& \frac{p}{\sqrt{2m\omega \mathrm{}}}\\ \frac{p+2im\omega x}{\sqrt{2m\omega \mathrm{}}}& 0\end{array}\right)`$ (10)
$`=\left(\begin{array}{cc}\frac{p^2}{2m}+\omega x\frac{d}{dx}+\frac{ϵ}{2}& \rho \frac{p}{\sqrt{2m\omega \mathrm{}}}\\ \rho \frac{p+2im\omega x}{\sqrt{2m\omega \mathrm{}}}& \frac{p^2}{2m}+\omega x\frac{d}{dx}\frac{ϵ}{2}\end{array}\right).`$
Then, the operator $`\stackrel{~}{H}`$ is typically QES because it preserves a finite dimensional vector spaces of polynomials namely $`𝒱_n=(P_{n1}(x),P_n(x))^t`$ with $`n\mathrm{IN}`$. Moreover $`\stackrel{~}{H}`$ is exactly solvable because $`n`$ does not have to be fixed (it can be any nonnegative integer).
Note that the above Hamiltonian $`\stackrel{~}{H}`$ is not invariant under simultaneous parity operator(P) and time reversal (T)reflection (i.e respectively $`xx`$ and $`ii`$) . Even if the operator $`\stackrel{~}{H}`$(therefore $`H`$)is nonhermitian and not PT invariant, it was pointed out that its spectrum is real. The reality of the eigenvalues of $`H`$ is a consequence of the unbroken $`P\sigma _z`$(i.e combined parity operator $`P`$ and Pauli matrice $`\sigma _z`$)invariance of $`H`$ (i.e $`[H,P\sigma _z]=0`$). In other words, the spectrum is real because $`H`$ is pseudo-hermitian with respect to $`\sigma _z`$(i.e $`\sigma _zH\sigma _z^1=H^{}`$) and also to the parity operator $`P`$ (i.e $`PHP^1=H^{}`$) . We would like to mention that it is not necessary to calculate the energy eigenvalues and their corresponding eigenvectors of $`H`$ because they have been determined in . In the following section, we will construct the spectrum of the generalized Hamiltonian of the one given by Eq.(1).
## 3 Family of exactly solvable Hamiltonians
The original JCM is defined by the Hamiltonian
$$H=\frac{ϵ}{2}\sigma _3+\mathrm{}\omega a^{}a+\rho (a\sigma _++a^{}\sigma _{}),$$
(11)
where $`\rho `$ is a real parameter(i.e it is a real coupling constant). Note here that the Hamiltonian $`H`$ is hermitian.
In the next, we consider an extension of the above JCM Hamiltonian in the form
$$H=\frac{ϵ}{2}\sigma _3+\mathrm{}\omega a^{}a+P(a^{}a)+\rho (a^k\sigma _++\varphi (a^{})^k\sigma _{}),$$
(12)
where $`\varphi =\pm 1`$ and $`P(a^{}a)`$ denotes a polynomial of degree $`d2`$, $`k`$ is an integer $`1`$ and $`\rho `$ is an arbitrary real parameter. In fact, the above Hamiltonian is nonhermitian( i.e for $`\varphi =1`$ ) and not PT invariant but it satisfies the pseudo-hermiticity with the operators $`P`$ (operator of parity) and $`\sigma _3`$ (Pauli matrice ). considering $`\varphi =+1`$ the Hamiltonian given by the Eq.(12) becomes hermitian. Both for these cases, it can be easily observed that the energy spectrum is entirely real. Thus, notice that the above Hamiltonian (12) is a generalization of the Hamiltonians given by the Eqs.(1) and (11). The matrix form of $`H`$ reads
$$\left(\begin{array}{cc}\mathrm{}\omega a^{}a+P(a^{}a)+\frac{ϵ}{2}& \rho a^k\\ \varphi \rho (a^{})^k& \mathrm{}\omega a^{}a+P(a^{}a)\frac{ϵ}{2}\end{array}\right)$$
(13)
which can be easily checked to preserve the vector spaces
$$𝒱_n=span\{\left(\begin{array}{c}n\\ 0\end{array}\right),\left(\begin{array}{c}0\\ n+k\end{array}\right)\},n\mathrm{IN}.$$
(14)
It means that the action of the operator $`H`$ on the vectors states $`\left(\begin{array}{c}n\\ 0\end{array}\right)`$ and $`\left(\begin{array}{c}0\\ n+k\end{array}\right)`$ can expressed as linear combinations of these same states. Here, we are allowed to conclude that $`H`$ is exactly solvable because it preserves the vector space $`𝒱_n`$ for any integer $`n`$.
The next step is to find the energy eigenvalues and their corresponding eigenvectors of the Hamiltonian $`H`$ for $`\varphi =1`$ and for $`\varphi =+1`$. For this purpose we recall the following identities
$`a^{}an,{\displaystyle \frac{1}{2}}m_s`$ $`=nn,{\displaystyle \frac{1}{2}}m_s,`$ (15)
$`\sigma _3n,{\displaystyle \frac{1}{2}}m_s`$ $`=m_sn,{\displaystyle \frac{1}{2}}m_s,`$
$`\sigma _+n,{\displaystyle \frac{1}{2}}`$ $`=0;\sigma _+n,{\displaystyle \frac{1}{2}}=n,{\displaystyle \frac{1}{2}},`$
$`\sigma _{}n,{\displaystyle \frac{1}{2}}`$ $`=0;\sigma _{}n,{\displaystyle \frac{1}{2}}=n,{\displaystyle \frac{1}{2}},`$
with $`n`$ and $`m_s=\pm 1`$ are respectively the eigenvalues of the number operator $`a^{}a`$ and the operator $`\sigma _3`$. It is readily seen that the state $`0,\frac{1}{2}`$ is a ground state of the operator $`H`$(i.e it is constructed by the lowest values of $`n`$ and $`m_s`$ which are respectively $`0`$ and $`1`$). We have now to consider the action of $`H`$ to the state $`0,\frac{1}{2}`$ in order to find its associated eigenvalue
$`H0,{\displaystyle \frac{1}{2}}`$ $`={\displaystyle \frac{ϵ}{2}}\sigma _30,{\displaystyle \frac{1}{2}}+\mathrm{}\omega a^{}a0,{\displaystyle \frac{1}{2}}+P(a^{}a)0,{\displaystyle \frac{1}{2}}`$ (16)
$`+\rho a^k\sigma _+0,{\displaystyle \frac{1}{2}}+\varphi \rho (a^{})^k\sigma _{}0,{\displaystyle \frac{1}{2}},`$
$`={\displaystyle \frac{ϵ}{2}}\sigma _30,{\displaystyle \frac{1}{2}},`$
$`={\displaystyle \frac{ϵ}{2}}0,{\displaystyle \frac{1}{2}}.`$
It is proved now that $`\frac{ϵ}{2}`$ is the eigenvalue of the ground state $`0,\frac{1}{2}`$. It is easily understood that the next state $`0,\frac{1}{2}`$ is not an eigenstate alone of the Hamiltonian $`H`$ because applying this operator to this state, we obtain a linear combination of two states $`0,\frac{1}{2}`$ and $`k,\frac{1}{2}`$,
$$H0,\frac{1}{2}=\frac{ϵ}{2}0,\frac{1}{2}\pm \rho \sqrt{k!}k,\frac{1}{2}.$$
(17)
The state $`k,\frac{1}{2}`$ under the action of $`H`$ leads to a linear combination also of two above states
$$Hk,\frac{1}{2}=(\mathrm{}\omega k+P(k)\frac{ϵ}{2})k,\frac{1}{2}+\rho \sqrt{k!}0,\frac{1}{2}.$$
(18)
The excited states $`0,\frac{1}{2}`$ and $`k,\frac{1}{2}`$ span an invariant subspace in the space of states so that the Hamiltonian matrix is written as follows
$$H_k=\left(\begin{array}{cc}\frac{ϵ}{2}& \rho \sqrt{k!}\\ \varphi \rho \sqrt{k!}& \mathrm{}\omega k+P(k)\frac{ϵ}{2}\end{array}\right)$$
(19)
In particular, note that for $`k=1`$, $`P(k)=0`$(i.e $`P(k)=k^d`$ , $`d2`$) and considering $`\varphi =1`$, $`H_k`$ becomes the matrix $`H_1`$ constructed in . In order to find the eigenvalues of the Hamiltonian matrix(19), we have to solve the following usual equation(i.e characteristic polynomial)
$`det(H_k\lambda 1\mathrm{I})`$ $`=0,`$ (20)
$`\left(\begin{array}{cc}\frac{ϵ}{2}\lambda & \rho \sqrt{k!}\\ \varphi \rho \sqrt{k!}& \mathrm{}\omega k+P(k)\frac{ϵ}{2}\lambda \end{array}\right)`$ $`=0,`$
$`4\lambda ^24(\mathrm{}\omega k+P(k))\lambda +2(\mathrm{}\omega k+P(k))ϵϵ^2+\varphi 4k!\rho ^2`$ $`=0.`$
After some algebra, the energy eigenvalues(i.e square-roots of the above equation in $`\lambda `$) of $`H_k`$ are
$`\lambda _k^I`$ $`={\displaystyle \frac{\mathrm{}\omega k+P(k)+\sqrt{(\mathrm{}\omega k+P(k)ϵ)^2+\varphi 4k!\rho ^2}}{2}},`$ (21)
$`\lambda _k^{II}`$ $`={\displaystyle \frac{\mathrm{}\omega k+P(k)\sqrt{(\mathrm{}\omega k+P(k)ϵ)^2+\varphi 4k!\rho ^2}}{2}}.`$
It is easily checked that for $`k=1`$, $`P(k)=0`$ and for $`\varphi =1`$), we obtain the eigenvalues $`\lambda _1^{I,II}`$ determined in. These are the energy eigenvalues of the Hamiltonian (1). The next step now is to calculate the associated eigenvectors of the above eigenvalues $`\lambda _k^{I,II}`$. Here, we propose to consider two cases : the first case for $`\varphi =1`$ and the second one for $`\varphi =+1`$.
### 3.1 The case $`\varphi =1`$
Considering $`\varphi =1`$, the eigenvalues (21) are given by
$`\lambda _k^I`$ $`={\displaystyle \frac{\mathrm{}\omega k+P(k)+\sqrt{(\mathrm{}\omega k+P(k)ϵ)^24k!\rho ^2}}{2}},`$ (22)
$`\lambda _k^{II}`$ $`={\displaystyle \frac{\mathrm{}\omega k+P(k)\sqrt{(\mathrm{}\omega k+P(k)ϵ)^24k!\rho ^2}}{2}}.`$
For the sake simplicity, we can impose $`P(k)=0`$ and the eigenvalues $`\lambda _k^{I,II}`$ have the form
$`\lambda _k^I`$ $`={\displaystyle \frac{\mathrm{}\omega k+\sqrt{(\mathrm{}\omega kϵ)^24k!\rho ^2}}{2}},`$ (23)
$`\lambda _k^{II}`$ $`={\displaystyle \frac{\mathrm{}\omega k\sqrt{(\mathrm{}\omega kϵ)^24k!\rho ^2}}{2}}.`$
The following relations are considered as in
$`|\mathrm{}\omega kϵ|`$ $`2\rho \sqrt{k!},`$ (24)
$`2\rho \sqrt{k!}`$ $`=(\mathrm{}\omega kϵ)\mathrm{sin}\theta _k`$
and the Hamiltonian matrix given by (19) reads
$`H_k`$ $`=\left(\begin{array}{cc}\frac{ϵ}{2}& \rho \sqrt{k!}\\ \rho \sqrt{k!}& \mathrm{}\omega k\frac{ϵ}{2}\end{array}\right),`$ (25)
$`=\left(\begin{array}{cc}\frac{ϵ}{2}& \frac{1}{2}(\mathrm{}\omega kϵ)\mathrm{sin}\theta _k\\ \frac{1}{2}(\mathrm{}\omega kϵ)\mathrm{sin}\theta _k& \mathrm{}\omega k\frac{ϵ}{2}\end{array}\right).`$
Taking account of the following equation
$$\left(\begin{array}{cc}\frac{ϵ}{2}& \frac{1}{2}(\mathrm{}\omega kϵ)\mathrm{sin}\theta _k\\ \frac{1}{2}(\mathrm{}\omega kϵ)\mathrm{sin}\theta _k& \mathrm{}\omega k\frac{ϵ}{2}\end{array}\right)\left(\begin{array}{c}A\\ B\end{array}\right)=\lambda _k^{I,II}\left(\begin{array}{c}A\\ B\end{array}\right),$$
(26)
the associated eigenvectors of $`\lambda _k^{I,II}`$ are determined
$`\psi _k^I`$ $`=\mathrm{sin}{\displaystyle \frac{\theta _k}{2}}0,{\displaystyle \frac{1}{2}}+\mathrm{cos}{\displaystyle \frac{\theta _k}{2}}k,{\displaystyle \frac{1}{2}},`$ (27)
$`for\lambda _k^I={\displaystyle \frac{\mathrm{}\omega k}{2}}(1+\mathrm{cos}\theta _k){\displaystyle \frac{ϵ}{2}}\mathrm{cos}\theta _k,`$
with $`A=\mathrm{sin}\frac{\theta _k}{2}`$ and $`B=\mathrm{cos}\frac{\theta _k}{2}`$.
$`\psi _k^{II}`$ $`=\mathrm{cos}{\displaystyle \frac{\theta _k}{2}}0,{\displaystyle \frac{1}{2}}+\mathrm{sin}{\displaystyle \frac{\theta _k}{2}}k,{\displaystyle \frac{1}{2}},`$ (28)
$`for\lambda _k^{II}={\displaystyle \frac{\mathrm{}\omega k}{2}}(1\mathrm{cos}\theta _k)+{\displaystyle \frac{ϵ}{2}}\mathrm{cos}\theta _k,`$
with $`A=\mathrm{cos}\frac{\theta _k}{2}`$ and $`B=\mathrm{sin}\frac{\theta _k}{2}`$.
In particular, for $`k=1`$, it is easily checked that $`\psi _k^I`$ and $`\psi _k^{II}`$ become respectively $`\psi _1^I`$ and $`\psi _1^{II}`$ which were determined in .
### 3.2 The case $`\varphi =+1`$
Taking account of $`\varphi =+1`$ and imposing $`P(k)=0`$ , the eigenvalues (21) read
$`\lambda _k^I`$ $`={\displaystyle \frac{\mathrm{}\omega k+\sqrt{(\mathrm{}\omega kϵ)^2+4k!\rho ^2}}{2}},`$ (29)
$`\lambda _k^{II}`$ $`={\displaystyle \frac{\mathrm{}\omega k\sqrt{(\mathrm{}\omega kϵ)^2+4k!\rho ^2}}{2}}.`$
The relations considered in Eq.(24) become
$`|\mathrm{}\omega kϵ|`$ $`2\rho \sqrt{k!},`$ (30)
$`2\rho \sqrt{k!}`$ $`=(\mathrm{}\omega kϵ)\mathrm{sinh}\theta _k.`$
Following the same method used in the previous case, the eigenvectors associated to above eigenvalues (29) are written as follows
$`\psi _k^I`$ $`=\mathrm{sinh}{\displaystyle \frac{\theta _k}{2}}0,{\displaystyle \frac{1}{2}}+\mathrm{cosh}{\displaystyle \frac{\theta _k}{2}}k,{\displaystyle \frac{1}{2}},`$ (31)
$`for\lambda _k^I={\displaystyle \frac{\mathrm{}\omega k}{2}}(1+\mathrm{cosh}\theta _k){\displaystyle \frac{ϵ}{2}}\mathrm{cosh}\theta _k,`$
$`\psi _k^{II}`$ $`=\mathrm{cosh}{\displaystyle \frac{\theta _k}{2}}0,{\displaystyle \frac{1}{2}}\mathrm{sinh}{\displaystyle \frac{\theta _k}{2}}k,{\displaystyle \frac{1}{2}},`$
$`for\lambda _k^{II}={\displaystyle \frac{\mathrm{}\omega k}{2}}(1\mathrm{cosh}\theta _k)+{\displaystyle \frac{ϵ}{2}}\mathrm{cosh}\theta _k,`$
For $`HH^{}`$ (i.e for $`\varphi =1`$), it may be easily observed that two states given in (27) and (28) are not orthogonal to each other. But one can prove that the states given by Eq.(31) (i.e for $`\varphi =+1`$, $`H=H^{}`$) are orthogonal.This property is a consequence of the hermiticity of $`H`$. In order to find the next excited states, one has to consider the next invariant subspace which is spanned by the vectors $`1,\frac{1}{2}`$ and $`k+1,\frac{1}{2}`$. The eigenvalues and eigenvectors for this doublet can be determined following the same method used previously.
### 3.3 The excited states
The next step is to generalize the previous results to the invariant subspace which is spanned by the vectors $`n,\frac{1}{2}`$ and $`n+k,\frac{1}{2}`$. Following the same technique used in the previous section and after some algebra, the Hamiltonian matrix for the above doublet is written as,
$$H_{n+k}=\left(\begin{array}{cc}\mathrm{}\omega n+P(n)+\frac{ϵ}{2}& \rho \sqrt{n+1}\mathrm{}\sqrt{n+k}\\ \varphi \rho \sqrt{n+1}\mathrm{}\sqrt{n+k}& \mathrm{}\omega (n+k)+P(n+k)\frac{ϵ}{2}\end{array}\right)$$
(32)
For the sake simplicity, we impose $`P(n)=P(n+k)=0`$ and $`H_{n+k}`$ is of the form
$$H_{n+k}=\left(\begin{array}{cc}\mathrm{}\omega n+\frac{ϵ}{2}& \rho \sqrt{n+1}\mathrm{}\sqrt{n+k}\\ \varphi \rho \sqrt{n+1}\mathrm{}\sqrt{n+k}& \mathrm{}\omega (n+k)\frac{ϵ}{2}\end{array}\right)$$
(33)
and its eigenvalues are
$`\lambda _{n+k}^I={\displaystyle \frac{\mathrm{}\omega (2n+k)+\sqrt{(\mathrm{}\omega kϵ)^2+\varphi 4\rho ^2(n+1)\mathrm{}(n+k)}}{2}},`$
$`\lambda _{n+k}^{II}={\displaystyle \frac{\mathrm{}\omega (2n+k)\sqrt{(\mathrm{}\omega kϵ)^2+\varphi 4\rho ^2(n+1)\mathrm{}(n+k)}}{2}},`$ (34)
In particular, putting $`k=1`$ and $`\varphi =1`$ only in (3.3), the above eigenvalues become the eigenvalues $`\lambda _{n+1}^{I,II}`$ associated to the operator $`H`$ given by the Eq.(1). These eigenvalues were determined in .
Now putting $`2\rho \sqrt{n+1}\mathrm{}\sqrt{n+k}=(\mathrm{}\omega kϵ)\mathrm{sin}\theta _{n+k}`$ and $`2\rho \sqrt{n+1}\mathrm{}\sqrt{n+k}=(\mathrm{}\omega kϵ)\mathrm{sinh}\theta _{n+k}`$ in Eq.(3.3) respectively for $`\varphi =1`$ and for $`\varphi =+1`$, we find the eigenvectors corresponding to the doublet $`n,\frac{1}{2}`$ and $`n+k,\frac{1}{2}`$.
First considering $`\varphi =1`$, the eigenvectors associated to this doublet are
$`\psi _{n+k}^I`$ $`=\mathrm{sin}{\displaystyle \frac{\theta _{n+k}}{2}}n,{\displaystyle \frac{1}{2}}+\mathrm{cos}{\displaystyle \frac{\theta _{n+k}}{2}}n+k,{\displaystyle \frac{1}{2}},`$ (35)
$`for\lambda _{n+k}^I=\mathrm{}\omega n+{\displaystyle \frac{\mathrm{}\omega k}{2}}(1+\mathrm{cos}\theta _{n+k}){\displaystyle \frac{ϵ}{2}}\mathrm{cos}\theta _{n+k},`$
$`\psi _{n+k}^{II}`$ $`=\mathrm{cos}{\displaystyle \frac{\theta _{n+k}}{2}}n,{\displaystyle \frac{1}{2}}+\mathrm{sin}{\displaystyle \frac{\theta _{n+k}}{2}}n+k,{\displaystyle \frac{1}{2}},`$
$`for\lambda _{n+k}^I=\mathrm{}\omega n+{\displaystyle \frac{\mathrm{}\omega k}{2}}(1\mathrm{cos}\theta _{n+k})+{\displaystyle \frac{ϵ}{2}}\mathrm{cos}\theta _{n+k}.`$
Finally considering $`\varphi =+1`$ for the Eq.(3.3), the eigenvectors for the doublet $`n,\frac{1}{2}`$ and $`n+k,\frac{1}{2}`$ are of the form
$`\psi _{n+k}^I`$ $`=\mathrm{sinh}{\displaystyle \frac{\theta _{n+k}}{2}}n,{\displaystyle \frac{1}{2}}+\mathrm{cosh}{\displaystyle \frac{\theta _{n+k}}{2}}n+k,{\displaystyle \frac{1}{2}},`$ (36)
$`for\lambda _{n+k}^I=\mathrm{}\omega n+{\displaystyle \frac{\mathrm{}\omega k}{2}}(1+\mathrm{cosh}\theta _{n+k}){\displaystyle \frac{ϵ}{2}}\mathrm{cosh}\theta _{n+k},`$
$`\psi _{n+k}^{II}`$ $`=\mathrm{cosh}{\displaystyle \frac{\theta _{n+k}}{2}}n,{\displaystyle \frac{1}{2}}\mathrm{sinh}{\displaystyle \frac{\theta _{n+k}}{2}}n+k,{\displaystyle \frac{1}{2}},`$
$`for\lambda _{n+k}^I=\mathrm{}\omega n+{\displaystyle \frac{\mathrm{}\omega k}{2}}(1\mathrm{cosh}\theta _{n+k})+{\displaystyle \frac{ϵ}{2}}\mathrm{cosh}\theta _{n+k}.`$
Note that all the discussions considered in the previous section are confirmed by these generalized results.
## 4 Quasi-exactly solvable Hamiltonians
In this section let us consider an extension of the Jaynes-Cummings Hamiltonian which includes two-photon interaction
$$H_2=\frac{ϵ}{2}\sigma _3+\mathrm{}\omega a^{}a+\rho (\sigma _+a^2+\sigma _{}a_{}^{}{}_{}{}^{2})$$
(37)
The matrix form of the above Hamiltonian reads
$$H_2=\left(\begin{array}{cc}\mathrm{}\omega a^{}a+\frac{ϵ}{2}& \rho a^2\\ \rho (a^{})^2& \mathrm{}\omega a^{}a\frac{ϵ}{2}\end{array}\right).$$
(38)
It is clear that this Hamiltonian $`H`$ is similar of the one reported in Ref. and is also a particular case of the Hamiltonian given in Eq.(13) (i.e if $`k=2,P(a^{}a)=0`$) and one can prove easily its exact solvability. However, if one would like to construct an JC-type Hamiltonian including both a one-photon and a two-photon interaction, the above Hamiltonian should be modified as follows
$$H_{12}=\left(\begin{array}{cc}\mathrm{}\omega a^{}a+\frac{ϵ}{2}& \rho a^2+\rho _1a\\ \rho (a^{})^2+\widehat{\rho }_1a^{}& \mathrm{}\omega a^{}a\frac{ϵ}{2}\end{array}\right).$$
(39)
where $`\rho ,\rho _1,\widehat{\rho }_1`$ are, a priori, arbitrary constants.
Unfortunately, the corresponding operator $`H_{12}`$ is not anylonger exactly solvable. Indeed, it is easy to show that it fails to admit any finite dimensional invariant vector spaces. Accordingly, it is impossible (to our knowledge) to find its energy spectrum by algebraic methods.
In order to restaure, at least partly, a certain algebraic solvability of $`H_{12}`$, one can attemp to supplement the Hamitonian $`H_{12}`$ with an appropriate interation term. After some algebra, one can convince oneself that adding an interaction term of the form
$$H_I=\frac{1}{n}\left(\begin{array}{cc}0& \rho _1aa^{}a\\ \widehat{\rho }_1a^{}aa^{}& 0\end{array}\right)$$
(40)
leads to a new Hamiltonian $`H_T=H_{12}+H_I`$ which is quasi-exactly solvable, as we will now demonstrate.
Assuming $`n`$ to be an integer and redefining $`c\frac{\rho _1}{n}`$, $`\widehat{c}\frac{\widehat{\rho _1}}{n}`$, the operator $`H_T`$ reads
$$H_T=\left(\begin{array}{cc}\mathrm{}\omega a^{}a+\frac{ϵ}{2}& \rho a^2+ca(a^{}an)\\ \varphi \rho (a^{})^2+\widehat{c}(a^{}an)a^{}& \mathrm{}\omega a^{}a\frac{ϵ}{2}\end{array}\right),$$
(41)
where that $`a^{}`$ and $`a`$ are respectively the usual creation and annihilation operator and $`ϵ`$ is chosen as previously according to $`ϵ=2\mu B_0`$.
The main idea now is to reveal that the above operator $`H_T`$ is quasi-exactly solvable(QES). In this purpose we construct a finite dimensional vector space which is invariant under the action of $`H_T`$. Let us apply now the Hamiltonian $`H`$ to the states $`\left(\begin{array}{c}N\\ 0\end{array}\right)`$ and $`\left(\begin{array}{c}0\\ M\end{array}\right)`$ with $`N,M\mathrm{IN}`$ as follows
$$H_T\left(\begin{array}{c}N\\ M\end{array}\right)=\left(\begin{array}{c}\left(\mathrm{}\omega N+\frac{ϵ}{2}\right)N+\rho \sqrt{M\left(M1\right)}M2+c\sqrt{M}\left(Mn\right)M1\\ \varphi \rho \sqrt{\left(N+1\right)\left(N+2\right)}N+2+\widehat{c}\sqrt{N+1}\left(N+1n\right)N+1+\left(\mathrm{}\omega M\frac{ϵ}{2}\right)M\end{array}\right).$$
(42)
In order to be in agreement with the invariance of the two vectors states $`\left(\begin{array}{c}N\\ 0\end{array}\right)`$ and $`\left(\begin{array}{c}0\\ M\end{array}\right)`$ under the action of the Hamiltonian $`H_T`$, we have to impose the value of the integer $`n`$ according to $`n=M=N+2`$ (i.e $`N=M2`$). Taking account of the above fixed value of $`n`$, we obtain
$$H_T\left(\begin{array}{c}N\\ M\end{array}\right)=\left(\begin{array}{c}\left[\left(\mathrm{}\omega N+\frac{ϵ}{2}\right)+\rho \sqrt{\left(N+2\right)\left(N+1\right)}\right]N\\ \left[\mathrm{}\omega \left(N+2\right)\frac{ϵ}{2}+\varphi \rho \sqrt{\left(N+1\right)\left(N+2\right)}\right]N+2\widehat{c}\sqrt{N+1}N+1\end{array}\right).$$
(43)
Finally the Hamiltonian $`H_T`$ is of the new form
$$H_T=\left(\begin{array}{cc}\mathrm{}\omega a^{}a+\frac{ϵ}{2}& \rho a^2+ca(a^{}a(N+2))\\ \pm \rho (a^{})^2+\widehat{c}(a^{}a(N+2))a^{}& \mathrm{}\omega a^{}a\frac{ϵ}{2}\end{array}\right).$$
(44)
As it is clear from the Eq.(43), the Hamiltonian $`H_T`$ preserves the finite dimensional vector space namely
$$𝒱_n=span\{\left(\begin{array}{c}j\\ 0\end{array}\right),\left(\begin{array}{c}0\\ k\end{array}\right),j=N,\mathrm{},0;k=N+2,\mathrm{},0\},$$
(45)
and $`n`$ is fixed according to $`n=N+2`$. From this, we conclude that the Hamiltonian $`H_T`$ is quasi-exactly solvable. Hence the terms of perturbation added to $`H_{12}`$ have broken its non solvability.
Notice that is also easily to reveal the quasi-exact solvability of the operator expressed in Eq.(41) by considering the matrix Hamiltonian Eq.(41) in terms of differential expressions. Here we have to replace the operators $`a^{}`$ and $`a`$ respectively by their differential expressions given by Eq.(3), performing the standard gauge transformation as,
$$\stackrel{~}{H}_T=exp(\frac{\omega x^2}{2})H_Texp(\frac{\omega x^2}{2}),$$
(46)
and thus, after some algebra, we obtain a matrix Hamiltonian which preserves the finite dimensional vector space of the form $`𝒱_k=(P_k(x),P_{k+2}(x))^t`$ with $`k\mathrm{IN}`$ and $`n=k+2`$ ( i.e $`n`$ which is expressed in Eq.(41)). This operator $`\stackrel{~}{H}_T`$ (therefore $`H_T`$) is quasi-exactly solvable because it is expressed in terms of the integer $`n`$ which is fixed according to $`n=k+2`$.
## 5 Spectral properties
In this section, we would like to emphasize a few properties of the spectrum of the Hamiltonian discussed above. First we stress that for given $`k`$ the JC model admits $`k`$ levels which are $`\rho `$-independant and which are not involved in the list given above. They are of the form
$$\psi _j=\left(\begin{array}{c}\stackrel{}{0}\\ |j\end{array}\right),0jk1,$$
where $`\stackrel{}{0}`$ denotes the null vector of the Hilbert space. The corresponding eigenvalue is $`E_j=j\frac{ϵ}{2}`$.
The spectrum of the JC model (and of its generalisations for $`k>1`$) varies considerably with the parameter $`\rho `$. In Fig. 1, we show the evolution of six levels in the $`k=2,\varphi =1`$ case. They correspond to the two $`\rho `$-independant eigenstates and the ones with $`n=0,1`$ in Eq.(34). In Fig. 1 and in the following we assume $`ϵ=1`$ for simplicity but the features pointed out below remain similar for $`ϵ1`$. The same levels corresponding to the non hermitian case $`\varphi =1`$ are reported on Fig. 2. The contrast with Fig.1 is obvious. Couples of eigenvalues regularly disappear at finite values of the coupling constants $`\rho `$. So that, at finite $`\rho `$ only a finite number of real eigenvalues subsist, the other being real. In this respect, the Hamiltonian is like a quasi exactly solvable operator.
The energy levels displayed on Fig.1 corresponds to the six lowest ones in the limit $`\rho =0`$. The figure clearly shows that they mix relatively quickly for increasing $`\rho `$ and that, for instance, eigenvectors involving two or more quanta become the ground state for $`\rho 1`$.
We have studied the evolution of the spectrum when the QES-extension of the model, $`H_{12}=\rho a^2+\theta a(1\frac{1}{N+2}a^{}a)`$ namely characterized by the new coupling constant $`\theta `$, is progressivel switched on. Notice that the vector $`\psi _0=(\stackrel{}{0},|0)^t`$ is an eigenvector with $`E=ϵ/2`$, irrespectively of $`\rho ,\theta `$
In the case $`\rho =0,N=1`$ the effect of the new term on the eigenvalues under consideration leads to
$$E=\frac{1}{2},\frac{1}{6}(3\pm 4\theta ),\frac{1}{6}(9\pm 2\sqrt{2})\theta ,\frac{5}{2}$$
These levels are indicated on Fig. 3 by the dotted lines and it is clearly seen that they also lead to numerous level crossing.
The evolution of the eigenvalues corresponding to the case $`\rho =1`$ is displayed by the dashed lines in Fig.3, supplemented by the black line $`E=1/2`$ which is present irrespectively of $`\rho `$. The figure clearly shows that the occurence of the new term induced only one level mixing, namely two levels cross at $`E=1/2`$ for $`\theta =1.5`$ For larger values of $`\rho `$, e.g. $`\rho =2`$, the analysis reveals that the algebraic eigenvalues depend only weakly of $`\theta `$.
## 6 Series expansion and Recurence relations
Here we would like to present another aspect of the QES Hamiltonian presented in the previous section. Following the ideas of we will construct the solution for energy $`E`$ under the form of a formal serie in the basic vector whose coefficients are polynomials in $`E`$. More precisely, we write the solution of the equation
$$H_T\psi =E\psi ,$$
(47)
in the form
$$\psi =\left(\begin{array}{c}_{j=0}^{\mathrm{}}p_j(E)j\\ _{j=2}^{\mathrm{}}q_j(E)j+2\end{array}\right)$$
(48)
and where $`H_T`$ is given by the Eq.(41). After some algebra it can be seen that the polynomials $`p_j(E),q_j(E)`$ obey the following recurence relations
$$A_{j+1}P_{j+1}+B_jP_j=0,$$
(49)
where
$`A_{j+1}`$ $`=\left(\begin{array}{cc}\rho \sqrt{(j+2)(j+3)}& (E(j+1)\frac{ϵ}{2})\\ 0& \widehat{c}(j+2n)\sqrt{j+2}\end{array}\right),`$ (50)
$`B_j`$ $`=\left(\begin{array}{cc}c(j+2n)\sqrt{j+2}& 0\\ (E(j+2)+\frac{ϵ}{2})& \rho \sqrt{(j+1)(j+2)}\end{array}\right),`$
$`P_j`$ $`=\left(\begin{array}{c}q_j\\ p_j\end{array}\right),j=2,1,0,1,\mathrm{}`$
These equations have to be solved with the initial conditions
$$q_2=0,q_1=𝒩$$
(51)
with $`𝒩`$ fixing the normalisation of the solution. Then the solution for $`q_j`$ turns out to be a polynomial of degree $`E^{2j}`$. The quasi-exact solvability of the system leads to the fact that $`A_{n1}`$ is not invertible and that $`p_{n1}`$ can be choosen arbitrarily. With the choice $`p_{n1}=0`$ it turns out that all polynomials $`p_j,q_j`$ with $`jn2`$ are proportional to $`q_{n3}(E)`$. As a consequence for fixed $`n`$ and for the values of $`E`$ such that $`q_{n3}(E)=0`$ the serie above is truncated and the set of algebraic eigenvectors are recovered. We would like to stress that series considered in this section are built with the basis vector of the harmonic oscillator and not on monomials in $`x`$ contrasting with the construction of Ref.. In the case of standard QES equations there it appears a three terms recurence relations which leads to sets of orthogonal relation. In the case of systems of QES equations adressed in the recurence relation is also three terms but the situation here is quite different. Actually, it is to our knowledge, an open question to know whether the set of polynomials $`(p_j(E),q_j(E))`$ are somehow orthogonal as it is the case for standard scalar equation.
## 7 Hidden algebraic structures
As pointed out in the previous sections, the different Hamiltonians studied here posses the property that their spectrum can be (partly or fully) computed. This property is deeply related to the fact that the corresponding operators are elements of the enveloping algebra of particular graded algebra in an appropriate finite dimensional representation. The classification of linear operators preserving the vector spaces $`𝒱(m,n)=(P_m(x),P_n(x))^t`$ was reported in . It is shown that these operators are the elements of the enveloping algebra of some non-linear graded algebra depending essentially of $`|mn|`$. Note that, in the present context, the difference $`|mn|`$ is nothing else but the parameter called $`k`$ in Sect. 3. The cases $`k=1`$ and $`k=2`$ are special because the underlying algebra is indeed a graded Lie algebra. In the case $`k=1`$, related to the conventional JC model, the Hamiltonian is an element of the enveloping algebra of $`osp(2,2)`$; in the representation constucted in . The generators involved in this relation do not depend explicitely on $`n`$, i.e. of the dimension of the representation, explaining that the Hamiltonian is exactly solvable. Finally, in the case $`k=2`$, the Hamiltonian is an element of the graded Lie algebra q(2), as shown in . This algebra possesses an sl(2)$`\times `$U(1) bosonic subalgebra and six fermionic operators splitted into three triplets of the sl(2) subalgebra. In the case of the JC model corresponding to $`k=2`$, the Hamiltonian is independant on the dimension of the representation $`n`$ and the model is exactly solvable. For the modified model of Sect. 4, the supplementary interaction term $`H_I`$ defined in (40) indeed depends on $`n`$ and the operator admit only the vector space $`𝒱_n`$ as finite dimensional invariant vector space.
## 8 Conclusions
In this letter, we have considered several extensions of JCM by adding to its original Hamiltonian the polynomial $`P(a^{}a)`$ of degree $`d2`$ and an arbitrary sign, say $`\varphi =\pm 1`$, in the non-diagonal interaction term. In fact, considering this sign $`\varphi =1`$, these extended Hamiltonians are nonhermitian and not PT invariant but they satisfy the pseudo-hermiticity with respect of different operators $`P`$ and $`\sigma _3`$. This new property reveals the reality of the energy spectrum which has been constructed algebraically. They become hermitian when one considers the sign $`\varphi =1`$. Notice that these Hamiltonians are completely solvable as it has been pointed out by the QES technique.
Several usual properties available with hermitian Hamiltonian are not kep with pseudo-hermitican. Namely the eigenstates given by (27) and (28)(i.e corresponding to the doublet $`0,\frac{1}{2}`$ and $`k,\frac{1}{2}`$) are not orthogonal to each other, but they are orthogonal to all eigenstates corresponding to other doublets. For example, the eigenstate (27) and the one given by Eq.(35)(i.e it corresponds to the doublet $`n,\frac{1}{2}`$ and $`n+k,\frac{1}{2}`$) are orthogonal to each other. The eigenstates of any particular doublet are orthogonal to each other only if $`\theta _m=m\pi `$ (i.e with $`m=0,1,2,\mathrm{},k,\mathrm{},n+k`$), this implies $`\rho =0`$ because it depends to $`\mathrm{sin}\theta _m`$. In fact, as the energy eigenvalues are entirely real, it is impossible to have all eigenstates orthogonal to each other. This is explained by the unbroken symmetry of the operator $`P\sigma _3`$. But for energy eigenvalues complex, the orthonormality condition is satisfied by all the associated eigenstates. All these discussions are the result of the scalar product applied to those eigenstates.
We manage to construct a JC-type Hamiltonian describing both one and two-photons interactions in terms of quasi exactly solvable operators. This involves a very specific interaction term of degree one in the creators and annihilators which can be seen as a perturbation of more conventional p-photons interacting term. Several properties of this new family of QES-operators have been presented. Namely, (i) they can be written in terms of the generators of the graded Lie algebra osp(2,2) in a suitable representation; (ii) when expressed as series, the formal solutions of $`H_T\psi =E\psi `$ leads to a different type of recurence relation between the different terms of the series.
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# Gravitational field of relativistic gyratons
## I Introduction
Studies of the gravitational fields of beams and pulses of light have a long history. Tolman To found a solution in the linearized approximation. Peres Pe1 ; Pe2 and Bonnor Bo obtained exact solutions of the Einstein equations for a pencil of light. These solutions belong to the class of pp-waves. The generalization of these solutions to the case where the beam of radiation carries angular momentum has been found recently in FrFu:05 . Such a solution corresponds to a pulsed beam of radiation with negligible radius of cross-section, finite duration in time, and which has finite both energy $`E`$, and angular momentum $`J`$. An ultra-relativistic source with these properties is called a gyraton.
The gyraton-type solutions are of general interest, since for example they allow one to address the question: What is the gravitational field of a photon or ultrarelativistic electron or proton? This question becomes important in the discussion of possible mini-black-hole production in future collider or cosmic ray experiments. In the absence of spin, one can use the Aichelburg-Sexl metric AiSe ; BaHo to describe the gravitational field of each of the colliding particles. Such an approach allows one to estimate the cross-section for mini-black-hole formation EaGi:02 ; YoNa:02 ; YoNa:03 ; YoRy (for a general review see Kant ). The metric obtained in FrFu:05 makes it possible to consider the gravitational scattering and mini-black-hole formation in the interaction of particles with spin. The estimates show FrFu:05 that the spin-spin and spin-orbit interaction may be important at the threshold energies for mini-black-hole formation.
In the present paper we study the gravitational field of gyratons. We start by discussing the general properties of the metric (1) describing the gravitational field of relativistic gyratons (Section II). First we show that this metric belongs to the class of metrics for which all the scalar invariants constructed from the curvature and its covariant derivatives vanish identically. For $`𝐀=0`$ this result was obtained by Amati and Klimcik AmatiKlimcik:89 and Horowitz and Steif HorowitzSteif:90 , who argued that such metrics are classical solutions to string theory. (For a general discussion of spacetimes with vanishing curvature invariants see Coley ; PPCM ; NSI ).
After this we show that the vacuum Einstein equations for metric (1) in a spacetime with arbitrary number $`D`$ of dimensions reduce to the linear problems for the gravitoelectric ($`\mathrm{\Phi }`$) and gravitomagnetic ($`𝐀`$) potentials in the $`(D2)`$dimensional Euclidean space. These linearized problems can be easily solved. The solutions obtained in FrFu:05 are characterized by the property that only lowest harmonics are present in the harmonic decomposition of $`\mathrm{\Phi }`$ and $`𝐀`$. For this reason one can consider the gyraton solutions presented in FrFu:05 as some ground state, while the more general solutions obtained in this paper are their excitations (or distortions). It should be emphasized that the vacuum solutions are valid only outside the region occupied by gyratons. In order to obtain a solution describing the total spacetime one needs to obtain a solution inside the gyraton. This solution depends on the gyraton structure. In the present paper we do not discuss concrete gyratons models. But since we obtain a general solution for the vacuum metric outside a gyraton, one can gurantee that for any model of the gyraton there exists a corresponding solution, so that the characteristics of the gyraton are ”encoded” in the parameters of the exterior vacuum metric.
After discussing the asymptotic properties of the gyraton metrics (Section III), we consider general solutions for $`4`$ (Section IV) and $`5`$dimensional (Section V) gyraton metrics. Section VI discusses the higher dimensional gyraton metrics. In Section VII we summarize the obtained results and discuss open problems.
## II Metric for relativistic gyratons
### II.1 Gyraton metric ansatz
Let us consider Brinkmann Br metric in $`D`$dimensional spacetime of the form
$$ds^2=g_{\mu \nu }dx^\mu dx^\nu $$
$$=2dudv+d𝐱^2+\mathrm{\Phi }du^2+2(𝐀,d𝐱)du,$$
(1)
$$\mathrm{\Phi }=\mathrm{\Phi }(u,𝐱),A_a=A_a(u,𝐱).$$
(2)
Evidently, $`l^\mu _\mu =_v`$ is the null Killing vector.
When $`\mathrm{\Phi }=𝐀=0`$, the coordinates $`x^1=v=(t+\xi )/\sqrt{2}`$ and $`x^2=u=(t\xi )/\sqrt{2}`$ are null. The coordinate $`u`$ remains null for the metric (1). The metric is generated by an object moving with the velocity of light in the $`\xi `$direction. The coordinates $`(x^3,\mathrm{},x^D)`$ are coordinates of an $`n`$dimensional space ($`n=D2`$) transverse to the direction of motion. We use bold-face symbols to denote vectors in this space. For example, $`𝐱`$ is a vector with components $`x^a`$ ($`a=3,\mathrm{},D`$). We denote by $`r`$ the length of this vector, $`r=|𝐱|`$. We also denote
$$d𝐱^2=\underset{a=3}{\overset{D}{}}(dx_a)^2,(𝐀,d𝐱)=\underset{a=3}{\overset{D}{}}A_adx^a,$$
(3)
$$=\underset{a=3}{\overset{D}{}}_a^2,\text{div}𝐀=\underset{a=3}{\overset{D}{}}A_{,a}^a.$$
(4)
Later we assume that the sum is taken over the repeated indices and omit the summation symbol. Working in the Cartesian coordinates we shall not distinguish between upper and lower indices.
The form of the metric (1) is invariant under the following (gauge) transformation
$$vv+\lambda (u,𝐱),A_aA_a\lambda _{,a},\mathrm{\Phi }\mathrm{\Phi }2\lambda _{,u}.$$
(5)
It is also invariant under rescaling
$$uau,va^1v,\mathrm{\Phi }a^2\mathrm{\Phi },𝐀a𝐀.$$
(6)
It is easy to show that for the metric (1)
$$\sqrt{g}=1,$$
(7)
and the inverse metric is
$$g^{\mu \nu }_\mu _\nu =(\mathrm{\Phi }𝐀^2)_v^22_u_v+2A_a_a_v+_a^2.$$
(8)
The non-vanishing components of the Christoffel symbol $`\mathrm{\Gamma }_{\mu ,\nu \lambda }`$ are
$$\mathrm{\Gamma }_{v,\mu \nu }=\mathrm{\Gamma }_{\mu ,\nu v}=\mathrm{\Gamma }_{a,bc}=0,\mathrm{\Gamma }_{u,uu}=\frac{1}{2}_u\mathrm{\Phi },$$
$$\mathrm{\Gamma }_{a,uu}=_uA_a\frac{1}{2}_a\mathrm{\Phi },\mathrm{\Gamma }_{u,ua}=\frac{1}{2}_a\mathrm{\Phi },$$
(9)
$$\mathrm{\Gamma }_{u,ab}=\frac{1}{2}(_aA_b+_bA_a),\mathrm{\Gamma }_{a,bu}=\frac{1}{2}F_{ab},$$
where
$$F_{ab}=_aA_b_bA_a.$$
(10)
We shall also need the Christoffel symbols $`\mathrm{\Gamma }_{\nu \lambda }^\mu =g^{\mu \alpha }\mathrm{\Gamma }_{\alpha ,\nu \lambda }`$. Their non-vanishing components are
$$\mathrm{\Gamma }_{uu}^v=\frac{1}{2}(\mathrm{\Phi }_{,u}+A^a\mathrm{\Phi }_{,a})+A^aA_{a,u},$$
$$\mathrm{\Gamma }_{ua}^v=\frac{1}{2}(\mathrm{\Phi }_{,a}F_{ab}A^b),\mathrm{\Gamma }_{ab}^v=\frac{1}{2}(A_{a,b}+A_{b,a}),$$
(11)
$$\mathrm{\Gamma }_{uu}^a=A_{,u}^a\frac{1}{2}\mathrm{\Phi }^{,a},\mathrm{\Gamma }_{ua}^b=\frac{1}{2}F_a^b.$$
It is easy to check that
$$l_{\mu ;\nu }=0.$$
(12)
It means that the null Killing vector $`𝐥`$ is covariantly constant. In the 4-dimensional case, space-times admitting a (covariantly) constant null vector field are called plane-fronted gravitational waves with parallel rays, or briefly pp-waves (see e.g. EK ; J ; ES ). Similar terminology is often used for higher dimensional metrics (see e.g. B ; O ).
### II.2 Curvature invariants
In the next section we derive conditions under which metric (1) is Ricci flat and hence obeys the vacuum Einstein equations. But before this let us prove that the metric (1) belongs to the class of metrics with vanishing curvature invariants. Namely, all the local scalar invariants constructed from the Riemann tensor and its covariant derivatives for the metric (1) vanish. This statement is valid off shell, that is the metric need not be a solution of the vacuum Einstein equations. This property is well known for 4-dimensional case, since pp-wave solutions are of Petrov type N. Generalization of this result to higher-dimensional metrics (1) with $`𝐀=0`$ was given in AmatiKlimcik:89 ; HorowitzSteif:90 . (For a general discussion of spacetimes with vanishing curvature invariants see Coley ; PPCM ; NSI ).
To demonstrate that curvature invariants vanish for the metric (1), let us consider a covariant tensor $`A_{\mu \mathrm{}\nu }`$. We shall call such a tensor degenerate if it has the following properties: It does not depend on $`v`$, and its components, which either contain at least one index $`v`$ or do not contain index $`u`$, vanish. Since $`_v`$ is the Killing vector, the Riemann curvature tensor does not depend on $`v`$. Using the expressions for the Christoffel symbols (11) one can show that the only non-vanishing components of the Riemann tensor are $`R_{[au][bu]}`$, $`R_{[ab][cu]}`$, and $`R_{[cu][ab]}`$. Hence it is degenerate. Let us demonstrate now that the action of a covariant derivative $`_\mu `$ on a degenerate tensor $`A_{\mu \mathrm{}\nu }`$ gives a tensor which is also degenerate. Really, since $`\mathrm{\Gamma }_{v\mu }^\alpha =0`$ one has
$$_vA_{\mu \mathrm{}\nu }=_vA_{\mu \mathrm{}\nu }\mathrm{\Gamma }_{v\mu }^\alpha A_{\alpha \mathrm{}\nu }\mathrm{}\mathrm{\Gamma }_{v\nu }^\alpha A_{\mu \mathrm{}\alpha }=0.$$
(13)
Thus the covariant differentiation of the degenerate tensor cannot have a non-vanishing $`v`$ component. Since $`\mathrm{\Gamma }_{\mu \nu }^u=0`$, the covariant differentiation cannot also produce a non-vanishing component which does not contain index $`u`$.
Consider now a scalar invariant constructed from any set of degenerate tensors and $`g^{\mu \nu }`$. The only non-vanishing component of $`g^{\mu \nu }`$ which contains an index $`u`$ is $`g^{uv}=1`$. Hence a scalar invariant constructed from degenerate tensors and metric always vanishes.
### II.3 Calculation of the Ricci tensor
In order to calculate the Ricci tensor for the metric (1) let us introduce the following vectors
$$V^\lambda =\frac{x^\lambda }{v},U^\lambda =\frac{x^\lambda }{u},e_{(a)}^\lambda =\frac{x^\lambda }{x^a}.$$
(14)
One has
$$V_{;\alpha }^\beta =\mathrm{\Gamma }_{v\alpha }^\beta =0,U_{;\alpha }^\beta =\mathrm{\Gamma }_{u\alpha }^\beta ,e_{(a);\beta }^\lambda =\mathrm{\Gamma }_{a\beta }^\lambda ,$$
(15)
$$U_{;\alpha }^\beta e_{(c);\beta }^\alpha =\mathrm{\Gamma }_{u\alpha }^\beta \mathrm{\Gamma }_{c\beta }^\alpha =\mathrm{\Gamma }_{ua}^b\mathrm{\Gamma }_{cb}^a+\mathrm{\Gamma }_{ua}^u\mathrm{\Gamma }_{cu}^a=0.$$
(16)
The last equality holds because $`\mathrm{\Gamma }_{cb}^a=\mathrm{\Gamma }_{ua}^u=0`$.
The Ricci identity implies
$$R_{\alpha \beta }Y^\alpha X^\beta =(X_{;\alpha }^\beta Y^\alpha )_{;\beta }X_{;\alpha }^\beta Y_\beta ^\alpha X_{;\beta \alpha }^\beta Y^\alpha .$$
(17)
From the relation $`V_{;\alpha }^\beta =0`$ it follows that $`R_{\alpha \beta }Y^\alpha V^\beta =0`$ and hence
$$R_{v\alpha }=0.$$
(18)
Let us set $`Y^\alpha =e_{(a)}^\alpha `$ and $`X^\beta =U^\beta `$, then using (16) one has
$$R_{au}=(U_{;\alpha }^\beta e_{(a)}^\alpha )_{;\beta }U_{;\beta \alpha }^\beta e_{(a)}^\alpha .$$
(19)
Using (7) one obtains
$$U_{;\beta }^\beta =\frac{1}{\sqrt{g}}_\beta (\sqrt{g}\delta _u^\beta )=0.$$
(20)
One also has
$$U_{;\alpha }^\beta e_{(a)}^\alpha =\mathrm{\Gamma }_{ua}^\beta =\frac{1}{2}\delta _b^\beta F_a^b\frac{1}{2}\delta _v^\beta (\mathrm{\Phi }_{,a}F_{ab}A^b).$$
(21)
Since $`\mathrm{\Phi }`$ and $`A_a`$ do not depend on $`v`$, and $`\mathrm{\Gamma }_{v\beta }^\alpha =0`$, one has
$$R_{au}=\frac{1}{2}_bF_a^b.$$
(22)
Similarly
$$R_{uu}=(U_{;\alpha }^\beta U^\alpha )_{;\beta }U_{;\alpha }^\beta U_{;\beta }^\alpha U_{;\beta \alpha }^\beta U^\alpha .$$
(23)
Relation (20) implies that the last term on the right hand side vanishes. Since $`U_{;\alpha }^\beta U^\alpha =\mathrm{\Gamma }_{uu}^\beta `$ using (11) one obtains
$$(U_{;\alpha }^\beta U^\alpha )_{;\beta }=_b(A_{,u}^b\mathrm{\Phi }^{,b}).$$
(24)
One also has
$$U_{;\alpha }^\beta U_{;\beta }^\alpha =\mathrm{\Gamma }_{u\alpha }^\beta \mathrm{\Gamma }_{u\beta }^\alpha =\mathrm{\Gamma }_{ua}^b\mathrm{\Gamma }_{ub}^a=\frac{1}{4}F_{ab}F^{ab}.$$
(25)
Combining these results one obtains
$$R_{uu}=_u\text{div}𝐀\frac{1}{2}\mathrm{\Phi }+\frac{1}{4}𝐅^2,$$
(26)
where
$$\text{div}𝐀=_aA^a,𝐅^2=F_{ab}F^{ab},\mathrm{\Phi }=_a^a\mathrm{\Phi }.$$
(27)
Finally, let us substitute $`X^\alpha =e_{(a)}^\alpha `$ and $`Y^\beta =e_{(b)}^\beta `$ into (17), then one has
$$R_{ab}=\left(e_{(b);\alpha }^\beta e_{(a)}^\alpha \right)_{;\beta }e_{(b);\alpha }^\beta e_{(a);\beta }^\alpha e_{(b);\beta \alpha }^\beta e_{(a)}^\alpha .$$
(28)
Using (11) it is easy to show that
$$e_{(b);\alpha }^\beta e_{(a)}^\alpha =\mathrm{\Gamma }_{ab}^\beta =\frac{1}{2}\delta _v^\beta (A_{a,b}+A_{b,a}),$$
(29)
$$e_{(b);\alpha }^\beta e_{(a);\beta }^\alpha =\mathrm{\Gamma }_{b\alpha }^\beta \mathrm{\Gamma }_{a\beta }^\alpha =0,e_{(b);\beta }^\beta =_\beta \delta _b^\beta =0.$$
(30)
Hence $`R_{ab}=0`$.
### II.4 Vacuum equations for gravitational field of a gyraton
To summarize, the metric (1) is a solution of vacuum Einstein equations if and only if the following equations are satisfied
$$_bF_a^b=0,$$
(31)
$$\mathrm{\Phi }2_u\text{div}𝐀=\frac{1}{2}𝐅^2.$$
(32)
In the next section it will be shown that for solutions describing a gyraton with finite energy and angular momentum the quantities $`\mathrm{\Phi }_{,a}`$ and $`F_{ab}`$ are vanishing at transverse space infinity. We assume that the homogeneous equations are valid everywhere outside the point $`𝐱=0`$ where a point-like source is located. It is easy to see that the left hand side of (32) is gauge invariant, that is invariant under the transformations (5). In the ‘Lorentz’ gauge $`A_{,a}^a=0`$ equation (32) takes the form
$$\mathrm{\Phi }=\frac{1}{2}𝐅^2.$$
(33)
Here $`𝐅^2=F_{ab}F^{ab}`$. Using the analogy of gravity with electromagnetism, one can say that the problem of solving the $`D`$dimensional vacuum Einstein equations for the gyraton metric is reduced to finding an electric potential $`\phi `$ and magnetic field $`F_{ab}`$ created by a local source in the $`(D2)`$dimensional Euclidean space. For these solutions the retarded time $`u`$ plays a role of an external parameter which enters through the dependence of point-like sources on $`u`$.
As we already mentioned in the Introduction, in physical applications there always exists a source of the gravitational field which generates the metric (1). We called this source a gyraton FrFu:05 . In order to obtain a solution describing the total system, one must obtain a solution inside the region occupied by the gyraton and to glue it together with a vacuum solution (1) outside it. In the present paper we study only solutions outside the gyraton. We shall obtain a general solution of the magnetostaic equations (31) for point-like currents localized at the point $`𝐱=0`$. Since this equation is linear this current can be written as a linear combination of $`\delta (𝐱)`$ and its derivatives. Similarly, one can write a general solution of the equation
$$\phi =0,$$
(34)
with a charge density, localized at $`𝐱=0`$, or, what is equivalent, with the charge-density proportional to $`\delta (𝐱)`$ and its derivatives. It is convenient to write $`\mathrm{\Phi }=\phi +\psi `$ , where
$$\psi =\frac{1}{2}𝐅^2.$$
(35)
After finding $`A_a(u,𝐱)`$ and $`\phi (u,𝐱)`$, one needs only to find the ‘induced’ potential $`\psi (u,𝐱)`$ determined by the equation (35). A formal solution of this problem can be obtained as follows. The Green function of the $`n`$dimensional Laplace operator
$$\mathrm{\Delta }𝒢_n(𝐱,𝐱^{})=\delta (𝐱𝐱^{})$$
(36)
is
$$𝒢_2(𝐱,𝐱^{})=\frac{1}{2\pi }\mathrm{ln}|𝐱𝐱^{}|,\text{if }n=2,$$
(37)
$$𝒢_n(𝐱,𝐱^{})=\frac{g_n}{|𝐱𝐱^{}|^{n2}},\text{if }n>2,$$
(38)
$$g_n=\frac{\mathrm{\Gamma }(\frac{n2}{2})}{4\pi ^{n/2}}.$$
(39)
Using these Green functions one can present the solution for $`\psi `$ in the form
$$\psi (u,𝐱)=\frac{1}{2}𝑑𝐱^{}𝒢_n(𝐱,𝐱^{})𝐅^2(u,𝐱^{}).$$
(40)
Let us emphasized that in the general case the solution (40) is only formal and may not have a well-defined sense. The reason is that for a point-like current, $`𝐅`$ has a singularity at $`𝐱=0`$. If one considers this singular function as a distribution, one needs to define what is the meaning of $`𝐅^2`$ in (40). This problem does not exist for a distributed source (gyraton). If we do not want to input an explicit form of the matter distribution within gyraton, we can proceed as follows remark .
Suppose $`𝐅`$ and $`\phi `$ are solutions with localized sources. Let us surround a point $`𝐱=0`$ by a $`(D3)`$dimensional surface $`\sigma `$. For example, one may choose $`\sigma `$ to be a round $`(D3)`$dimensional sphere of small radius $`ϵ`$. Denote by $`F_{ab}^{(\sigma )}=F_{ab}\vartheta (\sigma )`$, where $`\vartheta (\sigma )`$ is equal to 1 outside $`\sigma `$ and vanishes inside $`\sigma `$. The ”magnetic” field $`F_{ab}^{(\sigma )}`$ obeys the equation
$$_bF_a^{(\sigma )b}=n_bF_a^b\delta (\sigma ),$$
(41)
where $`𝐧`$ is the unit normal to $`\sigma `$ vector directed to the exterior of $`\sigma `$. In orther words, the field $`F_{ab}^{(\sigma )}`$ corresponds to the special case of an extended gyraton for which its angular momentum density is localised on $`\sigma `$. The value $`\psi ^{(\sigma )}`$ obtained for $`F_{ab}^{(\sigma )}`$ by using (40) is well difined. Certainly this function $`\psi ^{(\sigma )}`$ depends on the choice of $`\sigma `$. Suppose $`\sigma ^{}`$ is another surface, surrounding $`𝐱=0`$, and lying inside $`\sigma `$. It is easy to see that outside $`\sigma `$ one has
$$(\psi ^{(\sigma )}\psi ^{(\sigma )})=0.$$
(42)
That is outside $`\sigma `$ these two solutions $`\psi ^{(\sigma )}`$ and $`\psi ^{(\sigma )}`$ differ by a term which can be absorbed into the solution $`\phi `$.
For a distributed source (a gyraton) one can use a similar procedure. If one is interested in the gravitational field of the gyraton outside a surface $`\sigma `$ surrounding the matter distribution one can calculate $`\psi ^{(\sigma )}`$ and choose $`\phi `$ correspondingly. For a given distribution of the gyraton matter, the parameters of the vacuum solution outside the gyraton are uniquely specified. In sections IV and V we shall give explicit examples of the vacuum gyraton solutions.
## III Energy and angular momentum of a gyraton
### III.1 Weak field approximation
The asymptotics of functions $`\mathrm{\Phi }`$ and $`A_a`$ at the transverse-spatial infinity are related to the energy and angular momentum of a gyraton. In order to find these relations let us consider a linearized problem. Let us write the Minkowski metric in the form
$$ds_0^2=\eta _{\mu \nu }dx^\mu dx^\nu =2dudv+d𝐱^2.$$
(43)
Then its perturbation $`h_{\mu \nu }`$ generated by the stress-energy tensor $`T_{\mu \nu }`$ obeys the equation
$$\mathrm{}h_{\mu \nu }=\kappa \overline{T}_{\mu \nu },\overline{T}_{\mu \nu }=(T_{\mu \nu }\frac{1}{n}\eta _{\mu \nu }T),$$
(44)
where $`\kappa =16\pi G`$.
In the general case, if the metric (1) is a solution of the Einstein equations, then the stress-energy tensor which generates this solution possesses the following properties: (1) Its non-vanishing components are $`T_{uu}`$ and $`T_{ua}`$; (2) These components do not depend on $`v`$; and (3) It obeys the conservation law $`T_{\mu \nu }^{;\nu }=0`$. The latter condition in the linear approximations reduces to the relation
$$T_{ua,a}=0.$$
(45)
The metric perturbation $`h_{\mu \nu }`$ generated by such a stress-energy tensor also does not depend on $`v`$. Thus instead of $`D`$-dimensional $`\mathrm{}`$ operator in (44) one can substitute the $`n`$-dimensional flat Laplace operator $``$
$$h_{\mu \nu }=\kappa T_{\mu \nu }.$$
(46)
We omit the bar over $`T_{\mu \nu }`$ since its trace vanishes.
Using the Green function (36) one can write the following expression for $`h_{\mu \nu }(𝐱)`$
$$h_{\mu \nu }(u,𝐱)=\kappa 𝑑𝐱^{}𝒢_n(|𝐱_{}𝐱^{}|)\overline{T}_{\mu \nu }(u,𝐱^{}).$$
(47)
### III.2 Metric asymptotics
Consider $`T_{\mu \nu }(u,𝐱)`$ as a function of $`𝐱`$ and suppose that it vanishes outside of some compact region. We denote by $`l`$ the size (in the transverse direction) of this region. To determine the field at far distance $`rl`$ we use the following relation
$$|𝐱𝐱^{}|r\frac{(𝐱,𝐱^{})}{r}.$$
(48)
Thus if one point, $`𝐱`$, is at far distance from the source, while the other is close to it one has the following asymptotics for the Green functions
$$𝒢_2(𝐱,𝐱^{})=\frac{1}{2\pi }\mathrm{ln}r+\frac{(𝐱,𝐱^{})}{2\pi r^2}+\mathrm{},\text{if }n=2,$$
(49)
$$𝒢_n(𝐱,𝐱^{})=\frac{g_n}{r^{n2}}+\frac{g_n(n2)(𝐱,𝐱^{})}{r^n}+\mathrm{},\text{if }n>2,$$
(50)
where $`\mathrm{}`$ denote the terms of higher order in $`1/r`$. Similarly, one has
$$h_{\mu \nu }=\frac{\kappa }{2\pi }\mathrm{ln}r𝒯_{\mu \nu }+\frac{\kappa }{2\pi r^2}x^a𝒥_{a\mu \nu }+\mathrm{},\text{if }n=2,$$
(51)
$$h_{\mu \nu }=\frac{\kappa g_n}{r^{n2}}𝒯_{\mu \nu }+\frac{\kappa g_n(n2)}{r^n}x^a𝒥_{a\mu \nu }+\mathrm{},\text{if }n>2,$$
(52)
where
$$𝒯_{\mu \nu }=𝑑𝐱T_{\mu \nu },𝒥_{a\mu \nu }=𝑑𝐱x_aT_{\mu \nu }.$$
(53)
The structure of the stress-energy tensor implies that only components $`h_{uu}`$ and $`h_{ua}`$ do not vanish. In order to relate the coefficients, which enter the asymptotic expressions for these components, to physical quantities such as energy and angular momentum we use the following relations
$$𝑑𝐱T_{ua}=𝑑𝐱T_{uc}^{,c}x_a=0,$$
(54)
$$𝑑𝐱(T_u^ax^b+T_u^bx^a)=𝑑𝐱T_{u,c}^cx^ax^b=0.$$
(55)
We used here the conservation law (45).
The energy $`E`$ and the angular momentum $`J_{ab}`$ in the flat spacetime are defined by the relation
$$E=𝑑\xi 𝑑𝐱T_{tt},J_{ab}=𝑑\xi 𝑑𝐱(T_{ta}x_bT_{tb}x_a).$$
(56)
The integration is performed over a surface $`t=`$const. At this surface $`d\xi =\sqrt{2}du`$, thus
$$_{\mathrm{}}^{\mathrm{}}𝑑\xi (\mathrm{})=\sqrt{2}_{\mathrm{}}^{\mathrm{}}𝑑u(\mathrm{}).$$
(57)
One also has
$$T_{tt}=\frac{1}{2}T_{uu},T_{ta}=\frac{1}{\sqrt{2}}T_{ua}.$$
(58)
By combining these relations one obtains
$$E=𝑑u\epsilon (u),J_{ab}=𝑑uj_{ab}(u),$$
(59)
$$\epsilon (u)=\frac{1}{\sqrt{2}}𝑑𝐱T_{uu},j_{ab}(u)=2𝑑𝐱T_{ua}x_b.$$
(60)
Relation (55) shows that $`j_{ab}=j_{ba}`$. The function $`\epsilon (u)`$ describes the energy-density profile of the gyraton as a function of the retarded time $`u`$, while $`j_{ab}(u)`$ is similar profile functions for the components of the density of the angular momentum.
Using these results one obtains
$$\mathrm{\Phi }h_{uu}=\kappa \sqrt{2}\epsilon \{\begin{array}{c}\frac{1}{2\pi }\mathrm{ln}r,\text{if }n=2,\\ \\ \frac{g_n}{r^{n2}},\text{if }n>2,\end{array}$$
(61)
$$A_ah_{ua}=\frac{\kappa g_n(n2)j_{ab}x^b}{r^n}.$$
(62)
The latter relation is valid in the 4-dimensional spacetime (for $`n=2`$) if one substitutes $`1/(2\pi )`$ for $`g_n(n2)`$.
### III.3 Canonical form
If $`j_{ab}`$ were a time independent antisymmetric matrix then by making rotations
$$x^a=O_b^a\stackrel{~}{x}^b,x_a=\stackrel{~}{x}_cO_a^c$$
(63)
one would be able to bring $`h_{ua}`$ into a form where instead of $`j_{ab}`$ stands its block canonical form Gant
$$\stackrel{~}{j}_{ab}=\left(\begin{array}{ccccc}0& j_1& 0& 0& \mathrm{}\\ j_1& 0& 0& 0& \mathrm{}\\ 0& 0& 0& j_2& \mathrm{}\\ 0& 0& j_2& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(64)
In the presence of time dependence the situation is slightly more complicated. Let us consider transformations (63) with time dependent orthogonal matrix $`O_b^a(u)`$. Then
$$dx^a=O_b^ad\stackrel{~}{x}^b+\dot{O}_b^a\stackrel{~}{x}^bdu,dx_a=O_a^cd\stackrel{~}{x}_c+\dot{O}_a^c\stackrel{~}{x}_cdu,$$
(65)
where $`\dot{B}=_uB`$. Under these transformations the metric (1) preserves its form with
$$\stackrel{~}{A}_a=A_bO_a^b+B_{ab}\stackrel{~}{x}^b,$$
(66)
$$\stackrel{~}{\mathrm{\Phi }}=A_a\dot{O}_b^a\stackrel{~}{x}^b+C_{ab}\stackrel{~}{x}^a\stackrel{~}{x}^b,$$
(67)
where
$$B_{ab}=O_{ac}\dot{O}_b^c=\dot{O}_{ac}O_b^c,$$
(68)
$$C_{ab}=\dot{O}_{ac}\dot{O}_b^c=B_{ac}B_b^c.$$
(69)
It is easy to see that $`B_{ab}\stackrel{~}{x}^b`$ is itself a solution of the magnetostatic equations and correspond to a constant magnetic field with $`F_{ab}=B_{ab}`$. It means that the (linearly growing at infinity) terms generated by time dependent rigid rotations can be compensated by adding to $`A_a`$ a new solution corresponding to a constant magnetic field. This is a direct analogue of the Larmor theorem in gravitomagnetism Mash:93 .
To summarize, we demonstrated that by making a time dependent rotation and adding to $`A_a`$ a vector potential for a homogeneous time dependent magnetic field it is always possible to transform a solution (1) into the form where $`\mathrm{\Phi }`$ and $`A_a`$ have the asymptotics (61) and (62), where $`j_{ab}`$ is an antisymmetric matrix in its canonical block form (64).
## IV $`4`$dimensional gyratons
### IV.1 General solution
Before analyzing general gyraton-like solutions in an arbitrary number of spacetime dimensions, we consider special lower-dimensional cases.
Let us derive a gyraton metric in a 4-dimensional spacetime. In this case the number of transverse dimensions $`n=2`$ and our problem reduces to 2-D electro- and magnetostatics.
Let us consider the equation (31) for the magnetic field. Any antisymmetric tensor of the second order in a 2-dimensional space can be written as $`F_{ab}=Fe_{ab}`$, where $`e_{ab}`$ is the totally antisymmetric tensor. Substituting this representation into (31) one obtains that $`F=`$const. It is easy to see that the corresponding vector potential $`A_a`$ can be written as
$$A_3=\alpha x^4,A_4=\beta x^3,F=\beta \alpha .$$
(70)
The gauge transformation (5) with $`\lambda =\gamma x^3x^4`$ changes the coefficients $`\alpha \alpha \gamma `$ and $`\beta \beta \gamma `$ but preserves the value $`F`$.
Equation (35) takes the form
$$\psi =\frac{1}{2}F^2.$$
(71)
If $`F0`$, the solution $`\psi `$ does not vanish at infinity. We exclude this case. Thus we put $`F=0`$.
Let us choose a 2-dimensional contour surrounding the source at $`𝐱=0`$. When $`F_{ab}=0`$, the value of the integral
$$j(u)=\frac{2}{\kappa }_CA_a𝑑x^a,j(u)=\frac{1}{2}ϵ^{ab}j_{ab},$$
(72)
does not depend on the choice of the contour $`C`$. This quantity which enters the solution (1) has the meaning of angular momentum of the gyraton. In polar coordinates $`(r,\varphi )`$
$$x^3+ix^4=re^{i\varphi }.$$
(73)
the corresponding potential $`A_a`$ can be written as
$$A_r=0,A_\varphi =\frac{\kappa }{4\pi }j(u).$$
(74)
Let us consider now equation (34) for the 2-dimensional ‘electric’ potential $`\phi `$. A solution corresponding to a point-like charge is
$$\phi _0=\frac{\kappa \sqrt{2}}{2\pi }\epsilon (u)\mathrm{ln}r.$$
(75)
Any other solution of this equation decreasing at infinity can be written as
$$\phi =\phi _0+\underset{n=\mathrm{}}{\overset{{}_{}{}^{}\mathrm{}}{}}\frac{b_n}{r^{|n|}}e^{in\varphi },\overline{b}_n=b_n.$$
(76)
$`^{}`$ indicates that the term $`n=0`$ is excluded. In the electromagnetic analogy, the harmonics with $`n1`$ describe the field created by an electric $`n`$-pole. Since $`F=0`$, $`\mathrm{\Phi }=\phi `$ and the solution for a distorted gyraton in $`4`$dimensional spacetime is
$`ds^2`$ $`=`$ $`2dudv+dr^2+r^2d\varphi ^2+{\displaystyle \frac{\kappa }{2\pi }}j(u)dud\varphi `$ (77)
$`+`$ $`\left[{\displaystyle \frac{\kappa \sqrt{2}}{2\pi }}\epsilon (u)\mathrm{ln}r+\phi \right]du^2,`$
where $`\phi =\phi (u,r,\varphi )`$ is given by (76) with $`b_n=b_n(u)`$.
It should be emphasized, that the solution (77), which is a special case of (1), is the pp-wave metric. The properties of pp-wave metrics in 4-dimensional spacetime are well known (see e.g. J ; EK ; ES ). In particular, a vacuum pp-wave metric in a simply connected region can be written in the form where $`𝐀=0`$. In a case of a gyraton, because of the presence of a sigularity at $`𝐱=0`$, the gauge transformations (5) cannot be used to banish the potential $`𝐀`$ globally. The situation here is similar to the well known Aharonov-Bohm effect AB ; Ham . The topological invariant $`j(u)`$, which has the meaning of the density of the angular momentum of the gyraton, is similar to the magnetic flux in the Aharonov-Bohm case.
In study of the Aharonov-Bohm effect it is usually helpful to consider at first a tube of finite radius where the magnetic field is localized. Similarly, in order to obtain well defined and finite expression for the gyraton metric, one must consider a spreaded source of finite size. In the 4-dimensional case the procedure proposed in section II does not work. We describe now a simple model of an extended gyraton in the 4-dimensional case. Let us modify the expression for (74) as follows
$$A_r=0,A_\varphi =\frac{\kappa }{4\pi }j(u)\left[\frac{r^2}{r_0^2}\vartheta (r_0r)+\vartheta (rr_0)\right].$$
(78)
For this modified vector potential the field strength $`F_{ab}=e_{ab}F`$ is
$$F=\frac{\kappa }{2\pi }\frac{j(u)}{r_0^2}\vartheta (r_0r).$$
(79)
In other words, the field strength $`F`$ is constant inside a disc of radius $`r_0`$. Outside this disc the field vanishes, while the contour integral (72) is $`j(u)`$ as earlier.
The function $`\psi `$ for such an extended gyraton can be found by using (37) and (40). In polar coordinates one has
$$\psi =\frac{F^2}{4\pi }_0^{r_0}𝑑r^{}r^{}Q,$$
(80)
$$Q=_0^{2\pi }𝑑\varphi \mathrm{ln}(r^2+r_{}^{}{}_{}{}^{2}2rr^{}\mathrm{cos}\varphi ).$$
(81)
For $`r>r_0`$ one has $`Q=4\pi \mathrm{ln}r`$. Thus
$$\psi =\frac{1}{2}F^2r_0^2\mathrm{ln}r.$$
(82)
It means that outside the gyraton $`\psi `$ has the same form as $`\phi _0`$ and can be absorbed into the latter by renormalizing the function $`\epsilon (u)`$.
### IV.2 Boosted 4-$`D`$ NUT metric
As an aside, it is worth mentioning that the Aichelburg-Sexl boost AiSe of the NUT stationary vacuum geometry NUT is a particular member of the class (77).
A convenient symmetric form of the NUT metric is
$$ds^2=\frac{dr^2}{f(r)}+(r^2+a^2)(d^2\theta +\mathrm{sin}^2\theta d\varphi ^2)$$
$$f(r)(dt2a\mathrm{cos}\theta d\varphi )^2,$$
(83)
where
$$f(r)=\frac{r^22mra^2}{r^2+a^2}.$$
(84)
It is known Bonnor ; MaRu:05 that the (singular) source of this geometry consists of a pair of semi-infinite line sources along the axis ($`\theta =0`$ and $`\theta =\pi `$ respectively), endowed with equal and opposite average angular momenta $`\pm a/2`$ per unit length, and joined to a massive particle at the origin. The two line sources are massless to linear order in $`a`$, i.e., to within terms of the order of the gravitational potential energy, which cannot be localized unambiguously.
To boost the metric (83), it is sufficient to consider its linearized form
$$ds^2=(1+\frac{2m}{r})(d\rho ^2+\rho ^2d\varphi ^2+d\overline{z}^2)$$
$$+4a\frac{\overline{z}}{r}d\varphi d\overline{t}(1\frac{2m}{r})d\overline{t}^2,$$
(85)
where $`r^2=\rho ^2+\overline{z}^2`$. We apply a Lorentz transformation
$$\overline{z}=\frac{1}{2}(ve^\chi ue^\chi ),\overline{t}=\frac{1}{2}(ve^\chi +ue^\chi ),$$
(86)
where $`u=tz`$ and $`v=t+z`$. In the limit $`\chi \mathrm{}`$ this sends the originally static source moving along the path $`z=t`$ in the new frame. Noting that $`lim_\chi \mathrm{}\overline{z}/r=ϵ(u)=\pm 1`$, rescaling mass and angular momentum according to
$$m=\mu e^\chi ,a=\alpha e^\chi ,$$
(87)
and using the distributional identity
$$\underset{\chi \mathrm{}}{lim}\frac{e^\chi }{\sqrt{\rho ^2+u^2e^{2\chi }}}=\frac{1}{|u|}\delta (u)\mathrm{ln}|\rho /l|,$$
(88)
where $`l`$ is an arbitrary length scale, we readily obtain the limiting form
$$ds^2=d\rho ^2+\rho ^2d\varphi ^2dudv2\mu \delta (u)\mathrm{ln}|\rho /l|du^22\alpha ϵ(u)dud\varphi .$$
(89)
(In (89), we have absorbed the term $`1/|u|`$ from (88) by a transformation of $`v`$.)This is a special case of (77). It represents a pair of semi-infinite gyratons with equal and opposite angular momentum densities $`\alpha ϵ(u)`$, joined to the Aichelburg-Sexl boosted particle of energy $`\mu `$. Classification of the 4-dimensional pp-waves with an impulsive profile based on their symmetries can be found in AiBa .
## V $`5`$dimensional gyratons
### V.1 General solution
In order to obtain a solution for the gravitational field of a 5-dimensional gyraton one needs to analyze electro- and magnetostatics in a flat 3-dimensional space.
Let us consider first the ‘magnetic’ equation (31). Using the standard 3-dimensional notations one can write these equations in the form
$$𝐁=\text{curl}𝐀,\text{curl}𝐁=0.$$
(90)
The second equation implies that there exists a function $`\mathrm{{\rm Y}}`$, the magnetic scalar potential, such that the magnetic field $`B`$ is
$$𝐁=\mathrm{{\rm Y}}.$$
(91)
The first of the equations (90) implies that the magnetic potential obeys the following equation
$$\mathrm{{\rm Y}}=0.$$
(92)
Let $`(r,\theta ,\varphi )`$ be the spherical coordinates
$$x^3+ix^4=r\mathrm{sin}\theta e^{i\varphi },x^5=r\mathrm{cos}\theta .$$
(93)
Then the general solution of (92) decreasing at infinity can be written as follows
$$\mathrm{{\rm Y}}=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}a_{lm}\frac{Y_{lm}(\theta ,\varphi )}{r^{l+1}},$$
(94)
where the complex coefficients $`a_{lm}`$ obey the conditions $`\overline{a}_{lm}=a_{lm}`$. Here $`Y_{lm}(\theta ,\varphi )`$ are spherical harmonics
$$Y_{lm}(\theta ,\varphi )=\sqrt{\frac{(2l+1)(lm)!}{4\pi (l+m)!}}P_l^m(\mathrm{cos}\theta )e^{im\varphi }.$$
(95)
The magnetic induction vector $`𝐁`$ is
$$𝐁=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}a_{lm}\left(\frac{Y_{lm}}{r^{l+1}}\right).$$
(96)
In order to find the corresponding vector-potential $`𝐀`$ one needs to solve the following equation
$$\text{curl}𝐀=\mathrm{{\rm Y}}.$$
(97)
It can be done by using the properties of vector spherical harmonics. Let us denote
$$𝚿_{lm}(\theta ,\varphi )=rY_{lm}(\theta ,\varphi ),$$
(98)
$$𝚽_{lm}(\theta ,\varphi )=𝐫\times Y_{lm}(\theta ,\varphi ).$$
(99)
The vector spherical harmonics obey the following relations BaEsGi
$$\times \left(\frac{𝚽_{lm}}{r^{l+1}}\right)=\times \left(\frac{𝐫\times 𝚿_{lm}}{r^{l+2}}\right).$$
(100)
$$\left(\frac{𝚽_{lm}}{r^{l+1}}\right)=0.$$
(101)
Using the first of these relations one finds
$$𝐀=𝐀_0\underset{l=1}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\frac{a_{lm}}{l}\frac{𝚽_{lm}}{r^{l+1}}.$$
(102)
The relation (101) shows that the solution (102) obeys the following gauge condition
$$\text{div}𝐀=0.$$
(103)
We denote by $`𝐀_0`$ a vector potential for $`l=0`$ case which requires a special treatment, since in this case $`𝚽_0=0`$ and ratio $`𝚽_{lm}/l`$ is not determined.
A general solution of the equation (34) for $`\phi `$ can be written as
$$\phi =\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}b_{lm}\frac{Y_{lm}(\theta ,\varphi )}{r^{l+1}},$$
(104)
where the coefficients $`b_{lm}`$ obey the conditions $`\overline{b}_{lm}=b_{lm}`$. For a gyraton solution coefficients $`a_{lm}`$ and $`b_{lm}`$ are arbitrary functions of the retarded time $`u`$. To obtain $`\psi `$ one can use (40) with the
$$𝒢(𝐱,𝐱^{})=\frac{1}{4\pi |𝐱𝐱^{}|}.$$
(105)
### V.2 Monopole solution
As we mentioned, the case of a magnetic monopole ($`l=0`$) is special. Let us consider it in more details. The magnetic potential $`\mathrm{{\rm Y}}`$ for the magnetic monopole is
$$\mathrm{{\rm Y}}=\frac{\mu }{r},$$
(106)
where $`\mu `$ is an arbitrary function of $`u`$. The magnetic induction vector has components
$$B_r=\frac{\mu }{r^2},B_\theta =B_\varphi =0.$$
(107)
The corresponding vector potential is of the form
$$A_r=A_\theta =0,A_\varphi =\mu \mathrm{cos}\theta .$$
(108)
The potential obeys the condition $`\text{div}𝐀=0`$ and the potential $`\psi `$ is
$$\psi =\frac{\mu ^2}{4r^2}.$$
(109)
The corresponding monopole solution for the gyraton is
$$ds^2=2dudv+dr^2+r^2d\theta ^2+r^2\mathrm{sin}^2\theta d\varphi ^2$$
$$+(\phi +\frac{\mu ^2(u)}{4r^2})du^22\mu (u)\mathrm{cos}\theta dud\varphi .$$
(110)
Here $`\phi =\phi (u,r,\theta ,\varphi )`$ is a solution of the equation (34).
Similarly to the 4-D case this metric is related to the boosted NUT-like metric. Consider a metric ($`R^2=r^2+w^2`$)
$$ds^2=\left(1\frac{2m}{R^2}\right)dt^2+4a\mathrm{cos}\theta d\varphi dt+$$
$$\left(1+\frac{2m}{R^2}\right)(dr^2+r^2d\theta ^2+r^2\mathrm{sin}^2\theta d\varphi ^2+dw^2).$$
(111)
It is Ricci-flat to linear order and it is a linearized version of 5-D NUT spacetime. (An exact counterpart of this linearized metric does not seem to be known.) Applying the boost
$$w=\frac{1}{2}(ue^\chi ve^\chi ),t=\frac{1}{2}(ue^\chi +ve^\chi )$$
(112)
rescaling $`me^{2\chi }=\frac{1}{2}\mu ^2`$, $`ae^\chi =\mu `$, and noting
$$\underset{\chi \mathrm{}}{lim}R^2=r^2,$$
(113)
we recover (110) with $`\phi =1`$ in the limit $`\chi \mathrm{}`$.
### V.3 Dipole solution
The spherical harmonics for $`l=1`$ case are
$$Y_{10}=\alpha \frac{x^5}{r},\alpha =\sqrt{\frac{3}{4\pi }}$$
(114)
$$Y_{11}=\overline{Y}_{11}=\frac{\alpha }{\sqrt{2}}\frac{x^3+ix^4}{r}.$$
(115)
Notice that $`𝐫\times F(r)=0`$. Using this property we obtain
$$𝚽_{10}=\frac{\alpha }{r}(x^4,x^3,0),$$
(116)
$$𝚽_{11}=i\frac{\alpha }{\sqrt{2}r}(x^5,ix^5,x^3+ix^4).$$
(117)
Let us denote
$$\omega ^{ab}=x^adx^bx^bdx^a.$$
(118)
Then the expressions for $`(𝐀_{1m},d𝐱)`$ take the form
$$(𝐀_{10},d𝐱)=a_{10}\frac{\alpha }{r^3}\omega ^{34},$$
(119)
$$(𝐀_{11},d𝐱)=a_{11}\frac{\alpha }{r^3}(\omega ^{45}i\omega ^{35}).$$
(120)
The vector potential for a general dipole solution can be written as follows
$$(𝐀,d𝐱)=\frac{\kappa }{8\pi }\frac{j_{ab}x^bdx^a}{r^3},$$
(121)
where
$$j_{ab}x^bdx^a=\frac{8\pi \alpha }{\kappa }\left[a_{10}\omega ^{34}+\sqrt{2}\mathrm{}(a_{11})\omega ^{45}+\sqrt{2}\mathrm{}(a_{11})\omega ^{35}\right],$$
(122)
and $`a_{10}`$, $`a_{11}`$ are arbirary functions of $`u`$.
## VI Higher dimensional case
Let us discuss first the scalar (electrostatic) equation (34) in the $`n`$dimensional Euclidean space $`R^n`$
$$\phi =0.$$
(123)
To solve this equation it is convenient to decompose the potential $`\phi `$ into the scalar spherical harmonics RuOr
$$Y^l=r^l𝒴^l,𝒴^l=C_{c_1\mathrm{}c_{l1}}x^{c_1}\mathrm{}x^{c_{l1}},$$
(124)
where $`C_{c_1\mathrm{}c_{l1}}`$ is a symmetric traceless rank-$`l`$ tensor. It is easy to see that the number of linearly independent components of coefficients $`C_{c_1\mathrm{}c_{l1}}`$ is
$$d_0(n,l)=\frac{(l+n3)!(2l+n2)}{l!(n2)!}.$$
(125)
These harmonics are eigenfunctions of the invariant Laplace operator on a unit sphere $`S^{n1}`$ with eigenvalues $`l(n+l2)`$. For each $`l`$ there exists $`d_0(n,l)`$ linearly independent harmonics. We shall use an index $`q`$ to enumerate the independent harmonics. The functions $`Y^{lq}`$ form a complete set, so that any smooth function $`F`$ on $`S^{n1}`$ can be decomposed as
$$F=\underset{l=0}{\overset{\mathrm{}}{}}\underset{q}{}F_{lq}Y^{lq}.$$
(126)
Consider now a special mode $`F_{lq}(r)Y^{lq}`$. It is a decreasing-at-infinity solution of (123) if $`F_{lq}r^{(n+l2)}`$. This can be proved by using the properties of the scalar spherical harmonics. We demonstrate this directly by using the relations (124).
First, it is easy to check that
$$𝒴^l=0,x^d_d𝒴^l=l𝒴^l.$$
(127)
Using these relations one obtains
$$(f(r)𝒴^l)=(f^{\prime \prime }+\frac{(n+2l1)}{r}f^{})𝒴^l.$$
(128)
Thus for $`f=1/r^{n+2l2}`$ the mode functions $`f(r)𝒴^l`$ obey the equation (123). To summarize, a general solution of the electrostatic equation (123) can be written in the form
$$\varphi =\underset{l=0}{\overset{\mathrm{}}{}}\underset{q}{}\frac{𝒴^{lq}}{r^{n+2l2}}.$$
(129)
In the gyraton solution (1) $`d_0(n,l)`$ independent components of $`C_{c_1\mathrm{}c_{l1}}`$ are arbitrary functions of $`u`$.
In a similar way, one can obtain solutions of the equations of magnetostatics in $`n`$dimensional Euclidean space by using the vector spherical harmonics RuOr . Let us denote
$$A_a^l=f(r)𝒴_a^l,$$
(130)
$$𝒴_a^l=C_{abc_1\mathrm{}c_{l1}}x^bx^{c_1}\mathrm{}x^{c_{l1}}.$$
(131)
Here $`C_{abc_1\mathrm{}c_{l1}}`$ is a $`(l+1)`$-th-rank constant tensor which possesses the following properties: it is antisymmetric under interchange of $`a`$ and $`b`$, and it is traceless under contraction of any pair of indices RuOr .
First, let us demonstrate that $`A_a^l`$ obeys the gauge condition
$$^aA_a^l=0.$$
(132)
Notice that
$$_af(r)=f^{}(r)\frac{x^a}{r}.$$
(133)
Thus
$$^aA_a^l=f^a𝒴_a^l=0.$$
(134)
The latter equality follows from the fact that when $`_a`$ is acting on one of $`x`$ it effectively produces a contraction of two indices in $`C`$ which vanishes.
In the gauge (132) the magnetostatic field equations (31) reduce to the following equation
$$A_a^l=0.$$
(135)
It is easy to get
$$𝒴_a^l=0,x^b_b𝒴_a^l=l𝒴_a^l.$$
(136)
Using these relations one obtains
$$(f𝒴_a^l)=(f^{\prime \prime }+\frac{n+2l1}{r}f^{})𝒴_a^l.$$
(137)
Hence $`𝒴_a^l`$ is a solution of (135) if
$$f^{\prime \prime }+\frac{n+2l1}{r}f^{}=0.$$
(138)
Solving this equation we get $`f=1/r^{n+2l2}`$. Hence a general decreasing at infinity solution of the magnetostatic equations in the $`n`$dimensional space (31) can be written as
$$A_a=\underset{l=1}{\overset{\mathrm{}}{}}\underset{q}{}\frac{𝒴_a^{lq}}{r^{n+2l2}}.$$
(139)
Again, we use an index $`q`$ to enumerate different linearly independent vector spherical harmonics. The total number of these harmonics for given $`l`$ is RuOr
$$d_1(n,l)=\frac{l(n+l2)(n+2l2)(n+l3)!}{(n3)!(l+1)!}.$$
(140)
In the gyraton solution (1) the coefficients $`C_{abc_1\mathrm{}c_{l1}}`$ in the decomposition (139) are arbitrary functions of the retarded time $`u`$. For a given solution $`𝐀`$ relation (40) allows one to find $`\psi `$.
## VII Summary and discussions
The main result of this paper is that the vacuum Einstein equations for the gyraton-type metric (1) in an arbitrary number of spacetime dimensions $`D`$ can be reduced to linear problems in the Euclidean $`(D2)`$dimensional space. These problems are: (1) To find a static electric field $`\phi `$ created by a point-like source; (2) To find a magnetic field $`𝐀`$ created by a point-like source. The retarded time $`u`$ plays the role of an external parameter. One can include $`u`$-dependence by making the coefficients in the harmonic decomposition for $`\phi `$ and $`𝐀`$ to be arbitrary functions of $`u`$. After choosing the solutions of these two problems one can define $`\psi `$ by means of equation (40). By substituting $`\mathrm{\Phi }=\phi +\psi `$ and $`𝐀`$ into the metric ansatz one obtains a vacuum solution of the Einstein equations.
Such a gyraton-like solution has a singularity located at the spatial point $`𝐱=0`$ during some interval of the retarded time $`u`$. It means that the corresponding point-like source is moving with the velocity of light. Energy $`E`$ and angular momentum $`J_{ab}`$ are finite. It was demonstrated that for given energy and angular momentum the gyraton can also have other characteristics, describing the deviation of $`\mathrm{\Phi }`$ from spherical symmetry (in the transverse space $`R^n`$) and the presence of higher than dipole terms in the multipole expansion of $`𝐀`$. One can interpret such solutions as excitations or distortions of the gyraton solutions obtained earlier in FrFu:05 .
It should be emphasized that the point-like sources are certainly an idealization. In FrFu:05 it was shown that gyraton solutions can describe the gravitational field of beam-pulse spinning radiation. In such a description one uses the geometric optics approximation. For its validity the size of the cross-section of the beam must be much larger than the wave-length of the radiation. In the presence of spin $`J`$ one can expect additional restrictions on the minimal size of both, the cross-section size and the duration of the pulse. As usual in physics, one must have in mind that in the possible physical applications the obtained solution is valid only outside some region surrounding the immediate neighborhood of the singularity.
The gyraton solutions might be used for study the gravitational interaction of ultrarelativistic particles with spin. The gyraton metrics might be also interesting as possible exact solutions in the string theory.
## Acknowledgments
The authors are grateful to Alan Coley and Dmitri Fursaev for stimulating discussions and remarks. This work was supported by the Natural Sciences and Engineering Research Council of Canada and by the Killam Trust. The authors also kindly acknowledge the support from the NATO Collaborative Linkage Grant (979723).
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# Calculation of the Masses and the Running Masses of the Quarks and Leptons from Electroweak to Supersymmetric Grand Unification Mass
## I Introduction
Pati and Salama1 pioneered the idea that leptons are the fourth colour; quarks and leptons should be brought under the same umbrella of one group so that all forces ( except gravity) can be understood in terms of one unifying force parameter near the Planck scale. The six leptons should be treated at par with the six coloured quarks. The elementary constituents of matter would become twelve only. With the availability of enormous data from the high energy accelerators, phenomenological analysis backed by imaginative theories, it has been found that the familiar standard model, which is the product group $`SU_C(3)SU_L(2)U_Y(1)`$, can explain all of them successfully. This standard model is characterised by three coupling strengths of the weak, the electromagnetic and the strong interactions. It was conceived that all the three couplings should run, i.e., they change with energy, and eventually become one at the grand unified scale. This could be made possible by solving the relevant RGE. The coefficients of the theory are specified by the group structure alone. Many theoretical models were investigated and only a few years back, it has become essential to bring in supersymmetry. In the minimal version of the MSSM, the beta coefficients are such that they run the coupling strengths to a single value $`\frac{4\pi }{g_U^2}=\alpha _{GUT}\frac{1}{25}`$ at a mass of $`M_X=2.2\times 10^{16}`$ GeV. This has been the most attractive result in the current investigations in gauge theoriesa2 .
Next in order, the major challenge in particle physics, was and is the theoretical derivation of the mass spectrum of the quarks and leptons in the same successful MSSM theory. In this model, all the masses of the fermions and the mixing angles were being chosen arbitrarily to account for the 19 free parameters of the theory. In the absence of such a fundamental theory, it has been in vogue to pursue the method which has been known as ‘textures analysis’. After finding a suitable texture, one can investigate further and possibly obtain unification mass parameters for all the fermions as a generalisation of the hypothesis by Georgi et ala3 . Eventually, one can predict their individual masses by using the analytical MSSM group coefficients for a RGE for the masses of the fermions. As yet, this has not been successful.
One of the 13 parameters of the standard model was first predicted in 1974 by Gaillard and Leea4 . Then came the work of the popular mass matrix ansatz of Fritzscha5 . A complete listing of textures and their relevance to experimental findings were made by Ramond, Roberts and Ross ( RRR)a6 . Since they also took the help of RGE, we shall present the essential results of their work in this section. Before that, we present the MSSM Lagrangian
$$_M=\overline{Q}_LM_U\mathrm{\Phi }_uU_R+\overline{Q}_LM_D\mathrm{\Phi }_dD_R+\overline{l}_LM_E\mathrm{\Phi }_dE_R+\overline{\nu }_LM_N\mathrm{\Phi }_u\nu _R+h.c.$$
(1)
We ignore the sparticles and consider the rigid part of the Lagrangian given by Demir Demir04 .
In the renormalisation schemes, the Yukawa couplings $`M_F(t)`$ and v.e.v. $`v_F(t)`$ change with energy separately. The Dirac masses are given by
$$m_F(t)=v_F(t)M_F(t)$$
(2)
The masses considered above are not the masses of the flavor eigenstates of the model.
The one loop RGE, as written by Grzadkowski, Lindner and Theisen a9 for $`M_Y(t)`$ are
$$16\pi ^2\frac{dM_Y}{dt}=\left(G_Y+T_Y+S_Y\right)M_Y(t),$$
(3)
$`t=log(\mu /M_Z)`$, $`\mu `$ is the renormalisation point and $`M_Z`$ is the mass of the Z-boson. Here Y=U stands for U-quarks (F=1,2,3): up(u), charm(c) and top(t); Y=D stands for Down quarks (F=4,5,6): bottom(b), strange(s) and down(d); Y=N stands for Neutrinos (F=7,8,9): $`\nu _e,\nu _\mu `$, and $`\nu _\tau `$, and Y=E stands for electrons(F=10,11,12): e, $`\mu `$ and $`\tau `$. $`G_Y(t)`$ contains the gauge coupling terms, given in Section-2 and
$$T_U=Tr\left(3M_UM_U^{}+M_NM_N^{}\right),T_D=Tr\left(3M_DM_D^{}+M_EM_E^{}\right),T_U=T_N,T_E=T_D,$$
(4)
and
$$S_U=3M_UM_U^{}+M_DM_D^{},S_D=3M_DM_D^{}+M_UM_U^{},S_E=3M_EM_E^{}+M_NM_N^{},S_N=3M_NM_N^{}+M_EM_E^{}.$$
(5)
To find the couplings, we have to solve the twelve differential equations and determine all the couplings/masses from only one value of coupling at $`t=t_X`$ or $`M_F(M_X)`$. First, we turn our attention to the mass of the top quark. Incidentally, Pendleton and Ross pendleton , Faraggifaraggi , and some others have predicted the value of the mass of the top quark which was around 175 GeV, even before the top was discovered.
In this paper, all the one loop equations are solved following the same method used for gauge coupling RGE. A particular $`M_F(M_Z)`$ is obtained relating to $`M_F(M_X)`$ and other fermions. Considering the top, first we make ’heavy top integral’ approximation stated below. Eventhough we aim at single input value $`M_F(M_X)=M_U`$ , which is independent of F, the integrals of the solutions for the masses are such that
$$M_{top}^2(\tau )𝑑\tau M_{Qtop}^2(\tau )𝑑\tau M_{lepton}^2(\tau )𝑑\tau $$
(6)
Neglect of the integrals other than the top quark, gives a very good result for the $`M_{top}(M_X)M_U`$. Following this approximation procedure, we find from other eleven equations that $`M_{top}(M_X)M_F(M_X)M_U`$ 115 GeV for all the fermions.
We choose $`M_U`$ 115 GeV as the only input. Furthermore, MSSM has two Higgs.
$`<\varphi _U>=v_u(t)`$ $`=`$ $`v(t)Sin\beta (t)`$ (7)
$`<\varphi _D>=v_d(t)`$ $`=`$ $`v(t)Cos\beta (t)`$ (8)
$`v^2(t)`$ $`=`$ $`v_d^2(t)+v_u^2(t)`$ (9)
$`tan\beta (t)`$ $`=`$ $`{\displaystyle \frac{v_u(t)}{v_d(t)}}.`$ (10)
It is to be noted that $`v(t)`$ is the vacuum expectation value of single Higgs of the Standard Model. $`v(M_Z)`$=246 GeV, $`Sin\beta _{SM}(M_Z)=\frac{1}{\sqrt{2}}`$. We note that $`v_0=\frac{246}{\sqrt{2}}`$ GeV=174 GeV, which is close to $`M_{top}(M_Z)175`$ GeV. For simplicity, we shall use this as the unit of energy whenever unspecified.
The one loop RGE are
$`16\pi ^2{\displaystyle \frac{dv_u}{dt}}`$ $`=`$ $`\left({\displaystyle \frac{3}{20}}g_1^2+{\displaystyle \frac{3}{4}}g_2^2Tr(3M_UM_U^{})\right)v_u,`$ (11)
$`16\pi ^2{\displaystyle \frac{dv_d}{dt}}`$ $`=`$ $`\left({\displaystyle \frac{3}{20}}g_1^2+{\displaystyle \frac{3}{4}}g_2^2Tr(3M_DM_D^{})\right)v_d.`$ (12)
The calculation of
$$tan\beta (t)=\frac{v_u(t)}{v_d(t)},$$
(13)
in MSSM has attracted considerable attention, there exist extensive literature, most of them are given in referenceparida . We note that below the mass scale $`M_Z`$ of the normal SM, one Higgs v.e.v. satisfies the one loop equation
$$16\pi ^2\frac{dv_{SM}}{dt}=\left(\frac{3}{20}g_1^2+\frac{3}{4}g_2^2Tr(3M_UM_U^{}+3M_DM_D^{})\right)v_{SM}.$$
(14)
To maintain continuity, we assume that at $`t`$=0,
$$tan\beta (M_Z)=tan\beta _{SM}(M_Z)=1.$$
(15)
Using the approximation given in equation (6), we get from equations (11), (12) and (13),
$$tan\beta (t)exp\left(\frac{3}{16\pi ^2}_0^t\left[M_{top}^2(\tau )M_{bottom}^2(\tau )\right]𝑑\tau \right).$$
(16)
Using the results given in section-3, equation(113), we find that $`tan\beta (t)`$ is a slowly varying function of $`t`$ and drops from 1 at 91 GeV to 0.9 at $`10^{16}`$ GeV. So the calculation is much simplified if we take $`Sin\beta =\frac{1}{\sqrt{2}}=Cos\beta `$. This is not ruled out by experimenta9 .
The authors (RRR), noted that there are only six possible forms of symmetric mass matrices with three non-zero eigenvalues and three texture zeroes, capable of describing the hierarchy of up or down quark mass matrices. Those are
$`T_1=\left(\begin{array}{ccc}a_1& 0& 0\\ 0& b_1& 0\\ 0& 0& c_1\end{array}\right),T_2=\left(\begin{array}{ccc}0& a_2& 0\\ a_2& b_2& 0\\ 0& 0& c_2\end{array}\right),T_3=\left(\begin{array}{ccc}a_3& 0& 0\\ 0& 0& b_3\\ 0& b_3& c_3\end{array}\right),`$ (26)
$`T_4=\left(\begin{array}{ccc}0& 0& a_4\\ 0& b_4& 0\\ a_4& 0& c_4\end{array}\right),T_5=\left(\begin{array}{ccc}0& a_5& 0\\ a_5& 0& b_5\\ 0& b_5& c_5\end{array}\right)\text{and}T_6=\left(\begin{array}{ccc}0& a_6& b_6\\ a_6& 0& 0\\ b_6& 0& c_6\end{array}\right).`$ (36)
$`T_5`$ was the one first pioneered by Fritzscha5 . More recently Dimopoulos, Hall and Rabya7 have analysed and included the leptons following Georgia3 and solved the MSSM RGE with some degree of success. RRR tried to analyse all the cases and put them in a CKM matrix form, proposed by Wolfensteina8 and diagonalise them to the texture types.
$`V_{CKM}`$ $`=`$ $`\left(\begin{array}{ccc}c_1c_2s_1s_2e^{i\varphi }& s_1+c_1s_2e^{i\varphi }& s_2(s_3s_4)\\ c_1s_2s_1e^{i\varphi }& s_1s_2+(c_1c_2c_3c_4+s_1s_4)e^{i\varphi }& s_3s_4\\ s_1(s_3s_4)& c_1(s_3s_4)& (c_3c_4+s_1s_4)e^{i\varphi }\end{array}\right)`$ (40)
$`=`$ $`\left(\begin{array}{ccc}\lambda ^2/2& \lambda & A\lambda ^3(\rho +i\eta )\\ \lambda & 1\lambda ^2/2& A\lambda ^2\\ A\lambda ^3(1\rho +i\eta )& A\lambda ^2& 1\end{array}\right).`$ (44)
Here $`s_i,c_i`$ (i=1,..,4)are the sines and cosines of mixing angles. From the identities given by Dimopoulos, Hall and Rabya7 , following from equations (40) and (44)
$$\lambda =(s_1^2+s_2^2+s_1s_2Cos\varphi )^{1/2}.$$
(45)
$`\varphi `$ is the CKM phase angle and
$$s_1=\left(\frac{M_d}{M_s}\right)^{\frac{1}{2}}=\lambda ,s_2=\left(\frac{M_u}{M_c}\right)^{\frac{1}{2}}=\lambda ^2;s_4=\left(\frac{M_dM_s}{M_b^2}\right)^{\frac{1}{2}}=\lambda ^3.$$
(46)
As will be discussed later, the $`Cos\varphi `$ defined in a7 , we shall get
$$Cos\varphi =\lambda /2,$$
(47)
and
$$s_3s_4=\lambda ^2A(t)$$
(48)
where A(t) depends on $`t`$.
The small expansion parameter is $`\lambda `$0.2 and $`A0.9\pm 0.1`$. Olechowski and Poroskia9 were the first to write down the RG Equations for the parameters
$`16\pi ^2{\displaystyle \frac{d|J_{c\rho }|}{dt}}`$ $`=`$ $`3c(h_t^2+h_b^2)|J_{c\rho }|,`$ (49)
$`16\pi ^2{\displaystyle \frac{dA}{dt}}`$ $`=`$ $`{\displaystyle \frac{3}{2}}c(h_t^2+h_b^2)A,`$ (50)
$`{\displaystyle \frac{d\lambda }{dt}}`$ $`=`$ $`0,`$ (51)
$`{\displaystyle \frac{d\rho }{dt}}`$ $`=`$ $`0,`$ (52)
$`{\displaystyle \frac{d\eta }{dt}}`$ $`=`$ $`0.`$ (53)
Here $`c=2/3`$ for MSSM and $`J_{c\rho }`$ is the irreducible phase of the CKM matrix. To arrive at equation (50), we equate
$$A(t)=\left[\frac{M_{bottom}(t)M_{top}(t)}{M_{top}^2(M_X)}\right]^{1/7}.$$
(54)
Using one loop RG equations (3) to (5), we get
$$A(M_X)=1.00\text{and}A(M_Z)=1.474.$$
(55)
RRR have made a complete listing and analysis. They arrived at the value $`\lambda =`$=0.22 and the oft quoted result,
$`m_\tau :m_\mu :m_e=m_b:m_s:m_d=1:\lambda ^2:\lambda ^4`$
$`m_t:m_c:m_u=1:\lambda ^4:\lambda ^8,`$ (56)
which is very well satisfied by experimentally found masses. This has been thought to be like a miracle.
We shall take the masses of the twelve fermions to be real, wherever necessary. In section-2, we shall write the RG Equation with the MSSM coefficients and give the one loop exact solutions. In section-3, the unification mass for all fermions at GUT scale ranging from 113 GeV to 125 GeV as computed by Deo and Maharanaa10 , will be discussed as a follow up of their letter. An expression for Wolfenstein parameter in terms of RGE coefficients is given in section-4. The reason to raise this parameter by integers $`n_F`$ to obtain fermion masses except the top, $`M_F`$ equaling $`\lambda ^{n_F}`$ multiplied by the parametric unification mass 115 GeV, is given. In sections-5,6 and 7, we suggest an alternative method of calculation of experimental masses of all fermions by a suitable self contained procedure which includes the gauge couplings as well. The equation for the running of all the masses of all the fermions is given in section-8. The results are given in tables and are also shown graphically for greater clarity. The concluding remarks are given in the section-9.
## II Renormalisation Group Equations and Solutions
As stated, the gauge sector of the standard model is characterised by three coupling constants $`g_3`$, $`g_2`$ and $`g_1`$ of $`SU_C(3)SU_L(2)U_Y(1)`$, respectively. However, these couplings are not constants, they change with energy/mass values. The nature of variation is given by the solutions of the RG Equations. The coefficients are calculated by the specific nature of the Standard Model group. For MSSM, the three couplings at mass $`M_X=2.2\times 10^{16}`$ GeVa11 unite to a unified coupling constant $`g_U^2/4\pi `$=1/24.6. Here, the supersymmetry descends from $`M_X`$ down to $`M_Z`$ 91 GeV, as suggested by Witten. The RG Equations for the couplings in the lowest order are given by
$$16\pi ^2\frac{dg_i(t)}{dt}=c_ig_i^3(t),i=1,2,3.$$
(57)
The coefficients are $`c_1`$ = 6.6, $`c_2`$ =1, $`c_3`$ = -3. The first two coefficients are positive, indicating that $`U_Y(1)`$ and $`SU_L(2)`$ are not asymptotically free, whereas the $`SU_C(3)`$ colour group is free and makes the entire product group asymptotically free . This implies that in a perturbative formulation, the higher order contributions are small and can be neglected. Therefore, we shall use equation(57) only in the gauge sector, with two Higgsa11 .
The running parameter $`t`$ is defined as $`t=log_e\frac{\mu }{M_Z}`$ so that it varies from 0 to $`log_e\frac{M_X}{M_Z}33`$. The solution to RG Equation(57) is
$$\frac{4\pi }{g_i^2(t)}=\frac{4\pi }{g_i^2}\frac{c_i}{2\pi }t.$$
(58)
Here, $`g_i^2=g_i^2(0)`$ are the coupling strengths in the electroweak scale $`M_Z`$. Taking the value of $`M_X=2.2\times 10^{16}`$ GeV and $`\frac{4\pi }{g_U^2}=24.6`$, we calculate the values of $`\frac{4\pi }{g_1^2}=59.24`$, $`\frac{4\pi }{g_2^2}=29.85`$ and $`\frac{4\pi }{g_3^2}=8.85`$. These are consistent with the experimental results. Thus the three coupling strengths are descendants of one coupling constant $`g_U`$.
Taking the clue from equation(57), we rewrite the Yukawa sector SUSY RG Equations given earlier in equation (3), which had also been written by Babua12 following Georgi and Glashow, and Eitchen et ala13 in the following way;
$`16\pi ^2{\displaystyle \frac{dM_F(t)}{dt}}`$ $`=`$ $`A_FM_F^3(t)+[Y_F(t)G_F(t)]M_F(t)`$ (59)
$`=`$ $`A_FM_F^3(t)+Z_F(t)M_F(t).`$ (60)
The masses are in units of 175 GeV. We repeat for ready reference that the 12 fermions are suffixed as $`F`$ = 1,2,… 12. $`F`$ = 1,2,3 are the $`U`$-quarks: top($`t`$), charm($`c`$) and up($`u`$). Similarly, $`F`$=4,5,6 denote the $`D`$-quarks: bottom($`b`$), strange($`s`$) and down($`d`$); $`F`$=7,8,9 are the $`E`$-leptons: electron($`e`$), muon($`\mu `$) and tau($`\tau `$); and $`F`$=10,11,12 are the $`N`$-neutrinoes $`\nu _e,\nu _\mu ,\nu _\tau `$. Further, $`M_1=M_{top},M_2=M_{charm},\mathrm{},M_{12}=M_{tau}`$.
$`A_F`$ is a group theoretic factor whose value is ‘6’ for quarks, i.e., $`F`$=1,2,$`\mathrm{}`$,6 and ‘4’ for the leptons i.e., $`F`$=7,8,$`\mathrm{}`$ , 12. The positive values indicate the field theory containing Yukawa couplings only and may not be asymptotically free.
$`Y_F`$ is the mixing term which can be put in matrix form
$$Y_F=\underset{H}{}A_{FH}M_H^{}(t)M_H(t),H=1,2,\mathrm{},12.$$
(61)
In MSSM, the matrix $`A_{FH}`$ is specified by the 144 elements given below,
$`A_{FH}=\left(\begin{array}{cccccccccccc}0& 3& 3& 1& 0& 0& 0& 0& 0& 1& 1& 1\\ 3& 0& 3& 0& 1& 0& 0& 0& 0& 1& 1& 1\\ 3& 3& 0& 0& 0& 1& 0& 0& 0& 1& 1& 1\\ 1& 0& 0& 0& 3& 3& 1& 1& 1& 0& 0& 0\\ 0& 1& 0& 3& 0& 3& 1& 1& 1& 0& 0& 0\\ 0& 0& 1& 3& 3& 0& 1& 1& 1& 0& 0& 0\\ 0& 0& 0& 3& 3& 3& 0& 1& 1& 1& 0& 0\\ 0& 0& 0& 3& 3& 3& 1& 0& 1& 0& 1& 0\\ 0& 0& 0& 3& 3& 3& 1& 1& 0& 0& 0& 1\\ 3& 3& 3& 0& 0& 0& 1& 0& 0& 0& 1& 1\\ 3& 3& 3& 0& 0& 0& 0& 1& 0& 1& 0& 1\\ 3& 3& 3& 0& 0& 0& 0& 0& 1& 1& 1& 0\end{array}\right).`$ (74)
The diagonal elements of $`A_{FH}`$ have been taken out as the cubic term in equation(59); so they are zero.
As the model is minimal supersymmetric, the gauge factors $`G_F(t)`$, which are the sum of gauge couplings, are fixed, we take the values from referencea14 and a7 .
$$G_U(t)=\frac{13}{15}g_1^2(t)+3g_2^2(t)+\frac{16}{3}g_3^2(t)=\underset{i=1}{\overset{3}{}}K_U^ig_i^2(t).$$
(75)
$`F=`$1,2,3 stand for $`U`$ and they are degenerate electromagnetic gaugewise. Similarly,
$`G_D(t)`$ $`=`$ $`{\displaystyle \frac{7}{15}}g_1^2(t)+3g_2^2(t)+{\displaystyle \frac{16}{3}}g_3^2(t),F=4,5,6`$ (76)
$`G_E(t)`$ $`=`$ $`{\displaystyle \frac{9}{5}}g_1^2(t)+3g_2^2(t),F=7,8,9`$ (77)
$`andG_N(t)`$ $`=`$ $`{\displaystyle \frac{3}{5}}g_1^2(t)+3g_2^2(t),F=10,11,12.`$ (78)
Here $`K_N^3=K_E^3=0`$, as the leptons do not have the strong colour interaction. We shall need the integrals,
$$\frac{1}{8\pi ^2}_0^t𝑑\tau G_F(\tau )=\underset{i=1}{\overset{3}{}}\frac{K_i^F}{c_i}\mathrm{log}(1\frac{c_ig_i^2t}{8\pi ^2})$$
(79)
and
$`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^{t_X}}𝑑\tau G_F(\tau )`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{K_i^F}{c_i}}\mathrm{log}(1{\displaystyle \frac{c_ig_i^2t_X}{8\pi ^2}})`$ (80)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{K_i^F}{c_i}}\mathrm{log}{\displaystyle \frac{g_i^2}{g_U^2}}.`$ (81)
Deo and Maharanaa10 made a very important observation that equation(60), which has to be solved for a given fermion, does not contain the coeficients of the same mass in the matrix $`A`$ of equation (74). This fact has been overlooked by all previous authors. As a result, the calculational dedails become erroneous and the values obtained are unreliable. The present approach gives a hope of a simple method of entangling the mass due to finding the solution of 12 differentential equations. As such, the terms $`Z_F(t)`$ can be exponiented away. We introduce a subsidiary mass $`m_F(t)`$, through
$$M_F(t)=m_F(t)\mathrm{exp}\left(\frac{1}{16\pi ^2}_0^tZ_F(\tau )𝑑\tau \right),$$
(82)
such that $`M_F(M_Z)=m_F(M_Z)m_F(0)`$. They satisfy the equation
$$16\pi ^2\frac{dm_F(t)}{m_F^3}=A_F\mathrm{exp}\left(\frac{1}{8\pi ^2}_0^tZ_F(\tau )𝑑\tau \right)dt.$$
(83)
They look, astonishingly, similar to the gauge sector one loop RG equation (57) and can be solved exactly. Integrating equation(83) from $`M_Z`$ to $`M_X`$ i.e., from $`t`$=0 to $`t_X`$, we get
$$\frac{8\pi ^2}{m_F^2(M_Z)}=\frac{8\pi ^2}{m_F^2(M_X)}+A_F_0^{t_X}𝑑t\mathrm{exp}\left(\frac{1}{8\pi ^2}_0^tZ_F(\tau )𝑑\tau \right).$$
(84)
Putting back the exponential,
$$\frac{M_{top}^2(M_Z)}{M_F^2(M_Z)}=\frac{M_{top}^2(M_Z)}{M_F^2(M_X)}\mathrm{exp}\left(\frac{1}{8\pi ^2}_0^{t_X}Z_F(\tau )𝑑\tau \right)+\frac{A_F}{8\pi ^2}_0^{t_X}𝑑t\mathrm{exp}\left(\frac{1}{8\pi ^2}_0^tZ_F(\tau )𝑑\tau \right).$$
(85)
This is the exact one loop solution. $`M_F^2(M_Z)`$ is the mass of the fermions at $`M_Z`$.
The descent or ascent running of the masses from GUT $`M_X`$ to electroweak $`M_Z`$, can be obtained by integrating equation(83) from $`t=t_X`$ to $`t`$. The result which is not given in reference a10 is
$$\frac{8\pi ^2M_{top}^2}{M_F^2(t)}=8\pi ^2\frac{M_{top}^2}{M_F^2(M_X)}\mathrm{exp}\left(\frac{1}{8\pi ^2}_t^{t_X}Z_F(\tau )𝑑\tau \right)+A_F_t^{t_X}𝑑t_1\mathrm{exp}\left(\frac{1}{8\pi ^2}_t^{t_1}Z_F(\tau )𝑑\tau \right).$$
(86)
This is also one loop exact. By solving equations (85) and (86), we can find $`M_F(M_X)`$ and $`M_F(t)`$ repsectively.
## III Original mass of all fermions at $`M_X`$ ?
The gauge integrals over $`G_F(t)`$ are easily and acurately calculable. The most difficult task is to evaluate $`Y_F(t)`$. Even though, it does not contain $`M_F(t)`$, it is a sum of squares of moduli of the masses of all fermions. For example, from the matrix as given by equation(74),
$$Y_{top}(t)=3M_c^2(t)+3M_u^2(t)+M_b^2(t)+M_{\nu _e}^2(t)+M_{\nu _\mu }^2(t)+M_{\nu _\tau }^2(t).$$
(87)
There is mixing of six other fermions for the top.
In the ‘top heavy integral’ approximation, one retains only those terms containing $`M_{top}^2(t)`$ occuring in any integral with $`Y_F(t)`$ . Consider the top case. $`Y_{top}(t)`$ can be set equal to zero in the RG Equation for the top. The top mass is then given by a simple expression
$$1=\frac{M_{top}^2(M_Z)}{M_{top}^2(M_X)}C_{top}+D_{top},$$
(88)
where, using integrals (79) to (81),
$`C_{top}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}\left({\displaystyle \frac{g_i^2}{g_U^2}}\right)^{\frac{K_i^U}{c_i}}=0.086`$ (89)
and (90)
$`D_{top}`$ $`=`$ $`{\displaystyle \frac{6}{8\pi ^2}}{\displaystyle _0^{t_X=33}}𝑑t{\displaystyle \underset{i=1}{\overset{3}{}}}\left(1{\displaystyle \frac{g_i^2c_it}{8\pi ^2}}\right)^{\frac{K_i^U}{c_i}}=0.802`$ (91)
Putting these values, the original mass of the top, at an energy of $`2.2\times 10^{16}`$ GeV was, approximately,
$$M_{top}(M_X)=M_{top}(M_Z)\left(1D_{top}\right)^{1/2}C_{top}^{1/2}=114\text{GeV}.$$
(92)
This is Deo-Maharana result.
We can write a general formula for all fermions
$$\frac{M_{top}^2(M_Z)}{M_F^2(M_Z)}=\frac{M_{top}^2(M_Z)}{M_F^2(M_X)}C_F+D_F,$$
(93)
where
$`C_F`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}\left({\displaystyle \frac{g_i^2}{g_U^2}}\right)^{\frac{K_i^F}{C_i}}\mathrm{exp}\left({\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^{t_X}}Y_F(\tau )𝑑\tau \right)`$ (94)
$`andD_F`$ $`=`$ $`{\displaystyle \frac{A_F}{8\pi ^2}}{\displaystyle _0^{t_X}}𝑑t{\displaystyle \underset{i=1}{\overset{3}{}}}\left(1{\displaystyle \frac{c_ig_i^2t}{8\pi ^2}}\right)^{\frac{K_i^F}{c_i}}\mathrm{exp}\left({\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^t}Y_F(\tau )𝑑\tau \right).`$ (95)
It is not so easy to calculate $`Y_F`$ for other fermions even with ’heavy top integral’ approximation given by equation (6). So we consider the Yukawa-like coupling $`h_F(t)`$ for $`F`$=2,3, …, 12, which is related to $`M_F(t)`$ as
$$M_F(t)=h_F(t)\mathrm{exp}\left(\frac{1}{16\pi ^2}_0^tG_F(\tau )𝑑\tau \right).$$
(96)
The RG Equation for $`h_F(t)`$ is,
$$8\pi ^2d\left(\mathrm{log}h_F^2(t)\right)=A_Fh_F^2(t)\mathrm{exp}\left(\frac{1}{8\pi ^2}_0^tG_F(\tau )𝑑\tau \right)dt+Y_F(t)dt.$$
(97)
Here $`h_F(t)`$ contains the usual $`\beta `$-angle factors of the two Higgs system. As has been discussed in Section-1, we may not need this angle in our approach to the problem at hand. If $`F>1`$, the heaviest fermion next to the top is the bottom. The first term of equation(97), at the $`M_Z`$ scale, is $`1/1235`$ times smaller, whereas the last term contains one $`h_{top}`$ in $`h_{botm}`$. To a good approximation, and to begin with, we shall neglect this term and express $`h_F(t)`$ in terms of $`Y_F(t)`$’s. Then
$$_0^{t_X}Y_F(\tau )𝑑\tau =8\pi ^2\mathrm{log}\frac{h_F^2(M_X)}{h_F^2(M_Z)}.$$
(98)
Here, $`h_F^2(M_Z)=M_F^2(M_Z)`$ by definition. In the units of top mass 175 GeV,
$`h_F^2(M_X)=M_F^2(M_X)\mathrm{exp}\left({\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^{t_X}}G_F(\tau )𝑑\tau \right)`$ $`=`$ $`M_{top}^2(M_X)\mathrm{exp}\left({\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^{t_X}}G_F(\tau )𝑑\tau \right)`$ (99)
$`=`$ $`M_{top}^2(M_Z)=1.`$ (100)
So,
$`{\displaystyle _0^{t_X}}Y_F(\tau )𝑑\tau `$ $`=`$ $`8\pi ^2\mathrm{log}M_F^2(M_Z)`$ (101)
$`or\mathrm{exp}\left({\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^{t_X}}Y_F(\tau )𝑑\tau \right)`$ $`=`$ $`{\displaystyle \frac{M_{top}^2(M_Z)}{M_F^2(M_Z)}}.`$ (102)
This is true as long as the first term of equation(97) is negligible. Since there is $`M_F^3`$ in this term, it is much more justifiable to set $`D_{lepton}=0`$. The general equation (93) for the leptons gives
$$\frac{M_{top}^2(M_Z)}{M_{lepton}^2(M_Z)}=\frac{M_{top}^2(M_Z)}{M_{lepton}^2(M_X)}C_{lepton},$$
(103)
where
$$C_{lepton}=\underset{i=1}{\overset{3}{}}\left(\frac{g_i^2}{g_U^2}\right)^{\frac{K_i^{lepton}}{C_i}},$$
$`M_{lepton}^2(M_X)`$ $`=`$ $`M_{lepton}^2C_{lepton}`$ (104)
$`\text{or}M_{lepton}(M_X)`$ $`=`$ $`M_{lepton}^{expt.}C_{lepton}^{1/2}.`$ (105)
Putting the gauge constants and masses in the above, as experimentally reported, we get the masses at GUT scale for different leptons as given in Table-1,
At the grand unification mass, the electron, the muon and the tau climb to 116 GeV in this approximation. The calculation for the neutrinoes are not reliable as the isospin factors are uncertain and may be inaccurate.
We are now left with the five quarks. For them, we attempt to find the next leading order approximation, i.e., we first set the first term equal to zero and obtain
$$\mathrm{exp}\left(\frac{1}{8\pi ^2}_0^{t_X}Y_Q(\tau )𝑑\tau \right)=\frac{1}{h_Q^2(M_Z)}=\frac{M_{top}^2(M_Z)}{M_Q^2(M_Z)}.$$
(106)
This is used in the calculation for $`C_Q`$ of equation(94), which gives
$$C_Q=\underset{i=1}{\overset{2}{}}\left(\frac{g_i^2}{g_U^2}\right)^{\frac{K_i^Q}{C_i}}\frac{1}{M_Q^2(M_Z)}.$$
(107)
Using
$$\mathrm{exp}\left(\frac{1}{8\pi ^2}_0^tY_F(\tau )𝑑\tau \right)=\frac{h_Q^2(t)}{M_Q^2(M_Z)},$$
(108)
in equation(95), we get
$$D_Q=\frac{6}{8\pi ^2}_0^{t_X}\frac{h_Q^2(t)}{M_Q^2(M_Z)}\underset{i=1}{\overset{2}{}}\left(1\frac{g_i^2c_it}{8\pi ^2}\right)^{\frac{K_i^Q}{c_i}}dt.$$
(109)
We look for ways to calculate $`h_Q(t)`$. We can retain only the top quark in the RG equation in a slightly different way than what has been taken by RRRa6 and get,
$$8\pi ^2d\mathrm{log}h_Q^2(t)N_Qdt.$$
(110)
We have
$$N_{charm}=3,N_{up}=3,N_{bottom}=1,N_{strange}=N_{down}=0.$$
(111)
To use this value in $`Y_F`$, it is necessary to take an average as the couplings change very rapidly,
$$\overline{\mathrm{log}h_Q^2(t)}=\frac{1}{t_X}_0^t\frac{d}{d\tau }\mathrm{log}h_Q(\tau )𝑑\tau =\frac{1}{t_X}_0^t\frac{N_Q}{8\pi ^2}𝑑\tau =\frac{t}{t_X}\frac{N_Q}{8\pi ^2},$$
(112)
or
$$\overline{h_Q^2(t)}=\mathrm{exp}\left(\frac{t}{t_X}\frac{N_Q}{8\pi ^2}\right).$$
(113)
This expression for Yukawa coupling, is not much different from unity for all allowed $`t`$’s. We arrive at the following result,
$$D_Q=\frac{6}{8\pi ^2}_0^{t_X}𝑑t\mathrm{exp}\left(\frac{t}{t_X}\frac{N_Q}{8\pi ^2}\right)\underset{i=1}{\overset{3}{}}\left(1\frac{g_i^2c_it}{8\pi ^2}\right)^{\frac{K_i^Q}{c_i}}.$$
(114)
The unification mass is calculated numerically from
$$M_Q(M_X)=M_{top}\left(\frac{C_Q}{1D_Q}\right)^{1/2}.$$
(115)
The results, for the quarks, are given in Table-2. We note that $`M_{charm}(M_X)=M_{up}(M_X)M_{top}(M_X)`$ and $`M_{strange}(M_X)=M_{down}(M_X)M_{bottom}(M_X)`$.
Thus, we have shown that all fermions seem to originate at an energy $`2.2\times 10^{16}`$ GeV with equal mass of about 115 GeV. Perhaps, this is due to the equality of A$``$ 1 in Wolfenstein’s parametrisation of CKM matrix. In this perturbative method of solution for finding the unification mass, information about the masses of the 11 fermions has been lost due to cancellation of $`M^2(M_Z)`$ in both r.h.s. and l.h.s. of the equation (85) due to use of equations (102) and (108). The result of a common mass at the origin is atleast a hypothesis and has an approximation a10 .
The top mass decreases with energy but all the other quarks and leptons, starting from the value at $`M_X`$ acquire smaller and smaller values and become quite light in the electroweak scale. The descent or ascent equation(86) describing the ‘run’ can be put in a form like equation(93).
$$\frac{M_{top}^2(M_Z)}{M_F^2(t)}=\frac{M_{top}^2(M_Z)}{M_F^2(M_X)}C_F(t)+D_F(t),$$
(116)
where
$$C_F(t)=a_F(t)\mathrm{exp}\left(\frac{1}{8\pi ^2}_0^{t_X}Y_F(\tau )𝑑\tau \right),$$
(117)
$$D_F(t)=\frac{A_F}{8\pi ^2}_t^{t_X}𝑑t_1b_F(t_1)\mathrm{exp}\left(\frac{1}{8\pi ^2}_t^{t_1}Y_F(\tau )𝑑\tau \right),$$
(118)
$$a_F(t)=\underset{i=1}{\overset{3}{}}\left[\frac{\left(1\frac{g_i^2c_it_X}{8\pi ^2}\right)}{\left(1\frac{g_i^2c_it}{8\pi ^2}\right)}\right]^{\frac{K_i^F}{c_i}},$$
(119)
$$andb_F(t)=\underset{i=1}{\overset{3}{}}\left[\frac{\left(1\frac{g_i^2c_it_1}{8\pi ^2}\right)}{\left(1\frac{g_i^2c_it}{8\pi ^2}\right)}\right]^{\frac{K_i^F}{c_i}}.$$
(120)
For the top, we can take $`Y_F0`$ and calculate variation of its mass from $`M_U115`$ GeV to the top mass 175 GeV.
$$M_{top}(t)=\frac{M_{top}(M_X)}{\left(C_{top}(t)+\frac{M_U^2}{M_{top}^2}D_{top}(t)\right)^{1/2}}$$
(121)
The values are given in Table-5 of Section-8.
For other cases, we shall try to fit them into a scheme which is not only much simpler, but has better physical content. However, in the following section, we discuss a different route for solutions without specifying gauge factors completely.
## IV Deduction of Wolfenstein parameters and the rotational integers
Let us construct a function $`B_F(t)`$ such that
$$8\pi ^2\frac{d}{dt}\mathrm{log}B_F^2(t)=Z_F(t).$$
(122)
Then
$`\mathrm{exp}\left({\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^{t_X}}Z_F(\tau )𝑑\tau \right)`$ $`=`$ $`{\displaystyle \frac{B_F^2(M_X)}{B_F^2(M_Z)}},`$ (123)
$`\mathrm{exp}\left({\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _{t_X}^0}Z_F(\tau )𝑑\tau \right)`$ $`=`$ $`{\displaystyle \frac{B_F^2(M_Z)}{B_F^2(M_X)}},`$ (124)
$`\mathrm{exp}\left({\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _0^t}Z_F(\tau )𝑑\tau \right)`$ $`=`$ $`{\displaystyle \frac{B_F^2(t)}{B_F^2(M_Z)}},`$ (125)
$`\text{and}\mathrm{exp}\left({\displaystyle \frac{1}{8\pi ^2}}{\displaystyle _{t_X}^t}Z_F(\tau )𝑑\tau \right)`$ $`=`$ $`{\displaystyle \frac{B_F^2(t)}{B_F^2(M_X)}}.`$ (126)
Equation(85) reduces to
$$M_F(M_Z)=\frac{\frac{B_F(M_Z)}{B_F(M_X)}M_F(M_X)}{\left(1+\frac{M_F^2(M_X)}{M_{top}^2}\frac{A_F}{8\pi ^2}_0^{t_X}𝑑t\frac{B_F^2(t)}{B_F^2(M_Z)}\right)^{1/2}}.$$
(127)
As indicated in the last section, except the top, we neglect the terms with $`A_F`$. Then
$$M_F(M_Z)M_F(M_X)\frac{B_F(M_Z)}{B_F(M_X)}=M_F(M_X)\mathrm{exp}\left(\frac{I_F}{16\pi ^2}\right)=M_F(M_X)\lambda ^{n_F},$$
(128)
where the integral $`I_F`$ is
$$I_F=_0^{t_X}Z_F(t)𝑑t=_0^{t_X}\left(\underset{G}{}A_{FG}M_G^{}(t)M_G(t)G_F(t)\right)𝑑t,$$
(129)
with $`A_{1G}`$=0 to indicate that this is not for the top. Since the ratio of the two fermion masses can be put as $`\frac{M_{F_1}}{M_{F_2}}=\lambda `$, the parameter $`M_\nu `$ cancels out and $`\lambda `$ is the Wolfenstein parameter. The purpose of introducing the logarithms is to incorporate the possible multivaluedness of $`I_F`$ and obtain the integral powers of $`\lambda `$. From equation (129),
$`I_F`$ $`=`$ $`{\displaystyle _0^{t_X}}Z_F(t)𝑑t={\displaystyle _0^{t_X}}\left({\displaystyle \underset{G}{}}A_{FG}M_G^{}(t)M_G(t)G_F(t)\right)𝑑t`$ (130)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^{t_X}}\left({\displaystyle \underset{G}{}}A_{FG}M_G^{}(t)M_G(t)G_F(t)+{\displaystyle \underset{G}{}}A_{FG}M_G^{}(t)M_G(t)G_F(t)\right)𝑑t`$ (131)
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle _{t_X}^{t_X}}{\displaystyle \frac{dt}{dM_F}}𝑑M_F(t)\left({\displaystyle \underset{G}{}}A_{FG}M_G^{}(t)M_G(t)G_F(t)+{\displaystyle \underset{G}{}}A_{FG}M_G^{}(t)M_G(t)G_F(t)\right)`$ (132)
$`=`$ $`{\displaystyle \frac{16\pi ^2}{4}}{\displaystyle _{M_U}^{M_U}}{\displaystyle \frac{dM_F(t)}{M_F(t)}}{\displaystyle \frac{\left(_GA_{FG}M_G^{}(t)M_G(t)\frac{1}{2}(G_F(t)+G_F(t))\right)}{\left(A_FM_F^{}(t)M_F(t)G_F(t)\right)}},`$ (133)
where we have used equation (60). Let us set
$$M_F(t)=M_Ue^{i\theta _F(M_F(t))n_F}=M_Ue^{i\theta _F(t)n_F},$$
where $`n_F`$ is an integer and we will call it rotational integer. Using the above in (133), we get
$$I_F=in_F\frac{16\pi ^2}{2}_{t_X}^{t_X}𝑑\theta _F(t)\frac{_GA_{FG}}{A_F+_GA_{FG}\frac{G_F(t)}{M_U^2}}.$$
(134)
We have omitted inconsequential factor $`\frac{G_F}{M_U^2}`$ in the numerator. We note that for quarks $`A_F+_{G=1}^6`$=6+7=13 and for leptons 4+9=13. The integral is almost a constant and isolate the integer $`n_F`$, characterising F from the inegral. Using $`1=\frac{1}{12}_H`$ in equation (134), we have
$$I_F=in_F\frac{16\pi ^2}{2}\frac{1}{12}\underset{H}{}\underset{G}{}A_{HG}_{t_X}^{t_X}𝑑\theta _H(t)\frac{1}{A_H+_GA_{HG}\frac{G_H}{M_U^2}}.$$
(135)
In the above we have averaged over the twelve fermions. Retracing the steps and using $`M_H(t)=M_Ue^{i\theta _H(t)M_U^2}`$ in equation(135), we finally get
$$I_Fn_F\frac{1}{12}\underset{H}{}\underset{G}{}A_{HG}_0^{t_X}𝑑t=\frac{1}{12}n_Ft_X\underset{H}{}\underset{G}{}A_{HG}.$$
(136)
First, we shall be interested in the CKM matrix for quarks with lepton masses taken as zero. Then let G vary from 1 to 6, whereas F will be taking values from 2 to 12, since all of them contains 6 quarks. Only the coefficients $`A_{HG}`$ are needed to calculate $`I_F`$ of equation (136). The coefficients $`A_{FG}`$ have non-vanishing values for $`F`$=2,3,$`\mathrm{}`$,12 and $`G`$=1,2,$`\mathrm{}`$ ,6.
$$\text{For}F=2,\mathrm{},6:A_{F1}+A_{F2}+\mathrm{}+A_{F6}=7,$$
(137)
and
$$\text{For}F=7,\mathrm{},12:A_{F1}+A_{F2}+\mathrm{}+A_{F6}=9.$$
(138)
For the twelve fermions, the sum of the values of the coefficients $`A_{FG}`$ is $`(7\times 5)+(9\times 6)=89`$. The average of Z, i.e., $`\overline{Z}`$ is geven
$$\overline{Z}=89/12,\text{and}I_F=n_Ft_X\frac{89}{12}.$$
(139)
From equation(128), the Wolfenstein parameter $`\lambda `$ turns out to be
$$\lambda =\mathrm{exp}\left(\frac{t_X}{16\pi ^2}\frac{89}{12}\right)=0.219$$
(140)
This is an excellent result in spite of the approximate estimates. The masses of all the fermions due to quark-lepton equivalence other than the top is
$$M_F(M_Z)M_F(M_X)\lambda ^{n_F}\lambda ^{n_F}M_U$$
(141)
The table-3 identifies the particles. We have increased $`n_F`$ by neighbourhood integers which we have called the rotational integers.
## V Gauge contributions
### V.1 The top mass
It is easy to write down the equation for the top mass in terms of the mass $`M_F(M_X)`$ and gauge couplings.
$$M_{top}=M_{top}(M_X)\left(\frac{1d_{top}}{a_{top}}\right)^{1/2}.$$
(142)
We shall also need the equation for the descent of top mass
$$\frac{M_{top}^2}{M_{top}^2(t)}=\frac{M_{top}^2}{M_{top}^2(M_X)}a_{top}(t)+D_{top}(t).$$
(143)
The equation contains only the gauge factors.
### V.2 The masses of charm and up quarks.
The RG Equation for the top is non-linear and intermixed. Specifically it is
$$8\pi ^2\frac{d}{dt}\mathrm{log}M_{top}^2(t)=G_U+6M_{top}^2(t)+Y_{top}(t).$$
(144)
Following a texture analysis of the RG solutions, we introduce a function $`B(t)`$,( not to be confused with $`B_F(t)`$), which satisfies
$$8\pi ^2\frac{d}{dt}\mathrm{log}B^2(t)=G_U+6M_{top}^2(t).$$
(145)
Equation(145) does not specify the function $`B(t)`$ completely except that $`d\mathrm{log}B(t)=d\mathrm{log}(M_{top}(t))`$. This lone restriction gives us an infinite number of free choices. Because letting B$`\xi `$B, where $`\xi `$ is an arbitrary constant, does not change the top mass function $`M_{top}(t)`$. For simplicity, we take the second derivative of this function $`B(t)`$ to be zero so that
$$d\mathrm{log}B(t)=C_Bdt.$$
(146)
$`C_B`$ is essentially $`d\mathrm{log}M_{top}(t)`$, $`M_{top}(t)`$ changes by only 30 to 50 GeV as the mass $`\mu `$ goes from 91 GeV to $`2.2\times 10^{16}`$ GeV. $`C_B(175(123\text{ to}115))/1750.297\text{ to}0.342`$. We take $`C_B=0.3`$. The nonequiness due to nonlinearity can also be seen as follows:
$$dlogB(t)=dlogM(t)=dlog\frac{M_{top}}{M_o}=\frac{dt}{16\pi ^2}\left[G_U+6\left(\frac{M_{top}}{M_o}\right)^2\right].$$
$`M_o`$ is an arbitrary constant. We can choose $`M_o`$ suitably so that
$$\left[G_U+6\left(\frac{M_{top}}{M_o}\right)^2\right]/16\pi ^20.3=C_B.$$
Integrating from $`t=t_X(M_X)`$ to $`t=0(M_Z)`$, we have
$$\mathrm{log}\frac{B_Z}{B_U}=C_Bt_X,$$
(147)
and
$$\frac{B_Z}{B_U}=e^{C_Bt_X}=(0.192)^6.$$
(148)
Furthermore, integrating from $`t=t_X`$ to arbitrary $`t`$,
$$\mathrm{log}\frac{B(t)}{B_U}=C_B(tt_X),$$
(149)
or
$$\frac{B(t)}{B_U}=\mathrm{exp}\left(C_B(tt_X)\right)=\left[\left(\frac{B_Z}{B_u}\right)^{1/6}\right]^{(1t/t_X)}.$$
(150)
Let us examine the case for the charm. We have
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_c^2(t)`$ $`=`$ $`G_U+3M_t^2,`$ (151)
$`=`$ $`G_U+{\displaystyle \frac{1}{2}}[6M_t^2G_U]+{\displaystyle \frac{1}{2}}G_U,`$
$`=`$ $`{\displaystyle \frac{1}{2}}G_U+4\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2(t).`$
Integrating from $`t_X`$ to $`t`$,
$`\mathrm{log}{\displaystyle \frac{M_{charm}^2(t)}{M_{charm}^2(M_X)}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}{\displaystyle \frac{B^2(t)}{B_U^2}}{\displaystyle \frac{1}{16\pi ^2}}{\displaystyle _{t_X}^t}G_U(t)𝑑t,`$ (152)
$`M_{charm}(t)`$ $`=`$ $`M_U\left({\displaystyle \frac{B(t)}{B_U}}\right)^{1/2}\mathrm{exp}\left({\displaystyle \frac{1}{32\pi ^2}}{\displaystyle _{t_X}^t}G_U(\tau )𝑑\tau \right),`$ (153)
$`M_{charm}(M_Z)`$ $`=`$ $`M_U\left({\displaystyle \frac{B_Z}{B_U}}\right)^{1/2}C_{top}^{1/4},C_{top}^{1/4}=1.8466,`$ (154)
$`andM_c`$ $`=`$ $`\left[\left({\displaystyle \frac{B_Z}{B_U}}\right)^{1/6}C_{top}^{1/12}\right]^3M_{charm}(M_X).`$ (155)
From section-4, we now calculate
$$\lambda _{charm}=e^{c_Bt_X/6}C_{top}^{1/12}=0.221$$
(156)
Again this is a good result. The mass of the charm is 1.24 GeV from the equation (155) in very good agreement with the experimental value. The rotational integer $`n_F`$ is three as in Table-3.
We continue further and consider the up quark. The ‘heavy top integral’ approximation of the RGE for the up quark is
$$8\pi ^2\frac{d}{dt}\mathrm{log}M_{up}^2=G_U+3M_{top}^2(t).$$
(157)
Guessing from the general rotational integral parametrization, as shown in Table-3 as $`\lambda ^7`$, we write this equation as
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{up}^2`$ $`=`$ $`G_U+{\displaystyle \frac{1}{2}}(6M_{top}^2(t)),`$ (158)
$`=`$ $`{\displaystyle \frac{G_U}{2}}+8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2(t){\displaystyle \frac{1}{2}}8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{top}^2(t).`$
Integrating, we get
$`{\displaystyle \frac{M_{up}}{M_U}}`$ $`=`$ $`\left({\displaystyle \frac{B_Z}{B_u}}\right)\left({\displaystyle \frac{M_U}{M_{top}}}\right)^{1/2}a_U^{1/4},`$ (159)
$`\lambda _{up}`$ $`=`$ $`\left[\left({\displaystyle \frac{B_Z}{B_u}}\right)\left({\displaystyle \frac{M_U}{M_{top}}}\right)^{1/2}a_U^{1/4}\right]^{1/7}=0.220.`$ (160)
This gives $`M_{up}`$ 0.0029 GeV, with $`M_U=M_{top}(M_X)`$, a little lower value but quite close to the experimental result. The values depend crucially on $`C_B`$, which is itself not exact. It is also true that $`M_U`$ may be a little different from 115 GeV.
## VI Masses of the down quarks
The calculation of the mass of up-quark has set the method of finding solutions close to the experimental values by using the texture analysis function $`B(t)`$ ignoring the terms with coefficients $`A_F`$. We present a general recipe from the ‘heavy top integral’ approximation. Let the RG Equation for any $`F1`$ be
$$8\pi ^2\frac{d}{dt}\mathrm{log}M_F^2=G_F+N_FM_{top}^2.$$
(161)
Here $`N_F`$ can be zero as well.
$`N_FM_{top}^2(t)`$ $`=`$ $`{\displaystyle \frac{N_F}{6}}\left(6M_{top}^2(t)G_U\right)+{\displaystyle \frac{N_F}{6}}G_U,`$ (162)
$`=`$ $`{\displaystyle \frac{N_F+\eta _F}{6}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2(t)\right){\displaystyle \frac{\eta _F}{6}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{top}^2(t)\right)+{\displaystyle \frac{N_F}{6}}G_U.`$
$`\eta _F`$ can be any arbitrary coefficient. All of them will satisfy the RG Equation in the ‘heavy top integral’ approximation. But they should let $`n_F`$ and $`\lambda `$ be such that they are within rotational integers and gauge interaction contributions.
Integrating from the known values, $`t_X`$ to 0, we get
$$M_F(M_Z)=M_U\mathrm{exp}\left(\frac{1}{16\pi ^2}_0^{t_X}G_F(t)𝑑t\frac{N_F}{96\pi ^2}_0^{t_X}G_U(t)𝑑t\right)\left(\frac{B_Z}{B_U}\right)^{\frac{N_F+\eta _F}{6}}\left(\frac{M_U}{M_{top}}\right)^{\frac{\eta _F}{6}}.$$
(163)
For the down quarks,
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_b^2`$ $`=`$ $`G_D+M_{top}^2,`$ (164)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_s^2`$ $`=`$ $`G_D,`$ (165)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_d^2`$ $`=`$ $`G_D,`$ (166)
$`\text{So}8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_b^2`$ $`=`$ $`G_D+2M_{top}^2M_{top}^2,`$ (167)
$`=`$ $`G_D+{\displaystyle \frac{2.5}{6}}(6M_{top}^2G_U){\displaystyle \frac{1.5}{6}}(6M_{top}^2G_U)+{\displaystyle \frac{G_U}{6}},`$ (168)
$`=`$ $`G_D+{\displaystyle \frac{2.5}{6}}(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2){\displaystyle \frac{1.5}{6}}(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{top}^2)+{\displaystyle \frac{G_U}{6}}.`$ (169)
Integrating from $`t_X`$ to 0, we get
$`{\displaystyle \frac{M_b}{M_U}}`$ $`=`$ $`\left({\displaystyle \frac{B_Z}{B_U}}\right)^{2.5/6}\left({\displaystyle \frac{M_U}{M_{top}}}\right)^{1.5/6}a_D^{1/2}a_U^{1/12},`$ (170)
$`=`$ $`\left[\left({\displaystyle \frac{B_Z}{B_U}}\right)^{2.5/12}\left({\displaystyle \frac{M_U}{M_{top}}}\right)^{1.5/12}a_D^{1/4}a_U^{1/24}\right]^2.`$ (171)
The quantity in the square bracket is $`\lambda `$ 0.203. This gives the value of the bottom mass as to be 4.7 GeV.
For the strange we have,
$$8\pi ^2\frac{d}{dt}\mathrm{log}M_s^2=G_D+\frac{4.5}{6}\left(6\pi ^2\frac{d}{dt}\mathrm{log}B^2G_U\right)\frac{4.5}{6}\left(6\pi ^2\frac{d}{dt}\mathrm{log}M_{top}^2G_U\right).$$
(172)
This leads to
$$\frac{M_s}{M_U}=\left(\frac{B_Z}{B_U}\right)^{4.5/6}\left(\frac{M_U}{M_{top}}\right)^{4.5/6}a_D^{1/2}=\left[\left(\frac{B_Z}{B_U}\right)^{4.5/24}\left(\frac{M_U}{M_{top}}\right)^{4.5/24}a_D^{1/8}\right]^4=\lambda ^4.$$
(173)
Here $`\lambda `$ comes out to be near 0.196 and the strange mass is found to be 0.168 GeV.
Proceeding further to the down quark, we get
$$8\pi ^2\frac{d}{dt}\mathrm{log}M_d^2=G_D+\frac{6.5}{6}\left(6\pi ^2\frac{d}{dt}\mathrm{log}B^2G_U\right)\frac{6.5}{6}\left(6\pi ^2\frac{d}{dt}\mathrm{log}M_{top}^2G_U\right),$$
(174)
which gives
$$\frac{M_d}{M_U}=\left(\frac{B_Z}{B_U}\right)^{6.5/6}\left(\frac{M_U}{M_{top}}\right)^{6.5/6}a_D^{1/2}=\left[\left(\frac{B_Z}{B_U}\right)^{6.5/36}\left(\frac{M_U}{M_{top}}\right)^{6.5/36}a_D^{1/12}\right]^6=\lambda ^6.$$
(175)
This gives $`\lambda 0.19`$ and $`M_s`$=0.0053 GeV. The ratio $`M_s/M_d`$ 31. However, from equation(56), this ratio is
$$\frac{M_s}{M_d}=\lambda ^2=(0.2)^2=25.$$
(176)
These can be considered as gauge interaction corrections to the Wolfenstein parameter.
## VII The masses of the leptons
Ignoring the non-linear terms proportional to $`A_F`$=4, the RG Equation, retaining the top quark only, for the six leptons are
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_e^2`$ $`=`$ $`G_E,`$ (177)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_\mu ^2`$ $`=`$ $`G_E,`$ (178)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_\tau ^2`$ $`=`$ $`G_E,`$ (179)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{\nu _e}^2`$ $`=`$ $`G_N+3M_{top}^2,`$ (180)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{\nu _\mu }^2`$ $`=`$ $`G_N+3M_{top}^2,`$ (181)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{\nu _\tau }^2`$ $`=`$ $`G_N+3M_{top}^2.`$ (182)
For the first three electron-leptons, we write, with choices based on $`\lambda .22`$ and rotational integer $`n_F`$,
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_e^2`$ $`=`$ $`G_E+{\displaystyle \frac{30}{24}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2(t)G_U\right){\displaystyle \frac{30}{24}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{top}^2(t)G_U\right),`$ (183)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_\mu ^2`$ $`=`$ $`G_E+{\displaystyle \frac{17}{24}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2(t)G_U\right){\displaystyle \frac{17}{24}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{top}^2(t)G_U\right),`$ (184)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_\tau ^2`$ $`=`$ $`G_E+{\displaystyle \frac{11}{24}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2(t)G_U\right){\displaystyle \frac{11}{24}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{top}^2(t)G_U\right),`$ (185)
and obtain
$$\frac{M_e}{M_U}=\left(\frac{B_Z}{B_U}\right)^{5/4}\left(\frac{M_U}{M_{top}}\right)^{5/4}a_E^{1/2}=\left[\left(\frac{B_Z}{B_U}\right)^{5/32}\left(\frac{M_U}{M_{top}}\right)^{5/32}a_E^{1/16}\right]^8=\lambda _e^8,$$
(186)
which gives $`\lambda _e=0.209`$; $`M_e=4\times 10^4`$ GeV. Similarly
$$\frac{M_\mu }{M_U}=\left(\frac{B_Z}{B_U}\right)^{17/24}\left(\frac{M_U}{M_{top}}\right)^{17/24}a_E^{1/2}=\left[\left(\frac{B_Z}{B_U}\right)^{17/120}\left(\frac{M_U}{M_{top}}\right)^{17/120}a_E^{1/10}\right]^5=\lambda _\mu ^5,$$
(187)
which gives $`\lambda _\mu =0.238`$; $`M_\mu =`$ 0.09 GeV,
and
$$\frac{M_\tau }{M_U}=\left(\frac{B_Z}{B_U}\right)^{11/24}\left(\frac{M_U}{M_{top}}\right)^{11/24}a_E^{1/2}=\left[\left(\frac{B_Z}{B_U}\right)^{11/72}\left(\frac{M_U}{M_{top}}\right)^{11/72}a_E^{1/6}\right]^3=\lambda _\tau ^3,$$
(188)
which gives $`\lambda _\tau =0.23`$; $`M_\tau =`$ 1.36 GeV.
The masses of the neutrinoes have not yet been measured, only limits have been set. The reported values, as shown in Table-3, fall into a good pattern, namely, $`M_{\nu _e}=M_U\lambda ^{16}`$, $`M_{\nu _\mu }=M_U\lambda ^9`$ and $`M_{\nu _\tau }=M_U\lambda ^6`$, $`\lambda `$=0.22. So the equivalent, degeneracy lifting equations, showing the choices of $`n_F`$ are
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{\nu _e}^2`$ $`=`$ $`G_N+{\displaystyle \frac{5}{2}}\left(6M_{top}^2G_U\right)2\left(6M_{top}^2G_U\right)+{\displaystyle \frac{1}{2}}G_U,`$ (189)
$`=`$ $`G_N+{\displaystyle \frac{5}{2}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2\right)2\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{top}^2\right)+{\displaystyle \frac{1}{2}}G_U,`$
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{\nu _\mu }^2`$ $`=`$ $`G_N+{\displaystyle \frac{3}{2}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2\right)\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{top}^2\right)+{\displaystyle \frac{1}{2}}G_U,`$ (190)
$`8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{\nu _\tau }^2`$ $`=`$ $`G_N+\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}B^2\right){\displaystyle \frac{1}{2}}\left(8\pi ^2{\displaystyle \frac{d}{dt}}\mathrm{log}M_{top}^2\right)+{\displaystyle \frac{1}{2}}G_U.`$ (191)
The neutrino masses are obtained as
$`M_{\nu _e}=M_Ua_N^{1/2}a_U^{1/4}\left({\displaystyle \frac{B_Z}{B_U}}\right)^{5/2}\left({\displaystyle \frac{M_U}{M_{top}}}\right)^2\lambda _{\nu _e}^{}{}_{}{}^{16}3\times 10^9\text{GeV},`$ (192)
$`M_{\nu _\mu }=M_Ua_N^{1/2}a_U^{1/4}\left({\displaystyle \frac{B_Z}{B_U}}\right)^{3/2}\left({\displaystyle \frac{M_U}{M_{top}}}\right)\lambda _{\nu _\mu }^{}{}_{}{}^{9}1.4\times 10^4\text{GeV},`$ (193)
$`M_{\nu _\tau }=M_Ua_N^{1/2}a_U^{1/4}\left({\displaystyle \frac{B_Z}{B_U}}\right)\left({\displaystyle \frac{M_U}{M_{top}}}\right)^{1/2}\lambda _{\nu _\tau }^{}{}_{}{}^{6}1.3\times 10^2\text{GeV}.`$ (194)
This completes our calculation of the masses of the quarks and leptons in terms of known and calculable quantities $`\eta _F`$, satisfying the RG Equations. It is absolutely necessary to write down the solutions of the RG Equations for each quark and leptons separately, then alone, the contributions from the gauge interactions and the Yukawa mass coefficients $`A_{FH}`$ can be ascertained.
The values of the coefficients $`\eta _F`$ used above have been calculated, upto to the nearest fraction, from the cited rotational integers $`n_F`$ in Table-3 and the top coupling coefficients $`N_F`$ from the RG Equations. They are given in the Table-4.
## VIII Running masses of the fermions
The exact one loop solution of the RGE for change in mass values with energy is given by equation(84). With $`Z_F(t)=8\pi ^2\frac{d}{dt}\mathrm{log}B_F^2(t)`$, this equation reduces to
$$\frac{M_{top}^2}{M_F^2(t)}=\frac{M_{top}^2}{M_U^2}\frac{B_F^2(t_X)}{B_F^2(t)}+\frac{A_F}{8\pi ^2}_t^{t_X}𝑑t_1\frac{B_F^2(t_1)}{B_F^2(t)},$$
(195)
and
$$M_F(t)=M_F(M_X)\frac{B_F(t)/B_F(t_X)}{\left(1+\frac{M_U^2}{M_{top}^2}\frac{A_F}{8\pi ^2}_t^{t_X}𝑑t_1\frac{B_F^2(t_1)}{B_F^2(t_X)}\right)^{1/2}}.$$
(196)
If we take
$$B_F(t)=\lambda ^{n_F(1t/t_X)}=\mathrm{exp}\left(n_F(1t/t_X)\mathrm{log}_e\lambda \right),$$
(197)
and
$$M_F(t)=M_U\lambda ^{n_F}\left(1+\frac{M_U^2}{M_{top}^2}\frac{A_F}{8\pi ^2}_t^{t_X}𝑑t_1\lambda ^{2n_F(1t_1/t_X)}\right)^{1/2},$$
(198)
with
$`{\displaystyle _t^{t_X}}𝑑t_1\lambda ^{2n_F}t_1/t_X`$ $`=`$ $`{\displaystyle _t^{t_X}}𝑑t_1e^{2n_Ft_1/t_X\mathrm{log}_e\lambda }`$ (199)
$`=`$ $`{\displaystyle \frac{t_X}{2n_F}}{\displaystyle \frac{1}{\mathrm{log}_e\lambda }}\left[e^{2n_F\mathrm{log}_e\lambda }e^{2n_Ft/t_X\mathrm{log}_e\lambda }\right]`$
$`=`$ $`{\displaystyle \frac{t_X}{2n_F}}{\displaystyle \frac{1}{\mathrm{log}_e\lambda }}\left[\lambda ^{2n_F}\lambda ^{2n_Ft/t_X}\right],`$
we have, from equation (196)
$$M_F(t)=M_F(M_X)\frac{\lambda ^{n_F(1t/t_X)}}{\left[1+\frac{M_F^2(M_X)}{M_{top}^2}\frac{A_F}{8\pi ^2}\lambda ^{2n_F}\lambda ^{2n_F/t_X}\left(\lambda ^{2n_F}\lambda ^{2n_Ft/t_X}\right)\right]^{1/2}}$$
(200)
We shall not use this for constructing the tables. A simpler approximate form ignoring $`A_F`$ terms of equation (200),
$$M_F(t)=M_F(M_X)\lambda ^{n_F\left(1\frac{t}{33}\right)}$$
(201)
is presented in tabular form in Table-5 to Table-7. This gives the same CKM matrix. Now, we present our results in the following tables.
From the above values, we find that $`\frac{d^2M_{top}(t)}{dt^2}`$ is nearly zero. The top mass variation is approximately
$$M_{top}(t)=[175.\frac{t}{33}\times 60]\text{GeV}.$$
(202)
The graphs are much more revealing. The variation of mass for the Up quarks, including the top, have been shown in figure-1, for the Down in figure-2 , for the Electrons in figure-3 and for the Neutrinoes in figure- 4. All have been compressed in figure-5 to allow a glance at the totality of the descent and ascent of the masses to the unification mass of 115 GeV. We believe that an exact analysis will not differ much from those presented in these figures.
## IX Concluding Remarks
Deo and Maharanaa10 have already given exact solutions for the one loop RGEs in MSSM. They proved that all the fermions might have originated from a common mass $`M_U`$ 115 GeV at the GUT energy of $`M_X2.2\times 10^{16}`$ GeV in a perturbative scheme. This has not been confirmed by authors working in this field as they have not looked at the solutions which proved the same wayas one loop gauge sector renormalisation group equation.
As the energy diminishes, the mass of the top increases to 175 GeV at the $`Z`$-meson mass $`M_Z`$, whereas, masses of all other quarks and tiny leptonic masses increase from their value at $`M_Z`$ to $`M_U`$ of each 115 GeV. Hopefully, the exact nature of this variation could be obtained from RGE. By introducing an auxiliary function $`B(t)`$ through the mixing terms of RGE, a very simplified expression can be obtained for both the mass values at $`M_Z(t=0)`$ and their variation till $`M_X`$ when, the masses of the fermions which keep changing, attain the final value at $`M_X`$.
The success of the analysis of the Wolfenstein and the ratio parametrization of RRR given by equation(56) is clearly brought out by our solutions. An extension to running masses at linearity level, leads to a very simple formula for the fermions other than the top,
$$M_F(t)=M_U\left(\frac{M_F}{M_U}\right)^{(1\frac{t}{t_X})}=M_U\lambda ^{n_F(1\frac{t}{t_X})},$$
(203)
where $`M_F^{exp.}=M_F(M_Z)=\lambda ^{n_F}`$. The increase is exponential. Much more exact analytic studies are needed to deduce the values of the $`n_F`$ from RGE. As a first step, we follow up the texture analysis procedure by introducing a similar, but not the same function $`B(t)`$a8 , which is such that $`d\mathrm{log}B(t)=d\mathrm{log}M_{top}(t)`$. The resulting non-uniqueness is taken advantage of. Introducing known gauge couplings and the values of $`n_F`$, the effect of gauge interaction have been calculated.
In the foregoing sections 1 to 8, there has been two important omissions. We have not remarked neither about higher loop effects nor about the threshhold corrections. They may be very large due to very heavy top. There has been no mention of the CKM phases. But, we work with three generations only. So, in this case there is precisely, one phase angle $`\varphi `$. This is defined from the Wolfenstein parametrisation as
$$\lambda =\left(\frac{M_d}{M_s}+\frac{M_u}{M_c}+2\sqrt{\frac{M_d}{M_s}\frac{M_u}{M_c}}\mathrm{cos}\varphi \right)^{\frac{1}{2}},$$
(204)
and can be deduced by extending the Gatto-Sarton-Tonio-Oakes (GSTO)a14 to next leading order of solution of RGE, $`\lambda =(M_d/M_s)^{1/2}`$. Putting in the numbers $`n_F`$ from table-III, we get
$$\mathrm{cos}\varphi \frac{\lambda }{2}0.1,\varphi 95^0.$$
(205)
This result is reasonable.
It may be argued that there are too many parameters still in the guise of the rotation integers $`n_F`$ even though $`\lambda `$ could be determined. Let us recall the case of the unification of the strong, the weak, the electromagnetic and the gravitational interactions. One consciously omits gravity and lists the other three in order of their strengths. Here the top mass, with the largest value,has been deduced from the knowledge of the unification mass and gauge couplings quite accurately. For $`F`$=2, the two successive rotations i.e., 0 to $`4\pi `$ gives $`n_F`$=2 and $`M_U\lambda ^2`$ falls near the mass of the bottom. This process of complete successive rotations yields the masses of all the quarks and electrons, except $`n_F`$=5. However, there is a GUT prediction by Georgi and Jarlskoga3 , that the ratio $`M_s/M_d25\lambda ^2`$. This may be the reason why $`n_F`$ jumps from 4 to 6. There is also another accurate prediction
$$\frac{M_d/M_s}{(1M_d/M_s)^2}=9.\frac{M_e/M_\mu }{(1M_e/M_\mu )^2}.$$
(206)
With these GUT supplements, there is no other unknown additional parameter left. The unification mass $`M_U115`$ GeV is the only input.
Corroborating and in consonance with a lot of excellent work in neutrino physics have been done by many workers ,e.g. Babu, Mohapatra and Barra15 . The neutrino masses fall with the rotational integers, determined for the neutrinoes in the table-V in a very nice way.
The figure-5 illustrates the significance of this approach of solving the RGE. The similarity of figure-5 with gauge unification graph is striking. Thus, we prove, beyond doubt that the original ideas of Pati and Salama1 is correct as well as the SUSY standard model. We have given many extrapolatory ideas which can be successfully implemented so that SUSY standard model world can be characterised by only three parameters, the GUT mass $`M_X2.2\times 10^{16}`$, the GUT coupling strength $`\alpha _{GUT}\frac{1}{24.5}`$ and common-origin -fermion mass $`M_U115`$ GeV. This is a simplification beyond expectation and very much gratifying.
###### Acknowledgements.
We thank Prof. L. Maharana and Dr. P.K. Jena for elaborate discussions and Sri Sidhartha Mohanty and Sri Haraprasanna Lenka for help in computations and in preparing graphs respectively.
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# Lower bounds on the Calabi functional
## 1 Introduction
A cornerstone of Atiyah and Bott’s treatment of Yang-Mills theory over Riemann surfaces is a relation they discovered between the Yang-Mills functional and filtrations of a holomorphic bundle. The relation can be stated as follows. Let $`\mathrm{\Sigma }`$ be a compact Riemann surface with a fixed compatible metric, normalised to have area $`4\pi ^2`$. Let $`E`$ be a holomorphic vector bundle over $`\mathrm{\Sigma }`$ and consider a flag $``$ of sub-bundles
$$0=E_0E_1\mathrm{}E_q=E.$$
Recall that for any bundle $`V`$ over $`\mathrm{\Sigma }`$ the slope $`\mu (V)`$ of $`V`$ is defined to be the quotient of the degree of V by the rank. We say the flag $``$ is slope-decreasing if $`\mu (E_1)>\mu (E_2)\mathrm{}>\mu (E)`$. Equivalently the quotients $`Q_i=E_i/E_{i1}`$ have $`\mu (Q_1)>\mathrm{}>\mu (Q_q)`$. Define
$$\mathrm{\Phi }()=\left(\underset{i=1}{\overset{q}{}}\mu (Q_i)^2\mathrm{rank}(Q_i)\right)^{1/2}.$$
(1)
Atiyah and Bott relate this algebro-geometric quantity to the Yang-Mills functional—the $`L^2`$ norm of the curvature $`F_A`$—restricted to compatible unitary connections $`A`$ on $`E`$. Their result is
$$\underset{A}{inf}F_A=\underset{,\mathrm{dec}}{sup}\mathrm{\Phi }()$$
(2)
where on the right hand side the supremum runs over the slope-decreasing flags. In fact this supremum is attained by the canonical Harder-Narasimhan filtration of $`E`$. The infimum on the left hand side is not in general attained: this happens if and only if $`E`$ is a direct sum of stable bundles. Notice that there is an easy lower bound $`F_A\mu (E)\sqrt{\mathrm{rank}(E)}`$ deriving from the fact that $`\frac{i}{2\pi }\mathrm{Tr}(F_A)`$ represents $`c_1(E)`$, This trivial lower bound is just $`\mathrm{\Phi }(^0)`$ where $`^0`$ is the trivial flag (with $`q=1`$). So $`sup_{,\mathrm{dec}}\mathrm{\Phi }()\mu (E)\sqrt{\mathrm{rank}E}`$, and it is easy to see that equality holds if and only if the bundle $`E`$ is semi-stable, i.e if there is no nontrivial slope-decreasing flag.
The Atiyah-Bott result can be viewed as two statements: $`inf_Asup_{}`$ and $`inf_Asup_{}`$. The proof of the first of these involves a simple differential-geometric argument, turning on the principle that “curvature decreases in holomorphic sub-bundles and increases in holomorphic quotients”(, Proposition 8.13 and the remark in the second paragraph of page 575). The proof of the second relies on the theorem of Narasamihan and Seshadri on the existence of projectively flat unitary connections on stable bundles.
In , Calabi began the study of the $`L^2`$ norm of the scalar curvature of Kahler metrics, running over a fixed Kahler class on a compact Kahler manifold. This functional is equivalent to the $`L^2`$-norm of the full curvature tensor, in that the two differ by topological terms. The purpose of this paper is to establish an analogue of (one half of) the Atiyah-Bott result for this Calabi functional (and some variants for $`L^p`$ norms). Our result bears on the algebraic case, so we suppose that $`X`$ is a smooth complex projective variety, that $`LX`$ is a fixed ample line bundle and we consider Kahler metrics $`\omega `$ in the class $`c_1(L)`$. We use the notion, introduced in , of a test configuration $`𝒳`$ for $`X`$. The detailed definition will be reviewed below, but in essence this comprises a $`𝐂^{}`$-equivariant family $`\pi :𝒳𝐂`$ with generic fibre $`X_t=\pi ^1(t)`$ isomorphic to $`X`$, for $`t0`$. The central fibre $`X_0`$ need not be isomorphic to $`X`$, it may be a highly singular variety or even a scheme, but it has a $`𝐂^{}`$-action and the definition requires a lift of this action to a line bundle $``$. We will define a numerical invariant $`\mathrm{\Psi }(𝒳)`$ using the $`𝐂^{}`$ action on the vector spaces $`H^0(X_0,^k)`$, related to the generalised Futaki invariant. In stating our results, it is convenient to work with a quantity $`S(\omega )`$ defined to be $`\frac{1}{4\pi }`$ times the usual scalar curvature of the metric $`\omega `$. (Thus $`S(\omega )\frac{\omega ^n}{n!}=\frac{1}{2}\rho \frac{\omega ^{n1}}{(n1)!}`$, where $`\rho `$ is the Ricci form representing $`c_1(X)`$.) Then our result takes the form
$$\underset{\omega }{inf}S(\omega )_{L^2}\underset{𝒳}{sup}\mathrm{\Psi }(𝒳).$$
(3)
There is a rather trivial lower bound on the Calabi functional, deriving from the fact that the integral of $`S(\omega )`$ yields $`\frac{1}{2(n1)!}c_1(X)\omega ^{n1},[X],`$ so the average value $`\widehat{S}`$ is a topological invariant of the Kahler class. Thus
$$S(\omega )_{L^2}^2=S\widehat{S}_{L^2}^2+\widehat{S}^2\mathrm{Vol}(X)\widehat{S}^2\mathrm{Vol}(X).$$
(4)
The essential feature of our definition of $`\mathrm{\Psi }(𝒳)`$ is that we get a better bound than this precisely when $`X`$ is not “K-semistable” in the sense of , just as in the Atiyah-Bott case. (There is also an extension of this discussion to extremal metrics, see Section 2 below.) Thus one immediate consequence of our Theorem (as Richard Thomas pointed out to the author) is a new and simpler proof of the fact that the existence of a constant scalar curvature Kahler metric implies K-semistability. This argument does not need any special constructions (as in ) in the case when the manifold has holomorphic vector fields.
It is natural to conjecture that in fact
$$\underset{\omega }{inf}S(\omega )_{L^2}=\underset{𝒳}{sup}\mathrm{\Psi }(𝒳).$$
This would be a variant of the conjectures in ,, relating K-stability to the existence of constant scalar curvature metrics (and the extension of that discussion to extremal metrics in ). However neither conjecture immediately implies the other.
We will explain the analogy with the Atiyah-Bott result in somewhat more detail in Section 2. Both set-ups fit, formally, into the more general framework of “moment maps and stability” discussed in for example. The results can be seen as infinite-dimensional versions of part of the theory developed by Kirwan in finite dimensions. The method of proof we use here is that of finite-dimensional approximation, in the mould of , . That is, we derive our inequalities by studying the asymptotics of finite-dimensional problems, essentially of the kind considered by Kirwan. However the proof of the result in this paper is substantially simpler than that in . In particular we do not need to make use of the notion of “balanced” metrics. The essential ingredient in our proof is the asymptotic expansion for the “density of states” function due to Yau, Tian, Zelditch,Liu, Catlin and Ruan , ,, , . We also need some discussion of singular varieties and schemes to allow us to apply the general moment map theory. The author first saw that this could be done using the point of view developed in the work of Zhang , Phong and Sturm , involving the “Chow norm” and an action on Chow vectors. However it turns out that one can avoid appealing to these concepts explicitly and we only need some comparatively straightforward technical facts to handle these issues of singularities (notably a result related to the equivariant Riemann-Roch Theorem, Proposition 3 below).
In 1997, and intermittently since, the author has discussed with X-X. Chen the possibility of obtaining lower bounds on the Calabi functional using a similar pattern of argument, but replacing the finite-dimensional approximations by the use of “geodesic rays” in the space of Kahler potentials, in the manner of , . This would have independent interest and Chen and the author hope to discuss this work in a a future article.
The author is grateful to Richard Thomas and Xiu-Xiong Chen for discussions of these topics.
## 2 Test configurations
We begin by recalling the definition of a test configuration from . Given an ample line bundle $`LX`$ a test configuration for the pair $`(X,L)`$ consists of:
* a scheme $`𝒳`$ with a $`𝐂^{}`$ action;
* a flat $`𝐂^{}`$-equivariant map $`\pi :𝒳𝐂`$, with fibres $`X_t`$;
* an equivariant line bundle $`𝒳`$, ample on all fibres;
* for some $`r>0`$, an isomorphism of the pair $`(X_1,|_{X_1})`$ with the original pair $`(X,L^r)`$.
Thus we have a sequence of vector spaces $`U_k=H^0(X_0,|_{X_0}^k)`$ with $`𝐂^{}`$-actions. Let $`A_k:U_kU_k`$ be the endomorphisms generating these actions (so $`e^t𝐂^{}`$ acts as $`e^{tA_k}`$ on $`U_k`$). We are interested in the asymptotic behaviour, with respect to $`k`$, of the dimension of $`U_k`$ and the total weight of the action, i.e. the trace of the $`A_k`$. These are given by Hilbert polynomials:
$$\mathrm{dim}U_k=a_0k^n+a_1k^{n1}+\mathrm{},$$
(5)
$$\mathrm{Tr}A_k=b_0k^{n+1}+b_1k^n+\mathrm{},$$
(6)
for large $`k`$, as discussed in , . We go one step further and define a positive number $`Q`$ by the leading term in the polynomial defining the trace of the squares:
$$\mathrm{Tr}A_k^2Qk^{n+2}.$$
(7)
It is clear that that $`Q0`$, one way of seeing that $`Q`$ is strictly positive is to use (20) in Section 5 below. Notice that, from general theory, the $`a_i,b_i`$ and $`Q`$ are rational numbers.
Now we define
$$\mathrm{\Psi }(𝒳)=\frac{1}{r^{(n2)/2}}\frac{b_1}{\sqrt{Q}}$$
(8)
The normalisation by $`1/r^{(n2)/2}`$ means that we do not change $`\mathrm{\Psi }(𝒳)`$ if we replace $``$ by a positive power $`^s`$. Also this scaling weight is the same as that of the Calabi functional
$$r^{(n2)/2}S(r\omega )_{L^2(r\omega )}=S(\omega )_{L^2(\omega )}.$$
The upshot of this is that in proving our theorem we can always suppose $`r=1`$, which we do from now on. By the flatness of the family and the Riemann-Roch Theorem for $`X`$, we can identify the co-efficient $`a_0`$ with the volume of $`X`$ and $`a_1`$ with the integral of $`S(\omega )`$ for any metric $`\omega `$ representing $`c_1(L)`$. Thus $`\widehat{S}=a_1/a_0`$ and the trivial lower bound (4) is $`S|a_1|/\sqrt{a_0}`$.
There are two ways in which we can modify a test configuration. First, we can pull back $`𝒳`$ by a $`d`$-fold covering of the base. This changes $`A_k`$ to $`dA_k`$ and plainly does not affect $`\mathrm{\Psi }`$. Second we can change the $`𝐂^{}`$-action on the line bundle $``$ by a character $`\lambda \lambda ^\nu `$ of $`𝐂^{}`$. This changes $`A_k`$ to $`A_k+k\nu 1`$ and so $`b_i`$ to $`\stackrel{~}{b_i}=b_i+\nu a_i`$, and $`Q`$ to $`\stackrel{~}{Q}=Q+2b_0\nu +a_0\nu ^2`$. Thus for the new configuration $`\stackrel{~}{𝒳}`$ we have
$$\mathrm{\Psi }(\stackrel{~}{𝒳})=\frac{b_1+\nu a_1}{\sqrt{Q+2b_0\nu +a_0\nu ^2}}$$
(9)
We let consider the supremum of this expression over $`\nu `$, initially regarded as a real variable. The analysis of this supremum brings in the Futaki invariant $`F_𝒳`$ of the test configuration. This is defined to be
$$F_𝒳=b_1\frac{b_0a_1}{a_0}.$$
(This terminology differs by a factor of $`a_0`$ from that in , where only the sign of $`F_𝒳`$ was relevant.) Then, as the reader will easily verify, if $`F_𝒳<0`$ the expression in (9) above is maximised when
$$\nu =\frac{a_1Qb_0b_1}{a_0b_1a_1b_0}.$$
By taking a covering we may suppose that this is an integer, so the supremum is realised by some test configuration. On the other hand if $`F_𝒳0`$ the supremum is not attained for finite $`\nu `$ but occurs in the limit as $`\nu `$ tends to $`\mathrm{}`$ or $`\mathrm{}`$, depending on the sign of $`a_1`$. In this second case the supremum is just $`|a_1|/\sqrt{a_0}`$ which is just the trivial lower bound on $`S`$. In the first case, when the supremum is attained, a little calculation shows that the supremum is
$$\left(\frac{|a_1|^2}{a_0}+\widehat{\mathrm{\Psi }}_𝒳^2\right)^{1/2}$$
where we define
$$\widehat{\mathrm{\Psi }}_𝒳=\frac{F_𝒳}{\sqrt{Qb_1^2/a_0}}.$$
The denominator here has a natural interpretation. We write $`\underset{¯}{A}_k:U_kU_k`$ for the trace-free part of $`A_k`$, i.e.
$$\underset{¯}{A}_k=A_k\frac{\mathrm{Tr}A_k}{\mathrm{dim}U_k}1.$$
Then as $`k\mathrm{}`$,
$$\mathrm{Tr}\underset{¯}{A}_k^2(Qb_1^2/a_0)k^{n+2}.$$
With this discussion in place we can state our main theorem precisely.
###### Theorem 1
If $`𝒳`$ is a test configuration for the pair $`(X,L)`$ then for any Kahler metric $`\omega `$ in the class $`c_1(L)`$ we have
$$S(\omega )_{L^2}\mathrm{\Psi }(𝒳).$$
By the discussion above it is completely equivalent to prove that for any test configuration
$$S(\omega )\widehat{S}_{L^2}\frac{F_𝒳}{Q\frac{b_1^2}{a_0}},$$
and it is in this form that we shall prove the result. Notice that the inequality is vacuous unless $`F_𝒳<0`$, that is, unless $`(X,L)`$ is not K-semistable. We will prove a more general result dealing with $`L^q`$ norms (loosely analogous to the discussion of other norms in ). For an even integer $`p`$ we define $`N_p(𝒳)>0`$ by the leading term
$$\mathrm{Tr}\underset{¯}{A}_k^pN_p^pk^{n+p}$$
(10)
in the appropriate Hilbert polynomial. Thus, when $`p=2`$,
$$Q\frac{b_1^2}{a_0}=N_2^2.$$
Momentarily re-instating the integer $`r`$ in the definition of a test configuration, we define
$$\widehat{\mathrm{\Psi }}_p(𝒳)=\frac{1}{r^{(n/q)1}}\frac{F_𝒳}{N_p(𝒳)}.$$
(11)
Again, the power of $`r`$ is chosen so that this is unchanged if we replace $``$ by $`^s`$.
###### Theorem 2
If $`𝒳`$ is any test configuration for $`(X,L)`$ then for any metric $`\omega `$ in the class $`c_1(L)`$ we have
$$S(\omega )\widehat{S}_{L^q}\widehat{\mathrm{\Psi }}_p(𝒳),$$
where $`p`$ is any even integer and $`q`$ is the conjugate exponent to $`p`$ (i.e. $`p^1+q^1=1`$ ).
Just as before, one can check that that the scaling weight of the $`L^q`$ norm of $`S(\omega )`$ agrees with the power of $`r`$ in (11), so we can reduce to the case when $`r=1`$. By the preceding discussion, Theorem 2 (in the case when $`p=2`$) implies Theorem 1, and in the body of the paper below we prove Theorem 2. (We remark that, when $`p2`$, there is no simple exact relation between the $`L^q`$ norms of $`S`$ and $`S\widehat{S}`$. One can derive slightly different lower bounds for the $`L^q`$ norms of $`S`$ using the techniques of this paper, but we find it simpler to work with $`S\widehat{S}`$.)
If the automomorphism group of the pair $`(X,L)`$ contains a non-trivial compact connected subgroup $`G`$ then there is another relatively elementary lower bound on the Calabi functional, obtained from the Futaki invariant in its original, differential geometric, form. We consider a $`G`$-invariant metric on $`X`$. Then we get a Lie algebra homomorphism from $`\mathrm{Lie}(G)`$ to the functions on $`X`$, under Poisson bracket. Let $`\xi `$ be an element of $`\mathrm{Lie}(G)`$ and let $`H`$ be the corresponding Hamiltonian. Then the integrals
$$_XSH𝑑\mu _XH^p𝑑\mu ,$$
are topological invariants of the data: they do not change as we vary $`\omega `$ among $`G`$-invariant metrics in the same Kahler class. By Holder’s inequality we have
$$S\widehat{S}_{L^q}\frac{1}{H\widehat{H}_{L^p}}_X(S\widehat{S})H𝑑\mu $$
(12)
where $`\widehat{H}`$ is the average value of $`H`$, and the right hand side is a topological invariant of the data.
We will now explain how to derive this as a special case of our Theorem 2. By a density argument we can suppose that $`\xi `$ generates a circle action on $`(X,L)`$ which we extend to a holomorphic $`𝐂^{}`$ action. Then we define a test configuration by taking the product $`𝒳=X\times 𝐂`$ but using this non-trivial $`𝐂^{}`$ action. Thus in this case the central fibre $`(X_0,)`$ is just the original pair $`(X,L)`$, with the given $`𝐂^{}`$ action. The key point now is that
$$b_0=_XH𝑑\mu ,b_1=_XSH𝑑\mu ,$$
so that
$$F_𝒳=_X(S\widehat{S})H𝑑\mu .$$
Thus the algebro-geometric definition of the Futaki invariant reduces to the differential geometric one. This is explained in Prop. 2.2.2 of and we will obtain a generalisation of the fact in Proposition 3 below. Moreover, the same discussion shows that, for even integers $`p`$,
$$N_p=H\widehat{H}_{L^p}.$$
Changing the sign of $`\xi `$, if necessary, we can suppose that the Futaki invariant is negative. Then (12) follows as a special case of our Theorem 2. In fact our Theorem gives rather more even in this case since the lower bound is obtained for arbitrary metrics in the Kahler class, not just the $`G`$-invariant ones.
We return now to situation considered by Atiyah and Bott and explain how their result can be cast in a similar form to ours. Let $``$ be a flag of sub-bundles in $`E\mathrm{\Sigma }`$, as before. Let $`W=(w_1,\mathrm{},w_q)`$ be a vector of strictly increasing integers. Then one can construct a degeneration $`=(,W)`$ from this data. This is a $`𝐂^{}`$-equivariant bundle over $`𝐂^{}\times \mathrm{\Sigma }`$ which is isomorphic to $`E`$ on each slice $`\{t\}\times \mathrm{\Sigma }`$ for $`t0`$ and to $`E_0=Q_1Q_2\mathrm{}Q_q`$ when $`t=0`$, with the property that the $`𝐂^{}`$-action on $`E_0`$ has weight $`w_i`$ on the summand $`Q_i`$. We fix a square root $`K_\mathrm{\Sigma }^{1/2}`$ of the canonical bundle and any line bundle $`L\mathrm{\Sigma }`$ of degree $`1`$. Now we consider the $`𝐂^{}`$ action on the vector spaces
$$U_k=H^0(E_0K_\mathrm{\Sigma }^{1/2}L^k),$$
with generators $`A_k:U_kU_k`$. Following just the same pattern as before, we look at the large $`k`$ behaviour,
$$\mathrm{dim}U_k=\alpha _0k+\alpha _1,$$
$$\mathrm{Tr}A_k=\beta _0k+\beta _1,$$
$$\mathrm{Tr}A_k^2Qk,$$
and we define
$$\mathrm{\Psi }(,W)=\frac{\beta _1}{\sqrt{Q}}.$$
###### Lemma 1
For any bundle $`E\mathrm{\Sigma }`$ we have
$$\underset{,W}{sup}\mathrm{\Psi }(,W)=\underset{,\mathrm{dec}}{sup}\mathrm{\Phi }(),$$
where on the left hand side the supremum is taken over arbitrary flags $``$ and increasing weight vectors $`W`$ and on the right hand side the supremum is taken over slope-decreasing flags $``$.
Given a flag $``$ and weight vector $`W`$ we have
$$U_k=H^0(Q_iK_\mathrm{\Sigma }^{1/2}L^k),$$
with the action of weight $`w_i`$ on the the ith summand. It follows that
$$\alpha _0=r_i,\alpha _1=d_i,\beta _0=r_iw_i,\beta _1=d_iw_i,Q=r_iw_i^2,$$
where $`Q_i`$ has rank $`r_i`$ and degree $`d_i`$. So we obtain
$$\mathrm{\Psi }(,W)=\frac{d_iw_i}{\sqrt{r_iw_i^2}}.$$
(13)
Now suppose that $``$ is slope-decreasing. We maximise $`\mathrm{\Psi }(,W)`$ over weights $`w_i`$, initially regarded as arbitrary real numbers. The maximum occurs at $`w_i^{\mathrm{max}}=Cd_i/r_i`$ for any constant $`C`$. Taking a suitable $`C`$ we can suppose that $`w_i^{\mathrm{max}}`$ are integers, but more crucially the condition that $``$ was a slope-decreasing flag means that the $`w_i^{\mathrm{max}}`$ are increasing. Then we have
$$\mathrm{\Psi }(,W^{\mathrm{max}})=\mathrm{\Phi }(),$$
so we have established one half of the Proposition, i.e $`sup\mathrm{\Psi }sup\mathrm{\Phi }`$. We prove the other half by induction on the length $`q`$ of the flag. We take as inductive proposition the statement that for any flag $``$ of length $`q`$ and any weakly increasing weight vector $`W`$ (i.e with $`w_1\mathrm{}w_q)`$ we have $`\mathrm{\Psi }(,W)sup_{,\mathrm{dec}}\mathrm{\Phi }`$, where $`\mathrm{\Psi }(,W)`$ is defined by (13). This is clearly true for $`q=1`$. Suppose that $``$ is any flag and $`W^0`$ is a weakly-increasing weight vector. We maximise the expression for $`\mathrm{\Psi }(,W)`$ over all weakly-increasing weight vectors. If $``$ is slope-decreasing then we are in the same position as above and the maximum realises $`\mathrm{\Phi }()`$; so in this case
$$\mathrm{\Psi }(,W^0)\mathrm{\Psi }(,W^{max})sup\mathrm{\Phi }$$
as desired. If $``$ is not slope increasing then the maximum occurs at a vector $`W^{max}`$ with $`w_i^{max}=w_{i+1}^{max}`$ for some $`i`$. Then let $`^{}`$ be the flag obtained from $``$ by deleting $`E_i`$ and let $`W^{}`$ be the weight vector for $`^{}`$ obtained from $`W^{max}`$ in the obvious way. In the associated sum of quotients we are replacing the original pair $`Q_iQ_{i+1}`$ by a bundle $`Q_i^{}=E_{i+1}/E_{i1}`$. The rank of $`Q_i^{}`$ is the sum $`r_i+r_{i1}`$ and the degree of $`Q_i^{}`$ is $`d_i+d_{i1}`$. Thus $`\mathrm{\Psi }(^{},W^{})=\mathrm{\Psi }(,W^{max})`$. So we have
$$\mathrm{\Psi }(,W^0)\mathrm{\Psi }(,W^{max})=\mathrm{\Psi }(^{},W^{})sup\mathrm{\Phi },$$
where in the last inequality we use the inductive hypothesis. This completes the proof.
In sum, the version of the Atiyah-Bott result which is closer to our Theorem is the statement that
$$\underset{A}{inf}F(A)=\underset{,W}{sup}\mathrm{\Psi }(,W),$$
which is completely equivalent to the previous formulation by the elementary Lemma above. One could probably prove one direction here (i.e. that $`inf_Asup_{,W}`$) by using a finite-dimensional approximation argument in the manner of Wang , just as we will do for the Calabi functional. However there would not be much point to this, in view of the simple and direct differential geometric proof.
## 3 Differential Geometric asymptotics
In this section we relate the Calabi functional (and its $`L^q`$ variants) to the norm of a matrix $`\underset{¯}{M}`$ associated to a projective variety. Let $`z_\alpha `$ be standard homogeneous coordinates on $`𝐂𝐏^N`$ and let let $`h_{\alpha \beta }`$ be the function
$$h_{\alpha \beta }=\frac{z_\alpha \overline{z}_\beta }{\underset{¯}{z}^2}$$
on $`𝐂𝐏^N`$. For a smooth projective variety $`V𝐂𝐏^N`$ of dimension $`n`$ we define a self-adjoint matrix $`M=M(V)`$ with entries
$$M_{\alpha \beta }=_Vh_{\alpha \beta }𝑑\mu _{FS},$$
where $`d\mu _{FS}`$ is the standard volume form induced from the Fubini-Study metric (normalised so that the volume of $`V`$ is equal to its degree divided by $`n!`$). Let $`\underset{¯}{M}`$ be the trace-free part of $`M`$. Since the sum of the $`h_{\alpha \alpha }`$ is $`1`$ we have $`\underset{¯}{M}=M\frac{\mathrm{Vol}(V)}{N+1}1`$. We recall that for any $`q>1`$ the $`q`$-norm of a self-adjoint matrix $`T`$ is defined by
$$T_q^q=|\lambda _\alpha |^q,$$
where $`\lambda _\alpha `$ are the eigenvalues, repeated according to multiplicity.
Now let $`(X,L)`$ be an abstract polarised variety, as before. For large $`k`$ the sections of $`L^k`$ define a projective embedding of $`X`$, i.e a choice of basis of $`H^0(L^k)`$ yields a specific projective variety $`V`$.
###### Proposition 1
Let $`\omega _0`$ be a Kahler metric on $`X`$ in the class $`c_1(L)`$. Then for large enough $`k`$ there is a basis of $`H^0(L^k)`$ yielding a projective embedding $`XV_k𝐂𝐏^{N_k}`$ with for, any $`q1`$,
$$\underset{¯}{M}(V_k)_qk^{(n/q)1}S(\omega _0)\widehat{S}_{L^q}+O(k^{(n/q)2}).$$
The proof is a straightforward application of the asymptotic expansion for the density of states function (see references in Section 1) which we now recall. Let $`||_0`$ be a Hermitian metric on $`L`$ whose associated curvature form is $`2\pi i\omega _0`$. We write $`||_0`$ also for the induced metric on $`L^k`$. Then endow $`H^0(L^k)`$ with the standard $`L^2`$-norm defined by the volume form $`d\mu _0`$ of the fixed metric $`\omega _0`$ and the fibre metric $`||_0`$ on $`L^k`$. Define a function $`\rho _k`$ on $`X`$ by
$$\rho _k=|s_\alpha |^2,$$
for any orthonormal basis $`s_\alpha `$ of $`H^0(L^k)`$. Then the statement we need is that
$$\rho _k=k^n(1+k^1\eta _k)$$
where the functions $`\eta _k`$ converge (in $`C^{\mathrm{}}`$) to $`S(\omega _0)`$ as $`k\mathrm{}`$; in fact $`\eta _k=S+O(k^1)`$.
The construction of the projective embedding we need is the most obvious one. For each $`k`$, we take any orthonormal basis $`s_\alpha `$ of $`H^0(X,L^k)`$ and define a new fibre metric on the line bundle $`L`$ by
$$||_{}^2=\frac{1}{\rho _k^{1/k}}.$$
Then, denoting the induced metric on $`L^k`$ also by $`||_{}`$,
$$|s_\alpha |_{}^2=1.$$
The pull-back of the Fubini-Study metric under the embedding is $`k\omega `$ where
$$\omega =\omega _0+k^1\frac{i}{2\pi }\overline{}(\mathrm{log}(1+k^1\eta _k)).$$
Write the volume form of $`\omega `$ as
$$(1+k^2\nu _k)d\mu _0.$$
It is clear, from the fact that the sequence $`\eta _k`$ is bounded, that the $`\nu _k`$ are bounded. The functions $`h_{\alpha \beta }`$ pull back to under the projective embedding to $`(s_\alpha ,s_\beta )_{}`$. So we have
$$_{V_k}h_{\alpha \beta }𝑑\mu _{FS}=_X(s_\alpha ,s_\beta )_0\left(\frac{1+k^2\nu _k}{1+k^1\eta _k}\right)𝑑\mu _0.$$
We do not change the norm of $`\underset{¯}{M}(V_k)`$ if we apply a unitary transformation to $`𝐂^{N+1}`$, in other words if we make a different choice of orthonormal basis $`s_\alpha `$. Thus we can choose the basis so that $`M`$ is a diagonal matrix with diagonal entries $`m_\alpha `$ where
$$m_\alpha =_X|s_\alpha |_0^2\left(\frac{1+k^2\nu _k}{1+k^1\eta _k}\right)𝑑\mu _0.$$
Now the dimension $`N+1`$ is the integral of $`\rho _k`$ over $`X`$, and it follows that the trace-free part $`\underset{¯}{M}`$ of $`M`$ is the diagonal matrix with diagonal entries
$$\underset{¯}{m}_\alpha =m_\alpha \frac{1}{1+\widehat{\eta _k}k^1},$$
where $`\widehat{\eta }_k`$ is the average value of $`\eta _k`$. Then we obtain $`\underset{¯}{m}_\alpha =k^1(b_\alpha +ϵ_\alpha )`$ where
$$b_\alpha =_X|s_\alpha |_0^2(\widehat{\eta _k}\eta _k)𝑑\mu _0,$$
and $`ϵ_\alpha `$ is $`O(k^1)`$. Now write
$$|s_\alpha |_0^2|\eta _k\widehat{\eta _k}|=|s_\alpha |_0^{2/p}|s_\alpha |_0^{2/q}|\eta _k\widehat{\eta _k}|,$$
where $`p`$ is the index conjugate to $`q`$. Apply Holder’s inequality to get
$$|b_\alpha |\left(_X|s_\alpha |_0^2𝑑\mu _0\right)^{1/p}\left(_X|s_\alpha |_0^2|\eta _k\widehat{\eta _k}|^q𝑑\mu _0\right)^{1/q}.$$
But since the $`s_\alpha `$ are orthonormal this gives
$$|b_\alpha |^q_X|s_\alpha |_0^2|\eta _k\widehat{\eta _k}|^q𝑑\mu _0.$$
Summing over $`\alpha `$ and using the asymptotic statement again in the weak form
$$|s_\alpha |_0^2=k^n+O(k^{n1}),$$
we obtain
$$\underset{\alpha }{}|b_\alpha |^q(k^n+Ck^{n1})_X|\eta _k\widehat{\eta _k}|^q𝑑\mu _0.$$
Now
$$\underset{¯}{M}_qk^1(B_q+E_q)$$
where $`B,E`$ are the diagonal matrices with entries $`b_\alpha ,ϵ_\alpha `$ respectively. We have
$$B_qk^{n/q}(1+Ck^1)^{1/q}\eta _k\widehat{\eta _k}_{L^q}k^{n/q}(1+\frac{C}{q}k^1)\eta _k\widehat{\eta _k}_{L^q},$$
and $`E=O(k^{(n/q)1})`$ since the dimension $`N+1`$ is $`O(k^n)`$. This gives
$$\underset{¯}{M}_qk^{(n/q)1}\eta _k\widehat{\eta _k}_{L^q}+O(k^{(n/q)2}),$$
and our result follows from the fact that $`\eta _k\widehat{\eta _k}=(S(\omega _0)\widehat{S})+O(k^1)`$.
## 4 The finite-dimensional argument
Consider a $`𝐂^{}`$ action on $`𝐂𝐏^N`$ induced by a $`1`$-parameter subgroup $`\rho :𝐂^{}GL(N+1)`$. Let $`V`$ be a smooth $`n`$ dimensional projective variety and set $`V^t=\rho (t)(V)`$. Then it follows from standard theory that the $`V^t`$ converge as $`t0`$ to some algebraic cycle $`V^0`$. Thus $`V^0`$ is a formal sum
$$V^0=m_iW_i$$
where the mulptiplicities $`m_i`$ are positive integers and $`W_i`$ are irreducible $`n`$-dimensional projective varieties. This convergence can be understood at various levels, but the crucial point for us is that the $`V^t`$ converge to $`V^0`$ in the sense of currents. So for any smooth test form $`\varphi `$
$$_{V^t}\varphi _{V^0}\varphi =m_i_{W_i}\varphi .$$
Now suppose that $`\rho `$ maps $`S^1𝐂^{}`$ to the unitary group. Thus the infinitesimal generator $`A`$ of $`\rho `$ is a Hermitian matrix; as usual, we write $`\underset{¯}{A}`$ for the trace-free part of $`A`$. Let $`h`$ be the real valued function
$$h=\underset{\alpha \beta }{}A_{\alpha \beta }h_{\alpha \beta }$$
on $`𝐂𝐏^N`$. This is a Hamiltonian for the action of $`S^1`$ on the projective space. Let
$$I(A,V^0)=_{V(0)}h𝑑\mu _{FS}=\frac{1}{n!}_{V(0)}h\omega _{FS}^n,$$
(14)
where $`\omega _{FS}`$ is the Fubini-Study form, and set
$$FCh(A,V^0)=\frac{\mathrm{Vol}(V^0)}{N+1}\mathrm{Tr}(A)I(A,V^0),$$
where the volume of $`V^0`$ has the obvious meaning.
###### Proposition 2
Suppose that $`FCh(A,V^0)<0`$. Then for any conjugate indices $`p,q`$ we have
$$\underset{¯}{M}(V)_q\frac{FCh(A,V(0))}{\underset{¯}{A}_p}$$
To prove this we consider the function
$$f(t)=\mathrm{Tr}(\underset{¯}{A}\underset{¯}{M}(V^t)),$$
for $`t𝐑^{}`$. The crucial point is that this is increasing with $`t`$. From one point of view this follows from the fact that $`\underset{¯}{M}`$ is a moment map, and from the general theory of such maps, see Section 6.5.2 in , for example. If $`\mu `$ is a moment map for an isometric action of a group $`G`$ on a Kahler manifold we have, for any $`a\mathrm{Lie}(G)`$,
$$\frac{d}{ds}\mu (\mathrm{exp}(sa)X),a=|\frac{d}{ds}\mathrm{exp}(sa)X|^20.$$
Then our assertion follows (taking $`t=e^s`$) from the fact established in that $`\underset{¯}{M}`$ is a moment map for the $`SU(N+1)`$ action on the set of varieties projectively equivalent to $`V`$, with a suitable Kahler structure. An equivalent fact was used in the earlier work of Zhang , and the relation between this and the constructions of is explained in , . This monotonicity property is often stated in the literature as the convexity of a certain function, and other direct proofs are given in , Prop.1, and , Lemma 3.1, so we will not discuss the matter further here. Given this monotonicity property we argue as follows. We have
$$\mathrm{Tr}(AM(V^t))=_{V(e^t)}h_A𝑑\mu _{FS}$$
and $`\underset{¯}{M}=M\frac{\mathrm{Vol}(V)}{N+1}1`$. So
$$f(t)=\mathrm{Tr}(\underset{¯}{A}\underset{¯}{M}(V^t))=\mathrm{Tr}(A\underset{¯}{M}(V^t))=_{V^t}h_A𝑑\mu _{FS}\frac{\mathrm{Vol}}{N+1}\mathrm{Tr}A.$$
Hence the limit of $`f(t)`$ as $`t0`$ is $`FCh(A,V^0)`$, which is positive by hypothesis. Since the function $`f`$ is increasing
$$|f(t)|=f(t)FCh(A,V^0)$$
for all $`t`$. In particular, taking $`t=1`$, we have
$$\mathrm{Tr}(\underset{¯}{A}\underset{¯}{M}(V))FCh(A,V^0).$$
Now use the fact that for, any Hermitian matrices $`S,T`$
$$|\mathrm{Tr}(ST)|S_pT_q,$$
to obtain the required result.
We remark that all of this dicussion can be placed in the context of the action of $`GL(N+1)`$ on the space of Chow vectors, as explained in , , and the criterion $`FCh(A,V^0)<0`$ is the standard Hilbert-Mumford criterion for a destabilising $`1`$-parameter subgroup. On the other hand the criterion can also be placed in a dynamical, Riemannian geometry context. Let $`P`$ be a compact Riemannian manifold, $`SP`$ a submanifold and let $`h`$ be a real-valued function on $`P`$. Let $`\varphi _s`$ be the flow generated by the ($`h`$-decreasing) gradient of $`h`$ and suppose that the $`\varphi _s(S)`$ converge to some varifold $`S^{}`$ as $`s\mathrm{}`$. Say that $`S`$ has property (\*) with respect to $`h`$ if the mean value of $`h`$ over $`S^{}`$ is less than the mean value of $`h`$ on $`P`$. Then a variety $`V𝐂𝐏^N`$ is Chow stable if it has property (\*) with respect to all the functions $`h_A`$. This is because
$$FCh(A,V^0)=\mathrm{Vol}(V^0)\left(\frac{I}{\mathrm{Vol}(V^0)}\frac{\mathrm{Tr}(A)}{N+1}\right),$$
and $`\mathrm{Tr}(A)/N+1`$ is the mean value of $`h_A`$ on $`𝐂𝐏^N`$.
## 5 Algebro-geometric asymptotics
We now go back to a test configuration $`𝒳`$ for $`(X,L)`$. We can suppose that the parameter $`r`$ is $`1`$. The essential point is that for large enough $`k`$ the configuration can be realised by a $`𝐂^{}`$-action on an ambient projective space. That is to say, $`𝒳`$ can be embedded as a $`𝐂^{}`$-invariant subscheme in the product $`𝐏(U_k^{})\times 𝐂`$, extending the embedding of the central fibre $`X_0𝒳`$ by the complete linear system $`U_k`$ in $`𝐏(U_k^{})=𝐏(U_k^{})\times \{0\}.`$ This is explained in . In essence we consider the $`𝐂^{}`$-equivariant bundle $`𝒰=\pi _{}^k`$ over $`𝐂`$ and pick an equivariant trivialisation $`𝒰𝐂\times U_k`$. There is no loss of generality in supposing that $`^k`$ is very ample on all fibres for all $`k1`$; then we obtain the embedding of $`𝒳`$ from the fibrewise embeddings of $`X_t`$ in $`𝐏(𝒰_t^{})`$ under this trivialisation. Let $`|X_0|`$ denote the cycle in $`𝐏(U_k^{})`$ associated to the scheme $`X_0`$. Then we are precisely in the situation considered above with a $`1`$-parameter subgroup acting on $`𝐏^{N+1}`$, generated by $`A_k`$. If we write $`V^t`$ for the image of $`X_t\times \{t\}𝒳𝐏\times 𝐂`$ under the projection map to $`𝐏`$ then $`V^t=\rho (t)(V_1)`$ and the $`V^t`$ converge to the cycle $`|X_0|`$ as $`t0`$.
We need a general fact about the choice of equivariant trivialisation of the bundle $`𝒰`$.
###### Lemma 2
Let $`E`$ be $`𝐂^{}`$-equivariant bundle over $`𝐂`$ and let $`H`$ be a Hermitian metric on the fibre $`E_1`$. Then there is an equivariant trivialisation $`E𝐂\times E_0`$ which takes $`H`$ to a Hermitian metric on the central fibre $`E_0`$ which is preserved by the action of $`S^1𝐂^{}`$ on $`E_0`$.
To see this consider the weight spaces $`E_0=V_i`$ say, where $`𝐂^{}`$ acts with weight $`w_i`$ on $`V_i`$ and the ordering is chosen so that $`w_1<w_2<\mathrm{}`$. Let $``$ be the flag
$$V_1V_1V_2V_1V_2V_3\mathrm{},$$
in $`E_0`$ and let $`\alpha :E_0E_0`$ be a linear map which preserves the flag $``$. Thus $`\alpha `$ has a block matrix description $`\alpha _{ij}:V_iV_j`$ with $`\alpha _{ij}=0`$ for $`i>j`$. For any $`t𝐂`$ we define $`\alpha (t)`$ with blocks $`t^{w_jw_i}\alpha _{ij}`$. This defines a $`𝐂^{}`$ equivariant automorphism of the trivial bundle $`𝐂\times E_0`$ equal to $`\alpha `$ on the fibre $`\{1\}\times E_0`$. Conversely, all equivariant automorphisms arise in this way. What this means is that in the fibre $`E_1`$ of our equivariant bundle $`E`$ there is a canonical flag and a choice of equivariant trivialisation is equivalent to a choice of compatible direct sum splitting of $`E_1`$. Thus the proof of the Lemma is simply to take the direct sum splitting furnished by the succesive orthogonal complements in the flag in $`E_1`$ using the metric $`H`$, and then take the corresponding equivariant trivialisation.
Applying the Lemma to the bundle $`𝒰`$ we see that, given any metric on $`H^0(X,L^k)`$, we may choose our representation of $`𝒳`$ so that the $`1`$-parameter subgroup takes the circle to unitary transformations of $`U_k`$, with respect to the metric arising from the identification of $`U_k`$ with $`H^0(X,L^k)`$. (Of course we can choose a basis so that the given metric is identified with the standard one on $`𝐂^{N+1}`$.) Then we have a numerical invariant
$$I_k=I(A_k,|X_0|)=_{|X_0|}h𝑑\mu _{FS},$$
as above.
###### Proposition 3
The integral $`I_k`$ is equal to the leading term $`b_1k^{n+1}`$ in the Hilbert polynomial for $`\mathrm{Tr}(A_k)`$.
Assuming this for the moment we go on to complete the proof of our main result. Observe that it is equivalent to consider a fixed embedding $`XV=V^1`$ with a variable metric on the underlying vector space $`𝐂^{N+1}`$ (as we are doing here) or to consider a fixed metric on $`𝐂^{N+1}`$ and varying the embedding by projective transformations (as we considered in Section 3). Set
$$FCh_k=FCh(A_k,|X_0|)=\frac{\mathrm{Vol}(X,k\omega )}{\mathrm{dim}U_k}\mathrm{Tr}A_kI_k,$$
as in the previous section. Thus
$$FCh_k=\left(\frac{a_0k^n}{a_0k^n+a_1k^{n1}+\mathrm{}}\right)(b_0k^{n+1}+b_1k^n+\mathrm{})b_0k^{n+1}.$$
So
$$FCh_k=(b_1\frac{a_1b_0}{a_0})k^n+O(k^{n1})=F_𝒳k^n+O(k^{n1}).$$
(15)
Now Theorem 2 is vacuous if $`F_𝒳0`$, so suppose that $`F_𝒳<0`$. Then (15) means that $`FCh_k<0`$ for large $`k`$. Thus we can apply Proposition 2 to deduce that for the embedding $`XV=V(1)`$
$$\underset{¯}{M}(V)_q\frac{FCh_k}{\underset{¯}{A}_k_p}.$$
If $`p`$ is an even integer then $`A_k_p^p=\mathrm{Tr}\underset{¯}{A}^p`$, so
$$\frac{FCh_k}{\underset{¯}{A}_k_p}=k^{(n/q)1}\frac{F_𝒳}{N_p}+O(k^{n/q)2}).$$
Thus
$$\underset{¯}{M}(V)_qk^{(n/q)1}\frac{F_𝒳}{N_p}+O(k^{(n/q)2}).$$
(16)
Then Theorem 2 follows from this lower bound combined with Proposition 1, since if there was a Kahler metric $`\omega _0`$ with $`S\widehat{S}_{L^q}<\widehat{\mathrm{\Psi }}_p(𝒳)`$ we would get, for large $`k`$, an embedding $`V`$ with $`\underset{¯}{M}(V)_q<k^{(n/q)1}F_𝒳/N_p`$ in contradiction to (16).
### 5.1 Proof of Proposition 3
First observe that if we have proved that $`\mathrm{Tr}(A_k)k^{n+1}I_1`$ then replacing $`L`$ by $`L^s`$ it will follow that $`\mathrm{Tr}(A_{sk})(sk)^{n+1}I_s`$ and so that $`I_s=s^{n+1}I_1`$. (This can also be seen directly. If we asume that the powers of sections in $`H^0(X_0,)`$ generate $`H^0(X_0,^s)`$ then a choice of metric on the first space yields a natural metric on the second. With these metrics the integrand defining $`I_s`$ is $`s^{n+1}`$ times that defining $`I_1`$, pointwise on $`|X|`$.) To sum up, our task is to establish the asymptotic relation
$$\mathrm{Tr}(A_k)k^{n+1}I,$$
(17)
where we recall that
* $`X_0`$ is a $`𝐂^{}`$-equivariant subscheme of $`𝐂𝐏^N`$ for the action generated by $`A=A_1`$ on $`𝐂^{N+1}=H^0(X_0,𝒪(1))`$;
* $`A_k`$ is the generator of the induced action on $`H^0(X_0,𝒪(k))`$;
* We fix an $`S^1`$-invariant metric on $`𝐂^{N+1}`$. Then $`I`$ is the integral of $`h\omega _{FS}^n/n!`$ over the cycle $`|X_0|`$ associated to $`X_0`$, where $`\omega _{FS}`$ is the Fubini-Study metric on $`𝐂𝐏^N`$ and $`h`$ is the function on $`𝐂𝐏^N`$ associated to $`A`$.
If $`X_0`$ is smooth the relation (17) is rather standard. The key point is that the function $`h`$ is a Hamiltonian for the $`S^1`$ action on $`𝐂𝐏^N`$. The desired result is then obtained from the equivariant Riemann-Roch Theorem and the de Rham model for equivariant cohmology, as explained by Atiyah and Bott in ; see the discussion in Section 2 of . Our problem is to extend this discussion to the case when $`X_0`$ is a singular variety or scheme. What we do take as known is the corresponding non-equivariant result. That is, if $`Z𝐂𝐏^M`$ is an $`m`$-dimensional projective scheme then
$$\mathrm{dim}H^0(Z,𝒪(k))\frac{D}{m!}k^m,$$
(18)
where $`D`$ is the degree of $`Z`$, which is the integral over the cycle $`|Z|`$ of $`\omega _{FS}^m`$. (The proof of this can be reduced to the case $`m=0`$ by taking hyperplane sections. Thus the assertion is essentially that the notions of multiplicity defined algebraically or by currents agree.) We will explain how to prove the equivariant result by reducing to (18).
We consider the following general situation. Let $`PB,QF`$ be a pair of principle $`S^1`$-bundles over manifolds $`B,F`$. So we have $`S^1`$ actions $`\sigma _P,\sigma _Q`$ say and vector fields $`v_P,v_Q`$ on $`P,Q`$. Suppose in addition that $`QF`$ is an $`S^1`$-equivariant circle bundle, so we have another action $`\rho `$ on $`Q`$, commuting with $`\sigma _Q`$, and another vector field, $`w`$ say, on $`Q`$. We also denote the induced action on $`F`$ by $`\rho `$. Now we can form the associated bundle
$$M=P\times _{\sigma _P,\rho }F.$$
Thus $`M`$ is a bundle over $`B`$ with fibre $`F`$. We can also form
$$\mathrm{\Pi }=P\times _{\sigma _P,\rho }Q.$$
Then $`\mathrm{\Pi }`$ is a circle bundle over $`M`$. Suppose now that we have connections on the bundles $`PB,QF`$ and that the connection on $`Q`$ is preserved by $`\rho `$. Thus we have $`1`$-forms $`\alpha _P,\alpha _Q`$ on $`P,Q`$ and
$$\alpha _P(v_P)=1,\alpha _Q(v_Q)=1,\alpha _Q(w)=H,$$
where $`H`$ is an $`S^1\times S^1`$-invariant function on $`Q`$, which can also be regarded as a function on $`F`$ or on $`M`$. Now the connection on $`P`$ defines a splitting of the tangent bundle of the associated bundle $`M`$, which we express rather loosely as
$$TM=TFTB.$$
Via this decomposition, the curvature forms $`\omega _P,\omega _Q`$ can naturally be regarded as $`2`$-forms on $`M`$. The key point is that there is a connection on $`\mathrm{\Pi }M`$ with curvature
$$\mathrm{\Omega }=\omega _Q+H\omega _P.$$
To see this we pull everything back to $`P\times Q`$ where we consider the $`1`$-form
$$\beta =\alpha _Q+H\alpha _P,$$
(making an obvious simplification in notation). Then $`\beta `$ vanishes on the generator $`v_P+w`$ of the action $`(\sigma _P,\rho )`$ and so descends to a form $`\underset{¯}{\beta }`$ on $`\mathrm{\Pi }=P\times _{\sigma _P,\rho }Q`$. Since $`\beta (v_Q)=1`$ the equivariant $`1`$-form $`\underset{¯}{\beta }`$ furnishes a connection on $`\mathrm{\Pi }`$. We have (again making various abuses of notation)
$$d\beta =\omega _Q+H\omega _P+dH\alpha _P.$$
This gives the lift of the curvature form of the connection on $`\mathrm{\Pi }`$ to $`P\times Q`$. The definition of the horizontal subspace defining the splitting $`TM=TFTB`$ means that $`\alpha _P`$ vanishes on the vectors representing $`TM`$, so we see that the connection $`\underset{¯}{\beta }`$ has curvature $`\mathrm{\Omega }`$.
Now suppose that $`B`$ is a compact oriented $`2r`$-manifold and let $`S`$ be an $`S^1`$-invariant oriented $`2n`$-dimensional submanifold of $`F`$, not necessarily closed. This defines a corresponding submanifold $`\stackrel{~}{S}`$ in $`M`$ (so $`\stackrel{~}{S}`$ is a bundle over $`B`$ with fibre $`S`$). Then it is clear, by integrating over the fibres, that
$$\frac{1}{(n+r)!}_{\stackrel{~}{S}}\mathrm{\Omega }^{n+r}=\frac{1}{r!}_B\omega _P^r\frac{1}{n!}_SH^n\omega _Q^n.$$
(19)
We apply this to the case when $`B`$ is the Riemann sphere $`𝐂𝐏^1`$ and $`F`$ is $`𝐂𝐏^N`$ with the bundles $`P,Q`$ being the Hopf fibrations. We take the action $`\rho `$ to be that induced by the given action on $`𝐂^{n+1}`$. Then the function $`H`$ considered above becomes the function $`h`$ by the standard discussion of Hamiltonians and equivariant classes, as in . We then have a manifold $`M`$, which is just the projectivization of a vector bundle over $`𝐂𝐏^1`$, and a circle bundle $`\mathrm{\Pi }M`$ which clearly corresponds to a holomorphic line bundle say $`𝒱M`$. Notice that, going back to our original data, we could change the generator $`A`$ to $`A+\nu 1`$ for any integer $`\nu `$. This changes the function $`h`$ to $`h+\nu `$ and, by applying (18), does not affect the truth of the result we seek. So we can suppose that $`h`$ is positive and this means that $`𝒱`$ is a positive line bundle over $`M`$. Thus we can embed $`M`$ as a projective variety in some $`𝐂𝐏^\mu `$, with $`𝒱^s=𝒪(1)`$, and there will be no loss in generality in supposing that $`s=1`$.
Now consider our $`𝐂^{}`$-invariant scheme $`X_0𝐂𝐏^N`$. Clearly we get a corresponding scheme $`Z`$ inside $`M`$, fibering as $`\pi :Z𝐂𝐏^1`$ with fibre $`X_0`$. For any $`k`$ we can identify $`H^0(Z;𝒱^k)`$ with the sections of the vector bundle $`\pi _{}(𝒱^k)`$ over $`𝐂𝐏^1`$. Now take the eigenspace decomposition
$$H^0(X_0,𝒪(k))=V_i,$$
where $`A_k`$ acts as $`w_i`$ on $`U_i`$. Then
$$\pi _{}(𝒱^k)=U_i𝒪(w_i).$$
It follows that
$$\mathrm{dim}H^0(Z,𝒱^k)=\mathrm{dim}U_i(w_i+1)=\mathrm{Tr}A_k+\mathrm{dim}H^0(X_0,𝒪(k)).$$
We apply (18) to the projective scheme $`ZM𝐂𝐏^M`$. This shows that
$$\mathrm{Tr}A_kDk^{n+1}$$
where $`D`$ is the degree of $`Z`$. This degree is given by integrating $`\mathrm{\Omega }^{n+1}`$ over the cycle $`|Z|`$, for any smooth form $`\mathrm{\Omega }`$ on $`M`$ representing $`c_1(𝒱)`$. Taking the form given by our construction above, taking $`S`$ to be the smooth points in $`|Z|`$ and taking due account of multiplicity, we see that the degree is given by the integral $`I`$, and our result follows.
Note that the same argument (taking $`B`$ to be $`𝐂𝐏^r`$) shows that for any positive integer $`r`$
$$\mathrm{Tr}A_k^rk^{n+r}_{|X_0|}h^rd\mu _{FS},\mathrm{Tr}\underset{¯}{A}_k^rk^{n+r}_{|X_0|}(h\widehat{h})^rd\mu _{FS}.$$
(20)
Thus the invariant $`N_p`$ is the $`L^p`$ norm of $`h\widehat{h}`$ on $`|X_0|`$. One consequence of this is that we can extend Theorem 2 to the case $`q=1,p=\mathrm{}`$, defining $`N_{\mathrm{}}(𝒳)=h\widehat{h}_{L^{\mathrm{}}(|X_0|)}`$. It would be interesting to extend Theorem 2 to general real exponents $`p,q`$ with
$$N_p=h\widehat{h}_{L^p(|X_0|)}.$$
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# Spin wave contribution to the nuclear spin-lattice relaxation in triplet superconductors
## Abstract
We discuss collective spin wave excitations in triplet superconductors with an easy axis anisotropy for the order parameter. Using a microscopic model for interacting electrons we estimate the frequency of such excitations in Bechgaard salts and ruthenate superconductors to be one and twenty GHz respectively. We introduce an effective bosonic model to describe spin-wave excitations and calculate their contribution to the nuclear spin lattice relaxation rate. We find that in the experimentally relevant regime of temperatures, this mechanism leads to the power law scaling of $`1/T_1`$ with temperature. For two and three dimensional systems the scaling exponents are three and five respectively. We discuss experimental manifestations of the spin wave mechanism of the nuclear spin lattice relaxation.
Nuclear magnetic resonance (NMR) is a powerful tool for analyzing ordered electron states in solids. NMR analysis has been successfully applied to study magnetic insulators (see e.g. beeman68 and references therein) as well as several classes of unconventional superconductors, including high Tc cuprates X1 , heavy fermion materials tou03 , ruthenates X2 , and organic superconductors lefebvre00 . In particular, NMR experiments have been useful for analyzing the symmetry of the SC order parameter lee02 and for clarifying the structure of the phase diagram in systems with competing orders lee05 .
A common feature of the NMR experiments in certain families of triplet superconductors (TSC) is the power law temperature dependence of the nuclear spin lattice relaxation rate. Bechgaard salts lee2000 , ruthenates X2 , and heavy fermion materials X4 showed $`1/T_1T^3`$ at low temperatures and for small magnetic fields. Such behavior has usually been interpreted as a signature of nodes in the quasiparticle gap on the Fermi surface. Indeed, point and line nodes should lead to $`T^5`$ and $`T^3`$ scaling of $`1/T_1`$ respectively. In several cases, however, we have reasons to doubt the presence of nodes in the TSC order parameter. For example, in quasi one dimensional Bechgaard salts the natural order parameter has different signs on the two sheets of the Fermi surface and no nodal points X5 . Also in ruthenates, the order parameter that is consistent with spontaneous time reversal breaking X6 and the quasi two dimensional nature of these materials corresponds to a constant quasiparticle gap on the entire Fermi surface Mackenzie . In this paper we consider a mechanism of the nuclear spin lattice relaxation that is not due to Bogoliubov quasiparticles but due to collective spin wave (SW) excitations of the TSC order parameter. We demonstrate that in the experimentally relevant regime of temperatures, this mechanism also leads to the power law scaling of $`1/T_1`$ with temperature.
Our starting point is the Moriya relation moriya63 for the nuclear spin lattice relaxation rate
$`{\displaystyle \frac{1}{TT_1}}={\displaystyle \frac{2g_\mathrm{N}^2|A|^2}{(g_{\mathrm{eff}}\mu _B)^2}}{\displaystyle d^dq\frac{\chi _H^{\prime \prime }(q,\omega _N)}{\omega _N}}`$ (1)
Here $`A`$ describes the strength of hyperfine interactions between nuclear spins and conduction electrons, $`g_N`$ is a gyromagnetic ratio of the nucleus, $`g_{\mathrm{eff}}`$ is an effective gyromagnetic ratio of conducting electrons, $`\mu _B`$ is a Bohr magneton, $`\chi _H^{\prime \prime }(q,\omega _N)`$ is the imaginary part of the transverse (i.e. perpendicular to the magnetic field) electron spin susceptibility taken at the nuclear Larmor frequency $`\omega _N`$. In the case of a perfect spin SU(2) symmetry, linearly dispersing SW excitations exist down to arbitrarily small energies. In real materials there is always spin anisotropy which gives rise to a finite gap for spin excitations, $`\omega _0`$. Below, we estimate the value of $`\omega _0`$ to be tens of millikelvin for Bechgaard salts and hundreds of millikelvin for the ruthenates. This is much larger than the nuclear Larmor frequency, $`\omega _N`$, but smaller than the typical temperature used in experiments. When $`\omega _0`$ is much larger than $`\omega _N`$, creation and annihilation of individual SWs does not affect $`\chi ^{\prime \prime }(\omega _N)`$. However, there is a contribution due to the scattering of thermally excited SW excitations. Let $`\rho (E)`$ be the density of states for SW excitations and $`n(E)=(\mathrm{exp}(E/k_BT)1)^1`$ be the Bose distribution function. From the second order perturbation theory we have $`\chi _{zz}^{\prime \prime }(\omega _N)\rho (E)\rho (E+\omega _N)[n(E)n(E+\omega _N)]𝑑E`$. The characteristic energy scale in this integral is set by the temperature $`T`$. Since $`T\omega _0`$ we can assume linear dispersion of SW excitations and take $`\rho (E)E^{d1}`$, where $`d`$ is the number of spatial dimensions. Using $`\omega _NT`$, we have $`\chi _{zz}^{\prime \prime }(\omega _N)\omega _NE^{2d2}(n/E)𝑑ET^{2d2}`$. Combining this result with the Moriya relation (1) we obtain $`1/T_1T^{2d1}`$. This simple analysis does not take into account coherence factors in the expression for $`\chi ^{\prime \prime }`$. Below we demonstrate that coherence factors do not modify the scaling exponent of the nuclear spin lattice relaxation rate in TSC. This is in contrast to antiferromagnetic systems that also have linearly spin waves, but in which coherence factors contribute an additional $`1/T^2`$ factor to $`1/T_1`$ beeman68 .
For a detailed analysis of the nuclear spin lattice relaxation rate we introduce an effective model that captures the essence of triplet superconductivity and allows us to analyze collective excitations. A simple picture of the TSC state corresponds to binding electrons into Cooper pairs with spin one and momentum zero and Bose condensing such pairs. In an effective bosonic model one can neglect details of the orbital nature of Cooper pairs and consider them as “elementary” particles. Interactions are important for the correct description of collective excitations, thus we are led to the Hubbard type model for spin one bosons demler02
$``$ $`=`$ $`t{\displaystyle \underset{ij\sigma }{}}(a_{i\sigma }^{}a_{j\sigma }+a_{j\sigma }^{}a_{i\sigma })\delta r{\displaystyle \underset{i}{}}a_{i0}^{}a_{i0}`$ (2)
$`+`$ $`{\displaystyle \frac{U_0}{2}}{\displaystyle \underset{i}{}}n_i(n_i1)+{\displaystyle \frac{U_2}{2}}{\displaystyle \underset{i}{}}\stackrel{}{S}_i^2\mu {\displaystyle \underset{i}{}}n_i`$
Here $`a_{i\sigma }^{}`$ creates a boson on site $`i`$ with spin $`\sigma =\{1,0,1\}`$. Operators $`n_i`$ and $`\stackrel{}{S}_i`$ describe the number of atoms and the total spin on site $`i`$: $`n_i=_\sigma a_{i\sigma }^{}a_{i\sigma }`$, $`\stackrel{}{S}_i=_{\sigma ,\sigma ^{}}a_{i\sigma }^{}\stackrel{}{T}_{\sigma \sigma ^{}}a_{i\sigma ^{}}`$, where $`\stackrel{}{T}_{\sigma \sigma ^{}}`$ are the usual spin matrices for spin one particles. The first term in the Hamiltonian (2) describes tunneling of Cooper pairs between neighboring lattice sites $`i`$ and $`j`$. An important aspect of the model (2) is the presence of two types of interaction terms. The third term in (2) depends only on the number of particles on each site and is the same as for spinless bosons. The fourth term in (2) gives spin dependence to the interaction (without breaking the spin SU(2) symmetry) and is a novel feature of spinful Cooper pairs. The sign of $`U_2`$ determines the difference between unitary ($`U_2>0`$) and nonunitary ($`U_2<0`$) triplet superconductors. It is generally believed that triplet pairing between fermions leads to unitary Cooper pairs Mackenzie . Thus, from now on we assume that $`U_2`$ is positive. For concreteness, we consider a $`d`$ dimensional ($`d=2,3`$) hypercubic lattice. Our results, however, do not depend on the precise lattice structure. The second term in equation (2) introduces easy axis anisotropy for the order parameter by making the condensation of $`a_0`$ to be energetically favorable. A state with finite $`a_0`$ corresponds to the unitary state of Cooper pairs with the $`𝐝`$ vector pointing along the $`z`$ axis (see eq. (8) for the definition of the $`𝐝`$ vector).
In the mean field approximation, we take $`a_0=\mathrm{\Psi }_0`$ with $`|\mathrm{\Psi }_0|^2=(\mu +zt+\delta r)/U_0`$ and $`z=2d`$ being the coordination number. Without loss of generality, we can take $`\mathrm{\Psi }_0`$ to be real. Fluctuations in the phase of $`a_0`$ correspond to the density (Bogoliubov) mode. To find SW excitations we need to consider $`a_\pm `$ operators. In the Hamiltonian (2) we replace $`a_0`$ by its expectation value, take the terms quadratic in $`a_\pm `$, and perform the Bogoliubov rotation $`a_{i\alpha }=\frac{1}{\sqrt{N}}_\stackrel{}{k}a_{\stackrel{}{k}\alpha }e^{i\stackrel{}{k}\stackrel{}{r}_i}`$, $`a_{\stackrel{}{k}+}=v_k\gamma _{\stackrel{}{k}+}+u_k\gamma _\stackrel{}{k}^{},a_\stackrel{}{k}=u_k\gamma _{\stackrel{}{k}+}^{}+v_k\gamma _\stackrel{}{k}`$. Here $`N`$ is the total number of lattice sites, $`u_k^2+v_k^2=(\xi _k+\mathrm{\Delta })/E_k`$, $`2u_kv_k=\mathrm{\Delta }/E_k`$, $`\xi _k=2t_{n=1}^d\mathrm{cos}(k_nb)+U_0\mathrm{\Psi }_0^2\mu `$. In these equations $`b`$ is the lattice constant and $`\mathrm{\Delta }=U_2\mathrm{\Psi }_0^2`$. We obtain the diagonalized spin-wave Hamiltonian $`_0=_kE_k\left(\gamma _{k+}^{}\gamma _{k+}+\gamma _k^{}\gamma _k\right)`$ with $`E_k^2=\xi _k^2+2\mathrm{\Delta }\xi _k`$. Operators $`\gamma _{k\pm }^{}`$ create SW excitations with $`S_z=\pm 1`$. In the long wavelength limit we find $`E_k^2=\omega _0^2+v_s^2k^2`$ with $`\omega _0^2=\delta r^2+2\mathrm{\Delta }\delta r`$ and $`v_s^2=2tU_2\mathrm{\Psi }_0^2a^2`$.
We need to calculate the electron spin susceptibility in the direction perpendicular to the applied magnetic field. Let $`\theta `$ be the angle between the $`z`$ axis (i.e. the direction of the $`\stackrel{}{d}`$ vector) and the direction of the magnetic field (see Fig. 1). We have
$`\chi _H=\mathrm{sin}^2\theta \chi _{zz}+(1+\mathrm{cos}^2\theta )\chi _{xx}.`$ (3)
It is easy to see that $`\chi _{xx}`$ and $`\chi _{zz}`$ are given by the correlation function of SWs
$`\chi _{xx}(q,\omega )`$ $`=`$ $`2(g_{\mathrm{eff}}\mu _B\mathrm{\Psi }_0)^2(u_q+v_q)^2{\displaystyle \underset{\alpha \beta }{}}D_{\alpha \beta }^R(q,\omega )`$
$`\chi _{zz}(q,\omega )`$ $`=`$ $`(g_{\mathrm{eff}}\mu _B)^2{\displaystyle \underset{\alpha \beta }{}}{\displaystyle }{\displaystyle \frac{d^dkd\mathrm{\Omega }}{(2\pi )^{d+1}}}U_{\alpha \beta }^2(k,k+q)\times `$ (4)
$`D_{\alpha \alpha }^R(k+q,\mathrm{\Omega }+\omega )D_{\beta \beta }^A(k,\mathrm{\Omega })`$
where we introduced the Nambu-Gorkov type notations $`D_{\alpha \beta }^{R,A}(q,\omega )=𝑑te^{i\omega t}\theta (\pm t)\mathrm{T}\mathrm{\Psi }_{q\alpha }(t)\mathrm{\Psi }_{q\beta }^{}(0)`$, with $`\mathrm{\Psi }_{q\alpha }=\{\gamma _{q+},\gamma _q^{}\}^T,`$ and $`U_{\alpha \beta }(k,k^{})=\delta _{\alpha \beta }(v_kv_k^{}u_ku_k^{})+(1\delta _{\alpha \beta })(u_kv_k^{}v_ku_k^{}).`$
Note that there is a qualitative difference in calculating $`\chi _{zz}`$ and $`\chi _{xx}`$. A non-uniform magnetic field in the $`z`$ direction can scatter the existing thermally excited SWs. Thus we find finite imaginary part of $`\chi _{zz}`$ at small frequencies by considering the quadratic Bogoliubov Hamiltonian. We obtain
$`\chi _{zz}^{\prime \prime }(\omega _N,q)=2(g_{\mathrm{eff}}\mu _B)^2{\displaystyle \frac{\omega _N}{T}}b^d{\displaystyle \frac{d^dk}{(2\pi )^d}}`$ (5)
$`\times `$ $`n(E_k)(n(E_{k+q})+1)\delta (E_kE_{k+q}).`$
By contrast, a non-uniform magnetic field in the $`x`$ direction can only create or annihilate SW excitations. However, energies of these excitations cannot be smaller than $`\omega _0`$. Hence, if we limit ourselves to the quadratic Hamiltonian for SWs, we find that $`\chi _{xx}^{\prime \prime }`$ is identically zero for frequencies smaller than $`\omega _0`$ at any temperature. To get finite $`\chi _{xx}^{\prime \prime }`$ at small frequencies we need to consider interactions between SWs. Taking quartic terms in equation (2) and using definitions of SW operators, we obtain the interaction terms between SW excitations. These can be used to calculate self-energies for SW excitations as shown in Fig 2. We find unpublished
$`\chi _{xx}^{\prime \prime }(q,\omega _N)=(g_{\mathrm{eff}}\mu _B\mathrm{\Psi }_0)^2{\displaystyle \frac{\omega _N}{4\mathrm{\Delta }T}}(U_0+3U_2)^2b^{2d}`$ (6)
$`\times `$ $`{\displaystyle \frac{d^dk_1}{(2\pi )^d}\frac{d^dk_2}{(2\pi )^d}\frac{E_{k1}^2+E_{k2}^2+E_{k3}^2}{E_{k1}E_{k2}E_{k3}}}`$
$`\times `$ $`(n(E_{k1})+1)(n(E_{k2})+1)n(E_{k3})`$
$`\times `$ $`\delta (E_{k1}+E_{k2}E_{k3}).`$
where $`\stackrel{}{k}_3=\stackrel{}{k}_1+\stackrel{}{k}_2`$. We emphasize that equations (5) and (6) apply only in the low frequency limit $`\omega _N\omega _0`$ which is relevant for experiments.
Expressions (5) and (6) show that in equation (3) contributions to $`1/T_1`$ due to $`\chi _{zz}^{\prime \prime }`$ and $`\chi _{xx}^{\prime \prime }`$ scale as $`T^{2d1}`$ and $`T^{3d2}`$ respectively. For two and three dimensional systems the $`\chi _{zz}^{\prime \prime }`$ contribution dominates (one can check that this conclusion remains when we include prefactors). At first glance this result appears surprising. Firstly, the real static susceptibility is finite in the direction perpendicular to the $`𝐝`$ vector but is zero along it. Secondly, in the case of full SO(3) symmetry, creation and annihilation of individual SWs contribute to $`\chi _{xx}^{\prime \prime }`$ and $`\chi _{yy}^{\prime \prime }`$ at small frequencies, but have no effect on $`\chi _{zz}^{\prime \prime }`$. A crucial part of our analysis is the existence of a spin gap $`\omega _0`$ which is much larger than the nuclear Larmor frequency $`\omega _N`$. In this case $`\chi ^{\prime \prime }(\omega _N)`$ does not have contributions due to creation or annihilation of individual SWs. For $`\chi _{zz}^{\prime \prime }`$ we take thermally excited SWs and scatter them by the magnetic field. For $`\chi _{xx}^{\prime \prime }`$ we also need thermally excited SWs but in addition we must rely on interactions between them. SWs are pseudo-Goldstone modes. At low energies interactions between them are suppressed. This gives rise to the smallness of $`\chi _{xx}^{\prime \prime }`$ relative to $`\chi _{zz}^{\prime \prime }`$.
To summarize, for two and three dimensional systems, we find that as long as $`\theta `$ is not anomalously small, the SW contribution to the nuclear spin lattice relaxation rate is given by
$`{\displaystyle \frac{1}{T_1}}=\mathrm{sin}^2\theta |A|^2g_\mathrm{N}^2{\displaystyle \frac{b^d}{v_s^{2d}}}{\displaystyle \frac{S_d^2}{4\pi ^2}}|B_{2d2}|T^{2d1},`$ (7)
Here $`S_d=2\pi ^{d/2}/\mathrm{\Gamma }(d/2)`$ is the surface area of a unit sphere, and $`B_n`$ are Bernoulli numbers, $`B_2=1/6`$, $`B_4=1/30`$. We remind the readers that equation (7) applies when $`T\omega _0`$. At low temperatures, $`T\omega _0`$, we expect $`1/T_1`$ to start decreasing exponentially, reflecting the exponential suppression in the number of thermally excited SW excitations. Equation (7) also predicts that the nuclear spin lattice relaxation rate should be very sensitive to the direction of the magnetic field. We note, however, that this argument is valid only for magnetic fields that are smaller than the so-called spin-flop magnetic field, $`H_{\mathrm{flop}}`$. In magnetic fields larger than $`H_{\mathrm{flop}}`$, the order parameter $`\stackrel{}{d}`$ will always be perpendicular to the direction of the applied field unpublished .
It is useful to compare contributions to $`1/T_1`$ from magnons to the one from quasiparticles. For concreteness we consider a quasi two dimensional system with the TSC order parameter $`\stackrel{}{d}=\mathrm{\Delta }_0\widehat{z}k_x/k_F`$, that has a line of nodes along the $`\widehat{z}`$ axis. The quasiparticle contribution to the NMR relaxation rate in such a state sigrist91 is given by $`1/T_{1\mathrm{qp}}=\pi ^2|A|^2g_\mathrm{N}^2b^2m^2T^3/6\mathrm{\Delta }_0^2`$. For comparison, we take equation (7) and use the BCS expressions for the velocity of SW excitations (see below). We find $`T_{1\mathrm{mag}}/T_{1\mathrm{qp}}=\pi ^2\mathrm{\Delta }_0^2/ϵ_F^2`$. In a typical superconductor, the value of the quasiparticle gap is much smaller than the electron Fermi energy. Therefore, when we have both gapless quasiparticles and SWs, the quasiparticle contribution will strongly dominate. Only when the quasiparticles are fully gapped out do the magnons provide the dominant contribution to $`1/T_1`$.
Now we outline the key steps of the analysis that allowed us to estimate the characteristic frequency of SW excitations in TSC. We consider a phenomenological BCS type model for interacting electrons
$``$ $`=`$ $`{\displaystyle \underset{k\sigma }{}}ϵ_kc_{k\sigma }^{}c_{k\sigma }{\displaystyle \underset{aq}{}}V_a^1(q)d_a^{}(q)d_a(q)`$
$`d_a^{}(q)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k\alpha \beta }{}}V_a(q)f_kc_{k+q/2,\alpha }^{}(i\sigma _2\sigma _a)_{\alpha \beta }c_{k+q/2,\beta }^{}`$ (8)
Here $`d_a(q)`$ are Fourier components of the TSC order parameter in the direction $`a`$ ($`a=x,y,z`$), $`f_k`$ is an orbital wavefunction (e.g. $`f_k=\mathrm{sign}(k)`$ X5 ), and $`c_{k\sigma }^{}`$ are electron creation operators. Parameters $`V_a`$ describe electron interactions in the p-wave channel. A homogeneous unitary triplet superconducting state is the ground state of the Hamiltonian (8). Assuming easy axis anisotropy with $`V^z>V^x=V^y`$, we find $`d_z(q=0)=\mathrm{\Delta }_0`$ and $`d_x=d_y=0`$. To find SW excitations we consider fluctuations of the TSC order parameter $`d_x(r,t)=d_x^{}(r,t)`$. Integrating out fermions gives an effective action for SW excitations unpublished $`\mathrm{S}_{\mathrm{eff}}\{d_x\}=_{\omega q}D_x(q,\omega )|d_x(q,\omega )|^2`$, where $`q`$ and $`\omega `$ are the wavevector and the real frequency of the SW. For the spherically symmetric Fermi surface in $`d`$ dimensions, we have $`D_x(\omega ,q)=\frac{1}{4}N_0(\omega _0^2\omega ^2+v_s^2q^2)`$. Here $`N_0`$ is the density of states at the Fermi energy, $`v_s^2=v_F^2/d`$, $`v_F`$ is the Fermi velocity, and
$`\omega _0^2=4\mathrm{\Delta }_0^2N_0^1(V_x^1V_z^1)`$ (9)
Zeroes of the function $`D_x`$ correspond to SW excitations. In the case of spin symmetric interactions with $`V_z=V_x`$, SWs are Goldstone modes of broken spin symmetry and are gapless. They have linear dispersion with the velocity $`v_s`$. In the case of the easy axis anisotropy, SWs have a gap $`\omega _0`$.
As a concrete example, we consider the TSC state in Bechgaard salts. We assume that the spin anisotropic part of the interaction in this phase is the same as in the antiferromagnetic state of this material $`\mathrm{\Delta }_{\mathrm{anis}}=\delta J_z_{ij}S_i^zS_j^z`$. Here $`ij`$ corresponds to the nearest neighbor sites and the spin $`z`$ axis points along the crystallographic $`b^{}`$ axis. The antiferromagnetic resonance experiments torrance82 suggest $`\delta J_z=0.01K`$. We can express $`\mathrm{\Delta }_{\mathrm{anis}}`$ in the form similar to equation (8) with $`\delta V^z=\delta V^x=\delta V^y=\frac{1}{2}\delta J_zv_0`$. Here $`v_0`$ is the volume of the unit cell. Assuming that the anisotropic term is a small correction to the spin symmetric interaction, the total interaction entering equation (8) is $`V^a=V+\delta V^a`$. The value of $`V`$ can be estimated from the BCS equation for the transition temperature $`T_c=1.14\omega _{\mathrm{BOS}}e^{1/N_0V}`$. Here $`\omega _{\mathrm{BOS}}`$ is the characteristic frequency of bosons providing electron pairing. Combining all expressions, we find
$`{\displaystyle \frac{\omega _0}{\mathrm{\Delta }_0}}=2(\delta J_zN_0v_0)^{1/2}\mathrm{log}({\displaystyle \frac{1.14\omega _{BOS}}{T_c}})`$ (10)
Taking $`\delta J_z=0.01K`$, $`N_0=210^{33}erg^1cm^3`$, $`\omega _{\mathrm{BOS}}=1000`$K, $`T_c=1.4K`$, $`\mathrm{\Delta }_0=2.5K`$, and $`v_0=360\AA ^3`$ we obtain $`\omega _0`$ around one GHz. There is a simple qualitative argument that supports our result that $`\omega _0`$ for Bechgaard salts lies in the GHz range. This argument relies on a comparison of SW resonances in the antiferromagnetic and superconducting states of these materials. In quasi one-dimensional systems microscopic descriptions of the two states are similar and BCS type models and our analysis of SW excitations can be used in the antiferromagnetic phase as well. From the discussion above, we expect that the ratio of SW energies in the two phases is approximately proportional to the ratio of AF and TSC transition temperatures, i.e. around ten. In the AF phase the SW resonance frequency was measured in the tens of GHz range torrance82 . Hence, in the superconducting state it should be a factor of ten smaller, which brings us into the GHz range.
Similar analysis can be done for the TSC state in $`Sr_2RuO_4`$. Sigrist and coworkers (see e.g. Ref X9 ) showed that spin-orbit coupling in these materials leads to a difference in the transition temperature of various order parameters of the order of two percent. From the BCS expression for $`Tc`$ we have $`\delta T_c/T_c=\delta V/N_0V^2`$. Taking $`\mathrm{\Delta }_0=2.6`$K, from equation (9) we estimate $`\omega _0`$ to be around twenty GHz.
$`Sr_2RuO_4`$ is a quasi two dimensional material. Bechgaard salts have a mixed dimensionality schwartz98 ; lebed05 with easy, intermediate, and hard directions in transport. Hence, they may also exhibit properties of a quasi two dimensional system with respect to SW excitations. Both the ruthenates and Bechgaard salt superconductors exhibit the $`T^3`$ dependence of $`1/T_1`$ at low temperature. Therefore these materials should be good candidates for experimental investigation of the SW mechanism of nuclear spin lattice relaxation discussed in this paper.
In summary, we studied spin wave excitations in triplet superconductors with the easy axis spin anisotropy. We derived an explicit expression for the energy of such excitations and used it to estimate the spin wave energy in Bechgaard salts and ruthenate superconductors. We considered the effects of spin wave excitations on the nuclear spin lattice relaxation rate. We showed that in the experimentally relevant regime of temperatures, $`1/T_1`$ has a power law scaling with temperature, including the $`T^3`$ dependence for two dimensional systems. We showed that the spin wave mechanism predicts a dramatic decrease of $`1/T_1`$ for temperatures lower then the energy of spin wave excitations and leads to dependence of $`1/T_1`$ on the direction of the applied magnetic field.
We are grateful for useful discussions with E. Altman, D. Podolsky, and A. Polkovnikov. This work was supported by the NSF grant DMR-0132874, the Sloan Foundation, and KITP UCSB.
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# Intersecting Non-extreme 𝑝-Branes and Linear Dilaton Background
## I Introduction
Recently it has been shown Clement:2004ii that the general supergravity solution describing a charged $`p`$-brane without naked singularities can be either the standard black asymptotically flat (AF) $`p`$-brane, or the black brane which asymptotically approaches the linear dilaton background (LDB). The LDB is known to be the near-horizon limit of the extremal dilatonic branes; it is relevant for the description of (non-conformal) quantum field theory (QFT) in the Domain-Wall/QFT correspondence Boonstra:1998mp ; Behrndt:1999mk ; Cai:1999xg and (in the case of NS5 branes) little string theories (LST) Aharony:1998ub ; Aharony:1999ks . According to the standard reasoning, the same configuration endowed with an event horizon should describe the thermal version of the same QFT (LST). The asymptotically LDB black branes obtained must thus describe the class of quantum field theories in the thermal phase, or, in the particular case of NS5-branes, the thermodynamics of little string theory. This was extensively studied recently in the case of one-brane solutions HaOb ; BeRo ; BoRo , so it is interesting to present a more general intersecting brane framework.
The simplest example of the LDB endowed with an event horizon is the four-dimensional charged dilatonic black hole with the LDB asymptotics Chan:1995fr ; Clement:2002mb ; Clement:2004yr ; Cai:2004iy . This configuration was also identified as “the horizon plus throat” geometry Giddings:1992kn arising in the near-horizon limit of the near extremal dilatonic black hole Gibbons:1987ps ; Garfinkle:1990qj . The near-horizon limit of the BPS dilatonic black hole is the LDB itself, and in the theories admitting the 1/2 BPS branes the limiting configuration preserves the half of SUSY of the initial theory. The near horizon limit of black dilatonic black hole Clement:2002mb keeps the memory about the event horizon and turns out to be a non-supersymmetric configuration whose BPS limit is the LDB. The brane generalization of this construction has interesting particular cases in ten dimensions Clement:2004ii .
On the other hand, it has been also known that a more general class of solutions consisting of various kinds of black branes can be constructed within the same framework Ohta:1997wp ; Aref'eva:1997nz ; Ohta:1997gw ; Ohta:1997wd ; Ivashchuk:1998jg ; Ohta:2003rr ; Miao:2004bn ; Rama:2005bd ; Bai:2005jr (for the time-dependent branes see Chen:2002yq ; Ohta:2003uw ; Ohta:2003zh ). It is then natural to ask if the above results can be extended to such general solutions. The purpose of this paper is to generalize the construction to the case of intersecting branes (see for review Ga97 ; Lu98 ; Sm02 ). We show that the solutions are restricted to either the asymptotically flat black branes or asymptotically LDB ones and their mixed system if we impose the condition that there are no naked singularities. This generalization is interesting in that it opens the possibility of extending the Domain-Wall/QFT correspondence to more general configurations.
This paper is organized as follows. In the next section, we start with the action for the $`D`$-dimensional gravity coupled to the dilaton and an arbitrary number of form fields of various ranks. We summarize the field equations, metric ansatz and the background for forms. We then derive the most general static solutions in the theory requiring the usual intersection rules for the branes. The obtained solutions have a large number of parameters. In Sec. III, we fix some of them by the requirement that the solutions do not have naked singularities, and show that the resulting solutions consist of either AF black branes or the black branes which asymptotically approach the LDB. In Sec. IV, we transform the solutions to the more familiar coordinates, and show that these reduce to the known solutions of intersecting AF branes and/or branes in LDB. Of course, the latter can be obtained from the known intersecting AF branes by taking the near-horizon limit, but a naive application of this rule produces the solutions consisting only of intersecting branes with the LDB asymptotics; our additional solutions of mixed type of AF and asymptotically LDB black branes can not be simply obtained in such a limit. Not only this, our results show that these are the only solutions which can be obtained by imposing the requirement that the space-time is free of naked singularities.
## II Intersecting Non-extreme $`p`$-branes
We consider the action describing gravity coupled to a dilaton $`\varphi `$ and $`m`$ different $`n_A`$-form fields in $`D`$ dimensions
$$S=d^Dx\sqrt{g}\left(R\frac{1}{2}_\mu \varphi ^\mu \varphi \underset{A=1}{\overset{m}{}}\frac{1}{2n_A!}\mathrm{e}^{a_A\varphi }F_{[n_A]}^2\right).$$
(1)
All form fields are coupled to the unique dilaton with individual coupling constants $`a_A`$. This action is the simplest one to describe intersecting supergravity $`p`$-branes, it incorporates various (truncated) supergravity models for different choice of the parameters $`D,n_A,`$ and $`a_A`$. The corresponding equations of motion read
$`R_{\mu \nu }{\displaystyle \frac{1}{2}}_\mu \varphi _\nu \varphi {\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{\mathrm{e}^{a_A\varphi }}{2(n_A1)!}}\left[(F_{[n_A]}^2)_{\mu \nu }{\displaystyle \frac{n_A1}{n_A(D2)}}F_{[n_A]}^2g_{\mu \nu }\right]`$ $`=`$ $`0,`$ (2)
$`_\mu \left(\sqrt{g}\mathrm{e}^{a_A\varphi }F^{\mu \nu _2\mathrm{}\nu _{n_A}}\right)`$ $`=`$ $`0,`$ (3)
$`{\displaystyle \frac{1}{\sqrt{g}}}_\mu \left(\sqrt{g}^\mu \varphi \right){\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{a_A}{2n_A!}}\mathrm{e}^{a_A\varphi }F_{[n_A]}^2`$ $`=`$ $`0.`$ (4)
In addition, each form field satisfies the Bianchi identity
$$_{[\mu }F_{\nu _1\nu _2\mathrm{}\nu _{n_A}]}=0.$$
(5)
The notation $`(F_{[n_A]}^2)_{\mu \nu }`$ is defined as
$$(F_{[n_A]}^2)_{\mu \nu }:=F_{\mu \alpha _2\mathrm{}\alpha _{n_A}}F_\nu {}_{}{}^{\alpha _2\mathrm{}\alpha _{n_A}}.$$
(6)
The most general ansatz for the metric describing the set of intersecting non-extremal branes is
$$ds^2=\mathrm{e}^{2B}dt^2+\underset{\rho =1}{\overset{p}{}}\mathrm{e}^{2C_\rho }dx_\rho ^2+\mathrm{e}^{2A}dr^2+\mathrm{e}^{2D}d\mathrm{\Sigma }_{k,\sigma }^2,$$
(7)
where the overall transverse space is described by
$$d\mathrm{\Sigma }_{k,\sigma }^2=\overline{g}_{ab}dz^adz^b=\{\begin{array}{cc}d\psi ^2+\mathrm{sin}^2\psi d\mathrm{\Omega }_{k1}^2,\hfill & \sigma =+1,\hfill \\ d\psi ^2+\psi ^2d\mathrm{\Omega }_{k1}^2,\hfill & \sigma =0,\hfill \\ d\psi ^2+\mathrm{sinh}^2\psi d\mathrm{\Omega }_{k1}^2,\hfill & \sigma =1,\hfill \end{array}$$
(8)
satisfying
$$\overline{R}_{ab}=\sigma (k1)\overline{g}_{ab}.$$
(9)
The choice of $`\sigma `$ corresponds to different symmetries of the overall transverse space, namely, $`SO(k)`$ for $`\sigma =1`$, $`E(k)`$ for for $`\sigma =0`$, and $`SO(1,k1)`$ for $`\sigma =1`$. Two latter cases correspond to topologically non-trivial solutions which are mostly known in the (multidimensional) black hole case. In what follows we will derive the general solution valid for all $`\sigma `$, but later on we restrict to the case $`\sigma =1`$. For a discussion of the topological solutions (in the single brane case) see, e.g., Gr01 ; Gal'tsov:2004kn . The total number of dimensions occupied by the brane world-volumes is $`p`$. Each brane of the set is specified by the particular form field labeled by $`n_A`$.
The Ricci-tensor for this metric has the following non-vanishing components:
$`R_{tt}`$ $`=`$ $`\mathrm{e}^{2B2A}(B^{\prime \prime }+B^{}\mathrm{\Lambda }^{}),`$ (10)
$`R_{x_\rho x_\rho }`$ $`=`$ $`\mathrm{e}^{2C_\rho 2A}(C_\rho ^{\prime \prime }+C_\rho ^{}\mathrm{\Lambda }^{}),`$ (11)
$`R_{rr}`$ $`=`$ $`\mathrm{\Lambda }^{\prime \prime }A^{\prime \prime }+A^{}\mathrm{\Lambda }^{}+A^2B^2{\displaystyle \underset{\rho =1}{\overset{p}{}}}C_\rho ^2kD^2,`$ (12)
$`R_{ab}`$ $`=`$ $`\left[\mathrm{e}^{2D2A}(D^{\prime \prime }+D^{}\mathrm{\Lambda }^{})+\sigma (k1)\right]\overline{g}_{ab},`$ (13)
where
$$\mathrm{\Lambda }=A+B+\underset{\rho =1}{\overset{p}{}}C_\rho +kD.$$
(14)
We consider both the electrically and magnetically charges branes. In the electric case each form field is given by
$$F_{[n_A]}=f_A^{\mathrm{elec}}dtdx_{\rho _3}\mathrm{}dx_{\rho _{n_A}}dr,$$
(15)
which satisfies the Bianchi identity automatically. It supports the $`p_A(=n_A2)`$-brane in the set of intersecting branes. The equation of motion for the form fields is solved as
$$f_A^{\mathrm{elec}}=b_A\mathrm{exp}\left(A+B\underset{\rho =1}{\overset{p}{}}C_\rho kD+2\underset{\rho =1}{\overset{p}{}}C_\rho \delta _A^\rho a_A\varphi \right),$$
(16)
where $`\delta _A^\rho =1`$ for $`x^\rho `$ belonging to the world-volume of the $`p_A`$-brane.
The form fields for the magnetic branes read
$$F_{[n_A]}=b_Adx_{\rho _1}\mathrm{}dx_{\rho _{n_Ak}}\mathrm{vol}(\mathrm{\Sigma }_{k,\sigma }),$$
(17)
in which case the set $`x^{\rho _i}`$ does not belong to the world-volume of the $`p_A(=Dn_A2)`$-brane. It is easy to see that the field equations for these form fields are satisfied indeed.
To solve the Einstein equations we change the independent variable as
$$d\tau =\mathrm{e}^\mathrm{\Lambda }dr,$$
(18)
and present the radial part of the metric as
$$\mathrm{e}^{2A}dr^2=\mathrm{e}^{2𝒜}d\tau ^2,𝒜=\mathrm{\Lambda }+A.$$
(19)
The Einstein and dilaton equations to be solved are (with derivatives with respect to $`\tau `$)
$`B^{\prime \prime }`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{b_A^2(Dp_A3)}{2(D2)}}\mathrm{e}^{G_A},`$ (20)
$`C_\rho ^{\prime \prime }`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{b_A^2\delta _A^{(\rho )}}{2(D2)}}\mathrm{e}^{G_A},`$ (21)
$`D^{\prime \prime }`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{b_A^2(p_A+1)}{2(D2)}}\mathrm{e}^{G_A}+\sigma (k1)\mathrm{e}^{2𝒜2D},`$ (22)
$`\varphi ^{\prime \prime }`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{ϵ_Aa_Ab_A^2}{2}}\mathrm{e}^{G_A},`$ (23)
where
$$G_A=ϵ_Aa_A\varphi +2B+2\underset{\rho =1}{\overset{p}{}}C_\rho \delta _A^\rho ,$$
(24)
and we have defined
$`\delta _A^{(\rho )}=\{\begin{array}{c}Dp_A3\hfill \\ (p_A+1)\hfill \end{array}\mathrm{for}\{\begin{array}{c}x_\rho \text{ belonging to }p_A\text{-brane}\hfill \\ \mathrm{otherwise}\hfill \end{array}.`$ (29)
Note that this can be also written as
$`\delta _A^{(\rho )}=(D2)\delta _A^\rho (p_A+1).`$ (30)
The $`rr`$ component of Einstein equations gives
$$𝒜^{\prime \prime }𝒜^2+B^2+\underset{\rho =1}{\overset{p}{}}C_\rho ^2+kD^2=\frac{1}{2}\varphi ^2+\underset{A=1}{\overset{m}{}}\frac{b_A^2(Dp_A3)}{2(D2)}\mathrm{e}^{G_A}.$$
(31)
It is straightforward to check that the following combination separates:
$`(𝒜D)^{\prime \prime }`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{b_A^2}{2(D2)}}\left[Dp_A3+{\displaystyle \underset{\rho =1}{\overset{p}{}}}\delta _A^{(\rho )}(k1)(p_A+1)\right]\mathrm{e}^{G_A}+\sigma (k1)^2\mathrm{e}^{2(𝒜D)}`$ (32)
$`=`$ $`\sigma (k1)^2\mathrm{e}^{2(𝒜D)},`$
where we have taken into account the relation
$$Dp_A3+\underset{\rho =1}{\overset{p}{}}\delta _A^{(\rho )}=(k1)(p_A+1),$$
(33)
which follows from (29). The general solution of Eq. (32) is
$$H:=2(𝒜D)=\{\begin{array}{cc}\mathrm{ln}\left(\frac{\beta ^2}{4(k1)^2}\right)\mathrm{ln}\left[\mathrm{sinh}^2\left(\frac{\beta }{2}(\tau \tau _0)\right)\right]\hfill & \sigma =1,\hfill \\ \beta (\tau \tau _0),\hfill & \sigma =0,\hfill \\ \mathrm{ln}\left(\frac{\beta ^2}{4(k1)^2}\right)\mathrm{ln}\left[\mathrm{cosh}^2\left(\frac{\beta }{2}(\tau \tau _0)\right)\right]\hfill & \sigma =1.\hfill \end{array}$$
(34)
Differentiating Eq. (24) twice and substituting (20) – (23) into this, we arrive at the following coupled system of equations for the set of $`G_A`$:
$$G_A^{\prime \prime }=\underset{B}{}\left\{\frac{ϵ_Aa_Aϵ_Ba_B}{2}+\frac{Dp_A3}{D2}+\underset{\rho =1}{\overset{p}{}}\frac{\delta _B^{(\rho )}\delta _A^\rho }{D2}\right\}b_B^2\mathrm{e}^{G_B},$$
(35)
where the relation (30) is understood. The matrix on the right had side of this equation is non-diagonal generally. However if we impose the standard intersection rule for $`AB`$ Ohta:1997gw :
$$\underset{\rho =1}{\overset{p}{}}\delta _A^\rho \delta _B^\rho \overline{p}=\frac{(p_A+1)(p_B+1)}{D2}1\frac{ϵ_Aa_Aϵ_Ba_B}{2},$$
(36)
where $`\overline{p}`$ relates to the brane on which $`p_A`$\- and $`p_B`$-branes intersect, the non-diagonal part vanishes. In this case the system reduces to the set of separate Liouville equations, and it is integrable. Therefore, the intersection rule is a sufficient condition for integrability. It is an interesting question whether the intersection rules is a necessary condition for integrability, see the corresponding discussion in GaRy98 .
For the diagonal terms $`A=B`$, the expression in the parenthesis reduces to
$$\{\mathrm{}\}=\frac{\mathrm{\Delta }_A}{2},\mathrm{where}\mathrm{\Delta }_A=a_A^2+\frac{2(p_A+1)(Dp_A3)}{D2}.$$
(37)
With this assumption, one obtains the set of the decoupled Liouville equations for $`G_A`$:
$$G_A^{\prime \prime }=\frac{\mathrm{\Delta }_Ab_A^2}{2}\mathrm{e}^{G_A},$$
(38)
which leads to the solution
$$G_A=\mathrm{ln}\left(\frac{\alpha _A^2}{\mathrm{\Delta }_Ab_A^2}\right)\mathrm{ln}\left[\mathrm{sinh}^2\left(\frac{\alpha _A}{2}(\tau \tau _A)\right)\right].$$
(39)
Using this result, we can immediately integrate Eqs. (20-23) to obtain
$`B`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{Dp_A3}{(D2)\mathrm{\Delta }_A}}\left(G_A+g_A^{(1)}\tau +g_A^{(0)}\right)+B^{(1)}\tau +B^{(0)},`$ (40)
$`C_\rho `$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{\delta _A^{(\rho )}}{(D2)\mathrm{\Delta }_A}}\left(G_A+g_A^{(1)}\tau +g_A^{(0)}\right)+C_\rho ^{(1)}\tau +C_\rho ^{(0)},`$ (41)
$`\varphi `$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{ϵ_Aa_A}{\mathrm{\Delta }_A}}\left(G_A+g_A^{(1)}\tau +g_A^{(0)}\right)+\varphi ^{(1)}\tau +\varphi ^{(0)}.`$ (42)
It follows from Eq. (24) that the constants are connected by the following $`2m`$ relations:
$$g_A^{(0,1)}+2B^{(0,1)}+2\underset{\rho =1}{\overset{p}{}}C_\rho ^{(0,1)}\delta _A^\rho ϵ_Aa_A\varphi ^{(0,1)}=0,$$
(43)
From the solution (34) for $`𝒜D`$ and the gauge conditions (14) and (19), we can obtain the expressions for $`𝒜`$ and $`D`$:
$`𝒜`$ $`=`$ $`{\displaystyle \frac{k}{2(k1)}}H{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{p_A+1}{(D2)\mathrm{\Delta }_A}}\left(G_A+g_A^{(1)}\tau +g_A^{(0)}\right){\displaystyle \frac{B^{(1)}+C^{(1)}}{k1}}\tau {\displaystyle \frac{B^{(0)}+C^{(0)}}{k1}},`$ (44)
$`D`$ $`=`$ $`{\displaystyle \frac{1}{2(k1)}}H{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{p_A+1}{(D2)\mathrm{\Delta }_A}}\left(G_A+g_A^{(1)}\tau +g_A^{(0)}\right){\displaystyle \frac{B^{(1)}+C^{(1)}}{k1}}\tau {\displaystyle \frac{B^{(0)}+C^{(0)}}{k1}},`$ (45)
where the following constant parameters are introduced
$$C^{(1,0)}=\underset{\rho =1}{\overset{p}{}}C_\rho ^{(1,0)}.$$
(46)
Finally we have to fulfill the last equation (31). Using the intersection rules (36) and the constraints (43), this equation can be reduced to the following constraint equations on the parameters:
$$\frac{1}{2}\underset{A=1}{\overset{m}{}}\frac{\alpha _A^2(g_A^{(1)})^2}{\mathrm{\Delta }_A}+(B^{(1)})^2+\underset{\rho =1}{\overset{p}{}}(C_\rho ^{(1)})^2+\frac{1}{k1}(B^{(1)}+C^{(1)})^2+\frac{1}{2}(\varphi ^{(1)})^2\frac{k}{4(k1)}\beta ^2=0.$$
(47)
Let us count the number of free parameters. The total number of parameters appearing in the solutions is $`5m+2p+6`$. It consists of $`2`$ parameters of the function $`H`$ ($`\beta ,\tau _0`$), $`4m`$ parameters in $`G_A`$ ($`\alpha _A,\tau _A,g_A^{(0,1)}`$), $`2p+4`$ parameters entering $`B,C_\rho ,\varphi `$ ($`B^{(0,1)},C_\rho ^{(0,1)},\varphi ^{(0,1)}`$) and $`m`$ charge parameters $`b_A`$. These have to satisfy $`2m`$ constraints (43) and one constraint (47), which can fix, for example, $`g_A^{(0,1)}`$ and $`\varphi ^{(1)}`$. Thus the remaining number of independent parameters is $`3m+2p+5`$. However, not all of these $`3m+2p+5`$ parameters are physical. The coordinate transformations allow us to fix some of them: by rescaling $`t,x_\rho `$ one can absorb $`B^{(0)}`$ and $`C_\rho ^{(0)}`$. Also, since the system of equations was autonomous, we are free to shift the coordinate $`\tau `$ by a constant, so without loss of generality one can fix $`\tau _0=0`$. Therefore we are free to fix $`p+2`$ and leave $`3m+p+3`$ physical parameters, basically $`3m`$ of $`b_A,\alpha _A`$ and $`\tau _A`$ and $`p+3`$ of $`B^{(1)},C_\rho ^{(1)},\varphi ^{(0)}`$ and $`\beta `$. Apparently, we also have a freedom of rescaling $`r`$, but this is related to a change of the gauge function $`\mathrm{\Lambda }`$. So far we leave the constants in the expressions for the metric functions unfixed for later convenience.
## III Fixing the constants
From now on we will consider only the topologically simple case $`\sigma =1`$. We are interested in solutions (possibly) possessing an event horizon and not plagued with naked singularities. The simple, though incomplete, way to reveal the position of singularities and to impose the regularity condition on the horizon is to check the behavior of the Ricci scalar. Using the Einstein equations one can find the following expression for the Ricci scalar for the solution obtained:
$$R=\frac{(k1)^2}{2}\mathrm{e}^{2𝒜}\left(\varphi ^2+\underset{A=1}{\overset{m}{}}\frac{b_A^2(D2p_A4)}{(k1)^2(D2)}\mathrm{e}^{G_A}\right).$$
(48)
Using the explicit form of $`G_A`$ in (39), one can see that the points $`\tau =\tau _A`$ are singular, unless some of them are zero, in which case the singularity can be avoided by imposing further conditions on the parameters. Other special points are $`\tau \pm \mathrm{}`$ (we assume the $`\tau `$-coordinate to vary on the full real axis, and possibly to extend to the complex plane to ensure the change of signs of the exponential terms like $`\mathrm{e}^{2B}`$, see for details Ref. Gal'tsov:2004kn ). These can correspond to the horizons. At the horizons the metric coefficient $`g_{tt}`$ must vanish. Using Eq. (40) one can see this can be the case in the limits $`\tau \pm \mathrm{}`$ once suitable inequalities on the parameters are imposed. Combining this with the behavior of the Ricci scalar, we find, moreover, that one is free to choose for the regular event horizon any of these two limiting points, but then the other will be generically singular. We choose $`\tau \mathrm{}`$ as the horizon, then $`\tau \mathrm{}`$ will be generically a null singularity, except for some special choice of parameters.
Let us investigate the behavior of the metric functions at the event horizon. Assuming without loss of generality $`\beta 0,\alpha _A0`$, (we will not consider the possibility of imaginary values of these parameters which are also allowed by the overall reality of the solution), we find that, as $`\tau \mathrm{}`$, the functions $`G_A`$ and $`H`$ become linear in $`\tau `$:
$$G_A\alpha _A\tau ,H\beta \tau .$$
(49)
This ensures vanishing of the metric component $`\mathrm{e}^{2B}`$ at the horizon.
An important further information can be extracted from the constraint equation (31). It is convenient to rewrite it as follows:
$$𝒜^2+B^2+\underset{\rho =1}{\overset{p}{}}C_\rho ^2+kD^2+\frac{1}{2}\varphi ^2=\underset{A=1}{\overset{m}{}}\frac{b_A^2}{2}\mathrm{e}^{G_A}\sigma k(k1)\mathrm{e}^H.$$
(50)
From the radial geodesic equation, it follows that $`𝒜^{}=B^{}`$ at the horizon Gal'tsov:2004kn , so the first two terms on the left hand side cancel. The right hand side vanishes at the horizon, so one is left with the sum of positive definite terms. Therefore, all the derivatives $`C_\rho ^{},D^{},\varphi ^{}`$ must separately vanish in the limit $`\tau \mathrm{}`$. More precisely, the condition $`𝒜^{}=B^{}`$ gives
$$\frac{k}{2(k1)}\beta \underset{A=1}{\overset{m}{}}\frac{\alpha _A+g_A^{(1)}}{\mathrm{\Delta }_A}\frac{B^{(1)}+C^{(1)}}{k1}B^{(1)}=0.$$
(51)
Moreover, using Eqs. (41), (42) and (45), we find the following relations on the parameters involved, namely, from Eq. (41) we obtain $`p`$ relations
$$C_\rho ^{(1)}=\underset{A=1}{\overset{m}{}}\frac{\delta _A^{(\rho )}}{D2}\frac{\alpha _A+g_A^{(1)}}{\mathrm{\Delta }_A},$$
(52)
from Eq. (42)
$$\varphi ^{(1)}=\underset{A=1}{\overset{m}{}}ϵ_Aa_A\frac{\alpha _A+g_A^{(1)}}{\mathrm{\Delta }_A},$$
(53)
and from Eq. (45)
$$\frac{1}{2(k1)}\beta \underset{A=1}{\overset{m}{}}\frac{p_A+1}{D2}\frac{\alpha _A+g_A^{(1)}}{\mathrm{\Delta }_A}\frac{B^{(1)}+C^{(1)}}{k1}=0.$$
(54)
Consequently
$$C^{(1)}=\underset{A=1}{\overset{m}{}}\frac{D2k(p_A+1)}{D2}\frac{\alpha _A+g_A^{(1)}}{\mathrm{\Delta }_A}.$$
(55)
The above constraints (51)-(54) are consistent with (47). Combining (51) and (54), one can obtain
$$B^{(1)}=\frac{1}{2}\beta \underset{B=1}{\overset{m}{}}\frac{Dp_B3}{D2}\frac{\alpha _B+g_B^{(1)}}{\mathrm{\Delta }_B}.$$
(56)
Furthermore, the constraints (43), after substituting (52) and (53), become
$$B^{(1)}=\frac{1}{2}\alpha _A\underset{B=1}{\overset{m}{}}\frac{Dp_B3}{D2}\frac{\alpha _B+g_B^{(1)}}{\mathrm{\Delta }_B}.$$
(57)
Hence, all the parameters $`\alpha _A`$ must be equal:
$$\alpha _A=\beta .$$
(58)
Then the metric functions read
$`B`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{Dp_A3}{(D2)\mathrm{\Delta }_A}}\left(G_A\beta \tau +g_A^{(0)}\right)+{\displaystyle \frac{1}{2}}\beta \tau +B^{(0)},`$ (59)
$`C_\rho `$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{\delta _A^{(\rho )}}{(D2)\mathrm{\Delta }_A}}\left(G_A\beta \tau +g_A^{(0)}\right)+C_\rho ^{(0)},`$ (60)
$`\varphi `$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{ϵ_Aa_A}{\mathrm{\Delta }_A}}\left(G_A\beta \tau +g_A^{(0)}\right)+\varphi ^{(0)},`$ (61)
$`𝒜`$ $`=`$ $`{\displaystyle \frac{k}{2(k1)}}H{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{p_A+1}{(D2)\mathrm{\Delta }_A}}\left(G_A\beta \tau +g_A^{(0)}\right){\displaystyle \frac{\beta }{2(k1)}}\tau {\displaystyle \frac{B^{(0)}+C^{(0)}}{k1}},`$ (62)
$`D`$ $`=`$ $`{\displaystyle \frac{1}{2(k1)}}H{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{p_A+1}{(D2)\mathrm{\Delta }_A}}\left(G_A\beta \tau +g_A^{(0)}\right){\displaystyle \frac{\beta }{2(k1)}}\tau {\displaystyle \frac{B^{(0)}+C^{(0)}}{k1}}.`$ (63)
Though the parameters $`\alpha _A`$ are the same for all $`A`$, the functions $`G_A`$ differ in position of singularities $`\tau _A`$
### III.1 Asymptotically flat solutions
The asymptotic region is located at $`\tau \tau _0`$, and we have fixed the translational freedom by choosing $`\tau _0=0`$. The corresponding behavior of the function $`H`$ is
$$H\mathrm{ln}\frac{1}{\tau ^2}.$$
(64)
However, there are two different cases for $`G_A`$ depending on the value of $`\tau _A`$. For the case $`\tau _A0`$, the asymptotic value of $`G_A`$, for $`\alpha _A=\beta `$, is
$$G_AG_A^{(0)}:=\mathrm{ln}\left[\frac{\beta ^2}{b_A^2\mathrm{\Delta }_A\mathrm{sinh}^2(\frac{\beta }{2}\tau _A)}\right].$$
(65)
Therefore, asymptotically the solutions reduces to
$`B`$ $``$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{Dp_A3}{(D2)\mathrm{\Delta }_A}}\left(G_A^{(0)}+g_A^{(0)}\right)+B^{(0)},`$ (66)
$`C_\rho `$ $``$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{\delta _A^{(\rho )}}{(D2)\mathrm{\Delta }_A}}\left(G_A^{(0)}+g_A^{(0)}\right)+C_\rho ^{(0)},`$ (67)
$`\varphi `$ $``$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{ϵ_Aa_A}{\mathrm{\Delta }_A}}\left(G_A^{(0)}+g_A^{(0)}\right)+\varphi ^{(0)}.`$ (68)
The quantities $`B`$ and $`C_\rho `$ are constants which can be set to zero by rescaling of time and world-volume coordinates imposing the conditions
$`B^{(0)}`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{Dp_A3}{(D2)\mathrm{\Delta }_A}}\left(G_A^{(0)}+g_A^{(0)}\right),`$ (69)
$`C_\rho ^{(0)}`$ $`=`$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{\delta _A^{(\rho )}}{(D2)\mathrm{\Delta }_A}}\left(G_A^{(0)}+g_A^{(0)}\right).`$ (70)
For the dilaton, it is common to preserve the finite value $`\varphi _{\mathrm{}}`$ at infinity, and this can be ensured by the redefinition
$$\varphi ^{(0)}=\underset{A=1}{\overset{m}{}}\frac{ϵ_Aa_A}{\mathrm{\Delta }_A}\left(G_A^{(0)}+g_A^{(0)}\right)+\varphi ^{\mathrm{}}.$$
(71)
Then the relation (43) gives
$$\varphi ^{\mathrm{}}=\frac{G_A^{(0)}}{ϵ_Aa_A}.$$
(72)
This equation requires the quantities $`G_A^{(0)}/ϵ_Aa_A`$ to be identical for all $`A`$, which imposes a new set of relations between $`b_A`$ and $`\tau _A`$. The Ricci scalar diverges at $`\tau _A`$, so, to avoid naked singularities, one has to ensure that all $`\tau _A`$ lie inside the event horizon for black branes. If there is also an inner horizon, one can check that the real metrics correspond to location of singularities $`\tau _A`$ either outside the event horizon, or inside the inner horizon. The first possibility should be ruled out. The notable exception constitutes the case of $`\tau _A=0`$. In this case there is no singularity, and this point can be interpreted as spatial infinity. The physical nature of this will be clearer later when we express the solution in the Schwarzschild-type coordinate.
### III.2 Asymptotically LDB branes
For the case $`\tau _A=0`$, each $`G_A`$ has the same asymptotic behavior as $`H`$:
$$G_A\mathrm{ln}\frac{1}{\tau ^2}.$$
(73)
In such case, we have
$`B`$ $``$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{Dp_A3}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}{\displaystyle \frac{1}{\tau ^2}},`$ (74)
$`C_\rho `$ $``$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{\delta _A^{(\rho )}}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}{\displaystyle \frac{1}{\tau ^2}},`$ (75)
$`\varphi `$ $``$ $`{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{ϵ_Aa_A}{\mathrm{\Delta }_A}}\mathrm{ln}{\displaystyle \frac{1}{\tau ^2}}.`$ (76)
This solution is not asymptotically flat, but asymptotically approaching the LDB.
## IV Schwarzschild-type coordinates
To present our solution in a more familiar form, we change the coordinates similarly to Gal'tsov:2004kn . We map the horizon $`\tau =\mathrm{}`$ to $`r=r_H`$,
$$r_H=\mu ^{\frac{1}{k1}},\mu =\frac{\beta }{k1}.$$
(77)
and the internal horizon $`\tau =\mathrm{}`$ to $`r=0`$ by choosing the gauge function
$$\mathrm{\Lambda }=\mathrm{ln}(r^kf),\text{such that}d\tau =\frac{dr}{r^kf},$$
(78)
where
$$f=1\frac{\mu }{r^{k1}}.$$
(79)
This corresponds to the coordinate transformation
$$\tau =\frac{1}{(k1)\mu }\mathrm{ln}f.$$
(80)
Thus the region outside the event horizon $`r>r_H`$ corresponds to the half axis $`(\mathrm{},\mathrm{\hspace{0.33em}0})`$ of $`\tau `$. Then we will have
$$H=\mathrm{ln}\left(r^{2(k1)}f\right).$$
(81)
### IV.1 Asymptotically flat solutions
Introduce a new parameter $`q_A`$ instead of $`\tau _A`$ by
$$q_A=\frac{\mu }{\mathrm{e}^{\mu (k1)\tau _A}1}.$$
(82)
As we already noted, to ensure the absence of naked singularities, all $`\tau _A`$ should be taken non-negative. Here we assume that $`\tau _A`$ are strictly positive, so this definition leads to finite $`q_A`$ (the case $`\tau _A=0`$ will be considered below). The function $`G_A`$ then reads
$$G_A=\mathrm{ln}\left(\frac{4(k1)^2(\mu +q_A)q_A}{\mathrm{\Delta }_Ab_A^2}\frac{f}{h_A^2}\right),$$
(83)
where
$$h_A=1+\frac{q_A}{r^{k1}},$$
(84)
are the harmonic functions.
In terms of the new coordinates, the metric components (7) read
$`B`$ $`=`$ $`2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{Dp_A3}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}h_A+{\displaystyle \frac{1}{2}}\mathrm{ln}f,`$ (85)
$`C_\rho `$ $`=`$ $`2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{\delta _A^{(\rho )}}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}h_A,`$ (86)
$`\varphi `$ $`=`$ $`2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{ϵ_Aa_A}{\mathrm{\Delta }_A}}\mathrm{ln}h_A+\varphi ^{\mathrm{}},`$ (87)
$`A`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}f+2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{p_A+1}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}h_A,`$ (88)
$`D`$ $`=`$ $`\mathrm{ln}r+2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{p_A+1}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}h_A,`$ (89)
with a constraint following from (72):
$$b_A^2=\mathrm{e}^{ϵ_Aa_A\varphi ^{\mathrm{}}}\frac{4(k1)^2(\mu +q_A)q_A}{\mathrm{\Delta }_A}.$$
(90)
For the electric branes, the solution for the form field is
$$f_A^{\mathrm{elec}}=b_A\mathrm{e}^{a_A\varphi ^{\mathrm{}}}\frac{1}{r^kh_A^2}.$$
(91)
This solution corresponds to the one given in Ohta:1997gw . Free parameters of the solution are $`\mu ,q_A`$ and $`\varphi _{\mathrm{}}`$.
### IV.2 Asymptotically LDB branes
If $`\tau _A=0`$, $`G_A`$ takes the form
$$G_A=\mathrm{ln}\left(\frac{4(k1)^2}{\mathrm{\Delta }_Ab_A^2}r^{2(k1)}f\right).$$
(92)
Then the solution is, with the choice of parameters $`B^{(0)}=0`$ and $`C_\rho ^{(0)}=0`$,
$`B`$ $`=`$ $`2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{Dp_A3}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}{\displaystyle \frac{r^{k1}}{q_A}}+{\displaystyle \frac{1}{2}}\mathrm{ln}f,`$ (93)
$`C_\rho `$ $`=`$ $`2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{\delta _A^{(\rho )}}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}{\displaystyle \frac{r^{k1}}{q_A}},`$ (94)
$`\varphi `$ $`=`$ $`2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{ϵ_Aa_A}{\mathrm{\Delta }_A}}\mathrm{ln}{\displaystyle \frac{r^{k1}}{q_A}}+\varphi ^{\mathrm{}},`$ (95)
$`A`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}f2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{p_A+1}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}{\displaystyle \frac{r^{k1}}{q_A}},`$ (96)
$`D`$ $`=`$ $`\mathrm{ln}r2{\displaystyle \underset{A=1}{\overset{m}{}}}{\displaystyle \frac{p_A+1}{(D2)\mathrm{\Delta }_A}}\mathrm{ln}{\displaystyle \frac{r^{k1}}{q_A}},`$ (97)
where the parameters $`q_A`$ are defined as
$$q_A^2=\frac{4(k1)^2}{\mathrm{\Delta }_Ab_A^2}\mathrm{e}^{g_A^{(0)}}.$$
(98)
Our previous form strength parameter $`b_A`$ is related to $`q_A`$ via
$$b_A^2=\mathrm{e}^{ϵ_Aa_A\varphi ^{\mathrm{}}}\frac{4(k1)^2q_A^2}{\mathrm{\Delta }_A}.$$
(99)
The solution for the electric form field is
$$f_A^{\mathrm{elec}}=b_A\mathrm{e}^{a_A\varphi ^{\mathrm{}}}\frac{r^{k2}}{q_A^2}.$$
(100)
This second possibility for $`\tau _A`$ leads to an intersecting generalization of the solution of Ref. Clement:2004ii . The BPS limit of these solutions corresponds to $`\mu =0`$. It can be recognized that in this case one deals with a solution known earlier as the hear-horizon limit of the BPS intersecting branes. This corresponds to the usual rule of omitting constants in the harmonic functions describing the metric. Here we reproduce these metrics as a particular case of the general supergravity solution. Our result is, however, not just to reproduce them by another technique, but to prove that no other solutions without naked singularities exist within the class of metrics considered, which is fixed by their isometries, up to the the mixed intersections below.
The black versions of these solutions with $`\mu 0`$ should describe the thermal phase of QFT’s in the Domain-Wall/QFT correspondence associated to their BPS limit.
### IV.3 Mixed intersections
The mixed intersecting configurations, namely the black AF $`p`$-branes with the asymptotically LDB ones, can be obtained simply by picking up the corresponding term in the summation of $`A,B,C_\rho ,D`$ and $`\varphi `$ depending on whether $`p_A`$ is AF brane or LDB. Again, the BPS limit of these solutions was found previously through the near-horizon considerations. It was noted that one can drop the constant term in the harmonic functions describing the system of intersecting branes not in all harmonic functions at the same time, but in some of them Bo97 . This corresponds to the BPS limit of our mixed intersection. To our knowledge, the black version of these solutions is new.
## V Two intersecting branes
We give the examples of two intersecting branes. According to the intersection rules, we can have solutions listed below:
IIA $`:`$ $`22(0),44(2),66(4),04(0),24(1),26(2),46(3),48(4),68(5),`$
IIB $`:`$ $`33(1),55(3),77(5),13(0),15(1),35(2),37(3),57(4),59(5),`$
M $`:`$ $`\mathrm{M2}\mathrm{M2}(0),\mathrm{M2}\mathrm{M5}(1),\mathrm{M5}\mathrm{M5}(3),`$
NS $`:`$ $`\mathrm{NS1}\mathrm{NS5}(1),\mathrm{NS5}\mathrm{NS5}(3),\mathrm{D}p\mathrm{NS1}(0),\mathrm{\hspace{0.17em}0}p8,\mathrm{D}p\mathrm{NS5}(p1),\mathrm{\hspace{0.17em}1}p6,`$ (101)
where those branes indicated only by numbers are D-branes, and the number in the parenthesis denotes the dimensions of the overlapping coordinates. We give explicit solutions for M and IIB theories in the following.
### V.1 M2-M2 configurations
Let us consider the general solution for two M2-branes as an explicit example. In this case, the parameters have the values: $`D=11,k=5,n_A=4,a_A=0,p_A+1=3,Dp_A3=6`$ and $`\mathrm{\Delta }_A=4`$. The AF solutions have the metric
$$ds^2=h_1^{\frac{1}{3}}h_2^{\frac{1}{3}}\left[fh_1^1h_2^1dt^2+h_1^1(dx_1^2+dx_2^2)+h_2^1(dx_3^2+dx_4^2)+f^1dr^2+r^2d\mathrm{\Omega }_5^2\right],$$
(102)
where $`h_A=1+q_A/r^4`$, and the form fields are
$$f_A^{\mathrm{elec}}=\frac{b_A}{r^5h_A^2},b_A^2=16(\mu +q_A)q_A.$$
(103)
We recover the canonical intersecting two M2-branes.
For asymptotically LDB case, the solution we have derived corresponds to omitting the constant in the harmonic functions $`h_1`$ and $`h_2`$. It is
$$ds^2=\left(\frac{q_1}{r^4}\right)^{\frac{1}{3}}\left(\frac{q_2}{r^4}\right)^{\frac{1}{3}}\left[f\left(\frac{q_1}{r^4}\right)^1\left(\frac{q_2}{r^4}\right)^1dt^2+\left(\frac{q_1}{r^4}\right)^1(dx_1^2+dx_2^2)+\left(\frac{q_2}{r^4}\right)^1(dx_3^2+dx_4^2)+f^1dr^2+r^2d\mathrm{\Omega }_5^2\right],$$
(104)
and
$$f_A^{\mathrm{elec}}=\frac{b_Ar^3}{q_A^2},b_A^2=16q_A^2.$$
(105)
The mixed intersecting configurations of two M2-branes is given by
$$ds^2=h_1^{\frac{1}{3}}\left(\frac{q_2}{r^4}\right)^{\frac{1}{3}}\left[fh_1^1\left(\frac{q_2}{r^4}\right)^1dt^2+h_1^1(dx_1^2+dx_2^2)+\left(\frac{q_2}{r^4}\right)^1(dx_3^2+dx_4^2)+f^1dr^2+r^2d\mathrm{\Omega }_5^2\right],$$
(106)
where
$$f_1^{\mathrm{elec}}=\frac{b_1}{r^5h_1^2},b_1^2=16(\mu +q_1)q_1,f_2^{\mathrm{elec}}=\frac{b_2r^3}{q_2^2},b_2^2=16q_2^2.$$
(107)
### V.2 D1-D5 configurations
The other interesting case is the D1-D5 system. In this case, the parameters have the following values: $`D=10,k=3,n_A=3,a_A=1,ϵ_1=1,ϵ_2=1,p_1+1=2=Dp_23,p_2+1=6=Dp_13`$ and $`\mathrm{\Delta }_A=4`$. The AF solutions have the metric
$$ds^2=h_1^{\frac{1}{4}}h_2^{\frac{3}{4}}\left[h_1^1h_2^1(fdt^2+dx_1^2)+h_2^1(dx_2^2+\mathrm{}+dx_5^2)+f^1dr^2+r^2d\mathrm{\Omega }_3^2\right],$$
(108)
and the dilaton and the form fields
$$\mathrm{e}^{2\varphi }=\mathrm{e}^{2\varphi ^{\mathrm{}}}\frac{h_1}{h_2},f_1^{\mathrm{elec}}=\mathrm{e}^\varphi ^{\mathrm{}}\frac{b_1}{r^3h_1^2},b_A^2=\mathrm{e}^{ϵ_A\varphi ^{\mathrm{}}}4(\mu +q_A)q_A,$$
(109)
where $`h_A=1+q_A/r^2`$. We recover the canonical intersecting D1-D5-branes in the Einstein frame.
Again, for LDB, our solution corresponds to omitting the constant in the harmonic functions $`h_1`$ and $`h_2`$. It is
$$ds^2=\left(\frac{q_1}{r^2}\right)^{\frac{1}{4}}\left(\frac{q_2}{r^2}\right)^{\frac{3}{4}}\left[\left(\frac{q_1}{r^2}\right)^1\left(\frac{q_2}{r^2}\right)^1(fdt^2+dx_1^2)+\left(\frac{q_2}{r^2}\right)^1(dx_2^2+\mathrm{}dx_5^2)+f^1dr^2+r^2d\mathrm{\Omega }_3^2\right],$$
(110)
and
$$\mathrm{e}^{2\varphi }=\mathrm{e}^{2\varphi ^{\mathrm{}}}\frac{q_1}{q_2},f_1^{\mathrm{elec}}=\mathrm{e}^\varphi ^{\mathrm{}}\frac{b_1r}{q_1^2},b_A^2=\mathrm{e}^{ϵ_A\varphi ^{\mathrm{}}}4q_A^2.$$
(111)
There are also two mixed intersecting configurations of D1-D5-branes. One is given by
$$ds^2=h_1^{\frac{1}{4}}\left(\frac{q_2}{r^2}\right)^{\frac{3}{4}}\left[h_1^1\left(\frac{q_2}{r^2}\right)^1(fdt^2+dx_1^2)+\left(\frac{q_2}{r^2}\right)^1(dx_2^2+\mathrm{}+dx_5^2)+f^1dr^2+r^2d\mathrm{\Omega }_3^2\right],$$
(112)
where
$$\mathrm{e}^{2\varphi }=\mathrm{e}^{2\varphi ^{\mathrm{}}}\frac{r^2h_1}{q_2},f_1^{\mathrm{elec}}=\mathrm{e}^\varphi ^{\mathrm{}}\frac{b_1}{r^3h_1^2},b_1^2=\mathrm{e}^\varphi ^{\mathrm{}}4(\mu +q_1)q_1,b_2^2=\mathrm{e}^\varphi ^{\mathrm{}}4q_2^2,$$
(113)
and the other is
$$ds^2=\left(\frac{q_1}{r^2}\right)^{\frac{1}{4}}h_2^{\frac{3}{4}}\left[\left(\frac{q_1}{r^2}\right)^1h_2^1(fdt^2+dx_1^2)+h_2^1(dx_2^2+\mathrm{}+dx_5^2)+f^1dr^2+r^2d\mathrm{\Omega }_3^2\right],$$
(114)
where
$$\mathrm{e}^{2\varphi }=\mathrm{e}^{2\varphi ^{\mathrm{}}}\frac{q_1}{r^2h_2},f_1^{\mathrm{elec}}=\mathrm{e}^\varphi ^{\mathrm{}}\frac{b_1r}{q_1^2},b_1^2=\mathrm{e}^\varphi ^{\mathrm{}}4q_1^2,b_2^2=\mathrm{e}^\varphi ^{\mathrm{}}4(\mu +q_2)q_2.$$
(115)
## VI Conclusions
In this paper we obtained the general solution for the non-BPS intersecting supergravity $`p`$-branes delocalized in relative transverse dimensions. This was done by fully integrating the Einstein equations with a suitable ansatz for the metric, the antisymmetric forms and the dilaton. Our solution differs from earlier ones by a larger number of free parameters, which is maximal in the present case. This, in particular, opens a way to construct individual brane solutions which are either asymptotically flat, or asymptotically approaching the LDB. For intersecting branes, imposing the conditions on parameters which ensure the absence of naked singularities, we found that the resulting configuration is an intersection of some AF branes with some asymptotically LDB branes. The intersection rules remain the same as previously known. In the case when all branes are AF, we recover the solution of Ohta:1997gw . When all individual branes are asymptotically LDB, we obtain the solution which corresponds to omitting constant terms in the harmonic functions involved (the overall solution being black). Apart from these two extremes, there are intersections of a number of AF branes with a number of LDB branes within the intersection rules.
The linear dilaton backgrounds were extensively used in the context of DW/QFT (NS/LST) dualities Boonstra:1998mp ; Behrndt:1999mk ; Aharony:1998ub ; Aharony:1999ks . The LDB solutions endowed with the event horizon describe the thermal phase of the QFT (LST) associated with the linear dilaton background HaOb ; BeRo ; BoRo . Although this was not our intention here, our new intersecting solutions with LDB asymptotics may serve a basis for further studies along these directions.
The BPS versions of our solutions involving intersections with LDB branes correspond to partial omission of the constant terms in the harmonic functions involved. It has to be emphasized that our results are based on the complete integration of the Einstein equations for metric of chosen symmetry, so there are no other other solutions within this class which are free from naked singularities. Various suggestions for $`p`$-brane solutions with extra parameters which obey the same metric ansatz, but do not reduce to the standard BPS or black branes, generically suffer from naked singularities.
###### Acknowledgements.
CMC would like to thank Koryu Kyokai grant for supporting his visit to Japan and the hospitality of Yukawa Institute of Theoretical Physics and Osaka University during this work was proceeding. The work of CMC was supported by the NSC grant 93-2112-M-008-021 and 94-2119-M-002-001. The work of NO was supported in part by the Grant-in-Aid for Scientific Research Fund of the JSPS No. 16540250, and that of DG was supported in part by the RFBR grant 02-04-16949.
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# Exact renormalization group approach to a nonlinear diffusion equation
## Abstract
The exact renormalization group is applied to a nonlinear diffusion equation with a discontinuous diffusion coefficient. The generating functional of the solution for the initial-value problem of nonlinear diffusion equations is first introduced, and next a new regularization scheme is presented. It is shown that the renormalization of an action functional in the generating functional leads to an anomalous diffusion exponent in full order of the perturbation series with respect to a nonlinearity.
The renormalization group (RG) is a powerful tool to reveal universal behavior of various systems, including quantum field theories and statistical mechanics WilKog74 . Its basic idea lies in the coarse-graining of short-distance degrees of freedom, which causes redefinition of parameters governing the long-distance physics of the systems under investigation. In spite of its conceptual simplicity, there exist in the RG methods many techniques to solve with, and numerous applications have been made to equilibrium or near-equilibrium systems RGtextbook ; Car96 . In particular, exact RG (ERG) (also called nonperturbative RG or functional RG) techniques and some approximations based on them WilKog74 ; WegHou73 ; Pol84 ; Wet93 ; Mor94 have been attracting much renewed interest to reveal nonperturbative phenomena in field theories ERGrevFie , statistical mechanics ERGsta , and condensed matter physics ZanSch ; HalMet ; HonSal .
On the other hand, Goldenfeld et al. GMOL90 ; CGO have extended the RG methods to systems far from equilibrium. They have demonstrated that there exists a deep relationship between the RG and the intermediate asymptotics method Bar96 in the study of the nonlinear partial-differential equations for nonequilibrium systems. Their idea has attracted much interest, and the RG approach to nonlinear differential equations has been developed RGforNLDE .
In this paper, we apply the ERG method to a nonlinear diffusion equation called Barenblatt equation Bar96 . This equation has a discontinuous diffusion coefficient; this discontinuity makes perturbative expansion more complicated if one proceeds to higher order computations. We show in this paper that Polchinski equation, a version of the ERG equation, is very efficient even for such a nonlinear diffusion equation. It turns out that we can indeed solve the equation for all order in the perturbation series. The solution leads us to the full anomalous diffusion exponent.
Let us start with the following nonlinear diffusion equation called Barenblatt equation;
$`_tu(x,t)D(u)_x^2u(x,t)=0`$ (1)
with an initial condition $`u(x,t=0)=q\delta (x)`$, where $`D(u)\kappa [1+g\theta (\alpha _tu)]`$ denotes a nonlinear diffusion coefficient with $`\alpha `$ being a positive constant which makes $`\alpha _tu`$ dimensionless (below, we set $`\alpha =1`$, for simplicity). Here, $`\theta (x)=0(1)`$ for $`x<0(x>0)`$ stands for the step function. The dimensionless constant $`g`$ controls the nonlinearity of the diffusion coefficient. The Barenblatt equation describes the filtration of a compressible fluid through a compressible porous medium which can be irreversibly deformed. Goldenfeld et al. GMOL90 have obtained asymptotic behavior of the solution by solving this equation via an iteration scheme corresponding to the perturbative expansion with respect to $`g`$. This perturbation gives rise to divergences: Their basic idea is introducing a renormalization scheme which render the solution finite and deriving an anomalous diffusion exponent as an anomalous dimension in the RG language. Though they have successfully obtained the leading correction of the diffusion exponent, their method seems difficult to extend to higher order due to the discontinuous step-function nonlinearity. We will present a new renormalization scheme for the initial-value problems of nonlinear diffusion equations.
To be specific, we introduce a generating functional of the solution for the Barenblatt equation and a modified propagator with a short-time cutoff to render the solution finite. To this end, notice Car96 ; MSR73 that the solution of Eq. (1) can be written as $`u(x,t)=𝒟\varphi \varphi (x,t)\delta (_t\varphi D(\varphi )_x^2\varphi )\delta (\varphi (x,0)u(x,0))`$. This expression can be rewritten by a functional integral if the derivative with respect to $`t`$ is interpreted as a forward difference operator. Namely, using the Fourier transformation for the delta function, we reach $`u(x,t)=\varphi (x,t)=\frac{1}{𝒵}𝒟\varphi 𝒟\stackrel{~}{\varphi }\varphi (x,t)e^{iS}`$ with $`S`$ being an action functional
$`S={\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x\left[\stackrel{~}{\varphi }\mathrm{\Delta }^1\varphi g\stackrel{~}{\varphi }\theta (_t\varphi )\kappa _x^2\varphi \stackrel{~}{\varphi }J\right],`$ (2)
where $`\mathrm{\Delta }(x,t)e^{x^2/(4\kappa t)}/\sqrt{4\pi \kappa t}`$ denotes the diffusion propagator and the generating functional $`𝒵`$ is defined by $`𝒵𝒟\varphi 𝒟\stackrel{~}{\varphi }e^{iS}`$, as usual. The field $`J`$ in the last term is defined by $`J(x,t)u(x,0)\delta (t)`$ which controls the initial value of $`u`$. In what follows, we examine the case $`u(x,0)=q\delta (x)`$, as mentioned below Eq. (1), but it should be stressed that the generic initial-value problem can be treated similarly.
As discussed by Goldenfeld et al. GMOL90 , perturbative calculation diverges with the initial condition specified above. To regularize the solution, they have introduced an initial distribution with a finite width such as $`u(x,t=0)=e^{x^2/(2l^2)}/\sqrt{2\pi l^2}`$. We instead introduce the short-time cutoff $`\epsilon `$ for the propagator to formulate the ERG for the present nonlinear diffusion equation. To be specific, we define a modified propagator as
$`\mathrm{\Delta }_\epsilon (x,t)=\theta (t\epsilon )\mathrm{\Delta }(x,t).`$ (3)
One can easily check that this propagator indeed gives a finite solution in the perturbation theory, applying it to the calculations by Goldenfeld et al. GMOL90 . This regularization scheme can be used not only in the Barenblatt equation but also in generic diffusion problems.
Having defined the generating functional and the modified propagator, we next derive an ERG equation for the action functional Eq. (2). Using the propagator (3) with a cutoff $`\epsilon _0`$ and introducing a source term, we start with the generating functional,
$`𝒵[\stackrel{~}{J},J]={\displaystyle 𝒟\varphi 𝒟\stackrel{~}{\varphi }e^{i\stackrel{~}{\varphi }\mathrm{\Delta }_{\epsilon _0}^1\varphi iS_{\epsilon _0}[\stackrel{~}{\varphi },\varphi ]i\stackrel{~}{J}\varphi i\stackrel{~}{\varphi }J}},`$ (4)
where the bare action of the nonlinear term is $`S_{\epsilon _0}[\stackrel{~}{\varphi },\varphi ]=g\stackrel{~}{\varphi }\theta (_t\varphi )\kappa _x^2\varphi `$. At the end of the calculations, we must set $`J(x,t)=q\delta (x)\delta (t)`$ to obtain the solution for the initial-value problem of the present equation. Here, the symbol $`ab`$ implies $`ab=_0^{\mathrm{}}𝑑t_{\mathrm{}}^{\mathrm{}}𝑑xa(x,t)b(x,t)`$. Next, we introduce a new cutoff $`\epsilon (>\epsilon _0)`$ and divide the propagator into two parts $`\mathrm{\Delta }_{\epsilon _0}=\mathrm{\Delta }_>+\mathrm{\Delta }_<`$, where
$`\mathrm{\Delta }_>=\left[\theta (t\epsilon _0)\theta (t\epsilon )\right]\mathrm{\Delta },`$
$`\mathrm{\Delta }_<=\theta (t\epsilon )\mathrm{\Delta }.`$ (5)
Here, $`>`$ and $`<`$ imply the short-time and long-time modes, respectively. Separating also $`\varphi `$ and $`\stackrel{~}{\varphi }`$ into two fields $`\varphi =\varphi _>+\varphi _<`$ and $`\stackrel{~}{\varphi }=\stackrel{~}{\varphi }_>+\stackrel{~}{\varphi }_<`$ enables us to rewrite the generating functional as
$`𝒵[\stackrel{~}{J},J]{\displaystyle 𝒟\varphi _<𝒟\stackrel{~}{\varphi }_<e^{i\stackrel{~}{\varphi }_<\mathrm{\Delta }_<^1\varphi _<}𝒵_\epsilon [\stackrel{~}{J},J;\stackrel{~}{\varphi }_<,\varphi _<]},`$
$`𝒵_\epsilon [\stackrel{~}{J},J;\stackrel{~}{\varphi }_<,\varphi _<]={\displaystyle 𝒟\varphi _>𝒟\stackrel{~}{\varphi }_>e^{i\stackrel{~}{\varphi }_>\mathrm{\Delta }_>^1\varphi _>}}`$
$`\times e^{iS_{\epsilon _0}[\stackrel{~}{\varphi }_>+\stackrel{~}{\varphi }_<,\varphi _>+\varphi _<]i\stackrel{~}{J}(\varphi _>+\varphi _<)i(\stackrel{~}{\varphi }_>+\stackrel{~}{\varphi }_<)J},`$ (6)
up to a proportionality constant. The field $`\varphi _>`$, $`\stackrel{~}{\varphi }_>`$ and $`\varphi _<`$, $`\stackrel{~}{\varphi }_<`$ can be identified as fields describing short-time and long-time modes, respectively. Integrating out the short-time fields, we will next derive an effective action describing long-time modes.
Changing the integration variables $`\stackrel{~}{\varphi }_>`$ and $`\varphi _>`$ into $`\stackrel{~}{\varphi }_>=\stackrel{~}{\varphi }\stackrel{~}{\varphi }_<`$ and $`\varphi _>=\varphi \varphi _<`$, and integrating over the fields $`\stackrel{~}{\varphi }`$ and $`\varphi `$ in Eq. (6) yields
$`𝒵_\epsilon [\stackrel{~}{J},J;\stackrel{~}{\varphi }_<,\varphi _<]`$ (7)
$`={\displaystyle 𝒟\stackrel{~}{\varphi }𝒟\varphi e^{i(\stackrel{~}{\varphi }\stackrel{~}{\varphi }_<)\mathrm{\Delta }_>^1(\varphi \varphi _<)iS_{\epsilon _0}[\stackrel{~}{\varphi },\varphi ]i\stackrel{~}{J}\varphi i\stackrel{~}{\varphi }J}}`$
$`=e^{i\stackrel{~}{J}\mathrm{\Delta }_>Ji\stackrel{~}{\varphi }_<Ji\stackrel{~}{J}\varphi _<iS_\epsilon [\stackrel{~}{J}\mathrm{\Delta }_>+\stackrel{~}{\varphi }_<,\mathrm{\Delta }_>J+\varphi _<]},`$
where $`S_\epsilon `$ is defined by
$`e^{iS_\epsilon [\stackrel{~}{J}\mathrm{\Delta }_>+\stackrel{~}{\varphi }_<,\mathrm{\Delta }_>J+\varphi _<]}`$ (8)
$`e^{iS_{\epsilon _0}[i\delta /\delta J+\stackrel{~}{J}\mathrm{\Delta }_>+\stackrel{~}{\varphi }_<,i\delta /\delta \stackrel{~}{J}+\mathrm{\Delta }_>J+\varphi _<]}.`$
This equation implies that if we expand the exponential in r.h.s and make all the derivatives $`\delta /\delta J`$ and $`\delta /\delta \stackrel{~}{J}`$ in $`S_{\epsilon _0}`$ act on $`J`$ and $`\stackrel{~}{J}`$ in the right $`S_{\epsilon _0}`$, we reach some $`S_\epsilon `$ as a functional of $`\stackrel{~}{\mathrm{\Phi }}=\stackrel{~}{J}\mathrm{\Delta }_>+\stackrel{~}{\varphi }_<`$ and $`\mathrm{\Phi }=\mathrm{\Delta }_>J+\varphi _<`$. Instead of carrying out such calculations, however, we can alternatively determine the functional $`S_\epsilon `$ by noting that $`𝒵_\epsilon `$ obeys
$`{\displaystyle \frac{d𝒵_\epsilon }{d\epsilon }}=i\left(i{\displaystyle \frac{\delta }{\delta J}}\stackrel{~}{\varphi }_<\right){\displaystyle \frac{d\mathrm{\Delta }_>^1}{d\epsilon }}\left(i{\displaystyle \frac{\delta }{\delta \stackrel{~}{J}}}\varphi _<\right)𝒵_\epsilon ,`$ (9)
which follows from Eq. (6). Substituting Eq. (7) into Eq. (9), we obtain the following Polchinski RG equation,
$`{\displaystyle \frac{S_\epsilon }{\epsilon }}={\displaystyle \frac{\delta S_\epsilon }{\delta \mathrm{\Phi }}}{\displaystyle \frac{d\mathrm{\Delta }_>}{d\epsilon }}{\displaystyle \frac{\delta S_\epsilon }{\delta \stackrel{~}{\mathrm{\Phi }}}}+i\mathrm{tr}{\displaystyle \frac{d\mathrm{\Delta }_>}{d\epsilon }}{\displaystyle \frac{\delta ^2S_\epsilon }{\delta \stackrel{~}{\mathrm{\Phi }}\delta \mathrm{\Phi }}}.`$ (10)
Next task is to determine the functional $`S_\epsilon `$ by solving this equation. To this end, let us first consider the $`\stackrel{~}{\mathrm{\Phi }}`$ dependence of the functional $`S_\epsilon `$. The bare $`S_{\epsilon _0}`$ contains only the first order term in $`\stackrel{~}{\mathrm{\Phi }}`$, but in the process of the renormalization, Eq. (10) yields the zeroth order term in $`S_\epsilon `$. To be concrete, let us denote $`S_\epsilon [\stackrel{~}{\mathrm{\Phi }},\mathrm{\Phi }]=\stackrel{~}{\mathrm{\Phi }}H_\epsilon [\mathrm{\Phi }]+F[\mathrm{\Phi }]`$. Substituting this into Eq.(10), we find
$`{\displaystyle \frac{H_\epsilon }{\epsilon }}`$ $`=`$ $`{\displaystyle \frac{\delta H_\epsilon }{\delta \mathrm{\Phi }}}{\displaystyle \frac{d\mathrm{\Delta }_>}{d\epsilon }}H_\epsilon ,`$ (11)
$`{\displaystyle \frac{F_\epsilon }{\epsilon }}`$ $`=`$ $`{\displaystyle \frac{\delta F_\epsilon }{\delta \mathrm{\Phi }}}{\displaystyle \frac{d\mathrm{\Delta }_>}{d\epsilon }}H_\epsilon +i\mathrm{tr}{\displaystyle \frac{d\mathrm{\Delta }_>}{d\epsilon }}{\displaystyle \frac{\delta H_\epsilon }{\delta \mathrm{\Phi }}}`$ (12)
with the bare functions $`H_{\epsilon _0}[\mathrm{\Phi }]=g\theta (_t\mathrm{\Phi })\kappa _x^2\mathrm{\Phi }`$ and $`F_{\epsilon _0}[\mathrm{\Phi }]=0`$. The RG equation for the $`H_\epsilon `$-term, which determines the solution of the Barenblatt equation, is closed. Furthermore, it has no loop corrections. Nevertheless, the initial-value problems are still nontrivial, since after obtaining $`H_\epsilon `$, we must set $`J=q\delta (x)\delta (t)`$ and determine the $`\epsilon `$-dependence.
To solve the functional equation (11), we assume the form of $`H_\epsilon `$ as
$`H_\epsilon [\mathrm{\Phi }](x,t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑s{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑yV_\epsilon [_t\mathrm{\Phi }(xy,ts)]\kappa _y^2\mathrm{\Phi }(y,s)`$ (13)
$``$ $`V_\epsilon [\dot{\mathrm{\Phi }}]\kappa \mathrm{\Phi }^{\prime \prime }(x,t)`$
with a certain unknown function $`V_\epsilon `$, where we have denoted $`_t\mathrm{\Phi }=\dot{\mathrm{\Phi }}`$ and $`_x^2\mathrm{\Phi }=\mathrm{\Phi }^{\prime \prime }`$ for simplicity. Substituting this into Eq. (11), we have
$`_\epsilon V_\epsilon [\dot{\mathrm{\Phi }}]=V_\epsilon [\dot{\mathrm{\Phi }}]_\epsilon \kappa \mathrm{\Delta }_>^{\prime \prime }V_\epsilon [\dot{\mathrm{\Phi }}].`$ (14)
The bare function is given by $`V_{\epsilon _0}[\dot{\mathrm{\Phi }}](xy,ts)=g\theta (\dot{\mathrm{\Phi }}(x,t))\delta (xy)\delta (ts)`$. This equation can be solved if $`V_\epsilon `$ is expanded in power series of $`g`$ such that $`V_\epsilon =_{n=1}^{\mathrm{}}g^nV_\epsilon ^{(n)}`$. Actually, we find that each term obeys
$`_\epsilon V_\epsilon ^{(n)}[\dot{\mathrm{\Phi }}]={\displaystyle \underset{n_1+n_2=n}{}}V_\epsilon ^{(n_1)}[\dot{\mathrm{\Phi }}]_\epsilon \kappa \mathrm{\Delta }_>^{\prime \prime }V_\epsilon ^{(n_2)}[\dot{\mathrm{\Phi }}].`$ (15)
From this equation, it follows that $`_\epsilon V_\epsilon ^{(1)}=0`$ and hence, it turns out that $`V_\epsilon ^{(1)}`$ is not renormalized; that is, $`V_\epsilon ^{(1)}[\dot{\mathrm{\Phi }}]V[\dot{\mathrm{\Phi }}]`$, where
$`V[\dot{\mathrm{\Phi }}](xy,ts)=\theta (\dot{\mathrm{\Phi }}(x,t))\delta (xy)\delta (ts).`$ (16)
This enables us to calculate the higher order solutions by substituting Eq. (16) into Eq. (15) and by solving it successively order by order:
$`V_\epsilon ^{(n)}=V(\kappa \mathrm{\Delta }_>^{\prime \prime }V)^{n1}.`$ (17)
Thus, we have determined the functional $`H_\epsilon [\mathrm{\Phi }]`$ as infinite power series with respect to $`g`$. In passing, we briefly discuss the solution for $`F_\epsilon `$. Expanding similarly the RG equation (12) in power series of $`g`$, we find that the solution of Eq. (12) is $`F_\epsilon =0`$ because of the fact that the bare $`F_{\epsilon _0}=0`$ as well as that the second term in the r.h.s of Eq. (12) is zero due to the trace with respect to the time variable. Therefore, we end up with the renormalized action functional,
$`S_\epsilon `$ $`=`$ $`\stackrel{~}{\mathrm{\Phi }}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}g^nV[\dot{\mathrm{\Phi }}](\kappa \mathrm{\Delta }_>^{\prime \prime }V[\dot{\mathrm{\Phi }}])^{n1}\kappa \mathrm{\Phi }^{\prime \prime }`$ (18)
$``$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}g^nS_\epsilon ^{(n)}`$
with $`V`$ defined by Eq. (16). Thus, we have determined the renormalized action in full order in $`g`$. This shows the efficiency of the present approach.
To obtain the asymptotic solution of the Barenblatt equation, we must set $`J(x,t)=q\delta (x)\delta (t)`$ to specify the initial condition and calculate dominant parts with respect to $`\epsilon /\epsilon _0`$ in the action obtained so far. The action decomposed into some sectors by the substitution of $`\stackrel{~}{\mathrm{\Phi }}=\stackrel{~}{J}\mathrm{\Delta }_>+\stackrel{~}{\varphi }_<`$ and $`\mathrm{\Phi }=\mathrm{\Delta }_>J+\varphi _<`$. Then, we notice that among them dominant contributions are from the $`\stackrel{~}{\mathrm{\Phi }}(x,t)=\stackrel{~}{\varphi }_<(x,t)`$ and $`\mathrm{\Phi }(x,t)=q\mathrm{\Delta }(x,t)`$ sectors in the asymptotic region $`\epsilon t`$.
Let us start with the first order which Goldenfeld et al. have calculated within the perturbation in the iteration scheme. From Eq. (18), one can obtain
$`S_\epsilon ^{(1)}`$ $`=`$ $`q{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x\stackrel{~}{\varphi }_<(x,t)\theta (q\dot{\mathrm{\Delta }}_>(x,t))\kappa \mathrm{\Delta }_>^{\prime \prime }(x,t)`$ (19)
$`=`$ $`q\left[{\displaystyle \frac{1}{\sqrt{2\pi e}}}\mathrm{ln}{\displaystyle \frac{\epsilon }{\epsilon _0}}\right]\stackrel{~}{\varphi }_<(0,0)+qS_{\mathrm{reg}}`$
where the regular part of $`S^{(1)}`$ is defined as
$`S_{\mathrm{reg}}^{(1)}={\displaystyle _{\epsilon _0}^\epsilon }{\displaystyle \frac{dt}{t}}{\displaystyle _1^1}𝑑\omega \delta \stackrel{~}{\varphi }_<(\sqrt{2\kappa t}\omega ,t)f(\omega )`$ (20)
with $`\delta \stackrel{~}{\varphi }(\sqrt{2\kappa t}\omega ,t)\stackrel{~}{\varphi }_<(\sqrt{2\kappa t}\omega ,t)\stackrel{~}{\varphi }_<(0,0)`$ and $`f(\omega )\frac{\omega ^21}{2}\frac{e^{\omega ^2/2}}{\sqrt{2\pi }}`$. Here, the regular term in Eq. (19) is not involved with the renormalization of the action since we can safely set $`\epsilon _00`$, while the first term is relevant to the renormalization of $`q`$, the height of the initial distribution. Thus it turns out that at this order $`q`$ is indeed renormalized, whereas others, especially $`g`$, is not renormalized. This feature actually holds even in the next order, as will be checked below. This implies that the present system is always at a fixed point because $`g`$ is not renormalized, and hence, the anomalous dimension which is in general a scheme-dependent quantity is a physical observable in the present case. Considering these, we introduce the renormalization only to $`q`$ and define the renormalized $`q`$ as $`q_R=qZ`$. Expanding the renormalization constant $`Z`$ as $`Z=1+_{n=1}g^nZ^{(n)}`$ with $`Z^{(n)}=\gamma ^{(n)}\mathrm{ln}(\epsilon /\epsilon _0)`$, it turns out that the first order of $`Z`$ reads
$`\gamma ^{(1)}={\displaystyle \frac{1}{\sqrt{2\pi e}}},`$ (21)
from Eq. (19). This indeed reproduces the result of Goldenfeld et al. GMOL90 .
The renormalization of $`q`$ introduced above is indeed enough to render the solution of the Barenblatt equation finite also in higher order. To verify this, let us next calculate the second order renormalized action. Due to similar arguments to the first order, the action (18) yields
$`S_\epsilon ^{(2)}`$ $`=`$ $`q{\displaystyle _{2\epsilon _0}^{\epsilon _0+\epsilon }}𝑑t{\displaystyle _{\sqrt{2\kappa t}}^{\sqrt{2\kappa t}}}𝑑x{\displaystyle _{\epsilon _0}^{t\epsilon _0}}𝑑t_1{\displaystyle _{\sqrt{2\kappa t_1}}^{\sqrt{2\kappa t_1}}}𝑑x_1`$ (22)
$`\times \stackrel{~}{\varphi }_<(x,t)\kappa \mathrm{\Delta }^{\prime \prime }(xx_1,tt_1)\kappa \mathrm{\Delta }^{\prime \prime }(x_1,t_1)`$
$`+\text{reg. },`$
where reg. stands for regular parts of the renormalized action. Similar but a bit lengthy calculation leads to
$`S_\epsilon ^{(2)}`$ $`=`$ $`q{\displaystyle \frac{Z^{(1)2}}{2}}\stackrel{~}{\varphi }(0,0)+qZ^{(1)}S_{\mathrm{reg}}^{(1)}+qZ^{(2)}\stackrel{~}{\varphi }(0,0)`$ (23)
$`+\text{reg.}`$
where $`Z^{(2)}=\gamma ^{(2)}\mathrm{ln}(\epsilon /\epsilon _0)`$ with
$`\gamma ^{(2)}={\displaystyle _0^1}{\displaystyle \frac{d\tau _1}{\tau _1}}{\displaystyle _1^1}𝑑\omega _1f(\omega _1)`$ (24)
$`\times `$ $`\left[{\displaystyle \frac{1}{\sqrt{2\pi e}}}{\displaystyle \frac{1\sqrt{\tau _1}\omega _1}{1\tau _1}}{\displaystyle \frac{e^{(1\sqrt{\tau _1}\omega _1)^2/\left[2(1\tau _1)\right]}}{\sqrt{2\pi (1\tau _1)}}}\right].`$
This result indicates that the second order action correctly includes the contributions from the first order renormalized action. Namely, the renormalized action with the source term $`q\stackrel{~}{\varphi }(0,0)`$ satisfies $`gS_\epsilon ^{(1)}+g^2S_\epsilon ^{(2)}+q\stackrel{~}{\varphi }(0,0)=e^ZqS_{\mathrm{reg}}+e^Zq\stackrel{~}{\varphi }(0,0)`$ up to the second order of $`g`$. Therefore, we expect in general that $`q_R=q\left(\epsilon _0/\epsilon \right)^\gamma `$ where
$`\gamma ={\displaystyle \underset{n=1}{}}g^n\gamma ^{(n)}.`$ (25)
The constant $`\gamma `$ thus obtained indeed gives the anomalous dimension of the solution for the present diffusion equation.
To see this, notice that since $`\stackrel{~}{J}`$ is not renormalized, we have $`u(x,t;q,\epsilon _0)=u(x,t;q_R,\epsilon )`$, which tells that the solution is independent of $`\epsilon `$. Hence, the renormalized solution should satisfy the following RG equation;
$`\left(\epsilon {\displaystyle \frac{}{\epsilon }}\gamma q_R{\displaystyle \frac{}{q_R}}\right)u(x,t;q_R,\epsilon )=0,`$ (26)
where $`\gamma `$ is defined by Eq. (25), or alternatively by $`\gamma =\mathrm{ln}q_R/\mathrm{ln}\epsilon `$. On the other hand, the dimensional analysis requires $`u(x,t;q,\epsilon _0)=q/\sqrt{\kappa t}\mathrm{\Phi }(x/\sqrt{\kappa t},\epsilon _0/t)`$. Therefore, combining these observations, we can assume $`u(x,t;q_R,\epsilon )=q_R/\sqrt{\kappa t}\mathrm{\Phi }(x/\sqrt{\kappa t},\epsilon /t)`$. Substituting this into Eq. (26), it turns out that the relation $`\mathrm{\Phi }(x/\sqrt{\kappa t},\epsilon /t)=\left(\epsilon /t\right)^\gamma \stackrel{~}{\mathrm{\Phi }}\left(x/\sqrt{\kappa t}\right)`$ holds. Hence, the asymptotic behavior of the solution for the Barenblatt equation is indeed given by $`u1/t^{1/2+\gamma }`$.
So far we have derived the anomalous diffusion exponent up to second order with respect to $`g`$. The exponent of first order is the same as that obtained by Goldenfeld et al., whereas the exponent of second order (24) is estimated as $`\gamma ^{(2)}0.063546`$ by numerical integration. Though this value is different from the result in Ref. GMOL90 , it coincides with that obtained later by Cole and Wagner given by a different integral formula derived via different method ColWag96 . Since the anomalous diffusion exponent should be scheme-independent in the present case, we believe that our result as well as Cole and Wagner’s is correct.
The efficiency of our method lies in the fact that one can compute the higher order exponent in a similar way above. To present the exponent of $`n`$th order, we define a function,
$`g(\omega _1,\omega _2,\tau _2)={\displaystyle \frac{e^{(\omega _1\sqrt{\tau _2}\omega _2)^2/[2(1\tau _2)]}}{\sqrt{2\pi (1\tau _2)}}}{\displaystyle \frac{e^{\omega _2^2/2}}{\sqrt{2\pi }}}.`$ (27)
Then, Eq. (18) yields
$`\gamma ^{(n)}`$ $`=`$ $`{\displaystyle _0^1}{\displaystyle \underset{j=1}{\overset{n1}{}}}{\displaystyle \frac{d\tau _j}{\tau _j}}{\displaystyle _1^1}{\displaystyle \underset{j=1}{\overset{n1}{}}}d\omega _jf(\omega _{n1})`$ (28)
$`\times `$ $`\left[{\displaystyle \frac{1}{\sqrt{2\pi e}}}{\displaystyle \frac{1\sqrt{\tau _1}\omega _1}{1\tau _1}}{\displaystyle \frac{e^{(1\sqrt{\tau _1}\omega _1)^2/\left[2(1\tau _1)\right]}}{\sqrt{2\pi (1\tau _1)}}}\right]`$
$`\times {\displaystyle \underset{j=1}{\overset{n2}{}}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2g(\omega _j,\omega _{j+1},\tau _{j+1})}{d\omega _j^2}}.`$
We expect that this anomalous exponent is exact, since it has been derived from one formula (18), among which the first and second order coincides with those calculated by different methods. For references, we numerically estimate the exponent of third order, $`\gamma ^{(3)}=0.00314`$. The exponent obtained so far seems a good convergent series.
In summary, we have applied the ERG techniques to the initial-value problem of the Barenblatt equation, one of typical nonlinear diffusion equations. We have derived the anomalous diffusion exponent in full order with respect to the parameter controlling the nonlinearity. This implies that the ERG approach is efficient to systems far from equilibrium described by nonlinear partial-differential equations as well as to field theories and statistical mechanics.
The present approach would be useful to other types of initial-value problem for nonlinear diffusion equations. In particular, application to the critical dynamic of nonlinear traveling wave is of great interest. One of well known examples is the KPP equation KPP which shows an interesting universal behavior in the selection of the front velocity. This problem has been addressed by Paquette et al. PCGO94 , but more detailed analysis is needed if one wants to understand, e.g., the universal logarithmic corrections to the velocity in the pulled-front, from the RG point of view. Application of the ERG to such problems would be quite interesting.
One of the author (TF) would like to thank T. Ohta and T. Kunihiro for valuable discussions. This work was supported in part by Grant-in-Aid for Scientific Research from JSPS.
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# Comment on “Jamming at zero temperature and zero applied stress: The epitome of disorder”
## I What is “jammed”?
OSLN question whether the hard-sphere system is “physical” and therefore resort to studying particle systems with soft-sphere interactions to mimic hard-particle packings. The latter is inherently a geometrical problem. In fact, there is a simple and rigorous geometrical approach to jamming in hard-sphere systems that is not only well-defined, but, as we show below, is closely related to OSLN’s jamming point $`J`$.
Although the hard-sphere potential is an idealization, it is no less physical than any soft-sphere potential, especially in regards to jamming. Indeed, the singular nature of the hard-sphere potential is crucial because it enables one to be precise about the concept of “jamming.” Recently, three hierarchically ordered jamming categories have been introduced Torquato and Stillinger (2001): *local*, *collective* and *strict* *jamming*. Each successive category progressively relaxes the boundary conditions imposed on the particle displacements. These definitions are very intuitive and completely *geometric*, and are closely linked to definitions of “rigid” or “stable” packings appearing in the mathematics literature Connelly (1988); Connelly et al. (1998).
OSLN’s definition of jamming simply states that the configuration of particles is at a stable (strict) energy minimum. Such a definition is dependent on the particular interparticle potential, and thus it obscures the relevant packing geometry (exclusion-volume effects). Furthermore, the distinctions between different jamming categories is critical, especially if one is trying to determine the density of the MRJ state Donev et al. (2004). Specifically, this density will generally be higher the more demanding is the jamming category. Clearly, OSLN do not distinguish between different degrees or levels of jamming. We have recently demonstrated that the distinction between collective and strict jamming is important even for very large packings, especially in two dimensions Donev et al. (2004) .
For OSLN, a jammed configuration is one where there are no zero-frequency modes of the Hessian matrix of the total potential energy with respect to the positions of the particles (the dynamical matrix), while keeping the periodic unit cell *fixed*. Our definition of strict jamming relaxes this requirement and includes the lattice vectors as additional degrees of freedom Donev et al. (2003). As explained in detail in Ref. Connelly and Whiteley (1996), the Hessian consists of two parts, a negative definite *stress matrix* and a positive semidefinite *stiffness matrix*. OSLN’s definition of jammed means simply that the Hessian is positive definite at the energy minimum. A precise phrase for this is *a stable* or *strict (local) energy minimum*, and we see no point in redefining this elementary concept. In fact, according to OSLN, any stable energy minimum represents a jammed configuration, and it is not possible to relate this idea to *packing* concepts without numerous additional assumptions about the form of the pair potential and the interparticle distances at the energy minimum.
Although OSLN point out themselves that our definitions of jamming and MRJ are for *hard-sphere packings*, they claim to replace them with a “cleaner definition,” which applies *only* to systems of *soft spheres*. The two definitions cannot directly be compared as they apply to different systems. However, OSLN themselves clearly imply that their “jammed” soft-sphere systems and “jammed” hard-sphere packings are related, by referring to other works on hard-sphere systems. For example, they claim a direct relation between their special point $`J`$ and RCP of hard spheres in Section IID of Ref. O’Hern et al. (2003). The basic idea, as OSLN explain, is that one “can approach the hard sphere by making the potential harder and harder…\[to\] produce a limiting hard sphere value”. However, they question whether the hard-sphere limit is well-defined and “would argue that hard spheres are a singular limit and thus unphysical” and that “One should therefore concentrate on softer potentials for which unambiguous definitions can be constructed.”
To demonstrate that the limit is well-defined, let us first define a *collectively jammed sphere packing* to be *any nonoverlapping configuration of hard spheres in which no subset of particles can continuously be displaced so that its members move out of contact with one another and with the remainder set* (while maintaining nonoverlap) Torquato and Stillinger (2001). The following theorem Connelly and Whiteley (1996) shows that near the “jamming threshold” $`\varphi _c`$, as defined in Section IIB in Ref. O’Hern et al. (2003), the jamming of particle systems as defined by OSLN is directly related to this definition of collective jamming in hard-sphere packings:
Theorem: Consider an interparticle potential that is continuous and strictly monotonically decreasing around $`r_{ij}=D`$, and vanishes for $`r_{ij}>D+\delta `$. If in a finite configuration of particles interacting with such a potential, all interacting (i.e., closer then $`D+\delta `$) particles are a distance $`D`$ apart, and the configuration is a stable local energy minimum, then the configuration corresponds to a collectively jammed packing of hard spheres with diameter $`D`$.
If one relaxes the condition that all interacting particles are exactly at distance $`D`$ apart and instead asks only that the minimum interparticle distance be $`D`$, then for a sufficiently small $`\delta `$ one can prove Connelly (1982) that the above sphere packing is almost collectively jammed (i.e., it is trapped in a small neighborhood of the initial configuration Donev et al. (2003)). This theorem implies that the packings studied in Ref. O’Hern et al. (2003) that are very slightly above the “jamming threshold” $`\varphi _c`$ are indeed closely related to collectively jammed ideal packings of spheres of diameter $`D=\sigma `$ (polydispersity is trivial to incorporate). All of these considerations call into question the value of a definition of jamming that hinges on eigenvalues of dynamical matrices.
Finally, it is important to note that despite the fact that our definition of collective jamming above calls for virtual displacing (groups of) particles, one can in fact rigorously test for our hard-particle jamming categories using linear programming Donev et al. (2003), without what OSLN call “shifting particles,” even for very large disordered packings Donev et al. (2004). We have in fact communicated to OSLN the results foo (b) of our algorithm applied to several sample packings provided by them. In short, our algorithm verified that OSLN’s systems near $`\varphi _c`$ were indeed nearly collectively jammed (within a very small tolerance) when viewed as packings. However, they were not strictly jammed because OSLN keep the lattice vectors fixed during energy minimization.
## II What is “random”?
We agree that the maximum of an appropriate “entropic” metric would be a potentially useful way to characterize the randomness of a packing and therefore the MRJ state Kansal et al. (2002). However, as pointed out in Ref. Kansal et al. (2002), a substantial hurdle to be overcome is the necessity to generate all possible jammed states, or at least a representative sample of such states, in an unbiased fashion using a “universal” protocol in the large-system limit. Even if such a protocol could be developed, however, the issue of what weights to assign the resulting configurations remains. Moreover, there are other fundamental problems with entropic measures, as we discuss below, including its significance for two-dimensional monodisperse hard disk packings.
According to OSLN, maximally random is defined by “where the entropy of initial states is maximum” and imply that this is a universal measure of disorder. It is not clear exactly what the authors mean by entropy and how (or whether) it can be measured for “initial states”. It is not obvious that one can relate the “randomness” of the *final* configurations (which is what OSLN are analyzing) to that of the *initial* configurations. It appears OSLN’s rationale is that their algorithm goes to the “nearest” energy minimum from a given initial configuration. Does this process preserve “entropy” or randomness? Clearly if one used, for example, global energy minimization, one would obtain very different results. Furthermore, entropy is a concept inherently related to *distributions* of configurations. However, one classifies *particular* final configurations (packings) as random or disordered, and by considering a given configuration, one can devise a procedure for quantitatively measuring (using order metrics) how disordered or ordered it is. This distinction between distributions of configurations and particular configurations is an important one that OSLN do not make.
The MRJ state is defined in Torquato et al. (2000) as the jammed state which minimizes a given order metric $`\psi `$. OSLN suggest their interpretation of maximally random as superior because using order metrics “will always be subject to uncertainty since one never knows if one has calculated the proper order parameter.” Therefore, OSLN believe that they have identified the proper, unique, measure of order (related to entropy). We wish to stress the difference between *well-defined* and *unique*, as the two seem to be blurred in Ref. O’Hern et al. (2003). The MRJ state is well-defined in that for a particular choice of jamming category and order metric it can be identified unambiguously. For a finite system, it will consist of a discrete set (possibly one) of configurations, becoming more densely populated as the system becomes larger. At least for collective and strict jamming in three dimensions, a variety of sensible order metrics seem to produce an MRJ state near $`\varphi 0.64`$ Kansal et al. (2002), the traditionally accepted density of the RCP state.
However, the *density* of the MRJ state should not be confused with the MRJ *state* itself. It is possible to have a rather ordered packing at this very same density; for example, a jammed but diluted vacancy FCC lattice packing Kansal et al. (2002). This is why the two-parameter description of packings in terms of the density $`\varphi `$ and order metric $`\psi `$, as in Ref. Torquato et al. (2000), is not only useful, but actually necessary.
OSLN’s description of order implies a direct relation between probability densities and randomness, i.e., that the *most probable* foo (c) configurations represent the most disordered state. In this sense, one expects that the density of jammed configurations, when viewed as a three-dimensional plot over the $`\varphi \psi `$ plane will be very strongly peaked around the MRJ point for very large systems, just as the probability distribution curves in Fig. 6 in Ref. O’Hern et al. (2003) are very peaked around $`\varphi 0.64`$. As OSLN suggest themselves, this might explain why several different packing procedures yield similar hard-particle packings under appropriate conditions, historically designated as RCP.
However, this is far from being a closed question foo (d). Consider two-dimensional monodisperse circular disk packings as an example. It is well-known that two-dimensional analogs of three-dimensional computational and experimental protocols that lead to putative RCP states result in disk packings that are highly crystalline, forming rather large triangular domains (grains) foo (e). Because such highly ordered packings are the most probable for these *protocols*, OSLN’s entropic measure would identify these as the most disordered, a dubious proposition. An appropriate order metric, on the other hand, is capable of identifying a particular configuration (not an ensemble of configurations) of considerably lower density (e.g., a jammed diluted triangular lattice) that is consistent with our intuitive notions of maximal disorder. However, typical packing protocols would almost never generate such disordered disk configurations because of their inherent bias toward undiluted crystallization. This brings us to OSLN’s claim that they have devised an unbiased universal protocol, to which we now turn our attention.
## III Universal, Hard and Soft Algorithms
In this section, we focus on the algorithms used by OSLN and point out why they are neither universal nor superior to other procedures. We point out the close relations between OSLN’s algorithm for generating configurations near the onset of jamming and the Zinchenko hard-sphere packing algorithm Zinchenko (1994). Furthermore, we question OSLN’s implication that using one kind of interaction potential (with three different exponents) and one algorithm amounts to exploring the space of all jammed configurations in an unbiased manner. This puts into doubt the claimed universality of the point $`J`$.
By fixing the interaction potential, initial density and energy minimization (conjugate gradient) algorithm, OSLN obtain a well-defined collection of final configurations with well-defined (not unique) properties. In essence, OSLN make their algorithm devoid of tunable parameters by simply choosing specific and fixed values for them. Both the Zinchenko and OSLN algorithms are “dynamics independent”, in the sense that there is no tunable parameter for the rate of compression, which would be an analog of the cooling or quenching rate in molecular systems. Both also imply that this makes their algorithm universal or superior to other algorithms and that the (well-defined) results they obtain are somehow special. Most sensible algorithms will in fact produce a well-defined density in the limit of large systems given a choice of algorithmic parameters. For example, by changing the expansion rate in the Lubachevsky-Stillinger algorithm, one can achieve final densities for spheres anywhere in the range from $`0.64`$ (fast expansion) to $`0.74`$ (very slow expansion), as clearly illustrated in Fig. 2a in Ref. Torquato et al. (2000). Therefore, if we followed the logic of OSLN, we could claim that any number in that range represents a special point. In our opinion, a good packing algorithm should be capable of generating a variety of packings, in both density and the amount of order. How can one ascertain that the packings one produces are “most random” if there are no other jammed packings to compare to?
OSLN use two main procedures to generate final configurations. The first procedure is to choose a density and then use conjugate gradients to find a nearby energy minimum, starting from a randomly-generated initial configuration ($`T=\mathrm{}`$), as described in Section IIA in Ref. O’Hern et al. (2003). Using this procedure, OSLN sampled inherent structures Stillinger and Weber (1985) at fixed density to measure the fraction $`f_j(\varphi )`$ of states that had nonzero bulk and shear moduli, and showed that $`f_j`$ has a strong system-size dependence with its derivative becoming a delta-function in the large system limit. It is important to note that this procedure as such has little or nothing to do with hard-sphere packings, especially for the kind of soft potentials ($`\alpha 3/2)`$ that OSLN study. Many stable energy minima will be completely unrelated to packings, and especially not to those designated as MRJ states.
OSLN used a second procedure to study the mechanical and structural properties of systems near the onset of jamming $`\varphi _c`$. In this procedure, a configuration is compressed (or decompressed) using very small steps in density until the bulk and shear moduli vanished (or nonzero moduli develop), as described in Section IIB in Ref. O’Hern et al. (2003). We now demonstrate that this procedure is closely related to Zinchenko’s algorithm Zinchenko (1994) for generating hard-sphere packings. Start at low density with a set of nonoverlapping spheres of diameter $`\sigma `$. Both algorithms then slowly grow the particles (OSLN in small increments, Zinchenko continuously) while moving the particles to avoid overlap foo (f). In the Zinchenko algorithm, one strictly maintains the contact between particles as soon as they touch, which requires solving a system of ODE’s containing the rigidity matrix of the packing Donev et al. (2003) to find the necessary particle displacements. OSLN on the other hand, use conjugate gradients (CG) to reminimize the potential energy, which will simply push the particles just enough to keep them nonoverlapping, i.e., almost in contact. This procedure continues in both algorithms until no further densification is possible without inducing overlap.
Accordingly, it is not surprising the packing configurations close to $`\varphi _c`$ obtained in Ref. O’Hern et al. (2003) closely resemble (in packing fraction, amorphous character, coordination, etc.) packings generated via a variety of bonafide hard-sphere algorithms (and experiments Scott and Kilgour (1969)). In particular, very similar packings are produced with the Lubachevsky-Stillinger (LS) algorithm Lubachevsky and Stillinger (1990); Lubachevsky et al. (1991) (with sufficiently high expansion rates) and the Zinchenko algorithm Zinchenko (1994). OSLN criticize the LS algorithm for changing the density in a dynamic fashion. The stated advantage of the OSLN protocol is that one can “quench the system to the final state within a fixed energy landscape” since “the density is always held constant”. We are very puzzled by this last claim in light of their admission (in Section IIB of Ref. O’Hern et al. (2003)) that they slowly change the density of the packing to find $`\varphi _c`$. In fact, OSLN do not seem to clearly distinguish between the two rather different procedures they employ: the first for finding inherent structures (at a fixed density) and the second for generating packings at the jamming threshold (which searches in density). Fig. 6 of Ref. O’Hern et al. (2003), which supposedly represents the distributions of jamming thresholds $`\varphi _c`$, *defined* by the *second* procedure, is obtained by differentiating the distribution generated with the *first* procedure, with no clear justification.
Most problematic of all is OSLN’s claim that their results are universal. Despite the statement that “Starting with randomly generated $`T=\mathrm{}`$ states guarantees that we sample *all* \[emphasis added\] phase space equally”, all that their first algorithm manages to explore is the space of energy minima for the *particular* chosen interaction potential. By comparing three different exponents $`\alpha `$, OSLN conclude that the exact form of the potential is not important. However, a much more convincing picture would have been made if they instead tried *qualitatively* different kinds of interaction potentials, rather then simply changing the curvature of the potential at the contact point. Otherwise, why focus on continuous interaction potentials at all? Since it is geometry (i.e., the nonoverlap condition on the spherical cores) that is crucial, the hard-sphere system offers a far “cleaner” system to study when trying to understand the special point $`J`$.
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# Moduli evolution in the presence of flux compactifications
## I INTRODUCTION
String theory has associated with it many moduli fields, scalar fields whose presence are important in determining quantities such as the size of the internal dimensions, the gauge coupling constants, even the strength of the gravitational interactions. One of the problems that has faced physicists attempting to develop cosmology within string theory has been how to deal with these fields – basically they never want to stop evolving. In most models that have been developed to date, the moduli initially have no associated potential with which to trap them, the flat directions mean that there is nothing to prevent them from evolving for ever, leading to time varying coupling constants, decompactification of the internal dimensions and other nightmares. Non trivial potentials have been shown to emerge when non-perturbative features are included, such as gaugino condensates gaugino , but even then the local minima associated with these models are generally unstable in that the barrier height protecting it from the runaway solution is small compared to any other natural scale in the problem. Moreover, the value of the potential at its minima is negative, implying an Anti de Sitter Space solution racetrack . This issue has plagued string models for many years. Initially, in the context of weakly coupled heterotic string theory, this was pointed out by Dine and Seiberg dine-seiberg . In the context of cosmology Brustein and Steinhardt readdressed the issue, pointing out how difficult it would be to stabilise typical moduli fields emerging in these models because of the combination of the steepness of the potentials associated with the moduli, and the small size of the potential barrier separating the fields from the runaway regime brustein-steinhardt . Earlier, in ref. Kaloper:1991mq , Kaloper and Olive had studied the effect including a background of radiation in the equations that govern the dilaton evolution. In particular, they found that the dilaton can stop rolling if either its potential has a SUSY breaking minimum or radiation is dominating the energy density of the Universe. In Barreiro:1998aj this idea was further developed, in particular it was pointed out that, in the presence of an extra fluid component dominating the energy density, for example radiation, the effect of the fluid was to increase the effective friction experienced by the moduli field as it evolved down its steep potential. For a wide range of initial conditions of the field it was demonstrated that it would be slowed down, and its energy density would tend towards a scaling solution where it followed the evolution of the background radiation energy density. Moreover, it would then simply fall into the minimum of the potential and be stabilised there, even though the height of the potential barrier was so low. This analysis, which was initially used to demonstrate how the dilaton could be stabilised in weakly coupled heterotic string theory, was later extended to the case of heterotic M-theory, where including the evolution of more than one modulus field, it became evident how important the coupling between the various moduli fields could be in determining the basin of attraction into the true minimum of the potential heterotic-m-theory . Other attempts at stabilising the moduli fields in these models included adding specific temperature corrections to the moduli fields, in order to extend the range of initial values of the fields which would lead to stabilisation huey , although this met with limited success heterotic-m-theory .
Over the past few years there has been a renewed interest in the issue of moduli evolution. Giddings et al demonstrated, in the context of Type IIB string theory compactified on a Calabi-Yau manifold, how it is possible to use fluxes to stabilise all but the Kähler moduli, in particular the overall volume modulus GKP . This result was used to develop inflation models in string theory where brane inflation was embedded into one of these stable compactified models Kachru:2003aw ; Kachru:2003sx . However in Kachru:2003sx it was realised that there was a price to pay, and a degree of fine tuning was required in these models in order to obtain sufficient inflation. The particular problem they encountered was that, as the compactification volume modulus was being stabilised, the effect was to modify the inflaton potential, rendering it too steep for inflation to last for a sufficient length of time. The issue of stabilisation had raised its ugly head again. Recently Brustein et al Brustein:2004jp returned to the idea of using sources as a way of increasing the friction and thereby avoiding the runaway problem alluded to earlier. They did it in the context of a toy model with a single scalar field (the volume modulus) and scalar potential of the form suggested in Kachru:2003sx . They argued that the key reason they could stabilise the modulus was because of the way this field quickly entered a regime where its energy density was dominated by its kinetic energy. It therefore lost its energy far quicker than any other source, allowing the other sources (such as radiation) to eventually catch up and take over, slowing the field down even more and allowing it to come to rest at or near the local minimum of its potential. In this paper we extend the investigation into these models by allowing for the fact that the volume modulus is a complex field with a real and imaginary component. We therefore look at the system where more than one moduli field has to be dynamically stabilised in the presence of extra sources. We reproduce the results of Brustein:2004jp in the appropriate limits, but show that there are wider ranges of initial conditions which can lead to stability if we include the possibility of the moduli fields entering a period of scaling, where their energy density mimics that of the background energy density, more in the spirit of Barreiro:1998aj ; heterotic-m-theory . A related approach to the stabilisation issue has recently been investigated in refs. Dimopoulos:2002hm ; Kaloper:2004yj ; Easson:2005ug ; Berndsen:2005qq , where the authors have considered the presence of a gas of wrapped branes on the dynamics of some of the moduli fields.
The layout of this paper is as follows. In section II we introduce the system of equations we will be investigating in terms of a general moduli potential, and briefly describe the nature of the attractor solutions when exponential terms are present in the scalar potential. The particular example due to Kachru et al Kachru:2003sx is analysed in section III and a related but different example due to Kallosh and Linde Kallosh:2004yh is examined in section IV. Finally we conclude in section V.
## II BACKGROUND EQUATIONS OF MOTION
The class of models that we will be investigating can in general be described by an N=1, d=4 effective Supergravity (SUGRA) theory, where d is the number of non-compact spacetime dimensions. In this case the four-dimensional N=1 SUGRA action is of the form
$$S=\sqrt{g}\left(\frac{1}{2\kappa _P^2}R+K_{i\overline{j}}_\mu \mathrm{\Phi }^i^\mu \overline{\mathrm{\Phi }}^{\overline{j}}+V\right)d^4x,$$
(1)
where $`K_{i\overline{j}}=\frac{^2K}{\mathrm{\Phi }^i\overline{\mathrm{\Phi }}^{\overline{j}}}`$ is the Kähler metric; $`\mathrm{\Phi }^i`$ are complex chiral superfields; V($`\mathrm{\Phi }`$) is the scalar potential and $`\kappa _P`$ is the 4-dimensional Newton constant which, in terms of higher dimensional quantities, can be expressed as
$$\kappa _P^2=\frac{\kappa ^2}{2\pi \rho v}=8\pi G_N,$$
(2)
where $`\kappa `$ is the 11-dimensional Newton constant, $`v`$ is the volume of the six dimensional compact manifold, and $`\rho `$ is the radius of the eleventh spatial dimension.
The impact of the non-perturbative and flux effects considered in Kachru:2003aw ; Kachru:2003sx is to induce an effective scalar potential for the volume moduli whose most general form for 4-dimensional $`N=1`$ SUGRA is
$$V=e^K(K^{i\overline{j}}D_iW\overline{D_jW}3W\overline{W}).$$
(3)
where $`K^{i\overline{j}}`$ is the inverse Kähler metric and $`D_iW=_iW+\frac{K}{\mathrm{\Phi }^i}W`$ is the Kähler covariant derivative acting on the superpotential.
The equations of motion follow from the variation of the action (1). For simplicity we consider the case of homogeneous, time-dependent, fields in a spatially flat Friedmann-Robertson-Walker space time background. Given this ansatz, we find the following equations of motion for the complex superfields
$$\ddot{\mathrm{\Phi }}^i+3H\dot{\mathrm{\Phi }}^i+\mathrm{\Gamma }_{jk}^i\dot{\mathrm{\Phi }}^j\dot{\mathrm{\Phi }}^k+K^{i\overline{j}}_{\overline{j}}V=0,$$
(4)
where $`\mathrm{\Phi }^i`$ are the relevant moduli, $`\dot{\mathrm{\Phi }}^i=\mathrm{\Phi }^i/t`$, $`_{\overline{j}}V=V/\overline{\mathrm{\Phi }}^{\overline{j}}`$, and the connection on the Kähler manifold has the form
$$\mathrm{\Gamma }_{ij}^n=K^{n\overline{l}}\frac{K_{j\overline{l}}}{\mathrm{\Phi }^i}.$$
(5)
In addition, we obtain the Friedman equation for the Hubble factor $`H=\dot{a}/a`$, where $`a(t)`$ is the scale factor of the Universe,
$$3H^2=\kappa _P^2(\rho _\mathrm{\Phi }+\rho _b)=\kappa _P^2(K_{i\overline{j}}\dot{\mathrm{\Phi }}^i\dot{\mathrm{\Phi }}^{\overline{j}}+V+\rho _b),$$
(6)
with $`\rho _\mathrm{\Phi }`$ and $`\rho _b`$ are the energy density of the evolving moduli fields and background fluid respectively. The dynamics of the latter is given in terms of the scale factor and its background equation of state, $`\gamma 1p_b/\rho _b`$, where $`p_b`$ is the pressure of the fluid,
$$\rho _b=\rho _{b0}/a^{3\gamma }.$$
(7)
In what follows we set $`\kappa _P^2=1`$.
It is worth splitting the equations of motion for the complex chiral superfields into those for their real and imaginary parts
$$\ddot{\varphi }_R^i+3H\dot{\varphi }_R^i+\mathrm{\Gamma }_{jk}^i(\dot{\varphi }_R^j\dot{\varphi }_R^k\dot{\varphi }_I^j\dot{\varphi }_I^k)+\frac{1}{2}K^{i\overline{j}}_{j_R}V=0,$$
(8)
$$\ddot{\varphi }_I^i+3H\dot{\varphi }_I^i+\mathrm{\Gamma }_{jk}^i(\dot{\varphi }_I^j\dot{\varphi }_R^k+\dot{\varphi }_R^j\dot{\varphi }_I^k)+\frac{1}{2}K^{i\overline{j}}_{j_I}V=0,$$
(9)
where now $`\varphi _R^i`$ ($`\varphi _I^i`$) refers to the real (imaginary) part of the scalar fields and $`_{j_R}`$ ($`_{j_I}`$) are used to denote the derivative of the potential with respect to the real (imaginary) parts of the fields respectively.
## III KKLT model
The possibility of finding de Sitter vacua in string theory with a stabilized volume modulus, $`\sigma `$, was put forward in ref. Kachru:2003aw , and has been widely adopted in subsequent work. The key ingredient was to consider the combination of non perturbative effects and an additional flux term in the superpotential
$$W=W_0+A\mathrm{e}^{\alpha \sigma },$$
(10)
which, together with the usual Kähler potential
$$K=3\mathrm{l}\mathrm{n}(\sigma +\overline{\sigma }),$$
(11)
defines the F-part of the SUGRA potential, see eq. (3). It has been known for many years now that, in this context, it is possible to stabilize $`\sigma `$, although giving rise to an Anti de Sitter (AdS) vacuum. As pointed out in ref. Kachru:2003aw , if we include contributions from either anti-D3 or D7 branes, an additional D-type term is generated, of the form
$$V_D=\frac{C}{\sigma _r^3},$$
(12)
where we write $`\sigma =\sigma _r+i\sigma _i`$. By suitably tuning the value of $`C`$ one can move to a dS -or even Minkowski- vacuum.
In this section we are interested in studying the cosmological evolution of the field $`\sigma `$ as it rolls towards its minimum. Previous results addressing the same issue were published in ref. Brustein:2004jp , where only the evolution of the real part of $`\sigma `$, $`\sigma _r`$, was considered. Here we would like to extend this to study the behaviour of both $`\sigma _r`$ and $`\sigma _i`$.
The system of differential equations that one has to solve comes from eqs. (8,9) being applied to the present model, along with the evolution equation for the background fluid. Altogether, we have
$`\ddot{\sigma _r}+3H\dot{\sigma _r}{\displaystyle \frac{1}{\sigma _r}}(\dot{\sigma _r}^2\dot{\sigma _i}^2)+{\displaystyle \frac{2\sigma _r^2}{3}}_{\sigma _r}V`$ $`=`$ $`0,`$
$`\ddot{\sigma _i}+3H\dot{\sigma _i}{\displaystyle \frac{2}{\sigma _r}}\dot{\sigma _r}\dot{\sigma _i}+{\displaystyle \frac{2\sigma _r^2}{3}}_{\sigma _i}V`$ $`=`$ $`0,`$ (13)
$`\dot{\rho _b}+3H\gamma \rho _b`$ $`=`$ $`0,`$
subject to the Friedman constraint, see eq. (6),
$$3H^2=\frac{3}{4\sigma _r^2}(\dot{\sigma _r}^2+\dot{\sigma _i}^2)+V+\rho _b.$$
(14)
The scalar potential acquires a relatively simple form when written in terms of both real and imaginary parts of $`\sigma `$,
$`V`$ $`=`$ $`{\displaystyle \frac{\alpha A\mathrm{e}^{\alpha \sigma _r}}{2\sigma _r^2}}\left[A\left(1+{\displaystyle \frac{\alpha \sigma _r}{3}}\right)\mathrm{e}^{\alpha \sigma _r}+W_0\mathrm{cos}(\alpha \sigma _i)\right]`$ (15)
$`+`$ $`{\displaystyle \frac{C}{\sigma _r^3}}.`$
Given the above expression it is easy to see that the potential has an extremum in $`\sigma _i`$ for $`\alpha \sigma _i=n\pi `$, with $`n`$ an integer. Depending on the sign of $`W_0\mathrm{cos}(\alpha \sigma _i)`$ this can be either a maximum or a minimum. For example, in figure 1 we show
a contour plot of this scalar potential, $`V`$, in the ($`\varphi `$,$`\sigma _i`$) plane, for the same values used in ref. Brustein:2004jp , namely $`A=1.0`$, $`\alpha =0.1`$, $`C=3\times 10^{26}`$, and $`W_0`$ negative (with $`\mathrm{cos}(\alpha \sigma _i)=1`$) and such that the minimum at $`\sigma _i=0`$ is supersymmetric (see below). Following that same reference we work with the canonically normalized field $`\varphi =\sqrt{3/2}\mathrm{ln}\sigma _r`$, instead of $`\sigma _r`$ itself.
It is also worth mentioning the supersymmetric character of the minima shown in this plot. We have included the two supersymmetry-preserving conditions, $`\mathrm{Re}(F_\sigma )=\mathrm{Im}(F_\sigma )=0`$, and it can be clearly seen how both meet at the minima. Those correspond to $`\varphi =7.06`$ and even values of $`\alpha \sigma _i/\pi `$.
### III.1 One field evolution
In order to discuss the results of ref. Brustein:2004jp , we will first set the potential at a minimum in $`\sigma _i`$ and solve only the first and third equations of the system (13). This would correspond to taking the slice $`\alpha \sigma _i=0`$ in figure 1, and considering only the evolution along $`\varphi `$. As for the initial conditions, we set $`\dot{\sigma _r}_0=0`$, and we choose values for $`\sigma _{r0}`$ and the fractional energy density in the background fluid, $`\mathrm{\Omega }_b\kappa _P^2\rho _{b}^{}{}_{0}{}^{}/3H_0^2`$. We can then calculate the value of $`\rho _{b0}`$, using the relation
$$\rho _{b}^{}{}_{0}{}^{}=V(\sigma _{r0})\frac{\mathrm{\Omega }_b}{1\mathrm{\Omega }_b}.$$
(16)
In this section, the results are shown in terms of the initial abundance $`\mathrm{\Omega }_b`$ which will allow us to compare them with those in ref. Brustein:2004jp , where $`\mathrm{\Omega }_b=0.5`$.
In general one can identify up to five regions in the evolution of a scalar field with these types of potentials. In figure 2 we show a typical evolution going through the five regions although, of course, not all initial conditions will give rise to an evolution that will go through all of them.
First, if the energy density of the background dominates, the field is effectively frozen in its evolution. This can be seen in the plot before region 1, with an evolution similar to region 3. When the energy density in the background becomes comparable (region 1), the field starts rolling down its potential and eventually dominates the dynamics of the Universe. This happens when the potential is very shallow and is called the scalar field dominated solution. Eventually the potential becomes very steep, leading the evolution to a kinetically dominated solution (region 2), which ends when the frictional term in the equation of motion becomes dominant. The length of this period of ”kination” is related to the value of the ratio $`\rho _b/\rho _\varphi `$ at the end of region 1.
In region 3 the field is effectively frozen in the potential. The field restarts rolling down the potential once the background energy density has decayed to a value such that the frictional and potential terms in the equation of motion balance each other. At this stage the field evolves with a nearly constant ratio between kinetic and potential energy. This is called the scaling (or tracker) solution (region 4). Finally, if this ratio is sufficiently small the field does not possess enough kinetic energy to roll over the potential barrier and gets traped at the minimum (region 5). Details of all these phases of the evolution are explained in Appendix A.
In figure 3 we present the initial values of the variables ($`\varphi `$,$`\mathrm{\Omega }_r`$) for which there is stabilization of the field at the minimum of its potential in the presence of radiation<sup>1</sup><sup>1</sup>1We will use the subscript $`_r`$ when referring to a radiation background ($`\gamma =4/3`$), and the subscript $`_m`$ when referring to a matter background ($`\gamma =0`$)., i.e. $`\gamma =4/3`$. As expected, the smaller the initial background fraction (parametrized by $`\mathrm{\Omega }_r`$) is, the smaller the region of allowed initial values of $`\varphi `$ which lead to a late time stabilisation of the field at the minimum. The solid line along $`\mathrm{\Omega }_r=0.5`$ corresponds to the allowed region quoted in ref. Brustein:2004jp .
### III.2 Two field evolution
Let us now introduce $`\sigma _i`$, the imaginary part of the scalar field $`\sigma `$ in the evolution. The result of solving the full system given by eqs. (13) is presented in figure 4, where we show contour plots for the initial values of the pair ($`\varphi `$, $`\sigma _i`$), that lead to both fields ending up at their minimum. The different subplots show different contributions of the background (once again given by radiation).
A few comments are in order. As can be seen from figures 3, 4 and 5, there is both a minimum and a maximum value of $`\varphi `$ for which stabilisation at the minimum of the potential is successful.
Let us begin with the left bound on the value of $`\varphi `$. This comes from the type of potentials we are considering not being exactly an exponential potential. More precisely, they deviate most from an exponential shape for smaller values of $`\sigma _r`$. In these regions a scaling solution does not exist, and the evolution attractor is a scalar field dominated solution (the reasons for this are explained in more detail in Appendix A). Consider a typical evolution like the one in figure 2. Initially the field will have values that correspond to a scalar field dominated (region 1) evolution. Even if the background dominates initially (as in figure 2), the field will eventually take over. In other words, there is a saturation point for the initial value of $`\rho _b`$, above which all evolutions will be similar (they only differ in a time shift, corresponding to the time during which the field remains frozen initially). After some time the scalar field dominated solution becomes unstable and the field’s energy density becomes kinetically dominated (region 2 in figure 2). The background energy density will then dominate and the field will freeze at some constant value (region 3 in figure 2). Note that the larger the initial energy density for the field (that is, the smaller the initial value of $`\varphi `$), the longer it will take for the background to freeze the field. For very small initial values, the field will run past the minimum *before* being frozen by the background. The limiting case corresponds to the field freezing to a value nearly coincident with the position of the minimum (the case discussed in ref. Brustein:2004jp ). A matter ($`\gamma =1`$) type of background, having a slower redshift evolution, is much more effective at dominating the total energy density and freezing the field’s evolution. This can clearly be seen in a much smaller left bound on the initial values of $`\varphi `$ (see figure 5).
The bound arising on the maximum value for $`\varphi `$ is also related to the scaling solution. Naively one would think that, as the initial value approaches the minimum, the field would acquire less kinetic energy and it would be easier to stabilize its evolution in the minimum. The reason we still get a bound on the maximum allowed initial value of $`\varphi `$, is really due to the choice we made for the initial conditions.
If the initial energy density for the field is higher than the scaling one, the field will have to reach the scaling solution in a kinetic dominated regime, and it will be easy for it to overshoot the minimum if it starts too close to it. As can be seen from figure 2, for this type of potentials the scaling solution is essentially background dominated, so $`\mathrm{\Omega }_b`$ is close to unity. By choosing to fix the value of $`\mathrm{\Omega }_b<1`$ (which we did following ref. Brustein:2004jp ) we immediately force the initial conditions for the scalar field to be away from the scaling solution. Furthermore, since the potentials are not exact exponentials, the scaling value actually increases and goes assymptotically to 1 as we approach the minimum. Therefore, for most initial values of $`\mathrm{\Omega }_b`$, the field will start above the scaling solution, leading to a right bound on the initial values of $`\varphi `$. This bound will shift towards the minimum as we increase the initial $`\mathrm{\Omega }_b`$.
Of course, eventually, as $`\mathrm{\Omega }_b`$ approaches 1 we can get the right bound to coincide with the position of the minimum. This is much easier to see if we choose to fix the initial value of $`\rho _b`$, instead of $`\mathrm{\Omega }_b`$. In this case, the right bound always coincides with the position of the maximum in the potential (just before the “roll over” point). The results are shown in figure 6, where the right bound is always the same for changing values of the initial $`\rho _b`$.
Given the dependence of the scalar potential on $`\mathrm{cos}(\alpha \sigma _i)`$ these plots are symmetric under the inversion $`\sigma _i\sigma _i`$ (we only plot the positive $`\sigma _i`$ region). We can see that up to $`85\%`$ of the parameter space in the imaginary direction will lead to a stabilization of the field at the minimum.
## IV Kallosh-Linde model
This model generalizes the original version of the KKLT model by admiting one additional componet in the superpotential. More explicitly, we write
$$W=W_0+Ae^{\alpha \sigma }+Be^{\beta \sigma },$$
(17)
The particular example considered in ref. Kallosh:2004yh , sets parameters, $`A=1`$, $`B=1.03`$, $`C=0`$, $`\alpha =2\pi /100`$, $`\beta =2\pi /99`$ and $`W_0`$ such that there is a supersymmetric minimum with zero cosmological constant, i.e. $`W_0`$ is such that both $`W`$ and $`F_\sigma K_\sigma W+W_\sigma `$ vanish at some $`\sigma _r=\sigma _{\mathrm{crit}}`$, $`\sigma _i=0`$. Furthermore, there are a series of supersymmetric, AdS minima. All those are shown in figure 7,
where we plot the scalar potential for this model, in the ($`\sigma _r`$, $`\sigma _i`$) plane. The Minkowski solution is one of the two at $`\sigma _i=0`$, and corresponds to the smallest value of $`\sigma _r`$.
Unfortunately, this choice of parameters makes it almost impossible for the field to be stabilized if the initial value of the field is far from the minimum. In figure 8 we show the region of stabilization for a matter background, using a fixed $`\rho _m`$ for initial conditions. The reason for this is, again, related to the fact that the potential is very shallow for $`\varphi <\varphi _{\mathrm{min}}5.08`$, hence the scaling solution is never the dominant attractor.
Nonetheless, the mechanism of stabilising the field with a background fluid becomes impressive for models where the minumum lies within the range, in $`\varphi `$, of the scaling solution. One can see this by moving the minimum to larger values of $`\varphi `$, for example, by decreasing $`B`$ to $`B=1.5`$.
Following the approach of the previous section, in figures 9 and 10 we show the bounds on the initial position of the fields that lead to successful stabilisation. As in the examples of sections III.1 and III.2, the region of allowed initial conditions increases with the amount of initial energy density in the background. Also, as before, a smaller background equation of state will lead to better stabilization.
In the imaginary direction, this type of potential is much more difficult to stabilize than the ones used in the previous section. For these examples, we can only cover around $`0.6\%`$ of the available range of initial conditions. The range of values for $`\sigma _i`$ are similar in both cases, so this is only due to the fact that the examples in this section have a much larger period in the imaginary direction. This can be easily checked by comparing figures 7 and 1.
## V CONCLUSIONS
In this article we have studied moduli evolution in the context of the so-called KKLT-like models, that arise out of Type IIB string theory, where fluxes and the effect of anti-D3 or D7 branes have been taken into account. We have mainly looked at the allowed region of initial conditions of the moduli fields which leads to their eventual stabilisation at the supersymmetric minima of the scalar potential. Employing an approach first developed in Barreiro:1998aj , we found that a background perfect fluid has the ability to slow down the fields preventing them from running past the minimum, in agreement with Brustein:2004jp .
Moreover, we have extended the work of ref. Brustein:2004jp to include the dynamical evolution of the imaginary part of the moduli fields. We have evaluated the region of allowed initial conditions for the fields and confirmed that it increases for larger values of the initial background energy density $`\rho _0`$ and lower values of its equation of state. The effectiveness of this mechanism is also dependent on the existence of a scaling solution as the attractor. Regions in which the potential is considerably shallower than an exponential lead to a scalar field dominated attractor, working against successful stabilisation in the minimum.
In the second part of this paper we have turned our attention to analyse the Kallosh-Linde model of ref. Kallosh:2004yh . The conclusions are qualitatively equivalent but we note that only a small interval of initial conditions are allowed in the imaginary direction, making it more difficult to obtain stable moduli solutions in this case.
## Appendix A
In this appendix we summarise the key results presented in Barreiro:1998aj and Ng:2001hs , and describe how they can be applied to the discussion in the main text.
In the first of these papers Barreiro:1998aj , we considered the issue of stabilising the dilaton in the context of superstring cosmology without the explicit use of non-trivial three form fluxes, but including gaugino condensates and external sources such as radiation. Making use of the fact that in these models the scalar potentials arising out of gaugino condensates are generally exponential in nature, we demonstrated how scaling solutions associated with exponential potentials wetterich ; ferreira ; CLW could be slightly modified for these models, but basically had similar outcomes. In particular we showed how, due to the friction of the background expansion, the energy density in the scalar field could become, for a period, a fixed fraction of the background density , allowing the field to be trapped in the minimum of its potential, generally as it left its scaling regime. The key point was that previous scaling solutions available for pure exponential potentials were also approximately valid for quasi-exponential potentials. This was further reinforced in Ng:2001hs , where a phase space analysis of generic quasi-exponential potentials was carried out, looking for scaling solutions.
Consider a scalar field $`\varphi `$, with canonical kinetic terms, evolving in a FRW Universe containing a fluid with barotropic equation of state $`p_\gamma =(\gamma 1)\rho _b`$. It is useful to analyse the system in terms of the new variables $`x\frac{\dot{\varphi }}{\sqrt{6}H}`$ and $`y\frac{\sqrt{V}}{\sqrt{3}H}`$, as they allow for the determination of the scaling solutions. The evolution equations can then be written as CLW ; Ng:2001hs
$`H^{}`$ $`=`$ $`{\displaystyle \frac{3}{2}}H[2x^2+\gamma (1x^2y^2)]`$ (18)
$`x^{}`$ $`=`$ $`3x+\lambda \sqrt{{\displaystyle \frac{3}{2}}}y^2+{\displaystyle \frac{3}{2}}x[2x^2+\gamma (1x^2y^2)]`$ (19)
$`y^{}`$ $`=`$ $`\lambda \sqrt{{\displaystyle \frac{3}{2}}}xy+{\displaystyle \frac{3}{2}}y[2x^2+\gamma (1x^2y^2)]`$ (20)
$`\lambda ^{}`$ $`=`$ $`\sqrt{6}\lambda ^2(\mathrm{\Gamma }1)x,`$ (21)
where a prime denotes a derivative with respect to $`N\mathrm{ln}\frac{a}{a_{\mathrm{initial}}}`$, and we defined two new variables delaMacorra:1999ff ; Steinhardt:1999nw ,
$`\lambda {\displaystyle \frac{1}{\kappa _PV}}{\displaystyle \frac{dV}{d\varphi }},\mathrm{\Gamma }1{\displaystyle \frac{d}{d\varphi }}\left({\displaystyle \frac{1}{\kappa _P\lambda }}\right).`$ (22)
For the pure exponential potential, $`\lambda `$ is a constant and $`\mathrm{\Gamma }=1`$, so in a sense $`\mathrm{\Gamma }`$ is measuring how much our potential deviates from an exponential one. These equations can be solved easily in a number of regimes Barreiro:1998aj ; Ng:2001hs . The types of regimes available were already described in section III.1 and are shown in figure 2. Here we will present a more analytic description.
Let us start with the kinetic dominated regime, that is region 2 in figure 2. In this region, the evolution is dominated by the fields kinetic energy and so the potential energy, $`y`$, can be neglected in the equations of motion eqs. (1920). This leads to the solution Barreiro:1998aj
$$x=\left(1+\frac{1x_0^2}{x_0^2}e^{3(2\gamma )N}\right)^{1/2},$$
(23)
where $`x_0`$ is the initial condition for $`x`$ (taken at $`N=0`$, for simplicity). Notice that, since we are neglecting the potential terms ($`y`$), this solution is really potential independent and is valid as long as we are in a kinetic dominated regime. Clearly the kinetic energy of the field which is proportional to $`x^2`$ is decreasing rapidly with time. The solution for $`\varphi `$ follows and is given by
$`\varphi _I(N)`$ $`=`$ $`\varphi _0+{\displaystyle \frac{2\sqrt{6}}{3(2\gamma )}}[\mathrm{sinh}^1\left({\displaystyle \frac{x_0}{\sqrt{1x_0^2}}}\right)`$ (24)
$``$ $`\mathrm{sinh}^1\left({\displaystyle \frac{x_0}{\sqrt{1x_0^2}}}e^{3(2\gamma )N/2}\right)],`$
where $`\varphi _0`$ is the initial value of the field. At early times the solution corresponds precisely to that of eq. (8) of Brustein:2004jp , where the energy density is dominated by the kinetic energy of the scalar field. As time carries on though, and the field slows down, it eventually tends to a constant solution as can be seen in eq. (24). This constant value corresponds to region 3 of figure 2.
Let us now turn our attention to the scaling evolution, that is region 4 in figure 2. In Ng:2001hs an approximate equation of state for the field evolution in the scaling regime was derived. This is given by
$`\gamma _\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\gamma +(2\mathrm{\Gamma }1)\lambda ^2/3\right]`$
$``$ $`{\displaystyle \frac{1}{2}}\sqrt{\left[\gamma +(2\mathrm{\Gamma }1)\lambda ^2/3\right]^2+8\gamma (\mathrm{\Gamma }1)\lambda ^2/3}.`$
When $`\mathrm{\Gamma }1`$ (the exponential case), $`\gamma _\varphi `$ tends to $`\gamma `$, that is one recovers the usual result that the field evolution mimics the background. Moreover, this scaling solution is only stable if it satisfies the following conditions,
$`\gamma _\varphi `$ $`<`$ $`\sqrt{{\displaystyle \frac{\gamma \lambda ^2}{3}}},\mathrm{and}`$
$`\gamma _\varphi `$ $`<`$ $`{\displaystyle \frac{3\gamma 2\lambda ^2}{6}}\left[1\sqrt{1+{\displaystyle \frac{12\lambda ^2(2+\gamma )}{(3\gamma 2\lambda ^2)^2}}}\right].`$ (26)
In the particular case of double exponential potentials, $`V(\varphi )=V_0\mathrm{exp}(ae^{b\varphi })`$, it is possible to obtain an approximate solution for the field evolution. The most useful way to present it is in a recursive form Barreiro:1998aj ; Ng:2001hs
$`S_0`$ $`=`$ $`{\displaystyle \frac{1}{a}}\mathrm{ln}\left[{\displaystyle \frac{3}{2}}{\displaystyle \frac{\rho _0}{V_0}}{\displaystyle \frac{\gamma (2\gamma )}{a^2b^2}}\right],`$
$`S_1`$ $`=`$ $`S_0{\displaystyle \frac{3\gamma }{a}}N,`$ (27)
$`S_k`$ $`=`$ $`S_0{\displaystyle \frac{2}{a}}\mathrm{ln}(S_{k1}){\displaystyle \frac{3\gamma }{a}}N,`$
where $`S\mathrm{exp}(b\varphi )`$. It is easy to check that in the regions where the potentials discussed in sections III and IV are close to an exponential in $`\sigma _r`$, these solutions are a very good fit to the scaling evolution with $`S\sigma _r`$.
Finally, there is a third possible type of evolution, shown in region 1 of figure 2, the scalar field dominated solution. Its effective equation of state was shown in Ng:2001hs to be
$$\gamma _\varphi =\lambda _\varphi ^2/3,$$
(28)
where we have defined
$$\lambda _\varphi =\frac{3}{2}\left[\frac{1\sqrt{14(\mathrm{\Gamma }1)\lambda ^2/3}}{(\mathrm{\Gamma }1)\lambda }\right]$$
(29)
for $`\mathrm{\Gamma }1`$, and $`\lambda _\varphi =\lambda `$ otherwise. Again, a condition for the stability of the solution was derived, and is given by
$`6(\lambda _\varphi \lambda ){\displaystyle \frac{1}{\lambda _\varphi }}+{\displaystyle \frac{1}{2}}(\lambda _\varphi ^2+\lambda \lambda _\varphi 6\gamma )`$ $`<`$ $`0,`$ (30)
$`36{\displaystyle \frac{\lambda }{\lambda _\varphi }}+{\displaystyle \frac{1}{2}}\lambda _\varphi (3\lambda _\varphi 2\lambda )`$ $`<`$ $`0.`$ (31)
We can now apply this knowledge to the potentials we are dealing with in this paper. In figure 11 we compare the evolution of the equation of state resulting from a numerical integration with the scaling and scalar field dominated solutions given above, in their region of stability. The model used is the KKLT example presented in section III, also used to produce figure 2. We can clearly see from figure 11 that in this specific example the attractor solution only makes sense for $`\varphi >3.7`$. Therefore, if the field starts its evolution before this value, it is either in the scalar field dominated solution or in a kinetic energy dominated solution. This makes the stabilisation of the field more difficult to achieve for smaller initial values of the field.
###### Acknowledgements.
The authors wish to thank Ramy Brustein for useful discussions, and Nemanja Kaloper for pointing out the relevance of ref. Kaloper:1991mq to our work. EJC would like to acknowledge the Aspen Center for Physics for their support when part of this project was started. TB is supported by FCT grant SFRH/BPD/3512/2000. BdC is supported by PPARC. NJN is supported by the Department of Energy under contract DE-FG02-94ER40823 at the University of Minnesota.
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# Measurement of Time-Dependent 𝑪𝑷 Asymmetries and the 𝑪𝑷-Odd Fraction in the Decay 𝑩^𝟎→𝑫^{∗+}𝑫^{∗-}
BABAR-PUB-05/024
SLAC-PUB-11321, hep-ex/0506082
thanks: Deceased
The BABAR Collaboration
## Abstract
We present an updated measurement of time-dependent $`CP`$ asymmetries and the $`CP`$-odd fraction in the decay $`B^0D^+D^{}`$ using $`232\times 10^6B\overline{B}`$ pairs collected by the BABAR detector at the PEP-II $`B`$ factory. We determine the $`CP`$-odd fraction to be $`0.125\pm 0.044(\mathrm{stat})\pm 0.007(\mathrm{syst})`$. The time-dependent $`CP`$ asymmetry parameters $`C_+`$ and $`S_+`$ are determined to be $`0.06\pm 0.17(\mathrm{stat})\pm 0.03(\mathrm{syst})`$ and $`0.75\pm 0.25(\mathrm{stat})\pm 0.03(\mathrm{syst})`$, respectively. The Standard Model predicts these parameters to be 0 and $`\mathrm{sin}2\beta `$, respectively, in the absence of penguin amplitude contributions.
The time-dependent $`CP`$ asymmetry measurement in $`B^0D^+D^{}`$ decay provides an important test of the Standard Model (SM). In the SM, $`CP`$ violation arises from a complex phase in the Cabibbo-Kobayashi-Maskawa (CKM) quark-mixing matrix CKM . Measurements of $`CP`$ asymmetries by the BABAR Aubert:2002rg and BELLE Abe:2002px collaborations have firmly established this effect in the $`B^0J/\psi K_S^0`$ decay conjugate and related modes that are governed by the $`bc\overline{c}s`$ transition. The $`B^0D^+D^{}`$ decay is dominated by the $`bc\overline{c}d`$ transition. Within the framework of the SM, the $`CP`$ asymmetry of $`B^0D^+D^{}`$ is related to $`\mathrm{sin}2\beta `$ when the correction due to penguin diagram contributions are neglected. The penguin-induced correction has been estimated in models based on the factorization approximation and heavy quark symmetry and was predicted to be about 2% Pham:1999fy . A significant deviation of the measured $`\mathrm{sin}2\beta `$ from the one observed in $`bc\overline{c}s`$ decays would be evidence for a new $`CP`$-violating interaction. The enhanced sensitivity of $`B^0D^+D^{}`$ to such a process arises from its much smaller SM amplitude compared with that of the $`bc\overline{c}s`$ transition.
The $`B^0D^+D^{}`$ decay proceeds through the $`CP`$-even $`S`$ and $`D`$ waves and through the $`CP`$-odd $`P`$ wave. In this Letter, we present an improved measurement of the $`CP`$-odd fraction Aubert:2003uv ; Miyake:2005qb $`R_{}`$ based on a time-integrated one-dimensional angular analysis. We also present an improved measurement of the time-dependent $`CP`$ asymmetry Aubert:2003uv ; Miyake:2005qb , obtained from a combined analysis of time-dependent flavor-tagged decays and the one-dimensional angular distribution of the decay products.
The data used in this analysis comprise 232 million $`\mathrm{{\rm Y}}(4S)B\overline{B}`$ decays collected by the BABAR detector at the PEP-II storage ring. The BABAR detector is described in detail elsewhere Aubert:2001tu . We use a Monte Carlo (MC) simulation based on GEANT4 Agostinelli:2002hh to validate the analysis procedure and to study the relevant backgrounds.
We select $`B^0D^+D^{}`$ decay by combining two charged $`D^{}`$ candidates reconstructed in the modes $`D^+D^0\pi ^+`$ and $`D^+D^+\pi ^0`$. We include the $`D^+D^{}`$ combinations $`(D^0\pi ^+,\overline{D}{}_{}{}^{0}\pi _{}^{})`$ and $`(D^0\pi ^+,D^{}\pi ^0)`$, but not $`(D^+\pi ^0,D^{}\pi ^0)`$ because of the smaller branching fraction and larger backgrounds. To suppress the $`e^+e^{}q\overline{q}(q=u,d,s,\mathrm{and}\mathrm{c})`$ continuum background, we require the ratio of the second and zeroth order Fox-Wolfram moments Fox:1978vu to be less than 0.6.
Candidates for $`D^0`$ and $`D^+`$ mesons are reconstructed in the modes $`D^0K^{}\pi ^+`$, $`K^{}\pi ^+\pi ^0`$, $`K^{}\pi ^+\pi ^+\pi ^{}`$, $`K_S^0\pi ^+\pi ^{}`$ and $`D^+K^{}\pi ^+\pi ^+`$, $`K_S^0\pi ^+`$, $`K^{}K^+\pi ^+`$. The reconstructed mass of the $`D^0`$ ($`D^+`$) candidate is required to be within 20 $`\mathrm{MeV}/c^2`$ of its nominal mass Eidelman:2004wy , except for the $`D^0K^{}\pi ^+\pi ^0`$ candidate, where a looser requirement of 40 $`\mathrm{MeV}/c^2`$ is applied.
The $`K_S^0`$ candidates are reconstructed from two oppositely-charged tracks with an invariant mass within 20 $`\mathrm{MeV}/c^2`$ of the nominal $`K_S^0`$ mass. The $`\chi ^2`$ probability of the $`\pi ^+\pi ^{}`$ vertex fit must be greater than $`0.1\%`$. Charged kaon candidates are required to be inconsistent with the pion hypothesis, as inferred from the Cherenkov angle measured by the Cherenkov detector and the ionization energy loss measured by the charged-particle tracking system Aubert:2001tu . Neutral pion candidates are formed from two photons detected in the electromagnetic calorimeter Aubert:2001tu , each with energy above 30 $`\mathrm{MeV}`$. The mass of the pair must be within 30 $`\mathrm{MeV}/c^2`$ of the nominal $`\pi ^0`$ mass, and their summed energy is required to be greater than 200 $`\mathrm{MeV}`$. In addition, a mass-constrained fit is applied to the $`\pi ^0`$ candidates for further analysis.
The $`D^0`$ and $`D^+`$ candidates are subject to a mass-constrained fit prior to the formation of the $`D^+`$ candidates. A slow $`\pi ^+`$ from $`D^+`$ decay is required to have a momentum in the $`\mathrm{{\rm Y}}(4S)`$ center-of-mass (CM) frame less than 450 $`\mathrm{MeV}/c`$. A slow $`\pi ^0`$ from $`D^+`$ must have a momentum between $`70`$ and $`450\mathrm{MeV}/c`$ in the CM frame. No requirement on the photon-energy sum is applied to the $`\pi ^0`$ candidates from the $`D^+`$ decays.
For each $`B^0D^+D^{}`$ candidate, we construct a likelihood function Aubert:2002vn $`_{\mathrm{mass}}`$ from the masses and mass uncertainties of the $`D`$ and $`D^{}`$ candidates. The likelihood $`_{\mathrm{mass}}`$ is calculated as the product of the likelihoods for the $`D`$ and $`D^{}`$ candidates. The $`D`$ mass resolution is modeled by a Gaussian whose variance is determined on a candidate-by-candidate basis. The $`D^{}`$-$`D`$ mass difference resolution is modeled by a double-Gaussian distribution whose parameters are determined from simulated events. The values of $`_{\mathrm{mass}}`$ and the difference of the $`B^0`$ candidate energy $`E_B`$ from the beam energy $`E_{\mathrm{Beam}}`$, $`\mathrm{\Delta }EE_BE_{\mathrm{Beam}}`$, in the $`\mathrm{{\rm Y}}(4S)`$ CM frame are used to reduce the combinatoric background further. From the simulated events, the maximum allowed values of $`\mathrm{ln}_{\mathrm{mass}}`$ and $`|\mathrm{\Delta }E|`$ are optimized for each individual final state to obtain the highest expected signal significance using the previously measured $`B^0D^+D^{}`$ branching fraction Aubert:2003uv .
The energy-substituted mass, $`m_{\mathrm{ES}}\sqrt{E_{\mathrm{Beam}}^2p_B^2}`$, where $`p_B^{}`$ is the $`B^0`$ candidate momentum in the $`\mathrm{{\rm Y}}(4S)`$ CM frame, is used to extract the signal yield from the events satisfying the aforementioned selection. We select the $`B^0`$ candidates that have $`m_{\mathrm{ES}}5.23\mathrm{GeV}/c^2`$. In cases where more than one $`B^0`$ candidate is reconstructed in an event, the candidate with the smallest value of $`\mathrm{ln}_{\mathrm{mass}}`$ is chosen. A fit to the $`m_{\mathrm{ES}}`$ distribution with a probability density function (PDF) given by the sum of a Gaussian shape for the signal and an ARGUS Albrecht:1990cs function for the background yields $`391\pm 28(\mathrm{stat})`$ signal events. In the region of $`m_{\mathrm{ES}}>5.27\mathrm{GeV}/c^2`$, the signal purity is approximately 70 %.
In the transversity basis Dunietz:1990cj , we define the following three angles: the angle $`\theta _1`$ between the momentum of the slow pion from the $`D^{}`$ and the opposite direction of flight of the $`D^+`$ in the $`D^{}`$ rest frame; the polar angle $`\theta _{\mathrm{tr}}`$ and azimuthal angle $`\varphi _{\mathrm{tr}}`$ of the slow pion from the $`D^+`$ defined in the $`D^+`$ rest frame, where the opposite direction of flight of the $`D^{}`$ is chosen as the $`x`$-axis, and the $`z`$-axis is defined as the normal to the $`D^{}`$ decay plane.
The time-dependent angular distribution of the decay products is given in Ref. penguin . Taking into account the detector angular acceptance efficiency and integrating over the decay time and the angles $`\theta _1`$ and $`\varphi _{\mathrm{tr}}`$, we obtain a one-dimensional differential decay rate:
$`{\displaystyle \frac{1}{\mathrm{\Gamma }}}{\displaystyle \frac{d\mathrm{\Gamma }}{d\mathrm{cos}\theta _{\mathrm{tr}}}}`$ $`=`$ $`{\displaystyle \frac{9}{32\pi }}[(1R_{})\mathrm{sin}^2\theta _{\mathrm{tr}}`$ (1)
$`\times `$ $`\left\{{\displaystyle \frac{1+\alpha }{2}}I_0(\mathrm{cos}\theta _{\mathrm{tr}})+{\displaystyle \frac{1\alpha }{2}}I_{}(\mathrm{cos}\theta _{\mathrm{tr}})\right\}`$
$`+`$ $`2R_{}\mathrm{cos}^2\theta _{\mathrm{tr}}\times I_{}(\mathrm{cos}\theta _{\mathrm{tr}})],`$
where $`R_{}=|A_{}|^2/(|A_0|^2+|A_{}|^2+|A_{}|^2)`$, $`\alpha =(|A_0|^2|A_{}|^2)/(|A_0|^2+|A_{}|^2)`$, $`A_0`$ is the amplitude for longitudinally polarized $`D^{}`$ mesons, $`A_{}`$ and $`A_{}`$ are the amplitudes for parallel and perpendicular transversely polarized $`D^{}`$ mesons. The three efficiency moments, $`I_k`$ $`(k=0,,)`$, are defined as
$`I_k(\mathrm{cos}\theta _{\mathrm{tr}})={\displaystyle d\mathrm{cos}\theta _1d\varphi _{\mathrm{tr}}g_k(\theta _1,\varphi _{\mathrm{tr}})\epsilon (\theta _1,\theta _{\mathrm{tr}},\varphi _{\mathrm{tr}})},`$ (2)
where $`g_0=4\mathrm{cos}^2\theta _1\mathrm{cos}^2\varphi _{\mathrm{tr}}`$, $`g_{||}=2\mathrm{sin}^2\theta _1\mathrm{sin}^2\varphi _{\mathrm{tr}}`$, $`g_{}=\mathrm{sin}^2\theta _1`$, and $`\epsilon `$ is the detector efficiency. The efficiency moments are parameterized as second-order even polynomials of $`\mathrm{cos}\theta _{\mathrm{tr}}`$. Their parameter values are determined from the MC and are subsequently fixed in the likelihood fit to the differential decay distribution of $`\mathrm{cos}\theta _{\mathrm{tr}}`$. In fact, the three $`I_k`$ functions deviate only slightly from a constant, making the distribution, Eq. 1, nearly independent of the amplitude ratio $`\alpha `$.
The $`CP`$-odd fraction $`R_{}`$ is measured in a simultaneous unbinned maximum likelihood fit to the $`\mathrm{cos}\theta _{\mathrm{tr}}`$ and the $`m_{\mathrm{ES}}`$ distribution. The background shape is modeled as an even second-order polynomial in $`\mathrm{cos}\theta _{\mathrm{tr}}`$, while the signal PDF is given by Eq. 1. The finite detector resolution of the $`\theta _{\mathrm{tr}}`$ measurement is modeled as a double Gaussian plus a small tail component that accounts for misreconstructed events. The parameterization of the $`\theta _{\mathrm{tr}}`$ resolution function is fixed from the MC simulation and subsequently used to convolve the signal PDF in the maximum likelihood fit. Since the angle $`\theta _{\mathrm{tr}}`$ is calculated with the slow pion from the $`D^+`$, we categorize events into three types: $`D^+D^{}(D^0\pi ^+,\overline{D}{}_{}{}^{0}\pi _{}^{})`$, $`(D^0\pi ^+,D^{}\pi ^0)`$, and $`(D^+\pi ^0,\overline{D}{}_{}{}^{0}\pi _{}^{})`$, each with different signal-fraction parameters in the likelihood fit. Their angular efficiency moments and $`\mathrm{cos}\theta _{\mathrm{tr}}`$ resolutions are also separately determined from the MC simulation. The other parameters determined in the likelihood fit are the $`\mathrm{cos}\theta _{\mathrm{tr}}`$ background-shape parameter, three $`m_{\mathrm{ES}}`$ parameters ($`\sigma `$ and mean of the signal Gaussian, and the ARGUS shape parameter $`\kappa `$), as well as $`R_{}`$. The fit to the data yields
$$R_{}=0.125\pm 0.044(\mathrm{stat})\pm 0.007(\mathrm{syst}).$$
(3)
The projections of the fitted result onto $`m_{\mathrm{ES}}`$ and $`\mathrm{cos}\theta _{\mathrm{tr}}`$ are shown in Fig. 1.
In the fit described above, the value of $`\alpha `$ is fixed to zero. We estimate the corresponding systematic uncertainty by varying its value from $`1`$ to $`+1`$ and find negligible change (less than 0.002) in the fitted value of $`R_{}`$. Other systematic uncertainties arise from the parameterization of the angular resolution, the determination of the efficiency moments, and the background parameterization. The total systematic uncertainty on $`R_{}`$ is 0.007, significantly smaller than the statistical error.
We subsequently perform a combined analysis of the $`\mathrm{cos}\theta _{\mathrm{tr}}`$ distribution and the time dependence to extract the time-dependent $`CP`$ asymmetry, using the event sample described previously. We use information from the other $`B`$ meson in the event to tag the initial flavor of the fully reconstructed $`B^0D^+D^{}`$ candidate.
The decay rate $`f_+(f_{})`$ for a neutral $`B`$ meson accompanied by a $`B^0(\overline{B}{}_{}{}^{0})`$ tag is given by
$`f_\pm (\mathrm{\Delta }t,\mathrm{cos}\theta _{\mathrm{tr}})\mathrm{e}^{|\mathrm{\Delta }t|/\tau _{B^0}}\{G(1\mathrm{\Delta }\omega )(12\omega )`$
$`[F\mathrm{sin}(\mathrm{\Delta }m_d\mathrm{\Delta }t)+H\mathrm{cos}(\mathrm{\Delta }m_d\mathrm{\Delta }t)]\},`$ (4)
where $`\mathrm{\Delta }t=t_{\mathrm{rec}}t_{\mathrm{tag}}`$ is the difference between the proper decay time of the reconstructed $`B`$ meson ($`B_{\mathrm{rec}}`$) and that of the tagging $`B`$ meson ($`B_{\mathrm{tag}}`$), $`\tau _{B^0}`$ is the $`B^0`$ lifetime, and $`\mathrm{\Delta }m_d`$ is the mass difference determined from the $`B^0`$-$`\overline{B}^0`$ oscillation frequency Eidelman:2004wy . The average mistag probability $`\omega `$ describes the effect of incorrect tags, and $`\mathrm{\Delta }\omega `$ is the difference between the mistag rate for $`B^0`$ and $`\overline{B}^0`$. The $`G`$, $`F`$ and $`H`$ coefficients are defined as:
$`G`$ $`=`$ $`(1R_{})\mathrm{sin}^2\theta _{\mathrm{tr}}+2R_{}\mathrm{cos}^2\theta _{\mathrm{tr}},`$
$`F`$ $`=`$ $`(1R_{})S_+\mathrm{sin}^2\theta _{\mathrm{tr}}2R_{}S_{}\mathrm{cos}^2\theta _{\mathrm{tr}},`$ (5)
$`H`$ $`=`$ $`(1R_{})C_+\mathrm{sin}^2\theta _{\mathrm{tr}}+2R_{}C_{}\mathrm{cos}^2\theta _{\mathrm{tr}},`$
where we allow the three transversity amplitudes to have different $`\lambda _k=(q/p)(\overline{A}_k/A_k)`$ $`(k=0,,)`$ penguin due to possibly different penguin-to-tree amplitude ratios, and define the $`CP`$ asymmetry $`C_k=1|\lambda _k|^2/1+|\lambda _k|^2`$, $`S_k=2\mathrm{}(\lambda _k)/1+|\lambda _k|^2`$. Here we also have:
$$C_+=\frac{C_{}|A_{}|^2+C_0|A_0|^2}{|A_{}|^2+|A_0|^2},S_+=\frac{S_{}|A_{}|^2+S_0|A_0|^2}{|A_{}|^2+|A_0|^2}.$$
(6)
In the absence of penguin contributions, we expect that $`C_0=C_{}=C_{}=0`$, and $`S_0=S_{}=S_{}=\mathrm{sin}2\beta `$.
In Eq. 4, the small angular acceptance effects are not incorporated, but absorbed into the ”effective” value of $`R_{}`$, which is left free to vary in the final fit. No bias is seen in the resulting values of $`C_+`$, $`C_{}`$, $`S_+`$, and $`S_{}`$ in MC simulation.
The technique used to measure the $`CP`$ asymmetry is analogous to previous BABAR measurements as described in Ref. Aubert:2004zt . Only events with a $`\mathrm{\Delta }t`$ uncertainty less than $`2.5\text{ps}`$ and a measured $`|\mathrm{\Delta }t|`$ less than $`20\text{ps}`$ are accepted. We performed a simultaneous unbinned maximum likelihood fit to the $`\mathrm{cos}\theta _{\mathrm{tr}}`$, $`\mathrm{\Delta }t`$, and $`m_{\mathrm{ES}}`$ distributions to extract the $`CP`$ asymmetry. The signal PDF in $`\theta _{\mathrm{tr}}`$ and $`\mathrm{\Delta }t`$ is given by Eq. 4. The signal mistag probability is determined from a sample of neutral $`B`$ decays to flavor eigenstates, $`B_{\mathrm{flav}}`$. In the likelihood fit, the expression in Eq. 4 is convolved with an empirical $`\mathrm{\Delta }t`$ resolution function determined from the $`B_{\mathrm{flav}}`$ sample. The $`\theta _{\mathrm{tr}}`$ resolution is accounted for in the same way as described previously.
The background $`\mathrm{\Delta }t`$ distributions are parameterized with an empirical description that includes prompt and non-prompt components. We allow the non-prompt background to have two free parameters, $`C_{\mathrm{eff}}`$ and $`S_{\mathrm{eff}}`$, the effective $`CP`$ asymmetries, in the likelihood fit. The background shape in $`\theta _{\mathrm{tr}}`$ is modeled as an even second-order polynomial in $`\mathrm{cos}\theta _{\mathrm{tr}}`$, much as it is in the time-integrated angular analysis.
The fit to the data yields
$`C_+`$ $`=`$ $`0.06\pm 0.17(\mathrm{stat})\pm 0.03(\mathrm{syst}),`$
$`C_{}`$ $`=`$ $`0.20\pm 0.96(\mathrm{stat})\pm 0.11(\mathrm{syst}),`$
$`S_+`$ $`=`$ $`0.75\pm 0.25(\mathrm{stat})\pm 0.03(\mathrm{syst}),`$
$`S_{}`$ $`=`$ $`1.75\pm 1.78(\mathrm{stat})\pm 0.22(\mathrm{syst}).`$ (7)
Fig. 2 shows the $`\mathrm{\Delta }t`$ distributions and asymmetries in yields between $`B^0`$ and $`\overline{B}^0`$ tags, overlaid with the projection of the likelihood fit result. Because the $`CP`$-odd fraction is small, we have rather large statistical uncertainties for the measured $`C_{}`$ and $`S_{}`$ values. For comparison, we repeat the fit with the assumption that both $`CP`$-even and $`CP`$-odd states have the same $`CP`$ asymmetry. We find that $`C_+=C_{}=0.03\pm 0.13(\mathrm{stat})\pm 0.02(\mathrm{syst})`$, and $`S_+=S_{}=0.69\pm 0.23(\mathrm{stat})\pm 0.03(\mathrm{syst})`$. In both cases, the effective $`CP`$ asymmetries in the background are found to be consistent with zero within the statistical uncertainties.
The systematic uncertainties on $`C_+`$, $`C_{}`$, $`S_+`$ and $`S_{}`$ arise from the amount of possible backgrounds that tend to peak under the signal and their $`CP`$ asymmetry, the assumed parameterization of the $`\mathrm{\Delta }t`$ resolution function, the possible differences between the $`B_{\mathrm{flav}}`$ and $`B_{CP}`$ mistag fractions, knowledge of the event-by-event beam-spot position, and the possible interference between the suppressed $`\overline{b}\overline{u}c\overline{d}`$ amplitude and the favored $`bc\overline{u}d`$ amplitude for some tag-side decays Long:2003wq . It also includes the systematic uncertainties from the finite MC sample used to verify the fitting method. In general, all of the systematic uncertainties are found to be much smaller than the statistical uncertainties.
In summary, we have reported measurements of the $`CP`$-odd fraction and time-dependent $`CP`$ asymmetries for the decay $`B^0D^+D^{}`$. The measurement supersedes the previous BABAR result Aubert:2003uv , with more than $`50\%`$ reduction in the statistical uncertainty, and indicates that $`B^0D^+D^{}`$ is mostly $`CP`$-even. The time-dependent asymmetries are found to be consistent with the SM predictions within the statistical uncertainty.
We are grateful for the excellent luminosity and machine conditions provided by our PEP-II colleagues, and for the substantial dedicated effort from the computing organizations that support BABAR. The collaborating institutions wish to thank SLAC for its support and kind hospitality. This work is supported by DOE and NSF (USA), NSERC (Canada), IHEP (China), CEA and CNRS-IN2P3 (France), BMBF and DFG (Germany), INFN (Italy), FOM (The Netherlands), NFR (Norway), MIST (Russia), and PPARC (United Kingdom). Individuals have received support from CONACyT (Mexico), A. P. Sloan Foundation, Research Corporation, and Alexander von Humboldt Foundation.
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# Electronic Mach-Zehnder interferometer as a tool to probe fractional statistics
## I Introduction
One of the key features of the quantum Hall effect (QHE) is the fractional charge and statistics of quasiparticles. The seminal shot noise experiments of the Weizmann and Saclay groups weizmann ; Grenoble allowed for direct observations of fractional charges. A recent experiment on the mutual fractional statistics of two quasiparticles with different charges has been published in Ref. goldman2005, . It involved a setup consisting of an island of a fractional quantum Hall liquid embedded in a liquid with a different filling factor JKT . Theoretically, several approaches cfksw ; HBT ; kane have been proposed as possible tools to probe mutual statistics of identical quasiparticles. However, to this date, no experimental verification of the statistics of identical quasiparticles has been reported.
In this paper we propose a different approach to observing the statistics of identical quasiparticles. Our approach employs the electronic analogue of the Mach-Zehnder interferometer, recently designed at the Weizmann Institute MZ . This device has been used to observe the Aharonov-Bohm effect in the integer quantum Hall regime MZ . As is shown below, in higher magnetic fields this type of device would allow the observation of fractional statistics.
The proposed method of the observation of fractional statistics has a number of advantages in comparison with other set-ups. In Refs. HBT, ; kane, one needs to measure the noise or the current-current correlation function. In our approach it is sufficient to find the current through the interferometer. This resembles Ref. cfksw, . However, in Ref. cfksw, one has to control the number of quasiparticles trapped in the interferometer in order to probe fractional statistics. There is no such difficulty in our case. Besides, the interference pattern can be observed in the Mach-Zehnder interferometer at a larger interferometer size than in the standard geometry cfksw .
An electronic Mach-Zehnder interferometer is sketched in Fig. 1. Charge propagates along two quantum Hall edges and tunnels between the edges at the point contacts QPC1 and QPC2. Figs. 1a) and 1b) depict two possible setups: a) tunneling takes place between two different fractional quantum Hall puddles; and b) tunneling is between the edges of a single puddle. In the latter case, fractionally charged quasiparticles can tunnel at QPC1 and QPC2. In the former case, only electrons are allowed to tunnel wen . Pinching off the edges closer to each other, one can deform the system depicted in Fig. 1b) into the configuration Fig. 1a). The tunneling current between the edges depends on their voltage difference and the magnetic flux through the region A-QPC1-B-QPC2-A. As shown below, the current includes a flux-independent contribution $`I_0`$ and a contribution $`I_\mathrm{\Phi }`$ which oscillates as function of the magnetic flux with period $`\mathrm{\Phi }_0=hc/e`$. We will calculate $`I_\mathrm{\Phi }`$ and $`I_0`$ as functions of tunneling amplitudes $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ at QPC1 and QPC2. If no tunneling occurs at QPC2, i.e. $`\mathrm{\Gamma }_2=0`$, then the total current $`I(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)=I_0(\mathrm{\Gamma }_1,0)`$ is independent of the magnetic flux. At weak tunneling at QPC2, i.e. small $`\mathrm{\Gamma }_2\mathrm{\Gamma }_1`$, the flux-dependent component of the current scales as $`I_\mathrm{\Phi }(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)[I_0(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)I_0(\mathrm{\Gamma }_1,0)]^b`$, where the exponent $`b`$ depends on the statistics of tunneling particles. In the case of electron tunneling, $`b=1/2`$. In the case of the tunneling of fractional quasiparticles, the exponent depends on their statistics and is always greater than for electrons.
The outline of the paper is the following. First, we give a qualitative explanation of our results. Next, we introduce the model which is used in our calculations. Then, we calculate the $`IV`$ curves, respectively, for quasiparticle tunneling in the fractional QHE regime with filling factors $`\nu =1/(2m+1)`$, and electron tunneling between quantum Hall liquids with filling factors $`\nu _{1,2}=1/(2m_{1,2}+1)`$. In the latter case our calculations follow the standard route cfksw ; geller . In the former case one has to carefully treat the Klein factors describing quasiparticle statistics. The Appendices discuss some technical details of the perturbation theory employed.
## II Qualitative discussion
We study charge transport through the Mach-Zehnder interferometer in the limit of weak tunneling at QPC1 and QPC2. Thus, we assume that the edges of the fractional quantum Hall liquid are far from each other in comparison with the magnetic length $`l_0`$. Hence, $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2\mathrm{}/\tau _c`$, where the ultra-violet cut-off scale $`\tau _cl_0/v`$ and $`v`$ is the excitation velocity along the edges. The tunneling amplitudes can be controlled by gate voltages. We assume that the gate voltages are chosen such that one of the situations shown in Figs. 1a) and 1b) takes place. Other situations, e.g. weak quasiparticle tunneling at QPC1 and weak electron tunneling at QPC2, are also possible but will not be considered here.
We first address the electron tunneling case depicted in Fig. 1a). There are three contributions to the current $`I=I_1+I_2+I_{12}`$. The first contribution is due to the tunneling at QPC1, $`I_1=c_1|\mathrm{\Gamma }_1|^2`$, where $`c_1`$ depends on the temperature and voltage and is independent of the magnetic flux through the interferometer. The second contribution is due to the second quantum point contact, $`I_2=c_1|\mathrm{\Gamma }_2|^2`$. Finally, a contribution arises due to quantum interference between the electrons which tunnel at QPC1 and QPC2. This contribution equals $`I_{12}=c_2\mathrm{\Gamma }_1\mathrm{\Gamma }_2^{}\mathrm{exp}(i\varphi )+c.c.`$, where $`c_2`$ depends on the voltage, temperature and the distance between the point contacts, and $`\varphi `$ is the difference of the Aharonov-Bohm phases picked up by electrons propagating between QPC1 and QPC2 along two edges. The phase $`\varphi =\frac{e}{\mathrm{}c}[_{\mathrm{QPC1}\mathrm{A}\mathrm{QPC2}}\stackrel{}{A}𝑑\stackrel{}{l}_{\mathrm{QPC1}\mathrm{B}\mathrm{QPC2}}\stackrel{}{A}𝑑\stackrel{}{l}]=2\pi \mathrm{\Phi }/\mathrm{\Phi }_0`$, where $`\stackrel{}{A}`$ is the vector potential, $`\mathrm{\Phi }<0`$ the magnetic flux through the region A-QPC2-B-QPC1-A, $`e<0`$ the electron charge, and $`\mathrm{\Phi }_0`$ the magnetic flux quantum $`hc/e`$. The coefficients $`c_1`$ and $`c_2`$ will be calculated in section IV. Thus, the current exhibits a periodic dependence on the magnetic flux $`\mathrm{\Phi }`$ with period $`\mathrm{\Phi }_0`$. The amplitude of the flux-dependent contribution to the current $`I_\mathrm{\Phi }=\mathrm{max}I_{12}`$ is related to the flux-independent contribution $`I_0=I_1+I_2`$ by the equation
$$I_\mathrm{\Phi }(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)[I_0(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)I_0(\mathrm{\Gamma }_1,0)]^{1/2}.$$
(1)
Eq. (1) is derived rigorously in section IV.
A similar relation with a different exponent can be obtained for the set-up Fig. 1b) in which fractionally charged quasiparticles tunnel between the edges. Quasiparticles in a quantum Hall liquid with the filling factor $`\nu =1/(2m+1)`$ can be described as point charges $`q=\nu e`$ with attached solenoids anyons . Each solenoid carries one magnetic flux quantum $`\mathrm{\Phi }_0`$. When one quasiparticle makes a circle around another it picks up the Aharonov-Bohm phase
$$\theta =2\pi \nu $$
(2)
which describes fractional statistics. The total flux $`\stackrel{~}{\mathrm{\Phi }}`$ through the interferometer includes the contribution from the applied magnetic field $`\mathrm{\Phi }=\stackrel{}{B}𝑑\stackrel{}{S}`$ and the statistical contribution from the flux tubes attached to quasiparticles.
Fig. 2 is topologically equivalent to Fig. 1b). As is clear from Fig. 2, the solenoids attached to quasiparticles do not contribute to the magnetic flux through the dashed circle, if there is no tunneling between edges 1q and 2q. Thus, in the absence of tunneling the total magnetic flux $`\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Phi }`$. Each tunneling event changes the flux $`\stackrel{~}{\mathrm{\Phi }}`$ by one flux quantum. The flux decreases by $`|\mathrm{\Phi }_0|`$ when a quasiparticle tunnels from the outer edge to the internal edge, and increases by $`|\mathrm{\Phi }_0|`$ for the tunneling events from edge 2q to edge 1q. Hence, $`\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Phi }+n\mathrm{\Phi }_0`$, where $`n`$ is an integer.
We will assume that the electrochemical potential of source S1 is higher than the electrochemical potential of source S2 by $`eV`$, where $`e`$ is an electron charge. These electrochemical potentials are equal to the chemical potentials of edges 1q and 2q.
Let us first consider the simplest limit of zero temperature. In this case, quasiparticles can tunnel from the edge with the higher chemical potential to the edge with the lower chemical potential only and cannot tunnel from edge 2q to edge 1q. The tunneling probability can be derived in the same way as Eq. (1),
$$p(\stackrel{~}{\mathrm{\Phi }})=\stackrel{~}{c}_1(|\mathrm{\Gamma }_1|^2+|\mathrm{\Gamma }_2|^2)+(\stackrel{~}{c}_2\mathrm{\Gamma }_1\mathrm{\Gamma }_2^{}\mathrm{exp}(i\stackrel{~}{\varphi })+c.c.),$$
(3)
where $`\stackrel{~}{\varphi }=2\pi \nu \stackrel{~}{\mathrm{\Phi }}/\mathrm{\Phi }_0=2\pi \nu \mathrm{\Phi }/\mathrm{\Phi }_0+2\pi \nu n`$ is the Aharonov-Bohm phase accumulated by a quasiparticle with the charge $`\nu e`$ along the path A-QPC2-B-QPC1-A. The phase $`\stackrel{~}{\varphi }`$ is periodic in $`n`$ with period $`1/\nu `$. Let us assume that initially $`n=k/\nu `$. Then $`p(\stackrel{~}{\mathrm{\Phi }})=p(\mathrm{\Phi })`$ and the transfer of a quasiparticle requires the time $`t_1=1/p(\mathrm{\Phi })`$. After a quasiparticle tunneling event the statistical flux changes and the second quasiparticle tunnels after the time interval $`t_2=1/p(\mathrm{\Phi }+\mathrm{\Phi }_0)`$. After $`1/\nu `$ tunneling events we return to the situation with $`n`$ being a multiple of $`1/\nu `$. The time needed for the transfer of $`1/\nu `$ quasiparticles, i.e. a single electron charge, is $`t=_{l=1}^{1/\nu }t_l`$. Hence, the current
$$I=\frac{e}{t}=\frac{1/\nu }{_{l=1}^{1/\nu }\frac{1}{I_l}},$$
(4)
where $`I_l=\nu ep(\mathrm{\Phi }+[l1]\mathrm{\Phi }_0)`$. If the quasiparticles did not carry flux tubes the current through the Mach-Zehnder interferometer in the presence of the magnetic flux $`(\mathrm{\Phi }+[l1]\mathrm{\Phi }_0)`$ would equal $`I_l`$. Thus, the current of quasiparticles obeying fractional statistics is a harmonic average of $`1/\nu `$ currents corresponding to the imaginary situation of fractionally charged quasiparticles which do not obey fractional statistics.
If $`\mathrm{\Gamma }_2\mathrm{\Gamma }_1`$ then the current $`I`$ can be expanded in powers of $`\mathrm{\Gamma }_2`$. The term of order $`\mathrm{\Gamma }_2^u\mathrm{\Gamma }_2^w`$ is proportional to $`\mathrm{exp}(2\pi \nu i[uw]\mathrm{\Phi }/\mathrm{\Phi }_0)`$. The current (4) is a periodic function of the magnetic flux with period $`\mathrm{\Phi }_0`$. Hence, $`[uw]=k/\nu `$ in all non-zero contributions. Thus, the linear in $`|\mathrm{\Gamma }_2|`$ contribution to the current vanishes. Hence, the leading power of $`\mathrm{\Gamma }_2`$ in the expansion is $`|\mathrm{\Gamma }_2|^2`$. Such second order terms do not exhibit a magnetic flux dependence. Hence, the flux-independent contribution to the current $`I_0`$ satisfies the equation
$$I_0(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)I_0(\mathrm{\Gamma }_1,0)|\mathrm{\Gamma }_2|^2.$$
(5)
The leading flux-dependent contribution scales as
$$I_\mathrm{\Phi }\mathrm{const}\mathrm{\Gamma }_2^{1/\nu }\mathrm{exp}(2\pi i\mathrm{\Phi }/\mathrm{\Phi }_0)+c.c.$$
(6)
Finally, one gets from the comparison of Eqs. (5) and (6) at $`\mathrm{\Gamma }_2\mathrm{\Gamma }_1`$
$$I_\mathrm{\Phi }(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)[I_0(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)I_0(\mathrm{\Gamma }_1,0)]^{\pi /\theta },$$
(7)
where $`\theta `$ is the statistical angle (2). Eq. (7) is derived rigorously in section V. The flux-dependent and flux-independent contributions to the current can be expressed via the maximal and minimal values of the current as function of the magnetic flux: $`I_0=[\mathrm{max}I+\mathrm{min}I]/2`$, $`I_\mathrm{\Phi }=[\mathrm{max}I\mathrm{min}I]/2`$. The exponent $`\pi /\theta `$ is determined by the quasiparticle statistics. It equals $`1/2`$ for fermions, Eq. (1), and exceeds $`1/2`$ for fractional quasiparticles.
Luttinger liquid effects are know to give rise to power laws in the edge physics in quantum Hall systems KF . We would like to emphasize that Luttinger liquid physics is irrelevant for the above result, Eq. (7). The exponent $`\pi /\theta `$ emerges due to the statistical magnetic flux carried by quasiparticles, i.e. due to their fractional statistics.
We now discuss the finite temperature regime. In this case quasiparticles can tunnel both from edge 1q to edge 2q and from edge 2q to edge 1q. The tunneling probabilities are related by the principle of detailed balance
$$p_{}(\mathrm{\Phi }+[k\times \frac{1}{\nu }+r]\mathrm{\Phi }_0)=\gamma p_+(\mathrm{\Phi }+[k\times \frac{1}{\nu }+r1]\mathrm{\Phi }_0),$$
(8)
where $`\gamma =\mathrm{exp}(\nu eV/k_BT)`$; $`1r\nu `$; the convention $`r1=1/\nu `$ is used at $`r=1`$; $`p_+(\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Phi }+\mathrm{𝐬𝐭𝐚𝐭𝐢𝐬𝐭𝐢𝐜𝐚𝐥}\mathrm{𝐟𝐥𝐮𝐱}\mathrm{𝐛𝐞𝐟𝐨𝐫𝐞}\mathrm{𝐭𝐮𝐧𝐧𝐞𝐥𝐢𝐧𝐠})`$ and $`p_{}(\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Phi }+\mathrm{𝐬𝐭𝐚𝐭𝐢𝐬𝐭𝐢𝐜𝐚𝐥}\mathrm{𝐟𝐥𝐮𝐱}\mathrm{𝐛𝐞𝐟𝐨𝐫𝐞}\mathrm{𝐭𝐮𝐧𝐧𝐞𝐥𝐢𝐧𝐠})`$ denote the probabilities of the tunneling events from edge 1q to 2q and from edge 2q to edge 1q respectively. They depend quadratically on $`\mathrm{\Gamma }_{1,2}`$ and are calculated in section V. The tunneling probabilities depend on the total magnetic flux which includes the statistical contribution $`\mathrm{\Phi }_s=[k\times \frac{1}{\nu }+s]\mathrm{\Phi }_0`$, $`1s1/\nu `$. Since the Aharonov-Bohm phase due to the statistical flux $`\mathrm{\Phi }_s`$ equals $`\varphi _{AB}=(2\pi k+2\pi \nu s)`$, the probabilities depend on $`\mathrm{\Phi }`$ and $`s`$ only and are independent of $`k`$. Thus, the tunneling current is given by the equation
$$I=e\nu \underset{r=1}{\overset{1/\nu }{}}f_r[p_+(\mathrm{\Phi }+r\mathrm{\Phi }_0)p_{}(\mathrm{\Phi }+r\mathrm{\Phi }_0)],$$
(9)
where $`f_r`$ is the probability to find the system in one of the states with $`\mathrm{\Phi }_s=[k\times \frac{1}{\nu }+r]\mathrm{\Phi }_0`$, $`1r1/\nu `$ , $`k`$ being an arbitrary integer. The distribution function $`f_r`$ can be determined from the steady state condition
$$f_r[p_+(\mathrm{\Phi }+r\mathrm{\Phi }_0)+p_{}(\mathrm{\Phi }+r\mathrm{\Phi }_0)]=f_{r1}p_+(\mathrm{\Phi }+[r1]\mathrm{\Phi }_0)+f_{r+1}p_{}(\mathrm{\Phi }+[r+1]\mathrm{\Phi }_0)],$$
(10)
where we use the convention $`1/\nu +1=1`$. Using Eq. (8) the kinetic equation (10) can be rewritten as
$$p_+(\mathrm{\Phi }+r\mathrm{\Phi }_0)[f_r\gamma f_{r+1}]=p_+(\mathrm{\Phi }+[r1]\mathrm{\Phi }_0)[f_{r1}\gamma f_r];r=1,\mathrm{},1/\nu .$$
(11)
One finds from the above system of equations
$$f_r\gamma f_{r+1}=\frac{\alpha }{p_+(\mathrm{\Phi }+r\mathrm{\Phi }_0)},$$
(12)
where $`\alpha `$ is a constant independent of $`r`$. To calculate this constant we add Eqs. (12) with all possible $`r`$. Since $`f_r=1`$, one gets
$$\alpha =\frac{1\gamma }{_{r=1}^{1/\nu }\frac{1}{p_+(\mathrm{\Phi }+r\mathrm{\Phi }_0)}}.$$
(13)
Thus, the current, Eq. (9), equals
$$I=e\nu \underset{r=1}{\overset{1/\nu }{}}p_+(\mathrm{\Phi }+r\mathrm{\Phi }_0)[f_r\gamma f_{r+1}]=e\alpha =\frac{1/\nu }{_{r=1}^{1/\nu }\frac{1}{I_r^{}}},$$
(14)
where $`I_r^{}=\nu e(1\gamma )p_+(\mathrm{\Phi }+r\mathrm{\Phi }_0)`$. Similarly to $`I_l`$, Eq. (4), the currents $`I_r^{}`$ can be understood as the currents of fictitious fractionally charged quasiparticles which do not obey fractional statistics. A rigorous derivation of Eq. (14) is discussed in section V. Finally, the same analysis as in the zero-temperature case shows that Eq. (7) is satisfied at finite temperatures.
Eq. (7) is the main result of the article. An experimental test of this relation will allow the observation of fractional statistics. In such experiment one needs to change $`\mathrm{\Gamma }_2`$ at fixed $`\mathrm{\Gamma }_1`$. The tunneling amplitudes $`\mathrm{\Gamma }_{1,2}`$ are controlled by gate voltages. Generally, any change of gate voltages affects both tunneling amplitudes. They can be controlled independently only if QPC1 and QPC2 are far from each other. In most interferometer set-ups an increase in the distance between QPC1 and QPC2 would result in the suppression of the interference pattern. Fortunately, this is not the case for the Mach-Zehnder interferometer. The calculations for the electronic Mach-Zehnder interferometer in the integer quantum Hall regime FMZ1 ; FMZ2 show that the visibility of the interference pattern depends not on the distance between QPC1 and QPC2 but only on the difference of the distances between the point contacts along two edges. We confirm the same conclusion for the fractional quantum Hall systems in sections IV and V.
In order to determine the relation between the flux-dependent and flux-independent components of the current one has to vary the magnetic field and measure the current at different values of the field. The magnetic flux $`\mathrm{\Phi }`$ through the region QPC1-A-QPC2-B-QPC1, Fig. 1b), includes two contribution: the flux through the hole in the interferometer, i.e. the upper half of the region QPC1-A-QPC2-B-QPC1, and the flux through the lower half of the region QPC1-A-QPC2-B-QPC1 which is occupied by a quantum Hall liquid. If the magnetic field changes at the fixed density then the filling factor in the lower half deviates from $`\nu `$. As a result, quasiparticles can enter the region QPC1-A-QPC2-B-QPC1. Each of them brings one flux quantum. This does not change any of the results of the paper. Indeed, we predict that the current is a periodic function of the magnetic flux through the region QPC1-A-QPC2-B-QPC1 with period $`\mathrm{\Phi }_0`$. Hence, changing the flux by one flux quantum does not affect the current.
Several other quantum Hall interferometer set-ups have been discussed in the literature. The simplest set-up cfksw is illustrated in Fig. 3. In contrast to the Mach-Zehnder interferometer, the effective magnetic flux perceived by quasiparticles does not change after tunneling events in the set-up Fig. 3. Hence, the current is independent of the statistical phase $`\theta `$, Eq. (2) (if no quasiparticles are trapped between the quantum point contacts). On the other hand, the current exhibits a “fractional” Aharonov-Bohm periodicity with period $`\mathrm{\Phi }_0/\nu `$. On the technical level the set-up cfksw and our problem are described by very similar models (see section III). The main difference consists in the Klein factors which describe fractional statistics. They are present in the model of the Mach-Zehnder interferometer (section III) and are absent in the model of Ref. cfksw, . This difference between the models results in qualitatively different transport behavior. A set-up related to the Mach-Zehnder interferometer was studied in Ref. kane, , Fig. 4. The interferometer, Fig. 4, has the same topology as the Mach-Zehnder interferometer but includes three edges and three quantum point contacts. In the absence of QPC3 the set-up Fig. 4 is equivalent to Fig. 3. If weak tunneling at QPC3 is allowed the system exhibits strong telegraphic noise which carries information about fractional statistics. In contrast to Ref. kane, we investigate the set-up with two point contacts that was studied experimentally MZ . In our case not only the noise but also the average current is strongly sensitive to fractional statistics.
## III Model
### III.1 Effective action for electron tunneling case
We first consider the system represented by Fig. 1a). Its low-energy behavior can be described by the chiral Luttinger liquid model wen . Fig. 5 illustrates the model. The two edges 1 and 2 correspond to two chiral Luttinger liquids with the same propagation direction. The dashed lines describe quantum point contacts where electrons tunnel between the edges. Note that the respective distances $`L`$ and $`L+a`$ between the point contacts along the two edges are different. The Lagrangian assumes the standard form wen
$$L=\frac{\mathrm{}}{4\pi }𝑑x𝑑t\underset{k=1,2}{}[_t\varphi _k_x\varphi _k+v(_x\varphi _k)^2]𝑑t(T_1+T_2),$$
(15)
where $`T_1`$ and $`T_2`$ are tunneling operators, $`v`$ is the excitation velocity along the edges, and $`\{\varphi _k\}`$ represent two chiral Bose fields which satisfy the following commutation relations
$$[\varphi _l(x_l,t=0),\varphi _p(x_p,t=0)]=i\pi \delta _{lp}\mathrm{sign}(x_lx_p).$$
(16)
The Bose fields are related to the charge densities $`\rho _l`$ through
$$\rho _l=(\sqrt{\nu _l}e/2\pi )_x\varphi _l,$$
(17)
where $`e`$ is an electron charge and $`\nu _l=1/(2m_l+1)`$ are the filling factors of the QHE puddle defined by edge 2e and the QHE strip bounded by edge 1e. The tunneling operators wen
$`T_1=\mathrm{\Gamma }_1\mathrm{exp}(i[\varphi _1(0,t)/\sqrt{\nu _1}\varphi _2(0,t)/\sqrt{\nu _2}])+h.c.;`$
$`T_2=\mathrm{\Gamma }_2\mathrm{exp}(i[\varphi _1(L,t)/\sqrt{\nu _1}\varphi _2(L+a,t)/\sqrt{\nu _2}])+h.c`$ (18)
are proportional to the electron annihilation and creation operators $`\mathrm{exp}(\pm i\varphi _l/\sqrt{\nu _l})`$. The action also includes operators describing simultaneous tunneling of several electrons but they do not play a significant role at low energies.
In the presence of a magnetic flux $`\mathrm{\Phi }`$ through the region A-QPC2-B-QPC1-A, the tunneling amplitude $`\mathrm{\Gamma }_2`$ should be multiplied by the phase factor $`\mathrm{exp}(2\pi i\mathrm{\Phi }/\mathrm{\Phi }_0)`$, where $`\mathrm{\Phi }_0=hc/e`$ is the flux quantum footnote . This phase factor describes the difference of the phases $`e/(\mathrm{}c)_{\mathrm{QPC1}}^{\mathrm{QPC2}}\stackrel{}{A}𝑑\stackrel{}{l}`$ accumulated by the electrons moving along two edges between the point contacts. For one edge the integration path is QPC1-A-QPC2, and for the other edge the integration path is QPC1-B-QPC2 (Fig. 1).
The voltage bias $`V`$ results in the chemical potential difference $`\mu _1\mu _2=eV`$ between edges 1e and 2e. We would like to emphasize that the potential difference between the edges is determined by the voltage difference between the sources S1 and S2 in both set-ups Fig. 1a) and Fig. 1b); however, the shape of the edges 1e and 2e, Fig. 1a), is different from the shape of the edges 1q and 2q, Fig. 1b). We will use the interaction representation which makes both chemical potentials equal and introduces time-dependence into the tunneling amplitudes:
$$\mathrm{\Gamma }_1,\mathrm{\Gamma }_2\mathrm{exp}(\frac{ieVt}{\mathrm{}}).$$
(19)
The difference of the chemical potentials results also in the difference of the average densities between the edges. We assume that edge 2 is connected to the ground. Then the average density $`\rho _2_x\varphi _2`$ is zero in the limit of weak tunneling. It will be convienient for us to shift the Bose field on the first edge $`\varphi _1(x)\varphi _1(x)f(x)`$ by a function $`f(x)`$ of the coordinate in such a way that $`_x\varphi _1`$ becomes zero. The mean charge density on edge 1 is proportional to the voltage bias and can be found from the minimization of the Hamiltonian
$$H=\mathrm{}𝑑x\left[\frac{v}{4\pi }(_x\varphi _1)^2\frac{eV\sqrt{\nu _1}}{2\pi \mathrm{}}_x\varphi _1\right].$$
(20)
This yields $`q=_x\varphi _1=eV\sqrt{\nu _1}/(\mathrm{}v)`$. We next shift the field $`\varphi _1\varphi _1qx`$ such that $`_x\varphi _1`$ vanishes. At the same time $`\mathrm{\Gamma }_2`$ is multiplied by the factor $`\mathrm{exp}(ieVL/\mathrm{}v)`$.
Hereafter we assume that $`\nu _1\nu _2`$. The case $`\nu _1=\nu _2=1`$ corresponds to the free electron problem FMZ1 ; FMZ2 .
### III.2 Effective action for quasiparticle tunneling case
We now consider the system depicted in Fig. 1b). The model is given by the chiral Luttinger liquid action (15) with a different choice of the tunneling operators wen . Several modifications immediately follow from the fact that the quasiparticle charge $`\nu e`$ differs from the electron charge. 1) The tunneling operators should be expressed via the quasiparticle annihilation and creation operators $`\mathrm{exp}(\pm i\sqrt{\nu }\varphi _l)`$; 2) the flux-dependent phase factor in $`\mathrm{\Gamma }_2`$ is now $`\mathrm{exp}(2\pi \nu i\mathrm{\Phi }/\mathrm{\Phi }_0)`$; 3) the phase factor $`\mathrm{exp}(ie\nu VL/\mathrm{}v)`$ should be used instead of $`\mathrm{exp}(ieVL/\mathrm{}v)`$; and 4) in the interaction representation the time-dependence of the tunneling amplitudes becomes
$$\mathrm{\Gamma }_1,\mathrm{\Gamma }_2\mathrm{exp}(\frac{ie\nu Vt}{\mathrm{}}).$$
(21)
The fifth difference from Eqs. (III.1) consists in the introduction of two Klein factors $`\kappa _1`$ and $`\kappa _2`$ (a related model without Klein factors has been considered in Ref. tm, ):
$`T_1^q=\mathrm{\Gamma }_1\kappa _1\mathrm{exp}(i\sqrt{\nu }[\varphi _1(0,t)\varphi _2(0,t)])+h.c.;`$
$`T_2^q=\mathrm{\Gamma }_2\kappa _2\mathrm{exp}(i\sqrt{\nu }[\varphi _1(L,t)\varphi _2(L+a,t)])+h.c,`$ (22)
where the commutation relations are
$`\kappa _1\kappa _2=\mathrm{exp}(2\pi \nu i)\kappa _2\kappa _1;`$
$`\kappa _1\kappa _2^+=\kappa _2^+\kappa _1\mathrm{exp}(2\pi \nu i).`$ (23)
These commutation relations can be understood from a locality argument similar to Ref. kane, . Two tunneling operators affect two distant parts of the system. Hence, for any reasonable model $`[T_1^q,T_2^q]=0`$. One can calculate $`[T_1^q,T_2^q]`$ employing the commutation relations for the Klein factors, the commutation relations Eq. (16) for the Bose fields and the Baker-Hausdorff formula. This results in $`[T_1^q,T_2^q]=[T_1^q,(T_2^q)^+]=0`$ provided that Eq. (III.2) is satisfied. A different formulation of the same argument can be found in Ref. kane, . Ref. pa, discusses how Klein factors which ensure commutativity of tunneling operators can be derived using duality between weak quasiparticle tunneling and strong electron tunneling.
The Klein factors serve as a manifestation of fractional statistics and are absent in the case of fermion tunneling, section III.A. The importance of Klein factors in quantum Hall systems with more than two edges has been emphasized previously HBT ; kane ; pa ; nfll ; klein . In our problem, the Klein factors are necessary even though there are only two edges. Note that in the 1a) setup no Klein factors are needed. Indeed, the operators $`T_1`$ and $`T_2`$ defined by Eq. (III.1) commute. Commutativity of the quasiparticle tunneling operators is also ensured without Klein factors for the standard Aharonov-Bohm interferometer geometry cfksw , Fig. 3.
In our calculations we will use the following representation of the Klein factors by $`1/\nu \times 1/\nu `$ matrices:
$$\kappa _1=\left(\begin{array}{cccccc}0& 1& 0& 0& \mathrm{}& 0\\ 0& 0& 1& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0& 1& 0\\ 0& \mathrm{}& \mathrm{}& \mathrm{}& 0& 1\\ 1& 0& \mathrm{}& \mathrm{}& \mathrm{}& 0\end{array}\right);\kappa _2=\left(\begin{array}{cccccc}0& \psi & 0& 0& \mathrm{}& 0\\ 0& 0& \psi ^2& 0& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& \mathrm{}& 0& \psi ^{1/\nu 2}& 0\\ 0& \mathrm{}& \mathrm{}& \mathrm{}& 0& \psi ^{1/\nu 1}\\ 1& 0& \mathrm{}& \mathrm{}& \mathrm{}& 0\end{array}\right),$$
(24)
where $`\psi =\mathrm{exp}(2\pi \nu i)`$. One can easily check that the above matrices satisfy the commutation relations (III.2). The tunneling operators (III.2) with the Klein factors (24) can be understood as products of a quasiparticle creation operator, an annihilation operator and a phase factor which includes the statistical phase accumulated during the tunneling event.
The charge distribution on the edges is determined by the fields $`\{\varphi _l(x)\}`$, the total charge being determined by the zero modes of these fields. In the set-up depicted in Fig. 1a) the charge distribution completely describes all states in the Hilbert space on which the effective low-energy Hamiltonian acts. On the other hand, our discussion in section II shows that in the case of quasiparticle tunneling one needs to specify both the charge distribution and the effective statistical flux through the interferometer. Hence, the corresponding Hilbert space is the product of the space $`V_{\mathrm{charge}}`$ on which the Bose operators $`\varphi _l`$ act and the space $`V_{\mathrm{flux}}`$ on which the Klein factors act. As is clear from the size of the matrices (24) the dimensionality of the latter space is $`\mathrm{dim}V_{\mathrm{flux}}=1/\nu `$. This agrees with our discussion in section II where we found that the interferometer has $`1/\nu `$ classes of states characterized by different probabilities of quasiparticle tunneling.
Calculations based on the model (15) with the tunneling operators (III.2) confirm Eq. (7) of Section II. The Klein factors keep track of fractional statistics and are crucial for this result. In the absence of the Klein factors, i.e. for fractionally charged particles which do not obey fractional statistics, one gets qualitatively different results tm .
## IV Electron tunneling
We now study the electric current between the two edges. We consider the geometry depicted in Fig. 1a). The current operator
$`\widehat{I}={\displaystyle \frac{d}{dt}}\widehat{Q}_1={\displaystyle \frac{i}{\mathrm{}}}[\widehat{H},\widehat{Q}_1]={\displaystyle \frac{ie\mathrm{\Gamma }_1}{\mathrm{}}}\mathrm{exp}(ieVt/\mathrm{})\mathrm{exp}[i(\varphi _1(0)/\sqrt{\nu _1}\varphi _2(0)/\sqrt{\nu _2})]`$
$`+{\displaystyle \frac{ie\mathrm{\Gamma }_2}{\mathrm{}}}\mathrm{exp}(ieVL/\mathrm{}v+2\pi i\mathrm{\Phi }/\mathrm{\Phi }_0)\mathrm{exp}(ieVt/\mathrm{})\mathrm{exp}[i(\varphi _1(L)/\sqrt{\nu _1}\varphi _2(L+a)/\sqrt{\nu _2})]+h.c.,`$ (25)
where $`Q_1`$ is the total charge of the first edge and $`H`$ the Hamiltonian. For our non-equilibrium problem we employ the Keldysh technique Keldysh . To this end we assume that the tunneling amplitudes $`\mathrm{\Gamma }_{1,2}=0`$ at the moment of time $`t=\mathrm{}`$, and are subsequently turned on gradually. At $`t=\mathrm{}`$ the system is in thermal equilibrium at temperature $`k_BT`$ and chemical potential difference $`eV`$ between the edges. The initial equilibrium state determines the bare Keldysh Green functions which will be used in the perturbative calculations below.
The current at $`t=0`$ is
$$I=\mathrm{Tr}[\widehat{\rho }S(\mathrm{},0)\widehat{I}S(0,\mathrm{})],$$
(26)
where $`\widehat{\rho }`$ is the initial density matrix and $`S(0,\mathrm{})=\mathrm{T}\mathrm{exp}(i\widehat{H}𝑑t/\mathrm{})`$ the evolution operator. Expanding the latter to first order in the tunneling amplitudes one finds
$`I={\displaystyle \frac{e}{\mathrm{}^2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑t(|\mathrm{\Gamma }_1|^2+|\mathrm{\Gamma }_2|^2)\mathrm{exp}({\displaystyle \frac{ieVt}{\mathrm{}}})[F(0,0,t)F(0,0,t)]`$
$`+{\displaystyle \frac{e}{\mathrm{}^2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}dt\{\mathrm{\Gamma }_1\mathrm{\Gamma }_2^{}\mathrm{exp}({\displaystyle \frac{ieVL}{\mathrm{}v}}2\pi i\mathrm{\Phi }/\mathrm{\Phi }_0)\mathrm{exp}({\displaystyle \frac{ieVt}{\mathrm{}}})[F(L,(L+a),t)F(L,L+a,t)]+h.c.\},`$ (27)
where
$`F(b,c,t)=\mathrm{Tr}[\widehat{\rho }\mathrm{exp}(i\varphi _1(x=b,t)/\sqrt{\nu _1})\mathrm{exp}(i\varphi _1(x=0,0)/\sqrt{\nu _1})]`$
$`\times \mathrm{Tr}[\widehat{\rho }\mathrm{exp}(i\varphi _2(x=c,t)/\sqrt{\nu _2})\mathrm{exp}(i\varphi _2(x=0,0)/\sqrt{\nu _2})].`$ (28)
Let us first consider the zero temperature case. The correlation function is given by bos
$$F(b,c,t)=\frac{\tau _c^{1/\nu _1}}{[\delta +i(tb/v)]^{1/\nu _1}}\frac{\tau _c^{1/\nu _2}}{[\delta +i(tc/v)]^{1/\nu _2}},$$
(29)
where $`\delta `$ is an infinitesimal positive constant, and $`\tau _c`$ is the ultra-violet cut-off foot2 . With the above expression we find
$$I=I_0+I_\mathrm{\Phi },$$
(30)
where
$`I_0={\displaystyle \frac{2\pi e\tau _c}{\mathrm{}^2}}(|\mathrm{\Gamma }_1|^2+|\mathrm{\Gamma }_2|^2){\displaystyle \frac{[\tau _ceV/\mathrm{}]^{1/\nu _1+1/\nu _21}}{(1/\nu _1+1/\nu _21)!}},`$ (31)
$`I_\mathrm{\Phi }={\displaystyle \frac{2\pi ie\tau _c^{1/\nu _1+1/\nu _2}}{\mathrm{}^2}}(1)^{1/[2\nu _1]+1/[2\nu _2]}\times `$
$`\times \{\mathrm{\Gamma }_1\mathrm{\Gamma }_2^{}\mathrm{exp}(2\pi i\mathrm{\Phi }/\mathrm{\Phi }_0)[{\displaystyle \frac{1}{(1/\nu _11)!}}{\displaystyle \frac{d^{1/\nu _11}}{dz^{1/\nu _11}}}|_{z=0}{\displaystyle \frac{\mathrm{exp}(\frac{ieVz}{\mathrm{}})}{(z+a/v)^{1/\nu _2}}}+{\displaystyle \frac{1}{(1/\nu _21)!}}{\displaystyle \frac{d^{1/\nu _21}}{dz^{1/\nu _21}}}|_{z=a/v}{\displaystyle \frac{\mathrm{exp}(\frac{ieVz}{\mathrm{}})}{z^{1/\nu _1}}}]h.c.\}`$ (32)
The current oscillates as function of the magnetic flux with period $`\mathrm{\Phi }_0`$. At low voltages it follows the power law
$$IV^{1/\nu _1+1/\nu _21}$$
(33)
(see Appendix A; cf. Ref. geller, ). Fig. 6 illustrates the $`IV`$ curves for $`\nu _1=\nu _2=1`$, $`\nu _1=\nu _2=1/3`$ and $`\nu _1=1,\nu _2=1/3`$. In the case when $`\nu _1=\nu _2`$, similar expressions for other interferometer geometries were obtained in Refs. cfksw, ; geller, . One can easily see that Eq. (1) follows from Eqs. (30-IV).
Notice that only the difference $`a`$ of the lengths of the edges enters the above expression for the current while the total length $`L`$ of the first edge drops out. In the standard geometry of an Aharonov-Bohm interferometer cfksw the flux-dependent contribution to the current depends on the total interferometer size and decreases with the system’s size. Thus, quantum interference effects cannot be observed at large system’s sizes. In the Mach-Zehnder geometry, quantum interference can be observed for $`La`$ since only the difference of the phases accumulated by the particles moving along two edges is important.
The case of non-zero temperature is considered in Appendix B. The flux-dependent “interference” contribution $`I_\mathrm{\Phi }`$ to the current vanishes at large $`ak_BT/(hv)`$. In the opposite limit of $`ak_BT/(hv)1`$ as well as for $`ak_BThv`$, the flux-independent contribution $`I_0`$ and $`I_\mathrm{\Phi }`$ are related by Eq. (1). The linear conductance at low temperatures and low voltages $`eVk_BT`$ scales as
$$G=I/V[k_BT]^{1/\nu _1+1/\nu _22}.$$
(34)
## V Quasiparticle tunneling
We now consider the geometry depicted in Fig. 1b). The current through the interferometer oscillates as function of the magnetic flux. In the first subsection below we determine the oscillation period. It turns out to be the same as in the case of electron tunneling and equals one flux quantum $`\mathrm{\Phi }_0`$. Next, we use the perturbation theory to calculate the current as function of the voltage, temperature and the distance between the quantum point contacts. We confirm Eq. (14). The dependence of the current (14) on the tunneling amplitudes $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ is nonanalytic. Such dependence cannot be obtained in any finite order of the perturbation theory. Thus, we have to sum up an infinite set of diagrams. In subsection V.B this is made for the simplest case of zero temperature and low voltage, $`eVa\mathrm{}v`$, at the filling factor $`\nu =1/3`$. The general case is considered in the final subsection.
### V.1 Period of Aharonov-Bohm oscillations
The tunneling current operator can be found with the same method as in the previous section:
$`\widehat{I}={\displaystyle \frac{d}{dt}}\widehat{Q}_1={\displaystyle \frac{i}{\mathrm{}}}[\widehat{H},\widehat{Q}_1]={\displaystyle \frac{ie\kappa _1\nu \mathrm{\Gamma }_1}{\mathrm{}}}\mathrm{exp}(ie\nu Vt/\mathrm{})\mathrm{exp}[i\sqrt{\nu }(\varphi _1(0)\varphi _2(0))]`$
$`+{\displaystyle \frac{ie\kappa _2\nu \mathrm{\Gamma }_2}{\mathrm{}}}\mathrm{exp}({\displaystyle \frac{ie\nu VL}{\mathrm{}v}}+2\pi \nu i\mathrm{\Phi }/\mathrm{\Phi }_0)\mathrm{exp}(ie\nu Vt/\mathrm{})\mathrm{exp}[i\sqrt{\nu }(\varphi _1(L)\varphi _2(L+a))]+h.c.`$ (35)
The average current is given by Eq. (26). Since the tunneling amplitudes $`\mathrm{\Gamma }_l`$ are small, we will employ perturbation theory. To lowest (second) nonzero order in $`\mathrm{\Gamma }_l`$ the resulting contributions are proportional to $`|\mathrm{\Gamma }_1|^2`$ and $`|\mathrm{\Gamma }_2|^2`$. The cross-terms proportional to $`\mathrm{\Gamma }_1^{}\mathrm{\Gamma }_2`$ and $`\mathrm{\Gamma }_1\mathrm{\Gamma }_2^{}`$ vanish. Indeed, each term of the perturbative expansion is proportional to the product of the average of some function of the Bose-operators $`\{\varphi _l\}`$ and the average of some product of Klein factors. The averages are determined by the initial density matrix $`\widehat{\rho }`$, Eq. (26). The latter depends on the effective action Eq. (15) at $`t=\mathrm{}`$. That action does not contain Klein factors. Hence, the density matrix $`\widehat{\rho }=\widehat{\rho }_\varphi \widehat{\rho }_\kappa `$, where $`\widehat{\rho }_\varphi `$ and $`\widehat{\rho }_\kappa `$ act on the spaces $`V_{\mathrm{charge}}`$ and $`V_{\mathrm{flux}}`$ respectively, and $`\widehat{\rho }_\kappa `$ is proportional to the unit matrix at any finite temperature. Thus, the cross-terms are proportional to the expressions of the form $`\mathrm{Tr}[\kappa _l\kappa _m^+]`$, where $`lm`$. Such traces are zero. This is readily seen from the following argument. For any two linear operators $`\mathrm{Tr}AB=\mathrm{Tr}BA`$. Hence, $`\mathrm{Tr}[\kappa _1\kappa _2^+]=\mathrm{Tr}[\kappa _2^+\kappa _1]`$. At the same time it follows from Eq. (III.2) that $`\mathrm{Tr}[\kappa _1\kappa _2^+]=\mathrm{exp}(2\pi \nu i)\mathrm{Tr}[\kappa _2^+\kappa _1]`$. Hence, $`\mathrm{Tr}[\kappa _1\kappa _2^+]=0`$. It follows that there are no cross-terms in the second order perturbation theory.
The above result implies that the second-order terms are independent of the magnetic flux. The flux dependence of the current emerges only in higher orders of the perturbation theory. Each term of the perturbative expansion of Eq. (26) contains a product of Klein factors. Let the number of the Klein factors $`\kappa _l^+`$ in a given term $`I_\alpha `$ be $`n_l^+`$ and the number of the factors $`\kappa _l`$ be $`n_l^{}`$ . Each of the four numbers $`n_l^\pm `$ indicates the power in which the respective coefficient $`\mathrm{\Gamma }_1`$, $`\mathrm{\Gamma }_2`$, $`\mathrm{\Gamma }_1^{}`$ or $`\mathrm{\Gamma }_2^{}`$ enters in $`I_\alpha `$. Nonzero terms of the perturbation series describe the processes which do not change the edge charges. Hence, $`n_1^+n_1^{}=(n_2^+n_2^{})`$.
We next show that $`\nu (n_1^+n_1^{})^2`$ is an integer for any non-vanishing $`I_\alpha `$. Consider the trace $`W`$ of the product of the Klein factors in the term $`I_\alpha `$. We can move all operators $`\kappa _1^+,\kappa _1`$ to the left of all operators $`\kappa _2^+,\kappa _2`$. This will produce a phase factor $`\mathrm{exp}(i\gamma )`$. The trace can be represented as $`W=\mathrm{exp}(i\gamma )\mathrm{Tr}K_1K_2`$, where $`K_1`$ is a product of $`\kappa _1^+,\kappa _1`$ and $`K_2`$ is a product of $`\kappa _2^+,\kappa _2`$. We know that
$$\mathrm{Tr}K_1K_2=\mathrm{Tr}K_2K_1.$$
(36)
At the same time, one can move all Klein factors $`\kappa _2^+,\kappa _2`$ in the product $`K_1K_2`$ to the left of the operator $`K_1`$ using the commutation relations Eq. (III.2). This yields
$$K_1K_2=\mathrm{exp}[2\pi \nu i(n_1^+n_1^{})^2]K_2K_1.$$
(37)
Eqs. (36,37) show that $`I_\alpha \mathrm{Tr}K_1K_20`$ only if $`\nu (n_1^+n_1^{})^2`$ is an integer.
If $`\nu =1/[\mathrm{odd}\mathrm{prime}\mathrm{number}]`$, including experimentally relevant $`\nu =1/3`$ and $`\nu =1/5`$, the above result means that $`(n_1^+n_1^{})=(n_2^+n_2^{})=q/\nu `$, where $`q`$ is an integer. The term $`I_\alpha \mathrm{\Gamma }_2^{n_2^{}}(\mathrm{\Gamma }_2^{})^{n_2^+}\mathrm{exp}(2\pi \nu i\mathrm{\Phi }/\mathrm{\Phi }_0[n_2^{}n_2^+])`$. Hence, it follows that the current is a periodic function of the magnetic flux with period $`\mathrm{\Phi }_0`$. As shown in the following subsections the period is the same for any $`\nu =1/(2m+1)`$. This can also be verified by a direct calculation of the trace (36) using (24). Such periodicity agrees with the Byers-Yang theorem BY ; footnote2 . If fractionally charged quasiparticles did not obey fractional statistics, i.e. if there were no Klein factors, the period would be $`\mathrm{\Phi }^{}=\mathrm{\Phi }_0/\nu `$ (cf. Ref. tm, ).
Our effective hydrodynamic action (15,III.2) is applicable only for low temperatures and voltages. We are going to use the perturbation theory in $`\mathrm{\Gamma }_{1,2}`$. It turns out that the lowest order contribution to the current scales as $`IV\mathrm{\Gamma }_{1,2}^2\mathrm{max}(eV,k_BT)^\alpha `$, where $`\alpha =2\nu 2`$ is negative. The perturbation theory can be used for the calculation of the tunneling current only when the tunneling current is much smaller than the current $`\nu e^2V/h`$ incoming from the sources. Thus, our calculations are valid provided that $`\mathrm{\Gamma }_{1,2}\mathrm{max}(eV,k_BT)^{\nu 1}\left(\frac{\mathrm{}}{\tau _c}\right)^\nu `$.
We will see that interference effects can be observed only for small enough $`ahv/\mathrm{max}(eV,k_BT)`$. This condition is similar to the restriction on the total interferometer size in Ref. cfksw, . In our case the restriction on the total size $`L`$ is weaker. As is clear from section II our analysis is based on the assumption that the time $`\mathrm{\Delta }t`$ between tunneling events exceeds the time $`L/v`$ needed to a quasiparticle to travel from QPC1 to QPC2. Thus, $`L<v\mathrm{\Delta }t\frac{v\mathrm{}^2}{\mathrm{\Gamma }_{1,2}^2\tau _c}\left(\frac{\mathrm{}}{\mathrm{max}(eV,k_BT)\tau _c}\right)^{2\nu 1}`$
### V.2 $`\nu =\mathrm{𝟏}/\mathrm{𝟑}`$, $`T=\mathrm{𝟎}`$, $`\mathrm{𝑒𝑉𝑎}\mathrm{}v`$
We will use the expansion of the current, Eq. (26), in powers of the tunneling amplitudes $`\mathrm{\Gamma }_{1,2}`$. Only contributions proportional to even powers of the tunneling operators are non-zero. One might naively expect on the basis of power counting that the terms of order $`2n`$ in $`\mathrm{\Gamma }_{1,2}`$ scale as $`V|\mathrm{\Gamma }_1|^p|\mathrm{\Gamma }_2|^{2np}\times [\mathrm{max}(eV,T)]^{2n(\nu 1)}`$, if $`Lvh/\mathrm{max}(eV,k_BT)`$. This is however not the case beyond the second perturbative order. In fact, the 4th order contribution (as well as higher orders) is infinite. To demonstrate this we use an analogy between the power expansion of the average current and the partition function of a Coulomb gas CQ . Positively and negatively charged “particles” correspond to the tunneling events from edge 1q to edge 2q and from edge 2q to edge 1q respectively. The coordinates of the particles correspond to the times of the tunneling events. The charges can be located on both branches of the Keldysh contour, Fig. 7. The particles on the top and lower branches emerge from the expansions of $`S(0,\mathrm{})`$ and $`S(\mathrm{},0)`$ in Eq. (26) respectively. As discussed in Ref. CQ, , attraction between oppositely charged particles binds them in pairs. The typical pair size is of order $`\mathrm{}/\mathrm{max}(eV,k_BT)`$ at low temperatures and voltages footnote3 (if $`L>hv/\mathrm{max}(eV,k_BT)`$ then pairs of size $`L/v`$ are possible). When $`\frac{\mathrm{\Gamma }_{1,2}\tau _c}{\mathrm{}}\left[\frac{\mathrm{max}(eV,k_BT)\tau _c}{\mathrm{}}\right]^{\nu 1}1`$ and $`L\frac{v\mathrm{}^2}{\mathrm{\Gamma }_{1,2}^2\tau _c}\left(\frac{\mathrm{}}{\mathrm{max}(eV,k_BT)\tau _c}\right)^{2\nu 1}`$, the pairs are dilute and do not overlap. Hence, any term in the perturbation expansion reduces to the trace of a product of Klein factors times a product of the two-point correlation functions
$`F[b=x_1x_2;c=(x_1x_2)(L+a)/L;t=t_1t_2]=\mathrm{Tr}\{\widehat{\rho }\mathrm{exp}[i\sqrt{\nu }\varphi _1(x=x_1,t_1)]\mathrm{exp}[i\sqrt{\nu }\varphi _1(x=x_2,t_2)]\}\times `$
$`\mathrm{Tr}\{\widehat{\rho }\mathrm{exp}[i\sqrt{\nu }\varphi _2(x=x_1(L+a)/L,t_1)]\mathrm{exp}[i\sqrt{\nu }\varphi _2(x=x_2(L+a)/L,t_2)]\}`$ (38)
corresponding to the bound pairs. This product should be integrated over the times of all tunneling events, i.e. the positions of the charges. This integration can be separated into the product of the integrals over the dipole sizes $`(t_1t_2)`$, Eq. (V.2), and the integrals over the dipole positions $`(t_1+t_2)/2`$.
Fig. 7 illustrates the 4th order contribution to the current. One charge is located at $`t=0`$ and corresponds to the current operator in Eq. (26). It forms a dipole with an opposite charge at the point $`\tau _0\mathrm{}/\mathrm{max}(eV,k_BT)`$. Two more mutually opposite charges are located at the points $`t_1\pm \tau _1/2`$, where $`\tau _1`$ is the dipole size. The charges can reside on the same or different branches of the Keldysh contour as shown in Figs. 7a), 7b) and 7c). The integral over $`t_1`$ diverges. In the absence of the Klein factors the infinite integrals corresponding to the configurations Fig. 7a), 7b) and 7c) cancel but this is no longer the case when the Klein factors are included. Hence, the 4th order contribution is infinite. The same argument applies to higher-order contributions. Certainly, the sum of all perturbative orders must be finite but as is clear from section II it is not an analytic function of $`\mathrm{\Gamma }_{1,2}`$. Below we develop a method to sum up all orders of the perturbation theory.
We begin with the simplest case when $`\nu =1/3`$; $`T=0`$; $`eVa,eVL\mathrm{}v`$. The same approach will be applied to the general situation in section V.C. The first condition simplifies the structure of the Klein factors (24) which become $`3\times 3`$ matrices at filling factor $`1/3`$. The second condition allows us to use a simpler zero-temperature expression for the correlation function
$`F(b,c,t_1t_2)=\mathrm{Tr}[\widehat{\rho }\mathrm{exp}(i\sqrt{\nu }\varphi _1(x=b,t_1))\mathrm{exp}(i\sqrt{\nu }\varphi _1(x=0,t_2))]\times `$
$`\mathrm{Tr}[\widehat{\rho }\mathrm{exp}(i\sqrt{\nu }\varphi _2(x=c,t_1))\mathrm{exp}(i\sqrt{\nu }\varphi _2(x=0,t_2))]={\displaystyle \frac{\tau _c^\nu }{[\delta +i(t_1t_2b/v)]^\nu }}{\displaystyle \frac{\tau _c^\nu }{[\delta +i(t_1t_2c/v)]^\nu }}`$ (39)
The third condition makes it possible to neglect the distances $`L`$ and $`(L+a)`$ between the point contacts. Indeed, the tunneling operator $`T_2^q=:\mathrm{\Gamma }_2\kappa _2\mathrm{exp}(i\sqrt{\nu }[\varphi _1(L,t)\varphi _2(L+a,t)]):+h.c`$, Eq. (III.2), can be rewritten as
$`T_2^q=:\mathrm{\Gamma }_2\kappa _2\mathrm{exp}(i\sqrt{\nu }[\varphi _1(0,t)\varphi _2(0,t)]+i\sqrt{\nu }{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{^k\varphi _1}{x^k}}L^k{\displaystyle \frac{^k\varphi _2}{x^k}}(L+a)^k]/k!):+h.c.=`$
$`:\mathrm{\Gamma }_2\kappa _2\mathrm{exp}(i\sqrt{\nu }[\varphi _1(0,t)\varphi _2(0,t)])\{1+i\sqrt{\nu }[{\displaystyle \frac{\varphi _1}{x}}L{\displaystyle \frac{\varphi _2}{x}}(L+a)]+\mathrm{}\}:+h.c.,`$ (40)
where the ellipses denote higher order gradients of $`\varphi _l`$ and the colons denote normal ordering. After the substitution of Eq. (V.2) in the perturbative expansion of the current one can compare the contributions from the terms containing derivatives of $`\varphi _l`$ with the contributions from the terms which do not contain such derivatives. Power counting shows that the terms with derivatives are suppressed by the factors of order $`(LeV/\mathrm{}v)^k;([L+a]eV/\mathrm{}v)^k`$ and hence can be neglected. This conclusion agrees with the results of section V.C for arbitrary voltages, temperatures and interferometer sizes.
Thus, we can use an effective single impurity model. The tunneling operator $`O=(T_1^q+T_2^q)`$, Eqs. (III.2,24), assumes the form
$$O=O_{}+O_+;O_{}=O_+^{}=\kappa _{}A_{},$$
(41)
where
$$A_{}=A_+^{}=\mathrm{exp}(i\sqrt{\nu }[\varphi _1(0,t)\varphi _2(0,t)])$$
(42)
and the (non-unitary) operator $`\kappa _{}`$ is
$$\kappa _{}=\kappa _+^{}=\mathrm{\Gamma }_1\kappa _1+\mathrm{\Gamma }_2\mathrm{exp}(2\pi \nu i\mathrm{\Phi }/\mathrm{\Phi }_0)\kappa _2=\left(\begin{array}{ccc}0& C_1& 0\\ 0& 0& C_2\\ C_3& 0& 0\end{array}\right)$$
(43)
with $`C_1=\mathrm{\Gamma }_1+\mathrm{\Gamma }_2\mathrm{exp}(2\pi \nu i\mathrm{\Phi }/\mathrm{\Phi }_0)\psi `$, $`C_2=\mathrm{\Gamma }_1+\mathrm{\Gamma }_2\mathrm{exp}(2\pi \nu i\mathrm{\Phi }/\mathrm{\Phi }_0)\psi ^2`$, $`C_3=\mathrm{\Gamma }_1+\mathrm{\Gamma }_2\mathrm{exp}(2\pi \nu i\mathrm{\Phi }/\mathrm{\Phi }_0)`$ and $`\psi =\mathrm{exp}(2\pi i/3)`$. In what follows we will denote the basis vectors $`(1,0,0)`$, $`(0,1,0)`$ and $`(0,0,1)`$ as $`1|`$, $`2|`$ and $`3|`$ respectively. Since we neglect $`L`$ and $`(L+a)`$ in the rest of this section, we can use the correlation function (V.2) in the simplest limit $`b,c=0`$.
The contribution $`I_{2N}`$ of order $`2N`$ to the current corresponds to the charge configuration with $`N`$ dipoles. One dipole of size $`\tau _0`$ is located at $`t=0`$. The remaining dipoles of sizes $`\tau _1,\mathrm{},\tau _{N1}`$ are located at $`0>t_1>t_2>\mathrm{}>t_{N1}`$. One finds
$`I_{2N}={\displaystyle \underset{b_0=\pm 1}{}}{\displaystyle \underset{b_k^\pm =\pm 1}{}}{\displaystyle \underset{\sigma [0]=\pm }{}}\sigma [0]{\displaystyle \frac{\nu ei^{2N}}{\mathrm{}^{2N}}}{\displaystyle _{\mathrm{}}^0}𝑑t_1{\displaystyle _{\mathrm{}}^{t_1}}𝑑t_2\mathrm{}{\displaystyle _{\mathrm{}}^{t_{N2}}}𝑑t_{N1}{\displaystyle _{\mathrm{}}^0}𝑑\tau _0{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\tau _1\mathrm{}𝑑\tau _{N1}b_0\mathrm{\Pi }_kb_k^+b_k^{}`$
$`\mathrm{Tr}\left[\widehat{\rho }\mathrm{T}_cO_{\sigma [0]}(t=0)O_{\sigma [0]}(\tau _0;b_0)O_{}(t_1+{\displaystyle \frac{\tau _1}{2}};b_1^{})O_+(t_1{\displaystyle \frac{\tau _1}{2}};b_1^+)\mathrm{}O_{}(t_{N1}+{\displaystyle \frac{\tau _{N1}}{2}};b_{N1}^{})O_+(t_{N1}{\displaystyle \frac{\tau _{N1}}{2}};b_{N1}^+)\right],`$ (44)
where $`\mathrm{T}_c`$ denotes time-ordering along the Keldysh contour, $`b_0,b_k^+,b_k^{}=+1`$ correspond to the bottom branch of the Keldysh contour and $`b_0,b_k^+,b_k^{}=1`$ correspond to the top branch.
The trace in Eq. (V.2) can be factorized as
$$\mathrm{Tr}[\mathrm{}]=\mathrm{\Pi }_A\mathrm{\Pi }_\kappa ,$$
(45)
where $`\mathrm{\Pi }_A`$ stays for the trace of a product of operators $`A_\pm `$, and $`\mathrm{\Pi }_\kappa `$ denotes the trace of a time-ordered product of the Klein factors. The former trace can be further factorized into a product of the two-point correlation functions (V.2) corresponding to each dipole. We next simplify the expression for $`\mathrm{\Pi }_\kappa `$,
$$\mathrm{\Pi }_\kappa =\mathrm{Tr}\left[\widehat{\rho }_\kappa \mathrm{T}_c\kappa _{\sigma [0]}(t=0)\kappa _{\sigma [0]}(\tau _0;b_0)\kappa _{}(t_1+\frac{\tau _1}{2};b_1^{})\kappa _+(t_1\frac{\tau _1}{2};b_1^+)\mathrm{}\kappa _{}(t_{N1}+\frac{\tau _{N1}}{2};b_{N1}^{})\kappa _+(t_{N1}\frac{\tau _{N1}}{2};b_{N1}^+)\right].$$
(46)
Using the matrix elements of the Klein factors the above equation can be rewritten as
$$\mathrm{\Pi }_\kappa =\frac{1}{3}\underset{e_N=1,2,3}{}\underset{e_k^{t/b}=1,2,3;k=1,\mathrm{},N1}{}I(e_1^t,e_1^b)B(e_1^t,e_1^b;e_2^t,e_2^b)\mathrm{}B(e_{N2}^t,e_{N2}^b;e_{N1}^t,e_{N1}^b)B(e_{N1}^t,e_{N1}^b;e_N^t=e_N,e_N^b=e_N),$$
(47)
where the factor $`1/3`$ comes from the condition $`\mathrm{Tr}\widehat{\rho }_\kappa =1`$; $`I(e_1^t,e_1^b)=e_1^b|T_c\kappa _{\sigma [0]}(t=0)\kappa _{\sigma [0]}(\tau _0;b_0)|e_1^t`$ and
$`B(e_k^t,e_k^b;e_{k+1}^t,e_{k+1}^b)=e_{k+1}^b|e_k^be_k^t|T_c\kappa _{}(t_k+\tau _k/2;1)\kappa _+(t_k\tau _k/2;1)|e_{k+1}^t\mathrm{for}b_k^+=b_k^{}=1`$
$`B(e_k^t,e_k^b;e_{k+1}^t,e_{k+1}^b)=e_{k+1}^b|T_c\kappa _{}(t_k+\tau _k/2;+1)\kappa _+(t_k\tau _k/2;+1)|e_k^be_k^t|e_{k+1}^t\mathrm{for}b_k^+=b_k^{}=+1`$
$`B(e_k^t,e_k^b;e_{k+1}^t,e_{k+1}^b)=e_{k+1}^b|\kappa _+(t_k\tau _k/2;+1)|e_k^be_k^t|\kappa _{}(t_k+\tau _k/2;1)|e_{k+1}^t\mathrm{for}b_k^+=+1;b_k^{}=1`$
$`B(e_k^t,e_k^b;e_{k+1}^t,e_{k+1}^b)=e_{k+1}^b|\kappa _{}(t_k+\tau _k/2;+1)|e_k^be_k^t|\kappa _+(t_k\tau _k/2;1)|e_{k+1}^t\mathrm{for}b_k^+=1;b_k^{}=+1`$ (48)
It is clear from Eqs. (V.2) that for equal $`e_{k+1}^t=e_{k+1}^b=e_{k+1}`$, the expression $`B(e_k^t,e_k^b;e_{k+1},e_{k+1})`$ is nonzero only if $`e_k^t=e_k^b`$. Hence, Eq. (47) can be represented as a matrix product
$$\mathrm{\Pi }_\kappa =\underset{e_k}{}\stackrel{~}{I}(e_1)\stackrel{~}{B}(e_1,e_2)\mathrm{}\stackrel{~}{B}(e_{N2},e_{N1})\stackrel{~}{B}(e_{N1},e_N)\rho (e_N),$$
(49)
where $`\rho (e_N)=1/3`$, $`\stackrel{~}{B}(e_k,e_{k+1})=B(e_k,e_k;e_{k+1},e_{k+1})`$ and $`\stackrel{~}{I}(e_1)=I(e_1,e_1)`$. After the substitution of Eq. (49) in Eqs. (45) and (V.2), each matrix element $`\stackrel{~}{B}(e_k,e_{k+1})`$ multiplies by the correlation function $`i^2T_cb_k^{}A_{}(t_k+\tau _k/2;b_k^{})b_k^+A_+(t_k\tau _k/2;b_k^+)`$. The product should be integrated over $`d\tau _k`$. One finally obtains
$$I_{2N}=_{\mathrm{}}^0𝑑t_1_{\mathrm{}}^{t_1}𝑑t_2\mathrm{}_{\mathrm{}}^{t_{N2}}𝑑t_{N1}I|\widehat{D}^{N1}|\rho ,$$
(50)
where
$$\rho |=(1/3,1/3,1/3);$$
(51)
$`I|={\displaystyle \frac{2e\nu }{\mathrm{}^2}}{\displaystyle \underset{k=1,2,3}{}}\mathrm{Re}{\displaystyle _{\mathrm{}}^0}𝑑t\left[k|\kappa _+\kappa _{}|k{\displaystyle \frac{\mathrm{exp}(i\nu eVt/\mathrm{})\tau _c^{2\nu }}{(\delta +it)^{2\nu }}}k|\kappa _{}\kappa _+|k{\displaystyle \frac{\mathrm{exp}(i\nu eVt/\mathrm{})\tau _c^{2\nu }}{(\delta it)^{2\nu }}}\right]k|`$
$`={\displaystyle \frac{e\mathrm{\Gamma }(1/3)\tau _c^{2/3}}{\sqrt{3}\mathrm{}^2}}\left({\displaystyle \frac{eV}{3\mathrm{}}}\right)^{1/3}\times (|C_3|^2,|C_1|^2,|C_2|^2);`$ (52)
$$\widehat{D}=\frac{\sqrt{3}\mathrm{\Gamma }(1/3)\tau _c^{2/3}}{\mathrm{}^2}\left(\frac{eV}{3\mathrm{}}\right)^{1/3}\left(\begin{array}{ccc}|C_3|^2& |C_1|^2& 0\\ 0& |C_1|^2& |C_2|^2\\ |C_3|^2& 0& |C_2|^2\end{array}\right).$$
(53)
The total current
$$I=I_{2N}=I|\mathrm{exp}(\overline{t}\widehat{D})|\rho ,$$
(54)
where $`\overline{t}=+\mathrm{}`$ is the length of the Keldysh contour.
All elements of the matrix $`\widehat{D}`$ are real, all diagonal elements are negative, all nondiagonal elements are positive or zero and the sum of the elements in each column is zero. The Rohrbach theorem Rohrbach applies to such matrices. Since $`\rho |\widehat{D}=0`$, $`0`$ is an eigenvalue of the matrix $`\widehat{D}`$. According to the Rohrbach theorem, this eigenvalue is non-degenerate and the real parts of all other eigenvalues are negative. This allows for a simple calculation of $`\mathrm{exp}(\overline{t}\widehat{D})`$. Let $`\widehat{S}`$ be such a matrix that $`\widehat{\stackrel{~}{D}}=\widehat{S}\widehat{D}\widehat{S}^1`$ assumes the Jordan normal form. Without the loss of generality we can assume that the first column and the first string of the matrix $`\widehat{\stackrel{~}{D}}`$ are zero. Then the first string of the matrix $`\widehat{S}`$ can be chosen in the form $`(1,1,1)`$. Thus, the matrix exponent
$$\mathrm{exp}(\overline{t}\widehat{D})=\widehat{S}^1\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)\widehat{S}=\widehat{S}^1\left(\begin{array}{ccc}1& 1& 1\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)$$
(55)
Hence,
$$I=I|\widehat{S}^1|1=I|S_1^1,$$
(56)
where $`|S_1^1`$ denotes the first column of the matrix $`\widehat{S}^1`$.
In order to complete our calculation we have to determine the components $`f_1`$, $`f_2`$, and $`f_3`$ of the vector $`|S_1^1`$. From the condition $`\widehat{S}\widehat{S}^1=\widehat{E}`$ one finds that
$$f_1+f_2+f_3=1.$$
(57)
We also know that $`|S_1^1`$ is the eigenvector of $`\widehat{D}`$ with zero eigenvalue, i.e. $`\widehat{D}|S_1^1=0`$. Hence, for each $`k`$,
$$|C_{k1}|^2f_k=|C_k|^2f_{k+1},$$
(58)
where we use the convention $`3+1=1`$. The solution of the above equation is $`f_k=\alpha /|C_{k1}|^2`$, where $`\alpha =\frac{1}{1/|C_1|^2+1/|C_2|^2+1/|C_3|^2}`$ can be found from Eq. (57). Finally,
$$I=\frac{e\mathrm{\Gamma }(1/3)\tau _c^{2/3}}{\sqrt{3}\mathrm{}^2}\left(\frac{eV}{3\mathrm{}}\right)^{1/3}\frac{3}{\frac{1}{|C_1|^2}+\frac{1}{|C_2|^2}+\frac{1}{|C_3|^2}}.$$
(59)
This expression is equal to the harmonic average of the three tunneling currents in three systems with a single quantum point contact with the tunneling amplitude $`C_1`$, $`C_2`$ and $`C_3`$ respectively. In other words, Eq. (59) is equivalent to (4). Using the result of Appendix C, Eq. (59) can be represented as
$$I=\frac{e\mathrm{\Gamma }(1/3)\tau _c^{2/3}}{\sqrt{3}\mathrm{}^2}\left(\frac{eV}{3\mathrm{}}\right)^{1/3}\left[|\mathrm{\Gamma }_1|^6+|\mathrm{\Gamma }_2|^6+2|\mathrm{\Gamma }_1\mathrm{\Gamma }_2|^3\mathrm{cos}(3\alpha _0+2\pi \mathrm{\Phi }/\mathrm{\Phi }_0)\right]\frac{|\mathrm{\Gamma }_1|^2|\mathrm{\Gamma }_2|^2}{|\mathrm{\Gamma }_1|^6|\mathrm{\Gamma }_2|^6},$$
(60)
where $`\alpha _0=\mathrm{arg}[\mathrm{\Gamma }_2/\mathrm{\Gamma }_1]`$. This allows one to easily verify Eq. (7) in the case $`\nu =1/3`$, $`T=0`$, $`eVL,eVa\mathrm{}v`$.
### V.3 Tunneling current for arbitrary filling factors, temperatures and voltages
The calculations follow the same route as in the previous subsection. The contribution to the current of order $`2N`$, $`I_{2N}`$, expresses via the trace $`\mathrm{\Pi }`$ of the time-ordered product of $`2N`$ tunneling operators which form $`N`$ dipoles. The trace must be integrated over the size of each dipole and the positions of $`N1`$ dipoles, the remaining dipole being located at $`t=0`$. The trace $`\mathrm{\Pi }`$ factorizes as the product $`\mathrm{\Pi }=\mathrm{\Pi }_\kappa \mathrm{\Pi }_\varphi `$, where $`\mathrm{\Pi }_\kappa `$ stays for the trace of the product of the Klein factors $`\kappa _1`$, $`\kappa _2`$, $`\kappa _1^+`$ and $`\kappa _2^+`$; $`\mathrm{\Pi }_\varphi `$ stays for the trace of the product of the operators $`\mathrm{exp}(\pm i\sqrt{\nu }[\varphi _1(x=0)\varphi _2(x=0)])`$ and $`\mathrm{exp}(\pm i\sqrt{\nu }[\varphi _1(L)\varphi _2(L+a)])`$. The latter trace factorizes in the product of the two-point correlation functions bos corresponding to each dipole:
$`F(b,c,t_1t_2)=\mathrm{Tr}[\widehat{\rho }\mathrm{exp}(i\sqrt{\nu }\varphi _1(x=b,t_1))\mathrm{exp}(i\sqrt{\nu }\varphi _1(x=0,t_2))]\mathrm{Tr}[\widehat{\rho }\mathrm{exp}(i\sqrt{\nu }\varphi _2(x=c,t_1))\mathrm{exp}(i\sqrt{\nu }\varphi _2(x=0,t_2))]=`$
$`\left[{\displaystyle \frac{\pi k_BT\tau _c/\mathrm{}}{\mathrm{sin}(\pi k_BT[\delta +i(t_1t_2b/v)]/\mathrm{})}}\right]^\nu \left[{\displaystyle \frac{\pi k_BT\tau _c/\mathrm{}}{\mathrm{sin}(\pi k_BT[\delta +i(t_1t_2c/v)]/\mathrm{})}}\right]^\nu `$ (61)
The former trace $`\mathrm{\Pi }_\kappa `$ can be represented in the form similar to Eq. (49), where $`\stackrel{~}{B}(e_k,e_{k+1})`$ express via matrix elements of $`\kappa _1`$ and $`\kappa _2`$. Next, one can rewrite $`I_{2N}`$ in the form (50) with modified definitions of $`|\rho `$, $`I|`$ and $`\widehat{D}`$. $`|\rho `$ and $`|I`$ are now $`1/\nu `$-dimensional vectors; $`\widehat{D}`$ is a matrix of size $`1/\nu \times 1/\nu `$. Similar to the previous subsection, in order to obtain the matrix elements of $`\widehat{D}`$, one has to multiply $`\stackrel{~}{B}(e_k,e_{k+1})`$ by the correlation function (V.3) describing the dipole located at $`t=t_k`$ and then integrate the product over the dipole size. One finds
$$\rho |=(\nu ,\nu ,\mathrm{},\nu );$$
(62)
$$I|e_k=e\nu (I_{k(V)}I_{k(V)});$$
(63)
$$e_l|\widehat{D}|e_k=\delta _{k,l}(I_{k(V)}+I_{k(V)})\delta _{k,l1}I_{k(V)}\delta _{k,l+1}I_{k(V)},$$
(64)
where
$$I_{k(V)}=\frac{1}{\mathrm{}^2}\left[(|\mathrm{\Gamma }_1|^2+|\mathrm{\Gamma }_2|^2)j(V;0)+\mathrm{\Gamma }_1\mathrm{\Gamma }_2^{}(\psi ^{})^kj(V;a)+\mathrm{\Gamma }_1^{}\mathrm{\Gamma }_2\psi ^kj(V;a)\right];$$
(65)
$$I_{k(V)}=\frac{1}{\mathrm{}^2}\left[(|\mathrm{\Gamma }_1|^2+|\mathrm{\Gamma }_2|^2)j(V;0)+\mathrm{\Gamma }_1\mathrm{\Gamma }_2^{}(\psi ^{})^{k1}j(V;a)+\mathrm{\Gamma }_1^{}\mathrm{\Gamma }_2\psi ^{k1}j(V;a)\right],$$
(66)
where $`\psi =\mathrm{exp}(2\pi i\nu )`$, the phase factor $`\mathrm{exp}(2\pi i\nu \mathrm{\Phi }/\mathrm{\Phi }_0+ie\nu VL/\mathrm{}v)`$ should be included in $`\mathrm{\Gamma }_2`$ and
$`j(U,0)={\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑t\mathrm{exp}(i\nu eUt/\mathrm{})F(0,0,t);`$
$`j(U,\pm a)={\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑t\mathrm{exp}(i\nu eUt/\mathrm{})F(\pm L,\pm (L+a),t)`$ (67)
The total current is given by Eq. (54) of the previous section with the above definitions of $`I|`$, $`\widehat{D}`$ and $`|\rho `$.
One can easily check that $`j(U,x)=j^{}(U,x)`$. Hence, $`I_{k(\pm V)}`$ are real. Appendix D shows that
$$I_{k(V)}=\mathrm{exp}\left(\frac{\nu eV}{k_BT}\right)I_{(k+1)(V)}=\gamma I_{(k+1)(V)}$$
(68)
that is equivalent to the detailed balance condition (8). A comparison with the geometry Fig. 3, Ref. cfksw, , shows that the same integral which expresses the current $`I_{\mathrm{Fig}.3}`$ in the geometry Fig. 3 also equals the expression
$$I_{k(V)}I_{(k1)(V)}=(1\gamma )I_{k(V)}=I_{\mathrm{Fig}.3}^{(k)}/(\nu e)=\frac{2^{2g}\tau _c}{\mathrm{}^2}|\mathrm{\Gamma }_{\mathrm{eff}}|^2\left(\frac{\pi k_BT\tau _c}{\mathrm{}}\right)^{2\nu 1}\frac{|\mathrm{\Gamma }(\nu +\frac{i\nu eV}{2\pi k_BT})|^2}{\mathrm{\Gamma }(2\nu )}\mathrm{sinh}\frac{\nu eV}{2k_BT},$$
(69)
where
$`|\mathrm{\Gamma }_{\mathrm{eff}}|^2=2\pi {\displaystyle \frac{\mathrm{\Gamma }(2\nu )}{\mathrm{\Gamma }(\nu )}}{\displaystyle \underset{m,l}{}}\mathrm{\Gamma }_m\mathrm{\Gamma }_l^{}\mathrm{exp}(i[2\pi \nu \mathrm{\Phi }/\mathrm{\Phi }_0e\nu Va/(2\mathrm{}v)2\pi (k1)\nu ][\delta _{m,2}\delta _{l,2}])\times `$
$`{\displaystyle \frac{\mathrm{exp}(\nu \pi ak_BT/[\mathrm{}v])}{\mathrm{sinh}(eV\nu /2k_BT)}}\mathrm{Im}\left\{{\displaystyle \frac{\mathrm{exp}(ie\nu Va/[2\mathrm{}v])F(\nu ,\nu ie\nu V/2\pi k_BT;1ie\nu V/2\pi k_BT;e^{2\pi k_BTa/[\mathrm{}v]})}{\mathrm{\Gamma }(\nu +ie\nu V/2\pi k_BT)\mathrm{\Gamma }(1ie\nu V/2\pi k_BT)}}\right\},`$ (70)
$`F`$ being the hypergeometric function cfksw .
Since the particle current flows from the edge with the higher chemical potential to the edge with the lower potential, we know the sign of $`I_{\mathrm{Fig}.3}^{(k)}=dQ_1/dt>0`$ (electron charge $`e<0`$). Hence, $`I_{k(\pm V)}<0`$. One can also see that the sum of the elements in any column of the matrix $`\widehat{D}`$, Eq. (64), is zero. Thus, the Rohrbach theorem Rohrbach applies again and hence $`\mathrm{exp}(\overline{t}\widehat{D})=\widehat{S}^1\widehat{M}\widehat{S}`$, where the matrix $`\widehat{S}`$ reduces $`\widehat{D}`$ to the Jordan normal form, the matrix $`\widehat{M}`$ has only one non-zero element in its upper left corner and this element is equal to 1 just like in the preceding subsection. All elements of the first string of the matrix $`\widehat{S}`$ are equal to 1. This allows us to obtain the expression for the current in the form (56) following exactly the same steps as in section V.B.
The only thing left is the calculation of the components $`f_r`$ of the vector $`|S_1^1`$, i.e. the first column of the matrix $`\widehat{S}^1`$. From the condition $`\widehat{S}\widehat{S}^1=E`$ one finds
$$\underset{r=1}{\overset{1/\nu }{}}f_r=1.$$
(71)
From the condition $`\widehat{D}|S_1^1=0`$ one finds
$$I_{k(V)}f_kI_{(k1)(V)}f_{k1}=I_{(k+1)(V)}f_{k+1}I_{k(V)}f_k.$$
(72)
At the same time the current (63,56) is
$$I=\nu e\underset{k}{}(I_{k(V)}I_{k(V)})f_k$$
(73)
The system of equations (71-73) together with the detailed balance condition (68) is equivalent to the system (8-10) of section II with $`I_{k(\pm V)}`$ playing the role of the transition probabilities $`p`$. The components $`f_r`$ of $`|S_1^1`$ have the physical meaning of the distribution function. Thus, we can directly use the solution (14):
$$I=\frac{1/\nu }{_{k=1}^{1/\nu }I_k^{}[V,T,a]},$$
(74)
where the current $`I_k^{}=I_{\mathrm{Fig}.3}^{(k)}`$, Eq. (69), equals the tunneling current in the geometry Fig. 3 with the tunneling amplitudes $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2\mathrm{exp}(i[2\pi \nu \mathrm{\Phi }/\mathrm{\Phi }_0e\nu Va/(2\mathrm{}v)2\pi (k1)\nu ])`$.
In the limit $`T=0;aeVhv`$ the above result reduces to a simple generalization of Eq. (60):
$$I=\frac{2e\nu \mathrm{sin}(2\pi \nu )\mathrm{\Gamma }(12\nu )\tau _c^{2\nu }}{\mathrm{}^2}\left(\frac{eV\nu }{\mathrm{}}\right)^{2\nu 1}\frac{|\mathrm{\Gamma }_1|^2|\mathrm{\Gamma }_2|^2}{|\mathrm{\Gamma }_1|^{2/\nu }|\mathrm{\Gamma }_2|^{2/\nu }}\left[|\mathrm{\Gamma }_1|^{2/\nu }+|\mathrm{\Gamma }_2|^{2/\nu }+2|\mathrm{\Gamma }_1|^{1/\nu }|\mathrm{\Gamma }_2|^{1/\nu }\mathrm{cos}(2\pi \mathrm{\Phi }/\mathrm{\Phi }_0+\alpha _0/\nu )\right].$$
(75)
If $`|\mathrm{\Gamma }_1||\mathrm{\Gamma }_2|`$ the current never vanishes. At $`|\mathrm{\Gamma }_1|=|\mathrm{\Gamma }_2|`$ a “resonance” is reached when $`\mathrm{\Phi }=[n+1/2\alpha _0/(2\pi \nu )]\mathrm{\Phi }_0`$ and $`I=0`$.
## VI Conclusions
We have found that the tunneling current through the electronic Mach-Zehnder interferometer is a periodic function of the magnetic flux with period $`\mathrm{\Phi }_0`$. This result is valid both in the weak quasiparticle tunneling regime and in the weak electron tunneling regime. The relations between the flux-dependent and flux-independent components of the current, $`I_\mathrm{\Phi }`$ and $`I_0`$, are different in these respective regimes. In the electron tunneling case, $`I_\mathrm{\Phi }(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)[I_0(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)I_0(\mathrm{\Gamma }_1,0)]^{1/2}`$ at low voltages and temperatures, $`\mathrm{\Gamma }_2\mathrm{\Gamma }_1`$. In the quasiparticle tunneling case the flux-dependent contribution scales as $`I_\mathrm{\Phi }(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)[I_0(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2)I_0(\mathrm{\Gamma }_1,0)]^b`$, where $`b>1`$. The exponent in this power law contains information about quasiparticle statistics since the exponent derives from the algebra of the Klein factors. Thus, the Mach-Zehnder interferometer can be used to probe fractional statistics. Recently an interference pattern has been observed experimentally for the integer quantum Hall case in a Mach-Zehnder interferometer MZ . Higher magnetic fields would allow an investigation of a fractional quantum Hall liquid.
###### Acknowledgements.
We thank J. Chalker, M. P. A. Fisher, L. I. Glazman, M. Heiblum, I. Neder, X.-G. Wen and P. B. Wiegmann for useful discussions. DEF thanks KITP where a part of this work was completed for hospitality. This research was supported in part by the National Science Foundation under Grant No. PHY99-07949, by the U.S.-Israel Binational Science Foundation, and the ISF of the Israel Academy of Sciences.
## Appendix A Electron tunneling at low voltages and zero temperature
In this appendix we verify that the current, Eqs. (30-IV), satisfies the asymptotics (33) at $`V0`$. The flux-independent contribution to the current (31) scales as $`V^{g_1+g_21}`$, where $`g_1=1/\nu _1`$ and $`g_2=1/\nu _2`$. Hence, we need to check that the flux-dependent contribution $`I_\mathrm{\Phi }`$ does not exceed $`\mathrm{const}V^{g_1+g_21}`$ at low voltages. Expanding $`\mathrm{exp}(\frac{ieVz}{\mathrm{}})`$ in (IV) in powers of $`ieVz/\mathrm{}`$ one finds that the coefficient before $`V^k`$ in the Taylor expansion of $`I_\mathrm{\Phi }`$ equals
$$s_k=\mathrm{const}\left[\frac{1}{(g_11)!}\frac{d^{g_11}}{dz^{g_11}}|_{z=0}\frac{z^k}{(z+a/v)^{g_2}}+\frac{1}{(g_21)!}\frac{d^{g_21}}{dz^{g_21}}|_{z=\frac{a}{v}}z^{kg_1}\right].$$
(76)
We want to show that $`s_k=0`$ for all $`k<g_1+g_21`$.
Let us consider separately $`k>g_11`$ and $`kg_11`$.
1) $`k>g_11`$. This inequality can be satisfied for $`g_2>1`$ only since $`k<g_1+g_21`$. In this case $`\frac{d^{g_21}}{dz^{g_21}}|_{z=\frac{a}{v}}z^{kg_1}=0`$ since $`g_21>kg_10`$. Since $`\frac{d^p}{dz^p}|_{z=0}z^k=0`$ at $`p<k`$, one finds $`\frac{d^{g_11}}{dz^{g_11}}|_{z=0}\frac{z^k}{(z+a/v)^{g_2}}=_{p=0}^{g_11}C_{g_11}^p\frac{d^p}{dz^p}z^k\frac{d^{g_11p}}{dz^{g_11p}}\frac{1}{(z+a/v)^{g_2}}=0`$, where $`C_a^b`$ denote binomial coefficients. Thus, $`s_k=0`$ for $`g_1+g_21>k>g_11`$.
2) At $`kg_11`$ one finds
$$s_k=\mathrm{const}\left[\frac{1}{(g_11)!}C_{g_11}^k\frac{k!(1)^{g_11k}}{(a/v)^{g_1+g_2k1}}\frac{(g_1+g_2k2)!}{(g_21)!}+\frac{1}{(g_21)!}(a/v)^{kg_1g_2+1}\frac{(g_1+g_2k2)!}{(g_1k1)!}\right]=0.$$
(77)
Thus, $`s_k=0`$ for all $`k<g_1+g_21`$.
## Appendix B Perturbation theory at finite temperature
We want to calculate the current, Eq. (IV), at a finite temperature, $`k_BT>0`$. The correlation function, Eq. (IV), is bos
$$F(b,c,t)=\left[\frac{\pi k_BT\tau _c/\mathrm{}}{\mathrm{sin}(\pi k_BT[\delta +i(tb/v)]/\mathrm{})}\right]^{1/\nu _1}\left[\frac{\pi k_BT\tau _c/\mathrm{}}{\mathrm{sin}(\pi T[\delta +i(tc/v)]/\mathrm{})}\right]^{1/\nu _2}.$$
(78)
Since $`1/\nu _1=g_1`$ and $`1/\nu _2=g_2`$ are odd integers, $`F(L,(L+a),t)=F(L,L+a,t)`$ on the real axis except in the vicinity of the two real poles $`t=L/v`$ and $`t=(L+a)/v`$. $`F(0,0,t)=F(0,0,t)`$ except near $`t=0`$. Hence, the integral, Eq. (IV), reduces to the sum of the contour integrals along small circles around the real poles. A straightforward calculation yields
$$I=\frac{e\tau _c}{\mathrm{}^2}\left(\frac{\pi k_BT\tau _c}{\mathrm{}}\right)^{g1}(|\mathrm{\Gamma }_1|^2+|\mathrm{\Gamma }_2|^2)J_0+\frac{e\tau _c}{\mathrm{}^2}\left(\frac{\pi k_BT\tau _c}{\mathrm{}}\right)^{g1}[\mathrm{\Gamma }_1\mathrm{\Gamma }_2^{}\mathrm{exp}(2\pi i\mathrm{\Phi }/\mathrm{\Phi }_0)J_\mathrm{\Phi }+c.c.],$$
(79)
where $`g=g_1+g_2=1/\nu _1+1/\nu _2`$, and
$$J_0=\frac{2\pi }{i^{g+1}(g1)!}\frac{d^{g1}}{dz^{g1}}|_{z=0}\frac{\mathrm{exp}(ieVz/\pi k_BT)z^g}{\mathrm{sinh}^gz},$$
(80)
$$J_\mathrm{\Phi }=\frac{2\pi }{i^{g+1}(g_11)!}\frac{d^{g_11}}{dz^{g_11}}|_{z=0}\frac{\mathrm{exp}(\frac{ieVz}{\pi k_BT})z^{g_1}}{(\mathrm{sinh}z)^{g_1}(\mathrm{sinh}[z\frac{\pi ak_BT}{\mathrm{}v}])^{g_2}}+\frac{2\pi }{i^{g+1}(g_21)!}\frac{d^{g_21}}{dz^{g_21}}|_{z=\frac{\pi ak_BT}{\mathrm{}v}}\frac{\mathrm{exp}(\frac{ieVz}{\pi k_BT})(z\frac{\pi ak_BT}{\mathrm{}v})^{g_2}}{(\mathrm{sinh}z)^{g_1}(\mathrm{sinh}[z\frac{\pi ak_BT}{\mathrm{}v}])^{g_2}}$$
(81)
Note that the temperature enters the above expression in the combination $`ak_BT`$ but not in the combination $`Lk_BT`$. This can be understood from the picture of non-interacting electrons. For non-interacting particles the current is the sum of independent contributions from different electron energies. Each contribution depends on the phase difference between the two paths connecting the point contacts but not the total phase accumulated on any of those paths. The former is proportional to $`a`$, the latter to $`L`$.
One can easily extract the asymptotical behavior of the linear conductance $`G=\frac{dI}{dV}`$ at low temperatures (34) from equations (80,81). In the Loran expansion of $`G`$ in powers of $`a`$, each term of order $`a^n`$ is proportional to $`T^{1/\nu _1+1/\nu _22+n}`$ since $`a`$ enters the expression for the current in the combination $`(ak_BT/\mathrm{}v)`$ only. In the limit $`a0`$ the conductance must remain finite. Indeed, at $`L=0`$ this limit corresponds to a problem with a single tunneling contact. Hence, only positive and zero powers of $`a`$ are present in the Loran expansion and the leading contribution to the temperature dependence of the conductance scales as $`T^{1/\nu _1+1/\nu _22}`$.
## Appendix C Effective single impurity model
In the effective single impurity model, section V.B, the current is given by Eq. (59). The purpose of this appendix consists in the calculation of the coefficient in Eq. (59),
$$U=\frac{N}{_{k=1}^N\frac{1}{|C_k|^2}},$$
(82)
where $`C_k=\mathrm{\Gamma }_1+\mathrm{\Gamma }_2\mathrm{exp}(2\pi i\mathrm{\Phi }/[N\mathrm{\Phi }_0]2\pi ik/N)`$, Eq. (43), $`N=1/\nu `$. Without the loss of generality we can assume that $`\mathrm{\Gamma }_1=|\mathrm{\Gamma }_1|`$ is real. Let $`\mathrm{\Gamma }_2=|\mathrm{\Gamma }_2|\mathrm{exp}(i\alpha _0)`$. In the rest of the appendix we will use the notation $`\gamma _1=|\mathrm{\Gamma }_2|`$, $`\gamma _2=|\mathrm{\Gamma }_2|`$ and $`\varphi =2\pi \mathrm{\Phi }/[N\mathrm{\Phi }_0]+\alpha _0`$. $`U`$ can be represented as the ratio of two polynomials of $`\gamma _1`$ and $`\gamma _2`$:
$$U=N\frac{\mathrm{\Pi }_k|\gamma _1+\gamma _2\mathrm{exp}(i\varphi +2\pi ik/N)|^2}{_k\mathrm{\Pi }_{lk}^{}|\gamma _1+\gamma _2\mathrm{exp}(i\varphi +2\pi il/N)|^2},$$
(83)
where the prime after the product sign means that the term with $`l=k`$ is not included in the product. The nominator in Eq. (83) equals $`|P(\gamma _2\mathrm{exp}(i\varphi ))|^2`$, where the polynomial
$$P(z)=\mathrm{\Pi }_k[\gamma _1+z\mathrm{exp}(2\pi ik/N)]$$
(84)
All roots of the polynomial (84) coincide with the roots of the polynomial $`\gamma _1^N+z^N`$. Hence, from the basic theorem of algebra
$$P(z)=\gamma _1^N+z^N$$
(85)
and the nominator is
$$|P(\gamma _2\mathrm{exp}(i\varphi ))|^2=(\gamma _1^N+\gamma _2^N\mathrm{exp}(iN\varphi ))(\gamma _1^N+\gamma _2^N\mathrm{exp}(iN\varphi ))=\gamma _1^{2N}+\gamma _2^{2N}+2\gamma _1^N\gamma _2^N\mathrm{cos}N\varphi .$$
(86)
The denominator in Eq. (83) can be represented as
$`d={\displaystyle \underset{k}{}}\left|{\displaystyle \frac{P(\gamma _2\mathrm{exp}(i\varphi ))}{\gamma _1+\gamma _2\mathrm{exp}(i\varphi +2\pi ik/N)}}\right|^2={\displaystyle \underset{k}{}}|{\displaystyle \underset{p=0}{\overset{N1}{}}}(1)^p\gamma _1^{N1p}\gamma _2^p\mathrm{exp}(ip\varphi +2\pi pki/N)|^2=`$
$`{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle \underset{p=0}{\overset{N1}{}}}{\displaystyle \underset{r=0}{\overset{N1}{}}}\gamma _1^{N1p}\gamma _2^p(1)^p\mathrm{exp}(ip\varphi +2\pi kpi/N)\times \gamma _1^{N1r}\gamma _2^r(1)^r\mathrm{exp}(ir\varphi 2\pi kri/N).`$ (87)
The sum $`_{k=0}^{N1}\mathrm{exp}(2\pi k(pr)i/N)=0`$, if $`prnN`$. With the help of this property, Eq. (C) reduces to
$$d=N\underset{p=0}{\overset{N1}{}}\gamma _1^{2(N1)2p}\gamma _2^{2p}=N\frac{\gamma _1^{2N}\gamma _2^{2N}}{\gamma _1^2\gamma _2^2}.$$
(88)
Finally, the combination of Eqs. (86) and (88) yields
$$U=\frac{\gamma _1^2\gamma _2^2}{\gamma _1^{2N}\gamma _2^{2N}}\left[\gamma _1^{2N}+\gamma _2^{2N}+2\gamma _1^N\gamma _2^N\mathrm{cos}N\varphi \right]$$
(89)
## Appendix D Detailed balance
Here we derive Eq. (68). We need to show that
$$j(V;0)=\gamma j(V;0);j(V;\pm a)=\gamma j(V;a)$$
(90)
This is equivalent to the equation
$$_{\mathrm{}}^+\mathrm{}𝑑tF(b,c,t)\mathrm{exp}(\frac{ie\nu Vt}{\mathrm{}})=\gamma _{\mathrm{}}^+\mathrm{}𝑑tF(b,c,t)\mathrm{exp}(\frac{ie\nu Vt}{\mathrm{}}),$$
(91)
where $`F`$ is given by Eq. (V.3). Both integrals are taken over contour a) which goes below the real axis in Fig. 8. We can change the sign of $`t`$ in the first integral in Eq. (91). This also changes the integration contour into contour b), Fig. 8 (infinitesimal $`\delta `$ is positive). The integral over contour b) is equal to the integral over contour c). The latter integral equals $`\gamma _{\mathrm{}}^+\mathrm{}𝑑tF(b,c,t)\mathrm{exp}(\frac{ie\nu Vt}{\mathrm{}})`$ indeed.
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# Transport and mixing in the radiation zones of rotating stars: II-Axisymmetric magnetic field
## 1 Motivation
The need for improved stellar structure models, which go beyond the so-called standard model, is now well recognized. The radiation zones can no longer be treated as motionless regions, where no mixing occurs, but one has to include physical processes which lead to the transport of matter and angular momentum. So far the focus has been on the effect of rotation, which generates a thermally driven meridional circulation; by advecting angular momentum, that circulation renders the rotation non-uniform and prone to shear turbulence. This rotational mixing has been formulated by Zahn (1992) and Maeder & Zahn (1998), assuming that the angular velocity varies much less in the horizontal than in the vertical direction (shellular rotation). For rapidly rotating stars, models built along these lines are in much better agreement than standard models; this was first shown by Talon et al. (1997), and confirmed by Maeder and Meynet (2000), Maeder and Meynet (2001), who in addition refined the description for the mass loss by taking into account the latitudinal variation of the wind.
However, these models fail to correctly reproduce the flat rotation profile observed in the solar radiative interior (Matias & Zahn 1997). In solar-type stars, at least once they have been spun down through their magnetized wind, other physical processes are therefore responsible for the transport of angular momentum, and the most plausible candidates are now being investigated. One is the transport by internal gravity waves emitted at the base of the convection zones, associated with turbulence which smoothes the differential rotation; this mechanism was described by Kumar et al. (1999), Talon et al. (2002) and applied by Talon and Charbonnel (2003, 2004) to stars of various types. The other possibility is magnetic torquing, as advocated already by Mestel (1953): he showed that even a very modest field could enforce rigid rotation. In the present paper we shall address the latter effect, and examine how rotational mixing is affected by an axisymmetric field.
Moss (1974) was the first to map the thermally-driven meridional circulation in a star containing an axisymmetric poloidal field; when rotation was added, it was assumed constant and uniform. Various generalizations of this work followed, by Mestel and Moss (1977) and by Moss (1977, 1984, 1985, 1987). The a priori assumption of rigid rotation was relaxed by Tassoul and Tassoul (1986), who claimed that no magnetic field could render the rotation nearly uniform in presence of meridional circulation; however their treatment failed to correctly represent the generation of toroidal field through the shearing of poloidal field, which provides the major feed-back in this problem. This was pointed out by Mestel et al. (1988), who carried out numerical simulations that illustrated how a weak magnetic field could enforce nearly uniform rotation within a few Alfvén times; for sake of simplification, both the poloidal field and the meridional circulation were taken as given, and constant in time.
Recently Garaud (2002) treated the problem in a fully consistent way, and she applied it to the description of the solar tachocline: she let the field adjust through Ohmic diffusion and be advected by the meridional circulation. However, due to numerical limitations, her circulation was driven mainly through Ekman pumping, rather than through thermal diffusion; also, only stationary solutions were constructed.
The purpose of the present paper is to model the time-dependent rotational mixing in presence of an axisymmetric magnetic field, thus extending the results obtained by Mathis & Zahn (2004; from now on referred to as Paper I) in the purely hydrodynamical case. To achieve this goal:
* we allow for a general axisymmetric magnetic field, including all its effects,
* we expand all variables in spherical harmonics, formally to unlimited order,
* we treat explicitely the departures from shellular rotation, in the linear approximation, in order to capture the tachocline(s),
* we filter out the short times, and keep only the time derivative in the heat equation and in that describing the transport of angular momentum.<sup>1</sup><sup>1</sup>1All waves are thus filtered out (acoustic, gravity, inertial, Alfvén); the equation for the transport of angular momentum (eq. 44 below) still formally allows for mixed Alfvén-inertial waves, but the timestep of the evolutionary calculation (which uses an implicit scheme) will in general exceed the typical wave travel time, and suppress these waves.
In a forthcoming paper, we shall generalize the present work to a non-axisymmetric magnetic field.
## 2 Assumptions and expansions
We focus here our attention to magnetic fields of moderate strength, whose Alfvén speed $`V_A`$ does not exceed the rotation velocity $`R\mathrm{\Omega }`$, but is larger than the meridional circulation velocity $`V_M`$. Thus the field will not affect much the hydrostatic balance in the meridional plane, but it may play a major role in the transport of angular momentum, by tending to render the angular velocity $`\mathrm{\Omega }`$ constant along the field lines of the poloidal field (Ferraro 1937). We also assume that the field does not vary on a timescale which is shorter than the Alfvén time $`R/V_A`$; thus we shall deal with fields of primordial origin, and exclude possible dynamo fields produced in the adjacent convection zones, whose penetration would be damped anyway by skin effect (Garaud 1999).
Our treatment is not able to address the question whether the field is stable or not, because it is axisymmetric, and because we do not resolve the Alfvén waves. We will have to rely on our intuition, guided by 3D simulations such as recently performed by Braithwaite and Spruit (2004), to choose the magnetic configurations we implement in our code. As for instabilities which may lead to a turbulent state, for instance the magneto-rotational instability when the equator rotates slower than the poles (Balbus & Hawley 1994), their transport properties will have to be accounted for by a suitable parametrization, as was done previously in the hydrodynamical case.
As in Paper I, we assume that the rotating star is only weakly two-dimensional, and this for two reasons. The first is that the rotation rate and the magnetic field, which is taken axisymmetric, are sufficiently moderate to allow the centrifugal and the Lorentz forces to be considered as perturbations compared to gravity. The second reason rests on the less justified hypothesis that the shear instabilities due to the differential rotation give rise to turbulent motions which are strongly anisotropic due to the stable stratification, with much stronger transport in the horizontal directions than in the vertical. In the radiation zones, we expect it to smooth out the horizontal variations of angular velocity and of chemical composition, a property we shall use to discard certain non-linear terms. Let us emphasize also that the influence of the magnetic field on such turbulence is not taken into account, as seems to be allowed by the condition $`V_A<R\mathrm{\Omega }`$.
Hence, we consider an axisymmetric star, and assume that the horizontal variations of all quantities are small and smooth enough to allow their linearization and their expansion in a modest number of spherical harmonics. As reference surface, we choose either the sphere or the isobar, and write all scalar quantities either as:
$$X(r,\theta )=X_0(r)+\underset{l}{}\widehat{X}_l(r)P_l\left(\mathrm{cos}\theta \right)$$
(1)
or
$$X(P,\theta )=\overline{X}(P)+\underset{l>0}{}\stackrel{~}{X}_l(P)P_l\left(\mathrm{cos}\theta \right),$$
(2)
where $`\stackrel{~}{X}_l(r)`$ and $`\widehat{X}_l(r)`$ are related by (cf. Paper I, §2)
$$\stackrel{~}{X}_l(r)=\widehat{X}_l(r)\left(\frac{\mathrm{d}X_0}{\mathrm{d}P_0}\right)\widehat{P}_l(r),$$
(3)
with $`r`$ being here the mean radius of an isobar.
Likewise, we expand all axisymmetric vector fields using the method outlined by Rieutord (1987):
$$u(r,\theta )=\underset{l=0}{\overset{\mathrm{}}{}}\left\{u_0^l(r)R_l^0(\theta )+v_0^l(r)S_l^0(\theta )+w_0^l(r)T_l^0(\theta )\right\}.$$
(4)
The orthogonal axisymmetric vectorial spherical harmonics $`R_l^0\left(\theta \right)`$, $`S_l^0\left(\theta \right)`$, $`T_l^0\left(\theta \right)`$ are defined by:
$$R_l^0(\theta )=Y_l^0(\theta )\widehat{e}_r\text{}S_l^0(\theta )=_𝒮Y_l^0(\theta )\text{ and }T_l^0(\theta )=_𝒮R_l^0(\theta )$$
(5)
where $`Y_l^0\left(\theta \right)`$ is the classical spherical harmonic with $`m=0`$ (cf. Edmonds 1974) and $`_𝒮`$ is the horizontal gradient
$$_𝒮=\widehat{e}_\theta _\theta +\widehat{e}_\phi \frac{1}{\mathrm{sin}\theta }_\phi .$$
(6)
Their detailed properties are given in Appendix A.
Using (90) and (95), $`u`$ may be written explicitly on the usual vector basis in spherical coordinates:
$`u`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}𝒩_l^0\left[u_0^l\left(r\right)P_l\left(\mathrm{cos}\theta \right)\widehat{e}_r+v_0^l\left(r\right)_\theta P_l\left(\mathrm{cos}\theta \right)\widehat{e}_\theta w_0^l\left(r\right)_\theta P_l\left(\mathrm{cos}\theta \right)\widehat{e}_\phi \right]`$ (7)
$`=`$ $`𝒩_0^0u_0^0\left(r\right)\widehat{e}_r+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}𝒩_l^0\left[u_0^l\left(r\right)P_l\left(\mathrm{cos}\theta \right)\widehat{e}_rv_0^l\left(r\right)P_l^1\left(\mathrm{cos}\theta \right)\widehat{e}_\theta +w_0^l\left(r\right)P_l^1\left(\mathrm{cos}\theta \right)\widehat{e}_\phi \right];`$
$`\widehat{e}_r`$, $`\widehat{e}_\theta `$ and $`\widehat{e}_\phi `$ are the unit-vectors repectively in the $`r`$, $`\theta `$ and $`\phi `$ directions.
This expansion of the vector fields allows us to separate explicitly the spatial coordinates $`(r,\theta )`$ in the vectorial partial differential equations which govern the problem: using this decomposition, we cast the problem into partial differential equations in $`t`$ and $`r`$ only. This point is very important numerically, because the existing stellar evolution codes in which we have to introduce the transport equations are all 1-dimensional, and also because we need to achieve much higher accuracy in the radial direction than in the horizontal.
## 3 Axisymmetric magnetic field
We start with expansions and equations involving the magnetic field.
### 3.1 Definitions
The magnetic field, being divergenceless, is conveniently split in its poloidal and toroidal parts:
$$B(r,\theta )=B_P(r,\theta )+B_T(r,\theta )\left(\xi _P(r,\theta )\widehat{e}_r\right)+\left(\xi _T(r,\theta )\widehat{e}_r\right),$$
(8)
where $`\xi _P`$ and $`\xi _T`$ are respectively the poloidal and the toroidal magnetic stream-functions. Using the classical method introduced by Bullard and Gellman (1954) for the spectral treatment of the geodynamo problem (see also James 1973, 1974 and Serebrianaya 1988), we get the following expansion in radial functions and in spherical harmonics which are defined in §A.1.1.:
$$\xi _P(r,\theta )=\underset{l=1}{\overset{\mathrm{}}{}}\xi _0^l\left(r\right)Y_l^0\left(\theta \right)\text{ }\text{and}\xi _T(r,\theta )=\underset{l=1}{\overset{\mathrm{}}{}}\chi _0^l\left(r\right)Y_l^0\left(\theta \right).$$
(9)
Since $`R_l^0\left(\theta \right)=Y_l^0\left(\theta \right)\widehat{e}_r`$, as we have seen in the previous section where we have defined $`R_l^0\left(\theta \right),S_l^0\left(\theta \right),T_l^0\left(\theta \right)`$, the field may be written as
$$B_P=\underset{l=1}{\overset{\mathrm{}}{}}\left(\xi _0^lR_l^0\left(\theta \right)\right)\text{and}B_T=\underset{l=1}{\overset{\mathrm{}}{}}\left(\chi _0^lR_l^0\left(\theta \right)\right).$$
(10)
Next, we develop the curl operator as explained in the appendix (eqs. (121) and (122)) to obtain the following expansions for the magnetic field:
$$B(r,\theta )=B_P(r,\theta )+B_T(r,\theta )=\underset{l=1}{\overset{\mathrm{}}{}}\left\{\left[l(l+1)\frac{\xi _0^l}{r^2}\right]R_l^0\left(\theta \right)+\left[\frac{1}{r}_r\xi _0^l\right]S_l^0\left(\theta \right)\right\}+\underset{l=1}{\overset{\mathrm{}}{}}\left\{\left[\frac{\chi _0^l}{r}\right]T_l^0\left(\theta \right)\right\}.$$
(11)
We proceed likewise for the current density $`j`$, which in the framework of magnetohydrodynamics is related with $`B`$ through the Maxwell-Ampère equation, neglecting the displacement current:
$$j(r,\theta )=\frac{1}{\mu _0}B(r,\theta ),$$
(12)
$`\mu _0`$ being the magnetic permeability of vacuum.
Therefore, using expression (8) for the field, we get:
$$j=\frac{1}{\mu _0}B=\frac{1}{\mu _0}\left[\left(\xi _P\widehat{e}_r\right)+\left(\xi _T\widehat{e}_r\right)\right],$$
(13)
which, using the relations (122) and (125), we can project again on $`R_l^0\left(\theta \right)`$, $`S_l^0\left(\theta \right)`$ and $`T_l^0\left(\theta \right)`$, and split into its poloidal and toroidal parts:
$$j_P(r,\theta )=\frac{1}{\mu _0}B_T(r,\theta )=\frac{1}{\mu _0}\left(\xi _T\widehat{e}_r\right)=\frac{1}{\mu _0}\underset{l=1}{\overset{\mathrm{}}{}}\left\{\left[l(l+1)\frac{\chi _0^l}{r^2}\right]R_l^0\left(\theta \right)+\left[\frac{1}{r}_r\chi _0^l\right]S_l^0\left(\theta \right)\right\}$$
(14)
$$j_T(r,\theta )=\frac{1}{\mu _0}B_P(r,\theta )=\frac{1}{\mu _0}\left(\xi _P\widehat{e}_r\right)=\frac{1}{\mu _0}\underset{l=1}{\overset{\mathrm{}}{}}\left\{\left[\mathrm{\Delta }_l\left(\frac{\xi _0^l}{r}\right)\right]T_l^0\left(\theta \right)\right\}.$$
(15)
with $`\mathrm{\Delta }_l`$ being the laplacian operator
$$\mathrm{\Delta }_l=_{r,r}+\frac{2}{r}_r\frac{l(l+1)}{r^2}.$$
(16)
We are now ready to examine the aspects of transport and mixing in the radiation zones of rotating stars that are related to the presence of an axisymmetric magnetic field. To achieve this, we shall proceed in three steps
* first, we introduce the Lorentz force;
* next, we derive the transport equation for $`B`$, which is the classical induction equation;
* finally, we take into account the energy losses due to Ohmic heating.
### 3.2 Lorentz force
The Lorentz force $`_{}`$ plays a crucial role in a rotating star, since it acts to render the angular velocity constant on the field lines of the poloidal field $`B_P`$. In this section, we establish the formal expression of $`_{}`$ in terms of vectorial spherical harmonics.
Starting from the definition of the Lorentz force:
$$_{}(r,\theta )=jB=\left[\frac{1}{\mu _0}(B)\right]B,$$
(17)
we replace $`B`$ by its expansions (11), and use the algebra ruling the vector product of two general axisymmetric vectors in Appendix A, whereby we obtain the following formal projection:
$`_{}(r,\theta )`$ $`=`$ $`𝒳_{_{};0}(r)\widehat{e}_r+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{𝒳_{_{};l}(r)P_l(\mathrm{cos}\theta )\widehat{e}_r+𝒴_{_{};l}(r)P_l^1(\mathrm{cos}\theta )\widehat{e}_\theta +𝒵_{_{};l}(r)P_l^1(\mathrm{cos}\theta )\widehat{e}_\phi \right\}`$
$`=`$ $`\left[{\displaystyle \frac{𝒳_{_{};0}(r)}{𝒩_0^0}}\right]R_0^0(\theta )+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{\left[{\displaystyle \frac{𝒳_{_{};l}(r)}{𝒩_l^0}}\right]R_l^0(\theta )+\left[{\displaystyle \frac{𝒴_{_{};l}(r)}{𝒩_l^0}}\right]S_l^0(\theta )+\left[{\displaystyle \frac{𝒵_{_{};l}(r)}{𝒩_l^0}}\right]T_l^0(\theta )\right\}.`$
In Appendix B, we give the explicit expressions for the radial functions $`𝒳_{_{};l}\left(r\right)`$, $`𝒴_{_{};l}\left(r\right)`$, $`𝒵_{_{};l}\left(r\right)`$ in terms of the magnetic stream-functions $`\xi _0^l`$ and $`\chi _0^l`$, the normalization coefficient $`𝒩_l^0`$ being given in (91).
As we have done previously for the magnetic field $`B`$ and for the current $`j`$, we split the magnetic force into its poloidal and toroidal parts:
$`_{,P}`$ $`=`$ $`j_TB_P+j_PB_T=𝒳_{_{};0}\widehat{e}_r+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{𝒳_{_{};l}P_l(\mathrm{cos}\theta )\widehat{e}_r+𝒴_{_{};l}P_l^1(\mathrm{cos}\theta )\widehat{e}_\theta \right\}`$
$`=`$ $`\left[{\displaystyle \frac{𝒳_{_{};0}}{𝒩_0^0}}\right]R_0^0(\theta )+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{\left[{\displaystyle \frac{𝒳_{_{};l}}{𝒩_l^0}}\right]R_l^0(\theta )+\left[{\displaystyle \frac{𝒴_{_{};l}}{𝒩_l^0}}\right]S_l^0(\theta )\right\}`$
$$_{,T}=j_PB_P=\underset{l=1}{\overset{\mathrm{}}{}}𝒵_{_{};l}P_l^1(\mathrm{cos}\theta )\widehat{e}_\phi =\underset{l=1}{\overset{\mathrm{}}{}}\left\{\left[\frac{𝒵_{_{};l}}{𝒩_l^0}\right]T_l^0(\theta )\right\}.$$
(20)
As it was underlined by Mestel et al. (1988) and as we shall see in §4, §5 and §6, this decomposition is very useful: the poloidal Lorentz force $`_{,P}`$ operates on the hydrostatic balance and thus contributes to the thermal imbalance, while the toroidal component $`_{,T}`$ acts on the transport of angular momentum through its torque.
The expression of $`_{}`$ in the case where only the $`l=\{1,2,3\}`$ terms are kept in the expansion of $`B`$ can be found in Appendix E. One should note that if the expansion of $`B`$ is terminated at mode $`l_{\mathrm{max}}=N_m`$ in (11), then the expansion of $`_{}`$ involves terms up to $`l_{\mathrm{max}}=2N_m`$, due to selection rules (see Appendix B).
### 3.3 Induction equation
We now turn to the evolution of the magnetic field in stellar radiation zones, which is governed by the induction equation:
$$_tB=\left(VB\right)\left(\eta B\right),$$
(21)
where we have allowed for an anisotropic eddy diffusivity $`\eta `$. We recall (cf. Paper I) that the macroscopic velocity field $`V`$ is the sum of a zonal flow $`𝒰_\phi =r\mathrm{sin}\theta \mathrm{\Omega }(r,\theta )\widehat{e}_\phi `$ and of a meridional flow $`𝒰(r,\theta )`$:
$$V(r,\theta )=r\mathrm{sin}\theta \mathrm{\Omega }(r,\theta )\widehat{e}_\phi +𝒰(r,\theta ).$$
(22)
The latter can be split into a spherically symmetric part, which represents the contractions and dilatations of the star during its evolution, plus the thermally-driven circulation
$$𝒰(r,\theta )=\dot{r}\widehat{e}_r+𝒰_M(r,\theta ),$$
(23)
the meridional flow being expanded in spherical functions:
$$𝒰_M=\underset{l>0}{}\left[U_l(r)P_l\left(\mathrm{cos}\theta \right)\widehat{e}_r+V_l(r)\frac{\mathrm{d}P_l(\mathrm{cos}\theta )}{\mathrm{d}\theta }\widehat{e}_\theta \right].$$
(24)
The radial functions $`U_l(r)`$ and $`V_l(r)`$ are related by the continuity equation, i.e. $`(\rho 𝒰_M)=0`$ in the anelastic approximation:
$$V_l=\frac{1}{l(l+1)\rho r}\frac{\mathrm{d}}{\mathrm{d}r}\left(\rho r^2U_l\right).$$
(25)
We introduce the expanded form of the velocity field in the induction equation (21):
$$_tB\left(\dot{r}\widehat{e}_rB\right)=\left[\left(𝒰_\phi +𝒰_M\right)B\right]\left(\eta B\right),$$
(26)
and as before we project this equation on the vectorial spherical harmonics $`R_l^0\left(\theta \right)`$, $`S_l^0\left(\theta \right)`$, $`T_l^0\left(\theta \right)`$. The time derivative on the left-hand side of (26) is readily derived (cf. 11):
$$_tB=\underset{l=1}{\overset{\mathrm{}}{}}\left\{\left[l(l+1)\frac{_t\xi _0^l}{r^2}\right]R_l^0(\theta )+\left[\frac{1}{r}_{t,r}\xi _0^l\right]S_l^0(\theta )+\left[\frac{_t\chi _0^l}{r}\right]T_l^0(\theta )\right\}.$$
(27)
Then, using the algebra related to the $`R_l^0\left(\theta \right)`$, $`S_l^0\left(\theta \right)`$, $`T_l^0\left(\theta \right)`$ and the identity (121), the following expansion of the contractions/dilatations term is obtained:
$$\left(\dot{r}\text{ }\widehat{e}_rB\right)=\underset{l=1}{\overset{\mathrm{}}{}}\left\{\left[l(l+1)\frac{\dot{r}}{r^2}_r\xi _0^l\right]R_l^0\left(\theta \right)+\left[\frac{1}{r}_r\left(\dot{r}_r\xi _0^l\right)\right]S_l^0\left(\theta \right)+\left[\frac{1}{r}_r\left(\dot{r}\chi _0^l\right)\right]T_l^0\left(\theta \right)\right\}.$$
(28)
In the same way, we get for the dissipation term:
$`\left(\eta B\right)`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\{\left[\eta _h{\displaystyle \frac{1}{r}}\mathrm{\Delta }_l\left(l(l+1){\displaystyle \frac{\xi _0^l}{r}}\right)\right]R_l^0\left(\theta \right)+\left[{\displaystyle \frac{1}{r}}_r\left(r\eta _h\mathrm{\Delta }_l\left({\displaystyle \frac{\xi _0^l}{r}}\right)\right)\right]S_l^0\left(\theta \right)`$ (29)
$`+`$ $`[{\displaystyle \frac{1}{r}}_r\left(\eta _h_r\chi _0^l\right)\eta _vl(l+1){\displaystyle \frac{\chi _0^l}{r^3}}]T_l^0\left(\theta \right)\},`$
where we recall that $`\mathrm{\Delta }_l`$ is the laplacian operator: $`\mathrm{\Delta }_l=_{r,r}+(2/r)_rl(l+1)/r^2`$.
The task is more difficult with the non-linear expressions describing the stretching and advection of the magnetic field: $`\left[𝒰_\phi B\right]+\left[𝒰_MB\right]`$. The first term corresponds to the generation of toroidal field trough the shearing of the poloidal component by the differential rotation, while the second one represents the advection of the field by the meridional circulation.
The first step is to expand ($`𝒰_\phi +𝒰_M)`$ on the $`R_l^0\left(\theta \right)`$, $`S_l^0\left(\theta \right)`$ and $`T_l^0\left(\theta \right)`$:
$$𝒰_\phi +𝒰_M=\underset{l=0}{\overset{\mathrm{}}{}}\left\{u_0^l\left(r\right)R_l^0\left(\theta \right)+v_0^l\left(r\right)S_l^0\left(\theta \right)+w_0^l\left(r\right)T_l^0\left(\theta \right)\right\}.$$
(30)
From (24), we get immediately for the meridional part:
$$u_0^l\left(r\right)=\frac{U_l\left(r\right)}{𝒩_l^0}\text{ and }v_0^l\left(r\right)=\frac{V_l\left(r\right)}{𝒩_l^0}.$$
(31)
Next, we expand the rotation law as in Paper I:
$$\mathrm{\Omega }(r,\theta )=\overline{\mathrm{\Omega }}\left(r\right)+\widehat{\mathrm{\Omega }}(r,\theta )=\overline{\mathrm{\Omega }}\left(r\right)+\underset{l>0}{}\mathrm{\Omega }_l\left(r\right)Q_l\left(\theta \right)=\underset{l=0}{\overset{\mathrm{}}{}}\mathrm{\Omega }_l^{}\left(r\right)P_l\left(\mathrm{cos}\theta \right),$$
(32)
where $`\overline{\mathrm{\Omega }}(r)=_0^\pi \mathrm{\Omega }(r,\theta )\mathrm{sin}^3\theta \mathrm{d}\theta /_0^\pi \mathrm{sin}^3\theta \mathrm{d}\theta `$ the horizontal functions are given by $`Q_l(\theta )=P_l\left(\mathrm{cos}\theta \right)I_l`$, with
$`I_l=_0^\pi P_l\left(\mathrm{cos}\theta \right)\mathrm{sin}^3\theta d\theta /_0^\pi \mathrm{sin}^3\theta d\theta =\delta _{l,0}\frac{1}{5}\delta _{l,2}`$. Thus, we get:
$$\{\begin{array}{cc}\mathrm{\Omega }_0^{}\left(r\right)=\overline{\mathrm{\Omega }}\left(r\right)_{l>0}\mathrm{\Omega }_l\left(r\right)I_l=\mathrm{\Omega }_0\left(r\right)+\frac{1}{5}\mathrm{\Omega }_2\left(r\right)\hfill & \\ \mathrm{\Omega }_l^{}\left(r\right)=\mathrm{\Omega }_l\left(r\right)\text{ for }l>0.\hfill & \end{array}$$
(33)
Using the identity (98), we obtain the zonal flow $`𝒰_\phi =_{l>0}w_0^l\left(r\right)T_l^0\left(\theta \right)`$, where $`w_0^l\left(r\right)`$ is given by
$$w_0^l(r)=r\left[\frac{D_{l1}^0}{𝒩_{l1}^0}\mathrm{\Omega }_{l1}^{}(r)\frac{C_{l+1}^0}{𝒩_{l+1}^0}\mathrm{\Omega }_{l+1}^{}(r)\right].$$
(34)
The expression of the $`w_0^l`$ when we keep only the two first terms of the expansion (32), $`\overline{\mathrm{\Omega }}`$ and $`\mathrm{\Omega }_2`$ is given in (161).
It remains to project the vector product of the two axisymmetric vectors, ($`𝒰_\phi +𝒰_M`$) and $`B`$, on the vector harmonics $`R_l^0\left(\theta \right)`$, $`S_l^0\left(\theta \right)`$ and $`T_l^0\left(\theta \right)`$. This task is accomplished in §A.2.3.; applying it to the advection term, we get:
$`\left(𝒰_\phi +𝒰_M\right)B`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{𝒳_{\mathrm{𝐀𝐝};l}(r)P_l(\mathrm{cos}\theta )\widehat{e}_r+𝒴_{\mathrm{𝐀𝐝};l}(r)P_l^1(\mathrm{cos}\theta )\widehat{e}_\theta +𝒵_{\mathrm{𝐀𝐝};l}(r)P_l^1(\mathrm{cos}\theta )\widehat{e}_\phi \right\}`$
$`=`$ $`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{\left[{\displaystyle \frac{𝒳_{\mathrm{𝐀𝐝};l}(r)}{𝒩_l^0}}\right]R_l^0(\theta )+\left[{\displaystyle \frac{𝒴_{\mathrm{𝐀𝐝};l}(r)}{𝒩_l^0}}\right]S_l^0(\theta )+\left[{\displaystyle \frac{𝒵_{\mathrm{𝐀𝐝};l}(r)}{𝒩_l^0}}\right]T_l^0(\theta )\right\},`$
where the radial functions $`𝒳_{\mathrm{𝐀𝐝};l}\left(r\right)`$, $`𝒴_{\mathrm{𝐀𝐝};l}\left(r\right)`$ and $`𝒵_{\mathrm{𝐀𝐝};l}\left(r\right)`$, which are explicit functions of $`\overline{\mathrm{\Omega }}`$, $`\mathrm{\Omega }_l`$, $`U_l`$, $`\xi _0^l`$ and $`\chi _0^l`$, are given in Appendix C. Finally, using once again (121), we obtain:
$`\left[\left(𝒰_\phi +𝒰_M\right)B\right]`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\{\left[{\displaystyle \frac{l(l+1)}{𝒩_l^0}}{\displaystyle \frac{𝒵_{\mathrm{𝐀𝐝};l}}{r}}\right]R_l^0(\theta )+\left[{\displaystyle \frac{1}{𝒩_l^0}}{\displaystyle \frac{1}{r}}_r\left(r𝒵_{\mathrm{𝐀𝐝};l}\right)\right]S_l^0(\theta )`$ (36)
$`+`$ $`\left[{\displaystyle \frac{1}{𝒩_l^0}}({\displaystyle \frac{𝒳_{\mathrm{𝐀𝐝};l}}{r}}+{\displaystyle \frac{1}{r}}_r\left(r𝒴_{\mathrm{𝐀𝐝};l}\right))\right]T_l^0\left(\theta \right)\}.`$
We are now ready to put in their final form the one-dimensional advection/diffusion equations respectively for the poloidal $`\xi _0^l`$ and for the toroidal magnetic stream-functions $`\chi _0^l`$, by collecting the $`R_l^0\left(\theta \right)`$ and the $`T_l^0\left(\theta \right)`$ components of (27), (28), (29) and (36):
$$\frac{\mathrm{d}\xi _0^l}{\mathrm{d}t}=\frac{1}{𝒩_l^0}r𝒵_{\mathrm{𝐀𝐝};l}+\eta _hr\mathrm{\Delta }_l\left(\frac{\xi _0^l}{r}\right)$$
(37)
$$\frac{\mathrm{d}\chi _0^l}{\mathrm{d}t}+_r\left(\dot{r}\right)\chi _0^l=\frac{1}{𝒩_l^0}\left[𝒳_{\mathrm{𝐀𝐝};l}+_r\left(r𝒴_{\mathrm{𝐀𝐝};l}\right)\right]+\left[_r\left(\eta _h_r\chi _0^l\right)\eta _vl(l+1)\frac{\chi _0^l}{r^2}\right].$$
(38)
Note that these equations involve the Lagrangian time-derivative, which makes them suitable for their implementation in stellar evolution codes. They are the equivalent in our formalism of the two classical equations for $`B_P=A`$, $`A=A\widehat{e}_\phi `$ being the potential vector, and $`B_T=B_T\widehat{e}_\phi `$:
$$_tA+\frac{1}{s}𝒰_M\left(sA\right)=\eta \left(^2A\frac{A}{s^2}\right)\text{and}_tB_T+s𝒰_M\left(\frac{B_T}{s}\right)+B_T𝒰_P=sB_P\mathrm{\Omega }+\eta \left(^2B_T\frac{B_T}{s^2}\right),$$
(39)
where $`s=r\mathrm{sin}\theta `$ and $`\eta `$ is taken isotropic and uniform (Campbell 1997, Mestel 1999).
Finally, we note that in the ideal case, in the absence of shear, meridional circulation and Ohmic diffusion, (37) and (38) reduce to:
$$\frac{\mathrm{d}\xi _0^l}{\mathrm{d}t}=0\text{ and }\frac{\mathrm{d}}{\mathrm{d}t}\left(\frac{\chi _0^l}{r^2\rho }\right)=0\text{and thus to}\frac{\mathrm{d}}{\mathrm{d}t}\left(r^2B_r\right)=0\text{ and }\frac{\mathrm{d}}{\mathrm{d}t}\left(\frac{B_\phi }{r\rho }\right)=0.$$
(40)
These equations express the Lagrangian flux conservation of respectively
* the magnetic field through the sphere of radius $`r`$, as the star expands or contracts,
* the toroidal field in a homothetic contraction or expansion where the density varies as $`r^3`$ (cf. Cowling 1957).
In Appendix E.1. we give the explicit expressions of (37) and (38) for the dipole, the quadrupole and the octupole ($`l=\{1,2,3\}`$), retaining the lowest order terms in the rotation law: $`\mathrm{\Omega }(r,\theta )=\overline{\mathrm{\Omega }}\left(r\right)+\mathrm{\Omega }_2\left(r\right)Q_2\left(\theta \right)`$, and the associated meridional circulation.
### 3.4 Ohmic heating
The last point we shall consider for completeness is Ohmic heating, which contributes somewhat to the thermal balance/imbalance. In the case of an anisotropic eddy-magnetic diffusivity, its input rate is given by
$$𝒥(r,\theta )=\frac{1}{\mu _0}\left[\eta \left(B\right)\right]\left(B\right),$$
(41)
which reduces to the classical result $`𝒥=j^2/\sigma `$, with $`\sigma =1/\left(\mu _0\eta \right)`$ being the conductivity, when the diffusivity is isotropic.
The scalar product of two axisymmetric vectors is treated in the §A.2.3.; applying the result to $`𝒥`$, we get:
$$𝒥(r,\theta )=\underset{l>0}{}𝒥_l\left(r\right)P_l\left(\mathrm{cos}\theta \right),$$
(42)
where the radial functions $`𝒥_l(r)`$, which are functions of the magnetic stream-functions $`\xi _0^l`$ and $`\chi _0^l`$, are given in Appendix D, and in Appendix E for the special case where only the $`l=\{1,2,3\}`$ modes are kept. Note again here that if the expansion of $`B`$ is cut at mode $`l_{\mathrm{max}}=N_m`$ in (11), then the expansion of $`𝒥`$ involves terms up to $`l_{\mathrm{max}}=2N_m`$.
We are now ready to introduce the Lorentz force in the equations governing the transport of momentum and of heat.
## 4 Transport of angular momentum
We start from the momentum equation
$$\rho \left[_tV+\left(V\right)V\right]=P\rho \varphi +\tau +_{},$$
(43)
where $`\rho `$ is the density, $`\varphi `$ the gravitational potential, $`\tau `$ the turbulent stresses and $`_{}`$ the Lorentz force. We insert in it the expression for the macroscopic velocity $`V`$ given above in (22, 23), and take its azimuthal component, which yields an advection/diffusion equation for the angular momentum density:
$$\rho \frac{\mathrm{d}}{\mathrm{d}t}\left(r^2\mathrm{sin}^2\theta \mathrm{\Omega }\right)+\left(\rho r^2\mathrm{sin}^2\theta \mathrm{\Omega }𝒰_M\right)=\frac{\mathrm{sin}^2\theta }{r^2}_r\left(\rho \nu _vr^4_r\mathrm{\Omega }\right)+\frac{1}{\mathrm{sin}\theta }_\theta \left(\rho \nu _h\mathrm{sin}^3\theta _\theta \mathrm{\Omega }\right)+\mathrm{\Gamma }_{_{}}(r,\theta ).$$
(44)
Note that here again we have introduced the Lagrangian time derivative, making use of the anelastic continuity equation. This equation is of course identical to that given in Paper I (eq. 14), except for the magnetic torque $`\mathrm{\Gamma }_{_{}}(r,\theta )=r\mathrm{sin}\theta e_\phi _{}(r,\theta )`$. As in Zahn (1992), we assume that the effect of the turbulent stresses on the large scale flow is adequately described by an anisotropic eddy viscosity, whose components are $`\nu _v`$ and $`\nu _h`$ respectively in the vertical and horizontal directions. In the absence of circulation, turbulence and magnetic field, we retrieve the Lagrangian conservation of angular momentum.
It remains to project this equation on spherical harmonics; the expansions of the meridional circulation and the angular velocity were established in Paper I and they have been recalled above in (24, 32). We proceed likewise for the magnetic torque:
$$\mathrm{\Gamma }_{_{}}(r,\theta )=r\mathrm{sin}\theta _{,\phi }=r\mathrm{sin}\theta \underset{l=0}{\overset{\mathrm{}}{}}𝒵_{_{};l}(r)P_l^1\left(\mathrm{cos}\theta \right),$$
(45)
where $`_{,\phi }`$ is the azimuthal component of $`_{}`$ obtained from (LABEL:Lorentzexp). Then, using the following property of the associated Legendre functions:
$$P_l^1\left(\mathrm{cos}\theta \right)=\frac{\mathrm{d}}{\mathrm{d}\theta }P_l\left(\mathrm{cos}\theta \right)=\mathrm{sin}\theta \frac{\mathrm{d}}{\mathrm{d}\mu }P_l\left(\mu \right)\text{ where }\mu =\mathrm{cos}\theta $$
(46)
and the expression for $`dP_l\left(\mu \right)/d\mu `$ which is given in (134), we put $`\mathrm{\Gamma }_{_{}}`$ in its final form:
$$\mathrm{\Gamma }_{_{}}(r,\theta )=\underset{l=0}{\overset{\mathrm{}}{}}\mathrm{\Gamma }_l\left(r\right)\mathrm{sin}^2\theta P_l\left(\mathrm{cos}\theta \right)\text{where}\mathrm{\Gamma }_l\left(r\right)=\underset{k=0}{\overset{\mathrm{}}{}}r𝒵_{_{};k}\left(r\right)\left[\underset{j=0}{\overset{E\left[\frac{k1}{2}\right]}{}}\left(2k4j1\right)\delta _{l,k2j1}\right],$$
(47)
$`E\left[x\right]`$ is the integer part of $`x`$, and $`𝒵_{_{};k}`$ is defined in (LABEL:Lorentzexp) and spelled out in explicit form in Appendix B.
### 4.1 Evolution equation for the angular mean angular velocity
Taking the horizontal average of equation (44) over an isobar and using the assumption that $`\overline{\mathrm{\Omega }}(r)\mathrm{\Omega }_l(r)`$, we obtain the following vertical advection/diffusion equation for the mean angular velocity $`\overline{\mathrm{\Omega }}`$:
$$\rho \frac{\mathrm{d}}{\mathrm{d}t}(r^2\overline{\mathrm{\Omega }})=\frac{1}{5r^2}_r\left(\rho r^4\overline{\mathrm{\Omega }}U_2\right)+\frac{1}{r^2}_r\left(\rho \nu _vr^4_r\overline{\mathrm{\Omega }}\right)+\overline{\mathrm{\Gamma }}_{_{}}\left(r\right),$$
(48)
$$\text{where the mean Lorentz torque is}\overline{\mathrm{\Gamma }}_{_{}}\left(r\right)=\frac{_0^\pi \mathrm{\Gamma }_{_{}}(r,\theta )\mathrm{sin}\theta \mathrm{d}\theta }{_0^\pi \mathrm{sin}^3\theta \mathrm{d}\theta }=\left[\mathrm{\Gamma }_0\left(r\right)\frac{1}{5}\mathrm{\Gamma }_2\left(r\right)\right].$$
(49)
Note that only the $`l=2`$ component of the circulation is able to advect a net amount of angular momentum; the higher order components of $`𝒰_M`$ (for instance those induced in its tachocline by a differentially rotating convection zone) do not contribute to the vertical transport of angular momentum, as was pointed out in Spiegel and Zahn (1992). The explicit form of (48), keeping only the $`l=\{1,2,3\}`$ terms in the expansion of $`B`$ (cf. eq. (11)) is derived in Appendix E.2.1. .
### 4.2 Evolution equation for the differential rotation in latitude
We establish the equation governing the horizontal transport of angular momentum by multiplying eq. (48) through $`\mathrm{sin}^2\theta `$ and subtracting it from the original form (44):
$`\rho {\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}\left(r^2\mathrm{sin}^2\theta \widehat{\mathrm{\Omega }}\right)+\left(\rho r^2\mathrm{sin}^2\theta \overline{\mathrm{\Omega }}𝒰_M\right)+{\displaystyle \frac{\mathrm{sin}^2\theta }{5r^2}}_r\left(\rho r^4\overline{\mathrm{\Omega }}U_2\right)={\displaystyle \frac{\mathrm{sin}^2\theta }{r^2}}_r\left(\rho \nu _vr^4_r\widehat{\mathrm{\Omega }}\right)+{\displaystyle \frac{1}{\mathrm{sin}\theta }}_\theta \left(\rho \nu _h\mathrm{sin}^3\theta _\theta \widehat{\mathrm{\Omega }}\right)`$ (50)
$`+`$ $`\mathrm{\Gamma }_{_{}}\mathrm{sin}^2\theta \overline{\mathrm{\Gamma }}_{_{}}.`$
In the advection term we have again neglected the fluctuation $`\widehat{\mathrm{\Omega }}(r,\theta )`$ compared to the mean $`\overline{\mathrm{\Omega }}`$.
The next step is to replace $`\widehat{\mathrm{\Omega }}(r,\theta )`$ by its expansion (32) in the horizontal functions $`Q_l(\theta )`$. For $`l=2`$, the equation separates neatly into
$$\rho \frac{\mathrm{d}}{\mathrm{d}t}\left(r^2\mathrm{\Omega }_2\right)2\rho \overline{\mathrm{\Omega }}r\left[2V_2\alpha (r)U_2\right]=\frac{1}{r^2}_r\left(\rho \nu _vr^4_r\mathrm{\Omega }_2\right)10\rho \nu _h\mathrm{\Omega }_2+\mathrm{\Gamma }_2$$
(51)
$$\text{with}V_2=\frac{1}{6\rho r}\frac{\mathrm{d}}{\mathrm{d}r}\left(\rho r^2U_2\right)\text{ and }\alpha (r)=\frac{1}{2}\frac{\mathrm{d}\mathrm{ln}\left(r^2\overline{\mathrm{\Omega }}\right)}{\mathrm{d}\mathrm{ln}r}.$$
(52)
It can be simplified by assuming that the turbulent transport is more efficient in the horizontal than in the vertical direction (i.e. $`\nu _v\nu _h`$):
$$\rho \frac{\mathrm{d}}{\mathrm{d}t}\left(r^2\mathrm{\Omega }_2\right)2\rho \overline{\mathrm{\Omega }}r\left[2V_2\alpha U_2\right]=10\rho \nu _h\mathrm{\Omega }_2+\mathrm{\Gamma }_2.$$
(53)
In the asymptotic regime $`tr^2/\nu _h`$, a stationary state is reached:
$$\nu _h\mathrm{\Omega }_2=\frac{1}{5}r\left[2V_2\alpha U_2\right]\overline{\mathrm{\Omega }}+\frac{\mathrm{\Gamma }_2}{10\rho },$$
(54)
where horizontal diffusion balances horizontal advection and the torque due to the Lorentz force.
For $`l>2`$ the situation is more intricate, because there are couplings between terms of different $`l`$, which prevent a clean separation for each $`l`$. This is mainly due to the magnetic torque. Indeed, we recall that in the hydrodynamical case the hypothesis of Spiegel & Zahn (1992) allows such separation (cf. Paper I, §3.2.). Therefore, we choose here to stop the expansion of the rotation law at $`\mathrm{\Omega }_2`$. The explicit form of (53), keeping only the $`l=\{1,2,3\}`$ terms in the expansion of $`B`$ (cf. eq. (11)) is derived in Appendix E.2.1. .
## 5 Structural properties of the differentially rotating magnetic star
### 5.1 Baroclinic relation
As in Paper I, we consider a non-uniform and a non-cylindrical rotation law, and add here a general axisymmetric magnetic field. Then neither the centrifugal force nor the Lorentz force derive from a potential, and therefore the isobars and the surfaces of constant density in general do not coincide. In order to determine how the density varies on an isobar, we start from the hydrostatic equation:
$$\frac{1}{\rho }P=g=\varphi +_𝒞+\frac{_{,P}}{\rho },$$
(55)
where $`\varphi `$ is the gravitational potential and $`g`$ the local effective gravity, which includes both the centrifugal force
$`_𝒞=\frac{1}{2}\mathrm{\Omega }^2\left(r^2\mathrm{sin}^2\theta \right)`$ and the meridional magnetic force $`_{,P}`$. Taking the curl of this equation, we get
$$\frac{1}{\rho ^2}\rho P=\frac{1}{\rho }\rho g=\frac{1}{2}(\mathrm{\Omega })^2(r\mathrm{sin}\theta )^2+\left(\frac{_{,P}}{\rho }\right),$$
(56)
which to first order reduces to
$$\frac{\rho ^{^{}}g}{\overline{\rho }}=\left[_r(\mathrm{\Omega }^2)r\mathrm{cos}\theta \mathrm{sin}\theta _\theta (\mathrm{\Omega }^2)\mathrm{sin}^2\theta \right]\widehat{e}_\phi +\left(\frac{_{,P}}{\overline{\rho }}\right),$$
(57)
with $`\rho ^{}(r,\theta )`$ being the variation of the density on the isobar.
We now expand all terms in spherical harmonics. For the density fluctuation this is readily done:
$$\rho ^{^{}}(r,\theta )=\underset{l>0}{}\stackrel{~}{\rho _l}\left(r\right)P_l(\mathrm{cos}\theta ).$$
(58)
Then, differentiating the equation of state:
$$\frac{d\rho }{\rho }=\alpha \frac{dP}{P}\delta \frac{dT}{T}+\phi \frac{d\mu }{\mu }$$
(59)
we can write the modal amplitude of the density fluctuation as
$$\frac{\stackrel{~}{\rho _l}}{\overline{\rho }}=\mathrm{\Theta }_l=\phi \mathrm{\Lambda }_l\delta \mathrm{\Psi }_l.$$
(60)
For the centrifugal force, we recall that
$$\mathrm{\Omega }(r,\theta )=\overline{\mathrm{\Omega }}(r)+\underset{l>0}{}\mathrm{\Omega }_l(r)\left(P_l(\mathrm{cos}\theta )I_l\right),$$
(61)
which leads us to
$$\mathrm{\Omega }^2(r,\theta )=\left[\overline{\mathrm{\Omega }}^22\overline{\mathrm{\Omega }}\underset{l>0}{}\mathrm{\Omega }_lI_l\right]+2\overline{\mathrm{\Omega }}\underset{l>0}{}\mathrm{\Omega }_lP_l(\mathrm{cos}\theta )$$
(62)
keeping only the terms linear in $`\mathrm{\Omega }_l`$ (beside $`\overline{\mathrm{\Omega }}^2`$). Finally, we draw the expansion of the poloidal Lorentz force from (LABEL:Lorentzexp):
$$_{,P}(r,\theta )=\underset{l=1}{\overset{\mathrm{}}{}}\left\{𝒳_{_{};l}(r)P_l(\mathrm{cos}\theta )\widehat{e}_r+𝒴_{_{};l}(r)P_l^1(\mathrm{cos}\theta )\widehat{e}_\theta \right\}.$$
(63)
It remains to insert these expansions in eq. 57, noting that it projects itself only on the azimuthal vectorial spherical harmonics $`T_l^0\left(\theta \right)`$. Using the algebra related with these harmonics, we reach the following expression for the modal amplitudes of the relative fluctuation of density on an isobar:
$$\frac{\stackrel{~}{\rho _l}\left(r\right)}{\overline{\rho }}=\phi \mathrm{\Lambda }_l\delta \mathrm{\Psi }_l=\frac{r}{\overline{g}}\left[𝒟_l(r)+\frac{𝒳_{_{};l}\left(r\right)}{r\overline{\rho }\left(r\right)}+\frac{1}{r}\frac{\mathrm{d}}{\mathrm{d}r}\left(r\frac{𝒴_{_{};l}\left(r\right)}{\overline{\rho }\left(r\right)}\right)\right].$$
(64)
$`\overline{g}`$ is the horizontal average of the modulus of $`g`$ and $`𝒟_l\left(r\right)`$ has been derived in Paper I (eq. 47):
$`𝒟_l(r)`$ $`=`$ $`𝒩_l^0\{r_r[\overline{\mathrm{\Omega }}^2(r)2\overline{\mathrm{\Omega }}(r)\mathrm{\Omega }_2(r)I_2]{\displaystyle \frac{1}{3𝒩_2^0}}\delta _{l,2}`$ (65)
$`+`$ $`2r{\displaystyle \underset{s>0}{}}_r\left(\overline{\mathrm{\Omega }}(r)\mathrm{\Omega }_s(r)\right){\displaystyle \frac{1}{𝒩_s^0}}\left[A_s^0\left(C_{s1}^0\delta _{l,s2}+D_{s1}^0\delta _{l,s}\right)+B_s^0\left(C_{s+1}^0\delta _{l,s}+D_{s+1}^0\delta _{l,s+2}\right)\right]`$
$``$ $`2\overline{\mathrm{\Omega }}(r){\displaystyle \underset{s>0}{}}{\displaystyle \frac{\mathrm{\Omega }_s(r)}{𝒩_s^0}}[G_s^0(C_{s+1}^0\delta _{l,s}+D_{s+1}^0\delta _{l,s+2})H_s^0(C_{s1}^0\delta _{l,s2}+D_{s1}^0\delta _{l,s})]\},`$
where all the numerical coefficients ($`A_l^0,B_l^0,C_l^0,D_l^0,G_l^0,H_l^0`$) are given in Appendix A.
This baroclinic equation (64 \- 65) plays a key role in linking the density fluctuation on an isobar with the differential rotation and the magnetic field. It allows us to close the system formed by the induction equation, that for the transport of angular momentum, that for the transport of the chemical species and that for the transport of heat, which we shall establish in §6.
### 5.2 Effective gravity and perturbed potential
It remains to examine how the redistribution of mass we just described modifies the local gravity. In the absence of magnetic field, the perturbing force is just the centrifugal force $`_𝒞=\frac{1}{2}\mathrm{\Omega }^2(r^2\mathrm{sin}^2\theta )`$, whose $`r`$ and $`\theta `$ components are expanded as
$$_{𝒞,r}(r,\theta )=\underset{l}{}a_l(r)P_l(\mathrm{cos}\theta )_{𝒞,\theta }(r,\theta )=\underset{l}{}b_l(r)_\theta P_l(\mathrm{cos}\theta ).$$
(66)
where $`a_l(r)`$ and $`b_l(r)`$ are given by:
$`a_l(r)`$ $`=`$ $`{\displaystyle \frac{2}{3}}r\left[\overline{\mathrm{\Omega }}^2(r)2\overline{\mathrm{\Omega }}(r)\mathrm{\Omega }_2(r)I_2\right]\left(\delta _{l,0}\delta _{l,2}\right)`$ (67)
$`+`$ $`2r\overline{\mathrm{\Omega }}(r){\displaystyle \underset{s>0}{}}\mathrm{\Omega }_s(r)\left[{\displaystyle \frac{1}{𝒩_s^0}}\left\{C_s^0\left(G_{s1}^0𝒩_s^0\delta _{l,s}H_{s1}^0𝒩_{s2}^0\delta _{l,s2}\right)D_s^0\left(G_{s+1}^0𝒩_{s+2}^0\delta _{l,s+2}H_{s+1}^0𝒩_s^0\delta _{l,s}\right)\right\}\right],`$
$`b_l(r)`$ $`=`$ $`{\displaystyle \frac{1}{3}}r\left[\overline{\mathrm{\Omega }}^2(r)2\overline{\mathrm{\Omega }}(r)\mathrm{\Omega }_2(r)I_2\right]\delta _{l,2}`$ (68)
$`+`$ $`2r\overline{\mathrm{\Omega }}(r){\displaystyle \underset{s>0}{}}{\displaystyle \frac{\mathrm{\Omega }_s(r)}{𝒩_s^0}}\left\{A_s^0\left(C_{s1}^0𝒩_{s2}^0\delta _{l,s2}+D_{s1}^0𝒩_s^0\delta _{l,s}\right)+B_s^0\left(C_{s+1}^0𝒩_s^0\delta _{l,s}+D_{s+1}^0𝒩_{s+2}^0\delta _{l,s+2}\right)\right\}.`$
To this we have here to add the meridional components of the Lorentz force
$$_{,P;r}(r,\theta )=\underset{l}{}𝒳_{_{};l}\left(r\right)P_l\left(\mathrm{cos}\theta \right)_{,P;\theta }=\underset{l}{}𝒴_{_{};l}\left(r\right)_\theta P_l\left(\mathrm{cos}\theta \right).$$
(69)
In the expressions that we have derived in Paper I, it suffices then replace $`\rho _0a_l`$ by $`\rho _0a_l+𝒳_{_{};l}`$ and $`\rho _0b_l`$ by $`\rho _0b_l+𝒴_{_{};l}`$ to establish the version including the effect of the magnetic field. Thus for the gravity perturbation along an isobar we now get (cf. Paper I, eq. 72)
$$\frac{\stackrel{~}{g}_l}{\overline{g}}=\left[\frac{\mathrm{d}g_0}{\mathrm{d}r}\frac{1}{g_0^2}r\left(b_l+\frac{𝒴_{_{};l}}{\rho _0}\right)+\frac{1}{g_0}\left(a_l+\frac{𝒳_{_{};l}}{\rho _0}\right)\right]+\frac{\mathrm{d}}{\mathrm{d}r}\left(\frac{\widehat{\varphi }_l}{g_0}\right).$$
(70)
The last term involves the fluctuation of the gravity field along the sphere, $`\widehat{\varphi }_l`$, which is obtained by integrating the Poisson equation (cf. Paper I, eq. 60), modified along the same rules:
$$\frac{1}{r}\frac{\mathrm{d}^2}{\mathrm{d}r^2}\left(r\widehat{\varphi }_l\right)\frac{l(l+1)}{r^2}\widehat{\varphi }_l\frac{4\pi G}{g_0}\frac{\mathrm{d}\rho _0}{\mathrm{d}r}\widehat{\varphi }_l=\frac{4\pi G}{g_0}\left[\rho _0a_l+\frac{\mathrm{d}}{\mathrm{d}r}\left(r\rho _0b_l\right)+𝒳_{_{};l}+\frac{\mathrm{d}}{\mathrm{d}r}\left(r𝒴_{_{};l}\right)\right].$$
(71)
Let us recall that Sweet (1950) was the first to establish this result for the most general perturbing force.
## 6 Thermal imbalance and transport of heat
It remains to implement the magnetic field in the heat equation:
$$\rho T\left[\frac{S}{t}+𝒰S\right]=\left(\chi T\right)+\rho ϵF_h+𝒥,$$
(72)
where $`S`$ is the entropy per unit mass, $`\chi `$ the thermal conductivity, $`ϵ`$ the nuclear energy production rate per unit mass and $`F_h`$ the flux carried in the horizontal direction by the anisotropic turbulence. Note that in a medium of varying composition, we have to take into account the entropy of mixing (cf. Maeder & Zahn 1998). In the simplest case, applicable to main-sequence stars, where the stellar material can be approximated by a mixture of hydrogen and helium with a fixed abundance of metals, it can be expressed in terms of the mean molecular weight only; then we have
$$\mathrm{d}S=C_p\left[\frac{\mathrm{d}T}{T}_{\mathrm{ad}}\frac{\mathrm{d}P}{P}+\mathrm{\Phi }(P,T,\mu )\frac{\mathrm{d}\mu }{\mu }\right]$$
(73)
where $`_{\mathrm{ad}}`$ is the adiabatic gradient and $`\mathrm{\Phi }`$ is a function of the metal mass fraction and of $`\mu `$, the mean molecular weight.
The magnetic field manifests itself in the heat equation (72) through the Ohmic heating term $`𝒥`$, but also in the divergence of the thermal flux, because it involves the divergence of the perturbing force. This force includes both the centrifugal force $`_𝒞`$ and the Lorentz force per unit volume, $`_{,P}/\rho _0`$, and their divergence will again be expanded in spherical harmonics as
$$_𝒞=\overline{f}_𝒞+\underset{l>0}{}\stackrel{~}{f}_{𝒞,l}P_l\left(\mathrm{cos}\theta \right)\text{and}\left(\frac{_{,𝒫}}{\rho _0}\right)=\overline{f}_{}+\underset{l>0}{}\stackrel{~}{f}_{,l}P_l\left(\mathrm{cos}\theta \right).$$
(74)
In Paper I (eq. 85) we gave the expressions of $`\overline{f}_𝒞`$ and $`\stackrel{~}{f}_{𝒞,l}`$ in terms of $`a_l`$ and $`b_l`$:
$$\overline{f}_𝒞=\frac{1}{r^2}_r\left(r^2a_0\right)\text{ and }\stackrel{~}{f}_{𝒞,l}=\frac{1}{r^2}_r\left(r^2a_l\right)+l\left(l+1\right)\frac{b_l}{r}.$$
(75)
As above in §5.2., it suffices to replace $`a_l`$ by $`𝒳_{_{};l}/\rho _0`$ and $`b_l`$ by $`𝒴_{_{};l}/\rho _0`$ to obtain the equivalent expressions for the divergence of the Lorentz force:
$$\overline{f}_{}=\frac{1}{r^2}_r\left(r^2\frac{𝒳_{_{};0}}{\rho _0}\right)\text{ and }\stackrel{~}{f}_{,l}=\frac{1}{r^2}_r\left(r^2\frac{𝒳_{_{};l}}{\rho _0}\right)+l\left(l+1\right)\frac{𝒴_{_{;l}}}{r\rho _0}.$$
(76)
This leads us to the modal form of the heat equation, which can be implemented directly into a stellar structure code:
$$\overline{T}C_p\left[\frac{\mathrm{d}\mathrm{\Psi }_l}{\mathrm{d}t}+\mathrm{\Phi }\frac{\mathrm{d}\mathrm{ln}\overline{\mu }}{\mathrm{d}t}\mathrm{\Lambda }_l+\frac{U_l(r)}{H_p}\left(_{ad}\right)\right]=\frac{L(r)}{M(r)}𝒯_l(r)+\frac{𝒥_l\left(r\right)}{\overline{\rho }},$$
(77)
where $`𝒯_l`$ is given by:
$`𝒯_l`$ $`=`$ $`2\left[1{\displaystyle \frac{\overline{f}_𝒞+\overline{f}_{}}{4\pi G\overline{\rho }}}{\displaystyle \frac{\left(\overline{ϵ}+\overline{ϵ}_{\mathrm{grav}}\right)}{ϵ_m}}\right]{\displaystyle \frac{\stackrel{~}{g}_l}{\overline{g}}}+{\displaystyle \frac{\stackrel{~}{f}_{𝒞,l}+\stackrel{~}{f}_{,l}}{4\pi G\overline{\rho }}}{\displaystyle \frac{\overline{f}_𝒞+\overline{f}_{}}{4\pi G\overline{\rho }}}\left(\delta \mathrm{\Psi }_l+\phi \mathrm{\Lambda }_l\right)`$
$`+`$ $`{\displaystyle \frac{\rho _m}{\overline{\rho }}}\left[{\displaystyle \frac{r}{3}}_r\left(H_T_r\mathrm{\Psi }_l(1\delta +\chi _T)\mathrm{\Psi }_l(\phi +\chi _\mu )\mathrm{\Lambda }_l\right){\displaystyle \frac{l(l+1)H_T}{3r}}\left(1+{\displaystyle \frac{D_h}{K}}\right)\mathrm{\Psi }_l\right]`$
$`+`$ $`{\displaystyle \frac{\left(\overline{ϵ}+\overline{ϵ}_{grav}\right)}{ϵ_m}}\left\{\left(H_T_r\mathrm{\Psi }_l(1\delta +\chi _T)\mathrm{\Psi }_l(\phi +\chi _\mu )\mathrm{\Lambda }_l\right)+(f_ϵϵ_Tf_ϵ\delta +\delta )\mathrm{\Psi }_l+(f_ϵϵ_\mu +f_ϵ\phi \phi )\mathrm{\Lambda }_l\right\}.`$
We recall that $`L`$ is the luminosity, $`M`$ the mass, $`\overline{T}`$ the horizontal average of the temperature, $`C_p`$ the specific heat at constant pressure, and $``$ the radiative gradient. We have also introduced the temperature scale-height $`H_T=\left|\mathrm{d}r/\mathrm{d}\mathrm{ln}\overline{T}\right|`$, the thermal diffusivity $`K=\overline{\chi }/\overline{\rho }C_p`$, the horizontal eddy-diffusivity $`D_h`$ and $`f_ϵ=\overline{ϵ}/\left(\overline{ϵ}+\overline{ϵ}_{\mathrm{grav}}\right)`$, with $`\overline{ϵ}`$ and $`\overline{ϵ}_{\mathrm{grav}}`$ being respectively the mean nuclear and gravitational energy release rates, whereas $`ϵ_\mu `$ and $`\chi _\mu `$ are the logaritmic derivatives of $`ϵ`$ and of the radiative conductivity $`\chi `$ with respect to $`\mu `$, their derivatives with respect to $`T`$ being noted as $`ϵ_T`$ and $`\chi _T`$. Moreover, we have $`ϵ_m=L\left(r\right)/M\left(r\right)`$ and $`\rho _m`$ is the mean density inside the considered level surface.
Three remarks before we conclude this section.
* First, we have noted previously that if the expansion of $`B`$ is stopped at mode $`l_{\mathrm{max}}=N_m`$ in (11), then the expansions of $`_{}`$ and of $`𝒥`$ must include all terms up to $`l_{\mathrm{max}}=2N_m`$, due to selection rules. Therefore, keeping $`\mathrm{\Omega }(r,\theta )=\overline{\mathrm{\Omega }}\left(r\right)+\mathrm{\Omega }_2\left(r\right)Q_2\left(\theta \right)`$, the expansion of the meridional flow will extend to $`l_{𝒰_M;\mathrm{max}}=\mathrm{max}\{4,2N_m\}`$.
* Next, the relative importance of the terms contributed respectively by the magnetic force and by the centrifugal force depends on the ratio $`\left(V_A/V_\mathrm{\Omega }\right)^2`$ where $`V_A=B/\sqrt{\mu _0\rho _0}`$ is the Alfvén speed and $`V_\mathrm{\Omega }=R\mathrm{\Omega }`$, $`R`$ being the radius of the star. If the magnetic field is weak enough, this ratio will be small except just below the surface because of the decay of the density. In this case, for the implementation in existing stellar evolution codes of the equations related to the thermal imbalance and to the structural properties of the star (cf. §5), we keep only the modes $`l=\{2,4\}`$, which allows to describe the tachocline(s) circulation to first order (see Appendix E). However, in the case of a strong field higher order modes of the Lorentz force must be taken into account.
* Finally, the ratio between the term representing the Ohmic heating, $`𝒥_l\left(r\right)/\overline{\rho }`$, and the highest-order derivative term of $`\left(\chi T\right)`$ in (77-LABEL:tcal-final), namely $`(L/M)(\rho _m/\overline{\rho })(r/3)_r\left(H_T_r\mathrm{\Psi }_l\right)`$, is given by $`(\eta /K)\left(V_A/V_\mathrm{\Omega }\right)^2`$. Therefore the Ohmic term can be neglected since $`V_A<R\mathrm{\Omega }`$ and $`\eta K`$.
In conclusion, the link between the circulation and its cause, namely the thermal imbalance due to the rotation, the magnetic field and the chemical composition, is established through (77-LABEL:tcal-final) where the temperature fluctuation on an isobar is governed by an advection/diffusion equation from which we can derive the radial component of $`𝒰_M`$. These equations have been established for a rotation law which depends on $`r`$ and $`\theta `$, and that allows us to treat simultaneously the bulk of the radiation zones and the tachoclines in presence of a general axisymmetric magnetic field.
## 7 Boundary conditions
The system of equations is now complete: we have the induction equation, an advection/diffusion equation for the transport of angular momentum (mean and fluctuating) including the action of the magnetic torque, another for the temperature (mean and fluctuating), and the baroclinic relation which allows us to close the system. The equation for the transport of chemical species (or alternatively for the transport of molecular weight) is unchanged. It thus remains to specify the boundaries conditions of this system, to be applied on the limits of the radiation zone. To be specific, we consider a star with a radiation zone located between a convective core and an upper convection zone. We designate by $`r_b`$ and $`r_t`$ the respective radius of the base and of the top of that radiative zone. Of course, in a solar-type main-sequence star we have $`r_b=0`$, whereas for a massive main sequence star $`r_t=R`$.
The boundary conditions for the equations already present in the hydrodynamic case, which was examined in Paper I, are unchanged, except those to be applied to the equation of angular momentum transport, since it includes the Lorentz torque. The novelty here is the induction equation, which we have split in two equations, respectively for the poloidal and toroidal magnetic stream-functions $`\xi _0^l`$ (eq. 37) and $`\chi _0^l`$ (eq. 38). These equations are of second order in $`r`$, and therefore each of them requires two boundary conditions. In principle they can be obtained by solving the induction equation, together with the equation of momentum, etc. in each adjacent convective zone, but this represents a formidable task which is well beyond our scope here. Instead, we make the following - simplifying but reasonable - hypothesis, valid only in the convection zone(s):
* (i): we do not take into account the meridional velocity field. The only motion we retain is the zonal flow associated with the differential rotation: $`𝒰_\phi =r\mathrm{sin}\theta \mathrm{\Omega }_{CZ}(r,\theta )\widehat{e}_\phi `$, with $`\mathrm{\Omega }_{CZ}(r,\theta )`$ being the angular velocity in the convection zone;
* (ii): we put $`\eta _v=\eta _h=\eta _{\mathrm{turb}}`$ where $`\eta _{\mathrm{turb}}`$ is the magnetic eddy-diffusivity, which we assume constant inside the considered convective region;
* (iii): this eddy-diffusivity is high enough to allow for a stationary state (compared to the slow evolution in the radiation zone).
With these assumptions, the poloidal field is potential, and (37) reduces to:
$$\frac{\mathrm{d}^2\xi _0^l}{\mathrm{d}r^2}\frac{l\left(l+1\right)}{r^2}\xi _0^l=0;$$
(79)
it admits the general solution:
$$\xi _0^l=Ar^{l+1}+\frac{B}{r^l}.$$
(80)
When applying it to a convective core where we must take the non-singular solution at the origin:
$$r\frac{\mathrm{d}\xi _0^l}{\mathrm{d}r}\left(l+1\right)\xi _0^l=0\text{at}r=r_b,$$
(81)
which becomes $`\xi _0^l=0`$ if $`r_b=0`$. In the case of a convective envelope, surrounded by vacuum, we have
$$r\frac{\mathrm{d}\xi _0^l}{\mathrm{d}r}+l\xi _0^l=0\text{at}r=r_t.$$
(82)
With the same assumptions, the toroidal field obeys the simplified form of (38):
$$\frac{\mathrm{d}^2\chi _0^l}{\mathrm{d}r^2}\frac{l\left(l+1\right)}{r^2}\chi _0^l=\frac{𝒮_l\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)}{\eta _{\mathrm{turb}}}\text{ }\text{where}\text{ }𝒮_l\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)=\frac{1}{𝒩_0^l}\left[𝒳_{\mathrm{𝐀𝐝};l}\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)+_r\left(r𝒴_{\mathrm{𝐀𝐝};l}\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)\right)\right].$$
(83)
$`𝒳_{\mathrm{𝐀𝐝};l}\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)`$ and $`𝒴_{\mathrm{𝐀𝐝};l}\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)`$ are derived using the previous multipolar solutions for $`\xi _0^l`$ (cf. eq. 80), expanding $`\mathrm{\Omega }_{\mathrm{CZ}}`$ in spherical functions like $`\mathrm{\Omega }_{\mathrm{CZ}}(r,\theta )=_l\mathrm{\Omega }_l\left(r\right)P_l\left(\mathrm{cos}\theta \right)`$ where the $`\mathrm{\Omega }_l`$ are taken from the results of the inversion of helioseismological data (Corbard et al. 2002), using (34) and Appendix C. Then, the solution of (83) inside a convection zone is derived using Green’s functions:
$$\chi _0^l=\frac{1}{2l+1}\frac{1}{\eta _{\mathrm{turb}}}\left[r^{l+1}_r^{R_{\mathrm{extCZ}}}x^l𝒮_l\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)dx+r^l_{R_{intCZ}}^rx^{l+1}𝒮_l\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)dx\right]$$
(84)
where $`R_{intCZ}`$ and $`R_{\mathrm{extCZ}}`$ are respectively the radius at its base and at its top. Applying (84) to a convective core where $`R_{intCZ}=0`$ and $`R_{\mathrm{extCZ}}=r_b`$, we get:
$$\chi _0^l=\frac{1}{2l+1}\frac{1}{\eta _{\mathrm{turb}}}r_b^l_0^{r_b}x^{l+1}𝒮_l\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)dx\text{at}r=r_b.$$
(85)
Finally, for a convective envelope where $`R_{intCZ}=r_t`$ and $`R_{\mathrm{extCZ}}=R`$, we obtain:
$$\chi _0^l=\frac{1}{2l+1}\frac{1}{\eta _{\mathrm{turb}}}r_t^{l+1}_{r_t}^Rx^l𝒮_l\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)dx\text{at}r=r_t.$$
(86)
In the case where $`𝒮_l\left(\mathrm{\Omega }_{\mathrm{CZ}}\right)`$ is negligible, the previous results leads to:
$$\chi _0^l=0\text{ }\text{at}\text{ }r=r_b\text{and}r=r_t,\text{and thus to}B_T=0.$$
(87)
It remains to write down the boundary conditions to be applied on the transport of angular momentum. Since the equation for the mean angular velocity $`\overline{\mathrm{\Omega }}`$ (48) is of second order in $`r`$, it requires two boundary conditions. They are obtained by evaluating the budget of angular momentum in each adjacent convective zone:
$$\frac{\mathrm{d}}{\mathrm{d}t}\left[_0^{r_b}r^4\rho \mathrm{\Omega }dr\right]=\frac{1}{5}r^4\rho \overline{\mathrm{\Omega }}U_2_B(r_b)\text{at}r=r_b,\frac{\mathrm{d}}{\mathrm{d}t}\left[_{r_t}^Rr^4\rho \mathrm{\Omega }dr\right]=\frac{1}{5}r^4\rho \overline{\mathrm{\Omega }}U_2_\mathrm{\Omega }+_B(r_t)\text{at}r=r_t.$$
(88)
$`_\mathrm{\Omega }`$ is the (signed) flux of angular momentum which is lost at the surface by the stellar wind, and $`_B(r)`$ the flux carried by the magnetic field through the considered surface:
$$_B(r)=\frac{1}{\mu _0}r^3_0^\pi B_r(r,\theta )B_\phi (r,\theta )\mathrm{sin}^2\theta \mathrm{d}\theta =\frac{3}{2\mu _0}\underset{l>0}{}\left\{l\left(l+1\right)𝒩_l^0\xi _0^l\left[\frac{\left(l+1\right)\left(l+2\right)𝒩_{l+1}^0}{\left(2l+1\right)\left(2l+3\right)}\chi _0^{l+1}\frac{\left(l1\right)l𝒩_{l1}^0}{\left(2l1\right)\left(2l+1\right)}\chi _0^{l1}\right]\right\}.$$
(89)
Here again the perturbation $`\mathrm{\Omega }_2`$ obeys an evolution equation which does not include any derivative in $`r`$, and therefore it needs no boundary condition.
## 8 Conclusion
The work presented here is the continuation of Paper I (Mathis & Zahn 2004), where the star was assumed without magnetic field. Here we allow for a magnetic field of moderate strength in the radiation zone(s), with an Alfvén speed not exceeding the rotational velocity: $`(V_A)^2<(R\mathrm{\Omega })^2`$. Such a field has little impact on the hydrostatic balance, since we also assume, as in Paper I, that the centrifugal force is small compared to gravity: $`R\mathrm{\Omega }^2g`$. But it will compete with the centrifugal force in the baroclinic balance (56), and therefore it will participate in governing the meridional flow. Moreover, such a field tends to enforce uniform rotation along the field lines of the poloidal field (Ferraro’s law). The question then arises whether or not the poloidal field threads into the convection zone(s), where such differential rotation is probably maintained through the turbulent motions, as observed in the Sun. To answer that question one has to treat consistently the evolution of both angular velocity and magnetic field, with special care for the dynamics in the tachocline(s), and for this reason we expanded them in spherical harmonics to arbitrary order.
The overall problem of rotational mixing in a magnetized star is highly non-linear, with multiple feed-backs, as illustrated by the diagram displayed in fig. 1.
Mixing is achieved through the meridional circulation, which is due to the thermal imbalance caused by the centrifugal force and by the Lorentz force (mainly through the baroclinic balance), and also through the turbulence generated by the shear of differential rotation and possibly by the magnetic field. These motions modify both the rotation profile and the magnetic field through large-scale advection and turbulent diffusion. As the star evolves, molecular gradients build up through microscopic diffusion (including radiative levitation and gravitational settling), and through nuclear burning. Large-scale advection turns them into horizontal gradients, which react on the meridional circulation and are smoothed by the turbulence.
The boundary conditions play a crucial role: in a non-magnetic star, the circulation is driven by the loss (or gain) of angular momentum, as explained in Zahn (1992); in the absence of such loss, the circulation tends to vanish, as was shown by Busse (1982). It remains to be seen whether this is still the case when angular momentum is transported mainly through magnetic stresses, with the poloidal field anchored in the convection zone(s).
The main weakness of our modeling remains the description of the turbulence. As in Paper I, it is assumed to be anisotropic, due to the stable stratification, and that it tends to smooth angular velocity and chemical composition on horizontal surfaces. The action of magnetic field on such turbulence is not taken in account, which seems reasonable as long as the Alfvén speed does not exceed the rotational velocity. Moreover, we have deliberately ignored here the possibility of MHD instabilities generating their own turbulence (Balbus & Hawley 1994; Spruit 1999).
From the technical point of view, we adopted here the formalism developed in Paper I and introduced the spherical vectorial harmonics, which obey the Racah-Wigner angular momentum algebra widely used in quantum mechanics. This allowed us to separate the colatitude $`\theta `$ in the vectorial partial differential equations which govern the problem, and to reduce the problem to solve coupled partial differential equations in $`t\text{ and }r`$ only. It is then straightforward to introduce these equations in existing stellar structure codes, a task we have now undertaken.
Among the first applications we plan to model a solar-type star, to check whether a fossil magnetic field, which is able to extract angular momentum from the radiative interior as the star spins down, yields a profile of angular velocity that is compatible with the helioseismic data. Next we will undertake the modeling of magnetized massive stars, since they are progenitors of neutron stars which are the seat of strong magnetic fields (Maeder & Meynet 2003, 2004; Heger et al. 2004).
Most of the Appendix is dedicated to the algebra involving the spherical and vectorial spherical harmonics used in this paper. The numerical values of the coupling coefficients given in the Tables have been computed with Mathematica. In Appendix E, we state the equations in a form ready to be implemented in a stellar structure code, when the magnetic field is expanded up to the octupole ($`l=3`$) and only the first term $`\mathrm{\Omega }_2`$ of the differential rotation is retained.
###### Acknowledgements.
This work was supported by the Centre National de la Recherche Scientifique (Programme National de Physique Stellaire and GdR Dynamo).
## Appendix A Algebra related to the spherical harmonics
### A.1 Scalar quantities
#### A.1.1 Definition and basic properties
The spherical harmonics are defined by:
$$Y_l^m(\theta ,\phi )=𝒩_l^mP_l^{|m|}\left(\mathrm{cos}\theta \right)e^{im\phi },$$
(90)
where $`P_l^{|m|}\left(\mathrm{cos}\theta \right)e^{im\phi }`$ is the associated Legendre function, and the normalization coefficient being
$$𝒩_l^m=(1)^{\frac{\left(m+|m|\right)}{2}}\left[\frac{2l+1}{4\pi }\frac{(l|m|)!}{(l+|m|)!}\right]^{\frac{1}{2}}.$$
(91)
They obey the orthogonality relation:
$$_\mathrm{\Omega }\left(Y_{l_1}^{m_1}(\theta ,\phi )\right)^{}Y_{l_2}^{m_2}(\theta ,\phi )d\mathrm{\Omega }=\delta _{l_1,l_2}\delta _{m_1,m_2}$$
(92)
where $`\mathrm{d}\mathrm{\Omega }=\mathrm{sin}\theta \mathrm{d}\theta \mathrm{d}\phi `$ and where the complex conjugate spherical harmonic is given by:
$$\left(Y_l^m(\theta ,\phi )\right)^{}=\left(1\right)^mY_l^m(\theta ,\phi ).$$
(93)
Using these properties, every function $`f(\theta ,\phi )`$ can be expanded as:
$$f(\theta ,\phi )=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}f_m^lY_l^m(\theta ,\phi )\text{where}f_m^l=_\mathrm{\Omega }f(\theta ,\phi )\left(Y_l^m(\theta ,\phi )\right)^{}d\mathrm{\Omega }.$$
(94)
#### A.1.2 Special case of axisymmetric spherical harmonics
In that case we have:
$$Y_l^0(\theta ,\phi )=𝒩_l^0P_l\left(\mathrm{cos}\theta \right)\text{ and }_\theta Y_l^0(\theta ,\phi )=𝒩_l^0_\theta P_l\left(\mathrm{cos}\theta \right)=\{\begin{array}{cc}0\text{ for }l=0,\hfill & \\ 𝒩_l^0P_l^1\left(\mathrm{cos}\theta \right)\text{ for }l>0.\hfill & \end{array}$$
(95)
#### A.1.3 Linear differential and recursion relations
We recall that the spherical harmonics obey to the differential equation:
$$\frac{1}{\mathrm{sin}\theta }_\theta \left(\mathrm{sin}\theta _\theta Y_l^m(\theta ,\phi )\right)+\frac{1}{\mathrm{sin}^2\theta }_\phi ^2Y_l^m(\theta ,\phi )=l\left(l+1\right)Y_l^m(\theta ,\phi ).$$
(96)
In deriving the functions which are related to the centrifugal force in §5 (cf. (65-67-68)) we have used the following recursion relations for $`m=0`$:
$$\mathrm{cos}\theta Y_l^0(\theta )=A_l^0Y_{l1}^0(\theta )+B_l^0Y_{l+1}^0(\theta )\text{where}A_l^0=\frac{l}{\sqrt{(2l+1)(2l1)}}\text{and}B_l^0=\frac{(l+1)}{\sqrt{(2l+3)(2l+1)}},$$
(97)
$$\mathrm{sin}\theta Y_l^0(\theta )=C_l^0_\theta Y_{l1}^0(\theta )D_l^0_\theta Y_{l+1}^0(\theta )\text{where}C_l^0=\frac{1}{\sqrt{(2l+1)(2l1)}}\text{and}D_l^0=\frac{1}{\sqrt{(2l+3)(2l+1)}},$$
(98)
$$\mathrm{cos}\theta _\theta Y_l^0(\theta )=E_l^0_\theta Y_{l1}^0(\theta )+F_l^0_\theta Y_{l+1}^0(\theta )\text{where}E_l^0=\frac{l+1}{\sqrt{(2l+1)(2l1)}}\text{and}F_l^0=\frac{l}{\sqrt{(2l+3)(2l+1)}},$$
(99)
$$\mathrm{sin}\theta _\theta Y_l^0(\theta )=G_l^0Y_{l+1}^0(\theta )H_l^0Y_{l1}^0(\theta )\text{where}G_l^0=\frac{l(l+1)}{\sqrt{(2l+3)(2l+1)}}\text{and}H_l^0=\frac{l(l+1)}{\sqrt{(2l+1)(2l1)}}.$$
(100)
The identities (98) and (99) have been deduced from the two others (97-100) with the help of (96).
#### A.1.4 Expansion of products of the spherical harmonics
Using the normalization and the orthogonality of spherical harmonics (cf. (92-94)) and their complex conjugate (cf. (93)), we can write:
$$Y_{l_1}^{m_1}(\theta ,\phi )Y_{l_2}^{m_2}(\theta ,\phi )=(1)^{\left(m_1+m_2\right)}\underset{l=|l_1l_2|}{\overset{l_1+l_2}{}}_{l_1,l_2,l}^{m_1,m_2,\left(m_1+m_2\right)}Y_l^{m_1+m_2}(\theta ,\phi )$$
(101)
where we define the integral $`_{l_1,l_2,l}^{m_1,m_2,m}`$ like in Edmonds (1968) or Varshalovich and al. (1975):
$$_{l_1,l_2,l}^{m_1,m_2,m}=_\mathrm{\Omega }Y_{l_1}^{m_1}(\theta ,\phi )Y_{l_2}^{m_2}(\theta ,\phi )Y_l^m(\theta ,\phi )𝑑\mathrm{\Omega }=\sqrt{\frac{(2l_1+1)(2l_2+1)(2l+1)}{4\pi }}\left(\begin{array}{ccc}l_1& l_2& l\\ m_1& m_2& m\end{array}\right)\left(\begin{array}{ccc}l_1& l_2& l\\ 0& 0& 0\end{array}\right)$$
(102)
with the 3j-Wigner coefficients which are related to the classical Clebsch-Gordan coefficients by:
$$\left(\begin{array}{ccc}l_1& l_2& l\\ m_1& m_2& m\end{array}\right)=\frac{(1)^{l_1l_2m}}{\sqrt{2l+1}}C_{l_1,m_1,l_2,m_2}^{l,m}.$$
(103)
So finally, we get the following expansion for the product of two spherical harmonics:
$$Y_{l_1}^{m_1}(\theta ,\phi )Y_{l_2}^{m_2}(\theta ,\phi )=\underset{l=|l_1l_2|}{\overset{l_1+l_2}{}}c_{l_1,m_1,l_2,m_2}^lY_l^{m_1+m_2}(\theta ,\phi )$$
(104)
where the coefficient $`c_{l_1,m_1,l_2,m_2}^l`$ is given by:
$$c_{l_1,m_1,l_2,m_2}^l=\left(1\right)^{\left(m_1+m_2\right)}\sqrt{\frac{(2l_1+1)(2l_2+1)(2l+1)}{4\pi }}\left(\begin{array}{ccc}l_1& l_2& l\\ m_1& m_2& \left(m_1+m_2\right)\end{array}\right)\left(\begin{array}{ccc}l_1& l_2& l\\ 0& 0& 0\end{array}\right).$$
(105)
Then, using the initial definition of spherical harmonics (cf. 90-91), we deduce the expansion for the product of two associated Legendre functions:
$$P_{l_1}^{m_1}\left(\mathrm{cos}\theta \right)P_{l_2}^{m_2}\left(\mathrm{cos}\theta \right)=\underset{l=|l_1l_2|}{\overset{l_1+l_2}{}}d_{l_1,m_1,l_2,m_2}^lP_l^{m_1+m_2}\left(\mathrm{cos}\theta \right)$$
(106)
where
$$d_{l_1,m_1,l_2,m_2}^l=(1)^{\left(m_1+m_2\right)}\left(2l+1\right)\sqrt{\frac{\left(l_1+m_1\right)!\left(l_2+m_2\right)!\left(l\left(m_1+m_2\right)\right)!}{\left(l_1m_1\right)!\left(l_2m_2\right)!\left(l+\left(m_1+m_2\right)\right)!}}\left(\begin{array}{ccc}l_1& l_2& l\\ m_1& m_2& \left(m_1+m_2\right)\end{array}\right)\left(\begin{array}{ccc}l_1& l_2& l\\ 0& 0& 0\end{array}\right).$$
(107)
Note that the previous expression is symmetric; therefore
$$d_{l_1,m_1,l_2,m_2}^l=d_{l_2,m_2,l_1,m_1}^l.$$
(108)
### A.2 Vector fields
#### A.2.1 Definitions and basic properties
Following Rieutord (1987), we expand any vector field $`u(r,\theta ,\varphi )`$ in vectorial spherical harmonics as
$$u(r,\theta ,\varphi )=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\left\{u_m^l(r)R_l^m(\theta ,\phi )+v_m^l(r)S_l^m(\theta ,\phi )+w_m^l(r)T_l^m(\theta ,\phi )\right\},$$
(109)
where the vectorial spherical harmonics $`R_l^m(\theta ,\phi )`$, $`S_l^m(\theta ,\phi )`$, $`T_l^m(\theta ,\phi )`$ are defined as:
$$R_l^m(\theta ,\phi )=Y_l^m(\theta ,\phi )\widehat{e}_r\text{}S_l^m(\theta ,\phi )=_𝒮Y_l^m(\theta ,\phi )\text{ and }T_l^m(\theta ,\phi )=_𝒮R_l^m(\theta ,\phi ),$$
(110)
$$\text{with the horizontal gradient}_𝒮=\widehat{e}_\theta _\theta +\widehat{e}_\phi \frac{1}{\mathrm{sin}\theta }_\phi .$$
(111)
These vector functions obey the following orthogonality relations:
$$_\mathrm{\Omega }R_{l_1}^{m_1}S_{l_2}^{m_2}𝑑\mathrm{\Omega }=_\mathrm{\Omega }R_{l_1}^{m_1}T_{l_2}^{m_2}𝑑\mathrm{\Omega }=_\mathrm{\Omega }S_{l_1}^{m_1}T_{l_2}^{m_2}𝑑\mathrm{\Omega }=0,$$
(112)
$$_\mathrm{\Omega }R_{l_1}^{m_1}\left(R_{l_2}^{m_2}\right)^{}𝑑\mathrm{\Omega }=\delta _{l_1,l_2}\delta _{m_1,m_2}\text{ and }_\mathrm{\Omega }S_{l_1}^{m_1}\left(S_{l_2}^{m_2}\right)^{}𝑑\mathrm{\Omega }=_\mathrm{\Omega }T_{l_1}^{m_1}\left(T_{l_2}^{m_2}\right)^{}𝑑\mathrm{\Omega }=l_1(l_1+1)\delta _{l_1,l_2}\delta _{m_1,m_2}.$$
(113)
The vector function $`u`$ may also be projected on the classical spherical vectorial basis:
$$u=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\left\{u_m^l(r)Y_l^m(\theta ,\phi )\widehat{e}_r+\left[v_m^l(r)_\theta Y_l^m(\theta ,\phi )+w_m^l(r)\frac{im}{\mathrm{sin}\theta }Y_l^m(\theta ,\phi )\right]\widehat{e}_\theta +\left[v_m^l(r)\frac{im}{\mathrm{sin}\theta }Y_l^m(\theta ,\phi )w_m^l(r)_\theta Y_l^m(\theta ,\phi )\right]\widehat{e}_\phi \right\},$$
(114)
$`\widehat{e}_r`$, $`\widehat{e}_\theta `$ and $`\widehat{e}_\phi `$ beeing the unit-vectors respectively in the $`r`$, $`\theta `$ and $`\phi `$ directions. These expansions of vector fields allow us to separate explicitly the angular variables $`\theta `$ and $`\phi `$ in the vectorial partial differential equations which govern the problem. We thus reduce the problem to solve partial differential equations in $`t\text{ and }r`$ only, which are easy to implement in existing stellar structure codes. Note that $`R_l^m(\theta ,\phi )`$ and $`S_l^m(\theta ,\phi )`$ represent the poloidal part, and $`T_l^m(\theta ,\phi )`$ the toroidal part of $`u`$.
#### A.2.2 Expansions of differential operators
As stated in the previous section, the expansion of the vector fields in $`R_l^m(\theta ,\phi )`$, $`S_l^m(\theta ,\phi )`$ and $`T_l^m(\theta ,\phi )`$ allows us to separate the variables in the vectorial partial differential equations which govern the problem. This will prove particularly useful when dealing with magnetic field, and with the non-linear expressions where it is involved. We start by expanding the classical linear vectorial operators: gradient, divergence, curl and laplacian (scalar or vectorial).
Gradient:
Taking a scalar function $`f(r,\theta ,\phi )`$ expanded in the $`Y_l^m(\theta ,\phi )`$:
$$f(r,\theta ,\phi )=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}f_m^l(r)Y_l^m(\theta ,\phi ),$$
(115)
it is straightforward to derive its gradient:
$$f(r,\theta ,\phi )=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\left\{_rf_m^lR_l^m(\theta ,\phi )+\frac{f_m^l}{r}S_l^m(\theta ,\phi )\right\}.$$
(116)
Divergence:
Taking a vector field $`u`$ expanded as in (109), its divergence is given by:
$$u(r,\theta ,\phi )=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\left[\frac{1}{r^2}_r\left(r^2u_m^l\right)l(l+1)\frac{v_m^l}{r}\right]Y_l^m(\theta ,\phi ),$$
(117)
where whe note that, if $`u`$ is divergence-free, such as the magnetic field $`B`$ or the momentum density $`\rho 𝒰_M`$ in the anelastic approximation, we have the following relation between $`u_m^l`$ and $`v_m^l`$:
$$v_m^l=\frac{1}{l(l+1)}\frac{1}{r}_r\left(r^2u_m^l\right).$$
(118)
Laplacian of a scalar function:
With these expressions for the gradient of a scalar quantity and for the divergence of a vector field, one can easily derive the laplacian of a scalar function. Using the well-known property $`^2f=\left(f\right)`$, we get:
$$^2f=(f)=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\left[\frac{1}{r}_{r^2}\left(rf_m^l\right)l(l+1)\frac{f_m^l}{r^2}\right]Y_l^m(\theta ,\phi )=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\mathrm{\Delta }_lf_m^lY_l^m(\theta ,\phi )$$
(119)
where $`\mathrm{\Delta }_l`$ is the laplacian operator:
$$\mathrm{\Delta }_l=_{r,r}+\frac{2}{r}_r\frac{l\left(l+1\right)}{r^2}.$$
(120)
Curl:
This operator is found in the expression of the magnetic field in terms of a stream function, in that of the current density, etc. If we take $`u`$ expanded as in (109), we retrieve Rieutord’s (1987) result:
$$u=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\left\{\left[l(l+1)\frac{w_m^l}{r}\right]R_l^m(\theta ,\phi )+\left[\frac{1}{r}_r(rw_m^l)\right]S_l^m(\theta ,\phi )+\left[\frac{u_m^l}{r}\frac{1}{r}_r\left(rv_m^l\right)\right]T_l^m(\theta ,\phi )\right\}.$$
(121)
Taking this result, the expansion for the laplacian of a vector field is derived. First, we have:
$$\left(u\right)=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\left\{\left[\frac{l(l+1)}{r}\left(\frac{u_m^l}{r}\frac{1}{r}_r\left(rv_m^l\right)\right)\right]R_l^m(\theta ,\phi )+\left[\frac{1}{r}_ru_m^l\frac{1}{r}_{r,r}(rv_m^l)\right]S_l^m(\theta ,\phi )+\left[\mathrm{\Delta }_lw_m^l\right]T_l^m(\theta ,\phi )\right\}$$
(122)
and therefore
$`^2u=\left(u\right)\left(u\right)`$ (123)
$`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}\left\{\left[\mathrm{\Delta }_lu_m^l{\displaystyle \frac{2}{r^2}}\left(u_m^ll(l+1)v_m^l\right)\right]R_l^m(\theta ,\phi )+\left[\mathrm{\Delta }_lv_m^l+2{\displaystyle \frac{u_m^l}{r^2}}\right]S_l^m(\theta ,\phi )+\left[\mathrm{\Delta }_lw_m^l\right]T_l^m(\theta ,\phi )\right\},`$
result which becomes in the case where $`u`$ is divergence-free (cf. (118) ):
$`^2u=\left(u\right)`$ (124)
$`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}\left\{\left[{\displaystyle \frac{1}{r}}\mathrm{\Delta }_l(ru_m^l)\right]R_l^m(\theta ,\phi )+\left[{\displaystyle \frac{1}{r}}_r\left(r{\displaystyle \frac{\mathrm{\Delta }_l(ru_m^l)}{l(l+1)}}\right)\right]S_l^m(\theta ,\phi )+\left[\mathrm{\Delta }_lw_m^l\right]T_l^m(\theta ,\phi )\right\}.`$
Finally, using (121) once again, we get the ‘triple curl’ operator:
$$\left(\right)^3u=\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\left\{\left[l(l+1)\frac{\mathrm{\Delta }_lw_m^l}{r}\right]R_l^m(\theta ,\phi )+\left[\frac{1}{r}_r\left(r\mathrm{\Delta }_lw_m^l\right)\right]S_l^m(\theta ,\phi )+\left[\mathrm{\Delta }_lz_m^l\right]T_l^m(\theta ,\phi )\right\}$$
(125)
where:
$$z_m^l=\frac{1}{r}_r(rv_m^l)\frac{u_m^l}{r},$$
(126)
which becomes in the case where $`u`$ is divergence-free (cf. (118) ):
$`\left(\right)^3u={\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}\left\{\left[l(l+1){\displaystyle \frac{\mathrm{\Delta }_lw_m^l}{r}}\right]R_l^m(\theta ,\phi )\left[{\displaystyle \frac{1}{r}}_r\left(r\mathrm{\Delta }_lw_m^l\right)\right]S_l^m(\theta ,\phi )+\left[{\displaystyle \frac{1}{l(l+1)}}\mathrm{\Delta }_l\mathrm{\Delta }_l(ru_m^l)\right]T_l^m(\theta ,\phi )\right\}.`$
(127)
#### A.2.3 Products of axisymmetric vectorial spherical harmonics
Before we turn to the non-linear terms involving the magnetic field: the advection term in the induction equation and the Lorentz force, we shall first perform the projections of the scalar product of two general axisymmetric vectors and of their vector product, and derive the associated coupling coefficients.
Scalar product:
We take two general axisymmetric vectors $`X_1(r,\theta )`$ and $`X_2(r,\theta )`$, which we expand as in (109):
$$\{\begin{array}{cc}X_1(r,\theta )=_{l_1=0}^{\mathrm{}}\left\{𝒜_0^{l_1}(r)R_{l_1}^0\left(\theta \right)+_0^{l_1}(r)S_{l_1}^0\left(\theta \right)+𝒞_0^{l_1}(r)T_{l_1}^0\left(\theta \right)\right\}\hfill & \\ X_2(r,\theta )=_{l_2=0}^{\mathrm{}}\left\{𝒟_0^{l_2}(r)R_{l_2}^0\left(\theta \right)+_0^{l_2}(r)S_{l_2}^0\left(\theta \right)+_0^{l_2}(r)T_{l_2}^0\left(\theta \right)\right\},\hfill & \end{array}$$
(128)
and perform their scalar product. After some algebra involving the expansion of products of spherical harmonics (cf. §A.1.4.), we obtain
$$X_1(r,\theta )X_2(r,\theta )=\underset{l=0}{\overset{\mathrm{}}{}}𝒫_{\left(X_1X_2\right);l}\left(r\right)P_l\left(\mathrm{cos}\theta \right)$$
(129)
with the following expression for $`𝒫_{\left(X_1X_2\right);l}(r)`$:
$`𝒫_{\left(X_1X_2\right);l}(r)`$ $`=`$ $`{\displaystyle \underset{l_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_2=0}{\overset{\mathrm{}}{}}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left[𝒜_0^{l_1}\left(r\right)𝒟_0^{l_2}\left(r\right)\right]{\displaystyle \underset{j=I(l_1,0,l_2,0)}{\overset{l_1+l_2}{}}}d_{l_1,0,l_2,0}^j\delta _{l,j}+{\displaystyle \frac{2}{3}}\left[_0^{l_1}\left(r\right)_0^{l_2}\left(r\right)+𝒞_0^{l_1}\left(r\right)_0^{l_2}\left(r\right)\right]X_{l_1,l_2}^l\right\}.`$ (130)
Here $`X_{l_1,l_2}^l`$is defined in terms of the coupling coefficients $`d_{l_1,m_1,l_2,m_2}^l`$ (107):
$$X_{l_1,l_2}^l=\underset{j=I(l_1,1,l_2,1)}{\overset{l_1+l_2}{}}\left(d_{l_1,1,l_2,1}^j\underset{p=0}{\overset{E\left[\frac{j1}{2}\right]}{}}\left[\left(2j4p1\right)X\right]\right)$$
(131)
with
$$X=\underset{q=0}{\overset{E\left[\frac{\left(j2p1\right)1}{2}\right]}{}}\left\{\left[2\left(j2p1\right)4q1\right]\left(\delta _{l,\left[(j2p1)2q1\right]}\underset{r=I(2,0,\left[(j2p1)2q1\right],0)}{\overset{\left[(j2p1)2q1\right]+2}{}}\left[d_{2,0,\left[(j2p1)2q1\right],0}^r\delta _{l,r}\right]\right)\right\}.$$
(132)
We have used the classical notation $`\delta _{i,j}`$ for the usual Kronecker symbol, $`E\left[x\right]`$ is the integer part of $`x`$ and $`I(l_1,m_1,l_2,m_2)=\mathrm{max}(|l_1l_2|,m_1+m_2)`$.
Vector product:
We operate likewise for the vector product of two general axisymmetric vectors $`X_1(r,\theta )`$ and $`X_2(r,\theta )`$, again expanded as in (109). We reach the following result:
$`X_1(r,\theta )X_2(r,\theta )`$ $`=`$ $`𝒳_{\left(X_1X_2\right);0}(r)\widehat{e}_r+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{𝒳_{\left(X_1X_2\right);l}(r)P_l\left(\mathrm{cos}\theta \right)\widehat{e}_r+𝒴_{\left(X_1X_2\right);l}(r)P_l^1\left(\mathrm{cos}\theta \right)\widehat{e}_\theta +𝒵_{\left(X_1X_2\right);l}(r)P_l^1\left(\mathrm{cos}\theta \right)\widehat{e}_\phi \right\}`$ (133)
$`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left\{\left[{\displaystyle \frac{𝒳_{\left(X_1X_2\right);l}(r)}{𝒩_l^0}}\right]R_l^0\left(\theta \right)+\left[{\displaystyle \frac{𝒴_{\left(X_1X_2\right);l}(r)}{𝒩_l^0}}\right]S_l^0\left(\theta \right)+\left[{\displaystyle \frac{𝒵_{\left(X_1X_2\right);l}(r)}{𝒩_l^0}}\right]T_l^0\left(\theta \right)\right\},`$
where we have also used the following property of the Legendre function (Gradshteyn & Ryzhik 1967):
$$\frac{\mathrm{d}P_l(\mu )}{\mathrm{d}\mu }=\underset{k=0}{\overset{E\left(\frac{l1}{2}\right)}{}}\left(2l4k1\right)P_{l2k1}(\mu ).$$
(134)
The radial function $`𝒳_{\left(X_1X_2\right);l}\left(r\right)`$ is given by:
$$𝒳_{\left(X_1X_2\right);l}(r)=\frac{2}{3}\underset{l_1=0}{\overset{\mathrm{}}{}}\underset{l_2=0}{\overset{\mathrm{}}{}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left[𝒞_0^{l_1}(r)_0^{l_2}(r)_0^{l_1}(r)_0^{l_2}(r)\right]X_{l_1,l_2}^l\right\},$$
(135)
with the same coefficient $`X_{l_1,l_2}^l`$ derived above in (131). The two other radial functions $`𝒴_{\left(X_1X_2\right);l}\left(r\right)`$ and $`𝒵_{\left(X_1X_2\right);l}\left(r\right)`$ are given by
$$𝒴_{\left(X_1X_2\right);l}(r)=\underset{l_1=0}{\overset{\mathrm{}}{}}\underset{l_2=0}{\overset{\mathrm{}}{}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left[𝒞_0^{l_1}(r)𝒟_0^{l_2}(r)\right]\underset{j=I(l_1,1,l_2,0)}{\overset{l_1+l_2}{}}d_{l_1,1,l_2,0}^j\delta _{l,j}\left[𝒜_0^{l_1}(r)_0^{l_2}(r)\right]\underset{j=I(l_1,0,l_2,1)}{\overset{l_1+l_2}{}}d_{l_1,0,l_2,1}^j\delta _{l,j}\right\}$$
(136)
and
$$𝒵_{\left(X_1X_2\right);l}(r)=\underset{l_1=0}{\overset{\mathrm{}}{}}\underset{l_2=0}{\overset{\mathrm{}}{}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left[_0^{l_1}(r)𝒟_0^{l_2}(r)\right]\underset{j=I(l_1,1,l_2,0)}{\overset{l_1+l_2}{}}d_{l_1,1,l_2,0}^j\delta _{l,j}\left[𝒜_0^{l_1}(r)_0^{l_2}(r)\right]\underset{j=I(l_1,0,l_2,1)}{\overset{l_1+l_2}{}}d_{l_1,0,l_2,1}^j\delta _{l,j}\right\}.$$
(137)
## Appendix B Lorentz force
In this section we shall project the Lorentz force,
$$_{}=jB=\left[\frac{1}{\mu _0}(B)\right]B,$$
(138)
on spherical vectorial harmonics, and write the result in terms of the poloidal and the toroidal magnetic stream-functions, respectively $`\xi _0^l`$ and $`\chi _0^l`$. We have just seen in §A.2.3. how to expand the vector product of two axisymmetric vectors in the $`R_l^0(\theta )`$, $`S_l^0(\theta )`$, $`T_l^0(\theta )`$. We apply the method to (138), with (cf. (14), (15))
$`X_1(r,\theta )`$ $`=`$ $`j(r,\theta )={\displaystyle \frac{1}{\mu _0}}\left(B(r,\theta )\right)={\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}\left\{𝒜_0^{l_1}(r)R_{l_1}^0(\theta )+_0^{l_1}(r)S_{l_1}^0(\theta )+𝒞_0^{l_1}(r)T_{l_1}^0(\theta )\right\}`$
$`=`$ $`{\displaystyle \frac{1}{\mu _0}}{\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}\left\{\left[l_1(l_1+1){\displaystyle \frac{\chi _0^{l_1}}{r^2}}\right]R_{l_1}^0(\theta )+\left[{\displaystyle \frac{1}{r}}_r\chi _0^{l_1}\right]S_{l_1}^0(\theta )+\left[\mathrm{\Delta }_{l_1}\left({\displaystyle \frac{\xi _0^{l_1}}{r}}\right)\right]T_{l_1}^0(\theta )\right\}`$
and (cf. (11))
$`X_2(r,\theta )`$ $`=`$ $`B(r,\theta )={\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}\left\{𝒟_0^{l_2}(r)R_{l_2}^0(\theta )+_0^{l_2}(r)S_{l_2}^0(\theta )+_0^{l_2}(r)T_{l_2}^0(\theta )\right\}`$
$`=`$ $`{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}\left\{\left[l_2(l_2+1){\displaystyle \frac{\xi _0^{l_2}}{r^2}}\right]R_{l_2}^0(\theta )+\left[{\displaystyle \frac{1}{r}}_r\xi _0^{l_2}\right]S_{l_2}^0(\theta )+\left[{\displaystyle \frac{\chi _0^{l_2}}{r}}\right]T_{l_2}^0(\theta )\right\}.`$
So, we have, using the definition of §A.2.3.:
$$\{\begin{array}{cc}𝒜_0^{l_1}=\frac{1}{\mu _0}\left[l_1(l_1+1)\frac{\chi _0^{l_1}}{r^2}\right]\hfill & \\ _0^{l_1}=\frac{1}{\mu _0}\left[\frac{1}{r}_r\chi _0^{l_1}\right]\hfill & \\ 𝒞_0^{l_1}=\frac{1}{\mu _0}\left[\mathrm{\Delta }_{l_1}\left(\frac{\xi _0^{l_1}}{r}\right)\right]\hfill & \end{array}\text{and}\{\begin{array}{cc}𝒟_0^{l_2}=\left[l_2(l_2+1)\frac{\xi _0^{l_2}}{r^2}\right]\hfill & \\ _0^{l_2}=\left[\frac{1}{r}_r\xi _0^{l_2}\right]\hfill & \\ _0^{l_2}=\left[\frac{\chi _0^{l_2}}{r}\right].\hfill & \end{array}$$
(141)
Therefore, we get using (133-135-136-137):
$`_{}(r,\theta )=X_1(r,\theta )X_2(r,\theta )`$ (142)
$`=`$ $`𝒳_{_{};0}(r)\widehat{e}_r+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{𝒳_{_{};l}(r)P_l\left(\mathrm{cos}\theta \right)\widehat{e}_r+𝒴_{_{};l}(r)P_l^1\left(\mathrm{cos}\theta \right)\widehat{e}_\theta +𝒵_{_{};l}(r)P_l^1\left(\mathrm{cos}\theta \right)\widehat{e}_\phi \right\}`$
$`=`$ $`\left[{\displaystyle \frac{𝒳_{_{};0}}{𝒩_0^0}}\right]R_0^0(\theta )+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{\left[{\displaystyle \frac{𝒳_{_{};l}(r)}{𝒩_l^0}}\right]R_l^0\left(\theta \right)+\left[{\displaystyle \frac{𝒴_{_{};l}(r)}{𝒩_l^0}}\right]S_l^0\left(\theta \right)+\left[{\displaystyle \frac{𝒵_{_{};l}(r)}{𝒩_l^0}}\right]T_l^0\left(\theta \right)\right\},`$
where the radial functions $`𝒳_{_{};l}\left(r\right)`$, $`𝒴_{_{};l}\left(r\right)`$ and $`𝒵_{_{};l}\left(r\right)`$ are:
$$𝒳_{_{};l}=\frac{2}{3}\frac{1}{\mu _0}\underset{l_1=1}{\overset{\mathrm{}}{}}\underset{l_2=1}{\overset{\mathrm{}}{}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left(\left[\mathrm{\Delta }_{l_1}\left(\frac{\xi _0^{l_1}}{r}\right)\right]\left[\frac{1}{r}_r\xi _0^{l_2}\right]\left[\frac{1}{r}_r\chi _0^{l_1}\right]\left[\frac{\chi _0^{l_2}}{r}\right]\right)X_{l_1,l_2}^l\right\}$$
(143)
$$𝒴_{_{};l}=\frac{1}{\mu _0}\underset{l_1=1}{\overset{\mathrm{}}{}}\underset{l_2=1}{\overset{\mathrm{}}{}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left(\left[\mathrm{\Delta }_{l_1}\left(\frac{\xi _0^{l_1}}{r}\right)\right]\left[l_2(l_2+1)\frac{\xi _0^{l_2}}{r^2}\right]\right)\underset{j=I(l_1,1,l_2,0)}{\overset{l_1+l_2}{}}d_{l_1,1,l_2,0}^j\delta _{l,j}\left(\left[l_1(l_1+1)\frac{\chi _0^{l_1}}{r^2}\right]\left[\frac{\chi _0^{l_2}}{r}\right]\right)\underset{j=I(l_1,0,l_2,1)}{\overset{l_1+l_2}{}}d_{l_1,0,l_2,1}^j\delta _{l,j}\right\}$$
(144)
$$𝒵_{_{};l}=\frac{1}{\mu _0}\underset{l_1=1}{\overset{\mathrm{}}{}}\underset{l_2=1}{\overset{\mathrm{}}{}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left(\left[\frac{1}{r}_r\chi _0^{l_1}\right]\left[l_2(l_2+1)\frac{\xi _0^{l_2}}{r^2}\right]\right)\underset{j=I(l_1,1,l_2,0)}{\overset{l_1+l_2}{}}d_{l_1,1,l_2,0}^j\delta _{l,j}\left(\left[l_1(l_1+1)\frac{\chi _0^{l_1}}{r^2}\right]\left[\frac{1}{r}_r\xi _0^{l_2}\right]\right)\underset{j=I(l_1,0,l_2,1)}{\overset{l_1+l_2}{}}d_{l_1,0,l_2,1}^j\delta _{l,j}\right\}.$$
(145)
Explicit values for $`X_{l_1,l_2}^l`$ (131) are given in Table 1, and in Table 2 for the coupling coefficients $`d_{l_1,1,l_2,0}^l`$.
## Appendix C Advection term of the induction equation
We deal likewise with the advection term in the induction equation:
$$\left(𝒰_\phi +𝒰_M\right)B,$$
(146)
where $`𝒰_\phi (r,\theta )=r\mathrm{sin}\theta \mathrm{\Omega }(r,\theta )\widehat{e}_\phi `$ and $`𝒰_M`$ represents the meridional flow, whose expansions are found in (30). Therefore
$`X_1(r,\theta )`$ $`=`$ $`𝒰_\phi (r,\theta )+𝒰_M(r,\theta )={\displaystyle \underset{l_1=0}{\overset{\mathrm{}}{}}}\left\{𝒜_0^{l_1}(r)R_{l_1}^0(\theta )+_0^{l_1}(r)S_{l_1}^0(\theta )+𝒞_0^{l_1}(r)T_{l_1}^0(\theta )\right\}`$
$`=`$ $`{\displaystyle \underset{l_1=0}{\overset{\mathrm{}}{}}}\left\{u_0^{l_1}\left(r\right)R_{l_1}^0(\theta )+v_0^{l_1}\left(r\right)S_{l_1}^0(\theta )+w_0^{l_1}\left(r\right)T_{l_1}^0(\theta )\right\}`$
and (cf. (11))
$`X_2(r,\theta )`$ $`=`$ $`B(r,\theta )={\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}\left\{𝒟_0^{l_2}(r)R_{l_2}^0(\theta )+_0^{l_2}(r)S_{l_2}^0(\theta )+_0^{l_2}(r)T_{l_2}^0(\theta )\right\}`$
$`=`$ $`{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}\left\{\left[l_2(l_2+1){\displaystyle \frac{\xi _0^{l_2}}{r^2}}\right]R_{l_2}^0(\theta )+\left[{\displaystyle \frac{1}{r}}_r\xi _0^{l_2}\right]S_{l_2}^0(\theta )+\left[{\displaystyle \frac{\chi _0^{l_2}}{r}}\right]T_{l_2}^0(\theta )\right\}.`$
So, we have, expliciting the velocity field:
$$\{\begin{array}{cc}𝒜_0^{l_1}=u_0^{l_1}=\frac{U_{l_1}}{𝒩_{l_1}^0}\hfill & \\ _0^{l_1}=v_0^{l_1}=\frac{V_{l_1}}{𝒩_{l_1}^0}\hfill & \\ 𝒞_0^{l_1}=w_0^{l_1}=r\left[\frac{D_{l_11}^0}{𝒩_{l_11}^0}\mathrm{\Omega }_{l_11}^{}\frac{C_{l_1+1}^0}{𝒩_{l_1+1}^0}\mathrm{\Omega }_{l_1+1}^{}\right]\hfill & \end{array}\text{where }\{\begin{array}{cc}\mathrm{\Omega }_0^{}\left(r\right)=\mathrm{\Omega }_0\left(r\right)+\frac{1}{5}\mathrm{\Omega }_2\left(r\right)\hfill & \\ \mathrm{\Omega }_l^{}\left(r\right)=\mathrm{\Omega }_l\left(r\right)\text{ for }l>0\hfill & \end{array}\text{and }\{\begin{array}{cc}𝒟_0^{l_2}=\left[l_2(l_2+1)\frac{\xi _0^{l_2}}{r^2}\right]\hfill & \\ _0^{l_2}=\left[\frac{1}{r}_r\xi _0^{l_2}\right]\hfill & \\ _0^{l_2}=\left[\frac{\chi _0^{l_2}}{r}\right].\hfill & \end{array}$$
(149)
Hence we get, using (133-135-136-137):
$`\left(𝒰_\phi (r,\theta )+𝒰_M(r,\theta )\right)B(r,\theta )=X_1(r,\theta )X_2(r,\theta )`$ (150)
$`=`$ $`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{𝒳_{\mathrm{𝐀𝐝};\mathrm{l}}(r)P_l\left(\mathrm{cos}\theta \right)\widehat{e}_r+𝒴_{\mathrm{𝐀𝐝};\mathrm{l}}(r)P_l^1\left(\mathrm{cos}\theta \right)\widehat{e}_\theta +𝒵_{\mathrm{𝐀𝐝};\mathrm{l}}(r)P_l^1\left(\mathrm{cos}\theta \right)\widehat{e}_\phi \right\}`$
$`=`$ $`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left\{\left[{\displaystyle \frac{𝒳_{\mathrm{𝐀𝐝};\mathrm{l}}(r)}{𝒩_l^0}}\right]R_l^0\left(\theta \right)+\left[{\displaystyle \frac{𝒴_{\mathrm{𝐀𝐝};\mathrm{l}}(r)}{𝒩_l^0}}\right]S_l^0\left(\theta \right)+\left[{\displaystyle \frac{𝒵_{\mathrm{𝐀𝐝};\mathrm{l}}(r)}{𝒩_l^0}}\right]T_l^0\left(\theta \right)\right\},`$
where $`𝒳_{\mathrm{𝐀𝐝};\mathrm{l}}\left(r\right)`$, $`𝒴_{\mathrm{𝐀𝐝};\mathrm{l}}\left(r\right)`$ and $`𝒵_{\mathrm{𝐀𝐝};\mathrm{l}}\left(r\right)`$:
$$𝒳_{\mathrm{𝐀𝐝};\mathrm{l}}=\frac{2}{3}\underset{l_1=0}{\overset{\mathrm{}}{}}\underset{l_2=1}{\overset{\mathrm{}}{}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left(w_0^{l_1}\left[\frac{1}{r}_r\xi _0^{l_2}\right]v_0^{l_1}\left[\frac{\chi _0^{l_2}}{r}\right]\right)X_{l_1,l_2}^l\right\},$$
(151)
$$𝒴_{\mathrm{𝐀𝐝};l}=\underset{l_1=0}{\overset{\mathrm{}}{}}\underset{l_2=1}{\overset{\mathrm{}}{}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left(w_0^{l_1}\left[l_2(l_2+1)\frac{\xi _0^{l_2}}{r^2}\right]\right)\underset{j=I(l_1,1,l_2,0)}{\overset{l_1+l_2}{}}d_{l_1,1,l_2,0}^j\delta _{l,j}\left(u_0^{l_1}\left[\frac{\chi _0^{l_2}}{r}\right]\right)\underset{j=I(l_1,0,l_2,1)}{\overset{l_1+l_2}{}}d_{l_1,0,l_2,1}^j\delta _{l,j}\right\},$$
(152)
$`𝒵_{\mathrm{𝐀𝐝};l}`$ $`=`$ $`{\displaystyle \underset{l_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}𝒩_{l_1}^0𝒩_{l_2}^0\left\{\left(v_0^{l_1}\left[l_2(l_2+1){\displaystyle \frac{\xi _0^{l_2}}{r^2}}\right]\right){\displaystyle \underset{j=I(l_1,1,l_2,0)}{\overset{l_1+l_2}{}}}d_{l_1,1,l_2,0}^j\delta _{l,j}\left(u_0^{l_1}\left[{\displaystyle \frac{1}{r}}_r\xi _0^{l_2}\right]\right){\displaystyle \underset{j=I(l_1,0,l_2,1)}{\overset{l_1+l_2}{}}}d_{l_1,0,l_2,1}^j\delta _{l,j}\right\};`$ (153)
Explicit values for $`X_{l_1,l_2}^l`$ (131) are given in Table 3, and in Tables 4, 5, 6 for the coupling coefficients $`d_{l_1,1,l_2,0}^l`$.
## Appendix D Ohmic heating
For sake of completeness, we also calculate the Ohmic heating rate, which we shall express in terms of the magnetic stream functions $`\xi _0^l`$ and $`\chi _0^l`$. From §2.4, we get:
$$𝒥=\frac{1}{\mu _0}\left[\eta \left(B\right)\right]\left(B\right),$$
(154)
where we allow for different eddy-diffusivities ($`\eta _h,\eta _v`$) respectively in the horizontal and vertical direction. We apply again the method of §A.2.3., and identify the two vector functions which enter the scalar product:
$`X_1(r,\theta )`$ $`=`$ $`{\displaystyle \frac{1}{\mu _0}}\left[\eta \left(B(r,\theta )\right)\right]={\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}\left\{𝒜_0^{l_1}(r)R_{l_1}^0\left(\theta \right)+_0^{l_1}\left(r\right)S_{l_1}^0\left(\theta \right)+𝒞_0^{l_1}\left(r\right)T_{l_1}^0\left(\theta \right)\right\}`$
$`=`$ $`{\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{\eta _v}{\mu _0}}\left[l_1(l_1+1){\displaystyle \frac{\chi _0^{l_1}}{r^2}}\right]R_{l_1}^0\left(\theta \right)+{\displaystyle \frac{\eta _h}{\mu _0}}\left[{\displaystyle \frac{1}{r}}_r\chi _0^{l_1}\right]S_{l_1}^0\left(\theta \right)+{\displaystyle \frac{\eta _h}{\mu _0}}\left[\mathrm{\Delta }_{l_1}\left({\displaystyle \frac{\xi _0^{l_1}}{r}}\right)\right]T_{l_1}^0\left(\theta \right)\right\}`$
and
$`X_2(r,\theta )`$ $`=`$ $`B(r,\theta )={\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}\left\{𝒟_0^{l_2}(r)R_{l_2}^0\left(\theta \right)+_0^{l_2}\left(r\right)S_{l_2}^0\left(\theta \right)+_0^{l_2}\left(r\right)T_{l_2}^0\left(\theta \right)\right\}`$
$`=`$ $`{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}\left\{\left[l_2(l_2+1){\displaystyle \frac{\chi _0^{l_2}}{r^2}}\right]R_{l_2}^0\left(\theta \right)+\left[{\displaystyle \frac{1}{r}}_r\chi _0^{l_2}\right]S_{l_2}^0\left(\theta \right)+\left[\mathrm{\Delta }_{l_2}\left({\displaystyle \frac{\xi _0^{l_2}}{r}}\right)\right]T_{l_2}^0\left(\theta \right)\right\}.`$
We thus obtain, using (129-130):
$$𝒥(r,\theta )=X_1(r,\theta )X_2(r,\theta )=\underset{l=0}{\overset{\mathrm{}}{}}𝒥_l\left(r\right)P_l\left(\mathrm{cos}\theta \right)$$
(157)
with the following expression for $`𝒥_l\left(r\right)`$:
$`𝒥_l`$ $`=`$ $`{\displaystyle \frac{1}{\mu _0}}{\displaystyle \underset{l_1=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l_2=1}{\overset{\mathrm{}}{}}}𝒩_{l_1}^0𝒩_{l_2}^0\{\left({\displaystyle \frac{\eta _v}{\mu _0}}\left[l_1(l_1+1){\displaystyle \frac{\chi _0^{l_1}}{r^2}}\right]\left[l_2(l_2+1){\displaystyle \frac{\chi _0^{l_2}}{r^2}}\right]\right){\displaystyle \underset{j=I(l_1,0,l_2,0)}{\overset{l_1+l_2}{}}}d_{l_1,0,l_2,0}^j\delta _{l,j}`$
$`+`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{\eta _h}{\mu _0}}(\left[{\displaystyle \frac{1}{r}}_r\chi _0^{l_1}\right]\left[{\displaystyle \frac{1}{r}}_r\chi _0^{l_2}\right]+\left[\mathrm{\Delta }_{l_1}\left({\displaystyle \frac{\xi _0^{l_1}}{r}}\right)\right]\left[\mathrm{\Delta }_{l_2}\left({\displaystyle \frac{\xi _0^{l_2}}{r}}\right)\right])X_{l_1,l_2}^l\},`$
recalling that $`X_{l_1,l_2}^l`$ is given in (131), and in explicit form in Table 1.
## Appendix E Equations to be implemented in stellar evolution codes
In this section we shall give the equations ready to be implemented in a stellar structure code, with all coupling coefficients being replaced by their explicit value, in the special case where only the dipole ($`l=1`$), quadrupole ($`l=2`$) and octupole ($`l=3`$) are kept in the magnetic field, and where the differental rotation is reduced to its first term (l=2).
### E.1 Induction equation
The linear terms of this equation readily separate into their poloidal and toroidal components, and they project on a single multipole, whereas the advection term is the result of various couplings.
#### E.1.1 Equations for the dipole
The induction equation translates into two evolution equations for the magnetic stream functions:
$$\{\begin{array}{cc}\frac{\mathrm{d}\xi _0^1}{\mathrm{d}t}=2\sqrt{\frac{\pi }{3}}r𝒵_{\mathrm{𝐀𝐝};1}+\eta _hr\mathrm{\Delta }_1\left(\frac{\xi _0^1}{r}\right)\hfill & \\ \frac{\mathrm{d}\chi _0^1}{\mathrm{d}t}+_r\left(\dot{r}\right)\chi _0^1=2\sqrt{\frac{\pi }{3}}\left[𝒳_{\mathrm{𝐀𝐝};1}+_r\left(r𝒴_{\mathrm{𝐀𝐝};1}\right)\right]+\left[_r\left(\eta _h_r\chi _0^1\right)2\eta _v\frac{\chi _0^1}{r^2}\right]\hfill & \end{array}\text{ }\text{where}\text{ }\mathrm{\Delta }_1=_{r,r}+\frac{2}{r}_r\frac{2}{r^2},$$
(159)
and where the advective terms are given by:
$$\{\begin{array}{cc}𝒳_{\mathrm{𝐀𝐝};1}=\frac{2}{3}\frac{1}{2\pi }\left[\frac{13}{2}\sqrt{\frac{3}{5}}\left(𝒞_0^1_0^2_0^2_0^1\right)+\frac{183}{7}\sqrt{\frac{5}{7}}\left(𝒞_0^3_0^2_0^2_0^3\right)\frac{1}{35}\left(123\sqrt{3}_0^4_0^1+\frac{5391}{\sqrt{7}}_0^4_0^3\right)\right]\hfill & \\ 𝒴_{\mathrm{𝐀𝐝};1}=\frac{1}{4\pi }\left[\sqrt{\frac{3}{5}}\left(𝒜_0^2_0^1𝒞_0^1𝒟_0^2\right)+\frac{18}{\sqrt{35}}\left(𝒞_0^3𝒟_0^2𝒜_0^2_0^3\right)+\frac{6}{\sqrt{7}}𝒜_0^4_0^3\right]\hfill & \\ 𝒵_{\mathrm{𝐀𝐝};1}=\frac{1}{4\pi }\left[\sqrt{\frac{3}{5}}\left(3_0^2𝒟_0^1+𝒜_0^2_0^1\right)\frac{9}{\sqrt{35}}\left(_0^2𝒟_0^3+2𝒜_0^2_0^3\right)+\frac{2}{\sqrt{7}}\left(5_0^4𝒟_0^3+3𝒜_0^4_0^3\right)\right]\hfill & \end{array}$$
(160)
with
$$\{\begin{array}{cc}𝒜_0^2=2\sqrt{\frac{\pi }{5}}U_2\hfill & \\ 𝒜_0^4=\frac{2}{3}\sqrt{\pi }U_4\hfill & \\ _0^2=2\sqrt{\frac{\pi }{5}}V_2\hfill & \\ _0^4=\frac{2}{3}\sqrt{\pi }V_4\hfill & \\ 𝒞_0^1=w_0^1=2\sqrt{\frac{\pi }{3}}r\overline{\mathrm{\Omega }}\hfill & \\ 𝒞_0^3=w_0^3=\frac{2}{5}\sqrt{\frac{\pi }{7}}r\mathrm{\Omega }_2\hfill & \end{array}\text{,}\text{ }\{\begin{array}{cc}𝒟_0^1=\left[2\frac{\xi _0^1}{r^2}\right]\hfill & \\ _0^1=\left[\frac{1}{r}_r\xi _0^1\right]\hfill & \\ _0^1=\left[\frac{\chi _0^1}{r}\right]\hfill & \end{array}\text{,}\text{ }\{\begin{array}{cc}𝒟_0^2=\left[6\frac{\xi _0^2}{r^2}\right]\hfill & \\ _0^2=\left[\frac{1}{r}_r\xi _0^2\right]\hfill & \\ _0^2=\left[\frac{\chi _0^2}{r}\right]\hfill & \end{array}\text{and}\text{ }\text{ }\{\begin{array}{cc}𝒟_0^3=\left[12\frac{\xi _0^3}{r^2}\right]\hfill & \\ _0^3=\left[\frac{1}{r}_r\xi _0^3\right]\hfill & \\ _0^3=\left[\frac{\chi _0^3}{r}\right].\hfill & \end{array}$$
(161)
#### E.1.2 Equations for the quadrupole
Likewise
$$\{\begin{array}{cc}\frac{\mathrm{d}\xi _0^2}{\mathrm{d}t}=2\sqrt{\frac{\pi }{5}}r𝒵_{\mathrm{𝐀𝐝};2}+\eta _hr\mathrm{\Delta }_2\left(\frac{\xi _0^2}{r}\right)\hfill & \\ \frac{\mathrm{d}\chi _0^2}{\mathrm{d}t}+_r\left(\dot{r}\right)\chi _0^2=2\sqrt{\frac{\pi }{5}}\left[𝒳_{\mathrm{𝐀𝐝};2}+_r\left(r𝒴_{\mathrm{𝐀𝐝};2}\right)\right]+\left[_r\left(\eta _h_r\chi _0^2\right)6\eta _v\frac{\chi _0^2}{r^2}\right]\hfill & \end{array}\text{ }\text{where}\text{ }\mathrm{\Delta }_2=_{r,r}+\frac{2}{r}_r\frac{6}{r^2}$$
(162)
and
$$\{\begin{array}{cc}𝒳_{\mathrm{𝐀𝐝};2}=\frac{2}{3}\frac{1}{2\pi }\left[13\sqrt{\frac{3}{7}}\left(𝒞_0^3_0^1+𝒞_0^1_0^3\right)\frac{3}{2}𝒞_0^1_0^1+\frac{239}{7}𝒞_0^3_0^3\frac{93}{14}_0^2_0^2\frac{113}{7}\sqrt{5}_0^4_0^2\right]\hfill & \\ 𝒴_{\mathrm{𝐀𝐝};2}=\frac{1}{4\pi }\left[\sqrt{\frac{3}{7}}\left(4𝒞_0^3𝒟_0^1𝒞_0^1𝒟_0^3\right)+𝒞_0^1𝒟_0^1\frac{5}{7}𝒜_0^2_0^2+\frac{4}{7}\sqrt{5}𝒜_0^4_0^2+\frac{2}{3}𝒞_0^3𝒟_0^3\right]\hfill & \\ 𝒵_{\mathrm{𝐀𝐝};2}=\frac{1}{4\pi }\left[\frac{5}{7}\left(_0^2𝒟_0^2𝒜_0^2_0^2\right)+\frac{2}{7}\sqrt{5}\left(5_0^4𝒟_0^2+2𝒜_0^4_0^2\right)\right].\hfill & \end{array}$$
(163)
All functions $`𝒜_0^2`$ . . . $`_0^3`$ have the same meaning as in (161).
#### E.1.3 Equations for the octupole
The result is similar for the octupole:
$$\{\begin{array}{cc}\frac{\mathrm{d}\xi _0^3}{\mathrm{d}t}=2\sqrt{\frac{\pi }{7}}r𝒵_{\mathrm{𝐀𝐝};3}+\eta _hr\mathrm{\Delta }_3\left(\frac{\xi _0^3}{r}\right)\hfill & \\ \frac{\mathrm{d}\chi _0^3}{\mathrm{d}t}+_r\left(\dot{r}\right)\chi _0^3=2\sqrt{\frac{\pi }{7}}\left[𝒳_{\mathrm{𝐀𝐝};3}+_r\left(r𝒴_{\mathrm{𝐀𝐝};3}\right)\right]+\left[_r\left(\eta _h_r\chi _0^3\right)12\eta _v\frac{\chi _0^2}{r^2}\right]\hfill & \end{array}\text{ }\text{where}\text{ }\mathrm{\Delta }_3=_{r,r}+\frac{2}{r}_r\frac{12}{r^2}$$
(164)
and
$$\{\begin{array}{cc}𝒳_{\mathrm{𝐀𝐝};3}=\frac{2}{3}\frac{1}{\pi }\left[\frac{3}{4}\sqrt{\frac{3}{5}}\left(𝒞_0^1_0^2+_0^2_0^1\right)+12\sqrt{\frac{5}{7}}\left(𝒞_0^3_0^2_0^2_0^3\right)\frac{23}{5}\sqrt{3}_0^4_0^1\frac{387041}{3630}\frac{1}{\sqrt{7}}_0^4_0^3\right]\hfill & \\ 𝒴_{\mathrm{𝐀𝐝};3}=\frac{1}{4\pi }\left[\sqrt{\frac{3}{5}}\left(𝒞_0^1𝒟_0^2𝒜_0^2_0^1\right)+\sqrt{\frac{7}{5}}\left(𝒞_0^3𝒟_0^2𝒜_0^2_0^3\right)+\frac{1}{\sqrt{3}}𝒜_0^4_0^1\frac{\sqrt{7}}{11}𝒜_0^4_0^3\right]\hfill & \\ 𝒵_{\mathrm{𝐀𝐝};3}=\frac{1}{4\pi }\left[\sqrt{\frac{3}{5}}\left(2_0^2𝒟_0^1𝒜_0^2_0^1\right)+\frac{1}{\sqrt{3}}\left(5_0^4𝒟_0^1+𝒜_0^4_0^1\right)+\sqrt{\frac{7}{5}}\left(\frac{1}{3}_0^2𝒟_0^3𝒜_0^2_0^3\right)+\frac{\sqrt{7}}{11}\left(5_0^4𝒟_0^3𝒜_0^4_0^3\right)\right].\hfill & \end{array}$$
(165)
### E.2 Mean equations
#### E.2.1 Mean rotation rate
We recall the transport equation for the mean angular momentum (48):
$$\rho \frac{\mathrm{d}}{\mathrm{d}t}(r^2\overline{\mathrm{\Omega }})=\frac{1}{5r^2}_r\left(\rho r^4\overline{\mathrm{\Omega }}U_2\right)+\frac{1}{r^2}_r\left(\rho \nu _vr^4_r\overline{\mathrm{\Omega }}\right)+\overline{\mathrm{\Gamma }}_{_{}}$$
(166)
where the mean magnetic torque is given by:
$$\overline{\mathrm{\Gamma }}_{_{}}=\mathrm{\Gamma }_0\frac{1}{5}\mathrm{\Gamma }_2\text{ }\text{with}\text{ }\{\begin{array}{cc}\mathrm{\Gamma }_0=r\left(𝒵_{_{};1}+𝒵_{_{};3}+𝒵_{_{};5}\right)\hfill & \\ \mathrm{\Gamma }_2=5r\left(𝒵_{_{};3}+𝒵_{_{};5}\right)\hfill & \end{array}$$
(167)
with
$$\{\begin{array}{cc}𝒵_{_{};1}=\frac{1}{4\pi }\left\{\sqrt{\frac{3}{5}}\left[3\left(_0^2𝒟_0^1𝒜_0^1_0^2\right)\left(_0^1𝒟_0^2𝒜_0^2_0^1\right)\right]+\frac{9}{\sqrt{35}}\left[2\left(_0^3𝒟_0^2𝒜_0^2_0^3\right)\left(_0^2𝒟_0^3𝒜_0^3_0^2\right)\right]\right\}\hfill & \\ 𝒵_{_{};3}=\frac{1}{4\pi }\left\{\sqrt{\frac{3}{5}}\left[2\left(_0^2𝒟_0^1𝒜_0^1_0^2\right)+\left(_0^1𝒟_0^2𝒜_0^2_0^1\right)\right]+\sqrt{\frac{7}{5}}\left[\left(_0^3𝒟_0^2𝒜_0^2_0^3\right)+\frac{1}{3}\left(_0^2𝒟_0^3𝒜_0^3_0^2\right)\right]\right\}\hfill & \\ 𝒵_{_{};5}=\frac{1}{2\pi }\sqrt{\frac{5}{7}}\left[\left(_0^3𝒟_0^2𝒜_0^2_0^3\right)+\frac{2}{3}\left(_0^2𝒟_0^3𝒜_0^3_0^2\right)\right].\hfill & \end{array}$$
(168)
The $`𝒜_0^l`$, $`_0^l`$, $`𝒞_0^l`$, $`𝒟_0^l`$, $`_0^l`$ and $`_0^l`$ with $`l=\{1,2,3\}`$ are given by:
$$\{\begin{array}{cc}𝒜_0^1=\frac{1}{\mu _0}\left[2\frac{\chi _0^1}{r^2}\right]\hfill & \\ _0^1=\frac{1}{\mu _0}\left[\frac{1}{r}_r\chi _0^1\right]\hfill & \\ 𝒞_0^1=\frac{1}{\mu _0}\left[\mathrm{\Delta }_1\left(\frac{\xi _0^1}{r}\right)\right]\hfill & \\ 𝒟_0^1=\left[2\frac{\xi _0^1}{r^2}\right]\hfill & \\ _0^1=\left[\frac{1}{r}_r\xi _0^1\right]\hfill & \\ _0^1=\left[\frac{\chi _0^1}{r}\right]\hfill & \end{array}\text{,}\text{ }\{\begin{array}{cc}𝒜_0^2=\frac{1}{\mu _0}\left[6\frac{\chi _0^2}{r^2}\right]\hfill & \\ _0^2=\frac{1}{\mu _0}\left[\frac{1}{r}_r\chi _0^2\right]\hfill & \\ 𝒞_0^2=\frac{1}{\mu _0}\left[\mathrm{\Delta }_2\left(\frac{\xi _0^2}{r}\right)\right]\hfill & \\ 𝒟_0^2=\left[6\frac{\xi _0^2}{r^2}\right]\hfill & \\ _0^2=\left[\frac{1}{r}_r\xi _0^2\right]\hfill & \\ _0^2=\left[\frac{\chi _0^2}{r}\right]\hfill & \end{array}\text{and}\text{ }\{\begin{array}{cc}𝒜_0^3=\frac{1}{\mu _0}\left[12\frac{\chi _0^3}{r^2}\right]\hfill & \\ _0^3=\frac{1}{\mu _0}\left[\frac{1}{r}_r\chi _0^3\right]\hfill & \\ 𝒞_0^3=\frac{1}{\mu _0}\left[\mathrm{\Delta }_3\left(\frac{\xi _0^3}{r}\right)\right]\hfill & \\ 𝒟_0^3=\left[12\frac{\xi _0^3}{r^2}\right]\hfill & \\ _0^3=\left[\frac{1}{r}_r\xi _0^3\right]\hfill & \\ _0^3=\left[\frac{\chi _0^3}{r}\right].\hfill & \end{array}$$
(169)
From (167), the final form of (166) is immediately derived:
$$\rho \frac{\mathrm{d}}{\mathrm{d}t}(r^2\overline{\mathrm{\Omega }})=\frac{1}{5r^2}_r\left(\rho r^4\overline{\mathrm{\Omega }}U_2\right)+\frac{1}{r^2}_r\left(\rho \nu _vr^4_r\overline{\mathrm{\Omega }}\right)+r𝒵_{_{};1}.$$
(170)
#### E.2.2 Mean chemical composition
We restate the mean transport equation for the concentration of chemical species, which remains unchanged by the introduction of the magnetic field (Paper I, eq. B.2):
$$\rho \frac{\mathrm{d}}{\mathrm{d}t}\overline{c_i}+\frac{1}{r^2}_r\left[r^2\rho \overline{c_i}U_i^{\mathrm{diff}}\right]=\frac{1}{r^2}_r\left[r^2\rho (D_v+D_{\mathrm{eff}})_r\overline{c_i}\right].$$
(171)
### E.3 System for $`l=2`$
#### E.3.1 Meridional circulation
We split (77-LABEL:tcal-final) in two first order equations as:
$$U_2=\frac{L}{M\overline{g}}\left(\frac{P}{\overline{\rho }C_p\overline{T}}\right)\frac{1}{_{\text{ad}}}_2$$
(172)
where
$`_2`$ $`=`$ $`2\left[1{\displaystyle \frac{\overline{f}_𝒞+\overline{f}_{}}{4\pi G\overline{\rho }}}{\displaystyle \frac{\left(\overline{ϵ}+\overline{ϵ}_{grav}\right)}{ϵ_m}}\right]{\displaystyle \frac{\stackrel{~}{g}_2}{\overline{g}}}+{\displaystyle \frac{\stackrel{~}{f}_{𝒞,2}+\stackrel{~}{f}_{,2}}{4\pi G\overline{\rho }}}{\displaystyle \frac{\overline{f}_𝒞+\overline{f}_{}}{4\pi G\overline{\rho }}}\left(\delta \mathrm{\Psi }_2+\phi \mathrm{\Lambda }_2\right)+{\displaystyle \frac{\rho _m}{\overline{\rho }}}\left[{\displaystyle \frac{r}{3}}_r𝒜_2{\displaystyle \frac{2H_T}{r}}\left(1+{\displaystyle \frac{D_h}{K}}\right)\mathrm{\Psi }_2\right]`$
$`+`$ $`{\displaystyle \frac{\left(\overline{ϵ}+\overline{ϵ}_{grav}\right)}{ϵ_m}}\left\{𝒜_2+(f_ϵϵ_Tf_ϵ\delta +\delta )\mathrm{\Psi }_2+(f_ϵϵ_\mu +f_ϵ\phi \phi )\mathrm{\Lambda }_2\right\}+{\displaystyle \frac{M}{L}}\left[{\displaystyle \frac{𝒥_2}{\overline{\rho }}}C_p\overline{T}\left({\displaystyle \frac{\mathrm{d}\mathrm{\Psi }_2}{\mathrm{d}t}}+\mathrm{\Phi }{\displaystyle \frac{\mathrm{d}\mathrm{ln}\overline{\mu }}{\mathrm{d}t}}\mathrm{\Lambda }_2\right)\right]`$
and
$$𝒜_2=H_T_r\mathrm{\Psi }_2(1\delta +\chi _T)\mathrm{\Psi }_2(\phi +\chi _\mu )\mathrm{\Lambda }_2.$$
(174)
The coefficients related respectively to the centrifugal force and to the Lorentz force are given by:
$$\overline{f}_𝒞=\frac{1}{r^2}_r\left(r^2a_0\right)\text{ }\text{,}\text{ }\stackrel{~}{f}_{𝒞,2}=\frac{1}{r^2}_r\left(r^2a_2\right)+6\frac{b_2}{r}\text{,}$$
(175)
$$\overline{f}_{}=\frac{1}{r^2}_r\left(r^2\frac{𝒳_{_{};0}}{\rho _0}\right)\text{ and }\stackrel{~}{f}_{,2}=\frac{1}{r^2}_r\left(r^2\frac{𝒳_{_{};2}}{\rho _0}\right)+6\frac{𝒴_{_{;2}}}{r\rho _0}$$
(176)
where
$$\{\begin{array}{cc}a_0=\frac{2}{3}r\overline{\mathrm{\Omega }}^2\hfill & \\ a_2=\frac{2}{3}r\overline{\mathrm{\Omega }}^2+\frac{24}{35}r\overline{\mathrm{\Omega }}\mathrm{\Omega }_2\hfill & \\ b_2=\frac{1}{3}r\overline{\mathrm{\Omega }}^2+\frac{8}{35}r\overline{\mathrm{\Omega }}\mathrm{\Omega }_2\hfill & \end{array}$$
(177)
and
$$𝒳_{_{};0}=\frac{2}{3}\frac{1}{\pi }\left\{\frac{\sqrt{21}}{5}\left[\left(𝒞_0^3_0^1_0^3_0^1\right)+\left(𝒞_0^1_0^3_0^1_0^3\right)\right]+\frac{3}{4}\left(𝒞_0^1_0^1_0^1_0^1\right)+\frac{69}{20}\left(𝒞_0^2_0^2_0^2_0^2\right)+\frac{93}{10}\left(𝒞_0^3_0^3_0^3_0^3\right)\right\}$$
(178)
$`𝒳_{_{};2}`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{1}{2\pi }}\{13\sqrt{{\displaystyle \frac{3}{7}}}[(𝒞_0^3_0^1_0^3_0^1)+(𝒞_0^1_0^3_0^1_0^3)]{\displaystyle \frac{3}{2}}(𝒞_0^1_0^1_0^1_0^1)+{\displaystyle \frac{93}{14}}(𝒞_0^2_0^2_0^2_0^2)`$ (179)
$`+`$ $`{\displaystyle \frac{239}{7}}(𝒞_0^3_0^3_0^3_0^3)\}`$
$$𝒴_{_{};2}=\frac{1}{4\pi }\left\{\sqrt{\frac{3}{7}}\left[4\left(𝒞_0^3𝒟_0^1𝒜_0^1_0^3\right)\left(𝒞_0^1𝒟_0^3𝒜_0^3_0^1\right)\right]+\left(𝒞_0^1𝒟_0^1𝒜_0^1_0^1\right)+\frac{5}{7}\left(𝒞_0^2𝒟_0^2𝒜_0^2_0^2\right)+\frac{2}{3}\left(𝒞_0^3𝒟_0^3𝒜_0^3_0^3\right)\right\}.$$
(180)
The $`𝒜_0^l`$, $`_0^l`$, $`𝒞_0^l`$, $`𝒟_0^l`$, $`_0^l`$ and $`_0^l`$ with $`l=\{1,2,3\}`$ are given in (169).
The relative fluctuation of the effective gravity is given in (70); we apply it here to $`l=2`$:
$$\frac{\stackrel{~}{g}_2}{\overline{g}}=\left[\frac{\mathrm{d}g_0}{\mathrm{d}r}\frac{1}{g_0^2}r\left(b_2+\frac{𝒴_{_{};2}}{\rho _0}\right)+\frac{1}{g_0}\left(a_2+\frac{𝒳_{_{};2}}{\rho _0}\right)\right]+\frac{\mathrm{d}}{\mathrm{d}r}\left(\frac{\widehat{\varphi }_2}{g_0}\right)$$
(181)
where $`\widehat{\varphi }_2`$ is solution of the Poisson equation (cf. eq. 71):
$$\frac{1}{r}\frac{\mathrm{d}^2}{\mathrm{d}r^2}\left(r\widehat{\varphi }_2\right)\frac{6}{r^2}\widehat{\varphi }_2\frac{4\pi G}{g_0}\frac{\mathrm{d}\rho _0}{\mathrm{d}r}\widehat{\varphi }_2=\frac{4\pi G}{g_0}\left[\rho _0a_2+\frac{\mathrm{d}}{\mathrm{d}r}\left(r\rho _0b_2\right)+𝒳_{_{};2}+\frac{\mathrm{d}}{\mathrm{d}r}\left(r𝒴_{_{};2}\right)\right].$$
(182)
The Ohmic heating term is derived from (LABEL:ohmproject):
$`𝒥_2`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}\{[𝒜_0^1D_0^1(_0^1_0^1+𝒞_0^1_0^1)]+[3\sqrt{{\displaystyle \frac{3}{7}}}(𝒜_0^3𝒟_0^1+𝒜_0^1𝒟_0^3)+{\displaystyle \frac{52}{\sqrt{21}}}(_0^3_0^1+𝒞_0^3_0^1+_0^1_0^3+𝒞_0^1_0^3)]`$ (183)
$`+`$ $`{\displaystyle \frac{2}{7}}[5𝒜_0^2𝒟_0^2+31(_0^2_0^2+𝒞_0^2_0^2)]+{\displaystyle \frac{4}{3}}[𝒜_0^3D_0^3+{\displaystyle \frac{239}{7}}(_0^3_0^3+𝒞_0^3_0^3)]\}`$
where the $`𝒜_0^l`$, $`_0^l`$, $`𝒞_0^l`$, $`𝒟_0^l`$, $`_0^l`$ and $`_0^l`$ with $`l=\{1,2,3\}`$ are given by:
$$\{\begin{array}{cc}𝒜_0^1=\frac{\eta _v}{\mu _0}\left[2\frac{\chi _0^1}{r^2}\right]\hfill & \\ _0^1=\frac{\eta _h}{\mu _0}\left[\frac{1}{r}_r\chi _0^1\right]\hfill & \\ 𝒞_0^1=\frac{\eta _h}{\mu _0}\left[\mathrm{\Delta }_1\left(\frac{\xi _0^1}{r}\right)\right]\hfill & \\ 𝒟_0^1=\left[2\frac{\chi _0^1}{r^2}\right]\hfill & \\ _0^1=\left[\frac{1}{r}_r\chi _0^1\right]\hfill & \\ _0^1=\left[\mathrm{\Delta }_1\left(\frac{\xi _0^1}{r}\right)\right]\hfill & \end{array}\text{ }\text{,}\text{ }\{\begin{array}{cc}𝒜_0^2=\frac{\eta _v}{\mu _0}\left[6\frac{\chi _0^2}{r^2}\right]\hfill & \\ _0^2=\frac{\eta _h}{\mu _0}\left[\frac{1}{r}_r\chi _0^2\right]\hfill & \\ 𝒞_0^2=\frac{\eta _h}{\mu _0}\left[\mathrm{\Delta }_2\left(\frac{\xi _0^2}{r}\right)\right]\hfill & \\ 𝒟_0^2=\left[6\frac{\chi _0^2}{r^2}\right]\hfill & \\ _0^2=\left[\frac{1}{r}_r\chi _0^2\right]\hfill & \\ _0^2=\left[\mathrm{\Delta }_2\left(\frac{\xi _0^2}{r}\right)\right]\hfill & \end{array}\text{ }\text{and}\text{ }\{\begin{array}{cc}𝒜_0^3=\frac{\eta _v}{\mu _0}\left[12\frac{\chi _0^3}{r^2}\right]\hfill & \\ _0^3=\frac{\eta _h}{\mu _0}\left[\frac{1}{r}_r\chi _0^3\right]\hfill & \\ 𝒞_0^3=\frac{\eta _h}{\mu _0}\left[\mathrm{\Delta }_3\left(\frac{\xi _0^3}{r}\right)\right]\hfill & \\ 𝒟_0^3=\left[12\frac{\chi _0^3}{r^2}\right]\hfill & \\ _0^3=\left[\frac{1}{r}_r\chi _0^3\right]\hfill & \\ _0^3=\left[\mathrm{\Delta }_3\left(\frac{\xi _0^3}{r}\right)\right].\hfill & \end{array}$$
(184)
#### E.3.2 Baroclinic relation
We apply (64) to $`l=2`$:
$$\phi \mathrm{\Lambda }_2\delta \mathrm{\Psi }_2=\frac{r}{\overline{g}}\left[𝒟_2+\frac{𝒳_{_{};2}}{r\overline{\rho }}+\frac{1}{r}\frac{\mathrm{d}}{\mathrm{d}r}\left(r\frac{𝒴_{_{};2}}{\overline{\rho }}\right)\right]$$
(185)
where $`𝒟_2=\frac{1}{3}\left[r_r\overline{\mathrm{\Omega }}^2\right]+\frac{8}{35}\left[r_r\left(\overline{\mathrm{\Omega }}\mathrm{\Omega }_2\right)\right]+\frac{8}{7}\overline{\mathrm{\Omega }}\mathrm{\Omega }_2`$ and where $`𝒳_{_{};2}`$ and $`𝒴_{_{};2}`$ have been given in (179-180).
#### E.3.3 Horizontal fluctuation of the molecular weight
The evolution equation for the mean molecular weight (Paper I, eq. 39) is unchanged by the introduction of the magnetic field:
$$\frac{\mathrm{d}\mathrm{\Lambda }_2}{\mathrm{d}t}\frac{\mathrm{d}\mathrm{ln}\overline{\mu }}{\mathrm{d}t}\mathrm{\Lambda }_2=\frac{U_2}{H_p}_\mu \frac{6}{r^2}D_h\mathrm{\Lambda }_2.$$
(186)
### E.4 System for $`l=4`$
#### E.4.1 Horizontal shear
We have:
$$\rho \frac{\mathrm{d}}{\mathrm{d}t}\left(r^2\mathrm{\Omega }_2\right)2\rho \overline{\mathrm{\Omega }}r\left[\frac{1}{3\rho r}_r\left(\rho r^2U_2\right)\alpha U_2\right]=10\rho \nu _h\mathrm{\Omega }_2+\mathrm{\Gamma }_2$$
(187)
where the expression for $`\mathrm{\Gamma }_2`$ has been given in (167).
#### E.4.2 Meridional circulation
In the same way as for $`l=2`$, we recast (77-LABEL:tcal-final) in two first order equations as
$$U_4=\frac{L}{M\overline{g}}\left(\frac{P}{\overline{\rho }C_p\overline{T}}\right)\frac{1}{_{\text{ad}}}_4$$
(188)
where
$`𝒯_4`$ $`=`$ $`2\left[1{\displaystyle \frac{\overline{f}_𝒞+\overline{f}_{}}{4\pi G\overline{\rho }}}{\displaystyle \frac{\left(\overline{ϵ}+\overline{ϵ}_{grav}\right)}{ϵ_m}}\right]{\displaystyle \frac{\stackrel{~}{g}_4}{\overline{g}}}+{\displaystyle \frac{\stackrel{~}{f}_{𝒞,4}+\stackrel{~}{f}_{,4}}{4\pi G\overline{\rho }}}{\displaystyle \frac{\overline{f}_𝒞+\overline{f}_{}}{4\pi G\overline{\rho }}}\left(\delta \mathrm{\Psi }_4+\phi \mathrm{\Lambda }_4\right)+{\displaystyle \frac{\rho _m}{\overline{\rho }}}\left[{\displaystyle \frac{r}{3}}_r𝒜_4{\displaystyle \frac{20H_T}{3r}}\left(1+{\displaystyle \frac{D_h}{K}}\right)\mathrm{\Psi }_4\right]`$
$`+`$ $`{\displaystyle \frac{\left(\overline{ϵ}+\overline{ϵ}_{grav}\right)}{ϵ_m}}\left\{𝒜_4+(f_ϵϵ_Tf_ϵ\delta +\delta )\mathrm{\Psi }_4+(f_ϵϵ_\mu +f_ϵ\phi \phi )\mathrm{\Lambda }_4\right\}+{\displaystyle \frac{M}{L}}\left[{\displaystyle \frac{𝒥_4}{\overline{\rho }}}C_p\overline{T}\left({\displaystyle \frac{\mathrm{d}\mathrm{\Psi }_4}{\mathrm{d}t}}+\mathrm{\Phi }{\displaystyle \frac{\mathrm{d}\mathrm{ln}\overline{\mu }}{\mathrm{d}t}}\mathrm{\Lambda }_4\right)\right]`$
and
$$𝒜_4=H_T_r\mathrm{\Psi }_4(1\delta +\chi _T)\mathrm{\Psi }_4(\phi +\chi _\mu )\mathrm{\Lambda }_4.$$
(189)
The coefficients related respectively to the centrifugal force and to the Lorentz force are given by:
$$\stackrel{~}{f}_{𝒞,4}=\frac{1}{r^2}_r\left(r^2a_4\right)+20\frac{b_4}{r}\text{ }\text{and}\text{ }\stackrel{~}{f}_{,4}=\frac{1}{r^2}_r\left(r^2\frac{𝒳_{_{};4}}{\rho _0}\right)+20\frac{𝒴_{_{};4}}{r\rho _0}$$
(190)
where:
$$\{\begin{array}{cc}a_4=\frac{24}{35}r\overline{\mathrm{\Omega }}\mathrm{\Omega }_2\hfill & \\ b_4=\frac{6}{35}r\overline{\mathrm{\Omega }}\mathrm{\Omega }_2\hfill & \end{array}$$
(191)
and
$$\{\begin{array}{cc}𝒳_{_{};4}=\frac{2}{3}\frac{1}{5\pi }\left\{\frac{9}{2}\sqrt{\frac{3}{7}}\left[\left(𝒞_0^1_0^3_0^1_0^3\right)\left(𝒞_0^3_0^1_0^3_0^1\right)\right]+\frac{27}{7}\left(𝒞_0^2_0^2_0^2_0^2\right)+\frac{14449}{154}\left(𝒞_0^3_0^3_0^3_0^3\right)\right\}\hfill & \\ 𝒴_{_{};4}=\frac{1}{4\pi }\left\{\sqrt{\frac{3}{7}}\left(3𝒞_0^3𝒟_0^1𝒜_0^3_0^1+𝒞_0^1𝒟_0^33𝒜_0^1_0^3\right)+\frac{9}{7}\left(𝒞_0^2𝒟_0^2𝒜_0^2_0^2\right)+\frac{9}{11}\left(𝒞_0^3𝒟_0^3𝒜_0^3_0^3\right)\right\}.\hfill & \end{array}$$
(192)
We get the relative fluctuation of the effective gravity from (70):
$$\frac{\stackrel{~}{g}_4}{\overline{g}}=\left[\frac{\mathrm{d}g_0}{\mathrm{d}r}\frac{1}{g_0^2}r\left(b_4+\frac{𝒴_{_{};4}}{\rho _0}\right)+\frac{1}{g_0}\left(a_4+\frac{𝒳_{_{};4}}{\rho _0}\right)\right]+\frac{\mathrm{d}}{\mathrm{d}r}\left(\frac{\widehat{\varphi }_4}{g_0}\right)$$
(193)
where $`\widehat{\varphi }_4`$ obeys the Poisson equation:
$$\frac{1}{r}\frac{\mathrm{d}^2}{\mathrm{d}r^2}\left(r\widehat{\varphi }_4\right)\frac{20}{r^2}\widehat{\varphi }_4\frac{4\pi G}{g_0}\frac{\mathrm{d}\rho _0}{\mathrm{d}r}\widehat{\varphi }_4=\frac{4\pi G}{g_0}\left[\rho _0a_4+\frac{\mathrm{d}}{\mathrm{d}r}\left(r\rho _0b_4\right)+𝒳_{_{};4}+\frac{\mathrm{d}}{\mathrm{d}r}\left(r𝒴_{_{};4}\right)\right]\text{ }\text{ }\text{ }\text{(cf. eq. }\text{71}\text{)}.$$
(194)
We proceed as for $`l=2`$ to calculate the Ohmic heating:
$`𝒥_4`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}\{4\sqrt{{\displaystyle \frac{3}{7}}}[𝒜_0^3𝒟_0^1+𝒜_0^1𝒟_0^3{\displaystyle \frac{3}{5}}(_0^3_0^1+𝒞_0^3_0^1+_0^1_0^3+𝒞_0^1_0^3)]+{\displaystyle \frac{18}{7}}[𝒜_0^2𝒟_0^2{\displaystyle \frac{4}{5}}(_0^2_0^2+𝒞_0^2_0^2)]`$ (195)
$`+`$ $`{\displaystyle \frac{2}{11}}[9𝒜_0^3𝒟_0^3+{\displaystyle \frac{28898}{105}}(_0^3_0^3+𝒞_0^3_0^3)]\}`$
where the $`𝒜_0^l`$, $`_0^l`$, $`𝒞_0^l`$, $`𝒟_0^l`$, $`_0^l`$ and $`_0^l`$ with $`l=\{1,2,3\}`$ have been given in the previous section in (184).
#### E.4.3 Baroclinic relation
We apply (64) to $`l=4`$. We get:
$$\phi \mathrm{\Lambda }_4\delta \mathrm{\Psi }_4=\frac{r}{\overline{g}}\left[𝒟_4+\frac{𝒳_{_{};4}}{r\overline{\rho }}+\frac{1}{r}\frac{\mathrm{d}}{\mathrm{d}r}\left(r\frac{𝒴_{_{};4}}{\overline{\rho }}\right)\right]$$
(196)
where $`𝒟_4=\frac{6}{35}\left[r_r\left(\overline{\mathrm{\Omega }}\mathrm{\Omega }_2\right)2\overline{\mathrm{\Omega }}\mathrm{\Omega }_2\right]`$ and where $`𝒳_{_{};4}`$ and $`𝒴_{_{};4}`$ have been given in (192).
#### E.4.4 Horizontal fluctuation of the molecular weight
This equation is not affected by the introduction of the magnetic field:
$$\frac{\mathrm{d}\mathrm{\Lambda }_4}{\mathrm{d}t}\frac{\mathrm{d}\mathrm{ln}\overline{\mu }}{\mathrm{d}t}\mathrm{\Lambda }_4=\frac{U_4}{H_p}_\mu \frac{20}{r^2}D_h\mathrm{\Lambda }_4.$$
(197)
These equations are ready to be implemented in a stellar evolution code, together with the boundary conditions discussed in §7.
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# Asymptotic behavior of the ghost propagator in SU3 lattice gauge theory
## 1 Introduction
Whereas lattice gauge theory (LGT) has been initially formulated to study gauge-invariant quantities in the non-perturbative regime, it has long been recognized that LGT could be a useful tool for studying gauge-variant quantities such as Green functions, both in the non-perturbative and in the perturbative regimes. The SU(3) gluon propagator in momentum space was first considered to gain some insight into the physics of confinement. Much work was then devoted to the study of its infrared behavior (for a review see ). Subsequent studies were focused on the ultraviolet behavior and have been able to compare quantitatively the large momentum dependence of the lattice gluon propagator with perturbative predictions beyond one-loop order. The result for $`\mathrm{\Lambda }_{\overline{\text{MS}}}`$ was found to depend strongly upon the order of the perturbation theory and upon the renormalisation scheme used in the parametrization. This strong dependence raised the question whether the energy windows in these calculations were large enough for perturbative QCD to be a valid approximation.
On the other hand, as shown by Gribov , the infrared behavior of the gluon propagator is closely related to the singularity structure of the ghost propagator inferred from the gauge-fixing ambiguities. As is well-known, the Landau gauge, which is presently the only covariant gauge for which there exists effective local gauge-fixing algorithms on the lattice, suffers from these ambiguities. The comprehensive theoretical study by Zwanziger of the Faddeev-Popov operator on the lattice in Landau gauge spurred the first numerical study of the ghost propagator in SU(2) and SU(3) gauge theories. Most subsequent activity has been dedicated to the SU(2) lattice gauge theory in the infrared region, mainly for technical reasons as we shall explain below. There are relatively few numerical studies of the SU(3) ghost propagator which are either more focused on the infrared region and the Gribov copy problem or have only performed a qualitative perturbative description in the quenched approximation and, quite recently, in the unquenched case also .
It is important to make the study of the SU(3) ghost propagator in the ultraviolet region more quantitative for comparison purposes with the gluon propagator. Lattice results at small distances may be described by perturbation theory and the independent extraction of the $`\mathrm{\Lambda }_{\text{QCD}}`$ scale from the two propagators would provide a self-consistency test of the analysis and of the lattice approach. It would be particularly significant to confirm or not, from the study of the lattice propagators alone, the need for the non-perturbative power corrections found in the study of the three-gluon coupling on the lattice .
The paper is organized as follows. We will begin by recalling in section 2 the method used to relate lattice data for the ghost propagator to its perturbative renormalization description. Then we proceed by exhibiting in section 3 the salient features of our lattice calculation, particularly of our implementation of the Faddeev-Popov operator on the lattice. The following section outlines the general method that we devised previously to eliminate hypercubic artifacts from two-point functions and extrapolate the lattice data towards the continuum. This extrapolation is crucial to succeed in a quantitative description. The results are discussed in section 5 which contains several subsections where the analysis is performed in different renormalization schemes up to four-loop order. In particular the scheme dependence is thoroughly investigated and used to probe the effects of the truncation of the perturbative series. We conclude in section 6 with a comparison of the different methods to compute the $`\mathrm{\Lambda }_{\text{QCD}}`$ scale on the lattice.
## 2 Renormalization description of the ghost propagator
Let $`\mathrm{\Gamma }_B^{(n)}`$ be some gauge-fixed multiplicatively renormalizable one-particle irreducible $`n`$-point bare Green function defined in euclidean momentum space and in some regularization scheme with cut-off $`\mathrm{\Lambda }`$. Let $`s`$ denotes some polarization state and kinematical configuration of the external particles contributing to $`\mathrm{\Gamma }_B^{(n)}`$. Let $`p`$ denote a scale transformation on $`s`$ and $`g_B`$ denote the bare coupling. It is well known that, in any renormalization scheme $`R`$ defined by some renormalization conditions on state $`s`$ at the renormalization point $`p=\mu `$, we have
$`\mathrm{\Gamma }_B^{(n)}(p,s,g_B,\mathrm{\Lambda })=Z_{\mathrm{\Gamma },R}(\mu ,s,g_R,\mathrm{\Lambda })\mathrm{\Gamma }_R^{(n)}(p,s,g_R,\mu )+𝒪(\mathrm{\Lambda }^1)`$ (1)
where $`Z_{\mathrm{\Gamma },R}`$ is the renormalization constant in scheme $`R`$, $`\mathrm{\Gamma }_R^{(n)}`$ is the renormalized Green function and $`g_R(\mu )`$ is the renormalized coupling. We omit the dependence on the gauge parameter for simplicity of notation since we will specialize to Landau gauge.
The explicit dependence on $`\mu `$ drops out of the renormalized Green function $`\mathrm{\Gamma }_R^{(n)}`$ at the renormalization point $`p=\mu `$. It follows that
$`\begin{array}{cc}\hfill \underset{\mathrm{\Lambda }\mathrm{}}{lim}{\displaystyle \frac{d\mathrm{ln}\left(\mathrm{\Gamma }_B^{(n)}(\mu ,s,g_B,\mathrm{\Lambda })\right)}{d\mathrm{ln}\mu ^2}}& =\underset{\mathrm{\Lambda }\mathrm{}}{lim}{\displaystyle \frac{d\mathrm{ln}\left(Z_{\mathrm{\Gamma },R}(\mu ,s,g_R,\mathrm{\Lambda })\right)}{d\mathrm{ln}\mu ^2}}+{\displaystyle \frac{d\mathrm{ln}\left(\mathrm{\Gamma }_R^n(s,g_R)\right)}{d\mathrm{ln}\mu ^2}}\hfill \\ & \gamma _{\mathrm{\Gamma },R}(g_R)+{\displaystyle \frac{dg_R}{d\mathrm{ln}\mu ^2}}{\displaystyle \frac{\mathrm{ln}\mathrm{\Gamma }_R^n}{g_R}}\hfill \end{array}`$ (2)
The arbitrariness in the choice of the renormalization scheme $`R`$ has prompted attempts at determining the “best” schemes for describing the $`q^2`$-evolution of bare Green functions on the lattice. Clearly it is always possible to find a change of coupling which will be a best approximation of a set of data at a given order of perturbation theory, within some prescribed criteria. Rather than pursuing this route, we will follow the standard wisdom which consists in choosing renormalization conditions appropriate to the continuum quantity under scrutiny.
Momentum substraction schemes have long been used to define renormalization conditions befitted to the description of the renormalization dependence of “physical” quantities. They are defined by setting some of the 2 and 3-point functions to their tree values. In the $`\stackrel{~}{\text{MOM}}`$ schemes, for these Green functions, Eq. (2) simplifies to
$`\underset{\mathrm{\Lambda }\mathrm{}}{lim}{\displaystyle \frac{d\mathrm{ln}\left(\mathrm{\Gamma }_B^{(n)}(\mu ,s,g_B,\mathrm{\Lambda })\right)}{d\mathrm{ln}\mu ^2}}`$ $`={\displaystyle \frac{d\mathrm{ln}(Z_{\mathrm{\Gamma },MOM})}{d\mathrm{ln}\mu ^2}}=\gamma _{\mathrm{\Gamma },MOM}(g_{MOM})`$ (3)
Infinitely many MOM schemes can be defined which differ by the substraction point of the vertices. We have shown in that the $`\stackrel{~}{\text{MOMg}}`$ scheme defined by substracting the transversal part of the three-gluon vertex at the asymmetric point where one external momentum vanishes, appears to provide a much better estimate of the asymptotic behavior of the gluon propagator in Landau gauge than the $`\overline{\text{MS}}`$ scheme. For the study of the asymptotic behavior of the ghost propagator in Landau gauge, it seems therefore natural to use a $`\stackrel{~}{\text{MOMc}}`$ scheme defined by substracting the ghost-gluon vertex at the asymmetric point where the momentum of the external gluon vanishes. Comparison of the two $`\stackrel{~}{\text{MOM}}`$ schemes should provide us with an estimate of the systematic error entailed in the truncation of the perturbation theory.
The perturbative calculation of the gluon, ghost and quark self-energies and all 3-vertices appearing in the QCD Lagrangian have been done at three-loop order in the $`\overline{\text{MS}}`$ scheme and in a general covariant gauge at the asymmetric point with one vanishing momentum . These three-loop results allow one to relate the coupling constants of any $`\stackrel{~}{\text{MOM}}`$-like scheme to the $`\overline{\text{MS}}`$ scheme at three-loop order. For the $`\stackrel{~}{\text{MOMg}}`$ and $`\stackrel{~}{\text{MOMc}}`$ schemes defined above these relations read respectively in Landau gauge and in the quenched approximation ($`n_f=0`$), with $`h={\displaystyle \frac{g^2}{16\pi ^2}}`$:
$`\begin{array}{cc}\hfill h_{\stackrel{~}{\text{MOMg}}}& =h_{\overline{\text{MS}}}+{\displaystyle \frac{70}{3}}h_{\overline{\text{MS}}}^2+\left({\displaystyle \frac{516217}{576}}{\displaystyle \frac{153}{4}}\zeta _3\right)h_{\overline{\text{MS}}}^3\hfill \\ & +\left({\displaystyle \frac{304676635}{6912}}{\displaystyle \frac{299961}{64}}\zeta _3{\displaystyle \frac{81825}{64}}\zeta _5\right)h_{\overline{\text{MS}}}^4\hfill \end{array}`$ (4)
$`\begin{array}{cc}\hfill h_{\stackrel{~}{\text{MOMc}}}& =h_{\overline{\text{MS}}}+{\displaystyle \frac{223}{12}}h_{\overline{\text{MS}}}^2+\left({\displaystyle \frac{918819}{1296}}{\displaystyle \frac{351}{8}}\zeta _3\right)h_{\overline{\text{MS}}}^3\hfill \\ & +\left({\displaystyle \frac{29551181}{864}}{\displaystyle \frac{137199}{32}}\zeta _3{\displaystyle \frac{74295}{64}}\zeta _5\right)h_{\overline{\text{MS}}}^4\hfill \end{array}`$ (5)
The very large coefficients of these perturbative expansions explain the difficulties met by the $`\overline{\text{MS}}`$ scheme to approach asymptotic scaling below 10 GeV.
The recent calculation of the anomalous dimensions in the $`\overline{\text{MS}}`$ scheme of the gluon and ghost fields at four-loop order, together with the knowledge of the $`\beta `$-function , makes it possible to perform the analysis of the lattice data for the gluon and ghost propagators up to four-loop order also in the $`\stackrel{~}{\text{MOMg}}`$ and $`\stackrel{~}{\text{MOMc}}`$ schemes. The numerical coefficients of the $`\beta `$-function defined as
$`\beta (h)={\displaystyle \frac{dh}{d\mathrm{ln}\mu ^2}}={\displaystyle \underset{i=0}{\overset{n}{}}}\beta _ih^{i+2}+𝒪(h^{n+3})`$ (6)
are:
$`\beta _2^{\stackrel{~}{\text{MOMg}}}=2412.16,\beta _2^{\stackrel{~}{\text{MOMc}}}=2952.73,\beta _3^{\stackrel{~}{\text{MOMg}}}=84353.8,\beta _3^{\stackrel{~}{\text{MOMc}}}=101484.`$ (7)
For completeness we also give the expansion coefficients of the renormalisation constants of the gluon and ghost fields in the MOM schemes with respect to the renormalized coupling of the $`\overline{\text{MS}}`$ scheme up to four-loop order:
$`\begin{array}{cc}\hfill {\displaystyle \frac{d\mathrm{ln}(Z_{3,MOM})}{d\mathrm{ln}\mu ^2}}& ={\displaystyle \frac{13}{2}}h_{\overline{\text{MS}}}+{\displaystyle \frac{3727}{24}}h_{\overline{\text{MS}}}^2+\left({\displaystyle \frac{2127823}{288}}{\displaystyle \frac{9747}{16}}\zeta _3\right)h_{\overline{\text{MS}}}^3\hfill \\ & +\left({\displaystyle \frac{3011547563}{6912}}{\displaystyle \frac{18987543}{256}}\zeta _3{\displaystyle \frac{1431945}{64}}\zeta _5\right)h_{\overline{\text{MS}}}^4\hfill \end{array}`$ (8)
$`\begin{array}{cc}\hfill {\displaystyle \frac{d\mathrm{ln}(\stackrel{~}{Z}_{3,MOM})}{d\mathrm{ln}\mu ^2}}& ={\displaystyle \frac{9}{4}}h_{\overline{\text{MS}}}+{\displaystyle \frac{813}{16}}h_{\overline{\text{MS}}}^2+\left({\displaystyle \frac{157303}{64}}{\displaystyle \frac{5697}{32}}\zeta _3\right)h_{\overline{\text{MS}}}^3\hfill \\ & +\left({\displaystyle \frac{219384137}{1536}}{\displaystyle \frac{9207729}{512}}\zeta _3{\displaystyle \frac{221535}{32}}\zeta _5\right)h_{\overline{\text{MS}}}^4\hfill \end{array}`$ (9)
## 3 Lattice calculation
### 3.1 Faddeev-Popov operator on the lattice
The ghost propagator is defined on the lattice as
$`G(xy)\delta ^{ab}\left(M^1\right)_{xy}^{ab}`$ (10)
where the action of the Faddeev-Popov operator $`M`$ on an arbitrary element $`\eta `$ of the Lie algebra $`𝒮𝒰`$(N) of the gauge group SU(N), in a Landau gauge fixed configuration, is given by :
$`(M\eta )^a(x)`$ $`={\displaystyle \frac{1}{N}}{\displaystyle \underset{\mu }{}}\{G_\mu ^{ab}(x)(\eta ^b(x+\widehat{\mu })\eta ^b(x))(xx\widehat{\mu })`$
$`+{\displaystyle \frac{1}{2}}f^{abc}(\eta ^b(x+\widehat{\mu })A_\mu ^c(x)\eta ^b(x\widehat{\mu })A_\mu ^c(x\widehat{\mu }))\}`$ (11)
and where, with antihermitian generators $`T^a`$,
$`G_\mu ^{ab}(x)`$ $`={\displaystyle \frac{1}{2}}\mathrm{𝖳𝗋}\left(\{T^a,T^b\}\left(U_\mu (x)+U_\mu ^{}(x)\right)\right)`$ (12)
$`A_\mu ^c(x)`$ $`=\mathrm{𝖳𝗋}\left(T^c\left(U_\mu (x)U_\mu ^{}(x)\right)\right)`$ (13)
Most lattice implementations of the Faddeev-Popov operator have followed closely the component-wise Eqs. (3.1-13). But the derivation in shows that the Faddev-Popov operator can also be written as a lattice divergence:
$`M(U)={\displaystyle \frac{1}{N}}\stackrel{~}{D}(U)`$ (14)
where the operator $`\stackrel{~}{D}`$ reads
$`\stackrel{~}{D}_\mu (U)\eta (x)`$ $`={\displaystyle \frac{1}{2}}\left(U_\mu (x)\eta (x+\widehat{\mu })\eta (x)U_\mu (x)+\eta (x+\widehat{\mu })U_\mu ^{}(x)U_\mu ^{}(x)\eta (x)\right)`$ (15)
Using conversion routines between the Lie algebra and the Lie group, eqs. (14-15) allow for a very efficient lattice implementation, sketched in Table 1, which is based on the fast routines coding the group multiplication law.
### 3.2 Inversion of the Faddeev-Popov operator
Constant fields are zero modes of the Faddeev-Popov operator. This operator can be inverted only in the vector subspace $`K^{}`$ orthogonal to its kernel. If the Faddeev-Popov operator has no other zero modes than constant fields, then the non-zero Fourier modes form a basis of $`K^{}`$:
$`\eta (x)={\displaystyle \underset{p0}{}}c_pe^{ipx},\eta K^{}`$ (16)
The standard procedure has been to invert the Faddev-Popov operator with one non-zero Fourier mode as a source
$`S_p^a(x)=\delta ^{ab}e^{ipx}`$ (17)
and to take the scalar product of $`M^1S_p^a`$ with the source:
$`\left(S_p^a|M^1S_p^a\right)`$ $`={\displaystyle \underset{x,y}{}}\left(M^1\right)_{xy}^{aa}e^{ip(xy)}`$ (18)
$`=V\widehat{G}(p)`$ (19)
after averaging over the gauge field configurations. This method requires one matrix inversion for each value of the ghost propagator in momentum space. It is suitable only when one is interested in a few values of the ghost propagator.
However, the study of the ultraviolet behavior of the ghost propagator in the continuum requires its calculation at many lattice momenta to control the spacing artifacts, as we shall see in the next section. This can be done very economically by noting that
$`\delta (x,y)={\displaystyle \frac{1}{V}}+{\displaystyle \frac{1}{V}}{\displaystyle \underset{p0}{}}e^{ip(xy)}`$ (20)
and choosing as source:
$`S_0^a(x)=\delta ^{ab}\left(\delta (x,0){\displaystyle \frac{1}{V}}\right)`$ (21)
The Fourier transform of $`M^1S_0^a`$, averaged over the gauge configurations, yields:
$`{\displaystyle \underset{x}{}}e^{ipx}M^1S_0^a`$ $`={\displaystyle \underset{x}{}}e^{ipx}\left(M^1\right)_{x0}^{aa}{\displaystyle \frac{1}{V}}{\displaystyle \underset{x,y}{}}e^{ipx}\left(M^1\right)_{xy}^{aa}`$
$`={\displaystyle \underset{x}{}}e^{ipx}G(x){\displaystyle \frac{1}{V}}{\displaystyle \underset{x,y}{}}e^{ipx}G(xy)`$
$`=\widehat{G}(p)\delta (p){\displaystyle \underset{x}{}}G(x)`$ (22)
as a consequence of the translation invariance of the ghost propagator. Therefore, with this choice of source, only one matrix inversion followed by one Fourier transformation of the solution is required to get the full ghost propagator on the lattice.
There is of course a price to pay, as can be read off Eq. (22) which lacks the factor $`V`$ present in Eq. (19). The statistical accuracy with the source $`S_p^a`$ is better, especially at high momentum $`p`$. However the statistical accuracy with the source $`S_0^a`$ turns out to be sufficient for our purpose.
There is one final point we want to make and which has never beeen raised to the best of our knowledge. It is mandatory to check, whatever the choice of sources, that rounding errors during the inversion do not destroy the condition that the solution belongs to $`K^{}`$:
$`{\displaystyle \underset{x}{}}\left(M^1S\right)(x)=0`$ (23)
Indeed, if the zero-mode component of the solution grows beyond some threshold during the inversion of the Faddeev-Popov operator on some gauge configuration, then that component starts to increase exponentially and a sizeable bias is produced in other components as well. We have observed this phenomenon occasionally, about one gauge configuration every few hundreds, when using the implementation of the lattice Faddeev-Popov operator based on Eqs. (3.1-13). But the systematic bias which is induced on the averages over gauge field configurations can be uncomfortably close to those ascribed to Gribov copies.
Another virtue of the algorithm described in Table 1 is its numerical stability which is improved by several orders of magnitude. We have never observed sizeable deviations from Eq. (23) with this algorithm.
### 3.3 The simulation
We ran simulations of the $`SU(3)`$ lattice gauge theory with the Wilson action in the quenched approximation on several hypercubic lattices, whose parameters are summarized in Table 2. All lattices have roughly the same physical volume except the $`24^4`$ lattice at $`\beta =6.0`$ which has been included to check out finite-volume effects.
The SU(3) gauge configurations were generated using a hybrid algorithm of Cabibbo-Marinari heatbath and Creutz overrelaxation steps. 10000 lattice updates were discarded for thermalization and the configurations were analyzed every 100/200/500 sweeps on the $`16^4/24^4/32^4`$ lattices.
Landau gauge fixing was carried out by minimizing the functional
$`F_U[g]=\text{Re}{\displaystyle \underset{x}{}}{\displaystyle \underset{\mu }{}}\left(1{\displaystyle \frac{1}{N}}g(x)U_\mu (x)g^{}(x+\widehat{\mu })\right)`$ (24)
by use of a standard overrelaxation algorithm driving the gauge configuration to a local minimum of $`F_U[g]`$. We did not try to reach the fundamental modular region $`\mathrm{\Lambda }`$, defined as the set of absolute minima of $`F_U[g]`$ on all gauge orbits. Indeed there have been numerous studies, in SU(2) and in SU(3) , of the effect of Gribov copies on the ghost propagator. The consensus is that noticeable systematic errors, beyond statistical errors, are only found for the smallest $`p^2`$, much smaller than the squared momenta that we used to study the asymptotic behavior of the ghost propagator.
Then the ghost propagator $`G(p)`$ is extracted from Eq. (22) for all $`p0`$. The required matrix inversion, with a conjugate-gradient algorithm without any preconditioning, and the Fourier transform consume in average less than half the computing time of the Landau gauge fixing.
## 4 Hypercubic artifacts
The ghost propagator $`\widehat{G}(p)`$ is a scalar invariant on the lattice which means that it is invariant along the orbit $`O(p)`$ generated by the action of the isometry group $`H(4)`$ of hypercubic lattices on the discrete momentum $`p\frac{2\pi }{La}\times (n_1,n_2,n_3,n_4)`$ where the $`n_\mu `$’s are integers, $`L`$ is the lattice size and $`a`$ the lattice spacing. The general structure of polynomials invariant under a finite group is known from group-invariant theory. Indeed it can be shown that any polynomial function of $`p`$ which is invariant under the action of $`H(4)`$ is a polynomial function of the 4 invariants $`p^{[n]}=a^n_\mu p_\mu ^n,n=2,4,6,8`$ which index the set of orbits.
Our analysis program uses these 4 invariants to average the ghost propagator over the orbits of $`H(4)`$ to increase the statistical accuracy:
$`a^2G_L(p^{[2]},p^{[4]},p^{[6]},p^{[8]})={\displaystyle \frac{1}{O(p)}}{\displaystyle \underset{pO(p)}{}}\widehat{G}(p)`$ (25)
where $`O(p)`$ is the cardinal number of the orbit $`O(p)`$. By the same token, one should always take the following real source
$`\overline{S}_p^a(x)=\delta ^{ab}{\displaystyle \underset{pO(p)}{}}\mathrm{cos}(px)`$ (26)
rather than a single complex Fourier mode for studies of the ghost propagator in the infrared region. Indeed, after averaging over the gauge configurations and use of the translational invariance, one gets
$`\left(\overline{S}_p^a|M^1\overline{S}_p^a\right)`$ $`={\displaystyle \underset{p,p^{}O(p)}{}}{\displaystyle \underset{x,y}{}}\left(M^1\right)_{xy}^{aa}e^{ip^{}x+ipy}`$
$`=VO(p)a^2G_L(p^{[2]},p^{[4]},p^{[6]},p^{[8]})`$ (27)
By analogy with the free lattice propagator
$`G_0(p)={\displaystyle \frac{1}{_\mu \widehat{p}_\mu ^2}}={\displaystyle \frac{a^2}{p^{[2]}}}\left(1+{\displaystyle \frac{1}{12}}{\displaystyle \frac{p^{[4]}}{p^{[2]}}}+\mathrm{}\right),\mathrm{where}\widehat{p}_\mu ={\displaystyle \frac{2}{a}}\mathrm{sin}\left({\displaystyle \frac{ap_\mu }{2}}\right)`$ (28)
it is natural to make the hypothesis that the lattice ghost propagator is a smooth function of the discrete invariants near the continuum limit, when $`ap_\mu 1,\mu `$,
$`G_L(p^{[2]},p^{[4]},p^{[6]},p^{[8]})G_L(p^{[2]},0,0,0)+p^{[4]}{\displaystyle \frac{G_L}{p^{[4]}}}(p^{[2]},0,0,0)+\mathrm{}`$ (29)
and $`G_L(p^{[2]},0,0,0)`$ is nothing but the propagator of the continuuum in a finite volume, up to lattice artifacts which do not break $`O(4)`$ invariance. When several orbits exist with the same $`p^2`$, the simplest method to reduce the hypercubic artifacts is to extrapolate the lattice data towards $`G_L(p^{[2]},0,0,0)`$ by making a linear regression at fixed $`p^2`$ with respect to the invariant $`p^{[4]}`$ since the other invariants are of higher order in the lattice spacing. The range of validity of this linear approximation can be checked a posteriori from the smoothness of the extrapolated data with respect to $`p^2`$.
Choosing the variables $`\widehat{p}_\mu `$ appropriate to the parametrization of a lattice propagator with periodic boundary conditions provides an independent check of the extrapolation. Indeed we can write as well
$`G_L(p^{[2]},p^{[4]},p^{[6]},p^{[8]})\widehat{G}_L(\widehat{p}^{[2]},\widehat{p}^{[4]},\widehat{p}^{[6]},\widehat{p}^{[8]})`$ (30)
with the new invariants, again hierachically suppressed with respect to the lattice spacing,
$`\widehat{p}^{[n]}=a^n{\displaystyle \underset{\mu }{}}\widehat{p}_\mu ^n`$ (31)
$`G_L`$ and $`\widehat{G}_L`$ are two different parametrizations of the same lattice data, but near the continuum limit one must also have
$`\widehat{G}_L(\widehat{p}^{[2]},\widehat{p}^{[4]},\widehat{p}^{[6]},\widehat{p}^{[8]})\widehat{G}_L(\widehat{p}^{[2]},0,0,0)+\widehat{p}^{[4]}{\displaystyle \frac{\widehat{G}_L}{\widehat{p}^{[4]}}}(\widehat{p}^{[2]},0,0,0)+\mathrm{}`$ (32)
where $`G_L(p^{[2]},0,0,0)`$ and $`\widehat{G}_L(\widehat{p}^{[2]},0,0,0)`$ are the same function, the propagator of the continuum , of a different variable (again up to lattice artifacts which do not break $`O(4)`$ invariance).
If one wants to include in the data analysis the points with a single orbit at fixed $`p^2`$, one must interpolate the slopes extracted from Eqs (29) or (32). This interpolation can be done either numerically or by assuming a functional dependence of the slope with respect to $`p^2`$ based on dimensional arguments. The simplest ansatz is to assume that the slope has the same leading behavior as for a free lattice propagator:
$`{\displaystyle \frac{G_L}{p^{[4]}}}(p^{[2]},0,0,0)`$ $`={\displaystyle \frac{1}{\left(p^{[2]}\right)^2}}\left(c_1+c_2p^{[2]}\right)`$ (33)
The inclusion of $`O(4)`$-invariant lattice spacing corrections is required to get fits with a reasonable $`\chi ^2`$. The quality of such two-parameter fits to the slopes, and the extension of the fitting window in $`p^2`$, supplies still another independent check of the validity of the extrapolations.
We have used Eqs. (29) and (33) to extrapolate our lattice data towards the continuum and determined the range of validity in $`p^2`$ of the extrapolations from the consistency of the different checks within our statistical errors. The errors on the extrapolated points have been computed with the jackknife method. Tables 3 and 4 summarize the cuts that have been applied to the data for the estimation of the systematic errors in the analysis of the next section. We have repeated the analysis of the gluon propagator to study the sensitivity of the results with respect to the window in $`p^2`$ which has been enlarged considerably in our new data. The cuts for the lattice ghost propagator are stronger than for the gluon lattice propagator because the statistical errors of the former are two to three times larger which make the continuum extrapolations less controllable.
The number of distinct orbits at each $`p^2`$ increases with the lattice size and, eventually, a linear extrapolation limited to the single invariant $`p^{[4]}`$ breaks down. However there is a systematic way to include higher-order invariants and to extend the range of validity of the extrapolations. A much more detailed exposition of the controlling of systematic errors is in preparation, since our method has been largely ignored in the litterature where very empirical recipes are still in use.
## 5 Data analysis
The evolution equation of the renormalization constants of the gluon or ghost fields in a MOM scheme, with respect to the coupling constant $`h`$ in an arbitrary scheme $`R`$ (the index $`R`$ is omitted but understood everywhere), can be written generically up to four-loop order:
$`{\displaystyle \frac{d\mathrm{ln}(Z_{\mathrm{\Gamma },MOM})}{d\mathrm{ln}\mu ^2}}=\overline{\gamma }_0h+\overline{\gamma }_1h^2+\overline{\gamma }_2h^3+\overline{\gamma }_3h^4`$ (34)
and the perturbative integration of Eq. (34) yields, to the same order,
$`\begin{array}{cc}\hfill \mathrm{ln}\left({\displaystyle \frac{Z_{\mathrm{\Gamma },MOM}}{Z_0}}\right)& =\mathrm{log}(h){\displaystyle \frac{\overline{\gamma }_0}{\beta _0}}+h{\displaystyle \frac{\left(\beta _0\overline{\gamma }_1\beta _1\overline{\gamma }_0\right)}{\beta _0^2}}\hfill \\ & +h^2{\displaystyle \frac{\left(\beta _0^2\overline{\gamma }_2\beta _0\beta _1\overline{\gamma }_1(\beta _0\beta _2\beta _1^2)\overline{\gamma }_0\right)}{2\beta _0^3}}\hfill \\ & +h^3(\beta _0^3\overline{\gamma }_3\beta _0^2\beta _1\overline{\gamma }_2+(\beta _0\beta _1^2\beta _0^2\beta _2)\overline{\gamma }_1\hfill \\ & +(\beta _0^2\beta _3+2\beta _0\beta _1\beta _2\beta _1^3)\overline{\gamma }_0){\displaystyle \frac{1}{3\beta _0^4}}\hfill \end{array}`$ (35)
with the standard four-loop formula for the running coupling
$`\begin{array}{cc}\hfill h(t)& ={\displaystyle \frac{1}{\beta _0t}}\left(1{\displaystyle \frac{\beta _1}{\beta _0^2}}{\displaystyle \frac{\mathrm{log}(t)}{t}}+{\displaystyle \frac{\beta _1^2}{\beta _0^4}}{\displaystyle \frac{1}{t^2}}\left(\left(\mathrm{log}(t){\displaystyle \frac{1}{2}}\right)^2+{\displaystyle \frac{\beta _2\beta _0}{\beta _1^2}}{\displaystyle \frac{5}{4}}\right)\right)\hfill \\ & +{\displaystyle \frac{1}{(\beta _0t)^4}}\left({\displaystyle \frac{\beta _3}{2\beta _0}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\beta _1}{\beta _0}}\right)^3\left(2\mathrm{log}^3(t)+5\mathrm{log}^2(t)+\left(46{\displaystyle \frac{\beta _2\beta _0}{\beta _1^2}}\right)\mathrm{log}(t)1\right)\right)\hfill \end{array}`$ (36)
and $`t=\mathrm{log}\left({\displaystyle \frac{\mu ^2}{\mathrm{\Lambda }^2}}\right)`$.
We now consider in turn the three renormalization schemes $`\overline{\text{MS}}`$, $`\stackrel{~}{\text{MOMg}}`$ and $`\stackrel{~}{\text{MOMc}}`$ and fit the two parameters of Eqs. (35) and (36) to our extrapolated lattice data. Figure 2 illustrates the typical quality of such fits.
### 5.1 $`\overline{\text{MS}}`$ scheme
The analysis in the $`\overline{\text{MS}}`$ scheme is summarized in Table 5. The statistical error is at the level of 1% for the gluon propagator and 2-3% for the ghost propagator, whereas the systematic error due to the extrapolations is around 3-5% and 5-10% respectively. The values of $`\mathrm{\Lambda }_{\overline{\text{MS}}}`$ extracted from the gluon and the ghost propagators are consistent within these errors and within each order of perturbation theory.
However, the three-loop and four-loop values, which are displayed in Table 6 with the physical units of Table 2, clearly confirm our previous result that we are still far from asymptoticity in that scheme.
### 5.2 $`\stackrel{~}{\text{MOMg}}`$ scheme
Table 7, which summarizes the analysis in the $`\stackrel{~}{\text{MOMg}}`$ scheme, shows that, at the lower $`\beta `$’s, we were not able to describe both lattice propagators at four-loop order with reasonable cuts and $`\chi ^2`$. This could be interpreted as an hint that perturbation theory has some problems of convergence beyond three-loop order below 3-4 GeV.
If we select the three-loop result as the best perturbative estimate of $`\mathrm{\Lambda }_{\stackrel{~}{\text{MOMg}}}`$ and convert it to the $`\overline{\text{MS}}`$ scheme with the asymptotic one-loop formula, $`\mathrm{\Lambda }_{\overline{\text{MS}}}=0.346\mathrm{\Lambda }_{\stackrel{~}{\text{MOMg}}}`$, then we get the physical values quoted in Table 8 which agree completely with previous values .
### 5.3 $`\stackrel{~}{\text{MOMc}}`$ scheme
The results of the analysis in the $`\stackrel{~}{\text{MOMc}}`$ scheme are displayed in Table 9. We still find that the three-loop and four-loop values of $`\mathrm{\Lambda }_{\stackrel{~}{\text{MOMc}}}^{(3)}`$ are very much the same both for the gluon propagator and for the ghost propagator. Thus the perturbative series seems again to become asymptotic at three-loop order in that scheme.
Selecting the three-loop result as the best perturbative estimate of $`\mathrm{\Lambda }_{\stackrel{~}{\text{MOMc}}}`$ and converting it to the $`\overline{\text{MS}}`$ scheme with the asymptotic formula, $`\mathrm{\Lambda }_{\overline{\text{MS}}}=0.429\mathrm{\Lambda }_{\stackrel{~}{\text{MOMc}}}`$, we get the physical values quoted in Table 10.
### 5.4 Scheme dependence
The puzzling feature of Tables 6, 8 and 10 is the rather large dependence of the $`\mathrm{\Lambda }_{\text{QCD}}`$ scale upon the loop-order and the renormalisation scheme whereas, within any scheme, the values from the ghost and gluon propagators are rather consistent at each loop order and pretty independent of the lattice spacing.
Let us consider again the evolution equation of the renormalisation constants of the gluon or ghost fields in a MOM scheme, with respect to the coupling $`h_R`$ in an arbitrary scheme $`R`$. We have
$`{\displaystyle \frac{d\mathrm{ln}(Z_{\mathrm{\Gamma },MOM})}{d\mathrm{ln}\mu ^2}}=\overline{\gamma }_R(h_R)={\displaystyle \frac{1}{2}}{\displaystyle \frac{d\mathrm{ln}(Z_{\mathrm{\Gamma },MOM})}{d\mathrm{ln}\mathrm{\Lambda }_R}}`$ (37)
where $`\mathrm{\Lambda }_R`$ is the scale in scheme $`R`$. If we truncate the perturbative expansion at order $`n`$
$`\mathrm{ln}\left({\displaystyle \frac{Z_{\mathrm{\Gamma },MOM}}{Z_0}}\right)=c_{R,0}\mathrm{ln}(h_R)+{\displaystyle \underset{k=1}{\overset{n1}{}}}c_{R,k}h_R^k`$ (38)
the change in the effective scale $`\mathrm{\Lambda }_R^{(n)}`$, or equivalently, the change in the coupling $`h_R`$, induced by adding the contribution at order $`n+1`$ is typically
$`{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Lambda }_R^{(n)}}{\mathrm{\Lambda }_R^{(n)}}}{\displaystyle \frac{c_{R,n}h_R^n}{2\gamma _R(h_R)}}`$ (39)
Now the dependence of the effective scale $`\mathrm{\Lambda }_R`$ upon the coupling $`h_R`$ is given up to order 4
$`2\mathrm{ln}\mathrm{\Lambda }_R^{(4)}=\mathrm{ln}\mu ^2{\displaystyle \frac{1}{\beta _0h_R}}{\displaystyle \frac{\beta _1}{\beta _0^2}}\mathrm{ln}(\beta _0h_R){\displaystyle \frac{\beta _0\beta _2\beta _1^2}{\beta _0^3}}h_R{\displaystyle \frac{\beta _0^2\beta _32\beta _0\beta _1\beta _2+\beta _1^3}{2\beta _0^4}}h_R^2`$ (40)
and, denoting the coefficient of order $`h_R^{n2}`$ in that equation by $`\rho _{R,n1}`$, the effective scales which describe a same coupling at order $`n`$ and $`n+1`$ are related by
$`\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }_R^{(n+1)}}{\mathrm{\Lambda }_R^{(n)}}}{\displaystyle \frac{1}{2}}\rho _{R,n1}h_R^{n1}`$ (41)
Combining Eqs. (39) and (41) gives the relation between the effective scales which describe the renormalisation constants of the gluon or ghost fields in a MOM scheme at order $`n`$ and $`n+1`$
$`{\displaystyle \frac{\mathrm{\Lambda }_R^{(n+1)}}{\mathrm{\Lambda }_R^{(n)}}}=\mathrm{exp}{\displaystyle \frac{1}{2}}\left(\rho _{R,n1}+{\displaystyle \frac{c_{R,n}h_R}{\gamma _R(h_R)}}\right)h_R^{n1}`$ (42)
Figure 3 displays the behavior of this ratio for the gluon and ghost propagators in the three schemes as a function of the momentum $`p`$ for $`n=2`$ and $`n=3`$. The couplings are taken from the fits at $`\beta =6.4`$.
There is a pretty good qualitative agreement with Tables 6, 8 and 10, which confirms the overall consistency with perturbation theory of the lattice data for the gluon and ghost propagators within any renormalization scheme.
The scheme dependence of the $`\mathrm{\Lambda }_{\text{QCD}}`$ scale can also be analyzed with Eq. (40):
$`{\displaystyle \frac{\mathrm{\Lambda }_{R_2}^{(n)}}{\mathrm{\Lambda }_{R_1}^{(n)}}}=\mathrm{exp}\left\{{\displaystyle \frac{1}{2\beta _0}}\left({\displaystyle \frac{1}{h_{R_1}}}{\displaystyle \frac{1}{h_{R_2}}}\right)+{\displaystyle \frac{\beta _1}{2\beta _0^2}}\mathrm{ln}{\displaystyle \frac{h_{R_1}}{h_{R_2}}}+\mathrm{}\right\}`$ (43)
Figure 4 shows the behavior of the ratios $`\frac{\mathrm{\Lambda }_{\overline{\text{MS}}}^{(n)}}{\mathrm{\Lambda }_{\stackrel{~}{\text{MOMg}}}^{(n)}}`$ and $`\frac{\mathrm{\Lambda }_{\stackrel{~}{\text{MOMc}}}^{(n)}}{\mathrm{\Lambda }_{\stackrel{~}{\text{MOMg}}}^{(n)}}`$, as a function of the momentum $`p`$ at each order of perturbation theory. The couplings are taken from the fits of the gluon propagator at $`\beta =6.4`$.
Clearly, the limiting values of these ratios are not the asymptotic values. If we replace in Eq. (43) the coupling $`h_{R_2}`$ by its perturbative expansion with respect to $`h_{R_1}`$
$`h_{R_2}=h_{R_1}+{\displaystyle \underset{k=1}{\overset{n1}{}}}r_kh_{R_1}^{k+1}`$ (44)
then the ratios do of course tend towards the asymptotic values $`\mathrm{exp}\left\{{\displaystyle \frac{r_1}{2\beta _0}}\right\}`$. The disagreement with respect to the perturbative expansion is not a problem with the lattice data or with the numerical analysis. Indeed the fits do a very good job at extracting a well-behave coupling as illustrated in Fig. 5 which displays the dimensionless scales $`a\mathrm{\Lambda }_{\overline{\text{MS}}}^{(4)}`$, $`a\mathrm{\Lambda }_{\stackrel{~}{\text{MOMg}}}^{(4)}`$ and $`a\mathrm{\Lambda }_{\stackrel{~}{\text{MOMc}}}^{(4)}`$ as a function of the momentum $`p`$, using Eq. (40) with the fitted couplings at $`\beta =6.4`$ from the ghost and gluon propagators. $`Z_0`$, the other fitted parameter of Eq. (35), is nearly independent, within a few percent, of the renormalisation scheme as it should in the absence of truncations. It follows that the difficulty to reproduce the asymptotic ratios between the scales of different renormalization schemes, is mainly a consequence of the truncation of the perturbative series of the renormalization constants of the gluon and ghost propagators.
We can substantiate this claim, and estimate the rate of convergence, by the following exercise. We solve $`h_{R_2}`$ in terms of $`h_{R_1}`$ using Eq. (38) at four-loop order
$`\mathrm{ln}\left({\displaystyle \frac{Z_{\mathrm{\Gamma },MOM}}{Z_0}}\right)=c_{R_2,0}\mathrm{ln}(h_{R_2})+{\displaystyle \underset{k=1}{\overset{3}{}}}c_{R_2,k}h_{R_2}^k=c_{R_1,0}\mathrm{ln}(h_{R_1})+{\displaystyle \underset{k=1}{\overset{3}{}}}c_{R_1,k}h_{R_1}^k`$ (45)
Then we plug the solution into Eq. (43). Figure 6 shows the behavior of the corresponding ratios, $`\frac{\mathrm{\Lambda }_{\overline{\text{MS}}}^{(4)}}{\mathrm{\Lambda }_{\stackrel{~}{\text{MOMg}}}^{(4)}}`$ and $`\frac{\mathrm{\Lambda }_{\stackrel{~}{\text{MOMc}}}^{(4)}}{\mathrm{\Lambda }_{\stackrel{~}{\text{MOMg}}}^{(4)}}`$, as a function of the coupling $`h_{\overline{\text{MS}}}`$ and $`h_{\stackrel{~}{\text{MOMc}}}`$ respectively. The effect of the truncation of the perturbative series is manifest for the $`\overline{\text{MS}}`$ scheme and gives the right order of magnitude of what is actually measured in Tables 5 and 7.
## 6 Conclusion
We have shown that the lattice formulation of the ghost propagator has the expected perturbative behavior up to four-loop order from 2 GeV to 6 GeV. We have been able to go beyond the qualitative level and to produce quantitative results for the scale $`\mathrm{\Lambda }_{\overline{\text{MS}}}`$ which are pretty consistent with the values extracted from the lattice gluon propagator. We have understood the strong dependence of the effective $`\mathrm{\Lambda }_{\overline{\text{MS}}}`$ scale upon the order of perturbation theory and upon the renormalisation scheme used for the parametrisation of the data. The perturbative series of the $`\stackrel{~}{\text{MOM}}`$ schemes seem to be asymptotic at three-loop order in the energy range we have probed whereas the $`\overline{\text{MS}}`$ scheme converges very slowly. If we assume that all perturbative series remain well behaved beyond four-loop above 4 GeV, then we get $`\mathrm{\Lambda }_{\overline{\text{MS}}}320`$ MeV with a 10% systematic uncertainty. The statistical errors are at the 1% level. This value is also in pretty good agreement with the values of $`\mathrm{\Lambda }_{\overline{\text{MS}}}`$ extracted from the three-gluon vertex in a $`\stackrel{~}{\text{MOM}}`$ scheme at three-loop order , at the same $`\beta `$’s and with the same lattice sizes. On the other hand it exceeds by 20% the value obtained from the same vertex at $`\beta =6.8`$ on a $`24^4`$ lattice. This discrepancy motivated the introduction of power corrections which are successful in describing the combined data of the three-gluon vertex . We will show in a forthcoming paper how the power corrections can be unraveled from the lattice propagators alone.
The value quoted above exceeds also by about 30% the previous determinations of the QCD scale in the quenched approximation based on gauge-invariant definitions of the strong coupling constant (take note, for comparison purposes, that our physical unit corresponds to the force parameter $`r_0`$ set approximately to 0.53 fm). However there is also an uncertainty due to the use of the asymptotic one-loop relation between $`\mathrm{\Lambda }_{\overline{\text{MS}}}`$ and the $`\mathrm{\Lambda }_\text{L}`$’s. For illustration, let us consider the determination of $`\mathrm{\Lambda }_{\overline{\text{MS}}}`$ using lattice perturbation theory up to three-loop order with the Wilson action . It is possible to estimate the rate of convergence of the ratio $`\frac{\mathrm{\Lambda }_\text{L}^{(3)}}{\mathrm{\Lambda }_{\overline{\text{MS}}}^{(3)}}`$ as a function of the bare lattice coupling $`h_\text{L}=\frac{6}{(4\pi )^2\beta }`$ by inserting the perturbative expansion of $`h_{\overline{\text{MS}}}`$ into Eq. (43). Figure 7 displays the evolution of this ratio and also of the ratio $`\frac{\mathrm{\Lambda }_{\mathrm{}}^{(3)}}{\mathrm{\Lambda }_{\overline{\text{MS}}}^{(3)}}`$ for the so-called “boosted” lattice scheme which re-express the lattice perturbative series as a function of the coupling $`h_{\mathrm{}}=h_\text{L}/plaq`$. The mere truncation of the perturbative series introduces an uncertainty on the absolute scale of the lattice schemes which could be as large as 30% in the range of $`\beta `$ studied in these simulations.
No strategy can fix the scale $`\mathrm{\Lambda }_{\text{QCD}}`$ to an accuracy better than the uncertainty entailed by the truncation of the perturbative series in the conversion to the $`\overline{\text{MS}}`$ scheme. We have shown that this error can be larger than the main well-known sources of systematic errors which come from setting the scale $`a^1`$ and from the continuum extrapolation. If we aim at reducing below 10% the error in the conversion of the $`\stackrel{~}{\text{MOM}}`$ schemes to the $`\overline{\text{MS}}`$ scheme, then a look at Figure 6 shows that we need to apply a cut at 6 GeV. Such an analysis would require simulations at $`\beta =6.6`$ and $`\beta =6.8`$ on $`48^4`$ and $`64^4`$ lattices respectively, to work at fixed volume and minimize finite-size effects. The existence of several lattice observables, gluon propagator, ghost propagator, three-gluon vertex, from which one can extract independent values of the scale $`\mathrm{\Lambda }_{\text{QCD}}`$, an advantage of the Green function approach, should then allow to disentangle unambiguously the effects of the truncation of the perturbative series from the non-perturbative corrections, and to get a value of $`\mathrm{\Lambda }_{\overline{\text{MS}}}`$ at a true 10% accuracy.
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# 1 Introduction and Overview
## 1 Introduction and Overview
Black holes are important objects to investigate the nature of space-time beyond the description by the general relativity. For example, the study of the black hole entropy provides the holographic principle , which gives a strong constraint to the property of quantum gravity. The microscopic properties of black holes are also intensively investigated from the viewpoint of string theory. For example, many black hole (or brane) solutions are known in the low energy effective theory of the superstring theory. In particular, for BPS saturated solutions of Type II (or Type I) supergravity, it is well known that there is a correspondence between the classical solutions and the systems of the BPS D-branes . In this sense, it is natural to assume that there is a stringy origin for a general classical solution of Type II supergravity, even for non-BPS systems.
In our previous paper , we treated a general classical solution of Type II supergravity with the symmetry $`ISO(1,p)\times SO(9p)`$, which is called as the “three-parameter solution” . The solution had been thought to be the low energy counterpart of the D$`p`$$`p`$-brane system with a constant tachyon profile (vacuum expectation value (VEV)), where the three parameters had been supposed to correspond to the number of the D-branes, that of the anti D-branes and the VEV of the tachyon field . In , however, we re-examined this correspondence and found that the above conjecture is not correct and the three-parameter solution is the low energy counterpart of the D$`p`$$`p`$-brane system with a constant tachyon VEV only if one of the three parameters, $`c_1`$, is tuned to zero.
In this paper, we consider the stringy interpretation of missing $`c_1`$ and give some explicit examples of the stringy counterparts of the three-parameter solution. As we mentioned in the previous paper , $`c_1`$ is related to the “dilaton charge”, which characterizes the asymptotic behavior of the dilaton field. It is well known that black holes with horizons do not have the dilaton charge in general, which is supported by the no-hair theorem, but some exceptions are known and we have not understood this phenomena systematically . One of the purposes of this paper is to give an approach to this problem from the microscopic (string theoretical) point of view.
Our strategy is making use of the known correspondence between the BPS black $`p`$-brane solution in the supergravity and the BPS D$`p`$-branes in the string theory, and extending this correspondence by appropriately deforming the both sides. On the supergravity side, this leads to a new characterization of the solution. We will find that the BPS limit of the three-parameter solution consists of two extremal limit, that is, the ordinary BPS condition and the condition that the dilaton charge is proportional to the ADM mass. Then, the three-parameter solution is characterized by two non-extremality parameters (and the RR-charge). In particular, the parameter $`c_1`$ is related to the second non-extremality. On the string theory side, we will formally deform the boundary state for the BPS D$`p`$-brane, which respects the same symmetry of the three-parameter solution. By evaluating the long distance behavior of the solution of the supergravity and the massless emission from the deformed boundary state, we will directly compare the (microscopic) deformation parameters of the boundary state and the (macroscopic) parameters in the three-parameter solution. From this analysis, we show that the first kind of non-extremality corresponds to the non-BPS nature, while the second kind (hence $`c_1`$) is related to the deformation of the boundary condition from the ordinary Neumann/Dirichlet one.
As examples of the deformed boundary states, we consider three systems of D-branes. The first one is the system we consider in , which possesses only the first non-extremality since the D$`p`$$`p`$-system follows the ordinary boundary condition, that is, the Neumann boundary condition for the longitudinal directions and the Dirichlet boundary condition for the transverse directions. The second and the third example are systems of unstable D9-branes and unstable D-instantons, respectively, with an appropriate tachyon condensation to the BPS D$`p`$-branes, which actually have non-trivial dilaton charges, i.e. non-zero value of $`c_1`$. They possess the deformation of the boundary condition for the worldsheet, which represent in some sense fuzziness in transverse or longitudinal directions, respectively. From these examples, we discuss that deformed boundary states which respect the symmetry of the three-parameter solution can be obtained by considering D$`p`$$`p`$-systems on which tachyonic and/or massive excitations are turned, in general.
The organization of this paper is as follows. In the next section, we review the three-parameter solution and we propose a new characterization of the solution in terms of two non-extremality parameters. In the section 3, we give the string theoretical interpretation to the parameter $`c_1`$ and give three examples of the deformed boundary states. The section 4 is devoted to the conclusions and discussions.
## 2 The Three-Parameter Solution and Non-Extremalities
In this section, we first give a short review of the three-parameter solution. Then, by investigating the BPS limit of the solution, we propose a new characterization of the solution in terms of two non-extremality parameters.
First of all, we consider Type II supergravity in the following setting:
* Assume a fixed, $`(p+1)`$-dimensional object as a source carrying only RR $`(p+1)`$-form charge.
* Spacetime has the symmetry $`ISO(1,p)\times SO(9p)`$ and it is asymptotically flat.
Note that these ansatzes are the same as those for the BPS black $`p`$-brane solution (we will consider the region $`0p6`$ in this paper). Under these conditions, it is sufficient to consider the metric, the dilaton and the RR $`(p+1)`$-form field. The relevant part of the ten-dimensional action (in the Einstein frame) is given by
$$S=\frac{1}{2\kappa ^2}d^{10}x\sqrt{g}\left[R\frac{1}{2}(\varphi )^2\frac{1}{2(p+2)!}e^{\frac{3p}{2}\varphi }|F_{p+2}|^2\right],$$
(2.1)
where $`F_{(p+2)}`$ denotes the $`(p+2)`$-form field strength which is related to the $`(p+1)`$-form potential of the RR-field $`𝒜^{(p+1)}`$ as $`F_{(p+2)}=d𝒜_{(p+1)}`$. According to the symmetry $`ISO(1,p)\times SO(9p)`$, we should impose the following ansatz,
$`ds^2`$ $`=g_{MN}dx^Mdx^N`$
$`=e^{2A(r)}\eta _{\mu \nu }dx^\mu dx^\nu +e^{2B(r)}\delta _{ij}dx^idx^j`$
$`=e^{2A(r)}\eta _{\mu \nu }dx^\mu dx^\nu +e^{2B(r)}(dr^2+r^2d\mathrm{\Omega }_{(8p)}^2),`$
$`\varphi `$ $`=\varphi (r),`$
$`𝒜^{(p+1)}`$ $`=e^{\mathrm{\Lambda }(r)}dx^0dx^1\mathrm{}dx^p,`$ (2.2)
where $`M=(\mu ,i)`$, $`\mu ,\nu =0,\mathrm{},p`$ are indices of the longitudinal directions of the $`p`$-brane, $`i,j=p+1,\mathrm{},9`$ express the orthogonal directions, and $`r^2=x^ix_i`$.
The authors of gave a general solution of Type II supergravity of this system. As the solution includes three integration constants, it is called as the three-parameter solution, which is given by<sup>1</sup><sup>1</sup>1 In , the authors consider the $`D`$-dimensional gravity theory and have constructed a general solution with the less symmetry $`SO(p)\times SO(Dp1)`$ as the four-parameter solution. The three-parameter solution is a restricted one.
$`A(r)`$ $`={\displaystyle \frac{(7p)(3p)c_1}{64}}h(r){\displaystyle \frac{7p}{16}}\mathrm{ln}\left[\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))\right],`$ (2.3)
$`B(r)`$ $`={\displaystyle \frac{1}{7p}}\mathrm{ln}[f_{}(r)f_+(r)]`$
$`{\displaystyle \frac{(p+1)(3p)c_1}{64}}h(r)+{\displaystyle \frac{p+1}{16}}\mathrm{ln}\left[\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))\right],`$ (2.4)
$`\varphi (r)`$ $`={\displaystyle \frac{(p+1)(7p)c_1}{16}}h(r)+{\displaystyle \frac{3p}{4}}\mathrm{ln}\left[\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))\right],`$ (2.5)
$`e^{\mathrm{\Lambda }(r)}`$ $`=\eta (c_2^21)^{1/2}{\displaystyle \frac{\mathrm{sinh}(kh(r))}{\mathrm{cosh}(kh(r))c_2\mathrm{sinh}(kh(r))}},`$ (2.6)
where
$`f_\pm (r)`$ $`1\pm {\displaystyle \frac{r_0^{7p}}{r^{7p}}},`$ (2.7)
$`h(r)`$ $`\mathrm{ln}\left({\displaystyle \frac{f_{}}{f_+}}\right),`$ (2.8)
$`k`$ $`\pm \sqrt{{\displaystyle \frac{2(8p)}{7p}}{\displaystyle \frac{(p+1)(7p)}{16}}c_1^2}`$
$`\pm {\displaystyle \frac{\sqrt{(p+1)(7p)}}{4}}\sqrt{c_m^2c_1^2},\left(c_m=\sqrt{{\displaystyle \frac{32(8p)}{(p+1)(7p)^2}}}\right),`$ (2.9)
$`\eta `$ $`\pm 1,`$ (2.10)
where $`\eta `$ denotes the sign of the RR-charge. The three parameters<sup>2</sup><sup>2</sup>2 We have labeled $`c_3`$ of as $`c_2`$ and $`k`$ as $`k`$ according to ., $`r_0,c_1,c_2`$, are the integration constants that parametrize the solution. As in , the domain of the parameters in the solution (2.3)–(2.6) we take is
$`c_1`$ $`[0,c_m],`$ (2.11)
$`c_2`$ $`(\mathrm{},1][1,\mathrm{}),`$ (2.12)
$`r_0^{7p}`$ $`𝐑.`$ (2.13)
where we have already fixed the $`𝐙_2`$ symmetries of the solution,
$`(r_0^{7p},c_1,c_2,sgn(k),\eta )`$ $`(r_0^{7p},c_1,c_2,sgn(k),\eta ),`$
$`(r_0^{7p},c_1,c_2,sgn(k),\eta )`$ $`(r_0^{7p},c_1,c_2,sgn(k),\eta ),`$ (2.14)
by choosing $`c_10`$ and $`k0`$. Furthermore, we have a degree of freedom to choose the signs of $`r_0^{7p}`$ and $`c_2`$. As we mentioned in Sec.1, we will discuss the BPS black $`p`$-brane limit of the three-parameter solution and in order to take that limit consistently, we must choose one of the two branches,
$`\{\begin{array}{cc}r_0^{7p}0,c_2>0\hfill & \text{: branch}I_+,\hfill \\ r_0^{7p}0,c_2<0\hfill & \text{: branch}I_{}.\hfill \end{array}`$ (2.15)
From the viewpoint of the gravity theory, the three-parameter solution describes charged dilatonic black objects. Thus, the RR-charge $`Q`$, the ADM mass $`M`$ and the dilaton charge $`D`$ are natural quantities to characterize the solution. For convenience, we consider wrapping the spatial worldvolume directions on a torus $`T^p`$ of volume $`V_p`$. The RR-charge is given by an appropriate surface integral over the sphere-at-infinity in the transverse directions ;
$$|Q|=2(c_2^21)^{1/2}kN_pr_0^{7p},$$
(2.16)
where
$$N_p\frac{(8p)(7p)\omega _{8p}V_p}{16\kappa ^2},$$
(2.17)
and $`\omega _d=\frac{2\pi ^{(d+1)/2}}{\mathrm{\Gamma }((d+1)/2)}`$ is the volume of the unit sphere $`S^d`$. The ADM mass is defined in as
$$g_{00}=1+\frac{7p}{8}\frac{M}{N_p}\frac{1}{r^{7p}}+𝒪\left(\frac{1}{r^{2(7p)}}\right),$$
(2.18)
where (2.17) is used. This gives us
$$M=\left(2c_2k+\frac{3p}{2}c_1\right)N_pr_0^{7p},$$
(2.19)
Similarly, we define the dilaton charge by its asymptotic behavior in the transverse direction as <sup>3</sup><sup>3</sup>3 Their original definition is in the string frame but here we define it in the Einstein frame in the same manner.
$$\varphi =\frac{D}{N_p}\frac{1}{r^{7p}}+𝒪\left(\frac{1}{r^{2(7p)}}\right),$$
(2.20)
which gives
$$D=\left(\frac{3p}{2}c_2k\frac{(p+1)(7p)}{8}c_1\right)N_pr_0^{7p}.$$
(2.21)
Here we give some historical remarks on the dilaton charge. As we referred in , the dilaton charge is a familiar quantity in the general relativity and it has been investigated for a long time in various contexts, for example, the no-hair theorem , the cosmic censorship, the scalar field collapse . It is well known that the existence of the dilaton charge would change the structures of spacetimes drastically and generally cause naked singularities to appear<sup>4</sup><sup>4</sup>4There are some exceptions which have horizons even if there is non-zero dilaton charge as we mentioned in the introduction. These solutions may be related not to the three-parameter solution, but to the four-parameter solution. The three-parameter and the four-parameter solution coincide with each other in the case of $`D=4`$ and $`p=0`$, so they are not distinguishable in such a case. We will discuss the relation between the condition of the formation of horizons and the parameters of the supergravity solutions in the forthcoming paper.. The extensions to the higher-dimensional theories or the stabilities of such singular spacetimes have been also investigated by many authors . In spite of those investigations, we have not systematically understood the relation between the horizon and the dilaton charge, for example, what kind of the dilaton couplings to the RR-charge can avoid to make a spacetime singular.
Since the three-parameter solution is the most general solution with the symmetry $`ISO(1,p)\times SO(9p)`$, it naturally includes the BPS black $`p`$-brane solution;
$`ds^2`$ $`=`$ $`f(r)^{\frac{7p}{8}}\eta _{\mu \nu }dx^\mu dx^\nu +f(r)^{\frac{p+1}{8}}\delta _{ij}dx^idx^j,`$ (2.22)
$`e^\varphi (r)`$ $`=`$ $`f(r)^{\frac{3p}{4}},`$ (2.23)
$`e^{\mathrm{\Lambda }(r)}`$ $`=`$ $`\eta \left(f(r)^11\right),`$ (2.24)
where
$$f(r)=1+\frac{\mu _0}{r^{7p}}.$$
(2.25)
Note that $`\mu _0`$ is the only parameter of the solution, and therefore, all the quantities are labeled by it as
$$|Q|=M=N_p\mu _0,D=\frac{3p}{4}N_p\mu _0,$$
(2.26)
where $`\mu _0`$ takes an arbitrary non-negative value. In this paper, we will restrict ourselves to the charged case, i.e., $`\mu _0>0`$.
As we claimed in , the parametrization by $`(c_1,c_2,r_0)`$ is not suitable to take the BPS limit, and so we defined a set of parameters $`(c_1,\mu _0,\nu )`$ which was defined by
$$r_0^{7p}\frac{\nu \mu _0}{2k},c_2^21\frac{1}{\nu ^2}.$$
(2.27)
The domain of $`\nu `$ becomes
$`\{\begin{array}{cc}0\nu <\mathrm{}\hfill & \text{for the branch}I_+,\hfill \\ \mathrm{}<\nu 0\hfill & \text{for the branch}I_{},\hfill \end{array}`$ (2.28)
since we choose $`c_10,k0`$ and $`\mu _0>0`$. Using these parameters, the RR-charge, the ADM mass and the dilaton charge are rewritten as
$`|Q|`$ $`=N_p\mu _0,`$ (2.29)
$`M`$ $`=\left(\sqrt{1+\nu ^2}+{\displaystyle \frac{3p}{4}}{\displaystyle \frac{c_1\nu }{k}}\right)N_p\mu _0,`$ (2.30)
$`D`$ $`=\left({\displaystyle \frac{3p}{4}}\sqrt{1+\nu ^2}{\displaystyle \frac{(p+1)(7p)}{4}}{\displaystyle \frac{c_1\nu }{k}}\right)N_p\mu _0={\displaystyle \frac{3p}{4}}M{\displaystyle \frac{c_1\nu }{k}}N_p\mu _0.`$ (2.31)
Then, it is easy to show that the three-parameter solution (2.3)–(2.6) reduces to the BPS black $`p`$-brane solution (2.22)–(2.24) in the limit $`\nu 0`$ for arbitrary values of $`c_1`$ and $`\mu _0`$. Since $`|Q|`$ is unchanged in this limit, we can set $`\mu _0`$ to the same value as the black $`p`$-brane<sup>5</sup><sup>5</sup>5 For the neutral case, i.e., $`|Q|=0`$, another parametrization is suitable. . We emphasize here that since there are three parameters, the BPS limit consists of two extremal limits, i.e. $`M=|Q|`$ and $`D=\frac{3p}{4}M`$. The former is the ordinary BPS condition, which guarantees the preservation of the spacetime supersymmetry. Although the latter says that the dilaton charge is proportional to the ADM mass, its physical meaning is obscure at this stage. Nevertheless, if we take the BPS limit as a starting point, it is quite natural to characterize the solution by two non-extremality parameters, that represent the deviation from the BPS black $`p`$-brane solution.
From this observation, we introduce following quantities $`m`$ and $`d`$ defined by
$`m`$ $`=\sqrt{1+\nu ^2},`$ (2.32)
$`d`$ $`={\displaystyle \frac{c_1\nu }{k}},`$ (2.33)
which are related to the quantities above as
$`|Q|`$ $`=N_p\mu _0,`$ (2.34)
$`M`$ $`=\left(m+{\displaystyle \frac{3p}{4}}d\right)N_p\mu _0,`$ (2.35)
$`D`$ $`=\left({\displaystyle \frac{3p}{4}}m{\displaystyle \frac{(p+1)(7p)}{16}}d\right)N_p\mu _0.`$ (2.36)
The region of these quantities are determined as
$`\{\begin{array}{cc}m1,d0\hfill & \text{for the branch}I_+,\hfill \\ m1,d0\hfill & \text{for the branch}I_{}.\hfill \end{array}`$ (2.37)
The BPS solution lies in $`m=1`$ and $`d=0`$ in the overlap of $`I_+`$ and $`I_{}`$. Note that the ADM mass and the dilaton charge are not conserved Noether charges and are frame-dependent quantities. Although their definition is ambiguous in some sense, they can characterize the deviation from the Minkowski space. On the other hand, our quantities $`m`$ and $`d`$ represent the deviation from the BPS black $`p`$-brane solution. As opposed to the ADM mass $`M`$, $`m`$ is always greater than $`1`$. Hence, $`m`$ can be considered as a non-extremality parameter in the ordinary sense, while $`d`$ denotes the difference between the ADM mass $`M`$ and the dilaton charge $`D`$. Note that $`d`$ is related to the parameter $`c_1`$ directly. As we have shown in , the three-parameter solution with $`c_1=0`$ denotes the spacetime which is produced from the D$`p`$$`p`$-brane system with a constant tachyon VEV. Since a main purpose of this paper is to clarify the stringy meaning of $`c_1`$, our task is to find sources in the string theory which produces $`d0`$ structure.
At the end of this section, we give the long distance behavior of the three-parameter solution (2.3)–(2.6), which is given by the leading term of $`1/r`$ expansion. In terms of $`(\mu _0,m,d)`$ it is given by
$`e^{2A(r)}`$ $`=1{\displaystyle \frac{7p}{8}}\left(m+{\displaystyle \frac{3p}{4}}d\right){\displaystyle \frac{\mu _0}{r^{7p}}}+𝒪\left(1/r^{2(7p)}\right),`$ (2.38)
$`e^{2B(r)}`$ $`=1+{\displaystyle \frac{p+1}{8}}\left(m+{\displaystyle \frac{3p}{4}}d\right){\displaystyle \frac{\mu _0}{r^{7p}}}+𝒪\left(1/r^{2(7p)}\right),`$ (2.39)
$`\varphi (r)`$ $`=\left({\displaystyle \frac{3p}{4}}m{\displaystyle \frac{(p+1)(7p)}{16}}d\right){\displaystyle \frac{\mu _0}{r^{7p}}}+𝒪\left(1/r^{2(7p)}\right),`$ (2.40)
$`e^{\mathrm{\Lambda }(r)}`$ $`=\eta {\displaystyle \frac{\mu _0}{r^{7p}}}+𝒪\left(1/r^{2(7p)}\right).`$ (2.41)
We will compare them with the massless emission from the source in the string theory in the next section.
## 3 Excitations on the D$`p`$$`p`$-brane System
In this section, we investigate stringy counterparts of the three-parameter solution. We argue that the geometry expressed by the three-parameter solution is made by D-brane systems which are in general sources of closed strings in the bulk. We stress that we only consider static sources and do not take into account interactions of open strings on the D-brane. Therefore they can be consistent sources of supergravity fields, even if they are unstable system. As a first step, we deform the boundary state for BPS D$`p`$-branes appropriately, and compare the long distance behavior of the solution with the massless emissions from them. And then, we give some examples.
### 3.1 Deformations of Boundary States
To obtain the three-parameter solution, we have required that the solution has the symmetry $`ISO(1,p)\times SO(9p)`$ and it carries a RR $`(p+1)`$-form charge as explained in the previous section. Therefore the source of closed strings, which produces the three-parameter solution, should respect the following restrictions at least in the low energy regime;
$`1)`$ $`\text{has the symmetry }ISO(1,p)\times SO(9p),`$
$`2)`$ carries only the RR $`(p+1)`$-form charge, (3.1)
$`3)`$ has the $`\delta `$-function distribution in the transverse space.
Note that the third condition can be slightly relaxed (see below).
Needless to say, the system of coincident BPS D$`p`$-branes is such an example. $`N`$ BPS D$`p`$-branes are expressed by the boundary state,
$$|Dp=𝒫_{\mathrm{GSO}}\left(N|Bp_{\mathrm{NS}}+N|Bp_{\mathrm{RR}}\right),$$
(3.2)
with<sup>6</sup><sup>6</sup>6 In this expression, we have omitted the label “$`+`$” of the spin structure.
$$|Bp_{\mathrm{NS}(\mathrm{RR})}=\frac{T_p}{2}\mathrm{exp}\left[\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\alpha _n^MS_{MN}\stackrel{~}{\alpha }_n^N+i\underset{r>0}{\overset{\mathrm{}}{}}b_r^MS_{MN}\stackrel{~}{b}_r^N\right]|p_\mu =0,x^i=0_{\mathrm{NS}(\mathrm{RR})},$$
(3.3)
where $`T_p`$ is the tension for a single D$`p`$-brane, $`\mu =0,1,\mathrm{},p`$ are the directions longitudinal to the worldvolume, $`i=p+1,\mathrm{},9`$ are the directions transverse to the D$`p`$-branes, and $`\alpha _n^M`$ and $`b_r^M`$ ($`\stackrel{~}{\alpha }_n^M`$ and $`\stackrel{~}{b}_r^M`$) are the creation operators of the modes of the left-moving (right-moving) worldsheet bosons and fermions, respectively. $`S_{MN}=\mathrm{diag}(\eta _{\mu \nu },\delta _{ij})`$ gives Neumann (Dirichlet) boundary conditions on the string worldsheet in $`\mu (i)`$ directions, respectively. Here the number of the D$`p`$-branes $`N`$ is the only microscopic parameter. In this case, it is shown that massless emission from this boundary state agrees with the long distance behavior of the BPS black $`p`$-brane under the identification,
$`\mu _0={\displaystyle \frac{2\kappa NT_p}{(7p)\omega _{8p}}},`$ (3.4)
which relates the macroscopic parameter $`\mu _0`$ with the microscopic parameter $`N`$. Of course, the system with $`\overline{N}`$ $`\overline{\mathrm{D}}p`$-brane also satisfies the same ansatz, which is simply given by replacing $`N`$ to $`\overline{N}`$ in the RR-sector of (3.2)<sup>7</sup><sup>7</sup>7 In this case, (3.4) is same but $`\eta `$ replaced by $`1`$. . We can also consider the system with $`N`$ D$`p`$-branes and $`\overline{N}`$$`p`$-branes without tachyon field which is given by the linear combination of the both above.
We now take the BPS boundary state (3.2) as the starting point as mentioned in Sec.1 and consider possible deformations of it with keeping the ansatzes (3.1), which should be the counterpart of the discussion in Sec.2, that is, on the supergravity side. One might try to deform it with massless open string excitations on the brane. However, the symmetry $`ISO(1,p)\times SO(9p)`$ prevents the fluctuation of massless scalar fields, and there should be no gauge flux since otherwise it produces other NSNS/RR-charges. Therefore, the excitations on the branes should be at least tachyonic and/or massive<sup>8</sup><sup>8</sup>8 Strictly speaking, the linear combination of the boundary state and other closed string states is also possible, but we do not consider this option. .
In order to make the discussion transparent, we formally deform the boundary state as follows. In the NSNS-sector, we set
$$|B_p^{}_{NS}=CN\frac{T_p}{2}\mathrm{exp}\left[\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\alpha _n^MS_{MN}^{(n)}\stackrel{~}{\alpha }_n^N+i\underset{r>0}{}b_r^MS_{MN}^{(r)}b_r^N\right]|p_\mu =0,x^i=0_{NS},$$
(3.5)
where
$$S_{MN}^{(n,r)}=\mathrm{diag}(A^{(n,r)}\eta _{\mu \nu },B^{(n,r)}\delta _{ij}).$$
(3.6)
Here $`A^{(n,r)}`$, $`B^{(n,r)}`$ and $`C`$ are deformation parameters which deform the boundary state without changing its symmetry. When we take $`A^{(n,r)}=B^{(n,r)}=C=1`$, the boundary state (3.3) for the $`N`$ BPS D$`p`$-branes is reproduced. We also deform the RR-sector of the boundary state (3.3) as the same manner but $`C=1`$ in this case since we fix the RR-charge to be the same as the $`N`$ D$`p`$-branes. We give some comments in order. First, one might assume the source object is a $`\delta `$-function type located at the origin in the transverse space. But it will turn out to be too restrictive and can be relaxed slightly: the “width” of the source allowed to be less than the string length. We will come back to this point in the next subsection. Second, the effect of $`(A^{(n,r)},B^{(n,r)})`$ in $`S_{MN}^{(n,r)}`$ is the deformation of the boundary condition. Here we note that only the $`r=1/2`$ mode contributes the massless mode of closed strings emitted from the boundary state. In other words, the deformations for other modes, $`S_{MN}^{(n)}`$ and $`S_{MN}^{(r1/2)}`$, are arbitrary in our analysis below. However, they should be related among them by other consistency, such as the modular invariance although we do not discuss them at this stage. Third, $`C`$ denotes the difference of the overall normalization in the NSNS-sector from that in the RR-sector. For example, in the system of $`(N+M)`$ D$`p`$-branes and $`M`$ $`\overline{\mathrm{D}}p`$-branes, we take $`CN=N+2M`$, the total number of branes, while the RR-sector has a fixed charge proportional to $`N`$.
Then, we can calculate the massless emissions with momentum $`p_M`$ from this deformed state as we have done in . In the NSNS-sector there are the graviton $`h_{MN}`$ and the dilaton $`\varphi `$, that are
$`f(p)`$ $`=`$ $`f(p)\left|\mathrm{\Delta }\right|D^{}p=CN{\displaystyle \frac{T_p}{2}}{\displaystyle \frac{V_{p+1}}{p_i^2}}ϵ_{MN}^fS^{(1/2)MN},`$ (3.7)
where $`|D^{}p`$ denotes the GSO projected deformed state, $`\mathrm{\Delta }`$ is the closed string propagator and $`ϵ_{MN}^f`$ is the appropriate polarization tensor for $`f=h_{MN},\varphi `$ . There is a similar expression for the RR $`(p+1)`$-form $`e^\mathrm{\Lambda }`$. In this case, since we set $`C=1`$ for the RR-sector and since only the zero-mode part contributes to the massless fields, we obtain same result as the BPS D$`p`$-brane. After the Fourier transformation to the position space, we obtain
$`h_{MN}(r)`$ $`=({\displaystyle \frac{7p}{8}}\eta _{\mu \nu },{\displaystyle \frac{p+1}{8}}\delta ij){\displaystyle \frac{C(A+B)}{2}}{\displaystyle \frac{2\kappa NT_p}{(7p)\omega _{8p}}}{\displaystyle \frac{1}{r^{7p}}},`$ (3.8)
$`\varphi (r)`$ $`={\displaystyle \frac{C}{8}}\left[(p+1)A+(7p)B\right]{\displaystyle \frac{2\kappa NT_p}{(7p)\omega _{8p}}}{\displaystyle \frac{1}{r^{7p}}},`$ (3.9)
$`e^\mathrm{\Lambda }(r)`$ $`=\pm {\displaystyle \frac{2\kappa NT_p}{(7p)\omega _{8p}}}{\displaystyle \frac{1}{r^{7p}}},`$ (3.10)
where we write $`A=A^{(1/2)}`$ and $`B=B^{(1/2)}`$. Comparing these quantities (3.8)-(3.10) with the long distance behaviors (2.38)-(2.41), we obtain the relation between the macroscopic quantities and the microscopic ones as
$`\mu _0={\displaystyle \frac{2\kappa NT_p}{(7p)\omega _{8p}}},`$ (3.11)
$`m={\displaystyle \frac{C}{8}}\left[(p+1)A+(7p)B\right],`$ (3.12)
$`d={\displaystyle \frac{C}{2}}(AB).`$ (3.13)
We see from the eqs. (3.12) and (3.13) that the two non-extremality parameters $`m`$ and $`d`$ in the three-parameter solution correspond to the deformation parameter $`(A,B,C)`$ on the string theory side. The first non-extremality $`m1`$ can be realized even if $`A=B=1`$, by taking $`C1`$. It means that the breaking of the BPS condition can occur within the ordinary boundary condition. However, the second non-extremality $`d0`$ requires necessarily $`AB`$. Therefore, in order to have the three-parameter solution with nonzero $`c_1`$, it is necessary to deform the boundary condition from the ordinary one of the BPS D$`p`$-branes so as to satisfy $`AB`$. In the next subsection, we give examples of such stringy objects.
### 3.2 Examples
In this subsection, we give three typical examples realizing the deformed states in the string theory which we referred in the previous subsection by considering various types of the tachyon condensation. The first example is the case considered in . The other two examples are obtained from the tachyon condensation from higher/lower dimensional unstable D-brane systems. In general, deformations are interpreted as tachyonic/massive excitations on the D$`p`$$`p`$ system.
#### 3.2.1 Case 1: D$`p`$$`p`$ system with a tachyon VEV
The first example is the D$`p`$$`p`$ system with a constant tachyon profile. To be precise we consider the $`(N+M)`$ D$`p`$-branes and $`M`$ $`\overline{\mathrm{D}}p`$-branes. The gauge symmetry of the worldvolume theory is $`U(N+M)\times U(M)`$ and there is a complex tachyon field $`T(x)`$ in the bi-fundamental representation of the gauge group. Here we consider the case where the $`M`$ D$`p`$$`p`$-pairs vanish and $`N`$ D$`p`$-branes remain. Namely, we decompose the $`M\times (N+M)`$ matrix by $`M\times N`$ and $`M\times M`$ components and set the tachyon profile as
$`T(x)=\left(\begin{array}{c}0\\ \\ t\end{array}\right),`$ (3.16)
where $`t`$ is a constant $`M\times M`$ matrix. Note that other open string excitations, such as gauge fields, massless scalar fields and a non-constant tachyon, break the symmetry $`ISO(1,p)\times SO(9p)`$.
It is known that this system can be expressed by the off-shell boundary state on which the boundary interaction for the tachyon field is turned . For our case, the NSNS-sector of the boundary state is then given by
$`|B_p^{};t_{\mathrm{NS}}`$ $`=\left[N+2\mathrm{T}\mathrm{r}_Me^{|t|^2}\right]{\displaystyle \frac{T_p}{2}}e^{_{n=1}^{\mathrm{}}\frac{1}{n}\alpha _n^MS_{MN}\stackrel{~}{\alpha }_n^N+i_{r>0}b_r^MS_{MN}\stackrel{~}{b}_r^N}|p_\mu =0,x^i=0_{\mathrm{NS}},`$ (3.17)
where $`S_{MN}=(\eta _{\mu \nu },\delta _{ij})`$ and the trace is taken over $`M\times M`$ matrices. Note that it is the same as (3.3) except for the overall factor, $`\left[N+2\mathrm{T}\mathrm{r}_Me^{|t|^2}\right]`$, which simply comes from the tachyon potential $`V(T)\mathrm{Tr}e^{|T|^2}`$. It reduces to the $`N+2M`$ in the limit $`|t|0`$, corresponding to the $`(N+M)`$ D$`p`$-branes and $`M`$ $`\overline{\mathrm{D}}p`$-branes, while reduces to $`N`$ in the limit $`|t|\mathrm{}`$, corresponding to the $`N`$ D$`p`$-branes.
From this state, we can read off $`(A,B,C)`$ by comparing (3.17) with (3.5) as
$$A=B=1,C=1+\frac{2}{N}\mathrm{Tr}_Me^{|t|^2}.$$
(3.18)
Therefore, macroscopic quantities $`m`$ and $`d`$ are given by
$$m=1+\frac{2}{N}\mathrm{Tr}_Me^{|t|^2},d=0.$$
(3.19)
We see that the constant tachyon $`t`$ only contributes to the ADM mass. Therefore, this system corresponds to at most two-parameter subset of the full three-parameter solution as we mentioned in , but it is an trivial example for our present purpose.
#### 3.2.2 Case 2: Unstable D9-branes to BPS D$`p`$-branes
Next example is the tachyon condensation from unstable D$`9`$-branes to the BPS D$`p`$-branes. For simplicity, we consider Type IIB superstring theory, that is, the tachyon condensation from pairs of D9-branes and D̄9-branes to BPS D$`p`$-branes. For Type IIA superstring theory, we start from non-BPS D9-branes and the discussion is completely parallel. Since a single BPS D$`p`$-brane is obtained from a pair of D(p+2)-brane and D̄(p+2)-brane, it is clear that it is obtained from $`2^{\frac{7p}{2}}`$ pairs of D9-brane and D̄9-brane via the tachyon condensation. Then let us consider $`N\times 2^{\frac{7p}{2}}`$ pairs of D9D̄9-branes and the tachyon profile,
$$T(x)=u\gamma _ix^i1_N,$$
(3.20)
where $`i=p+1,\mathrm{},9`$ are the transverse directions to the resulting D$`p`$-brane, $`\gamma _i`$ are the $`SO(9p)`$ gamma matrices and $`u`$ is a real parameter of dimension of mass. It is known that this system reduces to $`N`$ BPS D$`p`$-branes in the limit $`u\mathrm{}`$<sup>9</sup><sup>9</sup>9For details we refer the literatures, which discussed in the context of the worldvolume field theory , of the boundary string field theory and of the boundary state . . If we replace $`T`$ to $`T`$, the final state becomes D̄$`p`$-branes.
Here the effect of the tachyon profile (3.20) with $`u0`$ is twofold: First, it breaks the original global symmetry $`ISO(1,9)`$ down to $`ISO(1,p)\times SO(9p)`$. Next, it produces the correct RR-charge of $`N`$ D$`p`$-branes which is independent of the value of $`u`$. Therefore, this system satisfies the first two items of the ansatzes (3.1) for an arbitrary value of $`u0`$. However, accurately speaking, it does not satisfy the third item. In fact, the intermediate state ($`0<u<\mathrm{}`$) is neither D$`9`$-branes nor D$`p`$-branes but its energy density in the transverse direction has the Gaussian-shaped distribution with a width $`1/u`$. (This is easily seen by noting the tachyon potential has the form $`V(T)e^{u^2x^ix_i}`$.) Therefore, if we allow the width of the source less than the string length, $`u`$ should be sufficiently larger than $`1/\sqrt{\alpha ^{}}`$. We will evaluate it more precisely below.
This system is again described by the off-shell boundary state of the D9D̄9 system with the tachyon profile (3.20) turned on . It is given by in the NSNS-sector,
$$|Bp^{};u_{\mathrm{NS}}=\frac{T_9}{2}[d𝐗^M]\mathrm{Tr}\widehat{\mathrm{P}}e^{{\scriptscriptstyle 𝑑\widehat{\sigma }𝐌(\widehat{\sigma })}}|𝐗^M_{\mathrm{NS}},$$
(3.21)
with
$$𝐌=\left(\begin{array}{cc}0& T(𝐗)\\ T^{}(𝐗)& 0\end{array}\right),$$
(3.22)
where $`T_9`$ is the tension of a single D$`9`$-brane, $`T(x)`$ is given by (3.20) and the trace is taken over the Chan-Paton space of size $`2^{\frac{9p}{2}}N`$. It represents the Neumann boundary state with arbitrary numbers of tachyon vertex operators attached. By evaluating the path integral and rewriting it with oscillators, we obtain (see Appendix for more details)
$`|Bp^{};u_{\mathrm{NS}}=`$ $`N{\displaystyle \frac{T_p}{2}}e^{\frac{1}{4u^2}{\scriptscriptstyle 𝐏_iD𝐏_i}}|𝐏_\mu =0,𝐗^i=0_{\mathrm{NS}}`$ (3.23)
$`=`$ $`N{\displaystyle \frac{T_p}{2}}F(y)^{9p}e^{_{n=1}^{\mathrm{}}\frac{1}{n}\alpha _n^MS_{MN}^{(n)}\stackrel{~}{\alpha }_n^N+i_{r>0}b_r^MS_{MN}^{(r)}\stackrel{~}{b}_r^N}{\displaystyle 𝑑p_ie^{\frac{p_i^2}{8\pi u^2}}|p_\mu =0,p_i_{\mathrm{NS}}},`$ (3.24)
where $`y4\pi \alpha ^{}u^2`$ is a dimensionless parameter, and the function $`F(y)`$ is given by
$$F(y)\frac{4^yy^{1/2}\mathrm{\Gamma }^2(y)}{\sqrt{4\pi }\mathrm{\Gamma }(2y)},$$
(3.25)
and
$$S_{MN}^{(n)}=(\eta _{\mu \nu },\frac{yn}{y+n}\delta _{ij}),S_{MN}^{(r)}=(\eta _{\mu \nu },\frac{yr}{y+r}\delta _{ij}).$$
(3.26)
The corresponding state for the RR-sector is similarly obtained.
It is instructive to see the behavior of this state under the change of the value $`u`$ $`(0<u<\mathrm{})`$. First of all, we can identify the state (3.24) with the deformed state (3.5) (apart from the zero-mode part) as
$$A^{(n,r)}=1,B^{(n)}=\frac{yn}{y+n},B^{(r)}=\frac{yr}{y+r},C=F(y)^{9p}.$$
(3.27)
$`A^{(n,r)}`$ is independent of $`y`$, which means that the tangential directions always satisfy the Neumann boundary condition. On the other hand, $`B^{(n,r)}`$ is a monotonically increasing function of $`y`$, starting from $`1`$ (Neumann) at $`y=0`$ and goes to $`1`$ (Dirichlet) at $`y\mathrm{}`$. It means that the transverse directions satisfy a mixed boundary condition of Dirichlet and Neumann. Roughly speaking, the boundary condition is Dirichlet-like for $`B^{(n,r)}<0`$ and Neumann-like for $`B^{(n,r)}>0`$. Probing only by the massless fields, this state has similar feature as the BPS D$`p`$-branes at least in the range sufficiently larger than $`y=1/2`$, although other higher modes can be Neumann-like in this region. $`C`$ behaves as a monotonically decreasing function of $`y`$ with $`C\mathrm{}`$ for $`y0`$ and $`C1`$ for $`y\mathrm{}`$. (We give the profile of the function $`F(y)`$ in Fig.1.)
Therefore the effective tension is always greater than that of the BPS D$`p`$-brane. Although this factor itself is not well-behaved at $`y0`$, combined with the contribution from the zero-mode, we recover the correct tension for the D$`9`$-brane (see Appendix). The zero-mode part depends directly on $`u`$ and means the Gaussian-shaped distribution in the transverse with a $`(\mathrm{width})^2=1/8\pi u^2`$. Strictly speaking, this is excluded from our assumption for the deformed state. However, since we only concern the massless emission, it is enough to demand the delta-function source at the low energy, that is, the value of $`u`$ such that $`1/8\pi u^2\alpha ^{}`$ is in fact well-localized in the limit $`\alpha ^{}0`$. This condition is equivalent to $`y1/2`$, which is the same region discussed previously. When the width is the order of the string scale, then the supergravity approximation is no longer valid and the $`\alpha ^{}`$-correction in the equation of motion and also the source should be taken into account. In this case, the zero-mode factor in (3.24) is nothing but the $`\alpha ^{}`$-corrected $`\delta `$-function source considered in in the context of the smearing of the singularity.
Thus, at least in the region $`y1/2`$, we can identify this microscopic state with the macroscopic quantities as
$$m=F(y)^{9p}\frac{8y+p3}{4(2y+1)},d=F(y)^{9p}\frac{1}{2y+1}.$$
(3.28)
The behavior of $`m`$ and that of $`d`$ as function of $`y`$ are shown in Fig.2 and Fig.3, respectively.
Then, we see that they represent a trajectory in the $`(m,d)`$ space and belong to the branch $`I_+`$: $`m1`$, $`d0`$. This gives two-parameter subset of the three-parameter solution, which possesses an non-extremal dilaton charge.
As a final remark, we note that the deformation here is also regarded as a massive excitation of open strings on the BPS D$`p`$-branes, since a tachyonic excitation on the top of the potential is equivalent to the massive excitation from the viewpoint of the bottom (See Fig.4.).
In fact, the original expression (3.22) represents the deformation of the D9D̄9 system by a tachyonic excitation, while in (3.23) the same system is described by the D$`p`$-branes with massive vertex operators $`𝑑\widehat{\sigma }𝐏_iD𝐏_i`$ turned on.
#### 3.2.3 Case 3: Unstable D-instantons to BPS D$`p`$-branes
The third example is the tachyon condensation from a system of unstable D-instantons to the BPS D$`p`$-branes. In this case, the BPS D$`p`$-branes can be constructed as bound states of infinitely many unstable D-instantons<sup>10</sup><sup>10</sup>10For details we refer the literatures discussed that issue from the viewpoint of the worldvolume theory and from the viewpoint of the boundary state . . For the concreteness, we consider a system of non-BPS D-instantons in Type IIA string theory. There are ten scalar fields $`\mathrm{\Phi }^M`$ and a tachyon field $`T`$ on them. They are regarded as self-adjoint operators acting on the infinite dimensional Hilbert space, since the worldvolume is zero-dimensional and the matrix-size is infinite. Then the configuration representing $`N`$ BPS D$`p`$-branes is given by
$$T=v\widehat{p}_\mu \gamma ^\mu 1_N,$$
(3.29)
$$\mathrm{\Phi }^\mu =\widehat{x}^\mu 11_N,\mathrm{\Phi }^i=0,$$
(3.30)
where $`\widehat{x}^\mu `$ and $`\widehat{p}_\mu `$ are operators on a Hilbert space satisfying the canonical commutation relation,
$$[\widehat{x}^\mu ,\widehat{p}_\nu ]=i\delta _\nu ^\mu ,$$
(3.31)
$`\mu =0,1,\mathrm{},p`$ are the worldvolume direction and $`i=p+1,\mathrm{},9`$ are the transverse direction and $`\gamma ^\mu `$ are Hermitian gamma matrices. $`v`$ is a real parameter with dimension of length and this configuration becomes an exact solution in the limit $`v\mathrm{}`$.
The intermediate state $`(0<v<\mathrm{})`$ can be interpreted as a $`(p+1)`$-dimensional object, which has a fuzzy worldvolume in some sense. This is most easily understood as follows. Recall that the each set of eigenvalues of the scalar fields $`\{\mathrm{\Phi }^M\}`$ represents the position of each individual D-instanton. In the absence of the tachyon profile $`(v=0)`$, they are distributed uniformly on the $`(p+1)`$-dimensional plane, since the spectrum of $`\widehat{x}^\mu `$ spans the real axis. Note that they are still the collection of the D-instantons. However, if we turn on the tachyon profile $`(v>0)`$, D-instantons become correlated among them. As seen from the tachyon potential $`V(T)\mathrm{Tr}e^{v^2\widehat{p}_\alpha \widehat{p}^\alpha }`$, the momentum distribution is localized around the origin of the momentum space with a width $`\mathrm{\Delta }p1/v`$. This means (in an appropriate way) that the position of each D-instanton becomes uncertain with an amount of $`\mathrm{\Delta }xv`$. Then, in the limit $`v\mathrm{}`$, it becomes $`\mathrm{\Delta }x=\mathrm{}`$ and we cannot observe the individual D-instantons and the worldvolume of D$`p`$-brane appears.
The off-shell boundary state describing this system is given by the Dirichlet boundary state with the scalar and tachyon fields turned on as the boundary interaction. In the NSNS-sector, after calculating the trace over the Chan-Paton Hilbert space, we have
$`|Bp^{};v_{\mathrm{NS}}=`$ $`N{\displaystyle \frac{T_p}{2}}{\displaystyle \left[d𝐗^\mu \right]e^{\frac{1}{4v^2}{\scriptscriptstyle D𝐗^\mu D^2𝐗^\mu }}|𝐗^\mu ,𝐗^i=0_{\mathrm{NS}}},`$ (3.32)
$`=`$ $`N{\displaystyle \frac{T_p}{2}}F(y)^{p+1}e^{_{n=1}^{\mathrm{}}\frac{1}{n}\alpha _n^MS_{MN}^{(n)}\stackrel{~}{\alpha }_n^N+i_{r>0}b_r^MS_{MN}^{(r)}\stackrel{~}{b}_r^N}|p_\mu =0,x^i=0_{\mathrm{NS}},`$ (3.33)
where $`y`$ is a dimensionless parameter $`yv^2/\pi \alpha ^{}`$, $`F(y)`$ is a function,
$$F(y)\frac{4^yy^{1/2}\mathrm{\Gamma }^2(y)}{\sqrt{4\pi }\mathrm{\Gamma }(2y)},$$
(3.34)
and
$$S_{MN}^{(m)}=(\frac{ym}{y+m}\eta _{\mu \nu },\delta _{ij}),S_{MN}^{(r)}=(\frac{yr}{y+r}\eta _{\mu \nu },\delta _{ij}),$$
(3.35)
which is very similar to (3.24). Then we can identify this state with the deformed state (3.5) with setting the parameters as
$$A^{(n)}=\frac{yn}{y+n},A^{(r)}=\frac{yr}{y+r},B^{(n,r)}=1,C=F(y)^{p+1}.$$
(3.36)
Note that the zero-mode part is exactly the delta-function in this case. The boundary condition in the worldvolume direction is now deformed to $`A^{(n,r)}`$, which connects $`A^{(n,r)}=1`$ (Dirichlet) for $`y0`$ and $`A^{(n,r)}=1`$ (Neumann) for $`y\mathrm{}`$. This is understood more easily from the viewpoint of the deformation from the $`N`$ BPS D$`p`$-branes, that is, the expression (3.32), on which a massive vertex operator $`\frac{1}{4v^2}D𝐗^\mu D^2𝐗_\mu `$ turned <sup>11</sup><sup>11</sup>11 See also recent discussions on the relevance of vertex operator here and fuzziness of the worldvolume . . Since such a vertex operator makes the end point of the open string massive, the freely moving endpoint at $`y\mathrm{}`$ becomes heavier and heavier as $`y`$ decreases, and finally it is completely frozen in the limit of $`y0`$. The behavior of $`C`$ as a function of $`y`$ is the same as the previous example. Here the divergence at $`y=0`$ is reflected simply by the original system has infinite number of non-BPS D-instantons.
Comparing the massless emission from this state and the long range behavior of the three-parameter solution, we can relate the macroscopic to the microscopic quantities as
$$m=F(y)^{p+1}\frac{8yp+3}{4(2y+1)},d=F(y)^{p+1}\frac{1}{2y+1},$$
(3.37)
and the profiles of these parameters are depicted in Fig.5 and Fig.6.
In this case, one might conclude that the description of the supergravity is valid for all region of $`y`$ since the configuration of D-instantons does not seem to break the ansatzes (3.1). However, as shown in Fig. 5, $`m`$ is smaller than $`1`$ if $`y`$ is too small, which is out of the parameter region of the three-parameter solution. This means that this system is consistent with the three-parameter solution only when $`y`$ is sufficiently large. This would be understood as follows. As mentioned above, when $`v`$ $`(y)`$ is small, the boundary of open strings on the “D-brane” is localized in the tangential directions. From the viewpoint of closed strings, this means that the size of loops that construct the (deformed) boundary state is localized when $`y`$ is small, that is, the translational invariance is spontaneously broken. This is also related to the uncertainty of the position of D-instantons in the tangential directions. In order to recover the translational invariance of closed string on the brane, the uncertainty must be far larger than the string scale, $`\mathrm{\Delta }xv\sqrt{\alpha ^{}}`$, that is, $`y1`$ where $`m`$ is larger than $`1`$. Since the other non-extremal parameter $`d`$ is always smaller than 0 (see Fig.6), we conclude that the tachyon condensation from unstable instantons to BPS saturated D$`p`$-branes is expressed by a three-parameter solution in the branch $`I_{}`$: $`m1,d0`$.
#### 3.2.4 Tachyonic/massive excitation on the D$`p`$$`p`$ system
We have illustrated three types of deformation above. Here we would like to capture the general feature from above examples and discuss a possible generalization.
First, we summarize the three cases. We considered the situation where the RR-charge is fixed as the same value as $`N`$ BPS D$`p`$-branes have. In Case 1, the D$`p`$$`p`$ system with an tachyonic excitation gives the three-parameter solution with $`d=0`$ at low energy, which is trivial in some sense. In Case 2, we considered the tachyon condensation on the D$`9`$$`9`$ system, which gives the solutions with $`0d`$. It is equivalent to the massive excitation on the BPS D$`p`$-branes. On the other hand, in Case 3, although the tachyon condensation on the unstable D-instanton system is also equivalent to the massive excitation on the BPS D$`p`$-branes, their low energy solutions have $`d0`$. The difference between Case 2 and Case 3 is the polarization of the massive vertex operators: The bosonic part of the vertex operator $`𝑑\widehat{\sigma }𝐏_iD𝐏^i`$ is $`\delta _{ij}𝑑\sigma _\tau X^i_\tau X^j`$, which describes the lowest massive mode of open strings, a rank $`2`$ tensor with zero momentum. On the other hand, the vertex operator $`𝑑\widehat{\sigma }D𝐗^\mu D^2𝐗_\mu =\eta _{\mu \nu }𝑑\sigma _\sigma X^\mu _\sigma X^\nu +\mathrm{}`$ has a longitudinal polarization. Physically, the effects of transverse (longitudinal) polarization is characterized by the fuzziness of the worldvolume in the transverse (longitudinal) direction, respectively. In the context of the scattering of massive open strings, their difference is argued for example in . We here observe that the difference in the polarization is seen as the sign of $`d`$ if we probe these systems by massless closed strings.
Next, we discuss the possible generalization of the process of the tachyon condensation. In Case 2 and 3, we have only considered the tachyon condensation to BPS D$`p`$-branes. However, we can also construct a D$`p`$$`p`$ system as long as they have the same RR-charge proportional to $`N`$, Therefore, recalling that a tachyon excitation on a system of unstable D9-branes or unstable D-instantons is regarded as a massive excitation on the resulting D$`p`$-branes (see Fig. 4), we can combine the three cases as the D$`p`$$`p`$ system with tachyonic and massive excitations. As we have seen above, the tachyonic excitation contributes to $`m`$ and the massive excitations contribute to both $`m`$ and $`d`$. Note that $`1m`$ in any case.
It is then quite natural to take into account for higher massive excitations. Let us consider vertex operators quadratic in $`X`$, say,
$$𝑑\sigma _\tau ^lX^i_\tau ^lX_i,\mathrm{or}𝑑\sigma _\sigma ^lX^\mu _\sigma ^lX_\mu ,$$
(3.38)
with $`l=1,2,\mathrm{}`$. It is easy to show that such an excitation gives rise to the deformation of the boundary state in the form (3.5). It follows that they also contribute to the three-parameter solution and gives another trajectory in the $`(m,d)`$ space. Note that the relation between the polarization and the sign of $`d`$ is unchanged. We can also consider vertex operators with higher power in $`X`$, unless they break the global symmetry. In this case, the resulting deformed state no longer has the form (3.5) but treating these vertex operators as perturbation, they contribute to the coupling to massless modes of closed strings. In any case, our conclusion is that the D$`p`$$`p`$ system with tachyonic and massive excitations are seen as the three-parameter solution at the low energy.
## 4 Conclusions and Discussions
In this paper, we discussed the stringy origin of the general solution of Type II supergravity with the symmetry $`ISO(1,p)\times SO(9p)`$, which is called as the three-parameter solution. This solution contains the BPS saturated black $`p`$-brane solution in the parameter space, whose source is BPS saturated D$`p`$-branes expressed by a boundary state. We characterized the three-parameter solution in terms of two non-extremality parameter $`m`$ and $`d`$ on the supergravity side. On the other hand, we discussed a class of the deformation of the boundary state on the string side. Then we determined the relation between the (macroscopic) non-extremality parameters of the classical solution and the (microscopic) deformation parameters by extending the correspondence between the BPS black $`p`$-brane solution and the boundary state. In particular, we showed that the dilaton charge is related to the deformation of the boundary condition.
We gave three examples of deformed boundary states by considering the tachyon condensation. The first example is a D$`p`$$`p`$-brane system with a constant tachyon VEV discussed in , the second and the third example are the tachyon condensation processes from the unstable D9-branes and the unstable D-instantons to the BPS D$`p`$-branes, respectively. In the latter two examples, the boundary condition in the longitudinal and the transverse directions are deformed, respectively, then the corresponding classical solution learns to possess a non-trivial dilaton charge. We also showed that the deformed systems are generally regarded as tachyonic and/or massive excitations of the open strings on D$`p`$$`p`$-brane systems.
Our method is also applicable to the charge-neutral case and/or more complicated (less symmetric) systems. For example, the intersecting D-brane system is the one with several RR-charges and less global symmetry . Another possible application is the study of the relationship between the stability of the supergravity solution and that of the D-brane system. In this paper, we only consider a static solution, thus the source is also static even if tachyon fields are excited. However, if we consider a perturbation from the solution, the unstable modes are expected to correspond to tachyonic excitations on the D-branes. This would give the stringy meaning of the instability of the supergravity solutions.
It is also interesting to investigate the properties of the dilaton charge further. As we repeated in this paper, the existence of the dilaton charge changes the structures of spacetimes, in particular, it generally makes the spacetimes have no horizons. Therefore, the study of the dilaton charge from the viewpoint of string theory might lead us to the understanding of the meaning of the horizons from the viewpoint of the string theory. However, since the three-parameter solution does not have the horizon in most region of the parameter space, we will have to treat the four-parameter solution which has the horizon in some parameter region , in order to play with the most interesting feature of black objects. For this purpose, our strategy is expected to be essentially applicable. If we can clarify this issue, it might be possible, for example, to understand the Schwarzschild black hole in terms of the string theory, which would become one of the points of contact for the string theory and the general relativity.
## Acknowledgments
The authors would like to thank T. Harada, M. Hayakawa, Y. Himemoto, D. Ida, Y. Ishimoto, G. Kang, H. Kawai, S. Kinoshita, H. Kudoh, Y. Kurita, K. Maeda, Sh. Matsuura, S. Mukohyama, K. Nakao, M. Natsuume, K. Ohta, N. Ohta, T. Onogi, M. Sakagami, N. Sasakura, J. Soda, K. Takahashi, T. Tada, T. Tamaki, S. Watamura, B. de Wit and J. Yokoyama for great supports and helpful discussions. This work of T.A. and S.M. is supported by Special Postdoctoral Researchers Program at RIKEN.
## Appendix A Construction of Boundary States
### A.1 Construction from D9-branes
In this appendix, we review the tachyon condensation from a D9D̄9 system to $`N`$ BPS D$`p`$-brane in Type IIB superstring theory .
We start with a off-shell boundary state corresponding to D9D̄9-brane system in the NSNS-sector on which the tachyon field is excited;
$$|Bp^{}_{\mathrm{NS}}e^{S_b(T)}|B9_{\mathrm{NS}},$$
(A.1)
where $`|B9_{\mathrm{NS}}`$ is the boundary state of a single D9-brane in the NSNS-sector (3.3) and
$$e^{S_b}=\mathrm{Tr}\widehat{\mathrm{P}}e^{{\scriptscriptstyle 𝑑\widehat{\sigma }𝐌(\widehat{\sigma })}},𝐌=\left(\begin{array}{cc}0& T(𝐗)\\ T^{}(𝐗)& 0\end{array}\right),$$
(A.2)
is a boundary interaction and $`\widehat{\mathrm{P}}`$ denotes the supersymmetric path ordered product. $`𝐗^M(\widehat{\sigma })`$ and $`𝐏^M(\widehat{\sigma })`$ denote the position boundary superfields and the conjugate momentum superfields on the boundary, respectively, and $`\widehat{\sigma }=(\sigma ,\theta )`$ is the boundary supercoordinate. For notation of the superfields and the supercoordinate, see . For construction of the boundary state, see . When $`𝐌`$ can be expanded by $`SO(2m)`$ gamma matrices $`\mathrm{\Gamma }^I`$ $`(I=1,\mathrm{},2m)`$ as
$$𝐌=\underset{k=0}{\overset{2m}{}}𝐌^{I_1\mathrm{}I_k}\mathrm{\Gamma }^{I_1\mathrm{}I_k},$$
(A.3)
it is convenient to rewrite the boundary interaction using fermionic superfields as <sup>12</sup><sup>12</sup>12 The $`\mathrm{TrP}`$ in the second line is necessary when $`𝐌_{I_1\mathrm{}I_k}`$ are also matrices.
$`e^{S_b}`$ $`={\displaystyle \left[d𝚪^I\right]\mathrm{TrP}\mathrm{exp}\left\{𝑑\widehat{\sigma }\left(\frac{1}{4}𝚪^ID𝚪^I+\underset{k=0}{\overset{2m}{}}𝐌_{I_1\mathrm{}I_k}𝚪^{I_1}\mathrm{}𝚪^{I_k}\right)\right\}}.`$ (A.4)
We fix the measure of the path integral so that the boundary interaction (A.4) gives the number of D9-branes in the absence of the tachyon field.
Suppose $`N2^{(9p)/2}`$-pairs of D9-brane and D̄9-brane and consider the tachyon profile,
$$𝐌=u\mathrm{\Gamma }^i𝐗^i1_N,$$
(A.5)
where $`\mathrm{\Gamma }^i`$ $`(i=p+1,\mathrm{},9)`$ are the $`SO(2(9p))`$ $`\gamma `$-matrices. Then (A.1) becomes
$`|Bp;u_{\mathrm{NS}}`$ $`={\displaystyle [d𝐗^Md𝚪^i]\mathrm{Tr}P\mathrm{exp}\left[𝑑\widehat{\sigma }\left(\frac{1}{4}𝚪^iD𝚪^i+u𝐗^i𝚪^i\right)\right]|𝐗^M_{\mathrm{NS}}},`$ (A.6)
where the measure $`[d𝐗^M]`$ is determined so that this state becomes the boundary state of $`N2^{(9p)/2}`$-pairs of D9-brane and D̄9-brane in the limit of $`u0`$. Since $`|𝐗=e^{{\scriptscriptstyle 𝑑\widehat{\sigma }i\mathrm{𝐗𝐏}}}|𝐗=0`$ using the conjugate momentum superfield $`𝐏`$, the boundary fermion fields $`𝚪^i`$ are replaced by the momentum superfields $`𝐏^i`$ by carrying out the functional integral for $`𝐗^i`$ and $`𝚪^i`$. Moreover, it is easy to check that this state imposes the Neumann boundary condition for the directions $`\mu =0,\mathrm{},p`$. Then (A.1) can be written as
$`|Bp^{};u_{\mathrm{NS}}`$ $`=N\mathrm{exp}\left[{\displaystyle 𝑑\widehat{\sigma }\left(\frac{1}{4u^2}𝐏_iD𝐏_i\right)}\right]|Bp_{\mathrm{NS}},`$ (A.7)
where we determine the constant factor of (A.1) so that $`|Bp^{};u_{\mathrm{NS}}`$ becomes the boundary state corresponding to the NSNS-sector of $`N`$ D$`p`$-branes in the limit of $`u\mathrm{}`$.
Using the explicit expression of $`𝐏_i(\widehat{\sigma })`$ by the string oscillators and the zeta-function regularization,
$`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{b}{m+a}}={\displaystyle \frac{\mathrm{\Gamma }(a+1)}{\sqrt{2\pi b}}},{\displaystyle \underset{r=1/2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{b}{r+a}}={\displaystyle \frac{\mathrm{\Gamma }(a+1/2)}{\sqrt{2\pi }}},`$ (A.8)
we obtain
$`|Bp^{};u_{\mathrm{NS}}=`$ $`N{\displaystyle \frac{T_p}{2}}\left[{\displaystyle \frac{4^yy^{1/2}\mathrm{\Gamma }^2(y)}{\sqrt{4\pi }\mathrm{\Gamma }(2y)}}\right]^{9p}\mathrm{exp}[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\alpha _n^MS_{MN}^{(n)}\stackrel{~}{\alpha }_n^N+i{\displaystyle \underset{r=1/2}{\overset{\mathrm{}}{}}}b_r^MS_{MN}^{(r)}\stackrel{~}{b}_r^N]\times `$
$`\times {\displaystyle }{\displaystyle \frac{_{i=p+1}^9dp_i}{(2\pi )^{9p}}}\mathrm{exp}({\displaystyle \frac{p_i^2}{8\pi u^2}})|p_\mu =0,p_i_{\mathrm{NS}},(y=4\pi \alpha ^{}u^2)`$ (A.9)
with
$$S_{MN}^{(n)}=(\eta _{\mu \nu },\frac{yn}{y+n}\delta _{ij}),S_{MN}^{(r)}=(\eta _{\mu \nu },\frac{yr}{y+r}\delta _{ij}).$$
(A.10)
We can easily show that (A.9) correctly becomes the boundary state of the NSNS-sector of the $`N`$ D$`p`$-branes in the limit of $`u\mathrm{}`$, which is apparent from the original definition (A.7). On the other hand, in the limit of $`u0`$, (A.9) becomes
$$|Bp^{};u0=2^{(9p)/2}N\frac{(4\pi ^2\alpha ^{})^{(9p)/2}T_p}{2}e^{_{n=1}^{\mathrm{}}\frac{1}{n}\alpha _n^M\stackrel{~}{\alpha }_n^M+i_{r=1/2}^{\mathrm{}}b_r^M\stackrel{~}{b}_r^M}|p_M=0_{\mathrm{NS}}.$$
(A.11)
Since $`(4\pi ^2\alpha ^{})^{(9p)/2}T_p=T_9`$, (A.9) correctly expresses the tachyon condensation from $`N2^{(9p)/2}`$-pairs of D9-brane and D̄9-brane to $`N`$ D$`p`$-branes.
For the RR sector, we can carry out the similar calculation and a deformed boundary state like (A.9) appears, that is, the boundary condition and the zero-mode part is deformed. However, in the case of the RR-sector, the normalization factor of the state does not depend on $`u`$ since the contribution from the bosonic oscillators and the fermionic oscillators completely cancel. This leads to the conservation of the RR charge under the tachyon condensation.
### A.2 Construction from D-instantons
Next, we review tachyon condensation of infinitely many non-BPS D-instantons to the $`N`$ D$`p`$-branes. For detail, see .
In Type IIA superstring theory, a state corresponding to non-BPS D-instantons with scalar profiles and a tachyon profile is given by
$$|Bp^{}=e^{S_b(\mathrm{\Phi }^\mu ,T)}|B(1)_{\mathrm{NS}},$$
(A.12)
where
$`e^{S_b}=\mathrm{Tr}\widehat{\mathrm{P}}e^{{\scriptscriptstyle 𝑑\widehat{\sigma }𝐌(\widehat{\sigma })}},𝐌=\left(\begin{array}{cc}i\mathrm{\Phi }^M𝐏_M& T\\ T& i\mathrm{\Phi }^M𝐏_M\end{array}\right),`$ (A.15)
is a boundary interaction.
The solution that corresponds to $`N`$ D$`p`$-branes is
$`T=v\widehat{p}_\mu 1_N\gamma ^\mu ,\mathrm{\Phi }^\mu =\widehat{x}^\mu 1_N1,\mathrm{\Phi }^i=0,`$ (A.16)
where $`\gamma ^\mu `$ are gamma matrices, and $`\widehat{x}^\mu `$ and $`\widehat{p}_\mu `$ are operators satisfying
$$[\widehat{x}^\mu ,\widehat{p}_\nu ]=i\delta _\nu ^\mu .$$
(A.17)
In the configuration (A.16), the matrix $`𝐌`$ is decomposed as
$$𝐌=\left(i\widehat{x}^\mu 𝐏_\mu \right)1+\left(v\widehat{p}_\mu \right)\mathrm{\Gamma }^\mu ,\mathrm{\Gamma }^\mu =\left(\begin{array}{cc}0& \gamma ^\mu \\ \gamma ^\mu & 0\end{array}\right)$$
(A.18)
then it is again convenient to use the boundary fermion as (A.4);
$`e^{S_b}`$ $`={\displaystyle \left[d𝚪^\mu \right]\mathrm{Tr}\widehat{\mathrm{P}}\mathrm{exp}\left\{𝑑\widehat{\sigma }\left(\frac{1}{4}𝚪^\mu D𝚪^\mu i\widehat{x}^\mu 𝐏_\mu +\left(v\widehat{p}_\mu \right)𝚪^\mu \right)\right\}}`$
$`=N{\displaystyle \left[d𝚪^\mu d𝐱^\mu d𝐩_\mu \right]\mathrm{exp}\left\{𝑑\widehat{\sigma }\left(\frac{1}{4}𝚪^\mu D𝚪^\mu +i𝐩_\mu D𝐱^\mu i𝐱^\mu 𝐏_\mu +v𝐩_\mu 𝚪^\mu \right)\right\}}`$
$`=N{\displaystyle \left[d𝚪^\mu d𝐱^\mu \right]\delta \left(D𝐱^\mu iv𝚪^\mu \right)\mathrm{exp}\left\{𝑑\widehat{\sigma }\left(\frac{1}{4}𝚪^\mu D𝚪^\mu i𝐱^\mu 𝐏_\mu \right)\right\}}`$
$`=N{\displaystyle \left[d𝐗^\mu \right]\mathrm{exp}\left\{𝑑\widehat{\sigma }\left(i𝐗^\mu 𝐏_\mu \frac{1}{4v^2}D𝐗^\mu D^2𝐗^\mu \right)\right\}}.`$ (A.19)
Here we have replaced the operators $`\widehat{x}^\mu `$ and $`\widehat{p}_\mu `$ by superfields $`𝐱^\mu (\widehat{\sigma })`$ and $`𝐩_\mu (\widehat{\sigma })`$ in the second line by adding the kinetic term $`i𝐩_\mu D𝐱^\mu `$. Then we performed the functional integral for $`𝐩_\mu `$ and $`𝚪^\mu `$. In the last line $`𝐱^\mu `$ is identified with the superfield $`𝐗^\mu `$ on the string worldsheet. Then the state corresponding to this solution is
$`|Bp^{};v_{\mathrm{NS}}`$ $`{\displaystyle \left[d𝐗^\mu \right]\mathrm{exp}\left\{𝑑\widehat{\sigma }\left(i𝐗^\mu 𝐏_\mu \frac{1}{4v^2}D𝐗^\mu D^2𝐗^\mu \right)\right\}|𝐗^M=0_{\mathrm{NS}}}`$
$`={\displaystyle \left[d𝐗^\mu \right]\mathrm{exp}\left\{𝑑\widehat{\sigma }\left(\frac{1}{4v^2}D𝐗^\mu D^2𝐗^\mu \right)\right\}|𝐗^\mu ,𝐗^i=0_{\mathrm{NS}}},`$ (A.20)
where the measure has been fixed so that this state expresses the boundary state of $`N`$ D$`p`$-branes in the $`v\mathrm{}`$ limit. This expression can again be evaluated using the zeta-function regularization as
$`|Bp^{};v_{\mathrm{NS}}`$ $`={\displaystyle \frac{NT_p}{2}}\left({\displaystyle \frac{4^yy^{1/2}\mathrm{\Gamma }^2(y)}{\sqrt{4\pi }\mathrm{\Gamma }(2y)}}\right)^{p+1}`$
$`\times \mathrm{exp}\left[{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}\alpha _n^MS_{MN}^{(n)}\stackrel{~}{\alpha }_n^N+i{\displaystyle \underset{r=1/2}{\overset{\mathrm{}}{}}}b_r^MS_{MN}^{(r)}\stackrel{~}{b}_r^N\right]|p_\mu =0,x^i=0_{\mathrm{NS}},`$ (A.21)
where $`y=v^2/\pi \alpha ^{}`$ and
$$S_{MN}^{(n)}=(\frac{yn}{y+n}\eta _{\mu \nu },\delta _{ij}),S_{MN}^{(r)}=(\frac{yr}{y+r}\eta _{\mu \nu },\delta _{ij}),$$
(A.22)
In the limit of $`y\mathrm{}`$,
$$\frac{4^yy^{1/2}\mathrm{\Gamma }^2(y)}{\sqrt{4\pi }\mathrm{\Gamma }(2y)}1$$
(A.23)
as shown in Fig.1.
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# 1 Introduction
## 1 Introduction
The Supersymmetric SO(10) GUT based on the $`\mathrm{𝟏𝟐𝟔}\overline{\mathrm{𝟏𝟐𝟔}}\mathrm{𝟐𝟏𝟎}`$ Higgs multiplets was first introduced over 20 years ago in. At that time the electroweak gauge parameters and fermion mass data were quite incomplete. The next significant development was the realization that R parity $`R_p=()^{3(BL)+2S}`$ formed a part of the gauge symmetry in LR symmetric models with gauged $`U(1)_{BL}`$ and was hence protected by it. In particular it would remain unbroken as long as no field with odd $`BL`$ received a vev. SO(10) is now the most favoured GUT gauge group because of the natural way in which it accommodates complete fermion families together with the superheavy right handed neutrinos required by the seesaw mechanism. In 1992, with LEP data in hand and CKM parameters largely known Babu and Mohapatra took up the task of fitting Fermion masses and mixing using the $`\mathrm{𝟏𝟎}\overline{\mathrm{𝟏𝟐𝟔}}`$ FM Higgs system introduced in and the Type I seesaw mechanism. They proposed that once the charged fermion masses in the SM had been fitted the $`\mathrm{𝟏𝟎}\overline{\mathrm{𝟏𝟐𝟔}}`$ matter fermion Yukawa couplings of the GUT would be completely determined and hence the model would predict the allowed neutrino masses. Although, based on the then current notions about neutrino masses/mixings, they concluded that their proposal failed, it would inspire much future work.
In the mid-1990’s the Supersymmetric LR model based on a parity odd singlet was found to have a charge breaking vacuum and this was construed as evidence of a low scale linked breaking of $`SU(2)_R,BL`$ and R-parity driven by the soft supersymmetry breaking terms. However an alternative analysis based on allowing this scale to lie anywhere in the range from $`M_S`$ and $`M_X`$ allowed the construction of consistent Minimal Left Right Supersymmetric models (MSLRMs) which were shown to naturally preserve R-parity and thus predict a stable LSP. Moreover these analyses brought clearly to the fore a theme that had been noticed from the beginning of the study of multiscale Susy GUTs : that generically in LR Susy GUTs there are light multiplets whose masses are protected by supersymmetry even though they are submultiplets of a GUT multiplet that breaks gauge symmetry at a high scale. As such they violate the conventional wisdom of the “survival principle” which estimates the masses of such submultiplets to be the same as the large vev. Thus it became clear that Susy GUTs required the use of calculated not estimated masses for RG analyses.
This ambling pace of theoretical development was forced by the epochal discovery of neutrino oscillations by Super-Kamiokande in 1997 and the rapid refinement of our knowledge of the parameters thereof. As is well known, the seesaw mass scale indicated by this discovery was in the range indicated by Grand Unification. Since LR models and GUTs containing them naturally incorporate the Seesaw mechanism, this gave a strong motivation for taking up the detailed study of SO(10) Susy Guts anew, particularly with the understanding provided by the natural class of MSLRMs developed earlier. A completely viable R-parity preserving Susy GUT based on the $`\mathrm{𝟏𝟎}\overline{\mathrm{𝟏𝟐𝟔}}`$ FM (for Fermion Mass )Higgs system and an additional $`\mathrm{𝟒𝟓}\mathrm{𝟓𝟒}\mathrm{𝟏𝟐𝟔}`$ AM (for Adjoint type Multiplet and euphony) system was then developed. The RG analysis carried out in this work already indicated that the use of calculated spectra would force together the various possible intermediate scales into a narrow range close to the GUT scale resulting in an effective “SU(5) conspiracy”. Moreover the problem of seeing how the various MSLRMs considered by us would fit into Susy GUTs had motivated a review of the various possibilities and this had again brought up the model based on the 210-plet as an important and interesting possibility.
It thus became clear that a detailed calculation of the full GUT spectrum and couplings in various SO(10) GUT models would be necessary if further progress was to be made and that this would require complete calculational control on the group theory of SO(10) at the level needed by practical field theory calculations which was still lacking in spite of signal contributions. These techniques –based on an explicit decomposition of SO(10) tensor and spinor labels into those of the maximal (“Pati-Salam”) sub-group – were duly if laboriously developed by us. At this time the model based on the $`\mathrm{𝟏𝟐𝟔}\overline{\mathrm{𝟏𝟐𝟔}}\mathrm{𝟐𝟏𝟎}`$ was again brought to the fore by us and its old claim to be called the minimal Susy GUT(MSGUT) was buttressed by an analysis of its parameter counting and the simplicity of its structure.
The techniques developed by us permitted first a partial calculation of the mass matrices of the MSSM doublets $`(1,2,\pm 1)`$ and triplets $`(3,1,\pm \frac{2}{3})`$ and then a complete calculation of all the couplings and mass matrices of the MSGUT. With the same motivations two calculations, one in parallel and cross checking with ours and another quite separate, had commenced after. These calculations both used a somewhat abstract method that permitted the calculation of the “clebsches” that entered the spectra of MSSM submultiplets of SO(10) tensor (but not – so far – spinor) chiral supermultiplets (but not their couplings). After some controversy concerning the consistency of these three calculations a consensus seems to have emerged on their compatibility not withstanding notable differences in normalizations and phase conventions. In we also provided the complete (chiral and gauge) spectra, neutrino mass matrices, all gauge and chiral couplings and the effective $`d=5`$ operators for Baryon violation in terms of GUT parameters. This laid the stage for a completely explicit RG based analysis of the MSGUT : an important example of which was the calculation of threshold effects based on these spectra described in Section 3.
Meanwhile the old theme of utilizing the very restricted structure of fermion Higgs couplings in the class of SO(10) GUTs that used only $`\mathrm{𝟏𝟎}\overline{\mathrm{𝟏𝟐𝟔}}`$ representations and renormalizable couplings to make “predictions” concerning the neutrino mass sector was taken up again. It led to a remarkably simple (2 generation) insight into the operation of the Type II mechanism that naturally generated large atmospheric neutrino mixing angles based on the observed approximate $`b\tau `$ unification in the MSSM extrapolated to $`M_X`$. This simple insight then provoked a detailed analysis of the fitting problem for dominant Type II seesaw mechanism and 3 generations which was quite successful. A certain tension in the details of these fits was ameliorated by later work and disappeared when a 120-plet was also allowed to perturb the fit of the $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟐𝟔}}`$ slightly. Finally, very recently, in fact just in time for PLANCK05 Babu announced that they had succeeded in obtaining viable Type I, Type II and mixed fits to the (large mixing angle) neutrino mass data: which possibility had been neglected since the suggestion of. These sudden reversals are due to the extreme delicacy and complexity of the task of fitting accurately the masses of particles differing by a factor of upto $`10^6`$ in a multidimensional complex parameter space where algorithms “shoot” obliviously past solution points.
Obviously these successful fits of all low energy fermion “quadratic” data (extrapolated up to $`M_X`$ ) in generic MSGUTs cry out to be married to a particular simple MSGUT in which all coefficients are specified and testable for viability. This is particularly so because only in a particular MSGUT, where the needed coefficients are fully specified can one test the viability of a given FM fit (such as the Type II dominant ones). Furthermore, as we shall see in detail below, the true output of the FM fit is really the embedding matrices that specify the relation between MSSM and MSGUT supermultiplets at the GUT scale. Given the phenomenologically expected quasi-identity of flavour mixing in the scalar and fermionic sector down to Electroweak scales, these matrices are the main missing ingredient for performing an informed and realistic calculation of Baryon decay and Lepton Flavour violation in the MSGUT. The combined cuts provided by gauge unification and fermion mass fit viability constrain the MSGUT and render it falsifiable . If it passed these tests with some region of its parameter space unscathed we would be in a position to actually specify the theory at the high scale: a task long thought to be impossible by those despondent at the continuing failure to detect proton decay. The key to this remarkable development is of course the remarkable uroborotic ( $`o\upsilon \rho o\beta o\rho o\varsigma `$ : the world-snake that swallows its tail) feature of the Seesaw which connects the physics of the ultralight and ultramassive particles.
In Section II we review the structure and spontaneous symmetry breaking of the MSGUT and describe its decomposition to Pati-Salam labels. This is accompanied by an Appendix containing the complete spectrum of MSSM multiplets in the MSGUT -including gauge submultiplets. In Section III we describe the RG analysis of threshold effects and the type of cuts these corrections impose on the parameter space. In Section IV we review the fitting of Fermion masses and mixings with a view to supporting our assertion that the true output of this procedure is the specification of the embedding matrices. In Section V we plot the behaviour of the parameters that control Type I vs Type II dominance in the MSGUT and show that Type II dominance may be quite unlikely. We also use particular examples of solutions provided by one of the groups that have performed the successful fits to plot the Type I and Type II contributions and again find indications that Type II is not easily dominant and that the overall size of neutrino masses tends to be too low except possibly in very special regions of parameter space. In the next section we transform our previously derived formula for the effective $`d=5`$ Lagrangian for Baryon decay in the MSGUT to the MSSM using the embedding matrices. This operator is then ready for use to calculate Baryon decay. We also briefly comment on the relevance of our considerations to $`d=6`$ and Lepton Flavour violation. We conclude with a discussion of the outlook and the directions that will be pursued in the forthcoming detailed paper on the application of the techniques indicated here to the pinning down of the MSGUT.
## 2 Essentials of MSGUTs
There are at present two main types of SO(10) Susy GUT models namely renormalizable GUTs (RGUTs) of the classic type which invoke only gauge symmetry and preserve R-parity by maintaining the sharp distinction between Matter and Higgs Supermultiplets, and non renormalizable GUTs (nRGUTs) which allow non renormalizable operators - particularly for generating fermion masses - but impose additional symmetries to perform the functions of R-parity and suppress unwanted behaviours allowed by the license granted by non renormalizability. Practitioners of the art are often sharply fixed in their preference for one type or the other and our own is obviously with the former. Many (talks by Babu,Malinsky,Mohapatra..) of the contributions on GUTs to PLANCK05 are of this type but the other type is also well represented (Pati,Raby…).
Other distinguishing features of these two types of models are that RGUTs allow large representations such as the $`\overline{\mathrm{𝟏𝟐𝟔}}`$ Higgs to generate charged and neutral (via type I and Type II Seesaw) fermion masses but claim minimality on grounds of maximal economy of parameters. In contrast nRGUTs, allow only small Higgs representations and use non renormalizable operators($`\mathrm{𝟏𝟔}_𝐅^\mathrm{𝟐}\overline{\mathrm{𝟏𝟔}}_𝐇^\mathrm{𝟐}`$ etc) to generate effective vevs in the $`\overline{\mathrm{𝟏𝟐𝟔}},\mathrm{𝟏𝟐𝟎}`$ channels. They achieve simplifications by invoking additional symmetries, whose rationale is however not always appreciated by non adherents. A very large number of works on nRGUTs have appeared and some of these have far advanced fitting programs : with claims to generate sample spectra with specific numbers for most (i.e $`10^2`$) low energy MSSM parameters. Moreover they also seek to satisfy additional motivations like workable models of gauge mediated supersymmetry breaking.
The RGUTs, on the other hand are mostly at the stage of seeking a semi quantitative non-disagreement with the data as available at present. For example, since no superpartner has as yet actually been observed, and speculations on their possible masses range all the way from the classic 1 TeV up to almost the GUT scale, the value of the low energy parameters renormalized up to the Susy breaking scale where the Susy RG equations take hold, can, realistically, be considered as known to at best an accuracy of $`510\%`$. Invoking the remarkable accuracy of experimental data at the scale $`M_Z`$ is of little avail since nothing is known of the size of the threshold corrections due to superpartners. It is this approach that we shall adopt in this paper and our quantitative analyses regarding quantities unprotected by any symmetry shall not lay claim to accuracy that may turn out to be quite spurious.
A very important but controversial distinction between the two classes of theories is due to the fact that the one loop beta functions of the $`\overline{\mathrm{𝟏𝟐𝟔}},\mathrm{𝟏𝟐𝟎}`$ chiral supermultiplets are so large (35 and 28 respectively) that they alone overcome the contribution (-24) of the SO(10) gauge supermultiplet and cause the gauge coupling to explode just above $`M_X`$ the perturbative unification scale. By analyzing the complete two loop RG equations for the MSGUT above $`M_X`$ we have shown that the large positive coefficient in the one loop gauge beta function precludes evasion of this difficulty by taking refuge in a weakly coupled fixed point before the gauge coupling explodes. Thus such theories effectively determine their own upper cutoff $`\mathrm{\Lambda }_X10^{17}GeV`$ . This feature is so inescapable that we proposed to face it by making it a signal of a deeper non-perturbative feature of such Susy GUTs: namely that they dynamically break the GUT symmetry down to a smaller symmetry (such as $`H=SU(5)\times U(1)`$ or even $`H=G_{123}`$ ) at the scale $`\mathrm{\Lambda }_X`$ via a Supersymmetric condensation of gauginos in the coset $`SO(10)/H`$ which drives a $`G_{123}`$ preserving Chiral scalar condensate. A new fundamental length $`\mathrm{\Lambda }_X^1`$ below which SM probes cannot delve thus arises, endowing the particles of the SM with “hearts” i.e impenetrable cores with sizes $`10^{30}cm`$. We have implemented the strongly coupled Supersymmetric dynamics required by this type of scenario in a toy model based on $`SU(2)`$ which can be easily generalized to at least the simple breaking $`SO(10)SU(5)\times U(1)`$. Various fascinating perspectives including dual unification and induced gravity then beckon. However in view of the controversial nature of this proposal we shall not treat of it further here. This talk is aimed at indicating the feasibility of a promotion of the MSGUT to a falsifiable theory using the deep linkages between low and high scale physics inherent in the SO(10) seesaw mechanism in combination with the traditional -relatively permissive- constraints of gauge unification.
The MSGUT is the renormalizable globally supersymmetric $`SO(10)`$ GUT whose chiral supermultiplets consist of “AM type” totally antisymmetric tensors : $`\mathrm{𝟐𝟏𝟎}(\mathrm{\Phi }_{ijkl})`$ $`\overline{\mathrm{𝟏𝟐𝟔}}(\overline{𝚺}_{ijklm}),`$ $`\mathrm{𝟏𝟐𝟔}(𝚺_{ijklm})(i,j=1\mathrm{}10)`$ which break the GUT symmetry to the MSSM, together with Fermion mass (FM) Higgs 10-plet($`𝐇_i`$). The $`\overline{\mathrm{𝟏𝟐𝟔}}`$ plays a dual or AM-FM role since it also enables the generation of realistic charged fermion and neutrino masses and mixings (via the Type I and/or Type II mechanisms); three 16-plets $`𝚿_𝐀(A=1,2,3)`$ contain the matter including the three conjugate neutrinos ($`\overline{\nu }_L^A`$).
If in addition to the $`\mathrm{𝟏𝟎},\overline{\mathrm{𝟏𝟐𝟔}}`$ FM Higgs we also include a $`\mathrm{𝟏𝟐𝟎}`$ -plet Higgs allowed by the SO(10) multiplication rule $`\mathrm{𝟏𝟔}\times \mathrm{𝟏𝟔}=\mathrm{𝟏𝟎}+\mathrm{𝟏𝟐𝟎}+\mathrm{𝟏𝟐𝟔}`$, the GUT scale symmetry breaking is unchanged since the $`\mathrm{𝟏𝟐𝟎}`$ contains no SM singlets. However it does contribute two additional pairs of $`(1,2,\pm 1)`$ doublets taking the MSSM Higgs doublet mass matrix $``$ from $`4\times 4`$ to $`6\times 6`$. The resulting theory is thus justly called the next to minimal Susy GUT (nMSGUT).
The superpotential is
$`W`$ $`=`$ $`{\displaystyle \frac{1}{2}}M_HH_i^2+{\displaystyle \frac{m}{4!}}\mathrm{\Phi }_{ijkl}\mathrm{\Phi }_{ijkl}+{\displaystyle \frac{\lambda }{4!}}\mathrm{\Phi }_{ijkl}\mathrm{\Phi }_{klmn}\mathrm{\Phi }_{mnij}+{\displaystyle \frac{M}{5!}}\mathrm{\Sigma }_{ijklm}\overline{\mathrm{\Sigma }}_{ijklm}`$ (1)
$`+`$ $`{\displaystyle \frac{\eta }{4!}}\mathrm{\Phi }_{ijkl}\mathrm{\Sigma }_{ijmno}\overline{\mathrm{\Sigma }}_{klmno}+{\displaystyle \frac{1}{4!}}H_i\mathrm{\Phi }_{jklm}(\gamma \mathrm{\Sigma }_{ijklm}+\overline{\gamma }\overline{\mathrm{\Sigma }}_{ijklm})`$
$`+`$ $`h_{AB}^{}\psi _A^TC_2^{(5)}\gamma _i\psi _BH_i+{\displaystyle \frac{1}{5!}}f_{AB}^{}\psi _A^TC_2^{(5)}\gamma _{i_1}\mathrm{}\gamma _{i_5}\psi _B\overline{\mathrm{\Sigma }}_{i_1\mathrm{}i_5}`$
The parameter counting is as follows : the 7 complex parameters $`m,M,M_H,\lambda ,\eta ,\gamma `$ and $`\overline{\gamma }`$ can be relieved of 4 phases using the arbitrariness in the phases of the 4 SO(10) multiplets $`\mathrm{𝟐𝟏𝟎},\mathrm{𝟏𝟐𝟔},\overline{\mathrm{𝟏𝟐𝟔}},\mathrm{𝟏𝟎}`$ leaving 10 real parameters. The 12 complex Yukawa contained in the symmetric $`3\times 3`$ matrices $`h^{},f^{}`$ can be reduced to 15 real parameters by diagonalizing one linear combination of the two to a real diagonal form. In addition there is the gauge coupling. In all the MSGUT thus has exactly 26 non-soft parameters . Incidentally the MSSM also has 26 non-soft couplings. In we have shown that this number of “hard” parameters is considerably less than any GUT attempting to perform the same tasks vis a vis fermion masses that the MSGUT accomplishes. The ‘levity of the cognoscenti’ provoked by the number 26 will recur again below ! Note however that one of these parameters, e.g $`M_H`$, is traded for the electroweak vev via the fine tuning condition that yields two light doublets and another Susy electro-weak parameter i.e $`\mathrm{tan}\beta `$ must be taken as an additional given of the analysis since it is likely determined by dynamics dependent on the supersymmetry breaking parameters.
The GUT scale vevs that break the gauge symmetry down to the SM symmetry (in the notation $`a,b=1\mathrm{}6;\stackrel{~}{\alpha }=7\mathrm{}10`$ of ) are $`(15,1,1)_{210}:\varphi _{abcd}=\frac{a}{2}ϵ_{abcdef}ϵ_{ef};`$ $`(15,1,3)_{210}:\varphi _{ab\stackrel{~}{\alpha }\stackrel{~}{\beta }}=\omega ϵ_{ab}ϵ_{\stackrel{~}{\alpha }\stackrel{~}{\beta }};(1,1,1)_{210}:\varphi _{\stackrel{~}{\alpha }\stackrel{~}{\beta }\stackrel{~}{\gamma }\stackrel{~}{\delta }}=pϵ_{\stackrel{~}{\alpha }\stackrel{~}{\beta }\stackrel{~}{\gamma }\stackrel{~}{\delta }};`$ $`(10,1,3)_{\overline{126}}:\overline{\mathrm{\Sigma }}_{\widehat{1}\widehat{3}\widehat{5}\widehat{8}\widehat{0}}=\overline{\sigma };(\overline{10},1,3)_{126}:\mathrm{\Sigma }_{\widehat{2}\widehat{4}\widehat{6}\widehat{7}\widehat{9}}=\sigma `$. The vanishing of the D-terms of the SO(10) gauge sector potential imposes only the condition $`|\sigma |=|\overline{\sigma }|`$. Except for the simpler cases corresponding to enhanced unbroken symmetry ($`SU(5)\times U(1),SU(5)`$ $`G_{3,2,2,BL},G_{3,2,R,BL}`$ etc) this system of equations is essentially cubic and can be reduced to the single equation for a variable $`x=\lambda \omega /m`$, in terms of which the vevs $`a,\omega ,p,\sigma ,\overline{\sigma }`$ are specified :
$$8x^315x^2+14x3=\xi (1x)^2$$
(2)
where $`\xi =\frac{\lambda M}{\eta m}`$. This exhibits the crucial importance of the parameter $`\xi `$.
Using the above vevs and the methods of we calculated the complete gauge and chiral multiplet GUT scale spectra and couplings for the 52 different MSSM multiplet sets falling into 26 different MSSM multiplet types of which 18 are unmixed while the other 8 types occur in multiple copies. (On the lighter note : the occurrences yet again of the ‘mystic’ String Theory number 26 demonstrates that one can do just as well without string theory when searching for the fundamental theory !). These spectra may be found in the Appendix.
Among the mass matrices exhibited is the all important $`4\times 4`$ Higgs doublet mass matrix $``$. To keep one pair of these doublets light one tunes $`M_H`$ so that $`Det=0`$. This matrix can then be diagonalized by a bi-unitary transformation yielding thereby the coefficients describing the proportion of the doublet fields in $`\mathrm{𝟏𝟎},\overline{\mathrm{𝟏𝟐𝟔}},\mathrm{𝟏𝟐𝟔},\mathrm{𝟐𝟏𝟎}`$ present in the light doublets : which proportions are important for many phenomena.
## 3 RG Analysis
If the serendipity of the MSSM gauge coupling unification at $`M_X^0`$ is to survive closer examination the MSGUT must answer the query :
Are the one loop values of $`Sin^2\theta _W`$ and $`M_X`$ generically stable against superheavy threshold corrections ?.
We follow the approach of Hall in which the mass of the lightest baryon number violating superheavy gauge bosons is chosen as the common “physical” superheavy matching point ($`M_i=M_X`$) in the equations relating the MSSM couplings to the SO(10) coupling :
$`{\displaystyle \frac{1}{\alpha _i(M_S)}}={\displaystyle \frac{1}{\alpha _G(M_X)}}+8\pi b_iln{\displaystyle \frac{M_X}{M_S}}+4\pi {\displaystyle \underset{j}{}}{\displaystyle \frac{b_{ij}}{b_j}}lnX_j4\pi \lambda _i(M_X)`$ (3)
See for details. In this approach, rather than enforcing unification at a point, it is recognized that above the scale $`M_X`$ the effective theory changes to a Susy $`SO(10)`$ model structured by the complex superheavy spectra which we have computed and which appears as unbroken SO(10) only at the crudest resolution – here surpassed. Thus we compute the corrections to the three parameters $`Log_{10}M_X,sin^2\theta _W(M_S),\alpha _G^1(M_X)`$ as a function of the MSGUT parameters and the answer to the question of stability of the perturbative unification is determined by the ranges of GUT parameters where these corrections are consistent with the known or surmised data on $`Log_{10}M_X,sin^2\theta _W(M_S)`$ and the consistency requirement that the SO(10) theory remain perturbative after correction. We find the corrections
$`\mathrm{\Delta }^{(th)}(Log_{10}M_X)`$ $`=`$ $`.0217+.0167(5\overline{b}_{}^{}{}_{1}{}^{}+3\overline{b}_{}^{}{}_{2}{}^{}8\overline{b}_{}^{}{}_{3}{}^{})Log_{10}{\displaystyle \frac{M^{}}{M_X^0}}`$ (4)
$`\mathrm{\Delta }^{(th)}(sin^2\theta _W(M_S))`$ $`=`$ $`.00004.00024(4\overline{b}_{}^{}{}_{1}{}^{}9.6\overline{b}_{}^{}{}_{2}{}^{}+5.6\overline{b}_{}^{}{}_{3}{}^{})Log_{10}{\displaystyle \frac{M^{}}{M_X^0}}`$ (5)
$`\mathrm{\Delta }^{(th)}(\alpha _G^1(M_X))`$ $`=`$ $`.1565+.01832(5\overline{b}_{}^{}{}_{1}{}^{}+3\overline{b}_{}^{}{}_{2}{}^{}+12\overline{b}_{}^{}{}_{3}{}^{})Log_{10}{\displaystyle \frac{M^{}}{M_X^0}}`$ (6)
Where $`\overline{b}_{}^{}{}_{i}{}^{}=16\pi ^2b_i^{}`$ are 1-loop $`\beta `$ function coefficients ($`\beta _i=b_ig_i^3`$) for multiplets with mass $`M^{}`$ (a sum over representations is implicit).
These corrections are to be added to the one loop values corresponding to the successful gauge unification of the MSSM :
Using the values
$`\alpha _G^0(M_X)^1=25.6;M_X^0=10^{16.25}GeV;M_S=1TeV`$
$`\alpha _1^1(M_S)=57.45;\alpha _2^1(M_S)=30.8;\alpha _3^1(M_S)=11.04`$ (7)
the two loop corrections are
$`\mathrm{\Delta }^{2loop}(log_{10}{\displaystyle \frac{M_X}{M_S}})`$ $`=`$ $`.08;\mathrm{\Delta }^{2loop}(sin^2\theta _W(M_S))=.0026`$
$`\mathrm{\Delta }^{2loop}\alpha _G^1(M_X)`$ $`=`$ $`.546`$ (8)
We see that in comparison with the large threshold effects to be expected in view of the number of heavy fields and their beta functions the 2 loop corrections are quite small.
A few remarks on the role of the parameters are in order. The parameter $`\xi =\lambda M/\eta m`$ is the only numerical parameter that enters into the cubic eqn.(2) that determines the parameter $`x`$ in terms of which all the superheavy vevs are given. It is thus the most crucial determinant of the mass spectrum . The dependence of the threshold corrections on the parameters $`\lambda ,\eta ,\gamma ,\overline{\gamma }`$ is comparatively mild except when coherent e.g when many masses are lowered together leading to $`\alpha _G`$ explosion. The parameter ratio $`m/\lambda `$ can be extracted as the overall scale of the vevs. Since the threshold corrections we calculate are dependent only on (logarithms of) ratios of masses, the parameter $`m`$ does not play any crucial role in our scan of the parameter space : it is simply fixed in terms of the mass $`M_V=M_X`$ of the lightest superheavy vector particles mediating proton decay.
With these formulae in hand we can explore the dependence of the threshold corrections on the “fast” parameter $`\xi `$ in response to which the vevs and thus all the threshold corrections can gyrate wildly (see plots for the real solution for real $`\xi `$ below ), and the “slow diagonal ” parameters $`\lambda ,\eta `$ whose lowering tends to make fields light and thus give large negative corrections to $`\alpha _G^1`$. In addition there are the “slow off diagonal” parameters $`\gamma ,\overline{\gamma }`$ whose effect seems quite mild.
$`Sin^2\theta _W`$ is now very accurately known at $`M_Z`$: $`\widehat{s}_Z^2=.23120\pm .00015`$ or an error of less than $`.065\%`$. Similarly the value $`\alpha _{em}=127.906\pm .0019`$ has an error of only $`.015\%`$. The main uncertainty in the data (besides $`M_H`$ ) at $`M_Z`$ is in $`\alpha _3(M_Z)`$ : $`1.5\%`$. However the same is obviously not true of the overall sparticle mass scale or the intra sparticle mass splittings : with values in current speculation ranging from $`1`$ TeV to $`10^{10}`$ TeV ! Thus at least until the first superpartners are observed the assumption that we know the values of the $`\alpha _i(M_S)`$ with anything like the precision at $`M_Z`$ is quite unjustified. A rough estimate of the uncertainty in $`Sin^2\theta _W(M_S)`$ gives numbers in the ball park of 1% to 10% and we shall not pretend to have any license to impose stronger constraints on our parameters. For the uncertainty in $`Log_{10}M_X`$ however there is a quite stringent $`M_S`$ independent constraint coming from the requirement that gauge mediated proton decay in channels like $`pe^+\pi ^0`$ obey the current bounds $`\tau _{pe^+\pi ^0}>10^{33}yr`$. Thus $`\mathrm{\Delta }Log_{10}M_X>1`$ is a very firm constraint that we may rely upon. Finally it is clear that the applicability of the perturbative formulae used in our treatment will need detailed scrutiny if the fractional change in $`\alpha _G(M_X)^1`$ is greater than about 25%(particularly if the change is negative). So $`|\mathrm{\Delta }\alpha _G(M_X)^1|<10`$(say) is a limitation we tentatively observe as a limit on the range of validity of our calculation ( rather than the theory itself : since a detailed examination may permit one to handle larger changes – particularly positive ones – in $`\alpha _G(M_X)^1`$). In view of these relatively loose bounds on the acceptable changes associated with threshold effects we cannot expect that gauge constraints can limit the allowed parameter space very precisely. Nevertheless, these constraints used in conjunction with additional very significant information arising from the fitting of Fermion masses and (in more model dependent ways) $`d=5`$ Baryon violation and Lepton flavour violation could allow us to obtain a complete semi-quantitative picture of viable regions, if any, of the MSGUT parameter space.
We now present examples showing how such a semi-quantitative mapping of the RG constrained topography of the GUT scale parameter space may be obtained using the formulae given above. An attempt at an exhaustive mapping would be premature till the recent available FM fits have been digested. Later we shall present our program of further constraining the parameter space using the constraints of the fit of Fermion masses and also due to the consequent specification of Baryon decay operators. It should be kept in mind that at any $`\xi `$ the cubic equation (2 ) has three solutions any of which is in principle exploitable for defining a vacuum.
Consider first the plots of the threshold changes $`\mathrm{\Delta }^{(th)}(\alpha _G^1(M_X))`$ and $`\mathrm{\Delta }^{(th)}(Log_{10}M_X)`$ vs $`\xi `$ after (arbitrarily) setting $`\{\lambda ,\eta ,\gamma ,\overline{\gamma }\}=\{.7,.5,.3,.2\}`$ which are shown as Fig. 1 for the real solution of the eqn(2) for real $`\xi `$.
The “twin towers” due to singularities at $`\xi =5,10`$ in the plot of $`\alpha _G(M_X)^1`$ arise from the $`SU(5)`$ invariant real vevs at these values of $`\xi `$. Similarly the negative spike at $`\xi =2/3`$ corresponds to a real solution with $`SU(5)_{flipped}\times U(1)`$ symmetry. As emphasized by us in the plots of the threshold corrections show very sharply defined features corresponding to the special behaviour near points of enhanced symmetry and thus provide a rapid way of scanning the topography of the parameter space. For these moderately large values of $`\lambda ,\eta `$ the lower cut at $`\mathrm{\Delta }\alpha _G(M_X)^1=10`$ essentially allows all possible values of $`\xi `$ though the upper cut does advise caution around the special symmetry points with $`SU(5)`$ symmetry. However the graph of $`\mathrm{\Delta }Log_{10}M_X`$ with a cut at $`\mathrm{\Delta }Log_{10}M_X=1`$ immediately rules out most of the region $`\xi (.25,8.6)`$ due to the ultrarapid gauge boson mediated Proton decay in that range of $`\xi `$.
When we lower the value of $`\lambda `$ or $`\eta `$ many particles became light so that the entire $`\alpha _G(M_X)^1`$ vs $`\xi `$ plot shifts downward. This is seen in Fig. 2 where we repeat the same Plot as Fig 1. but now with $`\lambda =.1`$. As a result larger $`\xi `$ values are allowed on grounds of limiting the change in $`\alpha _G(M_X)^1`$. On the other hand the behaviour of $`\mathrm{\Delta }Log_{10}M_X`$ in response to a decrease in the diagonal slow parameters $`\lambda ,\eta `$ is quite mild. Thus although lowering $`\lambda `$ does tend to make $`\mathrm{\Delta }Log_{10}M_X`$ less negative in $`\xi >0`$ region, the large positive change in $`\alpha _G(M_X)`$ would rule out the small $`\lambda ,\eta `$ region. The condition $`\mathrm{\Delta }Log_{10}M_X>1`$ required by proton stability rules out the region $`.25<\xi <8.6`$ for the real case (see Figs. 1., 2.).
Next we plot $`\mathrm{\Delta }Sin^2\theta _W(M_S)`$ versus $`\xi `$ for moderate and very small $`\lambda `$ (Fig 3.). The change in $`\lambda `$ has practically no effect. We also see that there are ranges of $`\xi `$ where the change in $`Sin^2\theta _W(M_S)`$ is less that $`10\%`$ and that these ranges are only slightly affected by the change in $`\lambda `$. However the region $`2.8<\xi <8`$ which was excluded by $`\mathrm{\Delta }Log_{10}M_X>1`$ is also excluded by the large change $`\mathrm{\Delta }Sin^2\theta _W(M_S)`$.
The further evolution of $`\mathrm{\Delta }^{(th)}(\alpha _G^1(M_X))`$ and $`\mathrm{\Delta }^{(th)}(Log_{10}M_X)`$ as $`\lambda `$ is decreased to $`.01`$ is similar as can be seen in Fig.4. Lowering $`\eta `$ has effects very similar to those of lowering $`\lambda `$ since the effect of smaller values for these diagonal couplings is to lower the masses of whole sets of multiplets and therefore raise $`\alpha _G(M_X)`$ sharply. The effects of the off diagonal couplings $`\gamma `$ and $`\overline{\gamma }`$ are much milder.
With $`\xi `$ real there are also two complex (mutually conjugate) solutions of the basic cubic equation (2). Examples of these are displayed as Fig. 5,6 for moderate and small $`\lambda `$. We observe that the behaviour of the complex solution is smoother than the real one and that apart from the spikes observed at $`\xi =5`$ there is hardly any sharp feature to be seen. The effect of decreasing $`\lambda `$ or $`\eta `$ on $`\mathrm{\Delta }^{(th)}(\alpha _G^1(M_X))`$ is however even stronger than in the case of real $`x`$ as may be seen in Fig 6, thus restricting the viable magnitudes of $`\lambda ,\eta `$ to moderately large values again. The corrections to $`Sin^2\theta _W(M_S)`$ for the complex solution are shown for moderate and very small values of $`\lambda `$ as Fig. 7. One can also consider complex values of the parameter $`\xi `$ and the three solutions in that case. The behaviour is quite similar to that shown for the complex solutions for real $`\xi `$ but needs detailed and comprehensive examination.
We have exhibited these graphs to give a sense of the structure visible once one ramps up the resolution of analysis to reveal the finestructure hidden within the bland impressiveness of “supersymmetric unification at a point”. As already noted long since the ambiguities associated with superheavy thresholds would not allow one to predict the effective scale of superpartner masses or the unification scale even if the low energy values $`\alpha _i(M_Z)`$ were known exactly. In fact, following Hall we have chosen not to treat the unification beyond leading order by imposing unification at a point but rather in terms of quantifying the ambiguities in $`\alpha _G^1(M_X)`$, $`Sin^2\theta _W(M_S)`$ and $`Log_{10}M_X`$ caused by the finestructure of unification scale mass spectra. The range of behaviours exhibited make it unlikely that the constraints of gauge unification alone will rule out this minimal Susy SO(10) GUT or fix its parameters. However when taken together with the rest of the low energy data, the MSGUT provides a well defined and calculable framework within which significant questions regarding the viable GUT scale structures can be posed and answered. Furthermore the process of fitting the highly structured fermion data to the relatively few parameters of the MSGUT excavates crucially important information regarding the embedding of the MSSM within the MSGUT. This information (concerning mixing angles of various sorts) is critical for determining the precise predictions of the MSGUT for both $`d=5,6`$ baryon number violating operators in the effective MSSM as well as the predictions for Lepton Flavour violation.
## 4 FM Fitting Frenzy
We have already reviewed the sequence of developments preceding the current focus of interest on the fitting of fermion mass and mixing data in the MSGUT. The fitting program itself has used only the form of the fermion mass formulae in the MSGUT (which follows from the use of only the $`\mathrm{𝟏𝟎}+\overline{\mathrm{𝟏𝟐𝟔}}`$ representations) rather than the specific formulae for the coefficients in the fermion masses dictated by the MSGUT superpotential. Our concern here is not with the actual values of the successful fits but rather their implications when combined with the structure of the MSGUT. We therefore review the fitting procedure from our particular viewpoint. To begin with the “Clebsch ” coefficients for the couplings of the $`\mathrm{𝟏𝟔}\times \mathrm{𝟏𝟔}`$ SO(10) chiral spinors to the $`\mathrm{𝟏𝟎},\overline{\mathrm{𝟏𝟐𝟔}}`$ irreps were calculated as a part of our explicit decomposition of $`SO(10)`$ in terms of Pati-Salam labels : One obtains :
$`W_{FM}^H`$ $`=`$ $`h_{AB}^{}\psi _A^TC_2^{(5)}\gamma _i^{(5)}\psi _BH_i`$ (9)
$`=`$ $`\sqrt{2}h_{AB}^{}\left[H_{\mu \nu }\widehat{\psi }_A^{\mu \dot{\alpha }}\widehat{\psi }_{B\dot{\alpha }}^\nu +\stackrel{~}{H}^{\mu \nu }\psi _{\mu A}^\alpha \psi _{\nu \alpha B}H^{\alpha \dot{\alpha }}(\widehat{\psi }_{A\dot{\alpha }}^\mu \psi _{\alpha \mu B}+\psi _{\alpha \mu A}\widehat{\psi }_{\dot{\alpha }B}^\mu )\right]`$
$`=`$ $`2\sqrt{2}h_{AB}^{}\overline{h}_1[\overline{d}_AQ_B+\overline{e}_AL_B]+2\sqrt{2}h_{AB}^{}h_1\left[\overline{u}_AQ_B+\overline{\nu }_AL_B\right]`$
$`+`$ $`\mathrm{}..`$
$`W_{FM}^{\overline{\mathrm{\Sigma }}}`$ $`=`$ $`{\displaystyle \frac{1}{5!}}\psi ^TC_2^{(5)}\gamma _{i_1}\mathrm{}..\gamma _{i_5}\chi \overline{\mathrm{\Sigma }}_{i_1\mathrm{}i_5}`$
$`=`$ $`4\sqrt{2}\overline{\mathrm{\Sigma }}_\nu ^{\mu \alpha \dot{\alpha }}(\widehat{\psi }_{\dot{\alpha }}^\nu \chi _{\alpha \mu }+\psi _{\mu \alpha }\widehat{\chi }_{\dot{\alpha }}^\nu )+4(\overline{\mathrm{\Sigma }}_{\mu \nu }^{\dot{\alpha }\dot{\beta }}\widehat{\psi }_{\dot{\alpha }}^\mu \widehat{\chi }_{\dot{\beta }}^\nu +\overline{\mathrm{\Sigma }}^{\mu \nu \alpha \beta }\psi _{\mu \alpha }\chi _{\nu \beta })\mathrm{}..`$
$`=`$ $`4\sqrt{2}f_{AB}^{}[{\displaystyle \frac{i}{\sqrt{3}}}\{\overline{h}_2(\overline{d}_AQ_B3\overline{e}_AL_B)h_2(\overline{u}_AQ_B3\overline{\nu }_AL_B)\}`$
$`+`$ $`4f_{AB}^{}[2iG_5\overline{\nu }_A\overline{\nu }_B)+\sqrt{2}\overline{O}L_AL_B)+\mathrm{}\mathrm{}`$
where the alphabetical naming convention regarding the subcomponents of the Higgs multiplets is given in the Appendix and in detail in. From the properties of the SO(10) Clifford algebra it follows that the Yukawa coupling matrices $`h_{}^{}{}_{AB}{}^{},f_{}^{}{}_{AB}{}^{}(A,B=1,2,3)`$ are symmetric ($`h_{}^{}{}_{AB}{}^{}=h_{}^{}{}_{BA}{}^{},f_{}^{}{}_{AB}{}^{}=f_{}^{}{}_{BA}{}^{}`$) complex matrices. Therefore the freedom to make unitary changes of basis allows one to chose one linear combination of the matrices $`h^{},f^{}`$ to be diagonal. We shall choose the basis where $`f^{}`$ is diagonal since it proves convenient when analyzing the Seesaw mechanism but our conclusions are independent of any such choice.
To obtain the formulae for the charged fermion masses from the above decomposition one first needs to define the ($`G_{321}(1,2,\pm 1)`$) multiplets $`H^{(1)},\overline{H}^{(1)}`$ which are the (light) MSSM Higgs doublet pair. This is achieved by imposing the condition
$$DetO(M_W)$$
on the doublet mass matrix $``$ which occurs in the quadratic terms of the superpotential when expanding around the superheavy vevs: $`W=\overline{h}h+\mathrm{}`$. This amounts to a fine tuning of(say) the mass parameter $`M_H`$ of the 10-plet Higgs.
The $`4\times 4`$ matrix $``$ can be diagonalized by a bi-unitary transformation : from the 4 pairs of Higgs doublets $`h^{(i)},\overline{h}^{(i)}`$ arising from the SO(10) fields to a new set $`H^{(i)},\overline{H}^{(i)}`$ of fields in terms of which the doublet mass terms are diagonal.
$`\overline{U}^TU`$ $`=`$ $`Diag(m_H^{(1)},m_H^{(2)},\mathrm{}.)`$
$`h^{(i)}`$ $`=`$ $`U_{ij}H^{(j)};\overline{h}^{(i)}=\overline{U}_{ij}\overline{H}^{(j)}`$ (11)
Then it is clear that in the effective theory at low energies the GUT Higgs doublets $`h^{(i)},\overline{h}^{(i)}`$ are present in $`H^{(1)},\overline{H}^{(1)}`$ in a proportion determined by the first columns of the matrices $`U,\overline{U}`$ :
$`E<<M_X:h^{(i)}`$ $``$ $`\alpha _iH^{(1)};\alpha _i=U_{i1}`$
$`\overline{h}^{(i)}`$ $``$ $`\overline{\alpha }_i\overline{H}^{(1)};\overline{\alpha }_i=\overline{U}_{i1}`$ (12)
Thus the formulae for the charged fermion masses in the MSGUT are :
$`M^d`$ $`=`$ $`r_1\widehat{h}+r_2\widehat{f}`$
$`M^l`$ $`=`$ $`r_1\widehat{h}3r_2\widehat{f}`$ (13)
$`M^u`$ $`=`$ $`\widehat{h}+\widehat{f}`$
where
$`\widehat{h}`$ $`=`$ $`2\sqrt{2}h^{}v\alpha _1\mathrm{sin}\beta `$
$`\widehat{f}`$ $`=`$ $`4\sqrt{{\displaystyle \frac{2}{3}}}if^{}\alpha _2\mathrm{sin}\beta `$ (14)
$`r_1`$ $`=`$ $`{\displaystyle \frac{\overline{\alpha _1}}{\alpha _1}}\mathrm{cot}\beta `$
$`r_2`$ $`=`$ $`{\displaystyle \frac{\overline{\alpha _2}}{\alpha _2}}\mathrm{cot}\beta `$
here ($`h^{},f^{})`$ are the couplings in the MSGUT superpotential.
Similarly the Majorana mass of the ‘right handed neutrinos ’ i.e of the fields $`\overline{\nu }\nu ^c`$ is read off from the decomposition given above :
$$M_{AB}^{\overline{\nu }}=4i\sqrt{2}f_{AB}^{}<\overline{\mathrm{\Sigma }}_{44}^{(R+)}>=4\sqrt{2}f_{AB}^{}\overline{\sigma }$$
(15)
and is of the order of the Unification scale or somewhat less. The left handed or SM neutrinos receive a direct Majorana mass from the so called Type II seesaw mechanism when the left handed triplet $`\overline{O}`$ contained in the $`\overline{\mathrm{𝟏𝟐𝟔}}`$ field obtains a vev. One obtains :
$$M_{AB}^\nu =4\sqrt{2}f_{AB}^{}<\overline{O}^{11}>=8if_{AB}^{}<\overline{O}_{}>$$
(16)
The vev $`<\overline{O}_{}>`$ : $`(\overline{10},3_L,1)_{\overline{\mathrm{\Sigma }}}`$ arises from a tadpole following $`SU(2)_L`$ breaking (see below).
The final component of the Seesaw is the Neutrino Dirac mass which links the left and right handed neutrinos :
Dirac mass :
$$m_{AB}^{\overline{\nu }D}=2\sqrt{2}h_{AB}^{}<h_2^{(1)}>+4i\sqrt{6}f_{AB}^{}<h_2^{(2)}>$$
(17)
To determine the Majorana mass terms of the left handed neutrinos in the effective MSSM we must eliminate the superheavy $`\overline{\nu }`$ and evaluate the Type II Seesaw tadpole. One then obtains (some factors of $`\sqrt{2}`$ have been corrected relative to eqns(77-79) of).
$$M_\nu ^I=\frac{1}{4}m_\nu ^DM_{\overline{\nu }}^1m_\nu ^D=r_4(\widehat{h}3\widehat{f})\widehat{f}^1(\widehat{h}3\widehat{f})r_4\widehat{n}$$
$$M_\nu ^{II}=8if_{AB}^{}<\overline{O}_{}>=r_3\widehat{f}$$
(18)
$`r_3`$ $`=`$ $`2i\sqrt{3}(\alpha _1\gamma +2\sqrt{3}\eta \alpha _2)({\displaystyle \frac{\alpha _4}{\alpha _2}})({\displaystyle \frac{v}{M_O}})Sin\beta `$
$`r_4`$ $`=`$ $`{\displaystyle \frac{i\alpha _2Sin\beta }{4\sqrt{3}}}{\displaystyle \frac{v}{\overline{\sigma }}}`$ (19)
$`M_O`$ $`=`$ $`2(M+\eta (3ap)))`$
Thus we see that the fermion mass formulae are completely determined in terms of the GUT scale breaking parameters($`\xi ,\lambda ,\eta ,\gamma ,\overline{\gamma },m`$), the $`10+\overline{126}`$ Yukawa couplings(15 parameters) and the low energy parameters $`v_{EW}=174GeV,tan\beta `$.
To perform the fit, we must match the fermion masses and mixings of the MSSM RG-extrapolated to the GUT scale $`M_X`$ :
$`L_{MSSM}^{FM}(M_X)`$ $`=`$ $`l^{c^T}D_ll+u^{c^T}D_uu+d^{c^T}D_dd`$ (20)
$`+`$ $`\overline{l}\mathit{}\nu +\overline{U}\mathit{}d^{}+\nu ^T(𝒫D_\nu 𝒫^T)\nu +\mathrm{}`$
with the effective theory derived from the MSGUT by integrating out the heavy fields at $`M_X`$. Here the $`D_{l,u,d,\nu }`$, C, $`𝒫:`$ are the diagonal fermion mass matrices (with mass eigenvalue components $`l_A,u_A,d_A,\nu _A`$), the CKM mixing matrix and the PMNS matrix (Majorana Neutrino mixing) at the scale $`M_X`$ in some fixed phase convention for the MSSM masses and mixings (either at $`M_Z`$ or at $`M_X`$) and $`d^{}=Cd`$. For example the convention that all the diagonal masses are real and positive. Currently each group of FM fitters uses idiosyncratic phase conventions. A standard format presentation of the data of the MSSM at $`M_X`$ is needed and is being pursued. The corresponding quadratic terms in the effective theory derived from the GUT are (GUT fields carry carets) :
$`L_{GUT}^{(2)}`$ $`=`$ $`\widehat{u^c}^T(\widehat{h}+\widehat{f})\widehat{u}+\widehat{d}^c(\widehat{h}r_1+\widehat{f}r_2)\widehat{d}+\widehat{l^c}^T(\widehat{h}r_13r_2\widehat{f})\widehat{l}`$
$`+`$ $`\widehat{\nu }^T\widehat{M}_\nu \widehat{\nu }+\overline{\widehat{l}}\mathit{}\widehat{\nu }+\overline{\widehat{u}}\mathit{}\widehat{d}`$
where
$$\widehat{M}_\nu =r_3\widehat{f}r_4(\widehat{h}+\widehat{f})\widehat{f}^1(\widehat{h}+\widehat{f})r_3\widehat{f}r_4\widehat{n}$$
(22)
When equating the two quadratic forms we must allow for Unitary transformations between the fields in the two Lagrangians :
$`L_{GUT}^{(2)}`$ $`=`$ $`L_{MSSM}^{FM}(M_X)`$
$`\left(\begin{array}{c}u\\ \\ d^{}\end{array}\right)`$ $`=`$ $`𝒬\left(\begin{array}{c}\widehat{u}\\ \\ \widehat{d}\end{array}\right);\left(\begin{array}{c}\nu \\ \\ l\end{array}\right)=\left(\begin{array}{c}\widehat{\nu }\\ \\ \widehat{l}\end{array}\right)`$ (29)
$`u^c`$ $`=`$ $`𝐕_{(u^c)}\widehat{u^c},d^c=𝐕_{(d^c)}\widehat{d^c}l^c=𝐕_{(l^c)}\widehat{l^c}.`$
The unitary matrices $`𝒬,,𝐕_{(𝐮^𝐜)},𝐕_{(𝐝^𝐜)},𝐕_{(𝐥^𝐜)}`$ describe the embedding of the extrapolated MSSM within the MSGUT. The ($`d=5,6`$) effective lagrangian ( $`_{eff}^{\mathrm{\Delta }B0}(\widehat{\psi })`$) must be transformed to the extrapolated MSSM basis in order to derive rates for such exotic processes. Thus it is clear that these matrices are neither unphysical nor conventional once the conventions of the MSSM parameters are fixed. In fact we argue that the crucial information given to us by the fitting procedure is not a prediction of the neutrino masses and mixings but rather information on these embedding matrices.
Since the $`\mathrm{𝟏𝟎},\overline{\mathrm{𝟏𝟐𝟔}}`$ Yukawas are symmetric it follows that only two (say $`𝒬,`$ ) of these matrices are independent while the others ($`𝐕_{(u^c)},𝐕_{(d^c)},𝐕_{(l^c)}`$) can be determined in terms of the two chosen to be independent plus arbitrary diagonal unitary matrices $`\mathrm{\Phi }_u,\mathrm{\Phi }_d,\mathrm{\Phi }_l`$. Once this is done we obtain :
$`\mathrm{\Phi }_u^{}𝐕_{(u^c)}`$ $`=`$ $`C\mathrm{\Phi }_d^{}𝐕_{(d^c)}=𝒬\mathrm{\Phi }_l^{}𝐕_{(l^c)}=`$
$`\widehat{h}+\widehat{f}`$ $`=`$ $`𝐕_{(u^c)}^TD_u𝒬=𝒬^TD_u^{}𝒬𝒬^T\mathrm{\Phi }_uD_u𝒬`$
$`\widehat{h}r_1+r_2\widehat{f}`$ $`=`$ $`𝐕_{d^c}^TD_dC^{}𝒬=^TD_d^{}^T\mathrm{\Phi }_dD_d`$
$`\widehat{h}r_13r_2\widehat{f}`$ $`=`$ $`𝐕_{l^c}^TD_l=^TD_l^{}^T\mathrm{\Phi }_lD_l`$
$`r_3\widehat{f}r_4\widehat{n}`$ $`=`$ $`^T𝒫D_\nu 𝒫^T`$
where we have defined $`=C^{}𝒬`$. The phase freedoms $`\mathrm{\Phi }_u,\mathrm{\Phi }_d`$ have been found to be important for arranging the tunings that underlie the successful Type I and Type II fits of the fermion masses. However the phases $`\mathrm{\Phi }_l`$ play no important role so far since the phases in the PMNS matrix $`𝒫`$ are unknown at present. If we reabsorb the phases in the equation from $`\widehat{M}_l`$ in $``$ then we also need to redefine the PMNS matrix $`𝒫`$ to reabsorb them in the neutrino mass equation (last of equations (4)). This ambiguity should be kept in mind when deducing predictions of Leptonic phases from the fit, but does not play any role at this stage.
We remind the reader that all the parameters $`r_i;i=1,2,3,4`$ are known in terms of the MSGUT parameters but the explicit form is not used or invoked when solving the fitting problem. Rather the parameters $`r_i`$ and the Yukawa couplings are determined in terms of the extrapolated experimental data. Thus the compatibility of the FM fits with the MSGUT remains to be verified for each prima facie viable fit since the ability of the MSGUT to reach the required values of the $`r_i`$ simultaneously while preserving the constraints of perturbative gauge unification is not obvious.
### 4.1 Solution of the fitting problem
The equations for the down fermion and charged lepton masses may be immediately solved to yield
$`\widehat{h}r_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}(^TD_l+3^TD_d^{})`$
$`\widehat{f}r_2`$ $`=`$ $`{\displaystyle \frac{1}{4}}(^TD_l+^TD_d^{})`$
we define
$$𝒟=^{}^T,\stackrel{ˇ}{u}=C^TD_u^{}C$$
(32)
for later convenience.
Consider first the case where the mixing matrices, (hatted) Yukawa couplings and parameters are all assumed real. Moreover the phases $`\mathrm{\Phi }_u,\mathrm{\Phi }_d`$ are simply signs and the CKM phase $`\delta `$ is also a sign which in fact is found to be minus in “successful” fits. Then solving the Up quark equations and eliminating $`r_1,r_2`$ using 23 component and trace we can put it in the form
$$XX_l𝒟D_l𝒟^T=\frac{X_{23}}{\stackrel{ˇ}{u}_{23}}(\stackrel{ˇ}{u}\frac{T_u^{}}{T_d^{}}D_d^{})+\frac{T_l}{T_d^{}}D_d^{}$$
(33)
Here $`T_f^{}=tr[D_f^{}]`$. Notice that the non diagonality of $`X_l`$ is completely driven by the non-diagonality of the matrix $`u^{}`$ which in turn follows from that of the CKM matrix. Following it is convenient to rescale each of the diagonal fermion mass matrices by the mass of the third generation fermion of that type to get a dimensionless (tilde-ed) form of the equations :
$$D_fm_{f_3}D_{\stackrel{~}{f}}$$
(34)
$$\stackrel{~}{X}_l𝒟D_{\stackrel{~}{l}}𝒟^T=\frac{\stackrel{~}{X}_{23}}{\stackrel{~}{\stackrel{ˇ}{u}}_{23}}(\stackrel{~}{\stackrel{ˇ}{u}}\frac{T_{\stackrel{~}{u}^{}}}{T_{\stackrel{~}{d}^{}}}D_{\stackrel{~}{d}^{}})+\frac{T_{\stackrel{~}{l}}}{T_{\stackrel{~}{d}^{}}}D_{\stackrel{~}{d}^{}}$$
(35)
The matrix $`X_{\stackrel{~}{l}}`$ has eigenvalues $`\stackrel{~}{l}_{1,2}`$ and $`\stackrel{~}{l}_3=1`$ by construction and hence it follows that we must have
$$det(X_{\stackrel{~}{l}}\stackrel{~}{l}_iI_3)=0$$
(36)
Thus we obtain three coupled non-linear equations for $`\stackrel{~}{X}_{23}`$ and two other quantities which are conveniently chosen to be $`\stackrel{~}{d}_{1,2}`$. We can solve numerically for $`\stackrel{~}{X}_{23}`$ and $`\stackrel{~}{d}_{1,2}`$, given any set of up quark and charged lepton masses together with CKM data. Obviously only solutions within the error bars(usually allowed to be 1-$`\sigma `$) for the down quark masses are accepted. Note however that inasmuch as there is a strong inter-generational hierarchy for the charged fermion masses the numerical solution of the three coupled non-linear equations is an extremely delicate operation requiring utmost care and diligence to find the correct solutions. Since brute shooting usually fails in multidimensional problems of this delicacy approximate analytic solutions obtained by expanding in the light fermion masses are used to guide the numerical search for solutions.
Assuming this has been done the matrix $`X_l`$ is completely determined so that on diagonalizing it one obtains the crucial matrix $`𝒟=^T`$. With $`𝒟`$ in hand one can proceed by choosing the convenient $`\widehat{f}`$-diagonal basis mentioned earlier. Rescaling
$$\stackrel{~}{\widehat{f}}=\widehat{f}/m_t=\stackrel{~}{\widehat{r}_2}^TX_{\stackrel{~}{f}}$$
(37)
where we have defined $`X_{\stackrel{~}{f}}`$ and unitary $``$ by
$$X_{\stackrel{~}{f}}=𝒟^TD_{\stackrel{~}{d}}^{}𝒟D_{\stackrel{~}{l}}(\frac{m_\tau }{m_b})D_{\stackrel{~}{f}}^T$$
(38)
Since $`\widehat{f}`$ is diagonal by choice of basis it immediately follows that
$$==𝒟$$
(39)
Now since $`,`$ are known and since $`\stackrel{~}{\widehat{r}}_{1,2}=\frac{m_b}{m_t}\widehat{r}_{1,2}`$ are calculable in terms of $`X_{\stackrel{~}{l}}`$ it is clear that we have also determined the Yukawa coupling $`\widehat{h}`$ :
$`\stackrel{~}{\widehat{h}}=\widehat{h}/m_t=\stackrel{~}{\widehat{r}_1}^T[3𝒟^TD_{\stackrel{~}{d}}^{}𝒟+D_{\stackrel{~}{l}}({\displaystyle \frac{m_\tau }{m_b}})]`$ (40)
In other words one finds that the fitting of charged fermion mass ratios requires tuning of the down quark mass ratios $`\stackrel{~}{d}_{1,2}`$ to less than one part in $`10^3`$ for given precise values of up quark and charged lepton masses together with CKM data and yields the dimensionless matrices $`\stackrel{~}{\widehat{h}},\stackrel{~}{\widehat{f}}`$ and the exotic embedding matrices $`,`$ and, given $`m_t`$, $`\widehat{h},\widehat{f}`$.
In the realistic case when the parameters are complex a similar but numerically even more difficult procedure is followed. Like the sign freedom in the real case the phase freedom of choosing $`\mathrm{\Phi }_u,\mathrm{\Phi }_d`$ is found to be crucial to obtaining a successful fit. The u equation still has the same form
$$X_l𝒟D_l𝒟^T=p\stackrel{ˇ}{u}+qD_d$$
(41)
but now $`TrX_lTrD_l`$.
To proceed we solve the 23, 33 components to eliminate $`p,q`$ and get
$`\stackrel{~}{X}_l𝒟D_{\stackrel{~}{l}}𝒟^T`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{X}_{23}}{\stackrel{~}{\stackrel{ˇ}{u}}_{23}}}\stackrel{~}{\stackrel{ˇ}{u}}+(\stackrel{~}{X}_{33}{\displaystyle \frac{\stackrel{~}{X}_{23}}{\stackrel{~}{\stackrel{ˇ}{u}}_{23}}}\stackrel{~}{\stackrel{ˇ}{u}}_{33})D_{\stackrel{~}{d}}`$ (42)
Then to determine $`𝒟`$ we must diagonalize
$$X_{\stackrel{~}{l}}X_{\stackrel{~}{l}}^{}𝒟D_{\stackrel{~}{l}}^2𝒟^{}$$
(43)
This matrix has eigenvalues $`\stackrel{~}{l}_1^2,\stackrel{~}{l}_2^2,1`$ and an inspection of its explicit form shows that it requires knowledge of $`|X_{23}|,|X_{33}|,\varphi _X=Arg(X_{23})Arg(X_{33})`$ and a choice of the quark mass phases $`\mathrm{\Phi }_u,\mathrm{\Phi }_d`$. The numerical results obtained by these authors are then a specification of the right hand side of eqn(42) consistent with some acceptable values of the charged lepton masses. We refer the reader to the original papers for the procedure for fixing the unknown phases. Once this rather horrendous numerical problem has been solved (its trickiness accounts for the fact that even 12 years after it’s solution was first explicitly attempted this system continues to throw out surprises) one can proceed essentially as before : diagonalizing the rhs of eqn(42) numerically thus yields the matrix $`𝒟`$.
Since $`\widehat{f}`$ is diagonal and real by convention one writes
$$\stackrel{~}{\widehat{f}}=\stackrel{~}{\widehat{f}}/m_t=\stackrel{~}{\widehat{r}_2}^TX_{\stackrel{~}{f}}$$
(44)
where
$$X_{\stackrel{~}{f}}=𝒟^{}D_{\stackrel{~}{d}^{}}𝒟^{}D_{\stackrel{~}{l}}(\frac{m_\tau }{m_b})D_{\stackrel{~}{f}}^T$$
(45)
Since $`X_{\stackrel{~}{f}}`$ is symmetric , $``$ is Unitary and $`D_{\stackrel{~}{f}}`$ is real like $`\stackrel{~}{\widehat{f}}`$, it follows that
$$\stackrel{~}{\widehat{f}}=|\stackrel{~}{\widehat{r}_2}|D_{\stackrel{~}{f}}$$
(46)
so that
$$=^{}e^{\frac{iArg(\stackrel{~}{\widehat{r}_2})}{2}}$$
(47)
where $``$ is found by diagonalizing $`X_{\stackrel{~}{f}}X_{\stackrel{~}{f}}^{}=D_{\stackrel{~}{f}}^2^{}`$
Finally
$$=𝒟^{};𝒬=C$$
(48)
In summary : for given up and charged lepton masses and given CKM mixing angles, by tuning $`\stackrel{~}{d}_{1,2},\delta _{CKM}`$ within their allowed ranges and for a certain choice of the phases $`\mathrm{\Phi }_u,\mathrm{\Phi }_d`$, one completely fits the charged fermion masses and determines the FM Yukawa couplings $`h^{},f^{}`$ of the GUT along with the exotic embedding matrices $`,𝒬,\mathrm{\Phi }_u,\mathrm{\Phi }_d`$.
### 4.2 Fitting Neutrino Masses
The current FM Fitting Furore was triggered by the remarkably simple observation of regarding naturalness of large atmospheric mixing angles in MSGUTs with dominant Type II seesaw mechanisms given the near equality of the $`\tau `$ lepton and bottom quark masses at $`M_X`$. If one assumes Type II domination i.e $`r_3>>r_4`$ the Seesaw formula simplifies to just
$$\stackrel{~}{\widehat{f}}=(\frac{1}{m_tr_3})^T𝒫D_\nu 𝒫^T$$
(49)
Since we chose a$`\widehat{f}`$-diagonal basis it immediately follows that
$$𝒫=^{}e^{\frac{iarg(r_3)}{2}};D_\nu =|r_3|\widehat{f}$$
(50)
Thus the neutrino mixing angles and ratio of mass squared splittings can be determined under these assumptions since we know $`\widehat{f},`$.
In the general case the neutrino mass fitting equations take the scaled form :
$$𝒩D_{\stackrel{~}{N}}𝒩^T\stackrel{~}{\widehat{f}}(\frac{r_4}{r_3})\stackrel{~}{n}=(\frac{1}{m_tr_3})^T𝒫D_\nu 𝒫^T$$
(51)
where
$$\stackrel{~}{n}=n/m_t=(\stackrel{~}{\widehat{h}}3\stackrel{~}{\widehat{f}})\stackrel{~}{\widehat{f}}^1(\stackrel{~}{\widehat{h}}3\stackrel{~}{\widehat{f}})$$
(52)
Thus the mixing matrix and neutrino masses are also completely determined :
$$𝒫=^{}𝒩e^{\frac{iArg(r_3)}{2}};D_\nu =m_t|r_3|D_{\stackrel{~}{N}}$$
(53)
These fitting problems have been formulated and solved with increasing refinement by a number of authors. It seems that Type II dominant as well as Type I and Type I plus Type II combined solutions (which are not perturbations of Type II dominant solutions ) can be found in the complex case. Further solutions probably still remain to be found and the possible solutions definitely still need to be compactly characterized and parameterized. The fits achieved so far already motivate a detailed examination of what type of Seesaw is actually allowed by the MSGUT in various regions of its parameter space. We have argued that in each case the exotic embedding matrices are -like the ‘philosophers stone’ - the so far unregarded true product of the fitting calculation. We emphasize that in practice one accepts solutions which give $`\nu `$ oscillation parameters which lie within the $`1\sigma `$ ranges around the central values of the fermion mass and mixing parameters so that the true output of the fitting calculation are the previously completely unknown embedding matrices $`𝒬,,\mathrm{\Phi }_u,\mathrm{\Phi }_d`$ which specify how an MSSM($`M_X`$) with given conventions lies within the MSGUT. These matrices are crucial for pinning down the prediction of $`\mathrm{\Delta }B0`$ processes in the MSGUT. Before considering that aspect however we turn to the question of what kinds of solutions are compatible with the MSGUT.
## 5 Scanning the MSGUT for Neutrino Masses and Mixings
The possibility of large PMNS mixing angles is well understood in the Type II dominant case , where it appears as a natural corollary of the approximate unification of the running bottom quark and tau lepton masses at scales $`O(M_X)`$ due to the 3 fold faster evolution of the bottom quark mass. As is evident from eqn.(4) the parameter which controls the strength of Type I versus Type II Seesaw in the MSGUT is the ratio of the coefficients $`r_4`$ and $`r_3`$. Complete Type II dominance requires $`r_4<<r_3`$. To illustrate how the MSGUT yields information on this ratio we work with an example of a quasi-realistic real Type II fit that ignores the CP violating phase, which was kindly provided to us by S.Bertolini and M. Malinsky. Although Type II and only semi-realistic it should be emphasized that the values of $`\widehat{h},\widehat{f}`$ given will be rather typical since they are fixed by the charged fermion mass data and further selection is then imposed based on compatibility with the neutrino data.
The data of the example solution ( at $`M_X`$ and for $`\mathrm{tan}\beta =10`$ ) is
$`D_u`$ $`=`$ $`\{0.785556,191.546,70000\}MeV`$
$`D_l`$ $`=`$ $`\{0.3585,75.7434,1290.8\}MeV`$
$`d_3`$ $`=`$ $`m_b=1138.07MeV`$
$`\{Sin\theta _{12}^c,Sin\theta _{23}^c,Sin\theta _{13}^c\}`$ $`=`$ $`\{0.2229,0.03652,0.00319\}`$ (54)
Notice the negative sign of $`u_3`$ and $`\theta _{13}^c`$ (this corresponds to taking $`\delta _{CKM}=\pi `$), moreover the fitting procedure gives
$`\stackrel{~}{x}_{23}`$ $`=`$ $`0.14325;\stackrel{~}{d}_1=.001105;\stackrel{~}{d}_2=.02747`$ (55)
so that the sign ambiguity $`\mathrm{\Phi }_d`$ is also fixed.
$$𝒟=\left(\begin{array}{ccc}0.98953& 0.136941& 0.045582\\ 0.128508& 0.979735& 0.153642\\ 0.0656981& 0.146176& 0.987075\end{array}\right)$$
(56)
Supplying $`m_b=1138.07MeV`$ allows one to calculate
$$==\left(\begin{array}{ccc}0.794035& 0.582472& 0.173881\\ 0.563923& 0.599057& 0.568437\\ 0.226934& 0.549415& 0.804142\end{array}\right)$$
(57)
Then the matrices $`,𝒬`$ follow :
$``$ $`=`$ $`\left(\begin{array}{ccc}0.87329& 0.469294& 0.130873\\ 0.415588& 0.577356& 0.702813\\ 0.254266& 0.668149& 0.699233\end{array}\right)`$
$`𝒬`$ $`=`$ $`\left(\begin{array}{ccc}0.943138& 0.330924& 0.0313091\\ 0.219729& 0.691351& 0.688297\\ 0.24942& 0.642279& 0.724753\end{array}\right)`$ (58)
The scaled Yukawa couplings are
$`\stackrel{~}{\widehat{f}}`$ $`=`$ $`\left(\begin{array}{ccc}0.000595629& 0& 0\\ 0& 0.0020326& 0\\ 0& 0& 0.00804198\end{array}\right)`$ (62)
$`\stackrel{~}{\widehat{h}}`$ $`=`$ $`\left(\begin{array}{ccc}0.0626836& 0.159778& 0.181181\\ 0.159778& 0.409183& 0.466796\\ 0.181181& 0.466796& 0.532013\end{array}\right)`$ (66)
The matrix $`\widehat{n}`$ multiplying $`r_4`$ in the Type I mass is
$`\stackrel{~}{\widehat{n}}`$ $`=`$ $`\left(\begin{array}{ccc}1.4996& 3.8747& 4.55331\\ 3.8747& 9.98038& 11.6875\\ 4.55331& 11.6875& 13.63\end{array}\right)`$ (70)
From these the Type II mixing angles and the ratio of the 23 and 12 sector mass squared splittings is found to be
$`\{Sin^2(2\theta _{12}^P),Sin^2(2\theta _{23}^P),Sin\theta _{13}^P\}`$ $`=`$ $`\{0.90981,0.88873,0.17387\}`$
$`\mathrm{\Delta }m_{12}^2/\mathrm{\Delta }m_{23}^2`$ $`=`$ $`0.06237`$ (71)
The Yukawa couplings obtained are typical of Type II fits. It is clear that unless $`r_4/r_3`$ is very small the Type I term will dominate completely. To see just how small $`r_4`$ should be we plot the neutrino mixing angles $`Sin^22\theta _{12},Sin^22\theta _{23}`$ versus this ratio as Fig 8.
From Fig. 8 it is evident that a value of $`|r_4/r_3|>.0003`$ causes a collapse of the large mixing angles of the Type II dominant solution. Thus we should not expect a Type II solution to work in the MSGUT unless $`R=r_4/r_3`$ (which is completely determined by the GUT as given in the previous section ) obeys $`|R|<<10^3`$. The MSGUT formulae for $`r_3,r_4`$ are given in eqn(4)
Using these we plot the ratio R versus $`\xi `$ and find that its typical value is $`10^1`$ or more, not $`10^3`$ or less. Illustrative plots for the real and complex solutions of the cubic eqn.(2), for real $`\xi `$ and $`\lambda =.7`$ are shown as Fig 9.
From the plots it is evident that in the real case R has a chance of being of the required small size only in the window near $`\xi =.7`$ and between about $`\xi =1`$ and $`\xi =4`$. However in the former region one finds that the corrections to $`Sin^2\theta _W`$ diverge, while the region $`\xi (1,4)`$ is not allowed by the requirement that $`\mathrm{\Delta }Log_{10}M_X>1`$. On the other hand, in the complex case, the ratio seems bounded below by 1 ! Thus we see that if Type II fits were the only allowed ones then the combination of the gauge unification requirements with those imposed by neutrino mass phenomenology tend to rule out the MSGUT based on this class of solutions. This conclusion however still requires a more thorough study of all possibilities. Moreover recent work has shown that in fact large mixing angles can be achieved even in Type I solutions and Type I-II combined solutions (which are far from pure Type II solutions). Our intent, in this talk, is not to provide an exhaustive survey of the possibilities but only to illustrate how the combined requirements of neutrino oscillations and baryon stability can severely constrain the MSGUT over its parameter space. A detailed survey of the MSGUT for each of the three solutions of the basic cubic equation to find in which regions, if any, the ratio R matches that required by the FM fit is quite feasible using our methods, and will be reported elsewhere.
Note that while the FM fits described in the previous section do not determine the over all mass scale of the neutrinos since the input data does not contain this information, the same is not true in the context of the MSGUT. Using the dimensionless versions $`\stackrel{~}{\widehat{f}}`$ and $`\stackrel{~}{\widehat{h}}`$ given by the real Type II FM fitting analysis as a guide to typical values of the Yukawa couplings allows one to compute the magnitudes of all three neutrino masses for each type of fit. When this is done another problem becomes apparent : except possibly for narrow ranges of $`\xi `$ the largest neutrino mass (Type I or II) is much smaller than the mass splitting known from atmospheric neutrino oscillations. Thus even if one can find a region where the ratio of mass splittings and mixing angles for neutrinos are in the allowed region the additional consistency constraint that $`|(m_\nu )_{max}|>|(\mathrm{\Delta }m)_{atmos}|`$ alone can exclude most of the parameter space ! An illustrative plot of the maximum type I and Type II neutrino masses for the Yukawa coupling matrices that arises in the Bertolini and Malinsky solution used by us for illustration is given for the real solution as Fig. 10. Due to the cut at $`\xi =8.6`$ imposed by the condition $`\mathrm{\Delta }Log_{10}M_X>1`$ we see that there is no region with large enough values of $`|(m_\nu )_{max}|`$. Similarly in the complex case we get the plots shown in Fig. 11. As expected the Type I dominates completely. However the mass magnitudes for Type I tend to be too small : as can be seen from Fig. 11 and the magnifications shown in Fig 12. Since the shortfall is only a factor of 10 and we have not yet used the complex Type I fits found in it would be premature to rule out Type I fits. Nevertheless a certain tension is apparent.
If the impossibility of large enough overall neutrino mass scale is borne out by a comprehensive analysis it would require a revision of the perturbative MSGUT. Similar arguments were used in (for the Type II FM fit case only) to motivate an extension of the model by introducing an addition $`\mathrm{𝟓𝟒}`$ plet so as raise the value of the $`\overline{\mathrm{𝟏𝟐𝟔}}`$ vev to a high scale to allow Type II to dominate and yield a large enough neutrino mass. Note however that in that case the minimality of the MSGUT is seriously diluted necessitating a complete reanalysis which has not been performed so far. An alternative to this extension may be to implement a non perturbative mechanism based on dynamical GUT symmetry breaking down to $`SU(5)\times U(1)`$ at $`\mathrm{\Lambda }_X10^{17}GeV`$ and then use a Type I fit with small $`\overline{\mathrm{𝟏𝟐𝟔}}`$ vev $`\overline{\sigma }`$ to achieve a larger value of the overall neutrino mass scale.
## 6 $`\mathrm{\Delta }B0`$ IN MSGUT
Finally we briefly indicate the $`d=5`$ baryon decay operators determined by the FM fit. If any FM fit is found to be consistent with all the constraints discussed above then the computation of the actual Baryon Decay predictions will become a worthwhile exercise.
The basis for this is the effective superpotential that arises when the superheavy Higgs triplets are integrated out from the theory
$`W_{eff}^{\mathrm{\Delta }B0}`$ $`=`$ $`\widehat{L}_{ABCD}({\displaystyle \frac{1}{2}}ϵ\widehat{Q}_A\widehat{Q}_B\widehat{Q}_C\widehat{L}_D)+\widehat{R}_{ABCD}(ϵ\widehat{\overline{e}}_A\widehat{\overline{f}u}_B\widehat{\overline{u}}_C\widehat{\overline{d}}_D)`$ (72)
$`\widehat{L}_{ABCD}`$ $`=`$ $`𝒮_1^1h_{AB}h_{CD}+𝒮_1^2h_{AB}f_{CD}+𝒮_2^1f_{AB}h_{CD}+𝒮_2^2f_{AB}f_{CD}`$ (73)
$`\widehat{R}_{ABCD}`$ $`=`$ $`𝒮_1^1h_{AB}h_{CD}𝒮_1^2h_{AB}f_{CD}𝒮_2^1f_{AB}h_{CD}+𝒮_2^2f_{AB}f_{CD}`$ (74)
$``$ $`i\sqrt{2}𝒮_1^4f_{AB}h_{CD}+i\sqrt{2}𝒮_2^4f_{AB}f_{CD}`$
Here $`𝒮=𝒯^1`$ where $`𝒯`$ is the mass matrix in the $`[3,1,\pm 2/3]`$-sector which is the representation type that mediates $`d=5`$ baryon decay. The Yukawa coefficients $`h_{AB},f_{AB}`$ are related to those in the superpotential by $`h_{AB}=2\sqrt{2}h_{AB}^{},f_{AB}=4\sqrt{2}f_{AB}^{}`$
Substituting for the GUT superfields (with carets) in terms of the embedding matrices $`𝒬,,`$ determined by the FM fit
$`\widehat{Q}_L`$ $`=`$ $`𝒬^{}Q_L\widehat{L}_L=^{}L_L`$ (75)
$`\widehat{\overline{u}}=\widehat{u}_c`$ $`=`$ $`𝒬^{}\mathrm{\Phi }_u^{}u_c\widehat{\overline{d}}=\widehat{d}_c=^{}\mathrm{\Phi }_d^{}d_c\widehat{\overline{l}}=\widehat{l}_c=^{}\mathrm{\Phi }_l^{}l_c`$
one obtains the coefficients of mass eigenstate MSSM fields in the $`\mathrm{\Delta }B0`$ superpotential to be
$`L_{ABCD}`$ $`=`$ $`𝒬_{AE}^{}𝒬_{BF}^{}𝒬_{CG}^{}_{DH}^{}\widehat{L}_{EFGH}`$
$`R_{ABCD}`$ $`=`$ $`_{AE}^{}𝒬_{BF}^{}𝒬_{CG}^{}(^{})_{DH}(\mathrm{\Phi }_l^{})_{EE}(\mathrm{\Phi }_u^{})_{FF}(\mathrm{\Phi }_u^{})_{GG}(\mathrm{\Phi }_d^{})_{HH}\widehat{R}_{EFGH}`$ (76)
A similar calculation can be done to determine the effective operators for $`d=6`$ baryon violating operators, which will be relevant even if the Supersymmetry breaking scale eventually turns out to be large enough ($`>100`$TeV) to exclude any observable $`d=5`$ baryon decay. Similar care needs to be exercised when studying lepton flavour violation.
## 7 Conclusions and Outlook
The MSGUT based on the $`\mathrm{𝟐𝟏𝟎}\mathrm{𝟏𝟐𝟔}\overline{\mathrm{𝟏𝟐𝟔}}\mathrm{𝟏𝟎}`$ Higgs system is the simplest Supersymmetric GUT that elegantly realizes the classic program of Grand Unification and which is prima facie compatible with all known data. Its symmetry breaking structure is so simple as to permit an explicit analysis of its mass spectrum at the GUT scale and an evaluation therefrom of the threshold corrections and mixing matrices relevant to various phenomenologically important quantities. Since it has the least number of parameters of any theory that accomplishes as much this theory currently merits the name of the minimal supersymmetric GUT or MSGUT. The same simplicity and analyzability of GUT scale structure also applies to the theory with an additional 120-plet, since it contains no standard model singlets, and thus justifies calling it the nMSGUT.
The small number of Yukawa couplings of the MSGUT makes the fit to the now well characterized fermion mass spectra very tight so that the combined constraints of this fit and the preservation of the gauge unification observed at one loop in the MSSM may well be enough to rule out most or even all of the parameter space of this theory. A corresponding investigation is also possible for the nMSGUT and both are in progress.
We have emphasized the need for clarity regarding the flavour basis used when performing the FM fit. When this is maintained it becomes obvious that the true outputs of the FM fitting calculation are really the embedding matrices that define how the MSSM, with fixed phase conventions, lies within the MSGUT. With these matrices determined by the FM fit in hand, one will be in a position to perform a much more reliable calculation of $`d=5,6`$ Baryon decay in the MSGUT with the Susy breaking scale as the chief remaining uncertainty. This realization underlines and emphasizes the organic connection between the physics of neutrino mass and Baryon violation i.e between the physics of ultraheavy and ultra-light particles which is the most intriguing implication of the discoveries of Super-Kamiokande.
Our discussion concerning embedding angles also has implications for a treatment of Lepton Flavour violation in the MSGUT which will be worth exploring in detail if a viable region of the parameter space of the perturbative MGUT emerges. On the other hand if no such candidate region is available then the remaining options that will need to be explored could consist on the one hand of the analogous calculations in the nMSGUT, namely when an additional 120-plet is introduced : which radically enlarges the possibilities as far as FM fitting is concerned. Or the MSGUT could be extended by engineering the models to ensure Type II dominance.
Another possibility is that the symmetry breaking at the GUT scale is primarily determined not by the perturbative superpotential but rather by a non-perturbative mechanism whereby gaugino condensation in the coset SO(10)/H (where H could be, e.g, $`SU(5)\times U(1)`$ or $`SU(5)`$) drives a H-singlet Chiral condensate of (say) the 210-plet field at a scale $`\mathrm{\Lambda }/\lambda ^\alpha >M_X`$. In that case the spectra given by us in terms of the vevs $`a,p,\omega `$ are still of use but they are no longer determined by the cubic equation(2). Rather after breaking the symmetry to the group H at a scale higher than the perturbative scale $`M_X`$ one would examine the symmetry breaking in the effective $`SU(5)`$ symmetric theory at lower energies $`M_X`$. Such a scenario would thus not only utilize the problematic strong coupling regime lying just above the perturbative unification scale for dynamical symmetry breaking of the GUT theory at a scale $`\mathrm{\Lambda }_X`$ ( which is determined by the low energy data and the Grand Unified structure and functions as an internally defined upper cutoff for the MSGUT) but also provide an explanation for the “SU(5) conspiracy” that seems to operate within SO(10) Susy GUTs when proper account is taken of neutrino data, superheavy Susy spectra and RG evolution. This kind of scenario would fit naturally with a Type I mechanism which favours small values of the $`\overline{\mathrm{𝟏𝟐𝟔}}`$ vev. Thus it contrasts sharply with the proposal of and should be distinguishable phenomenologically from it.
We conclude with a tentative proposal to reconcile String theory based models and the $`\overline{\mathrm{𝟏𝟐𝟔}}`$ based Type I and Type II seesaw mechanisms that occur in RGUTs. Recall that String theory, particularly with level 1 Kac-Moody algebras, does not favour the emergence of effective GUTs containing adjoint and larger representations as its massless sector. Even the use of higher KM levels permits the occurrence of only very restricted numbers and types of higher GUT representations. In particular, in the case of SO(10), one cannot enlist the type of combinations characteristic of RGUTs ($`\mathrm{𝟒𝟓}\mathrm{𝟓𝟒}\overline{\mathrm{𝟏𝟐𝟔}}\mathrm{𝟏𝟐𝟔}`$ or $`\mathrm{𝟐𝟏𝟎}\overline{\mathrm{𝟏𝟐𝟔}}\mathrm{𝟏𝟐𝟔}`$) etc. These difficulties led to a decline in attempts to SO(10) GUTs from string theory. However of late the growing appreciation of the naturalness of SO(10) unification in the light of discovered neutrino mass has motivated renewed effort in this direction. SO(10) type families are generated with gauge symmetry pre-broken to MSSM or somewhat larger. However the issue of implementing the seesaw mechanism whether Type I or Type II in the way achieved so naturally in RGUTs remains problematic due to the difficulty of in building models in which the $`\overline{\mathrm{𝟏𝟐𝟔}}\mathrm{𝟏𝟐𝟔}`$ representations remain massless in the string model.
A way out of this difficulty may perhaps be found by appreciating that in RGUTs the $`\overline{\mathrm{𝟏𝟐𝟔}}\mathrm{𝟏𝟐𝟔}`$ fields are superheavy and their neutral components have either a superlarge vev (corresponding to $`M_{BL}M_X`$ ) or very small vevs ($`M_W`$ or $`m_\nu `$ ). The former kind of vev is that of the right handed triplets that give rise to the right handed neutrino’s superlarge Majorana mass and thus a small Type I seesaw mass for the MSSM neutrino. Its large size is compatible with the “pre-broken” structure of “string derived GUTs” where the breaking of the SO(10)/GUT gauge symmetry still discernible in the matter super multiplet structure is accomplished already at the level of defining the light sector of the String theory and in an effective description corresponds to a vevs of superheavy fields in the appropriate large chiral representations (such as $`\mathrm{𝟒𝟓},\mathrm{𝟓𝟒},\mathrm{𝟐𝟏𝟎},\overline{\mathrm{𝟏𝟐𝟔}},\mathrm{𝟏𝟐𝟔}`$ etc ). On the other hand the light vev $`M_W<M_S1TeV`$ in MSSM doublet channels occurs only after breaking of supersymmetry and its small size is then thought to be naturally compatible with the status of Susy breaking as a tiny correction to the pre-broken String GUT picture which is supersymmetric, conformal etc and therefore its derivation can be legitimately postponed. Furthermore the Type II seesaw mass generating vev of the left handed triplet in the $`\overline{\mathrm{𝟏𝟐𝟔}}`$ is an even higher order effect that arises due to a tadpole in a superheavy field induced once EW symmetry breaking has taken place due to the coupling of this superheavy field with the doublets that get a small EW weak vev.
Thus it appears that the $`\overline{\mathrm{𝟏𝟐𝟔}}\mathrm{𝟏𝟐𝟔}`$ fields need no longer be sought in the massless sector of the String theory. Instead it is sufficient to investigate whether the superheavy “Higgs Channel” corresponding to $`\overline{\mathrm{𝟏𝟐𝟔}}`$ or the above mentioned relevant sub-representations in fact couple to the putative light fields in an appropriate way. To put it simply the implementation of the $`R_p`$ preserving seesaw mechanism in Stringy GUTs may require a small “leakage” connecting the superheavy “field” in the $`\overline{\mathrm{𝟏𝟐𝟔}}`$ channel to the matter fermions and light doublets and the availability of the control parameter $`M_W/M_X`$ provides a natural way to keep the destabilizing effects of such heavy-light couplings under control : given that some way has somehow been found to break Susy and generate the EW scale in the first place ! This refinement of the effective theory paradigm used to extract the low energy theory from string models is both novel and consonant with the characteristic and elegantly consistent tying together of very large and very small mass scales achieved by the seesaw mechanism.
## 8 Acknowledgments
I thank Sumit Garg for collaboration and technical help while studying FM fitting in the MSGUT and S. Bertolini and M.Malinsky for providing me with a complete set of data for one of their FM fits. It is a pleasure for to acknowledge stimulating discussions with K.Babu, Stefano Bertolini, I. Dorsner, Ilya Gogoladze, M. Malinsky, R.N. Mohapatra, S. Raby, G. Senjanović, F. Vissani, during the Conference Plank05 at Trieste May 23-28, 2005. I am grateful to to the G. Senjanovic and High Energy Group of the Abdus Salam ICTP, Trieste, K. Huitu and the Theory Group, Institute of Physics, Helsinki, and and the Theory Group, CERN, Geneva for hospitality during the writing of these proceedings and A. Faraggi and H.P. Nilles for discussions concerning the possibility of implementing R-parity preserving seesaw and $`SO(10)`$ in string theory.
Appendix : Tables of masses and mixings
In this appendix we collect our results for the chiral fermion/gaugino states, masses and mixing matrices for the reader’s convenience. Mixing matrix rows are labelled by barred irreps and columns by unbarred. Unmixed cases(i)) are given as Table I.
ii) Chiral Mixed states
a)$`[8,1,0](R_1,R_2)(\widehat{\varphi }_{\overline{\mu }}^{\overline{\nu }},\widehat{\varphi }_{\overline{\mu }(R0)}^{\overline{\nu }})`$
$$=2\left(\begin{array}{cc}(m\lambda a)& \sqrt{2}\lambda \omega \\ \sqrt{2}\lambda \omega & m+\lambda (pa)\end{array}\right)$$
(77)
$$m_{R_\pm }=|_\pm |=|2m[1+(\frac{\stackrel{~}{p}}{2}\stackrel{~}{a})\pm \sqrt{(\frac{\stackrel{~}{p}}{2})^2+2\stackrel{~}{\omega }^2}]|$$
(78)
The corresponding eigenvectors can be found by diagonalizing the matrix $`^{}`$.
b) $`[1,2,1](\overline{h}_1,\overline{h}_2,\overline{h}_3,\overline{h}_4)[1,2,1](h_1,h_2,h_3,h_4)`$ $`.(H_{\dot{2}}^\alpha ,\overline{\mathrm{\Sigma }}_{\dot{2}}^{(15)\alpha },\mathrm{\Sigma }_{\dot{2}}^{(15)\alpha },\frac{\varphi _{44}^{\dot{2}\alpha }}{\sqrt{2}})(H_{\alpha \dot{1}},\overline{\mathrm{\Sigma }}_{\alpha \dot{1}}^{(15)},\mathrm{\Sigma }_{\alpha \dot{1}}^{(15)},\frac{\varphi _\alpha ^{44\dot{1}}}{\sqrt{2}})`$
$`=\left(\begin{array}{cccc}M_H& +\overline{\gamma }\sqrt{3}(\omega a)& \gamma \sqrt{3}(\omega +a)& \overline{\gamma }\overline{\sigma }\\ \overline{\gamma }\sqrt{3}(\omega +a)& 0& (2M+4\eta (a+\omega ))& 0\\ \gamma \sqrt{3}(\omega a)& (2M+4\eta (a\omega ))& 0& 2\eta \overline{\sigma }\sqrt{3}\\ \sigma \gamma & 2\eta \sigma \sqrt{3}& 0& 2m+6\lambda (\omega a)\end{array}\right)`$ (83)
The above matrix is to be diagonalized after imposing the fine tuning condition $`Det=0`$ to keep one pair of doublets light.
c) $`[\overline{3},1,\frac{2}{3}](\overline{t}_1,\overline{t}_2,\overline{t}_3,\overline{t}_4,\overline{t}_5)[3,1,\frac{2}{3}](t_1,t_2,t_3,t_4,t_5)`$ $`.(H^{\overline{\mu }4},\overline{\mathrm{\Sigma }}_{(a)}^{\overline{\mu }4},\mathrm{\Sigma }_{(a)}^{\overline{\mu }4},\mathrm{\Sigma }_{R0}^{\overline{\mu }4},\varphi _{4(R+)}^{\overline{\mu }})(H_{\overline{\mu }4},\overline{\mathrm{\Sigma }}_{(a)\overline{\mu }4},\mathrm{\Sigma }_{\overline{\mu }4(a)},\overline{\mathrm{\Sigma }}_{\overline{\mu }4(R0)},\varphi _{\overline{\mu }(R)}^4)`$
$`𝒯=\left(\begin{array}{ccccc}M_H& \overline{\gamma }(a+p)& \gamma (pa)& 2\sqrt{2}i\omega \overline{\gamma }& i\overline{\sigma }\overline{\gamma }\\ \overline{\gamma }(pa)& 0& 2M& 0& 0\\ \gamma (p+a)& 2M& 0& 4\sqrt{2}i\omega \eta & 2i\eta \overline{\sigma }\\ 2\sqrt{2}i\omega \gamma & 4\sqrt{2}i\omega \eta & 0& 2M+2\eta p+2\eta a& 2\sqrt{2}\eta \overline{\sigma }\\ i\sigma \gamma & 2i\eta \sigma & 0& 2\sqrt{2}\eta \sigma & 2m2\lambda (a+p4\omega )\end{array}\right)`$ (89)
iii) Mixed gauge chiral.
a)$`[1,1,0](G_1,G_2,G_3,G_4,G_5,G_6)(\varphi ,\varphi ^{(15)},\varphi _{(R0)}^{(15)},\frac{\mathrm{\Sigma }_{(R)}^{44}}{\sqrt{2}},\frac{\overline{\mathrm{\Sigma }}_{44((R+)}}{\sqrt{2}},\frac{\sqrt{2}\lambda ^{(R0)}\sqrt{3}\lambda ^{(15)}}{\sqrt{5}})`$
$`𝒢=2\left(\begin{array}{cccccc}m& 0& \sqrt{6}\lambda \omega & \frac{i\eta \overline{\sigma }}{\sqrt{2}}& \frac{i\eta \sigma }{\sqrt{2}}& 0\\ 0& m+2\lambda a& 2\sqrt{2}\lambda \omega & i\eta \overline{\sigma }\sqrt{\frac{3}{2}}& i\eta \sigma \sqrt{\frac{3}{2}}& 0\\ \sqrt{6}\lambda \omega & 2\sqrt{2}\lambda \omega & m+\lambda (p+2a)& i\eta \sqrt{3}\overline{\sigma }& i\sqrt{3}\eta \sigma & 0\\ \frac{i\eta \overline{\sigma }}{\sqrt{2}}& i\eta \overline{\sigma }\sqrt{\frac{3}{2}}& i\eta \sqrt{3}\overline{\sigma }& 0& M+\eta (p+3a6\omega )& \frac{\sqrt{5}g\sigma ^{}}{2}\\ \frac{i\eta \sigma }{\sqrt{2}}& i\eta \sigma \sqrt{\frac{3}{2}}& i\eta \sqrt{3}\sigma & M+\eta (p+3a6\omega )& 0& \frac{\sqrt{5}g\overline{\sigma }^{}}{2}\\ 0& 0& 0& \frac{\sqrt{5}g\sigma ^{}}{2}& \frac{\sqrt{5}g\overline{\sigma }^{}}{2}& 0\end{array}\right)`$ (96)
b) $`[\overline{3},2,\frac{1}{3}](\overline{E}_2,\overline{E}_3,\overline{E}_4,\overline{E}_5)[3,2,\frac{1}{3}](E_2,E_3,E_4,E_5)`$ $`.(\overline{\mathrm{\Sigma }}_{4\alpha \dot{1}}^{\overline{\mu }},\varphi _{(s)\alpha \dot{2}}^{\overline{\mu }4},\varphi _{\alpha \dot{2}}^{(a)\overline{\mu }4},\lambda _{\alpha \dot{2}}^{\overline{\mu }4})(\mathrm{\Sigma }_{\overline{\mu }\alpha \dot{2}}^4,\varphi _{\overline{\mu }4\alpha \dot{1}}^{(s)},\varphi _{\overline{\mu }4\alpha \dot{1}}^{(a)},\lambda _{\overline{\mu }\alpha \dot{1}})`$
$`=\left(\begin{array}{cccc}2(M+\eta (a3\omega ))& 2\sqrt{2}i\eta \sigma & 2i\eta \sigma & ig\sqrt{2}\overline{\sigma }^{}\\ 2i\sqrt{2}\eta \overline{\sigma }& 2(m+\lambda (a\omega ))& 2\sqrt{2}\lambda \omega & 2g(a^{}\omega ^{})\\ 2i\eta \overline{\sigma }& 2\sqrt{2}\lambda \omega & 2(m\lambda \omega )& \sqrt{2}g(\omega ^{}p^{})\\ ig\sqrt{2}\sigma ^{}& 2g(a^{}\omega ^{})& g\sqrt{2}(\omega ^{}p^{})& 0\end{array}\right)`$ (101)
c)$`[1,1,2](\overline{F}_1,\overline{F}_2,\overline{F}_3)[1,1,2](F_1,F_2,F_3)`$ $`.(\overline{\mathrm{\Sigma }}_{44(R0)},\varphi _{(R)}^{(15)},\lambda _{(R)})(\mathrm{\Sigma }_{(R0)}^{44},\varphi _{(R+)}^{(15)},\lambda _{(R+)})`$.
$`=\left(\begin{array}{ccc}2(M+\eta (p+3a))& 2i\sqrt{3}\eta \sigma & g\sqrt{2}\overline{\sigma }^{}\\ 2i\sqrt{3}\eta \overline{\sigma }& 2(m+\lambda (p+2a))& \sqrt{24}ig\omega ^{})\\ g\sqrt{2}\sigma ^{}& \sqrt{24}ig\omega ^{}& 0\end{array}\right)`$ (105)
d) $`[\overline{3},1,\frac{4}{3}](\overline{J}_1,\overline{J}_2,\overline{J}_3,\overline{J}_4)[3,1,\frac{4}{3}](J_1,J_2,J_3,J_4)`$ $`.(\mathrm{\Sigma }_{(R)}^{\overline{\mu }4},\varphi _4^{\overline{\mu }},\varphi _4^{\overline{\mu }(R0)},\lambda _4^{\overline{\mu }})(\overline{\mathrm{\Sigma }}_{\overline{\mu }4(R+)},\varphi _{\overline{\mu }}^4,\varphi _{\overline{\mu }(R0)}^4,\lambda _{\overline{\mu }}^4)`$
$`𝒥=\left(\begin{array}{cccc}2(M+\eta (a+p2\omega ))& 2\eta \overline{\sigma }& 2\sqrt{2}\eta \overline{\sigma }& ig\sqrt{2}\sigma ^{}\\ 2\eta \sigma & 2(m+\lambda a)& 2\sqrt{2}\lambda \omega & 2ig\sqrt{2}a^{}\\ 2\sqrt{2}\eta \sigma & 2\sqrt{2}\lambda \omega & 2(m+\lambda (a+p))& 4ig\omega ^{}\\ ig\sqrt{2}\overline{\sigma }^{}& 2\sqrt{2}iga^{}& 4ig\omega ^{}& 0\end{array}\right)`$ (110)
e)$`[3,2,\frac{5}{3}](\overline{X}_1,\overline{X}_2,\overline{X}_3)[3,2,\frac{5}{3}](X_1,X_2,X_3).(\varphi _{\alpha \dot{1}}^{(s)\overline{\mu }4},\varphi _{\alpha \dot{1}}^{(a)\overline{\mu }4},\lambda _{\alpha \dot{1}}^{\overline{\mu }4})(\varphi _{\overline{\mu }4\alpha \dot{2}}^{(s)},\varphi _{\overline{\mu }4\alpha \dot{2}}^{(a)},\lambda _{\overline{\mu }4\alpha \dot{2}})`$
$`𝒳=\left(\begin{array}{ccc}2(m+\lambda (a+\omega ))& 2\sqrt{2}\lambda \omega & 2g(a^{}+\omega ^{})\\ 2\sqrt{2}\lambda \omega & 2(m+\lambda \omega )& \sqrt{2}g(\omega ^{}+p^{})\\ 2g(a^{}+\omega ^{})& \sqrt{2}g(\omega ^{}+p^{})& 0\end{array}\right)`$ (114)
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# Microscopic Model of Charge Carrier Transfer in Complex Media
## 1 INTRODUCTION
The percolation concept has been a key notion for understanding transport and conduction processes in a wide range of complex disordered media. A few stray examples are ionic conduction in polymeric, amorphous or glassy ceramic electrolytes, diffusion in biological tissues and permeability of disordered membranes .
Most of physical situations studied so far (see Refs. and references therein) concern systems in which complex disordered environment can be considered as ”frozen”; that is, the random environment in which a given transfer process takes place does not change in time. This is certainly the case in many instances, but it is not true in general. In fact, there are many experimental systems in which the static percolation picture is not valid since the structure of the host material reorganizes itself on a time scale comparable to that at which the transfer itself occurs. Such a reorganization happens, namely, in certain biomembranes , solid protonic conductors , oil-continuous microemulsions or polymer electrolytes .
In particular, ionic transport across a biomembrane, such as, e.g., gramicidin-$`A`$, proceeds by the motion of ions through molecular channels along which they encounter potential barriers whose heights fluctuate randomly in time. The fluctuations of potential barriers may impede the transport constituting an important controlling factor . In the case of protonic conduction by the Grotthus mechanism , site-to-site hopping of charge carriers takes place between neighboring $`H_2O`$ or $`NH_3`$ groups that have a favorable relative orientation. Here, the structural host-reorganization process interacting with the carrier motion occurs due to thermally activated rotation of the $`H_2O`$ or $`NH_3`$ groups. In a similar fashion, within oil-continuous microemulsions, the charge transport proceeds by charge being transferred from one water globule to another, as globules approach each other in their Brownian motion . Last but not least, in polymer electrolytes, such as, e.g. polyethylene oxide complexed non-stoichiometrically with the ionic salt $`NaSCN`$, the $`Na^+`$ ions are largely tetrahedrally coordinated by polyether oxygens, but at the same time that $`Na^+`$ ions hop from one fourfold coordination site to another, the oxygens themselves, along with the polymeric backbone, undergo large-amplitude wagging and even diffusive motion .
Clearly, all these examples involve two characteristic time scales: the typical time $`\tau `$ between two successive hops of the charge carrier and a typical renewal time $`\tau ^{}`$ of the environment itself; that is, the time at which the host complex medium re-organizes itself sufficiently enough to provide a new set of available pathways for charge carrier transfer. Accordingly, the conventional static percolation picture can be strictly valid when only the characteristic time $`\tau ^{}`$ gets infinitely large. Only in this case one may expect an anomalous large-scale dynamics. On contrary, when $`\tau ^{}`$ is finite, dynamic percolation picture has to be applied. In this situation, one encounters quite a different behavior when compared to the random environments with quenched disorder. As a result, one expects an Ohmic-type or Stokes-type linear velocity-force relation for the carrier’s terminal velocities as a function of the applied field, in contrast to the threshold behavior and anomalous dynamics predicted by the conventional static percolation theory.
The prefactor in the linear velocity-force relation may depend, however, in a non-trivial way on the system’s parameters and this dependence constitutes the main challenge for the theoretical analysis here. On the other hand, we note that in the above mentioned examples of the dynamic percolative environments quite different physical processes are responsible for the time evolution of the host medium. Consequently, one expects that the prefactor in the Stokes-type velocity-force relation should also be dependent on the precise mechanism which underlies the temporal re-organization of the environment.
Theoretical modeling of charge carrier transfer in dynamic percolative environments has followed basically two different lines of thought. Early models of dynamic percolation described the random environment within the framework of a standard bond-percolation model allowing the strength of each bond to fluctuate in time between zero and some finite value. In this approach, the dynamics of the host medium was accounted for by a series of instantaneous renewal events. These events were assumed to occur at random times, chosen from a renewal time distribution. In the renewal process the positions of all unblocked bonds are being reassigned, such that after each renewal event a carrier sees a newly defined network. This approach is thus characterized by a global dynamical disorder without global conservation laws and correlations, since the entire set of random hopping rates is renewed independently of the previous history. Another model characterized by a local dynamical disorder has been proposed in Refs. and , and subsequently generalized to the non-Markovian case in Ref.. The difference of this model in regard to the previous one is that here the hopping rates at different sites fluctuate independently of each other. That is, states of individual bonds, rather than that of the whole lattice change in the renewal events. To describe the dynamical behavior in the local dynamical disorder case, a dynamical mean-field theory has been proposed , based on the effective medium approximation introduced for the analysis of random walks on lattices with static disorder . Subsequently, this model has been generalized to include the possibility of multistate transformations of the dynamically random complex medium . More recently, several exactly solvable one-dimensional models with global and local dynamical disorder have been discussed .
The second approach to the problem emerged within the context of ionic conductivity in superionic solids. Here, the dynamical percolative environment has been considered as a multicomponent mixture of mobile species in which one or several neutral components block the carrier component . In particular, such a situation can be observed in a superionic conductor $`\beta ^{\prime \prime }`$-alumina, doped with two different ionic species (e.g. $`Na^+`$ and $`Ba^{2+}`$), where small $`Na^+`$ ions are rather mobile, while the larger $`Ba^{2+}`$ ions move essentially slower and temporarily block the $`Na^+`$ ions. Contrary to the previous line of thought, the dynamics of such a percolative environment has essential correlations, generated by hard-core exclusion interactions between the species involved. Moreover, a salient feature of these situations is that here the total number of the particles involved is conserved, i.e., dynamics of the environment is globally constrained by a certain ”conservation law”. The impact of this property on charge carrier transfer will become clearer as we proceed.
Further on, in Ref. the frequency-dependent ionic conductivity of the light species has been analyzed combining a continuous time random walk approach for the dynamical problem with an effective medium approximation describing the frozen environment of slow species. Next, as an explanation of the sharp increase of electrical conductivity transition in water-in-oil microemulsions when the volume fraction of water is increased towards a certain threshold value, in Refs. and it has been proposed that the charge carriers are not trapped in the finite water clusters, but rather a charge on a water globule can propagate by either hopping to a neighboring globule, when they approach each other, or via the diffusion of the host globule itself. This picture has been interpreted in terms of a model similar to that employed in Ref., with the only difference being that here the ”blockers” of Ref. play the role of the transient charge carriers. In the model of Refs. and , in which the host dynamics is influenced by spatial correlations and conservation of the number of the water globules involved, the conductivity depends on the rate of cluster rearrangement. Lastly, a similar problem of a carrier diffusion in an environment created by mobile hard-core lattice-gas particles has been analyzed in Ref. by using the developed dynamic bond percolation theory of Refs. and .
In this paper, we pursue our previous analysis of the charge carrier transfer in complex media considering a generalize dynamic percolation model which shares common features with both bond-fluctuating models of Refs. and those involving mobile blockers of Refs.. The system we consider consists of a host lattice, which here is a regular cubic lattice whose sites support at most a single occupancy, hard-core ”environment” particles, and a single, hard-core charge carrier particle subject to external electric field. The ”environment” particles move on the lattice by performing a random hopping between the neighboring lattice sites, which is constrained by the hard-core interactions, and may disappear from and re-appear (renewal processes) on the empty sites of the lattice with some prescribed rates. Contrary to the environment particles the charge carrier can not disappear spontaneously, and is subject to a constant external force $`\stackrel{}{E}`$. Hence, the carrier performs a biased random walk, which is constrained by the hard-core interactions with the ”environment” particles. In consequence, on may consider the carrier as some ”probe” designed to measure the response of the percolative environment to the internal perturbation or, in other words, the frictional properties of such a dynamical environment.
Now, a salient feature of our model, which makes it different to the previously proposed models of dynamic percolation, is that here the interaction between ”environment” particles and the carrier are included, such that the latter may influence the dynamics of the environment. This results, as we proceed to show, in the emergence of complicated density profiles of the ”environment” particles around the carrier. These profiles, as well as the terminal velocity $`V_c`$ of the carrier, are determined here explicitly, in terms of an approximate approach of Ref., which is based on the decoupling of the triple carrier-particle-particle correlation functions into the product of pair correlations.
We realize that the ”environment” particles tend to accumulate in front of the driven charge carrier creating a ”traffic jam”, which impedes its motion. The density profiles around the carrier are asymmetric: the local density of the ”environment” particles in front of the carrier is higher than the average and approaches the average value as an exponential function of the distance from the carrier. The characteristic length and the amplitude of the density relaxation function are calculated explicitly. On the other hand, behind the carrier, i.e. in its ”wake”, the local density is lower than the average. We find that, surprisingly, the functional form of the density profile might be very different depending on the condition whether the number of particles in the percolative environment is explicitly conserved or not; as a matter of fact, the local density behind the carrier may tend to the average value either as an exponential or even as an algebraic function of the distance, revealing in the latter case especially strong memory effects and strong correlations between the particle distribution in the environment and the carrier position. Further on, we find that the terminal velocity of the carrier particle depends explicitly on the excess density in the ”jammed” region in front of the carrier, as well as on the ”environment” particles density past the carrier. Both, in turn, are dependent on the magnitude of the velocity, as well as on the rate of the renewal processes and the rate at which the ”environment” particles can diffuse away from the carrier. The interplay between the jamming effect of the environment, produced by the carrier particle, and the rate of its homogenization due to diffusive smoothening and renewal processes, manifests itself as a medium-induced frictional force exerted on the carrier, whose magnitude depends on the carrier velocity. As a consequence of such a non-linear coupling, in the general case, (i.e. for arbitrary rates of the renewal and diffusive processes), $`V_c`$ can be found only implicitly, as the solution of a non-linear equation relating $`V_c`$ to the system parameters. This equation simplifies considerably in the limit of small applied external fields $`\stackrel{}{E}`$ and we find that the force-velocity relation to the field becomes linear. This implies that the frictional force exerted on the carrier particle by the environment is viscous. This linear force-velocity relation can be therefore interpreted as the analog of the Stokes formula for the dynamic percolative environment under study; in this case, the carrier velocity is calculated explicitly as well as the corresponding friction coefficient. Assuming that the Einstein relation between the carrier mobility and its diffusion coefficient holds , we then estimate the self-diffusion coefficient of the carrier in absence of external field. We show that when only diffusive re-arrangement of the percolative environment is allowed, while the renewal processes are suppressed, the general expression for the diffusion coefficient reduces to the one obtained previously in Refs. and . We note that the result of Refs. and is known to serve as a very good approximation for the self-diffusion coefficient in hard-core lattice-gases .
The paper is structured as follows: In Section II we formulate the model and introduce basic notations. In Section III we write down the dynamical equations which govern the time evolution of the ”environment” particles and of the carrier. Section IV is devoted to the analytical solution of these evolution equations in the limit $`t\mathrm{}`$. In Section V we derive explicit asymptotic results for the carrier terminal velocity in the limit of small applied external fields $`\stackrel{}{E}`$ and obtain the analog of the Stokes formula for such a percolative environment. Asymptotic behavior of the density profiles of the ”environment” particles around the carrier is discussed in Section VI. Finally, we conclude in Section VII with a brief summary and discussion of our results.
## 2 THE MODEL.
The model for charge carrier transfer in complex media consists of a three-dimensional simple cubic lattice of spacing $`\sigma `$, the sites of which are partially occupied by identical hard-core ”environment” particles and a single, hard-core, carrier particle. For both types of particles the hard-core interactions prevent multiple occupancy of the lattice sites; that is, no two ”environment” particles or the ”carrier” and an ”environment” particle can occupy simultaneously the same site, and particles can not pass through each other.
The occupation of the lattice sites by the ”environment” particles is characterized by the time-dependent occupation variable $`\eta (\stackrel{}{r})`$, $`\stackrel{}{r}`$ being the lattice-vector of the site in question. This variable assumes two values:
$$\eta (\stackrel{}{r})=\{\begin{array}{cc}1,\text{ if the site }\stackrel{}{r}\text{ is occupied}\hfill & \\ 0,\text{ if the site }\stackrel{}{r}\text{ is empty}\hfill & \end{array}$$
Next, the dynamics of the ”environment” particles is defined via the following rules The particles may, at a given rate, spontaneously disappear from the lattice, and may re-appear at random positions and random time moments, which is reminiscent of the host medium dynamics stipulated in Refs.. These two processes will be referred to, in general, as renewal processes.
Now, the environment particles move randomly by performing nearest-neighbor random walks constrained by the hard-core interactions, which is the main feature of the approach in Refs.. More specifically, we stipulate that any of the ”environment” particles waits a time $`\delta \tau `$, which has an exponential probability distribution with a mean $`\tau ^{}`$, and then chooses from a few possibilities: (a) disappearing from the lattice at rate $`g/\tau ^{}`$, which is realized instantaneously, or (b) attempting to hop, at rate $`l/6\tau ^{}`$, onto one of $`6`$ neighboring sites. The hop is actually fulfilled if the target site is not occupied at this time moment by any other particle; otherwise, the particle attempting to hop remains at its initial position, and (c) particles may re-appear on any vacant lattice site with rate $`f/\tau ^{}`$.
Note that, for simplicity, we assumed that the characteristic diffusion time and the renewal times of the ”environment” particles are equal to each other. These times, i.e. $`\tau _{dif}`$, mean creation time $`\tau _{cr}`$ and mean annihilation time $`\tau _{an}`$ may, however, be different, and can be restored in our final results by a mere replacement $`ll\tau ^{}/\tau _{dif}`$, $`ff\tau ^{}/\tau _{cr}`$ and $`gg\tau ^{}/\tau _{an}`$.
Note also that the number of particles is not explicitly conserved in such a dynamical model of the environment due to the renewal processes. However, in the absence of attractive particle-particle interactions and external perturbations, the particles distribution on the lattice is uniform and the average occupation $`\rho (t)=\overline{\eta (\stackrel{}{r})}`$ of the lattice tends, as $`t\mathrm{}`$, to a constant value, $`\rho _s=f/(f+g)`$. This relation can be thought of as the Langmuir adsorption isotherm .
Hence, the limit $`\tau _{dif}\mathrm{}`$ (or, $`l0`$) corresponds to the ordinary site percolation model with immobile blocked sites. The limit $`f,g0`$, ($`\tau _{cr},\tau _{an}\mathrm{}`$), while keeping the ratio $`f/g`$ fixed, $`f/g=\rho _s/(1\rho _s)`$, corresponds to the usual hard-core lattice-gas with the conserved particles number.
At time $`t=0`$ we introduce at the origin of the lattice the charge carrier, whose role is to probe the response of the environment modeled by dynamic percolation to an external perturbation. We stipulate that only the carrier out of all participating particles can not disappear from the system, and moreover, its motion is biased by some external constant force. As a physical realization, we envisage that the carrier is charged, while all other particles are neutral, and the system is exposed to constant external electric field $`\stackrel{}{E}`$.
The dynamics of the carrier particle is defined as follows: We suppose that the waiting time between successive jumps of the carrier has also an exponential distribution with a mean value $`\tau `$, which may in general be different from the corresponding waiting time of the environment particles. Attempting to hop, the carrier first chooses a hop direction with probabilities
$$p_\mu =exp\left[\frac{\beta }{2}\left(\stackrel{}{E}\stackrel{}{e}_\mu \right)\right]/\underset{\nu }{}exp\left[\frac{\beta }{2}\left(\stackrel{}{E}\stackrel{}{e}_\nu \right)\right],$$
(1)
where $`\beta `$ is the reciprocal temperature, $`\stackrel{}{e}_\nu `$ (or $`\stackrel{}{e}_\mu `$) stand for six unit lattice vectors, $`\nu ,\mu =\{\pm 1,\pm 2,\pm 3\}`$, connecting the carrier position with $`6`$ neighboring lattice sites, and $`(\stackrel{}{E}\stackrel{}{e}_\nu )`$ denotes the scalar product. We adopt the convention that $`\pm 1`$ corresponds to $`\pm X`$, $`\pm 2`$ corresponds to $`\pm Y`$ while $`\pm 3`$ stands for $`\pm Z`$. The jump is actually fulfilled when the target lattice site is vacant. Otherwise, the carrier remains at its position. For simplicity we assume in what follows that the external field is oriented along the $`X`$-axis in the positive direction, such that $`\stackrel{}{E}=(E,0,0)`$. Note also that for the choice of the transition probabilities as in Eq.(1), the detailed balance is naturally preserved.
## 3 EVOLUTION EQUATIONS.
We proceed by writing the evolution equations describing the dynamics of the system. Let $`P(\stackrel{}{R}_c,\eta ;t)`$ denote the joint probability that at time moment $`t`$ the charge carrier occupies position $`\stackrel{}{R}_c`$ and all ”environment” particles are in configuration $`\eta \{\eta (\stackrel{}{r})\}`$. Next, let $`\eta ^{\stackrel{}{r},\mu }`$ denote particles’ configuration obtained from $`\eta `$ by exchanging the occupation variables of the sites $`\stackrel{}{r}`$ and $`\stackrel{}{r}+\stackrel{}{e}_\stackrel{}{\mu }`$, i.e. $`\eta (\stackrel{}{r})\eta (\stackrel{}{r}+\stackrel{}{e}_\stackrel{}{\mu })`$, and $`\widehat{\eta }^\stackrel{}{r}`$ be the configuration obtained from $`\eta `$ by changing the occupation of the site $`\stackrel{}{r}`$ as $`\eta (\stackrel{}{r})1\eta (\stackrel{}{r})`$. Clearly, the first type of process appears due to random hops of the ”environment” particles, while the second one stems from the renewal processes, i.e. random creation and annihilation of the ”environment” particles. Then, summing up all possible events which can result in the configuration $`(\stackrel{}{R}_c,\eta )`$ or change this configuration for any other, we find that the temporal evolution of the system under study is governed by the following master equation:
$`_tP(\stackrel{}{R}_c,\eta ;t)={\displaystyle \frac{l}{6\tau ^{}}}{\displaystyle \underset{\mu }{}}{\displaystyle \underset{\stackrel{}{r}\stackrel{}{R}_c\stackrel{}{e}_\stackrel{}{\mu },\stackrel{}{R}_c}{}}\left\{P(\stackrel{}{R}_c,\eta ^{\stackrel{}{r},\mu };t)P(\stackrel{}{R}_c,\eta ;t)\right\}+`$ (2)
$`+`$ $`{\displaystyle \frac{1}{\tau }}{\displaystyle \underset{\mu }{}}p_\mu \left\{\left(1\eta (\stackrel{}{R}_c)\right)P(\stackrel{}{R}_c\stackrel{}{e}_\stackrel{}{\mu },\eta ;t)\left(1\eta (\stackrel{}{R}_c+\stackrel{}{e}_\stackrel{}{\mu })\right)P(\stackrel{}{R}_c,\eta ;t)\right\}+`$
$`+`$ $`{\displaystyle \frac{g}{\tau ^{}}}{\displaystyle \underset{\stackrel{}{r}\stackrel{}{R}_c}{}}\left\{\left(1\eta (\stackrel{}{r})\right)P(\stackrel{}{R}_c,\widehat{\eta }^\stackrel{}{r};t)\eta (\stackrel{}{r})P(\stackrel{}{R}_c,\eta ;t)\right\}+`$
$`+`$ $`{\displaystyle \frac{f}{\tau ^{}}}{\displaystyle \underset{\stackrel{}{r}\stackrel{}{R}_c}{}}\left\{\eta (\stackrel{}{r})P(\stackrel{}{R}_c,\widehat{\eta }^\stackrel{}{r};t)\left(1\eta (\stackrel{}{r})\right)P(\stackrel{}{R}_c,\eta ;t)\right\}.`$
Note that the terms in the first (resp. second) line of Eq.(2) describe random hopping motion of the ”environment” particles (resp. biased motion of the carrier) in terms of the Kawasaki-type particle-vacancy exchanges, while the terms in the third and the fourth lines account for the Glauber-type decay and creation of the ”environment” particles.
### 3.1 Mean velocity of the charge carrier and correlation functions.
The velocity of the charge carrier can be now readily determined from Eq.(2). To do it, we multiply both sides of Eq.(2) by $`(\stackrel{}{R}_c\stackrel{}{e}_1)`$ and sum over all possible configurations $`(\stackrel{}{R}_c,\eta )`$. This yields the following expression for the carrier mean velocity $`V_c(t)`$:
$$V_c(t)=\frac{\sigma }{\tau }\left\{p_1(1k(\stackrel{}{e}_1;t))p_1(1k(\stackrel{}{e}_1;t))\right\},$$
(3)
where $`k(\stackrel{}{\lambda };t)`$ denotes the carrier-”environment” particles pair correlation function:
$$k(\stackrel{}{\lambda };t)\underset{\stackrel{}{R}_c,\eta }{}\eta (\stackrel{}{R}_c+\stackrel{}{\lambda })P(\stackrel{}{R}_c,\eta ;t).$$
(4)
Note that $`k(\stackrel{}{\lambda };t)`$ is just the density distribution of the ”environment” particles, as seen by someone residing on the carrier which moves with velocity $`V_c(t)`$.
Note now that $`V_c(t)`$ depends explicitly on the local density of the ”environment” particles in the immediate vicinity of the carrier. Note also that if the ”environment” is perfectly homogeneous, i.e., $`k(\stackrel{}{\lambda };t)\rho _s`$, one has that $`\eta (\stackrel{}{R}_c+\stackrel{}{\lambda })`$ and $`P(\stackrel{}{R}_c,\eta ;t)`$ in Eq.(4) decouple, which entails, in virtue of from Eq.(3), a trivial mean-field-type result
$$V_c^{(0)}=(p_1p_1)(1\rho _s)\frac{\sigma }{\tau }.$$
(5)
The latter (trivial) expression implies that for the perfectly homogeneous, well-stirred environment the frequency of jumps of the carrier particles ($`\tau ^1`$) merely gets renormalized by a factor $`1\rho _s`$, which gives the fraction of successful jumps.
The situation, however, appears to be more complex and there emerge essential backflow effects: As a matter of fact, the carrier effectively perturbs the spatial distribution of the ”environment” particles so that stationary density profiles emerge. This can be contrasted to the earlier dynamic percolation models in which the carrier had no impact on the embedding medium and hence there was no re-arrangement of the host medium around the carrier particle. As a consequence, in our model $`k(\stackrel{}{\lambda };t)\rho _s`$, and $`k(\stackrel{}{\lambda };t)`$ approaches $`\rho _s`$ only at infinite separations from the carrier, i.e. when $`|\stackrel{}{\lambda }|\mathrm{}`$. Therefore, we rewrite Eq.(3) in the form
$$V_c(t)=V_c^{(0)}\frac{\sigma }{\tau }\{p_1(k(\stackrel{}{e}_1;t)\rho _s)p_1(\rho _sk(\stackrel{}{e}_1;t))\},$$
(6)
which shows explicitly the deviation of the mean velocity of the carrier from the mean-field-type result in Eq.(5) due to the formation of the density profiles.
### 3.2 Evolution equations of the pair correlation functions.
Equation (3) signifies that in order to obtain $`V_c(t)`$ we have to evaluate $`k(\stackrel{}{e}_{\pm 1};t)`$. Multiplying both sides of Eq.(2) by $`\eta (\stackrel{}{R}_c)`$ and summing over all configurations $`(\stackrel{}{R}_c,\eta )`$, we find that $`k(\stackrel{}{\lambda };t)`$ obeys
$`_tk(\stackrel{}{\lambda };t)`$ $`=`$ $`{\displaystyle \frac{l}{6\tau ^{}}}{\displaystyle \underset{\mu }{}}(_\mu \delta _{\stackrel{}{\lambda },\stackrel{}{e}_\mu }_\mu )k(\stackrel{}{\lambda };t){\displaystyle \frac{(f+g)}{\tau ^{}}}k(\stackrel{}{\lambda };t)+{\displaystyle \frac{f}{\tau ^{}}}+`$ (7)
$`+`$ $`{\displaystyle \frac{1}{\tau }}{\displaystyle \underset{\mu }{}}{\displaystyle \underset{\stackrel{}{R}_c,\eta }{}}p_\mu \left(1\eta (\stackrel{}{R}_c+e_\mu )\right)_\mu \eta (\stackrel{}{R}_c+\stackrel{}{\lambda })P(\stackrel{}{R}_c,\eta ;t),`$
where $`_\mu `$ denotes the ascending finite difference operator of the form
$$_\mu f(\stackrel{}{\lambda })f(\stackrel{}{\lambda }+\stackrel{}{e}_\mu )f(\stackrel{}{\lambda }),$$
(8)
and
$$\delta _{\stackrel{}{r},\stackrel{}{r}^{}}=\{\begin{array}{cc}1,\text{ if the site }\stackrel{}{r}=\stackrel{}{r}^{}\hfill & \\ 0,\text{ otherwise.}\hfill & \end{array}$$
The Kroneker-delta term $`\delta _{\stackrel{}{\lambda },\stackrel{}{e}_\mu }`$ signifies that the evolution of the pair correlations, Eq.(7), proceeds differently at large separations and in the immediate vicinity of the carrier. This stems from the asymmetric hopping rules of the carrier particle defined by Eq.(1).
Note next that the contribution in the second line in Eq.(7), which stems out of the bias acting on the carrier, is non-linear with respect to the occupation numbers. In consequence, the pair correlation function is effectively coupled to the evolution of the third-order correlations. That is, Eq.(7) is not closed with respect to the pair correlations but rather represents a first equation in the infinite hierarchy of coupled equations for higher-order correlation functions. One faces, therefore, the problem of solving an infinite hierarchy of coupled differential equations and needs to resort to an approximate closure scheme.
### 3.3 Decoupling Approximation.
Here we employ a simple non-trivial closure approximation, based on the decoupling of the third-order correlation functions into the product of pair correlations. More precisely, we assume that for $`\stackrel{}{\lambda }\stackrel{}{e}_\nu `$, the third-order correlation fulfil
$$\underset{\stackrel{}{R}_c,\eta }{}\eta (\stackrel{}{R}_c+\stackrel{}{\lambda })\eta (\stackrel{}{R}_c+\stackrel{}{e}_\nu )P(\stackrel{}{R}_c,\eta ;t)k(\stackrel{}{\lambda };t)k(\stackrel{}{e}_\nu ;t)$$
(9)
The approximate closure in Eq.(9) has been already used for studying related models of biased carrier diffusion in hard-core lattice gases and has been shown to provide quite an accurate description of both the dynamical and stationary-state behavior. The decoupling in Eq.(9) was first introduced in Ref. to determine the properties of a driven carrier diffusion in a one-dimensional hard-core lattice gas with a conserved number of particles, i.e. without an exchange of particles with the reservoir. Extensive numerical simulations performed in Ref. have demonstrated that such a decoupling is quite a plausible approximation for the model under study. Moreover, rigorous probabilistic analysis of Ref. has shown that for this model the results based on the decoupling scheme in Eq.(9) are exact. Furthermore, the same closure procedure has been recently applied to study spreading of a hard-core lattice gas from a reservoir attached to one of the lattice sites . Again, a very good agreement between the analytical results and the numerical data has been found. Next, the decoupling in Eq.(9) has been used in a recent analysis of a biased carrier dynamics in a one-dimensional model of an adsorbed monolayer in contact with a vapor phase , i.e. a one-dimensional version of the model to be studied here. Also in this case an excellent agreement has been observed between the analytical predictions and the Monte Carlo simulations data . We now show that the approximate closure of the hierarchy of the evolution equations in Eq.(9) allows us to reproduce in the limit $`f,g=0`$ and $`f/g=const`$ the results of Refs. and , which are known (see e.g. Ref.) to provide a very good approximation for the carrier diffusion coefficient in three-dimensional hard-core lattice gases with arbitrary particle density. We expect therefore that such a closure scheme will render a plausible description of the carrier dynamics in our three-dimensional generalized dynamic percolation model. We base our further analysis on this approximation.
Making use of Eq.(9), we find from Eq.(7) that the time evolution of pair correlations is governed by the following equations. For all $`\stackrel{}{\lambda }`$ except for $`\stackrel{}{\lambda }=\{\mathrm{𝟎},\stackrel{}{e}_{\pm 1},\stackrel{}{e}_{\pm 2},\stackrel{}{e}_{\pm 3}\}`$, we have that $`k(\stackrel{}{\lambda };t)`$ obeys
$$_tk(\stackrel{}{\lambda };t)=\frac{l}{6\tau ^{}}\stackrel{~}{L}k(\stackrel{}{\lambda };t)+\frac{f}{\tau ^{}},$$
(10)
where the operator $`\stackrel{~}{L}`$ and coefficients $`A_\nu (t)`$ are given explicitly by
$$\stackrel{~}{L}\underset{\mu }{}A_\mu (t)_\mu \frac{6(f+g)}{l},$$
(11)
and
$$A_\mu (t)1+\frac{6\tau ^{}}{l\tau }p_\mu (1k(\stackrel{}{e}_\mu ;t)),$$
(12)
the operator $`_\mu `$ being defined previously in Eq.(8), $`\mu =\{\pm 1,\pm 2,\pm 3\}`$. On the other hand, at the sites adjacent to the carrier one has
$$_tk(\stackrel{}{e}_\nu ;t)=\frac{l}{6\tau ^{}}\left(\stackrel{~}{L}+A_\nu (t)\right)k(\stackrel{}{e}_\nu ;t)+\frac{f}{\tau ^{}},$$
(13)
where $`\nu =\{\pm 1,\pm 2,\pm 3\}`$.
Now, several comments about equations (10) and (13) are in order. First of all, let us note that Eq.(13) represents, from the mathematical point of view, the boundary conditions for the general evolution equation (10), imposed on the sites in the immediate vicinity of the carrier. Equations (10) and (13) have a different form since in the immediate vicinity of the carrier its asymmetric hopping rules perturb essentially the ”environment” particles dynamics. Equations (10) and (13) possess some intrinsic symmetries and hence the number of independent parameters can be reduced. Namely, reversing the field, i.e. changing $`EE`$, leads to the mere replacement of $`k(\stackrel{}{e}_1;t)`$ by $`k(\stackrel{}{e}_1;t)`$ but does not affect $`k(\stackrel{}{e}_\nu ;t)`$ with $`\nu =\{\pm 2,\pm 3\}`$, which implies that
$$k(\stackrel{}{e}_1;t)(E)=k(\stackrel{}{e}_1;t)(E),\text{and}k(\stackrel{}{e}_\nu ;t)(E)=k(\stackrel{}{e}_\nu ;t)(E)\text{for}\nu =\{\pm 2,\pm 3\}.$$
(14)
Besides, since the transition probabilities in Eq.(1) obey $`p_2=p_2=p_3=p_3`$, one evidently has that $`k(\stackrel{}{e}_2;t)=k(\stackrel{}{e}_2;t)=k(\stackrel{}{e}_3;t)=k(\stackrel{}{e}_3;t)`$, and hence, by symmetry, $`A_2(t)=A_2(t)=A_3(t)=A_3(t)`$, which somewhat simplifies equations (10) and (13). Lastly, we note that despite the fact that using the decoupling scheme in Eq.(9) we effectively close the system of equations on the level of the pair correlations, the solution of Eqs.(10) and (13) still poses serious technical difficulties. Namely, these equations are strongly non-linear with respect to the carrier velocity, which introduces the gradient term on the rhs of the evolution equations for the pair correlation, and depends by itself on the values of the ”environment” particles densities in the immediate vicinity of the carrier. Below we discuss a solution to this non-linear problem, focusing on the limit $`t\mathrm{}`$.
## 4 SOLUTION OF THE DECOUPLED EVOLUTION EQUATIONS IN THE STATIONARY STATE.
Consider the limit $`t\mathrm{}`$ and suppose that the density profiles and the stationary velocity of the carrier have non-trivial stationary values: $`k(\stackrel{}{\lambda })lim_t\mathrm{}k(\stackrel{}{\lambda };t)`$, $`V_clim_t\mathrm{}V_c(t)`$ and $`A_\mu =lim_t\mathrm{}A_\mu (t)`$. As the next step, we define the local deviations of $`k(\stackrel{}{\lambda })`$ from the unperturbed density: $`h(\stackrel{}{\lambda })k(\stackrel{}{\lambda })\rho _s`$. This yields the following system of equations:
$$\stackrel{~}{L}h(\stackrel{}{\lambda })=0,$$
(15)
for $`\stackrel{}{\lambda }\{\mathrm{𝟎},\stackrel{}{e}_{\pm 1},\stackrel{}{e}_{\pm 2},\stackrel{}{e}_{\pm 3}\}`$, while for the special sites adjacent to the carrier one has
$$(\stackrel{~}{L}+A_\nu )h(\stackrel{}{e}_\nu )+\rho _s(A_\nu A_\nu )=0,$$
(16)
Equations (15) and (16) determine the spatial distribution of the deviation from the unperturbed density $`\rho _s`$ in the stationary state. Note also that in virtue of the symmetry relations $`h(\stackrel{}{e}_{\pm 2})=h(\stackrel{}{e}_{\pm 3})`$ and $`A_2=A_2=A_3=A_3`$.
To solve the coupled non-linear Eqs.(3),(15) and (16) we proceed in the following, standard manner: We first solve these equations supposing that the carrier stationary velocity is a fixed, given parameter (and hence, the functions $`A_\nu `$ entering Eqs.(15) and (16) are known). In doing so, we obtain $`h(\lambda )`$ in the parameterized form $`h(\stackrel{}{\lambda })=h(\stackrel{}{\lambda };A_{\pm 1},A_2)`$. Then, choosing particular values $`\stackrel{}{\lambda }=\{\stackrel{}{e}_{\pm 1},\stackrel{}{e}_{\pm 2},\stackrel{}{e}_{\pm 3}\}`$ and making use of the definition of $`A_\mu `$, we find a system of three linear equations with three unknowns of the form
$$A_\nu =1+\frac{6\tau ^{}}{l\tau }p_\nu \left(1\rho _sh(\stackrel{}{e}_\nu ;A_{\pm 1},A_2)\right),$$
(17)
where $`\nu =\{\pm 1,2\}`$, which will allow us to obtain a closure relation and hence, to define all $`A_\nu `$ explicitly (and hence, all $`h(\stackrel{}{e}_\nu ))`$. Finally, substituting the results into Eq.(3), which can be written down in terms of $`A_\nu `$ as
$$V_c=\frac{l\sigma }{6\tau ^{}}(A_1A_1),$$
(18)
we arrive at a closed-form equation determining implicitly the stationary velocity.
### 4.1 Density profiles in the dynamic percolative environment.
The general solution of Eqs.(15) and (16) can be obtained in a standard fashion by introducing the following generating function:
$$H(w_1,w_2,w_3)\underset{n_1,n_2,n_3}{}h(\stackrel{}{\lambda })w_1^{n_1}w_2^{n_2}w_3^{n_3},$$
(19)
where $`n_1`$,$`n_2`$ and $`n_3`$ are the components of the vector $`\stackrel{}{\lambda }`$, $`\stackrel{}{\lambda }=\stackrel{}{e}_1n_1+\stackrel{}{e}_2n_2+\stackrel{}{e}_3n_3`$. Multiplying both sides of Eqs.(15) and (16) by $`w_1^{n_1}w_2^{n_2}w_3^{n_3}`$ and performing summation, we find then that $`H(w_1,w_2,w_3)`$ is given explicitly by
$$H(w_1,w_2,w_3)=l\frac{_\nu \left(A_\nu (w_{|\nu |}^{\nu /|\nu |}1)h(\stackrel{}{e}_\nu )+\rho _s(A_\nu A_\nu )w_{|\nu |}^\nu \right)}{l_\nu A_\nu (w_{|\nu |}^{\nu /|\nu |}1)6(f+g)},$$
(20)
an expression which allows us to determine the stationary density profiles as seen from the carrier which moves with a constant velocity $`V_c`$.
Inverting next the generating function, Eq.(20), we get, after rather lengthy but straightforward calculations, the following explicit result for the local deviation from the unperturbed density:
$`h(\stackrel{}{\lambda })=\alpha ^1\left\{{\displaystyle \underset{\nu }{}}A_\nu h(\stackrel{}{e}_\nu )_\nu \rho _s(A_1A_1)(_1_1)\right\}F(\stackrel{}{\lambda }),`$ (21)
where $`F(\stackrel{}{\lambda })`$ is given by
$`F(\stackrel{}{\lambda })=\left({\displaystyle \frac{A_1}{A_1}}\right)^{n_1/2}{\displaystyle _0^{\mathrm{}}}e^x\mathrm{I}_{n_1}\left(2{\displaystyle \frac{\sqrt{A_1A_1}}{\alpha }}x\right)\mathrm{I}_{n_2}\left(2{\displaystyle \frac{A_2}{\alpha }}x\right)\mathrm{I}_{n_3}\left(2{\displaystyle \frac{A_2}{\alpha }}x\right)dx,`$ (22)
and
$`\alpha ={\displaystyle \underset{\nu }{}}A_\nu +{\displaystyle \frac{6(f+g)}{l}}=A_1+A_1+4A_2+{\displaystyle \frac{6(f+g)}{l}}`$ (23)
Consequently, the particles density distribution as seen from the carrier moving with a constant velocity $`V_c`$ obeys
$`k(\stackrel{}{\lambda })=\rho _s+\alpha ^1\{{\displaystyle \underset{\nu }{}}A_\nu h(\stackrel{}{e}_\nu _\nu \rho _s(A_1A_1)(_1_1)\}F(\stackrel{}{\lambda }),`$ (24)
where we have to determine three yet unknown parameters $`A_1`$, $`A_1`$ and $`A_2`$.
To determine these parameters, we set in Eq.(21) $`\stackrel{}{\lambda }=\stackrel{}{e}_1`$, $`\stackrel{}{\lambda }=\stackrel{}{e}_1`$ and $`\stackrel{}{\lambda }=\stackrel{}{e}_2`$, which results in the system of three closed-form equations determining the unknown functions $`A_\nu `$, $`\nu =\{\pm 1,2\}`$,
$$A_\nu =1+\frac{6\tau ^{}}{l\tau }p_\nu \left\{1\rho _s\rho _s(A_1A_1)\frac{det\stackrel{~}{C}_\nu }{det\stackrel{~}{C}}\right\},$$
(25)
where $`\stackrel{~}{C}`$ is a square matrix of the third order defined as
$$\left(\begin{array}{ccc}A_1_1F(\stackrel{}{e}_1)\alpha & A_1_1F(\stackrel{}{e}_1)& A_2_2F(\stackrel{}{e}_1)\\ A_1_1F(\stackrel{}{e}_1)& A_1_1F(\stackrel{}{e}_1)\alpha & A_2_2F(\stackrel{}{e}_1)\\ A_1_1F(\stackrel{}{e}_2)& A_1_1F(\stackrel{}{e}_2)& A_2_2F(\stackrel{}{e}_2)\alpha \end{array}\right)$$
(26)
while $`\stackrel{~}{C}_\nu `$ stands for the matrix obtained from $`\stackrel{~}{C}`$ by replacing the $`\nu `$-th column by a column vector $`\left((_1_1)F(\stackrel{}{e}_\nu )\right)_\nu `$. Equation (24), together with the definition of the coefficients $`A_\nu `$, constitutes the first general result of our analysis defining the density distribution in the percolative environment under study.
### 4.2 General force-velocity relation.
Substituting Eq.(25) into (18), we find that the stationary velocity of the carrier particle is defined implicitly as the solution of equation:
$$V_c=\frac{\sigma }{\tau }(p_1p_1)(1\rho _s)\left\{1+\rho _s\frac{6\tau ^{}}{l\tau }\frac{p_1det\stackrel{~}{C}_1p_1det\stackrel{~}{C}_1}{det\stackrel{~}{C}}\right\}^1,$$
(27)
where $`\stackrel{~}{C}_1`$ and $`\stackrel{~}{C}_1`$ are the following square matrices of the third order:
$$\stackrel{~}{C}_1=\left(\begin{array}{ccc}(_1_1)F(\stackrel{}{e}_1)& A_1_1F(\stackrel{}{e}_1)& A_2_2F(\stackrel{}{e}_1)\\ (_1_1)F(\stackrel{}{e}_1)& A_1_1F(\stackrel{}{e}_1)\alpha & A_2_2F(\stackrel{}{e}_1)\\ (_1_1)F(\stackrel{}{e}_2)& A_1_1F(\stackrel{}{e}_2)& A_2_2F(\stackrel{}{e}_2)\alpha \end{array}\right)$$
(28)
and
$$\stackrel{~}{C}_1=\left(\begin{array}{ccc}A_1_1F(\stackrel{}{e}_1)\alpha & (_1_1)F(\stackrel{}{e}_1)& A_2_2F(\stackrel{}{e}_1)\\ A_1_1F(\stackrel{}{e}_1)& (_1_1)F(\stackrel{}{e}_1)& A_2_2F(\stackrel{}{e}_1)\\ A_1_1F(\stackrel{}{e}_2)& (_1_1)F(\stackrel{}{e}_2)& A_2_2F(\stackrel{}{e}_2)\alpha \end{array}\right).$$
(29)
Equation (27) represents our second principal result defining the force-velocity relation in the dynamic percolative environment for an arbitrary field and arbitrary rates of the diffusive and renewal processes.
## 5 CARRIER VELOCITY, FRICTION AND DIFFUSION COEFFICIENTS.
Consider now the case when the applied external field $`E`$ is small. Expanding the transition probabilities $`p_1`$ and $`p_1`$ in the Taylor series up to the first order in powers of the external field, i.e., setting
$$p_{\pm 1}=\frac{1}{6}\pm \frac{\sigma \beta E}{12}+𝒪\left(E^2\right),$$
(30)
we find that $`V_c`$, Eq.(18), follows
$$V_c\frac{\sigma }{6\tau }\left\{\sigma \beta E(1\rho _s)(h(\stackrel{}{e}_1)h(\stackrel{}{e}_1))\right\}.$$
(31)
On the other hand, Eq.(21) entails that
$$h(\stackrel{}{e}_1)h(\stackrel{}{e}_1)=\frac{2\sigma \rho _s(1\rho _s)\tau ^{}}{l\tau \left(\alpha _0(2A_0/\alpha _0)A_0\right)+2\rho _s\tau }\beta E+𝒪\left(E^2\right),$$
(32)
where
$$A_0=\underset{E0}{lim}A_\nu =1+\frac{\tau ^{}}{l\tau }(1\rho _s),$$
(33)
and
$$\alpha _0=\underset{E0}{lim}\alpha =6\left(1+\frac{\tau ^{}(1\rho _s)}{l\tau }+\frac{f+g}{l}\right),$$
(34)
while
$`(x)\left\{{\displaystyle _0^{\mathrm{}}}e^t\mathrm{I}_0^2(xt)\left(\mathrm{I}_0(xt)\mathrm{I}_2(xt)\right)dt\right\}^1=\left\{P(\mathrm{𝟎};3x)P(2\stackrel{}{e}_1;3x)\right\}^1,`$ (35)
with $`P(\stackrel{}{r};\xi )`$ being the generating function,
$$P(\stackrel{}{r};\xi )\underset{j=0}{\overset{+\mathrm{}}{}}P_j(\stackrel{}{r})\xi ^j,$$
(36)
of the probability $`P_j(\stackrel{}{r})`$ that a walker starting at the origin and performing a Polya random walk on the sites of a three-dimensional cubic lattice will arrive on the $`j`$-th step to the site with the lattice vector $`\stackrel{}{r}`$ .
Consequently, we find that in the limit of a small applied field $`E`$ the force-velocity relation in Eq.(27) attains the physically meaningful form of the Stokes formula $`E=\zeta V_c`$, which signifies that the frictional force exerted on the carrier by the environment particles is viscous. The effective friction coefficient $`\zeta `$ is the sum of two terms,
$$\zeta =\zeta _0+\zeta _{coop}$$
(37)
where the first term represents a mean-field-type result $`\zeta _0=6\tau /\beta \sigma ^2(1\rho _s)`$ (see Eq.(7)), while the second one, $`\zeta _{coop}`$, obeys
$$\zeta _{coop}=\frac{12\rho _s\tau ^{}}{\beta \sigma ^2l(1\rho _s).\left(\alpha _0(2A_0/\alpha _0)A_0\right)}$$
(38)
The second contribution has a more complicated origin and is associated with the cooperative behavior - formation of a inhomogeneous stationary particle distribution around the carrier moving with constant velocity $`V_c`$. Needless to say, such an effect can not be observed within the framework of previous models of dynamic percolation, since there the carrier does not influence the host medium dynamics .
Let us now compare the relative importance of two contributions, i.e. $`\zeta _0`$ and $`\zeta _{coop}`$, to the overall friction. Straightforward analysis shows that the cooperative behavior dominates at small and moderate $`f`$ (which entails also small values of $`g`$), while for larger $`f`$, when $`\zeta /\zeta _0`$ tends to $`1`$, the mean-field behavior becomes most important. The cooperative behavior also appears to be more pronounced at larger densities $`\rho _s`$.
Consider next some analytical estimates. We start with the situation, in which diffusion of the environment particles is suppressed, i.e. when $`l=0`$. In this case, we get
$$\frac{\zeta _{coop}}{\zeta _0}=\frac{2\rho _s}{(1\rho _s)\left(\frac{2}{y}(y)1\right)},$$
(39)
where
$$y=\frac{1}{3}\left(1+\frac{\tau }{\tau ^{}}\frac{(f+g)}{(1\rho _s)}\right)^1.$$
(40)
Suppose first that $`\rho _s`$ is small, $`\rho _s1`$. Then, $`y1/3(1+\tau /\tau ^{}(f+g))`$ and we can distinguish between two situations: when $`\tau (f+g)/\tau ^{}`$, i.e. when the carrier moves faster than the environment re-organizes itself, and the opposite limit, $`\tau (f+g)/\tau ^{}`$, when the environment changes very rapidly compared to the motion of the carrier. In the former case we find that $`y1/3`$, which yields $`\zeta _{coop}/\zeta _02\rho _s/(6(1/3)1)`$, $`(1/3)0.7942`$, while in the latter case we have $`y\tau ^{}/3\tau (f+g)`$ and $`\zeta _{coop}/\zeta _0\rho _s\tau ^{}/3\tau (f+g)`$. Note, that in both cases the ratio $`\zeta _{coop}/\zeta _0`$ appears to be small, which signifies that at small densities $`\rho _s`$ the mean-field friction dominates. Such a result is not counterintuitive, of course, since in the absence of the particles’ diffusion, which couples effectively the density evolution at different lattice sites, no significant cooperative behavior can emerge at small densities. On the other hand, at relatively high densities $`\rho _s1`$ and $`\tau /(1\rho _s)\tau ^{}/(f+g)\tau `$, when the carrier moves at much faster rate than the host medium reorganizes itself, we find that $`\zeta _{coop}/\zeta _0\tau ^{}/3\tau (f+g)1`$. This result stems from the circumstance that in sufficiently dense environments modeled by dynamic percolation a highly inhomogeneous density profile emerges even in the absence of particles diffusion. Here, on the one hand, the carrier perturbs significantly the particle density in its immediate vicinity. On the other hand, the density perturbation created by the carrier does not shift the global balance between creation and annihilation events, i.e. the mean particle density still equals $`\rho _s`$. The latter constraint induces then appearance of essential correlations in particles distribution and hence, appearance of cooperative behavior.
Let us consider the opposite case when the renewal processes are not allowed, which means that the particles number is conserved and local density in the percolative environment evolves only due to particles diffusion. In this case we find
$$\frac{\zeta _{coop}}{\zeta _0}=\frac{2\tau ^{}\rho _s}{(l\tau +\tau ^{}(1\rho _s))\left(6(1/3)1\right)}$$
(41)
Here, the ratio $`\zeta _{coop}/\zeta _0`$ can be large and the ”cooperative” friction dominates the mean-field one when $`l\tau \tau ^{}(3\rho _s1)`$, which happens, namely, at sufficiently high densities and in the limit when the carrier moves at a much faster rate than the environment reorganizes itself. Otherwise, the mean-field friction prevails.
To estimate the carrier particle diffusion coefficient $`D_c`$ we assume the validity of the Einstein relation, i.e. $`\beta D_c=\zeta ^1`$ (see, e.g., Ref.). We find that, in the general case, the carrier diffusion coefficient $`D_c`$ reads
$$D_c=\frac{\sigma ^2(1\rho _s)}{6\tau }\left\{1\frac{2\rho _s\tau ^{}}{l\tau }\left(\alpha _0(2A_0/\alpha _0)1+\frac{\tau ^{}(3\rho _s1)}{l\tau }\right)^1\right\}$$
(42)
In the particular case of conserved particles number, when $`f,g0`$ but their ratio $`f/g`$ is kept fixed, $`f/g=\rho _s/(1\rho _s)`$, the latter equation reduces to the classical result
$$D_c^{NK}=\frac{\sigma ^2(1\rho _s)}{6\tau }\left\{1\frac{2\rho _s\tau ^{}}{l\tau }\left(6A_0(1/3)1+\frac{\tau ^{}(3\rho _s1)}{l\tau }\right)^1\right\},$$
(43)
obtained earlier in Refs. and by different analytical techniques. The result in Eq.(43) is known to be exact in the limits $`\rho _s1`$ and $`\rho _s1`$, and serves as a very good approximation for the self-diffusion coefficient in hard-core lattice gases of arbitrary density .
Finally, in the absence of particle diffusion (fluctuating-site percolation), our result for the carrier particle diffusion coefficient reduces to
$$D_c^{per}=\frac{\sigma ^2(1\rho _s)}{6\tau }\left\{12\rho _s\left(4[(1\rho _s)+(f+g)\tau /\tau ^{}](y)+3\rho _s1\right)^1\right\}$$
(44)
Note, however, that this result only applies when both $`f`$ and $`g`$ are larger than zero, such that the renewal processes take place. In fact, the underlying decoupling scheme is only plausible in this case. Similarly to the approximate theories in Refs. and , our approach predicts that in the absence of the renewal processes $`D_c^{per}`$ vanishes only when $`\rho _s1`$, which is an incorrect behavior.
## 6 ASYMPTOTIC BEHAVIOR OF THE DENSITY PROFILES.
The density profiles at large separations in front of and past the carrier can be readily deduced from the asymptotical behavior of the following generating function
$$N(w_1)\underset{n_1=\mathrm{}}{\overset{+\mathrm{}}{}}h(n_1,n_2=0,n_3=0)w_1^{n_1}.$$
(45)
Inversion of Eq.(20) with respect to the symmetric coordinates $`n_2`$ and $`n_3`$ yields then
$`N(w_1)={\displaystyle \frac{\left(A_1h(\stackrel{}{e}_1)+\rho _s(A_1A_1)\right)\left(w_11\right)+\left(A_1h(\stackrel{}{e}_1)\rho _s(A_1A_1)\right)\left(w_1^11\right)}{\alpha A_1w_1^1A_1w_1}}\times `$ (46)
$`\times `$ $`{\displaystyle _0^{\mathrm{}}}\mathrm{exp}[x]\mathrm{I}_0^2({\displaystyle \frac{2A_2}{\alpha A_1w_1^1A_1w_1}}x)\mathrm{d}x+{\displaystyle \frac{4A_2h(\stackrel{}{e}_2)}{\alpha A_1w_1^1A_1w_1}}\times `$
$`\times `$ $`{\displaystyle _0^{\mathrm{}}}\mathrm{exp}[x]\mathrm{I}_0({\displaystyle \frac{2A_2}{\alpha A_1w_1^1A_1w_1}}x)(\mathrm{I}_1({\displaystyle \frac{2A_2}{\alpha A_1w_1^1A_1w_1}}x)`$
$``$ $`\mathrm{I}_0({\displaystyle \frac{2A_2}{\alpha A_1w_1^1A_1w_1}}x))\mathrm{d}x`$
We notice now that $`N(w_1)`$ is a holomorphic function in the region $`𝒲_1<w_1<𝒲_2`$, where
$`𝒲_1={\displaystyle \frac{\alpha 4A_2}{2A_1}}\sqrt{\left({\displaystyle \frac{\alpha 4A_2}{2A_1}}\right)^2{\displaystyle \frac{A_1}{A_1}}}`$ (47)
and
$`𝒲_2={\displaystyle \frac{\alpha 4A_2}{2A_1}}+\sqrt{\left({\displaystyle \frac{\alpha 4A_2}{2A_1}}\right)^2{\displaystyle \frac{A_1}{A_1}}}`$ (48)
As a consequence, the asymptotic behavior of $`h(n_1,n_2=0,n_3=0)`$ in the limit $`n_1\mathrm{}`$ (resp. $`n_1\mathrm{}`$) is controlled by the behavior of $`N(w_1)`$ in the vicinity of $`w_1=𝒲_2`$ (resp. $`w_1=𝒲_1`$) (see, for example, the analysis of the generating function singularities developed in Ref.).
### 6.1 Asymptotics of the density profiles at large separations in front of the carrier.
Consider first the asymptotic behavior of the density distribution of the ”environment” particles at large separations in front of the carrier. We find then that in the limit $`w_1𝒲_2`$, the function $`N(w_1)`$ follows
$`N(w_1)`$ $``$ $`{}_{w_1𝒲_2}{}^{}[{\displaystyle \frac{\left(A_1h(\stackrel{}{e}_1)+\rho _s(A_1A_1)\right)\left(𝒲_21\right)}{4\pi A_2}}+`$ (49)
$`+`$ $`{\displaystyle \frac{\left(A_1h(\stackrel{}{e}_1)\rho _s(A_1A_1)\right)\left(𝒲_2^11\right)}{4\pi A_2}}\left]\mathrm{ln}\right(𝒲_2w_1)`$
Then, (cf, Ref.), we obtain the following asymptotical result
$$h(n_1,0,0)_{n_1\mathrm{}}\frac{K^+}{n_1}e^{n_1/\lambda _+},$$
(50)
where the characteristic length $`\lambda _+`$ is given explicitly by
$$\lambda _+\mathrm{ln}^1\left(\frac{\alpha /22A_2}{A_1}+\sqrt{\left(\frac{\alpha /22A_2}{A_1}\right)^2\frac{A_1}{A_1}}\right),$$
(51)
and the amplitude $`K^+`$ obeys
$`K^+`$ $`=`$ $`[{\displaystyle \frac{\left(A_1h(\stackrel{}{e}_1)+\rho _s(A_1A_1)\right)\left(𝒲_21\right)}{4\pi A_2}}+`$ (52)
$`+`$ $`{\displaystyle \frac{(A_1h(\stackrel{}{e}_1)\rho _s(A_1A_1)(𝒲_2^11)}{4\pi A_2}}]>0,`$
which signifies that the density of the ”environment” particles in front of the carrier is higher than the average value $`\rho _s`$ and approaches $`\rho _s`$ at large separations from the carrier as an exponential function of the distance.
### 6.2 Asymptotics of the density profiles at large separations behind the carrier.
We consider next the asymptotic behavior of the ”environment” particles density profiles past the carrier particle, which turns out to be very different depending on whether the dynamics of the percolative environment obeys the strict conservation of the ”environment” particles number or not (the renewal processes are suppressed or allowed).
#### 6.2.1 Non-conserved particles number.
In the case when particles may disappear and re-appear on the lattice, one has that the root $`𝒲_1<1`$. We find then, following essentially the same lines as in the previous subsection, that
$`N(w_1)`$ $``$ $`{}_{w_1𝒲_1}{}^{}[{\displaystyle \frac{\left(A_1h(\stackrel{}{e}_1)+\rho _s(A_1A_1)\right)\left(𝒲_11\right)}{4\pi A_2}}+`$ (53)
$`+`$ $`{\displaystyle \frac{\left(A_1h(\stackrel{}{e}_1)\rho _s(A_1A_1)\right)\left(𝒲_1^11\right)}{4\pi A_2}}\left]\mathrm{ln}\right({\displaystyle \frac{1}{w_1𝒲_1}}).`$
Hence, in the non-conserved case the approach to the unperturbed value $`\rho _s`$ is also exponential when $`n_1\mathrm{}`$, and follows
$$h_{n_1,0,0}_{n_1\mathrm{}}\frac{K^{}}{|n_1|}e^{|n_1|/\lambda _{}},$$
(54)
where
$$\lambda _{}\mathrm{ln}^1\left(\frac{\alpha /22A_2}{A_1}\sqrt{\left(\frac{\alpha /22A_2}{A_1}\right)^2\frac{A_1}{A_1}}\right)$$
(55)
and
$`K^{}`$ $`=`$ $`[{\displaystyle \frac{\left(A_1h(\stackrel{}{e}_1)+\rho _s(A_1A_1)\right)\left(𝒲_11\right)}{4\pi A_2}}+`$ (56)
$`+`$ $`{\displaystyle \frac{\left(A_1h(\stackrel{}{e}_1)\rho _s(A_1A_1)\right)\left(𝒲_1^11\right)}{4\pi A_2}}]<0`$
which implies that the particles density past the carrier is lower than the average. Note that, in the general case, $`\lambda _+<\lambda _{}`$, which means that the depleted region past the carrier is more extended in space than the traffic-jam-like region in front of the carrier. The density profiles are therefore asymmetric with respect to the origin, $`n_1=0`$.
#### 6.2.2 Conserved particles number.
Finally, we turn to the analysis of the shape of the density profiles of the percolative environment behind the carrier in the particular limit when the host medium evolves only due to diffusion, while creation and annihilation of particles are completely suppressed. In this case, in which the particles number is explicitly conserved, one has that for arbitrary value of the field and particles’ average density, the root $`𝒲_11`$ and, consequently, the form of the generating function is qualitatively different from that in Eqs.(49) and (53),
$`N(w_1)_{w_11^+}\left[{\displaystyle \frac{(A_1h(\stackrel{}{e}_1)A_1h(\stackrel{}{e}_1)}{4\pi A_2}}+{\displaystyle \frac{2\rho _s(A_1A_1))}{4\pi A_2}}\right](w_11)\mathrm{ln}\left({\displaystyle \frac{1}{w_11}}\right).`$ (57)
Equation (57) implies that in the limit when the particle number is conserved the large-$`n_1`$ asymptotic behavior of $`h_{n_1,0,0}`$ is described by an algebraic function of $`n_1`$ with a logarithmic correction; that is,
$$h_{n_1,0,0}\frac{K_{}\mathrm{ln}(|n_1|)}{n_1^2},$$
(58)
where $`K_{}`$ is an $`n_1`$-independent constant. Remarkably, the power-law decay of correlations implies existence of a quasi-long-range order in the percolative environment past the carrier. In the conserved case the mixing of the three-dimensional percolative environment is not very efficient and there are considerable memory effects - the host medium remembers the passage of the carrier on large space and time scales.
## 7 CONCLUSIONS
To conclude, we have presented a microscopic model describing the dynamics of a charge carrier, driven by an external field $`\stackrel{}{E}`$ in a three-dimensional complex medium modeled by dynamic percolation, i. e. represented as a cubic lattice partially filled with mobile, hard-core ”environment” particles which can spontaneously disappear and reappear (renewal processes) in the system with some prescribed rates. Our analytical description of the transfer process in such a medium has been based on the master equation, describing the time evolution of the system, which has allowed us to evaluate a system of coupled dynamical equations for the charge carrier velocity and a hierarchy of correlation functions. To solve these coupled equations, we have invoked an approximate closure scheme based on the decomposition of the third-order correlation functions into a product of pairwise correlations, which has been first introduced in Ref. for a related model of a driven carrier dynamics in a one-dimensional lattice gas with conserved particles number. Within the framework of this approximation, we have derived a system of coupled, discrete-space equations describing evolution of the density profiles of the environment, as seen from the moving charge carrier, and its velocity $`V_c`$. We have shown that $`V_c`$ depends on the density of the ”environment” particles in front of and past the carrier. Both densities depend on the magnitude of the velocity, as well as on the rate of the renewal and diffusive processes. As a consequence of such a non-linear coupling, in the general case, (i.e. for an arbitrary driving field and arbitrary rates of renewal and diffusive processes), $`V_c`$ has been found only implicitly, as the solution of a non-linear equation relating its value to the system parameters. This equation, which defines the force-velocity relation for the dynamic percolation under study, simplifies considerably in the limit of small applied field $`\stackrel{}{E}`$. We find that in this limit it attains the physically meaningful form of the Stokes formula, which implies, in particular, that the frictional force exerted on the carrier by the environment modeled by dynamic percolation is viscous. In this limit, the carrier velocity and the friction coefficient are calculated explicitly. In addition, we determined the self-diffusion coefficient of the carrier in the absence of the field and show that it reduces to the well-know result of Refs. and in the limit when the particles number is conserved. Further more, we have found that the density profile around the carrier becomes strongly inhomogeneous: the local density of the ”environment” particles in front of the carrier is higher than the average and approaches the average value as an exponential function of the distance from the carrier. On the other hand, behind the carrier the local density is lower than the average, and depending on whether the number of particles is explicitly conserved or not, the local density past the carrier may tend to the average value either as an exponential or even as an algebraic function of the distance. The latter reveals especially strong memory effects and strong correlations between the particle distribution in the environment and the carrier position.
## 8 Acknowledgments.
G.O. acknowledges the financial support from the Alexander von Humboldt Foundation via the Bessel Research Award.
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# Near infrared and the inner regions of protoplanetary disks
## 1. Introduction
Processes in the immediate vicinity of young pre-main-sequence stars influence the initial stellar evolution and the formation of terrestrial planets. Since small scales of several AU around a star are difficult to resolve, we still lack a clear understanding of processes such as disk accretion, launching of bipolar jets and winds, and dynamics and reprocessing of dust in the inner hot disk regions. The dust geometry is one of the basic ingredients needed for constraining theoretical models of these processes. Traditionally, this geometry has been deduced from the spectral energy distribution (SED), which is dominated at infrared wavelengths by dust emission.
A widely popular geometrical description is the two-layered flared disk model developed by Chiang & Goldreich (1997) (CG hereafter). The model gives a simple method for estimating the flux from the optically-thin surface layer of an optically-thick disk directly exposed to the stellar radiation, and from the colder disk interior heated by the warmer surface. The simplicity of the method, together with evidence for the existence of disks based on radio imaging, made this model a dominant description of T Tau and Herbig Ae/Be stars (intermediate-mass, 1.5$`M_{}/M_{}`$10, counterparts of T Tau; Haebes hereafter).
Although the CG model successfully explains the observed SEDs, advances in imaging techniques revealed shortcomings of this model. Analyzing images at scattering and dust emission wavelengths, Miroshnichenko et al. (1999) concluded that disks alone cannot explain the imaging observations, at least not for some Haebes. Instead they modeled the SED with an optically thin halo surrounding an optically thick disk, and emphasized that only multi-wavelength imaging can distinguish between this and the CG model. Subsequent detailed modeling of imaging data in numerous systems revealed the existence of dusty halos around the putative flared disks (Gómez & D’Alessio, 2000; Schneider et al., 2003; Stapelfeldt et al., 2003). The inadequacy of the SED as the sole analysis tool in determining the geometry was further demonstrated by Vinković et al (2003) (V03 hereafter). They showed that the mathematical expression for the SED calculation in the CG model can be transformed into that for the halo-embedded-disk and vice versa. This has far reaching consequences for all studies based solely on SEDs. If not supported by imaging at various wavelengths, SED models can lead to erroneous conclusions about the spatial distribution of dust.
The disk inner region in Haebes (within $``$10 AU from the star) proved to be more complicated than the original CG model. Thermal dust emission from this region peaks at short wavelengths, creating a near IR bump ($`1\mu m<\lambda <8\mu m`$) in the SED of many Haebes (Hillenbrand et al., 1992). Chiang et al. (2001) noticed that the CG model did not produce enough near IR flux to explain the bump. This implies that the disk flaring, which increases the emitting volume of the optically thin disk surface, is too small at the inner radii. Since the disk geometry is constrained by vertical hydrostatic equilibrium, an additional hot dust component is required for explaining the near IR bump. To solve this problem, Dullemond, Dominik, & Natta (2001) (DDN hereafter) proposed to modify the CG geometry without introducing an additional component. They noted that the disk vertical height is increased (puffed up) at its inner rim because there the disk interior is directly exposed to the stellar radiation and hotter than in the CG model at the same radius. The rim is the hottest region of the disk and with its increased size it is possible to boost the near IR flux. This puffing of the rim is equivalent to the disk extra flaring that was identified by Chiang et al. (2001) as missing in the CG model.
Evidence in support of the DDN model was garnered from SED modeling of a large sample of Haebes (Dominik et al., 2003), and recently also of T Tau stars (Muzerolle et al., 2003). Still, the inner disk geometry remains controversial. Recent advances in near IR interferometry provide imaging data of this region, and the first results from a large sample of Haebes show that many of these objects appear close to circular symmetry (Millan-Gabet, Schloerb, & Traub, 2001; Eisner et al., 2004). This is an unusual result if disk inclinations are random. It also creates a new set of problems when interpreted as almost face-on disks because that often conflicts with outer disk inclinations derived from other imaging observations (HST, radio). This is difficult to accommodate in disk-only models, but is easily explained by halo-embedded disks (V03).
In this paper we reexamine the DDN model and the theoretical approach behind it, and identify some unresolved issues in its description of the rim emission. We employ exact radiative transfer calculations of the rim’s brightness and show that the concept of puffed up rim requires some fine tuning of the model parameters in order to produce enough flux to explain the observations (e.g. the dust must be perfectly gray).
Various independent observations indicate the existence of compact halos ($``$10 AU) around the disk inner regions (see V03 and references within), and we find that such halos readily explain also the observed near-IR excess. Furthermore, the halos also resolve the puzzling relationship noted by Monnier & Millan-Gabet (2002) between luminosity and the interferometric inner radii of disks.
## 2. Emission from the inner wall
A distinct feature of the near IR bump is its anomalously high flux $`F_\lambda `$ as compared with the stellar emission. To quantify this effect we introduce the flux ratio $`F_{2\mathrm{\mu m}}/F_{1\mathrm{\mu m}}`$ as a measure of the strength of the near IR bump; this ratio increases when the prominence of the bump becomes larger. These wavelengths are chosen because the 2$`\mu `$m flux is dominated by the dust, while the 1$`\mu `$m flux is dominated by the star.
Figure 1 summarizes the observed values of the $`F_{2\mathrm{\mu m}}/F_{1\mathrm{\mu m}}`$ flux ratio for a sample of well-observed stars, with data compiled from the following references: Low (1970); Gillett & Stein (1971); Strom et al. (1972); Allen (1973); Cohen (1973a, b, c); Glass & Penston (1974); Cohen (1975); Cohen & Schwartz (1976); Kolotilov et al. (1977); Cohen (1980); Bouchet & Swings (1982); Lorenzetti et al. (1983); Tjin A Djie et al. (1984); Kilkenny et al. (1985); The et al. (1985); Olnon et al. (1986); Berrilli et al. (1987); Strom et al. (1989); Hu et al. (1989); Lawrence et al. (1990); Fouque et al. (1992); Berrilli et al. (1992); Hutchinson et al. (1994); Li et al. (1994); Prusti et al. (1994); Sylvester et al. (1996); Garcia-Lario et al. (1997); Malfait et al. (1998); Herbst & Shevchenko (1999); van den Ancker et al. (2000); de Winter et al. (2001) and A. S. Miroshnichenko (2005, private communication).
Dust extinction at 1$`\mu `$m is larger than at 2$`\mu `$m and could enhance the observed strength of the near IR bump by $``$20% for $`A_V=1`$, therefore only objects with $`A_V1`$ were considered. Since the reddening correction is negligible, the uncorrected data displayed in the figure represent the true range of near IR bump strength in Herbig Ae stars. The underlying stars of all objects have temperatures of about 10,000 K, which gives $`F_{2\mathrm{\mu m}}/F_{1\mathrm{\mu m}}`$ = 0.09. Yet in all objects this ratio exceeds 0.25, reflecting a large NIR excess from hot dust emission (Hillenbrand et al. 1992). The luminosity of each object is displayed together with its name in figure 1 and it ranges from $``$5$`L_{}`$ to $``$80$`L_{}`$. Luminosity does not show any correlation with the near IR bump strength, reaffirming our conclusion that these data can be used as a general description of the near IR bump strength in Herbig Ae stars.
### 2.1. General Description of the Rim Emission
At the inner rim, gas that is typically part of the disk cold interior becomes directly exposed to the stellar radiation and expands to higher scale heights. According to DDN, emission from such a puffed-up rim can explain the near IR bump in the spectrum of Herbig Ae/Be stars. The rim geometry is sketched in figure 2. The rim is modeled as a cylinder of radius $`R_{rim}`$ and height 2$`H_{rim}`$ centered on the star. The basic assumption of this model is that the rim is optically thick in the near IR and shorter wavelengths. This maximizes the rim energy output.
The original DDN model (Dullemond, Dominik, & Natta, 2001) successfully explained the data, but it was based on an approximate treatment of the rim height and emission. More realistic models were calculated by the authors of the DDN model in their subsequent work. Dullemond (2002) used a 2D radiative transfer model for gray dust combined with the hydrostatic equilibrium. The obtained near IR bump strength is shown in figure 1 (solid line). The maximum strength is still too low to explain all the data, but it can accommodate the majority of observed near IR bump strengths.
A dramatic reduction of the DDN model efficacy happens when a mixture of small and big grains is introduced. Dullemond & Dominik (2004)(DD04 hereafter) combined 2mm (big, gray grains) and 0.1$`\mu `$m (small) grains in various ratios and performed 2D radiative transfer calculations coupled to the equation of vertical hydrostatics and dust settling. The model fails to explain the data even when 99.999% of the dust mass is in big grains (see figure 1). The behavior of this result is unexpected; decrease in the small grain population leaves more gray dust grains in the mix, which should move the whole solution closer to the gray dust result of Dullemond (2002).
A closer inspection of obtained results shows that the temperature of big grains in such a multi-grain mixture is lower than in the pure gray model. In a mixture, both small and big grains absorb a fraction of the local energy density and participate in providing the local diffuse heating. But, as shown in the next section (see also equation A8 in the appendix), small grains are a very inefficient source of local diffuse heating, resulting in less efficient heating of the big grains then in the pure gray model. With such a temperature decrease, the vertical hydrostatic equilibrium cannot produce disk puffing comparable to the gray model. While in the gray model the puffed-up disk rim height is close to $`H_{rim}/R_{rim}=0.2`$, multi-grain models have only $`H_{rim}/R_{rim}<0.15`$. Since the observed rim emission scales with rim height, this is the major reason behind the failure of the multi-grain models to explain the data. The presented model with the lowest fraction of small grains yields the largest discrepancy because it suffers the largest reduction in the small grain contribution to gas heating and rim puffing while still having enough small grains to suppress heating of the big grains. In quantifying this effect, the mass ratio between big and small grains that DD04 used is not the most illustrative choice. A more appropriate quantity would be the “equivalent” grain size of the grain mixture. In the case of two grain populations with sizes $`a_{big}`$ and $`a_{small}`$ and fractional number densities $`X_{big}`$ and $`X_{small}`$ (such that $`X_{big}+X_{small}=1`$), the average grain size obeys
$$a^2=X_{big}a_{big}^2+X_{small}a_{small}^2.$$
(1)
The number fractions can be deduced from the reported parameters of the DD04 models: fixed inner and outer disk radius, fixed total disk mass, and total dust mass in big and small grains. The model with 99.999% of the dust mass in big grains has only $`X_{big}=1.25\times 10^8`$, and $`X_{small}1`$, yielding $`a^2^{1/2}0.25\mu m`$. This grain size is too small to be considered equivalent to the gray dust model.
It is important to note that this equivalent grain is just an indicator of the overall solution and cannot be used as a general replacement (average or synthetic) grain for the radiative transfer calculation. As already shown by Wolf (2003), the approximation of an averaged single grain as a replacement for a dust mixture breaks down at the surface of a dust cloud (or in this case the rim surface). A more detailed study of multi-grain disk models will be presented in a separate publication, while in the next sections we explore the limits of possible DDN model applicability in the context of single dust grains.
### 2.2. Approximate Solution for the Rim Emission
Denote by $`R_{}`$ and $`T_{}`$ the stellar radius and temperature, respectively. At distance $`d`$ and direction $`i`$ where the star is free of rim obscuration (see figure 2), the overall observed flux at wavelength $`\lambda `$ is $`(R_{}/d)^2\pi B_\lambda (T_{})+F_\lambda ^{rim}(i)`$. If $`I_\lambda ^{rim}(i)`$ is the rim surface brightness in the observer’s direction then
$$F_\lambda ^{rim}(i)=\frac{1}{d^2}I_\lambda ^{rim}(i)\mathrm{\hspace{0.17em}4}H_{rim}R_{rim}\mathrm{sin}i.$$
(2)
Here the cylindrical visible surface is replaced with a flat rectangle. This approximation maximizes the flux since curvature decreases the projected area of portions of the visible surface, reducing the observed flux. We also assume that the stellar illumination is perpendicular to all portions of the rim. This too maximizes the observed flux.
The observed rim flux in equation 2 is determined by the rim height, surface brightness and radius. Our 2D radiative transfer calculations described in §2.3 confirm that the rim emission is indeed proportional to the rim height, therefore we maximize the rim emission in this study by using $`H_{rim}=0.2R_{rim}`$, the maximum height allowed before the rim starts to shadow large portions of the disk (DDN). The solution for any other rim height can be derived from our models by a simple scaling of the rim emission.
The surface brightness of a gray dust rim can be approximated with $`B_\lambda (T_{rim})`$, where $`T_{rim}`$ is the dust sublimation temperature. The description of a non-gray surface must take into account the spectral variation of optical depth of the emitting optically-thin surface layer. This was done by Chiang & Goldreich (1997). According to their model, the surface layer vertical optical thickness is unity at visual (a characteristic wavelength of the stellar radiation absorption) $`\tau _V=1`$, therefore at all other wavelengths it is $`\tau _\lambda =\sigma _\lambda ^{abs}/\sigma _V^{abs}q_\lambda `$. The rim emits at near IR where $`q_\lambda <1`$ (the dust NIR opacity is smaller than at visual), thus the surface layer is optically thin at these wavelengths and its emission is reduced accordingly. Therefore, the rim surface brightness becomes $`q_\lambda B_\lambda (T_{rim})/\mathrm{sin}i`$ and the observed rim flux is
$$F_\lambda \frac{4}{d^2}B_\lambda (T_{rim})\frac{H_{rim}}{R_{rim}}R_{rim}^2\times \{\begin{array}{cc}q_\lambda \hfill & \text{non-gray dust}\hfill \\ & \\ \mathrm{sin}i\hfill & \text{gray dust}\hfill \end{array}$$
(3)
This result shows that a non-gray rim creates a smaller IR excess than a gray opacity rim. In addition, non-gray opacity removes the angle dependence from the rim emission (we expect this approximation to break down at very small inclination angles where $`q_\lambda \mathrm{sin}i`$).
The rim radius is derived from radiative equilibrium, which gives (e.g. Ivezić & Elitzur, 1997)
$$R_{rim}=\frac{1}{2}R_{}\left(\frac{T_{}}{T_{rim}}\right)^2\left[\frac{\overline{\sigma }(T_{})}{\overline{\sigma }(T_{rim})}\psi \left(1+\frac{H_{rim}}{R_{rim}}\right)\right]^{1/2}$$
(4)
Here $`\overline{\sigma }(T)`$ is the Planck average of $`\sigma _\lambda ^{abs}`$ at temperature $`T`$, $`\psi `$ describes the correction for diffuse heating from the rim interior<sup>1</sup><sup>1</sup>1Note that Ivezić & Elitzur (1997) used $`\mathrm{\Psi }=\frac{\overline{\sigma }(T_{})}{\overline{\sigma }(T_{rim})}\psi `$ and $`1+H_{rim}/R_{rim}`$ is a correction (described by DDN) for self-irradiation from the other side of the rim. In appendix A we derive an approximate solution which shows that gray dust, with $`\overline{\sigma }(T_{})/\overline{\sigma }(T_{rim})=1`$, has $`\psi 4`$ and that non-gray dust, with $`\overline{\sigma }(T_{})/\overline{\sigma }(T_{rim})>1`$, has $`\psi 1`$. Note that for gray dust this makes equation 4 identical to the original DDN expression (their equation 14). The approximate near-IR bump strength is given in equation A10, yielding
$$\frac{F_{2\mu m}}{F_{1\mu m}}\{\begin{array}{cc}0.23\hfill & \text{non-gray dust}\hfill \\ & \\ 0.09+0.52\mathrm{sin}i\hfill & \text{gray dust}\hfill \end{array}$$
(5)
for $`T_{}`$=10,000 K, $`T_{rim}`$=1,500 K and $`H_{rim}/R_{rim}`$=0.2. Comparison of this result with the data in figure 1 shows that the NIR bump of non-gray dust is too small to explain the observations. Therefore, interpretation of the NIR bump in Herbig Ae stars with inner disk puffing places a strong constraint on dust evolution in this region. The dust must grow to a size greatly exceeding the initial interstellar size distribution, and small grains must be depleted to such a large extent that the inner disk opacity can be considered gray. In the next subsection we employ exact 2D radiative transfer code to obtain accurate values for $`\psi `$ and place more precise constrains on the DDN model.
### 2.3. Exact Models for Single-Size Grains
To examine the validity of conclusions based on our approximate solution we performed full 2D radiative transfer calculations for an optically thick torus centered on a 10,000 K star. The torus cross section is sketched in figure 2; it is a square with side-length of $`2H_{rim}`$, where $`H_{rim}=0.2R_{rim}`$. This configuration is the same as described by DDN, where the puffed-up disk rim is a cylindrical surface directly exposed to stellar radiation, while the rest of the inner disk is in its shadow. The dust has sublimation temperature $`T_{sub}=`$1,500 K and constant density everywhere in the torus, with horizontal and vertical optical depths $`\tau _V=10,000`$ in visual. Different density structures do not change our results as long as the $`\tau _V=1`$ layer on the illuminated surface is geometrically much smaller than $`H_{rim}`$.
Radiative transfer modeling was conducted with our code LELUYA (http://www.leluya.org) that works with axially symmetric dust configurations. It solves the integral equation of the formal solution of radiative transfer including dust scattering, absorption and thermal emission. The solution is based on a long-characteristics approach to the direct method of solving the matrix version of the integral equation (Kurucz, 1969).
The results are shown in figure 3 together with the original DDN solution (dashed line). Our 2D model results for gray dust without scattering (solid line) are very close to the Dullemond (2002) results, shown in figure 1, which also included vertical hydrostatics equilibrium. This model has $`\psi =`$4. Its rim radius (49$`R_{}`$) and flux are essentially the same as the original DDN model, confirming that a puffed-up rim of gray dust is capable of explaining the near IR bump.
Since realistic dust is not gray at all wavelengths, we calculated models for silicate dust with different grain radii, employing optical constants from Dorschner et al. (1995) ($`x`$ = 0.4 olivine). Figure 3 shows results for three representative grain radii, with the corresponding rim properties summarized in table 1. Our model results for 0.1$`\mu `$m grains are almost identical to the Dullemond (2002) results for purely small grains (see figure 1). As is evident from figure 3, the model can explain the data when the grain radii are 2$`\mu `$m and 0.5$`\mu `$m, but it starts to fail as a general explanation of the near IR bump when the grain radius drops below $`0.1\mu `$m.
A decrease in grain size has two opposing effects on the rim flux. On one hand, the ratio $`\overline{\sigma }(T_{})/\overline{\sigma }(T_{sub})`$ is increasing, leading to a larger rim radius and emitting area and thus enhancing the rim emission. On the other, the rim surface brightness is declining because $`q_\lambda `$ is decreasing, reducing the rim emission. The net result is that maximum rim emission occurs at grain radius of $``$0.5$`\mu `$m, which, as is evident from figure 4, corresponds to the transition between gray and non-gray opacity in the near IR region. This is predominantly a grain size effect; the dust chemistry introduces only second order corrections.
When the grain radius drops below 0.5$`\mu `$m the dust opacity becomes non-gray in the near IR and the puffed-up rim model begins to fail. The flux of the 0.1$`\mu `$m grain model, which is almost angle independent as predicted by equation A10, can reproduce only the weakest near-IR bumps. Therefore the DDN model can explain the near IR bump in Herbig Ae stars only when both of the following conditions are met:
1) the rim dust opacity is gray in the near IR (grain radius $``$ 0.5$`\mu `$m),
and
2) the disk is puffed to a height $`H_{rim}/R_{rim}0.15`$.
Figure 1 shows that for these conditions to be satisfied, the DDN model requires the complete absence of small grains in the disk inner region. Therefore, for this model to work, the rim dust must undergo substantial growth that also fully depletes the population of small grain. At the same time, this process cannot be so extreme in the rest of the disk because the mid IR spectrum of Herbig Ae stars displays the dust features of small grain emission (van Boekel et al., 2005).
## 3. The near IR bump and imaging explained with a dusty halo
A dusty halo around the disk inner regions ($``$ 10 AU) has been invoked to explain polarimetric measurements (Yudin, 2000) and correlations between variabilities in the optical and near IR (Eiroa et al., 2002). Such small regions are not yet accessible to direct imaging but have been resolved in near IR interferometry by Millan-Gabet, Schloerb, & Traub (2001) who also favor the halo geometry, although the interpretation of these visibility data is still model-dependent. Direct imaging is currently available only for larger scales, and these observations have revealed larger halos, $``$ 100 AU, around some objects (V03). The relation between the inner and outer halos, whether they are simply the inner and outer regions of the same circumstellar component, remains an open question. However, at the phenomenological level this issue is not relevant because the two can be treated as separate circumstellar components if both are optically thin. The inner halo is then radiatively decoupled from the cooler outer halo, simplifying the study of inner halos.
Here we explore the contribution of the inner halo to the near IR emission. The halo precise geometry is not particularly important. It could be elongated, clumpy or inhomogeneous, but as long as it is optically thin it can be approximated with spherical geometry. The reason is that the temperature of optically thin dust is dominated by the stellar heating, resulting in spherically symmetric isotherms and circularly symmetric images at wavelengths where the dust thermal emission dominates over scattering (V03). Optically thin halos are also transparent to the disk emission and we can ignore the disk effect on the halo. The exact image shape ultimately depends on detailed dust density and grain properties, telescope resolution and sensitivity, observational wavelength and the intrinsic ratio between the disk and halo surface brightness. Various observations of R Mon vividly illustrate these effects (see Weigelt et al., 2002).
If the halo optical depth at visual wavelengths $`\tau _V`$ is larger than $`\frac{1}{4}H/R`$, where $`H/R`$ is the disk flaring at the halo outer radius, then the halo dominates the SED coming from the dust within radius $`R`$ around the star (see V03 for details). At near IR wavelengths, this condition is satisfied for the halo optical depths of interest here ($`\tau _V0.1`$).
### 3.1. Theoretical Examples
Our models consist of a star surrounded by a spherical halo with radial density profile $`\eta r^p`$. The halo extends from inner radius $`R_{in}`$, set by the dust sublimation temperature $`T_{sub}`$, to outer radius $`R_{out}`$.
The dust chemistry is $`x`$=0.4 olivine from Dorschner et al. (1995), with grain size distribution n($`a`$)$`a^q`$ between the minimum grain radius $`a_{min}`$ and maximum $`a_{max}`$. We use $`q=2`$, $`a_{min}=0.01\mu `$m and vary $`a_{max}`$. The radial optical depth of the halo is specified at $`\lambda `$ = 0.55$`\mu `$m as $`\tau _V`$. The radiative transfer problem is solved with the code DUSTY (Ivezić, Nenkova, & Elitzur, 1999), which takes advantage of the scaling properties of the radiative transfer problem for dust absorption, emission and scattering (Ivezić & Elitzur, 1997).
Figures 5 and 6 show some SED examples for dusty halos around 10,500K and 5,000K stars, representative of Herbig Ae and T Tauri stars, respectively. The stellar spectrum is taken from Kurucz models. In addition to the strength parameter $`F_{2\mathrm{\mu m}}/F_{1\mathrm{\mu m}}`$, the flux ratio $`F_{4\mathrm{\mu m}}/F_{2\mathrm{\mu m}}`$ can be used to characterize the NIR bump shape. Both the strength and shape parameters are influenced by changes in the dust sublimation temperature, maximum grain size, halo outer radius and optical depth. Comparison of the data with halo model results for the strength and shape parameters is shown in figure 7 for the same objects as in figure 3. Models for $`p=1`$ and $`p=2`$ halos around a 10,500K star are dispersed all over the diagram. Arrows show how the model results move in the diagram as the model parameters are varied, indicating that various degeneracies are possible. The observed levels of bump strength and shape are readily reproduced with plausible values of the model parameters.
We briefly summarize the effect of various halo parameters on the strength and shape of the near IR bump.
Optical depth: A larger optical depth results in a stronger near IR bump. This reflects the dependence of flux on the total mass of emitting dust (equations A7 and A12 in V03). The dust sublimation radius $`R_{in}`$ is only slightly affected, as expected in the optically thin limit where the diffuse radiation is negligible.
Grain size: Larger grains shift the near IR bump toward longer wavelengths and make it appear more flat. With increased grain size the opacity becomes more similar to gray dust, resulting in a $`r^{0.5}`$ temperature profile since the geometrical dilution of stellar heating is the only cause of temperature variation. Smaller grain sizes create steeper radial temperature profiles. Therefore, for a given density profile smaller grains emit relatively more radiation at shorter wavelengths than larger grains. In practice, grain sizes come in mixtures and sublimate at different radial distances, greatly adding to the complexity of the problem. The SED models are therefore prone to various model degeneracies.
Sublimation temperature: With a higher dust sublimation temperature, the near IR bump shifts to shorter wavelengths, reflecting the shift of the emission peak.
Outer radius: The halo size can affect the near IR bump in two ways. On one hand, reducing the outer radius while keeping the dust distribution fixed reduces also the total optical depth. The near IR bump then starts to decrease when the dust removal reaches the NIR emission regions at radial distance $`10R_{in}`$ (temperatures $``$ 500 K). On the other hand, reducing the outer radius at a fixed halo optical depth is equivalent to redistributing the dust within the halo. The bump then becomes stronger as the outer radius is reduced because more dust is shifted toward smaller radii and higher temperatures.
Stellar temperature: As its temperature decreases, the emission from the star starts to blend with that from the halo, and the near IR bump disappears. Only a careful analysis can then separate the stellar from the diffuse flux in the near IR and reveal the bump. For comparison with T Tau stars, figure 7 shows also 5,000 K models (marked with T). In spite of the large variations in halo parameters, these models display only a limited range of bump strengths and shapes close to the naked star values. This explains why the near IR bump was not originally recognized in T Tau stars while easily detected in Herbig Ae/Be stars.
### 3.2. Observational Examples
We show three examples which illustrate different circumstellar dust configurations: AB Aur, HD 100546 and HD 163296. HST imaging suggests that AB Aur and HD 100546 have large halos at radii $`1`$”, while HD 163296 shows only a disk (Grady et al., 2003). Irrespective of the existence of a large halo, all three objects show a near IR bump, with the strongest bump in AB Aur. Since the focus of this study is the near IR bump, the large scale halos are irrelevant here and we only consider a small halo within $``$10 AU around the star.
Our fits to the data are shown in figure 8, with the model parameters listed in table 2. The halo outer radius is 10 times the dust sublimation radius in all models. Since our focus is the near IR bump, our model consists only of the star and the inner halo, and only data at wavelengths shorter then 6$`\mu `$m were employed in the fitting. The derived model parameters are not unique since various degeneracies exist in model results for the near IR flux (see §3.1). For example, the “hot component” in the Bouwman et al. (2000) models can be interpreted as a small-scale halo with dust properties different from those in our study.
#### 3.2.1 AB Aur
The emission from AB Aur has been resolved at various wavelengths and interpreted as a disk with vastly different estimates for the inclination angle, as follows:
| visual: | $`i45^{}`$ | Grady et al. (1999) |
| --- | --- | --- |
| near IR: | $`i30^{}`$ | Eisner et al. (2003), (2004) |
| | $`i=30\pm 5^{}`$ | Fukagawa et al. (2004) |
| mid IR: | $`i=55\pm 10^{}`$ | Liu et al. (2005) |
| millimeter: | $`i=17{}_{}{}^{}{}_{}{}^{+6}_3`$ | Semenov et al. (2004) |
| | $`i=21{}_{}{}^{}.5_{0.3}^{+0.4}`$ | Corder et al. (2005) |
| | $`i=33\pm 10^{}`$ | Piétu et al. (2005) |
| | $`i76^{}`$ | Mannings & Sargent (1997) |
Such a disparity is expected in halo-embedded disks (see figure 7 in V03) because, as noted by Miroshnichenko et al. (1999), the halo dominates the images at wavelengths extending to $``$ 100 µm or so, and the disk emerges only at longer wavelengths. Interpretation of molecular line images, too, must be done carefully to avoid confusion between the halo and disk contributions.
A general conclusion about the AB Aur inner halo is that it must have a radial density profile between $`1/r`$ and $`1/r^2`$; this differs from the outer halo, which has a $`1/r`$ density profile as deduced from the $`1/r^2`$ radial brightness profile of the HST image (Grady et al., 1999, see also equation A10 in V03). Conclusions regarding the properties of the dust grains in the inner halo are less firm. Near IR interferometry suggests the presence of dust close to the star, implying large grains that can survive at small distances. An example of a big grains model for AB Aur is shown in figure 8 with thick dashed line (see also table 2). The grain size and chemistry might be subject to radial variations, as is indicated by comparison between the HST (Grady et al., 1999) and Subaru images (Fukagawa et al., 2004), further complicating the modeling.
#### 3.2.2 HD 100546
The HST image of this source (Grady et al., 2001) shows a very tenuous large scale nebulosity, whose low surface brightness implies an optical depth of only $`\tau _V`$ 0.015. This component of the dust distribution does not contribute significantly to the IR emission and can be ignored in the current analysis. The HST image, which is produced purely by scattered light, reveals also a prominent disk with inclination angle $`49\pm 4^{}`$. Near IR (Augereau et al., 2001) and mid IR (Liu et al., 2003) imaging give similar results for the disk even though the latter is produced purely by dust emission and the former contains a mixture of both scattering and emission. The HST brightness contours are symmetric, with the brightness declining as $`1/r^3`$. These are the signatures of scattering from the CG layer of a flat disk (see V03). However, for the CG model to be applicable, every point on the scattering surface, which extends to a distance of $``$10″ from the star, must have a direct line of sight to the stellar surface. This is impossible in the case of a flat disk, since it would have to maintain a thickness smaller than the stellar radius for hundreds of AU. Therefore, the only self-consistent explanation of the HST imaging is with an optically thin halo whose dimensions are unrelated to the stellar size. The HST image implies that the halo has a flattened geometrical shape, and its $`1/r^3`$ brightness profile implies that it has a $`1/r^2`$ radial density profile (V03). This flattened halo is outlined as region A in figure 9. Since the halo dominates the imaging, the geometry of the optically thick disk structure, outlined as region B in the figure, remains unknown.
The HST imaging does not constrain the inner-halo geometry at radii $``$ 10 AU. The surface density must be reduced in that region because the near IR bump in HD 100546 is significantly smaller then in AB Aur even though otherwise the two stars are rather similar. Indeed, the fit to the near IR bump yields a $`1/r`$ radial density profile (figure 8 and table 2), shallower than in the region resolved by HST. The fit was further improved by an increased contribution from large grains ($`a_{max}=0.5\mu `$m and $`q=2`$) and a reduced fraction of carbon dust in the mix. Observations by Grady et al. (2005) show that a constant density profile, creating $`1/r`$ brightness profile, might be more appropriate in the region between 20 and 50 AU. We find that a constant density model could also fit the near IR spectrum if the sublimation temperature were increased to 1700 K. All these results point toward large structural differences between the inner and outer regions of HD 100546.
#### 3.2.3 HD 163296
The model properties of the inner halo in this source are very similar to AB Aur, except that a shallower density profile of $`p=1`$ is preferred (see table 2). A similar general conclusion is that the halo radial density profile is between $`1/r`$ and $`1/r^2`$, with uncertainties in the dust properties. Significantly, in this source the inner halo also fits the 10$`\mu `$m feature all by itself (figure 8). No other optically thin components are required for explaining the mid IR dust features, indeed none are observed. The HST image, which is incapable of resolving the inner halo, shows no evidence of a large scale structure other then the disk, with an inclination of $`60\pm 5^{}`$ (Grady et al., 2000) in agreement with $`58^{}`$ derived from millimeter observation (Mannings & Sargent, 1997). Another noteworthy feature of the HST image is a bipolar jet. The process responsible for jet formation could perhaps also lift up dust above the disk and create the small scale halo responsible for the near IR bump. Such possible correlation between jets and the near IR bump can be studied further when more high resolution data from a larger sample of objects become available.
### 3.3. Size-luminosity Correlation
The milli-arcsecond resolution reached in near IR interferometry enables studies of the immediate environment of young stars, down to 0.1 AU (Malbet, 2003). Unfortunately, current visibility data are not yet capable of reproducing the full 2D image of an object, instead requiring a model of the geometry for their analysis. One simple and often used model of the circumstellar geometry is a flat dust ring of uniform surface brightness. This ad hoc model did not arise from some specific radiative transfer modeling but rather chosen as a simple approach to the visibility fitting procedure.
Fitting the visibility data of a number of objects with this ring model, Monnier & Millan-Gabet (2002) discovered that the size of the ring inner radius increased with the stellar luminosity $`L_{}`$. This is the expected result when dust sublimation controls the size of the dust-free region around the star. Since radiative transfer is scale invariant (Ivezić & Elitzur, 1997), inner radii of rings would be expected to scale as $`L_{}^{1/2}`$ if their dust properties were the same. However, Monnier & Millan-Gabet (2002) do not find such a trend. Instead, at a fixed luminosity the derived radii vary by almost a factor of ten, which they refer to as scatter in the size-luminosity diagram. This scatter indicates either that the disk inner regions have vastly different properties, with the sublimation temperature varying from $``$1000 K to $``$2000 K, or that the ring model is not a proper description of the actual dust distribution. Monnier & Millan-Gabet (2002) also noted that some highly luminous objects ($`L_{}10^3L_{}`$) had smaller than expected inner ring radii, thus requiring even higher dust sublimation temperatures. New interferometric data by Monnier et al. (2005) slightly reduce the scatter in the ring-radius–luminosity relation, but the remaining scatter still implies a large range of sublimation temperatures, and very luminous objects still display abnormally small radii.
Instead of the ring model we have analyzed the interferometry results with the inner-halo model, performing simultaneous fits of both the near-IR bump and visibility data. Preliminary results are shown in figure 10. It is highly significant that there are no objects in the forbidden region below the indicated lower limits. The correlation of overall bolometric luminosity with inner radius is much tighter than in the ring model, the small remaining scatter arises from variation in halo optical depth and grain size. In contrast with the ring model, the sublimation temperature rarely differs from 1500K (it is 1800K in a couple of objects). The high luminosity object MWC297 which was especially troubling in the Monnier & Millan-Gabet (2002) analysis is now consistent with 1500K sublimation temperature. It is striking how some of the objects that were highly scattered in the diagram by Monnier & Millan-Gabet (2002) now settle on the same $`L_{}^{1/2}`$ size-luminosity relation (dotted line in figure 10), indicating similarities in the halo properties of all these stars, in turn pointing toward a common physical mechanism of halo formation.
## 4. Conclusions
An examination of the puffed-up disk rim model (DDN) shows that it has rather limited capabilities in explaining the near IR bump of Herbig Ae/Be stars. The observed level of near IR excess implies a certain emitting volume of optically thin puffed-up disk rim surface for given dust properties. The volume derived from the DDN model falls short of this observational limit, unless the disk is made of perfectly gray dust. The puffed-up rim produces enough near IR flux only when the inner disk consists purely of dust grains larger than $``$ 0.5 $`\mu `$m and the disk puffing reaches values of $`H_{rim}/R_{rim}0.15`$. Models by Dullemond & Dominik (2004) show that even traces of small grains inhibit the disk puffing, eliminating the DDN model as a viable explanation of the near IR bump. Since the 10$`\mu `$m emission feature indicates the presence of small grains in the circumstellar dust, additional mechanisms must be invoked to remove all small grains from the inner disk and keep the DDN model viable.
From fits to the SED of a number of Haebes, Dominik et al. (2003) conclude that the infrared excess in these stars is produced by disks alone without the need for additional circumstellar components. This conclusion is invalidated by the mathematical proof that a fit to the SED cannot distinguish between the surface of a flared disk and an optically thin halo (V03). Fits to the SED alone are not a conclusive proof of a particular dust geometry.
We find that that the optically thin dusty halos around the disk inner regions whose existence has been inferred in various observations readily explain the strength and shape of the near IR bump. The halo is not limited by the disk properties. Hence, it can extend above the disk surface and accommodate the emitting optically thin dust volume required by the near IR flux observations. The required halo is rather small, less than several AU in size, and its optical depth in visual is less than $``$ 0.4. Despite its small optical depth, the halo dominates the near IR spectrum and hides the disk near IR signature. However, detailed properties of the halo, such as its exact shape, grain properties or dust density profile, are not uniquely constrained by the SED since different combinations of the parameters can produce the same flux. These degeneracies can be broken only with imaging capable of resolving the disk inner regions.
Inner halos not only explain the near IR bump but also successfully resolve the puzzle presented by the relations between luminosities and near-IR interferometric sizes (Monnier & Millan-Gabet, 2002). In addition to their near IR emission, the halos contribute also to the mid IR flux. HD 163296 is an extreme example where the halo in itself fully explains the mid IR dust features without the need for additional extended components (figure 8). The absence of such components in the HST image of this source is another success of the inner-halo model. In general, though, the inner halo emission is not expected to dominate the mid IR but still make a significant contribution that must be included in fits to the overall SED for reliable modelling of the rest of the circumstellar material. Recently van Boekel et al. (2003) suggested that differences in the strength and shape of the mid IR silicate feature in Haebes are evidence for dust settling in the disk. However, these differences could instead reflect halo evolution, with the most active stars showing the strongest mid IR signature of the inner disk halo. High resolution imaging is necessary for definite conclusions about evolution either of the dust or the circumstellar disk. Such imaging will soon become available at the VLTI, which offers milli-arcsecond resolution at near IR.
We thank C. P Dullemond, C. Dominik and A. Natta for fruitful discussions on the physics of the DDN model. We also thank A. S. Miroshnichenko for help with the data compilation. DV thanks B. Draine and R. Rafikov for useful comments. Support by the NSF grant PHY-0070928 (DV) is gratefully acknowledged. This work was also supported by National Computational Science Alliance under AST040006 and utilized the NCSA’s Xeon Linux Cluster. DV also thanks the Institute for Advanced Study (IAS) for time on their Linux cluster. TJ acknowledges the hospitality and financial support of IAS during his visit to the Institute. ME acknowledges partial support by NSF and NASA.
## Appendix A Approximate solution for the rim radius
The observed emission from a puffed up inner disk depends on the rim radius, height and surface brightness (see §2.2). Deriving the rim radius requires a proper treatment of the rim temperature structure. As already noted by Dullemond (2002), in gray dust the diffuse radiation creates a temperature inversion — the dust temperature is maximum in the rim interior (that is, at $`R>R_{rim}`$), not on the rim surface. Here we derive an approximate solution for the dust temperature $`T_0`$ on the rim surface and the temperature $`T_1`$ inside the rim at depth $`\tau _V1`$ from the surface. The solution demonstrates the inversion effect for gray dust and shows that it does not exist in the non-gray case.
### A.1. Dust temperature $`T_0`$ on the rim surface
Consider a dust grain on the rim surface. It is heated by the stellar flux $`F_{}`$ and the diffuse flux $`F_{out}`$ coming out from the rim interior. For large optical depths these two fluxes are balanced: $`F_{}=F_{out}`$. Stellar flux absorbed by the grain is $`\sigma _VF_{}`$, where $`\sigma _V`$ is the dust cross section in visual. Absorbed diffuse flux is $`2\sigma _{IR}F_{out}`$, where $`\sigma _{IR}`$ is the cross section in the near IR and the factor 2 accounts for absorption from 2$`\pi `$ steradian. The grain emits into 4$`\pi `$ steradian, so that the energy balance is
$$\sigma _VF_{}+2\sigma _{IR}F_{out}=4\sigma _{IR}\sigma _{SB}T_0^4$$
(A1)
where $`\sigma _{SB}`$ is the Stefan-Boltzmann constant. Using $`F_{}=F_{out}`$ we get
$$\sigma _{SB}T_0^4=\frac{F_{}}{4}\left(2+\frac{\sigma _V}{\sigma _{IR}}\right)$$
(A2)
### A.2. Dust temperature $`T_1`$ at $`\tau _V1`$ from the surface
Now consider a dust grain at distance $`\tau _V1`$ from the surface into the rim. This grain is heated by the attenuated stellar flux $`F_{}\mathrm{exp}(\tau _V)=F_{}\mathrm{exp}(1)`$, by the diffuse flux from the surface dust between $`\tau _V=0`$ and $`\tau _V1`$, and by the diffuse flux from the rim interior. The absorbed stellar flux is $`\sigma _VF_{}\mathrm{exp}(1)`$. Diffuse contribution from the surface dust layer is $`2\sigma _{IR}\sigma _{SB}T_0^4\tau _{IR}`$, where $`\tau _{IR}\sigma _{IR}/\sigma _V`$ is the infrared optical depth of this surface layer. Diffuse heating from the rim interior, described by temperature $`T_1`$, is $`2\sigma _{IR}\sigma _{SB}T_1^4`$. This is a good approximation for gray dust and an overestimate for non-gray dust, where that temperature decreases rapidly with optical depth. The energy balance is
$$\sigma _VF_{}\mathrm{exp}(1)+2\sigma _{IR}\sigma _{SB}T_0^4\tau _{IR}+2\sigma _{IR}\sigma _{SB}T_1^4=4\sigma _{IR}\sigma _{SB}T_1^4$$
(A3)
Using $`T_0`$ from equation A2 we get the interior temperature
$$\sigma _{SB}T_1^4=\frac{F_{}}{4}\left(2\frac{\sigma _{IR}}{\sigma _V}+1+2\frac{\sigma _V}{\sigma _{IR}}e^1\right)$$
(A4)
Note that the ratio $`T_0/T_1`$ depends only on $`\sigma _{IR}/\sigma _V`$ and is independent of $`F_{}`$.
### A.3. Gray and non-gray regimes
We consider two distinct opacity regimes: gray when $`\sigma _{IR}/\sigma _V1`$ and non-gray when $`\sigma _{IR}/\sigma _V1`$. The ratio of the rim surface temperature $`T_0`$ and the interior temperature $`T_1`$ in these two regimes is
$$T_0/T_1=0.95\text{when}\sigma _{IR}/\sigma _V=1\text{(gray dust)}$$
(A5)
$$T_0/T_11.08\text{when}\sigma _{IR}/\sigma _V1\text{(non-gray dust)}$$
(A6)
The gray opacity creates a temperature inversion with the temperature in the rim interior higher than on its surface. This inversion does not appear in non-gray dust where the temperature decreases monotonically with distance from the rim surface. If the maximum dust temperature is 1,500K (sublimation temperature) then gray dust has $`T_1`$=1,500 K and $`T_0`$1,400 K, while non-gray dust has $`T_0`$=1,500 K and $`T_1`$1,400 K.
These approximate expressions are in reasonable agreement with the results of exact 2D radiative transfer calculations (see §2.3), which yield
$`T_0`$=1387 K and $`T_1`$=1490 K for gray dust
$`T_0`$=1500 K and $`T_1`$=1236 K for 0.1$`\mu `$m grains
The transition between these two regimes occurs at grain radius 0.5$`\mu `$m which yields $`T_0`$=1474 K, $`T_1`$=1492 K and the maximum temperature of 1,500 K at $`\tau _V`$=0.34.
### A.4. Disk rim radius and near-IR bump strength
Since the rim optical depth is large, we can assume that both temperatures $`T_0`$ and $`T_1`$ are located at essentially the same distance from the star. If we set $`T_1`$ to the dust sublimation temperature $`T_{sub}`$ then based on equation A4 the rim radius is
$$R_{rim}=\frac{1}{2}R_{}\left(\frac{T_{}}{T_{sub}}\right)^2\sqrt{2\frac{\sigma _{IR}}{\sigma _V}+1+2\frac{\sigma _V}{\sigma _{IR}}e^1}\sqrt{1+H_{rim}/R_{rim}}$$
(A7)
where we used $`F_{}=L_{}\sqrt{1+H_{rim}/R_{rim}}/4\pi R_{rim}^2`$ (the factor $`\sqrt{1+H_{rim}/R_{rim}}`$ is correction for rim self-irradiation, already introduced by DDN). Comparison with equation 4 for the rim radius gives
$$\psi =\frac{2}{e}+\frac{\sigma _{IR}}{\sigma _V}\left(1+2\frac{\sigma _{IR}}{\sigma _V}\right).$$
(A8)
The two extreme opacity regimes yield
$$\psi \{\begin{array}{cc}1\hfill & \text{when }\sigma _{IR}/\sigma _V1\text{ (non-gray dust)}\hfill \\ & \\ 4\hfill & \text{when }\sigma _{IR}/\sigma _V1\text{ (gray dust)}\hfill \end{array}$$
(A9)
This result can also be derived by setting $`T_0T_{sub}`$. Combining this result with equations 3 and 4 and dividing the overall observed flux at 2$`\mu `$m by the stellar flux at 1$`\mu `$m yields the near-IR bump strength
$$\frac{F_{2\mu m}}{F_{1\mu m}}=\frac{B_{2\mu m}(T_{})}{B_{1\mu m}(T_{})}+\left(\frac{T_{}}{T_{rim}}\right)^4\frac{B_{2\mu m}(T_{rim})}{B_{1\mu m}(T_{})}\frac{H_{rim}}{\pi R_{rim}}\left(1+\frac{H_{rim}}{R_{rim}}\right)\times \{\begin{array}{cc}1\hfill & \text{non-gray dust}\hfill \\ & \\ 4\mathrm{sin}i\hfill & \text{gray dust}\hfill \end{array}$$
(A10)
where we used the approximation $`q_\lambda \sigma _{IR}/\sigma _V\overline{\sigma }(T_{rim})/\overline{\sigma }(T_{})`$. This solution shows that non-gray dust gives angle-independent bump strength in addition to reducing its magnitude from the gray dust result.
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# Definition 1.1
1. Introduction
In the description of isolated gravitational system in General relativity a space-like time-slice has the structure of a complete Riemannian 3-manifold with an asymptotically flat end. Such asymptotically flat end is diffeomorphic to $`R^3B_1(0)`$ and the metric on it asymptotically approaches the Euclidean metric near the infinity:
$$g_{ij}=(1+\frac{2m}{r})\delta _{ij}+O(r^2),$$
where $`r`$ is the Euclidean distance in $`R^3`$. The constant $`m`$ can be interpreted as the total mass of the isolated system and is referred to as ADM mass in literature \[ADM\]. It has also been established in \[B\] that with reasonable conditions ADM mass can be geometrically defined independent of the choices of coordinate system at infinity.
Often it is better to consider an asymptotically flat end as a perturbation of the static time-slice of the Schwarzchild space-time. Let us start with a precise definition of asymptotically flat 3-manifolds adopted from \[HY\] for our discussions in this note as follows:
###### Definition 1.1
A complete Riemannian 3-manifold $`(M,g)`$ is said to be an asymptotically flat 3-manifold with mass $`m`$ if there is a compact domain $`K`$ of $`M`$ such that $`MK`$ is diffeomorphic to $`R^3B_1(0)`$ and the metric $`g`$ in this coordinate system is given as
$$g_{ij}(x)=(1+\frac{m}{2|x|})^4\delta _{ij}+T_{ij}(x),$$
for all $`xR^3B_1(0)`$ with a constant $`C`$ such that
$$|^lT_{ij}|(x)C|x|^{2l},1l4,$$
$`1.1`$
where $``$ denotes partial derivatives with respect to the Euclidean coordinates.
The existence of a unique foliation of spheres of constant mean curvature near the end in an asymptotically flat manifold is very important question. Among many applications, the unique foliation of spheres of constant mean curvature can be used to construct a geometrically canonical coordinate system at the infinity of asymptotically flat end. It can also be used to define a geometric center of mass for an isolated gravitational system (cf. \[HY\]). The existence of such unique foliation of spheres of constant mean curvature at the asymptotically flat end is also helpful to the study of Penrose inequality regarding the mass (cf. \[Br\]). In this note we show that indeed outside a given compact subset in an asymptotically flat 3-manifold with positive mass there is a unique foliation of stable spheres of constant mean curvature. Our main theorem<sup>2</sup><sup>2</sup>2The uniqueness problem addressed here was referred as the global uniqueness of stable CHC surfaces in \[HY\]. Their result on this global uniqueness was stated in Theorem 5.1 in \[HY\]. They proved that for $`q>\frac{1}{2}`$, if $`H`$ sufficiently small, there is a unique stable constant mean curvature surface of mean curvature $`H`$ outside $`B_{H^q}(0)`$. It has been a long-standing question whether stable constant mean curvature surfaces are unique outside a fixed compact subset. In the paragraph after Theorem 5.1 on page 301, Huisken and Yau stated: “it is an open question whether stable constant mean curvature surfaces are actually completely unique outside a fixed compact subset.” Our main theoerm gives an affirmative answer to this question. is
###### Theorem 1.1
Suppose $`(M,g)`$ is an asymptotically flat 3-manifold with positive mass. Then there exists a compact domain $`K`$ such that stable spheres of given constant mean curvature which separates the infinity from the compact domain $`K`$ are unique. Hence there exists a unique foliation of stable spheres of constant mean curvature outside the compact domain $`K`$ in $`M`$.
The existence of a foliation of stable spheres of constant mean curvature near asymptotically flat ends was established by Huisken and Yau in \[HY\] (also see \[Ye\]). Some uniqueness results with additional assumptions were also proven in \[Br\] \[HY\] \[Ye\]. The major difficulty of establishing the uniqueness of spheres of given constant mean curvature is that possible drifting of the spheres of constant mean curvature presents a hurdle to any useful global a priori estimates on the curvature. As a matter of fact, the uniqueness is known if one assumes no drifting (cf. \[HY\] \[Ye\]). Moreover, it was proven in \[HY\] that, if the drifting was somehow mild, then the uniqueness holds (cf. Theorem 5.1 on page 301 in \[HY\]).
Our main technical contributions can be summarized as follows: First, as a sharp contrast to the Euclidean space, similar to (5.13) in \[HY\], we find the following scale invariant integral which detects the nonzero mass. Suppose that $`N`$ is a surface of constant mean curvature in an asymptotically flat end $`(R^3B_1(0),g)`$ with positive mass $`m`$. Then
$$\frac{1}{8\pi }_N\frac{H}{|x|}\nu b𝑑\sigma +\frac{1}{4\pi }_N\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma Cm^1r_0^1,$$
$`1.2`$
where $`d\sigma `$ is induced from the Euclidean metric, $`C>0`$ is some constant, $`b`$ is any vector in $`R^3`$, $`\nu `$ is the unit out-going normal vector of $`N`$ in $`R^3`$ with respect to the Euclidean metric, and
$$r_0=\mathrm{min}\{|x|:xNR^3B_1(0)\},$$
$`1.3`$
provided that
$$_NH^2𝑑\mu <\mathrm{},$$
$`1.4`$
where $`d\mu `$ is induced from $`g`$. Secondly<sup>3</sup><sup>3</sup>3In \[HY\], a global estimate was sought after (cf. Lemma 5.6 in \[HY\]), with a compromise to assume that the inner radius is not smaller than $`H^q`$ for $`q>\frac{1}{2}`$. They stated in the paragraph after Theorem 5.1 (page 21, \[HY\]) that their assumption on inner radius “seems to be optimal from a technical point of view”. While in this paper we do different estimates in three different scales. Particularly we establish some decay estimate for the intermediate scales by using an asymptotic analysis developed in \[QT\]., we are able to obtain estimates (cf. Corollary 4.4 and Corollary 4.5 in Section 4), which are beyond one individual scale in the blow-down analysis, via an asymptotic analysis used in an early work of us \[QT\]. The blow-down for a surface $`N`$ of constant mean curvature $`H`$ with the scale $`H`$ is defined as,
$$\stackrel{~}{N}=\{\frac{1}{2}Hx:xNR^3B_1(0)\}R^3.$$
$`1.5`$
The use of the asymptotic analysis introduced in Section 4 is the key which allows us to obtain some finer estimates and untangle the problem that uniform roundness and non-drifting of spheres of constant mean curvature hinge on each other. More precisely, to eliminate the possible drifting, one carefully calculates the two integrals in left-hand side of (1.3) for $`\stackrel{~}{N}`$,
$$\frac{1}{4\pi }_{\stackrel{~}{N}}\frac{1}{|x|}\nu b𝑑\sigma +\frac{1}{4\pi }_{\stackrel{~}{N}}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma $$
$`1.6`$
with some particular choice of $`b`$. If drifting happened, then the rescaled surface $`\stackrel{~}{N}`$ would approach the origin. Then one evaluates the integrals over three different regions: 1) the part of $`\stackrel{~}{N}`$ that is any fixed distance away from the origin; 2) the part of $`\stackrel{~}{N}`$ that is near the origin in the scale of $`Hr_0`$; 3) the transition between the above two. We will employ Corollary 4.5 in Section 4 to show the integrals on third region contribute something negligible. Consequently we are able to prove that the drifting of stable spheres of constant mean curvature does not happen at all in an asymptotically flat 3-manifold with positive mass. Then using the early existence and uniqueness results in \[HY\] and \[Ye\], for instance, Theorem 5.1 in \[HY\], we may conclude our main theorem.
It is worthwhile to note that the uniqueness of spheres of given constant mean curvature outside the horizon in the Schwarzchild space is an interesting open problem. In his thesis \[Br\], Bray proved the coordinate spheres are the unique minimizing surfaces of given constant mean curvature outside the horizon in Scwarzchild space, in an attempt to prove the Penrose inequality regarding the mass by the foliation of constant mean curvature surfaces. Theorem 1.1 in the above particularly implies that the coordinate spheres are the only stable sphere of constant mean curvature near the infinity of the Schwarzchild space which separates the infinity from the horizon.
The paper is organized as follows: In Section 2 we will obtain the curvature estimates based on the Simons’ identity and the smallness of the integral of the traceless part of the second fundamental form. In Section 3 we introduce the blow-down analysis in all scales. In Section 4 we recall the asymptotic analysis from \[QT\] and prove a technical proposition. Finally in Section 5 we introduce a sense of the center of mass and prove our main theorem.
2. Curvature estimates
First let us recall the Simons’ identity \[SSY\] \[Sj\] for a hypersurface $`N`$ in a Riemannian manifold $`(M,g)`$ (cf. Lemma 1.3 in \[HY\]):
$`\mathrm{\Delta }h_{ij}`$ $`=_i_jH+Hh_{ik}h_{jk}|A|^2h_{ij}+HR_{3i3j}h_{ij}R_{3k3k}`$ $`2.1`$
$`+h_{jk}R_{klil}+h_{ik}R_{kljl}2h_{lk}R_{iljk}+_jR_{3kik}+_kR_{3ijk}`$
where $`A=(h_{ij})`$ is the second fundamental form for $`N`$ in $`M`$, $`H=\text{Tr}A`$ is mean curvature, and $`R_{ijkl}`$ and $`R_{ijkl}`$ are curvature and covariant derivatives of curvature for $`(M,g)`$. When $`N`$ is a constant mean curvature hypersurface, we rather like to rewrite it as an equation for the traceless part $`\AA `$ of $`A`$, i.e. $`\AA =A\frac{1}{2}H`$.
$`\mathrm{\Delta }\AA _{ij}`$ $`=H\AA _{ik}\AA _{jk}{\displaystyle \frac{1}{2}}H|\AA |^2\delta _{ij}(|\AA |^2+{\displaystyle \frac{1}{2}}H^2)\AA _{ij}`$ $`2.2`$
$`+HR_{3i3j}{\displaystyle \frac{1}{2}}HR_{3k3k}\delta _{ij}\AA _{ij}R_{3k3k}`$
$`+\AA _{jk}R_{klil}+\AA _{ik}R_{kljl}2\AA _{lk}R_{likj}`$
$`+_jR_{3kik}_kR_{3ikj}.`$
###### Lemma 2.1
Suppose that $`N`$ is a constant mean curvature surface in an asymptotically flat end $`(R^3B_1(0),g)`$. Then
$`|\AA |\mathrm{\Delta }|\AA |`$ $`|\AA |^4+CH|\AA |^3+CH^2|\AA |^2`$ $`2.3`$
$`C|\AA |^2|x|^3+CH|\AA ||x|^3+C|\AA ||x|^4.`$
Note that, in an asymptotically flat end (cf. Definition 1.1 in Section 1),
$$|R_{ijkl}|C|x|^3,|R_{ijkl}|C|x|^4.$$
$`2.4`$
We refer readers to \[HY\] for the calculations of curvature of the Schwarzchild space and asymptotically flat ends.
###### Lemma 2.2
Suppose that $`N`$ is a constant mean curvature surface in an asymptotically flat end $`(R^3B_1(0),g)`$. Then $`_NH_e^2𝑑\sigma `$ is bounded if and only if $`_NH^2𝑑\mu `$ is bounded, provided that $`r_0`$ is sufficiently large.
###### Demonstration Proof
First one may calculate
$$H_e=(1+\frac{m}{2r})^2H+2(1+\frac{m}{2r})^1\frac{m}{r^3}x\nu +O(r^3),$$
$`2.5`$
where $`H_e`$ is the mean curvature of $`NR^3`$ with respect to the Euclidean metric (cf. Lemma 1.4 in \[HY\]). Hence
$$H_e^2=H^2+O(r^1)H^2+O(r^2)H+O(r^3).$$
Following Lemma 5.2 in \[HY\] and the fact that $`g`$ is quasi-isometric to the Euclidean metric $`|dx|^2`$, we have:
$`{\displaystyle _N}H_e^2𝑑\sigma `$ $`C{\displaystyle _N}H_e^2𝑑\mu C{\displaystyle _N}H^2𝑑\mu +C({\displaystyle _N}H^2𝑑\mu )^{\frac{1}{2}}({\displaystyle _N}r^4𝑑\mu )^{\frac{1}{2}}+C{\displaystyle _N}r^3𝑑\mu `$
$`C{\displaystyle _N}H^2𝑑\mu `$
and
$$(1Cr_0^1)_NH^2𝑑\mu C_NH_e^2𝑑\sigma .$$
Thus the lemma is proved.
Therefore, following Lemma 1 in \[Si\], we have
###### Lemma 2.3
Suppose that $`N`$ is a constant mean curvature surface in an asymptotically flat end $`(R^3B_1(0),g)`$ with $`r_0(N)`$ sufficiently large, and that
$$_NH^2𝑑\mu C.$$
Then
$$C_1H^1\text{diam}(N)C_2H^1.$$
We would like to point out that, if the surface $`N`$ separates the infinity from the compact part, i.e. the origin is inside $`NR^3`$, then the above lemma implies
$$C_1H^1r_1(N)C_2H^1,$$
$`2.6`$
where the outer radius $`r_1(N)`$ is defined as
$$r_1(N)=\mathrm{max}\{|x|:xNR^3B_1(0)\}.$$
Based on Michael and Simon \[MS\], one has the following Sobolev inequality (cf. Lemma 5.6 in \[HY\]).
###### Lemma 2.4
Suppose that $`N`$ is a constant mean curvature surface in an asymptotically flat end $`(R^3B_1(0),g)`$ with $`r_0(N)`$ sufficiently large, and that
$$_NH^2𝑑\mu C.$$
Then
$$(_Nf^2𝑑\mu )^{\frac{1}{2}}C(_N|f|𝑑\mu +_NH|f|𝑑\mu ).$$
$`2.7`$
Now we are ready to state and prove the main curvature estimates:
###### Theorem 2.5
Suppose that $`(R^3B_1(0),g)`$ is an asymptotically flat end. Then there exist positive numbers $`\sigma _0`$, $`ϵ_0`$ and $`\delta _0`$ such that for any constant mean curvature surface in the end, which separates the infinity from the compact part, we have
$$|\AA |^2(x)C|x|^2_{B_{\delta _0|x|}(x)}|\AA |^2𝑑\mu +C|x|^4,$$
$`2.8`$
provided that
$$_N|\AA |^2𝑑\mu ϵ_0$$
and $`r_0(N)\sigma _0`$. And the corresponding a priori estimates for all covariant derivatives of curvature also hold consequently.
###### Demonstration Proof
Recall that
$`|\AA |\mathrm{\Delta }|\AA |`$ $`|\AA |^4+CH|\AA |^3+CH^2|\AA |^2`$
$`C(|\AA |^2|x|^3+CH|\AA ||x|^3+C|\AA ||x|^4).`$
Multiply the two sides with $`\varphi ^3`$, where $`\varphi `$ is an appropriate cutoff function of small support, and integrate,
$`{\displaystyle _N}\varphi ^3|\AA |\mathrm{\Delta }|\AA |d\mu `$ $`{\displaystyle _N}\varphi ^3|\AA |^4𝑑\mu +C{\displaystyle _N}H\varphi ^3|\AA |^3𝑑\mu +C{\displaystyle _N}H^2\varphi ^3|\AA |^2𝑑\mu `$
$`+Cr_0^1{\displaystyle _N}\varphi ^3(|\AA |^2|x|^2+CH|\AA ||x|^2+C|\AA ||x|^3)𝑑\mu `$
where
$`{\displaystyle _N}\varphi ^3|\AA |`$ $`\mathrm{\Delta }|\AA |d\mu ={\displaystyle _N}(\varphi ^3|\AA |)|\AA |d\mu `$
$`={\displaystyle _N}\varphi ^2(\varphi |\AA |)|\AA |+{\displaystyle _N}2\varphi ^2|\AA |\varphi |\AA |d\mu `$
$`={\displaystyle _N}\varphi |(\varphi |\AA |)|^2𝑑\mu +{\displaystyle _N}\varphi |\AA |(\varphi |\AA |)\varphi d\mu +{\displaystyle _N}2\varphi ^2|\AA |\varphi |\AA |𝑑\mu `$
$`{\displaystyle \frac{3}{4}}{\displaystyle _N}\varphi |(\varphi |\AA |)|^2𝑑\mu C{\displaystyle _N}\varphi |\AA |^2|\varphi |^2𝑑\mu +{\displaystyle _N}2\varphi ^2|\AA |\varphi |\AA |𝑑\mu `$
$`{\displaystyle \frac{1}{2}}{\displaystyle _N}\varphi |(\varphi |\AA |)|^2𝑑\mu C{\displaystyle _N}\varphi |\AA |^2|\varphi |^2𝑑\mu `$
$$_N\varphi ^3|\AA |^4𝑑\mu (_{\text{supp}(\varphi )}|\AA |^2𝑑\mu )^{\frac{1}{2}}(_N(\varphi |\AA |)^6𝑑\mu )^{\frac{1}{2}}$$
and
$$_NH\varphi ^3|\AA |^3𝑑\mu (_{\text{supp}(\varphi )}H^2𝑑\mu )^{\frac{1}{2}}(_N(\varphi |\AA |)^6𝑑\mu )^{\frac{1}{2}}.$$
For other terms
$$_NH^2\varphi ^3|\AA |^2𝑑\mu Cr_M^2_{\text{supp}(\varphi )}|\AA |^2𝑑\mu C|x_0|^2_{\text{supp}(\varphi )}|\AA |^2𝑑\mu ,$$
$$_N\varphi ^3|x|^2|\AA |^2𝑑\mu C|x_0|^2_{\text{supp}(\varphi )}|\AA |^2𝑑\mu ,$$
$$_N\varphi ^3H|x|^2|\AA |𝑑\mu C|x_0|^2(_{\text{supp}(\varphi )}|\AA |^2𝑑\mu )^{\frac{1}{2}},$$
and
$$_N\varphi ^2|x|^3|\AA |𝑑\mu C|x_0|^2(_{\text{supp}(\varphi )}|\AA |^2𝑑\mu )^{\frac{1}{2}}.$$
Note that, for a given point $`x_0`$, we may choose the cutoff function $`\varphi `$ so that it has the suppose of a disk of radius, say, $`\delta _0|x_0|`$ ($`\delta _0`$ to be determined). Now, combining all terms, we have
$`{\displaystyle _N}\varphi |`$ $`(\varphi |\AA |)|^2d\mu 2({\displaystyle _N}|\AA |^2d\mu )^{\frac{1}{2}}({\displaystyle _N}(\varphi |\AA |)^6d\mu )^{\frac{1}{2}}+`$
$`C({\displaystyle _{\text{supp}(\varphi )}}H^2𝑑\mu )^{\frac{1}{2}}({\displaystyle _N}(\varphi |\AA |)^6𝑑\mu )^{\frac{1}{2}}+C|x_0|^2({\displaystyle _{\text{supp}(\varphi )}}|\AA |^2𝑑\mu )^{\frac{1}{2}}.`$
Applying the Sobolev inequality with $`f=\varphi ^3g^3`$ where $`g=|\AA |`$, we have
$`({\displaystyle _N}`$ $`(\varphi g)^6d\mu )^{\frac{1}{2}}C(3{\displaystyle _N}(\varphi g)^2|(\varphi g)|d\mu +{\displaystyle _N}H(\varphi g)^3)`$
$`C({\displaystyle _N}\varphi ^3g^4𝑑\mu )^{\frac{1}{2}}({\displaystyle _N}|(\varphi g)|^2\varphi 𝑑\mu )^{\frac{1}{2}}+({\displaystyle _{\text{supp}(\varphi )}}H^2𝑑\mu )^{\frac{1}{2}}({\displaystyle _N}(\varphi g)^6𝑑\mu )^{\frac{1}{2}}`$
$`C({\displaystyle _N}g^2𝑑\mu )^{\frac{1}{2}}({\displaystyle _N}(\varphi g)^6𝑑\mu )^{\frac{1}{2}}+C({\displaystyle _{\text{supp}(\varphi )}}H^2𝑑\mu )^{\frac{1}{2}}({\displaystyle _N}(\varphi g)^6𝑑\mu )^{\frac{1}{2}}`$
$`+C{\displaystyle _N}|(\varphi g)|^2\varphi 𝑑\mu .`$
Thus
$$(_N(\varphi |\AA |)^6𝑑\mu )^{\frac{1}{2}}C|x_0|^2(_{\text{supp}(\varphi )}|\AA |^2𝑑\mu )^{\frac{1}{2}},$$
$`2.9`$
which implies
$$_N(\varphi |\AA |)^4𝑑\mu C|x_0|^2_{\text{supp}(\varphi )}|\AA |^2𝑑\mu .$$
$`2.10`$
Note that we have chosen $`\delta _0`$ small enough so that
$$C_{\text{supp}(\varphi )}H^2𝑑\mu \frac{1}{8}$$
and $`_N|\AA |^2𝑑\mu ϵ_0`$, where $`ϵ_0`$ is small enough so that
$$C_N|\AA |^2𝑑\mu \frac{1}{8}.$$
Now we proceed to get the point-wise estimates. First, if we take $`f=u^2`$ in the Sobolev inequality, then
$`({\displaystyle _N}u^4𝑑\mu )^{\frac{1}{2}}`$ $`C(2{\displaystyle _N}|u||u|𝑑\mu +{\displaystyle _N}Hu^2𝑑\mu )`$
$`C({\displaystyle _N}u^2𝑑\mu )^{\frac{1}{2}}({\displaystyle _N}|u|^2𝑑\mu )^{\frac{1}{2}}+C({\displaystyle _{\text{supp}(u)}}H^2𝑑\mu )^{\frac{1}{2}}({\displaystyle _N}u^4𝑑\mu )^{\frac{1}{2}}.`$
When $`u`$ has the support as the cutoff function $`\varphi `$, we have
$$(_Nu^4𝑑\mu )^{\frac{1}{2}}C(_Nu^2𝑑\mu )^{\frac{1}{2}}(_N|u|^2𝑑\mu )^{\frac{1}{2}}.$$
$`2.11`$
To finish the point-wise estimates we use the following rather standard estimate:
###### Lemma 2.6
Suppose that a nonnegative function $`v`$ in $`L^2`$ solves
$$\mathrm{\Delta }vfv+h$$
$`2.12`$
on $`B_{2R}(x_0)`$, where
$$_{B_{2R}(x_0)}f^2𝑑\mu CR^2$$
and $`hL^2(B_{2R}(x_0))`$. And Suppose that
$$(_Nu^4𝑑\mu )^{\frac{1}{2}}C(_Nu^2𝑑\mu )^{\frac{1}{2}}(_N|u|^2𝑑\mu )^{\frac{1}{2}}$$
holds for all $`u`$ with support inside $`B_{2R}(x_0)`$. Then
$$\underset{B_R(x_0)}{sup}vCR^1v_{L^2(B_{2R}(x_0))}+CRh_{L^2(B_{2R}(x_0))}.$$
###### Demonstration Proof
We will simply use the Moser iteration method. For convenience, we may rescale so that we are working on $`B_2`$. The correct scales would be
$$v_R(x)=v(Rx),f_R(x)=R^2f(Rx),\text{and}h_R=R^2h(Rx).$$
Let $`k=h_{L^2(B_2)}`$ and $`\overline{v}=v+k`$. Multiply the equation with $`\varphi ^2\overline{v}^{p1}`$ on the both sides
$$|(\varphi \overline{v}^{\frac{p}{2}})|^2\frac{2}{p}f\varphi ^2\overline{v}^p+h\varphi ^2\overline{v}^{p1}+C|\varphi |^2\overline{v}^p.$$
Set $`\overline{f}=\frac{h}{k}+f`$ we have
$$|(\varphi \overline{v}^{\frac{p}{2}})|^2\frac{2}{p}\overline{f}\varphi ^2\overline{v}^p+C|\varphi |^2\overline{v}^p.$$
Note that $`\overline{f}_{L^2(B_2)}1+f_{L^2(B_2)}`$. By the assumed Sobolev inequality, we
$$((\varphi \overline{v}^{\frac{p}{2}})^4)^{\frac{1}{2}}C(p|\overline{f}|\varphi ^2\overline{v}^p)^{\frac{1}{2}}(\varphi ^2\overline{v}^p)^{\frac{1}{2}}+C(|\varphi |^2v^p)^{\frac{1}{2}}(\varphi ^2v^p)^{\frac{1}{2}}.$$
To handle the first term, we apply Hólder inequality
$`(p{\displaystyle |\overline{f}|\varphi ^2\overline{v}^p})^{\frac{1}{2}}({\displaystyle \varphi ^2\overline{v}^p})^{\frac{1}{2}}`$ $`p^{\frac{1}{2}}({\displaystyle |\overline{f}|^2})^{\frac{1}{4}}({\displaystyle \varphi ^2\overline{v}^p})^{\frac{1}{2}}({\displaystyle (\varphi \overline{v}^{\frac{p}{2}})^4})^{\frac{1}{4}}`$
$`{\displaystyle \frac{1}{2C}}({\displaystyle (\varphi \overline{v}^{\frac{p}{2}})^4})^{\frac{1}{2}}+Cp({\displaystyle |\overline{f}|^2})^{\frac{1}{2}}{\displaystyle \varphi ^2\overline{v}^p}.`$
Hence
$$((\varphi \overline{v}^{\frac{p}{2}})^4)^{\frac{1}{2}}C(p\overline{f}_{L^2}\varphi ^2\overline{v}^p+|\varphi |^2\overline{v}^p).$$
Now, for $`i=1,2,\mathrm{},`$ let $`p=2^i`$ and
$$\varphi =\{\begin{array}{cc}\hfill 1& xB_{1+2^i}\hfill \\ \hfill 0& xB_{1+2^{i+1}}\hfill \end{array}$$
Then
$$(_{B_{1+2^i}}\overline{v}^{2^{i+1}})^{2^{i1}}C^{2^i}2^{i2^i}(_{B_{1+2^{i+1}}}\overline{v}^{2^i})^{2^i}.$$
Thus
$$\underset{B_1}{sup}v\underset{B_1}{sup}\overline{v}C^{_{i=1}2^i}2^{_{i=1}i2^i}(_{B_2}\overline{v}^2)^{\frac{1}{2}}C(v_{L^2(B_2)}+h_{L^2(B_2)}),$$
whose scaled version gives the lemma.
To get curvature estimates, we write the equation in such way as (2.12) that we may apply the above lemma for
$$f=C(|\AA |^2+H|\AA |+H^2+r^3)\text{and}h=C(Hr^3+r^4),$$
in the light of (2.9) and (2.10).
3. Blow-down analysis
In order to understand a surface of constant mean curvature $`N`$ in an asymptotically flat end $`(R^3B_1(0),g)`$, we will need to blow down the surface in different scales. We first consider, the blow-down by the scale $`H`$,
$$\stackrel{~}{N}=\frac{1}{2}HN=\{\frac{1}{2}Hx:xN\}.$$
$`3.1`$
Suppose that there is a sequence of constant mean curvature surfaces $`\{N_i\}`$ such that
$$\underset{i\mathrm{}}{lim}r_0(N_i)=\mathrm{}\text{and}\underset{i\mathrm{}}{lim}_{N_i}H^2𝑑\mu =16\pi .$$
$`3.2`$
Then, by a similar argument to the proof of Lemma 2.2 in the previous section, we have
$$\underset{i\mathrm{}}{lim}_{N_i}H_e^2𝑑\sigma =16\pi .$$
$`3.3`$
Hence, by the curvature estimates established in the previous section combining the proof of Theorem 1 in \[Si\], we have
###### Lemma 3.1
Suppose that $`\{N_i\}`$ is a sequence of constant mean curvature surfaces in a given asymptotically flat end $`(R^3B_1(0),g)`$ and that
$$\underset{i\mathrm{}}{lim}r_0(N_i)=\mathrm{}\text{and}\underset{i\mathrm{}}{lim}_{N_i}H^2𝑑\mu =16\pi .$$
And suppose that $`N_i`$ separates the infinity from the compact part. Then, there is a subsequence of $`\{\stackrel{~}{N}_i\}`$ which converges in Gromov-Hausdorff distance to a round sphere $`S_1^2(a)`$ of radius $`1`$ and centered at $`aR^3`$. Moreover, the convergence is in $`C^{\mathrm{}}`$ sense away from the origin.
¿From the above lemma, the difficulty will be to study the possibility of having the origin lying on the sphere $`S^2(a)`$, that is,
$$\underset{i\mathrm{}}{lim}r_0(N_i)=\mathrm{},\text{and}\underset{i\mathrm{}}{lim}r_0(N_i)H(N_i)=0.$$
$`3.5`$
Then, in the light of the curvature estimates we obtained in the previous section, we may use the smaller scale $`r_0(N_i)`$ to blow down the surface
$$\widehat{N}=r_0(N)^1N=\{r_0^1x:xN\}.$$
$`3.6`$
###### Lemma 3.2
Suppose that $`\{N_i\}`$ is a sequence of constant mean curvature surfaces in a given asymptotically flat end $`(R^3B_1(0),g)`$ and that
$$\underset{i\mathrm{}}{lim}r_0(N_i)=\mathrm{}\text{and}\underset{i\mathrm{}}{lim}_{N_i}H^2𝑑\mu =16\pi .$$
And suppose that
$$\underset{i\mathrm{}}{lim}r_0(N_i)H(N_i)=0.$$
Then there is a subsequence of $`\{\widehat{N}_i\}`$ converges to a 2-plane at distance $`1`$ from the origin. Moreover the convergence is in $`C^{\mathrm{}}`$ in any compact set of $`R^3`$.
As one would expect, the real difficulty is to understand the behavior of the surfaces $`N_i`$ in the scales between $`r_0(N_i)`$ and $`H^1(N_i)`$. To start we consider the intermediate scales $`r_i`$ such that
$$\underset{i\mathrm{}}{lim}\frac{r_0(N_i)}{r_i}=0\text{and}\underset{i\mathrm{}}{lim}r_iH(N_i)=0$$
$`3.7`$
and blow down the surfaces
$$\overline{N_i}=r_i^1N=\{r_i^1x:xN\}.$$
$`3.8`$
###### Lemma 3.3
Suppose that $`\{N_i\}`$ is a sequence of constant mean curvature surfaces in a given asymptotically flat end $`(R^3B_1(0),g)`$ and that
$$\underset{i\mathrm{}}{lim}r_0(N_i)=\mathrm{}\text{and}\underset{i\mathrm{}}{lim}_{N_i}H^2𝑑\mu =16\pi .$$
And suppose that $`\{r_i\}`$ are such that
$$\underset{i\mathrm{}}{lim}\frac{r_0(N_i)}{r_i}=0\text{and}\underset{i\mathrm{}}{lim}r_iH(N_i)=0.$$
Then there is a subsequence of $`\{\overline{N}_i\}`$ converges to a 2-plane at the origin in Gromov-Hausdorff distance. Moreover the convergence is $`C^{\mathrm{}}`$ in any compact subset away from the origin.
4. Asymptotic analysis
In this section we would like to apply the asymptotic analysis used in \[QT\] to obtain some estimate that holds over the whole transition region between the scales $`r_0(N_i)`$ and $`r_1(N_i)`$. But first, let us revise Proposition 2.1 in \[QT\] as follows. Let us denote
$$u_i^2=_{[(i1)L,iL]\times S^1}|u|^2𝑑t𝑑\theta .$$
###### Lemma 4.1
Suppose $`uW^{1,2}(\mathrm{\Sigma },R^k)`$ satisfies
$$\mathrm{\Delta }u+Au+Bu=h\text{in }\mathrm{\Sigma },$$
where $`\mathrm{\Sigma }=[0,3L]\times S^1`$. And suppose that $`L`$ is given and large. Then there exists a positive number $`\delta _0`$ such that, if
$$h_{L^2(\mathrm{\Sigma })}\delta _0\underset{1i3}{\mathrm{max}}\{u_i\}$$
and
$$A_{L^{\mathrm{}}(\mathrm{\Sigma })}\delta _0,B_{L^{\mathrm{}}(\mathrm{\Sigma })}\delta _0,$$
then,
(a). $`u_3e^{\frac{1}{2}L}u_2`$ implies $`u_2<e^{\frac{1}{2}L}u_1`$,
(b). $`u_1e^{\frac{1}{2}L}u_2`$ implies $`u_2<e^{\frac{1}{2}L}u_3`$, and
(c). If both $`_{L\times S^1}u𝑑\theta `$ and $`_{2L\times S^1}u𝑑\theta \delta _0\mathrm{max}_{1i3}\{u_i\}`$, then
either $`u_2<e^{\frac{1}{2}L}u_1`$ or $`u_2<e^{\frac{1}{2}L}u_3`$.
###### Demonstration Proof
We refer to the proof of Proposition 2.1 in \[QT\] for more details. In the proof by contradiction argument one needs to make sure that the sequence of normalized $`u_k`$ converges to a non-zero harmonic function $`u`$ and the non-zero harmonic function $`u`$ violates one of (a)-(c). Interior elliptic estimates give the strong convergence in the middle section $`I_2=[L,2L]\times S^1`$, which implies that $`u`$ is not trivially zero. Because, with the assumption of the proof by contradiction, the middle one is the largest. Finally $`u`$ indeed induces a contradiction due to the Fatou lemma.
We would like to point out that Proposition 2.1 in \[QT\] is overstated since it is not correct for $`l>3`$. But, in the proof of Proposition 3.1 in \[QT\], where Corollary 2.2 is used, one may replace the shifting cylinder with length $`3L`$ instead of $`5L`$. The proof still works the same, which is, one push to the direction of growth the cylinder of length $`3L`$ when Corollary 2.2 in \[QT\] applies and it gives the estimates regardless of where one is stopped applying Corollary 2.2.
Given a surface $`N`$ in $`R^3`$, recall from, for example, (8.5) in \[Ka\], that
$$\mathrm{\Delta }\nu +|\nu |^2\nu =H_e$$
$`4.2`$
where $`\nu `$ is the Gauss map from $`NS^2`$. For the constant mean curvature surfaces in the asymptotically flat end $`(R^3B_1(0),g)`$,
$$|H_e|(x)C|x|^3.$$
$`4.3`$
Therefore we consider that the Gauss map of the constant mean curvature surfaces in the asymptotically flat end $`(R^3B_1(0),g)`$ is an almost harmonic map. Hence we are in a situation which is very similar to that in \[QT\]. We will refer readers to \[QT\] for rather elementary yet involved analysis since the proof we present here is some modifications from the proof in \[QT\]. We will not carry the indices for the surfaces $`N_i`$ if it does not cause any confusion. Set
$$A_{r_1,r_2}=\{xN:r_1|x|r_2\}.$$
$`4.4`$
$`A_{r_1,r_2}^0`$ stands for the standard annulus in $`R^2`$. We are concerned with the behavior of $`\nu `$ on the part $`A_{Kr_0(N),sH^1(N)}`$ of $`N`$ where $`K`$ will be fixed large and $`s`$ will be fixed small. The first difference from \[QT\] is that, while we had a fixed domain in \[QT\], we need the following lemma in order to be in the position to use Lemma 4.1 in the above.
###### Lemma 4.2
Suppose that $`N`$ is a constant mean curvature surfaces in a given asymptotically flat end $`(R^3B_1(0),g)`$ . Then, for any $`ϵ>0`$ and $`L`$ fixed, there are $`ϵ_0,s`$ and $`K`$ such that, if
$$_N|\AA |^2𝑑\sigma ϵ_0$$
and $`Kr_0(N)<r<sH^1(N)`$, then $`(r^1A_{r,e^{4L}r},r^2g_e)`$ may be represented as $`(A_{1,e^{4L}}^0,\overline{g})`$ where
$$\overline{g}|dx|^2_{C^1(A_{1,e^{4L}}^0)}ϵ.$$
$`4.5`$
In other words, in the cylindrical coordinates $`(S^1\times [\mathrm{log}r,4L+\mathrm{log}r],\overline{g}_c)`$,
$$\overline{g}_c(dt^2+d\theta ^2)_{C^1(S^1\times [\mathrm{log}r,4L+\mathrm{log}r])}ϵ.$$
$`4.6`$
This is a consequence of Lemma 3.3 in the previous section. Another difference from \[QT\] is that, we are considering maps with tension fields possibly blowing up at point but no energy concentration, while in \[QT\] we were considering almost harmonic maps with concentration of energy but tension fields uniformly bounded in $`L^2`$. In cylindrical coordinates, the tension fields
$$|\tau (\nu )|=r^2|H_e|Cr^1=Ce^t$$
$`4.7`$
for $`t[\mathrm{log}(Kr_0),\mathrm{log}(sH^1]`$. Thus,
$$_{S^1\times [t,t+L]}|\tau (\nu )|^2𝑑t𝑑\theta Ce^{2t}$$
$`4.8`$
which decays as needed in the arguments in \[QT\]. But to get the growth (or decay) of the energy along the cylinder we first can only have the estimate (3.8) in \[QT\]. Then we need to use the Hopf differential
$$\mathrm{\Phi }=|_t\nu |^2|_\theta \nu |^22\sqrt{1}_t\nu _\theta \nu $$
and the stationary property, in complex variable $`z=t+\sqrt{1}\theta `$,
$$\overline{}\mathrm{\Phi }=\nu \tau (\nu )$$
$`4.9`$
to bound $`|_t\nu |^2`$ by $`|_\theta \nu |^2`$ (cf. \[QT\] \[DT\]) as follows:
$$_{S^1\times [t,t+L]}|_t\nu |^2𝑑t𝑑\theta _{S^1\times [t,t+L]}|\mathrm{\Phi }|𝑑t𝑑\theta +_{S^1\times [t,t+L]}|_\theta \nu |^2𝑑t𝑑\theta .$$
By the elliptic estimates (cf. \[DT\]), we have
$$_{S^1\times [t,t+L]}|\mathrm{\Phi }|𝑑t𝑑\theta _{NB_r^c}|\mathrm{\Phi }|𝑑t𝑑\theta C(_{NB_r^c}|\nu |^2𝑑\sigma )^{\frac{1}{2}}(_{NB_r^c}|\tau (\nu )|^2𝑑\sigma )^{\frac{1}{2}},$$
where $`NB_r^c`$ is the part of $`N`$ which is outside of $`B_r`$ and is a disk since $`N`$ is a sphere topologically. Hence, we have
$$_{S^1\times [t,t+L]}|\mathrm{\Phi }|𝑑t𝑑\theta C(_{NB_r^c}|\tau (\nu )|^2𝑑\sigma )^{\frac{1}{2}}Ce^t.$$
$`4.10`$
Notice that in \[QT\] we instead used the fact that the tension fields is uniformly bounded in $`L^2`$ inside $`B_\delta `$ (cf. lines between (3.8) and (3.9) in \[QT\]). The rest of the proof of Proposition 3.1 in \[QT\] works with little modifications. Thus we have
###### Proposition 4.3
Suppose that $`\{N_i\}`$ is a sequence of constant mean curvature surfaces in a given asymptotically flat end $`(R^3B_1(0),g)`$ and that
$$\underset{i\mathrm{}}{lim}r_0(N_i)=\mathrm{}\text{and}\underset{i\mathrm{}}{lim}_{N_i}H^2𝑑\mu =16\pi .$$
And suppose that
$$\underset{i\mathrm{}}{lim}r_0(N_i)H(N_i)=0.$$
Then there exist a large number $`K`$, a small number $`s`$ and $`i_0`$ such that, when $`ii_0`$,
$$\underset{I_j}{\mathrm{max}}|\nu |C(_{B_{sH^1(N_i)}N_i}|\nu |^2𝑑\sigma +r_0^1)(e^{\frac{1}{4}jL}+e^{\frac{1}{4}(n_ij)L}),$$
$`4.11`$
where
$$I_j=S^1\times [\mathrm{log}(Kr_0(N_i))+jL,\mathrm{log}(Kr_0(N_i))+(j+1)L]$$
and
$$j[0,n_i]\text{and}\mathrm{log}(Kr_0(N_i))+(n_i+1)L=\mathrm{log}(sH(N_i)^1).$$
This finer analysis improves our understanding of the blow-downs that we discussed in the previous section. Namely,
###### Corollary 4.4
Assume the same conditions as Proposition 4.3. Then the limit plane in Lemma 3.2 and the limit plane in Lemma 3.3 are all orthogonal to the vector $`a`$. In fact, we may choose $`s`$ small and $`i`$ large enough so that,
$$|\nu (x)+a|ϵ$$
for all $`xN_i`$ and $`|x|sH^1(N_i)`$.
And we have
###### Corollary 4.5
Asume the same condition as Proposition 4.3. Let $`\nu _i=\nu (p_i)`$ for some $`p_iI_{\frac{n_i}{2}}`$. Then
$`\underset{I_j}{\mathrm{max}}`$ $`|\nu \nu _i|`$ $`4.12`$
$`C{\displaystyle \frac{1}{1e^{\frac{1}{4}L}}}({\displaystyle _{B_{sH^1(N_i)}{\scriptscriptstyle N_i}}}`$ $`|\nu |^2d\sigma +r_0^1)(e^{\frac{1}{4}jL}+e^{\frac{1}{8}n_iL})`$
for $`j[0,\frac{1}{2}n_i]`$ and
$`\underset{I_j}{\mathrm{max}}`$ $`|\nu \nu _i|`$ $`4.13`$
$`C{\displaystyle \frac{1}{1e^{\frac{1}{4}L}}}({\displaystyle _{B_{sH^1(N_i)}{\scriptscriptstyle N_i}}}`$ $`|\nu |^2d\sigma +r_0^1)(e^{\frac{1}{8}n_iL}+e^{\frac{1}{4}(n_ij)L})`$
for $`j[\frac{1}{2}n_i,n_i]`$.
The two corollary above will be the key for us to calculate the integrals in next section to prove our main theorem.
5. Center of mass
First let us recall that, for any embedded surface $`N`$ in $`R^3`$ and any given vector $`bR^3`$,
$$_NH_e\nu b𝑑\sigma =0.$$
$`5.1`$
One may consider this as the first variation of the area of surface $`N_t=N+tbR^3`$. On the other hand, if $`N`$ is a constant mean curvature surface in the asymptotically flat end $`(R^3B_1(0),g)`$, then
$$_NH\nu b𝑑\sigma =H_N\nu b𝑑\sigma =0.$$
$`5.2`$
Since the flux is zero across any surface for a given constant velocity $`b`$. Thus, for any constant mean curvature surface in the asymptotically flat end,
$$_N(H_eH)\nu b𝑑\sigma =0.$$
$`5.3`$
One may calculate and find
$$H_eH=m(\frac{H}{|x|}+2\frac{\nu x}{|x|^3})+O(|x|^2)H+O(|x|^3).$$
$`5.4`$
###### Lemma 5.1
Suppose $`N`$ is a surface of constant mean curvature in the asymptotically flat end with positive mass $`m0`$. And suppose that
$$_NH^2𝑑\mu <\mathrm{}$$
and $`r_0(N)`$ is sufficiently large. Then for any given $`b`$ and for some $`C>0`$,
$$\frac{1}{8\pi }_N\frac{H}{|x|}\nu b𝑑\sigma +\frac{1}{4\pi }_N\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma Cm^1r_0^1.$$
$`5.5`$
###### Demonstration Proof
Simply multiply $`\nu b`$ to the both sides of (5.4) and integrate over the surface $`N`$, we have,
$$\frac{1}{8\pi }_N(H_eH)\nu b𝑑\sigma =\frac{m}{8\pi }_N\frac{H}{|x|}\nu b𝑑\sigma +\frac{m}{4\pi }_N\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma +O(r_0^1).$$
Here we used the Lemma 5.2 in \[HY\]. Then the lemma is proved due to (5.3).
Now, we are ready to state and prove our main theorem in this note as follows:
###### Theorem 5.2
Suppose that $`\{N_i\}`$ is a sequence of spheres of constant mean curvature in a given asymptotically flat end with positive mass $`m0`$ and that
$$\underset{i\mathrm{}}{lim}r_0(N_i)=\mathrm{}\text{and}\underset{i\mathrm{}}{lim}_{N_i}H^2𝑑\sigma =16\pi .$$
And suppose that $`N_i`$ separates the infinity from the compact part. Then
$$\underset{i\mathrm{}}{lim}\frac{r_0(N_i)}{r_1(N_i)}=1.$$
$`5.6`$
###### Demonstration Proof
We may apply Lemma 3.1 for the blow-down
$$\stackrel{~}{N}=\frac{1}{2}HN=\{\frac{1}{2}Hx:xN\}.$$
If the surfaces $`\stackrel{~}{N}_i`$ stay away from the origin, i.e.
$$0<CH_i^1r_0(N_i)$$
for some positive constants $`C`$, then a subsequence of $`\stackrel{~}{N}_i`$ converges to a sphere $`S^2(a)`$ radius $`1`$ and centered at $`aR^3`$ in $`C^{\mathrm{}}`$ by the curvature estimates Theorem 2.5 in Section 2. Also notice that (2.6) implies that the blow-down surfaces $`\stackrel{~}{N}_i`$ always stay within a bounded region in $`R^3`$. On one hand, by (5.5) in Lemme 5.1, we have
$$\frac{1}{4\pi }_{S^2(a)}\frac{\nu b}{|x|}𝑑\sigma +\frac{1}{4\pi }_{S^2(a)}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma =0,$$
$`5.7`$
for any $`b`$. On the other hand, if take $`b=\frac{a}{|a|}`$, then
$$\frac{1}{4\pi }_{S^2(a)}\frac{\nu b}{|x|}𝑑\sigma +\frac{1}{4\pi }_{S^2(a)}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma =|a|$$
$`5.8`$
due to an explicit calculation when the origin is inside. Therefore
$$a=0\text{and}\underset{i\mathrm{}}{lim}\frac{1}{2}r_0(N_i)H(N_i)=\underset{i\mathrm{}}{lim}\frac{r_0(N_i)}{r_1(N_i)}=1.$$
To conclude this is all that can happen we need only to exclude the case when
$$\underset{i\mathrm{}}{lim}H_i^1r_0(N_i)=0.$$
$`5.9`$
Assume otherwise, according to Lemma 3.1, the blow-down sequence $`\stackrel{~}{N}_i`$ converges to a unit round sphere $`S^2(a)`$ centered at $`aR^3`$ with $`|a|=1`$ in Hausdorff topology. We will take $`b=\frac{a}{|a|}`$. From Lemma 5.1, we know
$$\underset{i\mathrm{}}{lim}(_{\stackrel{~}{N}_i}\frac{\nu b}{|x|}𝑑\sigma +_{\stackrel{~}{N}_i}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma )=0.$$
$`5.10`$
But, we claim, on the other hand,
$$\underset{i\mathrm{}}{lim}(_{\stackrel{~}{N}_i}\frac{\nu b}{|x|}𝑑\sigma +_{\stackrel{~}{N}_i}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma )=4\pi $$
$`5.11`$
which gives us the contradiction. First, we have from explicit calculations
$$_{S^2(a)}\frac{\nu b}{|x|}𝑑\sigma =\frac{4}{3}\pi ,_{S^2(a)}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma =\frac{2}{3}\pi .$$
$`5.12`$
The first term in (5.11) is an easy term because the uniform integrability
$$\underset{i\mathrm{}}{lim}_{\stackrel{~}{N}_i}\frac{\nu b}{|x|}𝑑\sigma =_{S^2(a)}\frac{\nu b}{|x|}𝑑\sigma =\frac{4}{3}\pi .$$
$`5.13`$
To deal with the second term in (5.11), we break up the integral into three parts. For any fixed small number $`s>0`$ and large number $`K>0`$,
$`{\displaystyle _{\stackrel{~}{N}_i}}{\displaystyle \frac{(\nu x)(\nu b)}{|x|^3}}𝑑\sigma ={\displaystyle _{\stackrel{~}{N}_i{\scriptscriptstyle B_s^c(0)}}}`$ $`{\displaystyle \frac{(\nu x)(\nu b)}{|x|^3}}d\sigma `$ $`5.14`$
$`+{\displaystyle _{\stackrel{~}{N}_i{\scriptscriptstyle (B_s(0)B_{KHr_0}(0))}}}{\displaystyle \frac{(\nu x)(\nu b)}{|x|^3}}𝑑\sigma `$ $`+{\displaystyle _{\stackrel{~}{N}_i{\scriptscriptstyle B_{KHr_0}(0)}}}{\displaystyle \frac{(\nu x)(\nu b)}{|x|^3}}𝑑\sigma .`$
Then
$$\underset{i\mathrm{}}{lim}_{\stackrel{~}{N}_i{\scriptscriptstyle B_s^c(0)}}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma =_{S^2(a){\scriptscriptstyle B_s^c}}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma $$
$`5.15`$
and
$$\underset{i\mathrm{}}{lim}_{\stackrel{~}{N}_i{\scriptscriptstyle B_{KH\left(N_i\right)r_0\left(N_i\right)}(0)}}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma =_{P{\scriptscriptstyle B_K(0)}}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma ,$$
where $`P`$ is the limit plane in Lemma 3.2. By Corollary 4.4, we know
$$_P\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma =2\pi .$$
due to a simple calculation. Notice that
$$_{\stackrel{~}{N}_i}\frac{\nu x}{|x|^3}𝑑\sigma =4\pi $$
$`5.16`$
for any $`i`$ and
$$_{S^2(a)}\frac{\nu x}{|x|^3}𝑑\sigma =2\pi $$
$`5.17`$
because the origin is on the sphere $`S^2(a)`$. Since
$$\underset{i\mathrm{}}{lim}_{\stackrel{~}{N}_i{\scriptscriptstyle B_s^c(0)}}\frac{\nu x}{|x|^3}𝑑\sigma =_{S^2(a){\scriptscriptstyle B_s^c(0)}}\frac{\nu x}{|x|^3}𝑑\sigma ,$$
$`5.18`$
$$\underset{i\mathrm{}}{lim}_{\stackrel{~}{N}_i{\scriptscriptstyle B_{KHr_0}(0)}}\frac{\nu x}{|x|^3}𝑑\sigma =_{P{\scriptscriptstyle B_K(0)}}\frac{\nu x}{|x|^3}𝑑\sigma $$
$`5.19`$
and
$$_P\frac{\nu x}{|x|^3}𝑑\sigma =2\pi ,$$
$`5.20`$
we know
$$\underset{i\mathrm{},s0,K\mathrm{}}{lim}_{\stackrel{~}{N}_i{\scriptscriptstyle }(B_s^c(0)B_{KHr_0}(0)}\frac{\nu x}{|x|^3}𝑑\sigma =0.$$
$`5.21`$
Now we are ready to handle the difficult term: the integral over the transition region in (5.14). Our goal is to show that
$$\underset{i\mathrm{},s0,K\mathrm{}}{lim}_{\stackrel{~}{N}_i{\scriptscriptstyle (B_s(0)B_{KHr_0}(0))}}\frac{(\nu x)(\nu b)}{|x|^3}𝑑\sigma =0.$$
$`5.22`$
The key point is to use Corollary 4.5 to prove (5.22) from (5.21). Let $`\nu _i`$ be chosen as in Corollary 4.5. Then
$`{\displaystyle _{\stackrel{~}{N}_i{\scriptscriptstyle (B_sB_{KHr_0})}}}`$ $`{\displaystyle \frac{(\nu x)(\nu b)}{|x|^3}}=`$ $`5.23`$
$`(\nu _ib){\displaystyle _{\stackrel{~}{N}_i{\scriptscriptstyle (B_sB_{KHr_0})}}}{\displaystyle \frac{\nu x}{|x|^3}}`$ $`+{\displaystyle _{\stackrel{~}{N}_i{\scriptscriptstyle (B_sB_{KHr_0})}}}{\displaystyle \frac{(\nu x)((\nu \nu _i)b)}{|x|^3}}.`$
Hence we only need to deal with the second term on the right side of the above (5.23). We are better now to use the cylindrical coordinates used in Section 4.
$`{\displaystyle _{\stackrel{~}{N}_i{\scriptscriptstyle (B_s(0)B_{KHr_0}(0))}}}{\displaystyle \frac{(\nu x)((\nu \nu _i)b)}{|x|^3}}𝑑\sigma `$ $`5.24`$
$`={\displaystyle \underset{j=1}{\overset{n_i}{}}}{\displaystyle _{I_j}}{\displaystyle \frac{(\nu x)((\nu \nu _i)b)}{|x|^3}}A(t)𝑑\theta 𝑑t`$
$`C{\displaystyle \underset{j=1}{\overset{n_i}{}}}L\underset{I_j}{\mathrm{max}}|\nu \nu _i|`$
$`=C{\displaystyle \underset{j=1}{\overset{n_i/2}{}}}L\underset{I_j}{\mathrm{max}}|\nu \nu _i|+C{\displaystyle \underset{j=n_i/2+1}{\overset{n_i/2}{}}}L\underset{I_j}{\mathrm{max}}|\nu \nu _i|.`$
¿From (4.12) and (4.13), we have
$`{\displaystyle _{\stackrel{~}{N}_i{\scriptscriptstyle (B_s(0)B_{KHr_0}(0))}}}{\displaystyle \frac{(\nu x)((\nu \nu _i)b)}{|x|^3}}𝑑\sigma `$ $`5.25`$
$`C\eta ({\displaystyle \underset{j=1}{\overset{n_i/2}{}}}(e^{\frac{1}{4}Lj}+e^{\frac{1}{8}Ln_i})+{\displaystyle \underset{j=1}{\overset{n_i/2}{}}}(e^{\frac{1}{8}Ln_i}+e^{\frac{1}{4}(n_ij)L}))`$
$`C\eta (n_ie^{\frac{1}{8}n_iL}+2),`$
where
$$\eta =_{B_{sH^1(N_i)}N_i}|\nu |^2𝑑\sigma +r_0^1$$
and $`\eta `$ can be arbitrarily small as long as $`s0`$ and $`r_0(N_i)\mathrm{}`$. Thus (5.22) is proved and the proof of the theorem is completed.
###### Corollary 5.3
Suppose $`(R^3B_1(0),g)`$ is an asymptotically flat end with positive mass. Then there exist a large number $`K>0`$ and a small number $`ϵ>0`$ such that, for any $`H<ϵ`$, there exists a unique stable spheres $`N`$ of constant mean curvature $`H`$ with $`NR^3B_K(0)`$ and separates the infinity from the compact part. Hence there exists a unique foliation of stable spheres of constant mean curvature near the infinity.
###### Demonstration Proof
In the light of Proposition 5.3 in \[HY\] we know
$$\underset{i\mathrm{}}{lim}_{N_i}H^2𝑑\mu =16\pi ,$$
$`5.26`$
provided each $`N_i`$ is a stable sphere of constant mean curvature. Thus Corollary 5.3 follows from Theorem 5.2 in the above, Theorem 4.1 in \[HY\] and the proof of Theorem 5.1 in \[HY\].
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# Shock waves in a one-dimensional Bose gas: from a Bose-Einstein condensate to a Tonks gas
## I Introduction
Physics of one dimensional (1D) Bose gases attracts more and more attention due to both very interesting phenomena that appear in these systems korepin , and a very recent experimental progress toward realization of Tonks gases esslinger ; paredes ; raizen . In this paper we provide a unified description of basic shock waves properties in 1D Bose gases interacting with arbitrary strength. Previous works on this subject were limited to either weakly interacting gases bodzio\_b ; gammal ; konotop ; menotti ; el or the strongly interacting ones bodzio\_f . We also describe influence of depleted atoms on dynamics of a Bose-Einstein condensate (BEC) shocks. Finally, we derive a simple expression for speed of propagation of arbitrarily shaped density pulses.
Sec. II describes theoretical basics of an approximate hydrodynamic approach. Sec. III discusses shock-wave solutions of this approach on specific examples. Sec. IV describes propagation of arbitrarily shaped pulses. Sec. V explains experimental methods used for generation of shock structures, while Sec. VI presents quantum corrections to shock solutions of the Gross-Pitaevskii equation. Finally, Sec. VII provides a summary of this paper.
## II The model
We consider delta interacting Bose gas in a 1D box. Its experimental realization is possible due to recent experimental progress in trapping of small samples of ultracold bosonic atoms in a box-like optical trap raizen . Another exciting experimental setup may come from a paper of Gupta et al gupta , where successful realization of a ring shaped magnetic trap filled with a BEC was recently reported.
The Hamiltonian of our system in dimensionless quantities (see Appendix A for units) reads:
$$\widehat{H}=\frac{1}{2}\underset{i=1}{\overset{N}{}}\frac{^2}{x_i^2}+a\underset{i<j}{}\delta (x_ix_j),$$
(1)
where $`a>0`$ is the interaction coupling olshanii . It turns out that there is only one parameter that determines system static properties:
$$\gamma =\frac{a}{\rho },$$
(2)
where $`\rho `$ is atomic density normalized to $`N1`$. In the limit of $`\gamma 0`$ our system is a Bose-Einstein condensate, while when $`\gamma +\mathrm{}`$ it is a Tonks gas, whose properties are captured by the Fermi-Bose mapping theorem (FBMT) bfmt . To see how big is a BEC-Tonks crossover one can look at sound velocity and compare it to BEC and Tonks predictions. Exact expression for sound velocity lieb (see also Appendix A) is
$$v_s=\rho \sqrt{3e(\gamma )2\gamma \frac{de}{d\gamma }+\frac{1}{2}\gamma ^2\frac{d^2e}{d\gamma ^2}},$$
(3)
where $`e(\gamma )`$ is defined by Lieb and Liniger in liebliniger . The $`v_s/\rho `$ is depicted in Fig. 1 – notice how slowly the system enters a Tonks regime. A large BEC-Tonks crossover, $`\gamma (1,50)`$, provides us motivation for studies of shock waves outside BEC bodzio\_b ; gammal ; konotop ; menotti ; el and Tonks bodzio\_f limits.
Many-body solutions of (1) at zero absolute temperature have been analyzed in classic papers liebliniger ; lieb . These solutions are based on the Bethe ansatz and allow for analytical extraction of different static system properties. Unfortunately, they are too complicated for description of system dynamics.
To simplify the problem a proper hydrodynamical approach can be worked out olshani2 ; kohn and leads to the following set of equations
$$\frac{\rho }{t}+\frac{}{x}\left(v\rho \right)=0,$$
(4)
$$\frac{v}{t}+\frac{}{x}\left(\frac{1}{2}v^2\right)+\frac{}{x}\left(\mu (\rho )+V_l\frac{1}{2}\frac{_x^2\sqrt{\rho }}{\sqrt{\rho }}\right)=0$$
(5)
where
$$\mu (\rho )=\frac{1}{2}\rho ^2\left(3e(\gamma )\gamma \frac{de}{d\gamma }\right),$$
(6)
and $`V_l`$ is an external potential acting on atoms. A nice property of these equations is that they exactly reproduce sound velocity (3), while the problems with their usage for description of shock propagation (not formation) come from their derivation valid in the long-wavelength limit, where the quantum pressure (QP) term,
$$\frac{1}{2}\frac{_x^2\sqrt{\rho }}{\sqrt{\rho }},$$
(7)
is negligible.
To get more insight into physics described by (5) and validity of a QP term we look at BEC and Tonks limits. In the BEC case, $`\mu (\rho )a\rho `$ and Eqs. (4,5) can be obtained by time-dependent variational principle applied to a product wave function
$$\mathrm{\Psi }(x_1,\mathrm{},x_N,t)=\sqrt{N}\varphi (x_1,t)\mathrm{}\varphi (x_N,t),$$
(8)
$`dx|\varphi (x,t)|^2=1`$. It results in time-dependent Gross-Pitaevskii equation (29). Then substitution of $`\varphi =\sqrt{\rho }\mathrm{exp}(i\chi )`$ and $`v=_x\chi `$ into (29) gives a BEC version of (4,5). In other terms, $`\rho `$ in (4,5) is a single particle density defined in a many-body theory as
$$\rho (x,t)=dx_2\mathrm{}dx_N|\mathrm{\Psi }(x,x_2,\mathrm{},x_N,t)|^2,$$
(9)
which equals $`N|\varphi (x,t)|^2`$ in the BEC limit. A quantum pressure term in this limit is rigorously derived dalfovo .
In the Tonks regime one has $`\mu (\rho )\rho ^2\pi ^2/2`$, which was first found by renormalization group approach in kolomeisky ; kolomeisky1 , and then used in a number of papers, e.g., bodzio\_f ; proukakis ; tosi . Now derivation of (5), do not involve any sort of product simplification (8) of a wave-function. Indeed, a Fermi-Bose mapping theorem bfmt implies, e.g., that if $`\mathrm{\Psi }_F(x_1,\mathrm{},x_N)`$ is a ground state wave-function of noninteracting fermions placed in the same potential as a Tonks gas, then a ground state Tonks wave-function is $`|\mathrm{\Psi }_F(x_1,\mathrm{},x_N)|`$. Therefore, even a simple description of a Tonks gas based on the single particle density (9) involves $`N`$ orthogonal single particle orbitals instead of a single one, $`\varphi (x)`$, used in the BEC case.
In Tonks limit the QP term was shown to lead to unphysical density oscillations inter ; bodzio\_f . From knowledge that the quantum pressure term is present in the BEC limit $`\gamma 0`$ and absent in Tonks regime $`\gamma +\mathrm{}`$ it is clear that Eqs. (4,5) can not describe shock propagation for arbitrary $`\gamma `$. Nonetheless, shock formation from density perturbations that are initially wide can be successfully done, and will be described below. In this case the QP term is unimportant roughly up to shock formation. Due to lack of theoretical concepts for getting exact time-dependent solutions for the system of interest, we consider future experimental results as the best verification of our calculations based on hydrodynamic equations.
## III Shock wave solutions
The quantum pressure term (7) is important only when density changes occur on length scales smaller then the characteristic length given by
$$\xi (\rho )=\frac{1}{\sqrt{2\mu (\rho )}}.$$
(10)
In a BEC limit $`\xi (\rho )`$ is called a healing length, and we propose to use the same name regardless of $`\gamma `$.
Now we want to solve Eqs. (4,5,6). Assuming that the perturbation under consideration is, at least initially, broad compared to $`\xi `$, the quantum pressure term in (5) can be neglected. Following standard methods landau one gets
$$\rho (x,t)=f\left(x\left(𝑑\rho \sqrt{\frac{1}{\rho }\frac{\mu }{\rho }}+\sqrt{\rho \frac{\mu }{\rho }}\right)t\right),$$
(11)
$$v(x,t)=𝑑\rho \sqrt{\frac{1}{\rho }\frac{\mu }{\rho }}.$$
(12)
To proceed further with analytical calculations we use the following relations:
$$0<\gamma <\gamma _c:\sqrt{\frac{1}{\rho }\frac{\mu }{\rho }}\sqrt{\gamma }\frac{\gamma }{4\pi }\frac{\gamma ^{3/2}}{32\pi ^2}\frac{\gamma ^2}{128\pi ^3},$$
(13)
$$\gamma _c\gamma :\sqrt{\frac{1}{\rho }\frac{\mu }{\rho }}\frac{\pi }{(1+2/\gamma )^2},$$
(14)
where $`\gamma _c14.5`$ from the requirement that both expansions match at $`\gamma _c`$. Expression (13) is extracted from Lieb’s observation lieb that for $`\gamma <10`$
$$\sqrt{\frac{1}{\rho }_\rho \mu }\sqrt{\gamma \gamma ^{3/2}/(2\pi )}.$$
(15)
Since (15) complicates further calculations, we expanded it into a series around $`\sqrt{\gamma }/(2\pi )`$ equal to zero. Expression (14) is taken directly from Lieb and Liniger paper liebliniger – see Eqs. (3.32) and Appendix A.
To test above approximations we compare the shock wave solutions that neglect the QP term and use approximate expression for $`\sqrt{\frac{1}{\rho }\frac{\mu }{\rho }}`$, to the full numerical solution of hydrodynamic equations (4,5,6). Initially we determine velocity field from (12) and choose
$$\rho (x,0)=\rho _0+\frac{\eta \rho _0}{\mathrm{cosh}(x/\sigma )^2},\eta >0,$$
(16)
with $`\sigma \xi (\rho )`$ (calculations for $`\eta <0`$ can be easily repeated). The background density $`\rho _0`$ is found from the normalization condition: $`_l^l𝑑x\rho (x,0)=N`$, where we have assumed that periodic box has boundaries at $`\pm l`$. For well localized perturbations being of interest from now on one gets $`\rho _0=\frac{N}{2(l+\eta \sigma )}`$. Finally, it is convenient to define relative density
$$\varrho =\frac{\rho }{\rho _0}.$$
Now we consider separately $`\gamma >\gamma _c`$ and $`\gamma <\gamma _c`$ cases.
### III.1 $`\gamma \gamma _c`$ case
The healing length is found from (10) and Eqs. (3.32) of liebliniger subjected to rescalings of Appendix A
$$\xi (\rho )=\frac{1}{\pi \rho }\left(1+\frac{8}{3\gamma }+\frac{2}{3\gamma ^2}\right),$$
(17)
where the first term corresponds to the Fermi length of a non-interacting Fermi gas – a result that may be expected from the FBMT bfmt .
Rewriting implicit solution (11) to the form
$$\varrho (x,t)=f\left(xc_{\mathrm{}}w_{\mathrm{}}(\varrho )t\right),$$
where $`c_{\mathrm{}}=\pi \rho _0`$ is a background sound velocity of a Tonks gas, one gets from (11) and (14)
$$w_{\mathrm{}}(\varrho )=\frac{2\varrho +s\varrho ^21}{(1+s\varrho )^2(1+s)},s=\frac{2}{\gamma _0}.$$
Taking (16) as an initial density profile, the implicit shock-wave solution reads:
$$\varrho (x,t)=1+\frac{\eta }{\mathrm{cosh}[(xc_{\mathrm{}}w_{\mathrm{}}(\varrho )t)/\sigma ]^2}.$$
(18)
Although, the explicit form of $`\varrho (x,t)`$ can not be extracted analytically, important properties of pulse dynamics can be determined. First of all, impulse amplitude, $`1+\eta `$, is constant in time. Second, speed of impulse maximum for any $`\eta >0`$ equals
$$𝒱(\eta )=c_{\mathrm{}}w_{\mathrm{}}(1+\eta ),$$
and applicability of this expression is not limited to $`1/\mathrm{cosh}(x/\sigma )^2`$ perturbations only – Sec. IV. Third, the width of the impulse at given density is constant during propagation- a property that we missed in earlier papers bodzio\_b ; bodzio\_f . Fourth, impulse tails propagate roughly with the background sound velocity equal to $`c_{\mathrm{}}w_{\mathrm{}}(1)`$, while the impulse maximum moves with the speed $`c_{\mathrm{}}w_{\mathrm{}}(1+\eta )`$. Since $`w(\varrho )`$ monotonically increases with $`\varrho `$, the impulse deforms its shape so that a shock wave front forms, i. e., $`|_x\rho (x,t)|=+\mathrm{}`$ at one point.
Time and position of shock-wave creation can be extracted from equations landau :
$$_\rho x(\rho )=0,_\rho ^2x(\rho )=0.$$
(19)
Their solution gives: density $`\rho _s`$ at which density profile becomes locally vertical, and time $`t_s`$ at which this occurs:
$$\varrho _s=\frac{3+4s\eta +9s\sqrt{9(1+s)^2+16s^2\eta ^2}}{6s},$$
(20)
$$t_s=\frac{\sigma \sqrt{\eta }}{4c_{\mathrm{}}}\frac{(1+s\varrho _s)^3}{(\varrho _s1)\sqrt{1+\eta \varrho _s}}.$$
(21)
These expressions close to a Tonks gas limit take a simple form:
$$\rho _s=1+\frac{2}{3}\eta +𝒪\left(\frac{1}{\gamma }\right),t_s=\frac{3\sqrt{3}}{8}\frac{\sigma }{\eta c_{\mathrm{}}}+𝒪\left(\frac{1}{\gamma }\right).$$
Comparison of analytical solution (18) and full numerical one based on hydrodynamic equations (4,5,6) is presented in Fig. 2. As easily seen, there is a good agreement between hydrodynamical solution and a shock-wave one until the moment of shock creation. Then discrepancy increases due to appearance of density oscillations triggered by the QP term (7) neglected in derivation of (18). Since presence of the QP term for large $`\gamma `$, e.g., $`\gamma \gamma _c`$, is questionable, it is an interesting open question whether density oscillations in the form presented in Fig. 2 survive or not for any $`\gamma \gamma _c`$ considered here.
### III.2 $`\gamma <\gamma _c`$ case
The healing length reads
$$\xi (\rho )=\frac{1}{\sqrt{2a\rho }}\left(1+\frac{\sqrt{\gamma }}{2\pi }+\frac{3\gamma }{8\pi ^2}\right),$$
(22)
derived by extraction of $`\mu (\rho )`$ from (15) and subsequent Taylor expansion of (10). Naturally, the first term in (22) corresponds to the BEC healing length.
The implicit shock-wave solution has the form:
$$\varrho (x,t)=f\left(xc_0w_0(\varrho )t\right),$$
(23)
where $`c_0=\sqrt{a\rho _0}`$ is the background speed of sound in the limit of $`\gamma 0`$. Combining (11) and (13) one gets
$`w_0(\varrho )=`$ $`3\sqrt{\varrho }2{\displaystyle \frac{\sqrt{\gamma _0}}{4\pi }}\left(\mathrm{ln}\varrho +1\right)+{\displaystyle \frac{\gamma _0}{32\pi ^2}}\left(\sqrt{{\displaystyle \frac{1}{\varrho }}}2\right)`$ (24)
$`{\displaystyle \frac{\gamma _0^{3/2}}{128\pi ^3}}.`$
To verify accuracy of (23) we have simulated dynamics for the system with $`\gamma _0=a/\rho _0=5`$, i. e., for $`\gamma _0`$ outside the BEC mean-field regime (see inset of Fig. 1), but small enough to stay clearly within $`\gamma <\gamma _c`$ parameter range. Agreement between full numerical solution of Eqs. (4,5) and implicit shock solution is good before shock formation. Close to time of shock creation discrepancies show up in the oscillatory region missed in (23). Qualitatively, the plot that presents these results is the same as Fig. 2.
Since dynamics of initial density profile is qualitatively the same as in the $`\gamma \gamma _c`$ case, we notice that impulse maximum moves now for any $`\eta >0`$ with velocity
$$𝒱(\eta )=c_0w_0(1+\eta ).$$
The explicit solution of shock equations (19) can be found in the limit of $`\gamma _00`$, but still arbitrary $`\eta >0`$:
$$\rho _s=\frac{1+\eta +\sqrt{(1+\eta )(9+\eta )}}{4},$$
$$t_s=\frac{\sigma \sqrt{\eta \rho _s}}{3c_0(\rho _s1)\sqrt{1+\eta \rho _s}}.$$
Interestingly, for a gaussian like initial impulse discussed in bodzio\_b explicit expressions for $`\rho _s`$ and $`t_s`$ in this limit were beyond the reach. Finally, we note that the QP term for $`\gamma 0`$ is certainly present, but it is hard to say whether it survives in (5) for any $`\gamma <\gamma _c`$. Since the QP term affects shock dynamics strongly, an experiment should clarify this uncertainty.
## IV Propagation of pulses of arbitrary shape
In this section we describe some general properties of shock-wave solutions (11,12). To this aim we consider density profiles satisfying
$$\rho (x,t)=f(xW(\rho )t)=f(\zeta ).$$
(25)
Suppose that there are $`n`$ extrema (minima, maxima), placed on a background density $`\rho _0`$, in the initial density profile $`\rho (x,0)`$. Let us denote positions of these extrema at $`t=0`$ as $`x_i(0)`$ and set $`\rho _i=\rho (x_i(0),0)`$. It means that $`_xf(x)|_{x_i(0)}=0`$ and $`_x^2f(x)|_{x_i(0)}0`$. At $`t>0`$ one gets from (25)
$$_x\rho (x,t)=\frac{\frac{f(\zeta )}{\zeta }}{1+_\rho W(\rho )\frac{f(\zeta )}{\zeta }t},$$
which obviously equals zero iff $`\zeta =x_i(0)`$. It implies that $`x_i(t)=x_i(0)+W(\rho _i)t`$. Furthermore, we have $`_x^2\rho (x,t)|_{x_i(t)}=_x^2\rho (x,0)|_{x_i(0)}0`$, which proves that extrema do not change into saddle points in the course of time evolution. This means that velocity of impulse extrema equals $`W(\rho _i)`$, where $`\rho _i=\rho (x_i(t),t)=\rho (x_i(0),0)`$.
Since $`W(\rho )`$ is shape independent, we conclude that for any $`\gamma `$ density extrema propagate with constant amplitude and speed equal to $`W(\rho _i)`$. This speed is well approximated as
$$0<\gamma <\gamma _c:c_0w_0(\rho _i/\rho _0),$$
(26)
$$\gamma >\gamma _c:c_{\mathrm{}}w_{\mathrm{}}(\rho _i/\rho _0),$$
(27)
for $`\gamma `$’s a little off $`\gamma _c`$.
Finally, it is important to stress that above statements concerning the amplitude and velocity of density extrema are valid at least as long as the implicit solution works, i. e., up to the time of shock creation. After that time, they are valid in regions far enough from shock structures.
## V Experimental realization
As discussed in bodzio\_b ; bodzio\_f , experimental creation of matter-wave packets undergoing shock-wave dynamics is straightforward. The idea is to cool atoms in an additional external potential created by a well-detuned laser beam, and then suddenly turn the beam off. The term suddenly, means on time scale much smaller than time of sound propagation through the perturbation.
The initial density of atoms resulting from an external laser potential $`V_l`$ is determined as follows. Combining density $`\rho `$ and velocity field $`v=_x\chi `$ into $`\varphi =\sqrt{\rho }\mathrm{exp}(i\chi )`$ the Eqs. (4,5) can be rewritten to the standard form
$$i_t\varphi =\frac{1}{2}_x^2\varphi +V_l(x)\varphi +\mu (\rho )\varphi ,$$
(28)
which time-independent version reads
$$\frac{1}{2}_x^2\varphi +V_l(x)\varphi +\mu (\rho )\varphi =\stackrel{~}{\mu }\varphi ,$$
with $`\stackrel{~}{\mu }`$ being a chemical potential. Assuming that the laser potential induces density changes on length scales larger than the healing length (10), one determines an initial density profile from
$$V_l(x)+\mu (\rho )=\stackrel{~}{\mu },$$
with $`\stackrel{~}{\mu }`$ found from $`𝑑x\rho =N`$. Naturally, velocity field equals zero in an initial state considered here.
A straightforward generalization of results from bodzio\_b ; bodzio\_f shows that if the initial density profile is $`\rho (x,0)=\rho _0+h(x)`$ with $`\mathrm{max}|h(x)|\rho _0`$ then $`\rho (x,t)=\rho _0+h\left(xv_s(\rho _0)t\right)/2+h\left(x+v_s(\rho _0)t\right)/2`$, with $`v_s(\rho _0)`$ being background sound velocity (3). Therefore, there are left and right moving perturbations. Due to mirror symmetry, we can forget about the left moving one once the pulses are well separated. The right-moving pulse is described by $`\rho _0+h(xv_s(\rho _0)t)/2`$ as long as $`tt_s`$, since even arbitrarily small density perturbations experience a shock deformation after long enough evolution.
In BEC and Tonks limits qualitatively the same splitting process takes place even when $`\mathrm{max}|h(x)|/\rho _01`$ bodzio\_b ; bodzio\_f . When the initial density profile is $`\rho (x,0)=\rho _0+h(x)`$, the right moving pulse is described by $`\rho (x,t)\rho _0+f(xW(\rho )t)`$ with $`W(\rho )`$ being the same as in (25) and $`f(x)=h(x)/2`$. We checked by solving hydrodynamic equations (4,5) that above described splitting process works qualitatively the same way for any $`\gamma `$, as it works for $`\gamma 0,+\mathrm{}`$ bodzio\_b ; bodzio\_f . This observation explains how we imagine practical generation of matter-wave pulses described in previous sections.
## VI Applicability of the Gross-Pitaevskii mean-field approach to shock problems
In the BEC limit system dynamics in the mean-field approximation (8) is described by the following version of Eq. (28)
$$i_t\varphi =\frac{1}{2}_x^2\varphi +V_l(x)\varphi +aN|\varphi |^2\varphi H_{GP}\varphi ,$$
(29)
where $`dx|\varphi (x,t)|^2=1`$. This is the Gross-Pitaevskii (GP) equation, which can be rigorously derived by different means – see, e.g., dalfovo ; castin .
To go beyond the mean-field approximation we employ a second quantization formalism, i. e., we transform Hamiltonian (1) with an additional potential term, $`V_l(x)`$, into the form
$$\widehat{H}=𝑑x\widehat{\mathrm{\Psi }}^{}\left(\frac{1}{2}\frac{^2}{x^2}+V_l(x)+\frac{a}{2}\widehat{\mathrm{\Psi }}^{}\widehat{\mathrm{\Psi }}\right)\widehat{\mathrm{\Psi }},$$
(30)
where $`\widehat{\mathrm{\Psi }}`$ is a field operator. In the number conserving Bogolubov approach castin it reads
$$\widehat{\mathrm{\Psi }}=\widehat{a}_0\varphi +\delta \widehat{\mathrm{\Psi }},$$
with $`\widehat{a}_0`$ annihilating a particle from a condensate, and $`\delta \widehat{\mathrm{\Psi }}`$ annihilating a particle from modes orthogonal to the condensate one. Naturally, $`\varphi `$ is a condensate mode, i. e., it is an eigenstate of a single particle density matrix, $`\widehat{\mathrm{\Psi }}^{}(x^{})\widehat{\mathrm{\Psi }}(x)`$, to the highest eigenvalue ($`\lambda _0N`$):
$$dx^{}\widehat{\mathrm{\Psi }}^{}(x^{})\widehat{\mathrm{\Psi }}(x)\varphi (x^{})=\lambda _0\varphi (x).$$
The single particle density of atoms, $`\widehat{\mathrm{\Psi }}^{}(x)\widehat{\mathrm{\Psi }}(x)`$, equals
$$\widehat{a}_0^{}\widehat{a}_0|\varphi (x)|^2+\underset{k}{}|v_k(x)|^2,$$
(31)
where the first term accounts for condensate density while the second one corresponds to density of depleted atoms. In further calculations we need both $`u_k`$ and $`v_k`$, e. g., for finite temperature calculations. Indeed, density of atoms at sufficiently low temperature $`T`$ ($`\lambda _0(T)/N1`$) equals castin
$$\widehat{a}_0^{}\widehat{a}_0|\varphi (x)|^2+\underset{k}{}\widehat{b}_k^{}\widehat{b}_ku_k|u_k+\underset{k}{}\widehat{b}_k^{}\widehat{b}_k+1v_k|v_k,$$
(32)
where $`\widehat{b}_k^{}\widehat{b}_k=[\mathrm{exp}(\omega _k/(k_BT)1]^1`$. To get $`u_k`$, $`v_k`$ and $`\omega _k`$ one first finds a condensate mode from the time-independent GP equation, $`H_{GP}\varphi =\stackrel{~}{\mu }\varphi `$, and then constructs the matrix
$$=\left(\begin{array}{cc}H_{GP}+aNQ|\varphi |^2\stackrel{~}{\mu }& aNQ\varphi ^2\\ aNQ^{}\varphi ^2& H_{GP}aNQ^{}|\varphi |^2+\stackrel{~}{\mu }\end{array}\right),$$
where $`Q`$ is a projector to space orthogonal to a condensate mode: $`Q\psi =\psi \varphi \varphi |\psi `$. The modes and frequencies are found from the eigen equation
$$\left(\begin{array}{c}u_k\\ v_k\end{array}\right)=\omega _m\left(\begin{array}{c}u_k\\ v_k\end{array}\right),$$
while their dynamics is captured by
$$i\frac{}{t}\left(\begin{array}{c}u_k\\ v_k\end{array}\right)=\left(\begin{array}{c}u_k\\ v_k\end{array}\right),$$
where $`\stackrel{~}{\mu }`$ is set to zero in $``$, and the calculation of operator $``$ at every time step requires simultaneous solution of the time-dependent Gross-Pitaevskii equation. Expressions for densities of atoms (31) and (32) are unchanged in a time dependent case. For a list of $`u`$, $`v`$ modes properties, and details of the number-conserving Bogolubov approach see castin .
It is important to realize what are limitations of the Bogolubov approach. To simplify notation we discuss $`T=0`$ case now (the results for $`T0`$ can be easily obtained). The Bogolubov method fails in the following two situations. First, the global breakdown, shows up when total number of depleted atoms, i. e., $`𝑑x_k|v_k(x)|^2`$, becomes comparable to number of atoms in the system. This limitation comes directly from derivation of Bogolubov Hamiltonian castin . Second, local breakdown, happens when density of depleted atoms becomes comparable to condensate density on a length scale of the order of the healing length:
$$_{xf\xi }^{x+f\xi }𝑑x^{}N|\varphi (x^{})|^2_{xf\xi }^{x+f\xi }𝑑x^{}\underset{k}{}|v_k(x^{})|^2,$$
(33)
with $`f=𝒪(1)`$. To clarify (33), we notice that in the Bogolubov approach an impact of noncondensed atoms on condensed ones is neglected. Indeed, to get Bogolubov dynamics one solves a time-dependent GP equation and then puts this solution into equations of motion of noncondensed atoms. In this way noncondensed particles are influenced by condensed ones but not other way round. Next order corrections take into account change of condensate density as a result of repulsive interactions between condensed and noncondensed clouds, and lead to modifications of the GP equation. The correction to the right hand side of (29) is qualitatively of the form $`V_{\mathrm{depl}}\varphi `$ (Eqs. (95,96) in castin ), where
$$V_{\mathrm{depl}}(x)2a\underset{k}{}|v_k(x)|^2,$$
(34)
so that $`V_{\mathrm{depl}}(x)`$ can be regarded as an external potential acting on a condensate. Extension of changes of condensate density induced by an external, localized potential, is given by the healing length dalfovo . Therefore, impact of depleted atoms on a condensate is best given by comparison of the new term, $`V_{\mathrm{depl}}`$ averaged over a few healing lengths, to previously dominating term in the GP equation, $`aN|\varphi (x)|^2`$, averaged on the same set of points. It gives the condition (33). Appearance of such a local breakdown of Bogolubov approach will be shown after breaking of a dark shock density profile resulting in soliton production.
Finally, let us comment when condensate depletion might affect mean-field predictions. Since density of depleted atoms depends on the product of $`a`$ and $`N`$ only, it is unaffected by the limit
$$N+\mathrm{},a0,aN=\mathrm{fixed},$$
while the condensed part in total atomic density grows as $`\widehat{a}_0^{}\widehat{a}_0N`$ (31,32). It means that in this limit, being to some extent possible with the help of Feshbach resonances cornish , density of depleted atoms does not affect density measurements. The exception might happen when there is dynamical instability leading to fast increase of depleted fraction. Then sooner or later density of depleted atoms will start to have an important impact on condensate dynamics.
In a generic experimental situation, however, a fraction of depleted atoms is non-negligible. To compare a mean-field prediction to quantum many-body dynamics, we have chosen a particular set of experimentally accessible parameters – see Appendix B for details.
In the following we discuss evolution of shock structures resulting from white ($`\eta >0`$) and dark ($`\eta <0`$) initial density profiles, where $`\eta `$ stands for relative perturbation amplitude: as in Eq. (16). Since corrections to mean field equations are very different in these two cases we discuss them separately.
White shock waves: We have done different calculations changing not only the size of an initial perturbation but also $`aN`$, $`N`$, etc. In all cases we observe that total depletion of a condensate negligibly increases from beginning of time evolution to appearance of well developed shock structures. Density of depleted atoms around shocks becomes at most $`𝒪(1)`$ larger than depletion density far away from perturbations. Therefore, if density and distribution of depleted atoms was such that the Bogolubov approach was initially applicable, it will also work during shock creation and propagation.
The results for experimentally relevant choice of parameters are depicted in Fig. 3. As easily seen from differences in scales between Figs. 3(a), (c) and Figs. 3(b), (d), the corrections to the mean-field result are minor. Repulsion between condensed and depleted atoms, results in localization of depleted atoms around condensate density minima: Fig. 3(e). This was predicted by us in bodzio\_b , however, the amount of depleted atoms, turns out to be insufficient to fill condensate ripples.
We have also done finite temperature calculations according to (32). At $`T0`$ not only quantum depletion but also thermal one shows up, which leads to increase of total number of depleted atoms in comparison to the $`T=0`$ situation. We have considered low temperatures, i. e., $`T30`$nK (see Appendix B), to see whether thermal effects qualitatively affect shock dynamics. The only difference with respect to $`T=0`$ case is that total depletion is larger by a factor of $`𝒪(10)`$ for $`T30`$nK, which is still not enough for getting significant corrections to BEC shocks from depleted atoms.
These calculations suggest that the Gross-Pitaevskii equation captures correctly physics of white shock waves in a BEC. Finite temperature effects does not seem to destroy the qualitative picture of shock formation and propagation given by the mean-field approach.
Dark shock waves: Dark density profiles undergo two distinct stages of evolution. In the first one, shock waves form, while in the second one shock density profile breaks into a train of solitons that move, according to the mean-field approach, with constant velocity and without changes of shape. The same was independently observed numerically in budde and theoretically in umarov . Theoretical predictions in both papers are based solely on the basis of mean-field equations.
As in white shocks calculations, we have fixed the parameters to those experimentally relevant– see Appendix B. The formation of shock wave structures is depicted in Figs. 4(a), (b), (c), (d). Now a rear impulse edge self-steepens instead of a front one. This comes from the fact that now impulse tail moves faster than impulse center, i. e., density minimum. In the case considered in Fig. 4, a shock profile breaks before complete separation. Obviously, if the impulse would be more shallow the breakdown would happen later. As easily seen, the correction to total density coming from depleted atoms can be neglected at this stage of time evolution.
After shock breakdown, solitons are produced as depicted in Fig. 4(e). Then density of depleted atoms increases by orders of magnitude. It results in significant corrections to total density of atoms and causes local, i. e., around soliton minima, breakdown of the Bogolubov approach according to (33). Notice that still fraction of depleted atoms, less than $`0.7\%`$ in Fig. 4(f), is so small that the global condition of Bogolubov approach applicability is well satisfied. As depicted in Fig. 5 both peak density of depleted atoms at soliton minimum (dots) and total depletion of a condensate (squares) follow approximately power law increase from the moment of soliton train formation.
It is easy to predict qualitatively what are corrections to a Bogolubov result coming from depleted atoms. The next order corrections to the Bogolubov approach come from inclusion of repulsive interactions between condensed and non-condensed atoms. It is qualitatively accounted for by introducing a new potential term, Eq. (34), into the Gross-Pitaevskii equation. It means that soliton structures in condensate density will become wider, so that depleted atoms will have more space to distribute themselves inside solitons. Therefore, increase of peak density of depleted atoms will certainly slow down. This observation is supported by nonperturbative calculation of Dziarmaga jacek , who found that total density of atoms at soliton minimum approaches background density for large times.
Such a fast increase in density of depleted atoms around soliton minima is rather a generic phenomenon. Indeed, it was first theoretically discussed for phase-imprinted solitons in harmonic traps kark ; solitony as a mechanism for fast disappearance of solitons in the Hannover experiment hannover . The experimental time scale for soliton disappearance was found both theoretically solitony ; jacek and experimentally hannover to be $`𝒪(10\mathrm{m}\mathrm{s})`$. Our calculation, based on the same value of nonlinear parameter $`aN=7500`$, gives the same estimation for instant when corrections to dynamics of dark-shock-originated solitons become important. Indeed, the time instant $`0.75`$ in Fig. 4(e), corresponds to $`12`$ms– see Appendix B for a unit of time.
To end this section we recall that predictions presented here apply to a 1D BEC system. There was recently an attempt to experimentally verify applicability of the mean-field approach to shock phenomena simula . The conclusion was that the ripples in a white shock wave front are not filled with depleted atoms. Our present work supports this observation assuming that 1D results are qualitatively correct in 2D rapidly rotating array of vortices on top of which a white shock wave front propagates.
## VII Summary
This paper presents a complete hydrodynamical description of basic properties of shock-waves in a delta interacting 1D Bose gas. Our predictions can be directly verified experimentally. In principle, a proper experimental setup should answer at what Lieb-Liniger parameter $`\gamma `$ the quantum pressure term starts to lead to unphysical results. Notice that density oscillations in front of a shock wave are present due to this term for small enough $`\gamma `$, while for large $`\gamma `$ they are absent in exact solution and shocks propagate in a very different way bodzio\_f . In fact, a change of this kind in shock dynamics can be a nice experimental signature of a BEC-Tonks crossover.
We also discussed quantum many-body corrections to Gross-Pitaevskii shock-wave solutions. This way we clarified the role of depleted atoms in shock dynamics. This important point was missed in previous studies of BEC shock waves bodzio\_b ; gammal ; konotop ; menotti ; el .
From experimental point of view, tools for verification of our predictions seem to be available either right now or in the nearest future. Indeed, a first experiment on sound propagation in a BEC was done years ago ketterle . Since we do not discuss time-of-flight measurements the best comparison between our theory and experiment should rely on in situ measurements ketterle . The most prospective experimental setup for observation of shock dynamics in different $`\gamma `$ parameter ranges is probably provided by a box-like optical trap raizen , where values of $`\gamma 1`$ approaching the BEC-Tonks crossover (Fig. 1), have been already achieved. Other interesting experimental setups include atom chips jorg and circular atom waveguides gupta . The latter one, being best suited for studies of shock collisions.
From theoretical side, there are at least two possible interesting extensions of this work. First of all, one can analyze shock dynamics in quasi-1D trapping geometries, where the system is a 3D waveguide with adjustable harmonic transverse confinement. This can be easily done using results from parola . Indeed, instead of Eq. (6) one can consider Eq. (5) of parola and repeat subsequent calculations. In the limit of tight transversal confinement, i. e., for a 1D Bose gas, both expressions for $`\mu (\rho )`$ lead to physically identical results. The second interesting extension of this work includes studies of possible outcomes of a single shot density measurement. In fact, it should be stressed that predictions based on hydrodynamical and Bogolubov approaches apply to averages over different experimental measurements done on the system prepared many times in the same quantum state. Needless to say, averages may differ from a single shot outcomes – see juha for an illustrative example.
I would like to acknowledge discussions with Zbyszek Karkuszewski. This work was started in the Institute for Theoretical Physics in Hannover. I’m grateful to both the Alexander von Humboldt Foundation for support of this work in Germany, and to the US Department of Energy for support of this research at the LANL.
## Appendix A Units in a box
Since this paper uses extensively results of liebliniger ; lieb it is useful to link dimensionless quantities used by us to those of liebliniger ; lieb .
The eigen equation of a 1D Hamiltonian expressed in terms of dimensional quantities denoted by primes is
$$\frac{\mathrm{}^2}{2m}\underset{i}{}\frac{^2}{x_{i}^{}{}_{}{}^{2}}\mathrm{\Psi }+A^{}\underset{i<j}{}\delta (x_i^{}x_j^{})\mathrm{\Psi }=E^{}\mathrm{\Psi },$$
(35)
to get dimensionless Hamiltonian (1) one introduces: $`x_i^{}=x_il_0`$, $`E^{}=E\mathrm{}^2/(ml_0^2)`$, $`A^{}=a\mathrm{}^2/(ml_0)`$, where $`l_0`$ is an arbitrary length scale. Consideration of a time-dependent Schrödinger equation leads to $`t^{}=tml_0^2/\mathrm{}`$.
The dimensionless Hamiltonian used by Lieb and Liniger has the form
$$\widehat{H}=\underset{i}{}\frac{^2}{x_i^2}+2c\underset{i<j}{}\delta (x_ix_j),$$
and can be obtained from (35) after rescalings: $`x_i^{}=x_il_0`$, $`E^{}=E\mathrm{}^2/(2ml_0^2)`$, $`A^{}=c\mathrm{}^2/(ml_0)`$, $`t^{}=t\mathrm{\hspace{0.17em}2}ml_0^2/\mathrm{}`$.
Therefore, there are the following relations between our dimensionless quantities and Lieb and Liniger ones marked by LL:
$$\gamma =\gamma _{LL},\mu =\frac{\mu _{LL}}{2},v=\frac{v_{LL}}{2},$$
where $`v`$ is sound velocity. Finally we note, that there is a misprint in the first line of expression (1.4) providing sound velocity lieb . There should be $`2\gamma de/d\gamma `$ (as used in Sec. II) instead of $`\gamma de/d\gamma `$. The second line of (1.4) in lieb is correct.
## Appendix B A box approximation of a 3D system in the BEC limit
We aim at getting a 1D box approximation of a full 3D BEC gas confined in a 3D harmonic potential. The system in the mean-field approximation satisfies the following 3D Gross-Pitaevskii equation
$`i\mathrm{}{\displaystyle \frac{\mathrm{\Psi }^{}}{t^{}}}=`$ $``$ $`{\displaystyle \frac{\mathrm{}^2}{2m}}\stackrel{}{}^2\mathrm{\Psi }^{}+(m\omega ^2x^2/2+m\omega _{}^2r^2/2)\mathrm{\Psi }^{}`$
$`+`$ $`g_{3D}^{}|\mathrm{\Psi }^{}|^2\mathrm{\Psi }^{},`$
where $`\omega _{}\omega `$, $`g_{3D}^{}=4\pi \mathrm{}^2a_{sc}N/m`$ with $`a_{sc}`$ being interatomic scattering length, and $`d^3x^{}|\mathrm{\Psi }^{}|^2=1`$. First, we rescale all quantities using the harmonic oscillator units in the $`x^{}`$ direction. Denoting by symbols without primes dimensionless quantities,
$$(x^{},y^{},z^{})=(x,y,z)\sqrt{\frac{\mathrm{}}{m\omega }},\mathrm{\Psi }^{}=\mathrm{\Psi }\left(\frac{m\omega }{\mathrm{}}\right)^{3/4},t^{}=\frac{t}{\omega },$$
we arrive at the following 3D GP equation
$$i\frac{\mathrm{\Psi }}{t}=\frac{1}{2}\stackrel{}{}^2\mathrm{\Psi }+(x^2/2+\lambda ^2r^2/2)\mathrm{\Psi }+g_{3D}|\mathrm{\Psi }|^2\mathrm{\Psi },$$
where $`\lambda `$ is trap aspect ratio $`\omega _{}/\omega `$ while $`g_{3D}=4\pi a_{sc}N\sqrt{m\omega /\mathrm{}}`$. This equation can be derived from the following energy functional:
$$[\mathrm{\Psi }]=d^3x\left[\frac{1}{2}|\stackrel{}{}\mathrm{\Psi }|^2+(x^2/2+\lambda ^2r^2/2)|\mathrm{\Psi }|^2+\frac{g_{3D}}{2}|\mathrm{\Psi }|^4\right].$$
(36)
To proceed further we assume a simple form of a variational wave-function,
$$\mathrm{\Psi }(x,r)=\varphi (x)\mathrm{exp}\left(\frac{r^2}{2l_{}^2}\right)\frac{1}{\sqrt{\pi }l_{}},$$
(37)
where $`dx|\varphi (x)|^2=1`$ and $`l_{}`$ is a variational parameter. Substituting (37) into (36) one gets
$$[\varphi ]=𝑑x\left[\frac{1}{2}|_x\varphi |^2+\frac{x^2}{2}|\varphi |^2+\frac{g_{3D}}{4\pi l_{}^2}|\varphi |^4\right]+\frac{l_{}^2\lambda ^2}{2}+\frac{1}{2l_{}^2},$$
(38)
which supports the following 1D GP equation
$$i_t\varphi =\frac{1}{2}_x^2\varphi +\frac{x^2}{2}\varphi +aN|\varphi |^2\varphi ,$$
(39)
with $`aN=g_{3D}/(2\pi l_{}^2)`$. To determine $`l_{}`$, we substitute the Thomas-Fermi solution of (39),
$$aN|\varphi (x)|^2=\frac{1}{2}\left(\frac{3aN}{2}\right)^{2/3}\frac{x^2}{2},$$
(40)
into (38) getting $`(l_{})`$, which minimization leads to
$$(g_{3D}l_{})^{2/3}=(250\pi ^2/9)^{1/3}(\lambda ^2l_{}^41).$$
In the limit of $`g_{3D}0`$ one gets $`l_{}=1/\sqrt{\lambda }`$, which is a noninteracting cloud width. In a typical experiment, a cloud width is much larger than $`1/\sqrt{\lambda }`$ dalfovo , so
$$l_{}\left(\frac{9}{250\pi ^2}\right)^{1/10}\frac{g_{3D}^{1/5}}{\lambda ^{3/5}}.$$
For a cigar shaped cloud one can assume, e.g., $`\omega =2\pi 10`$Hz, $`\omega _{}=2\pi 569`$Hz. Taking also $`N=1.510^5`$, $`a_{sc}=5.210^9`$m (scattering length of <sup>87</sup>Rb in the $`|F=2,m_F=2`$ state), $`m=89.91\times 1.6610^{27}`$kg (atomic mass of <sup>87</sup>Rb), one can find that the units of length and time are $`3.4\mu m`$ and $`16`$ms, respectively. For these parameters $`aN7500`$ and the unit of temperature, $`\mathrm{}\omega /k_B`$, equals $`0.48`$nK. By using all these results it is easy to transform dimensionless plots from Sec. VI, into dimensional ones.
A box approximation of dimensionally reduced harmonically trapped cloud described above is the following. We place exactly the same number of atoms in the box as in the quasi-1D configuration described by (40). We assume that, in the absence of external laser potential, density in a box extending from $`[l,l]`$ is exactly the same as at a center of a harmonic trap: $`[3/(4\sqrt{2})]^{2/3}/(aN)^{1/3}=1/(2l)`$, which gives $`l=(aN)^{1/3}(2/3)^{2/3}.`$
For parameters defined above $`2l30`$ (about $`0.1`$mm), which can be compared to the Thomas-Fermi size of the harmonically trapped cloud equal here to $`44`$ (about $`0.15`$mm). The box-like approximation leads to results being in good qualitative agreement with calculations done in a harmonically trapped case. It concerns both length and time scales of shock dynamics.
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# QCD THEORY AT THE XL RENCONTRE DE MORIOND: FISH EYES AND PHYSICS
## 1 Introduction
The theory of the strong interactions is well understood in one sense: we are confident that the simple lagrangian of Quantum Chromodynamics (QCD) provides the correct description of the strong interactions. After all, it has passed many experimental tests that could have proven it wrong. Yet we face challenges in applying that theory. Predictions for bound state masses are difficult, as are predictions for decays of mesons containing a heavy quark, where we need the QCD part of the theory in order to use experimental results to get at the electroweak part. There appears to be a deep simplicity in the behavior of densely packed gluons, as probed at “small x,” but this behavior is not well understood. Interesting regularities in the development of the final state in heavy ion collisions have been observed, but have not been susceptible to a fully satisfactory interpretation in terms of quark-gluon interactions. Finally, we seek more powerful tools for using the theory in the context of high $`p_T`$ reactions at the upcoming Large Hadron Collider (LHC). Surely, we imagine, new physics signals will be found there, but understanding what those signals mean in terms of new particles or forces – maybe even new dimensions of space-time – will not be easy. We will want to use all of the theoretical methods we can muster.
These issues and more came up at the conference. In this talk, I review a selection of the talks that were either presented by theorists or were of special theoretical interest. I use the talks to provide some assessment of where we stand with respect to the problems and opportunities facing QCD theory. You have to read to the end to find out about the fish part.
## 2 b-quark production
In the bad old days, around Moriond XXX, the data on b-quark production indicated a problem, as illustrated on the left in Fig. 1. The results varied somewhat according to year, what was measured, and how the measurement was performed, but always indicated the observation of more $`b`$-quarks than theory predicted. One could hide behind the excuse that the $`P_T`$ values were never more than about 20 GeV and the mass of the $`b`$-quark itself is only 5 GeV, so that the effective $`\alpha _s`$ was not so small and the effectiveness of the perturbative QCD approach could be doubted. Nevertheless, one was left with the suspicion that the $`b`$ quark had some kind of anomalous behavior that lay beyond the Standard Model. In the last few years, the situation has been changing. At this conference, M. d’Onofrio presented the latest CDF and D0 data on this subject.$`^\mathrm{?}`$ I present one of the results on the right in Fig. 1. We see that the theory errors are still rather substantial, reflecting the difficulties mentioned above. The experimental errors are improved. And now data and theory agree.
There was some discussion at the conference of what had changed. First, we have better theory that sums logs of $`p_T/m_b`$. Second, we have better parton distributions and accompanying parton distributions. (See d’Onofrio’s talk.) Third, we measure a quantity whose definition does not involve too much theory, for example we measure cross sections for B-hadron production instead of b-quark production as was sometimes done in the past. A recent trend has been to measure the cross section for jets containing a b-quark. Examples of this were shown, although so far the NLO theory calculations are lacking and summing logarithms of $`p_T/m_b`$ will be needed to supplement a purely NLO calculation.
## 3 b-quark decays
There were several talks concerning the theory of b-quark decays. These talks concerned perturbative calculations of the effective lagrangian, lattice gauge theory calculations of decay matrix elements, and theoretical analyses of the decay spectrum in inclusive decays.
The issue of the effective lagrangian is illustrated in Fig. 2. Loop diagrams create an effective lagrangian of the form
$$=C_i𝒪_i$$
(1)
for the process $`bse^+e^{}`$ (for example). The $`𝒪_i`$ are operators such as $`\overline{s}(x)\gamma ^\mu b(x)\overline{e}(x)\gamma _\mu e(x)`$. The coefficients $`C_i`$ are to be calculated from the loop diagrams. My example diagram is too simple; what we need now are two loop diagrams. T. Huber reported his calculation of the $`C_i`$ for $`bse^+e^{}`$ from two loop diagrams in the Standard Model.$`^\mathrm{?}`$ S. Schilling described a calculation of the $`C_i`$ for $`bse^+e^{}`$ in the case of the two higgs doublet model.$`^\mathrm{?}`$
If we wish to calculate a completely inclusive decay rate, not much is needed beyond the effective lagrangian for the decay. But if we want an exclusive decay rate, we need a matrix element of the effective lagrangian between initial and final states, as illustrated in Fig. 3. M. Okamoto presented results for the matrix elements of appropriate weak decay operators between an initial heavy meson state and a final light meson state.$`^\mathrm{?}`$ The results are needed for a range of momentum transfers $`q`$ to the quark system. With appropriate limiting procedures, such matrix elements can be evaluated in lattice QCD. In the calculations described by Okamoto, enough results can be obtained to extract the complete CKM matrix from the corresponding experimental results.
I turn now to the decay spectrum in inclusive decays. Consider, in particular, the decay $`BX_s\gamma `$, where $`X_s`$ indicates any state that has an $`s`$-quark in it. Let $`x=2E_\gamma /M_b`$. From a theoretical point of view, by far the simplest cross section to calculate is the inclusive cross section $`𝑑x𝑑\mathrm{\Gamma }/𝑑x`$. However, the experimental acceptance typically does not include all $`x`$. For this reason, we need theory for $`d\mathrm{\Gamma }/dx`$ as a function of $`x`$. In the simplest approximation, the $`b`$ quark in a $`B`$ meson carries all of the four-momentum of the meson and decays into a two body state $`s\gamma `$. Then $`d\mathrm{\Gamma }/dx=\delta (1x)`$. In a more realistic picture, we expect the photon spectrum to be spread out, as illustrated in the left-hand part of Fig. 4. If one calculates Feynman diagrams like the simple one in the right-hand part of Fig. 4, one gets contributions that are singular at $`x1`$ (before smearing by the wave function). In the region $`(1x)1`$, one can sum a series containing more and more powers of $`\mathrm{log}(1x)`$. However, this series does not converge well. E. Gardi $`^\mathrm{?}`$ reported on this. He blamed the bad behavior of the perturbative series on an infrared renormalon, which is associated with a factor of $`n!`$ at $`n`$th order in perturbation theory. The simple way to understand this is that diagrams like that illustrated on the right in Fig. 4 get contributions from the region in which the loop momentum $`l`$ is smaller than, say, 1 GeV. The $`n!`$ factor is the price we pay for applying perturbation theory in a region in which perturbation theory does not really apply.
Of course, this kind of behavior is generic in QCD perturbation theory. However, the renormalon behavior is particularly bothersome in this instance. Gardi reported that the bad behavior stems from the use of the “pole mass” in the heavy quark propagator. With a reorganized calculation, he reported that the integration is less sensitive to the infrared region and the predictive power of the theory can be improved.
At this conference, J. Walsh reported new results from BaBar and Belle on $`d\mathrm{\Gamma }/dx`$ for this process. Andersen and Gardi $`^\mathrm{?}`$ were able to take this data and compare it to their theoretical results. They compared theory and experiment for the shape of the spectrum as represented by
$$\frac{1}{\mathrm{\Gamma }}_{E_0}^{\mathrm{}}𝑑E_\gamma E_\gamma \frac{d\mathrm{\Gamma }}{dE_\gamma }$$
as a function of the lower endpoint $`E_0`$ of the integral. The result is shown in Fig. 5. Theory and experiment appear to agree pretty well.
## 4 Pentaquark states
There was some discussion at the conference of the theoretical description of the state $`\mathrm{\Theta }^+(1540)`$, which in a quark description would be a $`udud\overline{s}`$ state. At the time of the conference, the experimental evidence for the existence of this state was mixed. This state is often discussed as a state consisting of two diquarks and an antiquark (see the talk of A. Vainshtein $`^\mathrm{?}`$ on diquarks). An alternative view, reviewed by M. Praszalowicz,$`^\mathrm{?}`$ derives the state in a chiral quark soliton model. This model is based on nature being near the chiral limit, in which the $`\pi `$ and $`K`$ mesons have zero mass, and near the limit of having an infinite number, $`N_c`$, of colors. Praszalowicz argued that for the $`\overline{10}`$ representation of flavor SU(3) to which the $`\mathrm{\Theta }^+(1540)`$ would belong, we are far from the $`N_c\mathrm{}`$ limit. Thus QCD theory does not require the existence of this state. That’s good, because the existence of this state is looking more doubtful: on the day this conference ended, the COMPASS experiment reported its non-observation of the $`\mathrm{\Theta }^+(1540)`$.$`^\mathrm{?}`$
## 5 Lattice studies
I have mentioned already lattice studies of the matrix elements for weak interaction decays. We heard two other kinds of lattice studies. One concerned the gluon propagator, the other was about one of the signals for the quark gluon plasma in heavy ion collisions.
Over the years, there have been quite a lot of studies in QCD of gluon two-point functions and three-point functions, and more generally $`n`$-point functions for quarks and gluons. Such studies begin with perturbation theory since they are based on the Schwinger-Dyson equations, but they go beyond any fixed order of perturbation theory by solving the equations within some approximation scheme. Since the $`n`$-point functions are gauge dependent, such studies must pick a gauge. F. de Soto Borrero reported on studies of the gluon two- and three-point function in QCD without dynamical quarks,$`^\mathrm{?}`$ based not on the Schwinger-Dyson equations but instead on a direct simulation in lattice gauge theory. One interesting way to express the results is to plot the three point function divided by the cube of the two point function, all of these suitably renormalized. That gives the one-particle-irreducible three-point function and thus an effective coupling that is a function of the momentum $`p`$ on each leg of the graph. I should caution that the effective coupling thus defined is a convenient object to use in discussing the theory, but is not directly observable in nature. In Fig. 6, I display the result, showing that the effective squared coupling $`\alpha _{\mathrm{MOM}}`$ thus defined increases as the momentum $`p`$ decrease, then reaches a maximum and decreases. The nice thing is that one can understand the large $`p`$ behavior using the operator product expansion and the small $`p`$ behavior using a picture involving the classical solutions of the (Euclidean) equations of motion known as instantons.
One signal for the quark-gluon plasma that one hopes to see in heavy ion collisions is the melting of $`J/\mathrm{\Psi }`$ states. At relatively low temperatures, hadronic matter allows the existence of a $`J/\mathrm{\Psi }`$ resonance. But at higher temperatures in the plasma phase, the $`J/\mathrm{\Psi }`$ can no longer maintain itself. R. Petreczky presented lattice results that test this prediction. The results are shown in Fig. 7. We see that the resonance is present at low temperatures but for temperatures well above the plasma phase transition temperature $`T_c`$, it has melted away. The strength of the resonance has substantially diminished at $`2.25T_c390\mathrm{MeV}`$.
## 6 Jet Physics
A number of talks at this conference dealt with the physics of jets. I will touch on a few topics that seemed to me to be of particular interest.
In Run I at the Fermilab Tevatron, jet cross sections were typically measured using a cone algorithm. The idea is that a jet consists of particles whose momentum vectors lie in a cone centered on a jet axis. This sounds simple. However, it is not so simple when one considers what to do with overlapping cones. For Run II, the cone algorithm has been improved with respect to sensitivity to effects from soft particles, but the complications remain. An alternative is the $`k_T`$-algorithm, which is modelled on that used in electron-positron annihilation. This algorithm is simple enough to state completely in a short paragraph. The idea is that one successively combines “nearby” subjets, thus capturing the characteristic of QCD that there are jets within jets over a wide range of transverse momentum scales. At Run II at the Tevatron, CDF and D0 have been experimenting with the use of this algorithm. A. Kupco presented data from CDF showing that the $`k_T`$-algorithm works well in a practical sense.$`^\mathrm{?}`$ The data is displayed in Fig. 8.
It has recently become possible to calculate three-jet quantities in hadron collisions at next-to-leading order. An interesting example of this is the correlation in azimuthal angle $`\mathrm{\Delta }\varphi `$ between two observed jets (out of three or more jets in the event). A. Kupco $`^\mathrm{?}`$ showed results from D0 on this quantity. The NLO theory does well in predicting the result over a wide angular range as seen in Fig. 9. (In the region close to $`\mathrm{\Delta }\varphi =\pi `$ we are close to the two jet region, so that fixed order perturbation is not expected to work.) The NLO theory is evidently significantly better than the LO three-jet theory, which inevitably fails as one gets near $`\mathrm{\Delta }\varphi =\pi /2`$ because the two jets with the largest transverse momenta cannot have $`\mathrm{\Delta }\varphi =\pi /2`$ unless they recoil against at least two other jets. A separate graph in this talk illustrated that the theoretical prediction from the Pythia Monte Carlo event generator is sensitive to the available tuning parameters in the $`\mathrm{\Delta }\varphi \pi /2`$ region. Presumably Pythia would have done better if it matched to the exact tree level four-jet matrix element.
Jet physics took a major step forward in the 1980s with the analysis of three-jet events at the PETRA accelerator at DESY. The data is still useful, as illustrated in the talk of S. Kluth.$`^\mathrm{?}`$ He and collaborators used data from the Jade experiment at PETRA to extract the four jet rate at several values of $`\sqrt{s}`$ using the Durham (or $`k_T`$) jet algorithm. It is the progress of QCD theory that has made this difficult job worthwhile: neither the Durham algorithm nor a calculation of four jet rates at next-to-leading order were available at the time of the Jade experiment. The results provide a new way to extract $`\alpha _s(\sqrt{s})`$ from the Jade data for $`\sqrt{s}`$ in the PETRA energy range.
## 7 Perturbative calculations
In the prior sections we have seen examples of the progress in QCD theory over the years, as in the comparison of four-jet data from the Jade experiment in the early 1980s to NLO four jet calculations that have only been accomplished in the past few years. Continuing progress in theory was reported at this conference.
For a “high value” cross section like Higgs production via gluon-gluon fusion, it has been on everyone’s wish list to have a next-to-next-to-leading order calculation. Furthermore, one would like to be able to calculate the cross section for an arbitrary infrared safe observable involving the Higgs boson and other partons created in the process. This is a very difficult problem, and in my opinion the best way to do this ultimately will be to have a system of subtractions that take care of all the singularities in the partonic matrix element – based on the general singularity structure of QCD. However, progress in this direction has been slow. C. Anastasiou discussed a program that does just this task based not on knowing the singularity structure of QCD but on letting a computer find the singularities.$`^\mathrm{?}`$ For three final state partons, one maps the momenta $`\{p_1,p_2,p_3\}`$ constrained to have the sum of the transverse momenta vanish into seven variables $`\{x_1,x_2,x_3,x_4,x_5,x_6,x_7\}`$ in a seven dimensional hypercube. The mapping should be such that the singularities are at the edges and faces of the hypercube. Then the computer is asked to find the singularities and make the appropriate subtractions. I can illustrate the idea in a simple fashion if I use only two variables $`x`$. Consider an integral of the form
$$I=_0^1𝑑x_1x_1^ϵ_0^1𝑑x_2\frac{f(x_1,x_2)}{x_1},$$
(2)
where the factor $`x_1^ϵ`$ mimics the result of using dimensional regularization and $`f(x_1,x_2)`$ is a smooth function that could be quite complicated. We would like to separate this into a pole term proportional to $`1/ϵ`$ and a term that is finite as $`ϵ0`$. To this end, we ask our computer to notice the singularity at $`x_10`$ and write the integral as
$$\begin{array}{cc}\hfill I=& _0^1𝑑x_1x_1^ϵ_0^1𝑑x_2\left\{\frac{f(0,x_2)}{x_1}+\frac{f(x_1,x_2)f(0,x_2)}{x_1}\right\}\hfill \\ \hfill & \frac{1}{ϵ}_0^1𝑑x_2f(0,x_2)+_0^1𝑑x_1_0^1𝑑x_2\frac{f(x_1,x_2)f(0,x_2)}{x_1}.\hfill \end{array}$$
(3)
The first term is the pole term, the second is the finite term. Both integrals can be computed by numerical integration. In real life, the situation is much more complicated, but it appears from the results presented in $`^\mathrm{?}`$ that this is a practical way to proceed.
By going from NLO calculations to NNLO calculations, one reduces the estimated theory error for the prediction of a certain class of cross sections. There is also another needed direction for improvement of calculations.
Conventional NLO calculations do not do well at predicting the probability to emit soft particles or at describing the inner structure of jets. That is why one restricts their use to the prediction of those cross sections that are “infrared safe” in the sense that they are not sensitive to soft particles in the final state or to how jets are divided into subjets. For instance, suppose that one were to take a program that gives the cross section for $`p+\overline{p}W+\mathrm{𝑗𝑒𝑡}+X`$ at NLO. Now imagine using the program to predict the distribution of masses for the jet. The result would be a certain number of events with jet masses near 10 GeV, more events with jet masses near 5 GeV, many more with jet masses near 1 GeV, and yet more with yet smaller jet masses. This all adds up to an infinite number of events, but the situation is saved by having an infinitely negative number of events with jet mass zero. Evidently, this is not a satisfactory description of nature. Nor is it satisfactory that the jets are made of one or two partons rather than many hadrons.
Clearly it would be better to let the partons fragment to form parton showers and then let the parton showers form hadrons in the style of event generator Monte Carlo programs like Pythia and Herwig. It is, however, not so easy to do this while maintaining the next-to-leading order accuracy of the calculation for the infrared safe quantities that the NLO program was designed to get right.
There has been some progress in this area in recent years,$`^\mathrm{?}`$ but the presently existing programs are either tied to a specific Monte Carlo event generator or a limited class of processes and are not specified in a simple general algorithm that authors of NLO calculations could easily use.
At this conference, Z. Nagy $`^\mathrm{?}`$ presented an algorithm for extending an NLO calculation in a fashion that would allow partonic events from the modified NLO program to be fed to a Monte Carlo event generator for showering and hadronization, on the condition that the showering be started for each event with suitable initial conditions. The algorithm is presented for jet production in electron-positron annihilation, but it should be possible to extend it to lepton-hadron and hadron-hadron collisions. The idea is to base the algorithm on the Catani-Seymour dipole subtraction scheme $`^\mathrm{?}`$ that is widely used in constructing NLO programs. In addition, the algorithm makes use of the $`k_T`$ jet-matching scheme $`^\mathrm{?}`$ for switching descriptions among, for example, $`p+\overline{p}W+X`$, $`p+\overline{p}W+1\mathrm{𝑗𝑒𝑡}+X`$, and $`p+\overline{p}W+2\mathrm{𝑗𝑒𝑡𝑠}+X`$. This is so far an algorithm but not yet code. Perhaps we will see a working program at the XLI Rencontre de Moriond.
## 8 QCD in the high energy/small $`x`$ limit
At this conference there was quite a lot of discussion of how QCD behaves when seen with experiments that, in one view, probe the gluon distribution at small momentum fraction $`x`$. Relevant experiments include deeply inelastic scattering at small Bjorken $`x`$ and heavy ion collisions at RHIC.
An outsider like me can understand at least part of this discussion with the aid of a few of simple observations. First, the whole subject hinges on how systems with vastly different rapidities can interact. For instance, if one is interested in the total cross section for scattering two nuclei at high energy, the rapidity difference $`\mathrm{\Delta }y=\mathrm{log}(s/M^2)`$ is large. In deeply inelastic scattering, the rapidity difference between the virtual photon and the proton is $`\mathrm{\Delta }y=\mathrm{log}(Q/m)+\mathrm{log}(1/x)`$ and is large if $`x`$ is small (even if $`Q/m`$ is not so large). Second, systems with a large rapidity difference can communicate easily by exchanging gluons, so that the whole subject hinges on gluons. Third, experiment shows that lots of gluons are to be found in the proton carrying a small fraction $`x`$ of the proton’s momentum. Thus the small $`x`$ gluons appear to be densely packed inside the transverse size of a proton or nucleus. With respect to a probe that might measure this gluon distribution by scattering from it, the hadron looks black. This suggests a system that is purely non-perturbative and thus resistant to theoretical analysis. However, it is useful to think of our probe as a fast moving color dipole. For instance, in deeply inelastic scattering, the virtual photon creates a quark-antiquark pair with net color zero – that is a color dipole. This leads to the fifth observation: If the dipole is small enough, the dense gluonic system is transparent to the dipole. Thus there is some size $`1/Q_s(x)`$ for which the gluonic system just starts to scatter the dipole. At small $`x`$, the density of gluons is so high $`1/Q_s(x)`$ is small. That is, the “saturation scale” $`Q_s(x)`$ is big. This means that there is naturally a big transverse momentum scale in the problem and perturbation theory can be helpful. This analysis does not carry us very far into the subtleties of the theory, but it is perhaps useful in getting us started.
One approach to the high energy limit of QCD is the much studied BFKL equation. There is a leading order version of this equation and, more recently, next-to-leading order corrections. However, it has appeared that the general approach is unstable, with the corrections leading to big changes in the behavior of the solutions. At this conference, L. Schoeffel $`^\mathrm{?}`$ discussed ways of modifying the equation so as to avoid some of the problems and showed results for the solutions of the resulting equation. J. Andersen $`^\mathrm{?}`$ took a different approach, arguing that the widely used Mellin transform method of solving the BFKL equation leads to problems, so that one should instead use an iterative numerical approach. If one defines an effective pomeron intercept $`\alpha _{IP}`$ by $`xg(x)x^{\alpha _{IP}}`$, then Andersen finds a sensible behavior for $`\alpha _{IP}`$ as a function of rapidity, as shown in Fig.!10.
Another approach to understanding the high energy limit of QCD is to go beyond the BFKL equation. There were separate talks on this by E. Iancu,$`^\mathrm{?}`$ by D. Triantafyllopoulos $`^\mathrm{?}`$ and by G. Soyez.$`^\mathrm{?}`$ Here, I follow mostly Iancu and Triantafyllopoulos. One considers the interaction of a color dipole with a hadron to be created by the exchange of $`n`$ gluon ladders, or “pomerons.” Letting $`T_n`$ be contribution from for $`n`$ pomerons, one writes an equation for the variation of $`T_n`$ with the rapidity $`y`$ of the dipole relative to the hadron,
$$dT_n/dy=K_{\mathrm{BFKL}}T_n+K_{\mathrm{split}}T_{n1}+K_{\mathrm{merge}}T_{n+1}.$$
(4)
Here the $`T_n`$ are functions of transverse position arguments and the various $`K`$ are kernels in an integral equation. This general scheme incorporates the linear evolution of a gluon ladder, the splitting of one ladder to make two, and the joining of two gluon ladders to make one. The first term, for $`n=1`$ is the leading order BFKL equation. Triantafyllopoulos argued that the splitting and merging terms are more important than the next-to-leading order corrections to the BFKL equation. Iancu argued that one has some hope of making progress with the splitting-joining picture, even though it is complicated, since the problems are related to problems that have been studied in statistical physics.
## 9 Automating the scientific method
Part of the purpose of this talk was to provide some assessment of where we stand with respect to the problems and opportunities facing QCD theory. Of particular interest is the upcoming beginning of the LHC era. One anticipates that what is seen may not be in agreement with the Standard Model. Since the incoming particles and many of the outgoing particles in an LHC event carry the strong interactions, one is going to have to incorporate QCD theory into the interpretation of whatever signals may be found. Are we ready?
B. Knutsen $`^\mathrm{?}`$ presented computer tools that can help in the interpretation of the data. First, there are tools for finding discrepancies between the data and Standard Model predictions. Then, there are tools for testing hypotheses against the data. Finally, there is a tool for generating hypotheses and testing them. One might call this last step automating theorists.
I am a bit skeptical about this last part. Of course, it is not the theorists who construct computer programs for doing calculations who are in danger of being replaced by automatons: their programs would just be added to the array of programs that does hypothesis testing. Rather, it is the model building theorists whom we perhaps don’t need. However, my guess is that we will be faced with a more difficult problem than simply filling in the blanks for the Minimal Supersymmetric Standard Model, so that some clever ideas from people like those who built the MSSM will be required.
Overall, it seems to me that we ought to automate whatever we can. However, I was struck by Knutsen’s comment that high energy physics is behind astrophysics, heavy ion physics, and biology with respect to automation. Accordingly, I would like to present a real example from current developmental biology that provides an interesting comparison to high energy physics.
The example concerns the zebrafish, a small tropical fish that is easily raised in tanks. This fish has the wonderful feature that its embryos are quite transparent, so that the parts of an embryo can be studied under a microscope as the embryo develops. Study of how the embryo develops can help us figure out the mechanisms for development. One finds $`^\mathrm{?}`$ that in the early embryo there is a band of cells that are able to be either pituitary gland cells or else cells in the lens of the fish’s eye (Fig. 11). How, then, do these cells decide what to do? Experiments $`^\mathrm{?}`$ show that there are some nearby cells that express a signalling molecule called Hedgehog. The dual purpose cells that get a strong Hedgehog signal become pituitary gland cells, while those that get a weaker signal become lens cells (Fig. 12). One can test this by using a mutant fish that does not produce Hedgehog. This amounts to changing the DNA of the starting fish egg.
The experiments in biology that lead to such discoveries are difficult. Things would be easier if biologists could automate the process by producing a computer program, which might be called Icthia, that could simulate fish development (Fig. 13). The input would be the fish genome. Then the program would simulate the complete development of the fish, including the signalling pathways just discussed. In the end, we would have a model final state fish to compare to an experimental fish.
Of course, there is no such program. Biology is too hard for that. But physics is easier. We do have programs, for instance Pythia, that take the lagrangian of the Standard Model or various extensions as input and simulate the entire development of high energy physics events, from an initial state shower to a hard interaction such as, say, a quark and antiquark producing a squark and antisquark, followed by parton showering and decays, followed by hadronization, producing final states of events that can be compared to experimental events (Fig. 14).
We thus have, I submit, theoretical tools for the analysis of the upcoming LHC experiments that are extraordinarily powerful. The event generator Monte Carlo programs just described have a very wide scope. There are also well automated tools for creating tree level cross sections for an even wider variety of partonic processes. Where more accuracy is needed for suitably inclusive cross sections, we have NLO programs and even some NNLO programs. These tools have been developed over the last couple of decades and are now, I think, adequate to the task at hand. Furthermore, these theory tools are currently being improved. It was a pleasure to hear about some of the progress at this conference. The more improvements we have, the better off we will be in trying to understand what we find at the LHC.
## References
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# Vacuum-stimulated cooling of single atoms in three dimensions
## Abstract
Taming quantum dynamical processes is the key to novel applications of quantum physics, e.g. in quantum information science. The control of light-matter interactions at the single-atom and single-photon level can be achieved in cavity quantum electrodynamics, in particular in the regime of strong coupling where atom and cavity form a single entity. In the optical domain, this requires a single atom at rest inside a microcavity Ye99 ; McKeever03 ; McKeever04:2 ; Boca04 ; Maunz04 ; Maunz05 ; Mossberg91 ; Doherty97 ; Horak97 ; Vuletic00 ; Vuletic01 ; Domokos04 ; Murr03 . We have now discovered that an orthogonal arrangement of a cooling laser, a trapping laser and the cavity vacuum gives rise to a unique combination of friction forces that act along all three directions. This novel combination of cooling forces is applied to catch and cool a single atom in a high-finesse cavity. Very low temperatures are reached, and an average single-atom trapping time of 17 seconds is observed, which is unprecedented for a strongly coupled atom under permanent observation.
Cooling and trapping of single atoms in a micro-cavity is difficult, mainly because of the limited access and the complexity of the setup. Long trapping times have been achieved for ions, but not in the small cavities required for strong coupling Guthoerlein01 ; Mundt02 . Neutral atoms, in contrast, have been stored in the potential wells of a standing-wave dipole laser resonant with the micro-cavity Ye99 ; McKeever03 ; McKeever04:2 ; Boca04 ; Maunz04 ; Maunz05 . In these experiments, the surprisingly short trapping times originated mainly from the axial geometry of the laser-cavity system. We have now changed this geometry and use a standing-wave dipole laser oriented perpendicular to the cavity axis. Moreover, an additional pump laser induces rapid three-dimensional cooling, an effect not anticipated for a deep dipole trap. Our findings result in a deterministic strategy for assembling a permanently bound and strongly-coupled atom-cavity system.
As sketched in Fig. 1, our technique employs a standing-wave dipole-force trap and a pump beam that cross in the centre of a high-finesse optical cavity, and run perpendicular to the cavity axis. The pump beam is near-resonant with the cavity, so that an atom in the crossing point scatters pump light into the cavity. The momentum kicks the atom experiences when scattering these photons lead to cooling of the atomic motion along the pump and cavity directions. This process is strongly enhanced by the Purcell effect Purcell46 and has the unique advantage that cooling is effective over a large range of atomic transition frequencies. It therefore allows one to catch a free atom on its flight through the cavity and to cool it down to the bottom of a deep potential well of the dipole trap, even though the average trap-induced AC-Stark shift of the atom increases as the localization improves. Strong cooling forces also act along the standing-wave axis and are caused by the delayed response of the atomic excitation and the intra-cavity field to a change in the atomic position. In our experiment, such a cold atom is well localized at an antinode of the standing-wave trap.
The origin of the cooling forces can be understood from a simple model based on a two-level atom. Starting from a master equation which describes the coupling of the atom to the cavity mode, the pump, and the dipole trap, one obtains four velocity-dependent forces, that we will discuss in more detail elsewhere. The first two cooling forces below have been theoretically predicted by Vuletić et al. Vuletic00 ; Vuletic01 and Domokos et al. Domokos04 . We now add a dipole trap that allows for cooling along an additional third direction. Consider an atom that is exposed to a retro-reflected pump beam of photon momenta $`\pm \mathrm{}𝒌_P`$, with a frequency $`\omega _P`$ close to the resonance $`\omega _C`$ of a surrounding cavity. The cavity provides a means for removing kinetic energy from the atom if the cavity resonance is blue detuned with respect to the pump $`(\mathrm{\Delta }_C=\omega _C\omega _P>0)`$. In this case, friction along the pump beams is caused by preferential absorption of photons traveling in the opposite direction than the atom. The Doppler effect shifts these photons towards the blue by $`|𝒌_P𝒗|`$, such that the coupled atom-cavity system becomes resonant with counter-propagating pump photons as soon as $`\mathrm{\Delta }_C+𝒌_P𝒗0`$. This gives rise to a friction force,
$$𝑭_P=4\mathrm{}𝒌_P(𝒌_P𝒗)\frac{\kappa \mathrm{\Delta }_C}{(\mathrm{\Delta }_C^2+\kappa ^2)^2}g^2P_E,$$
(1a)
along each pump beam. For a fixed low occupation probability of the atom’s excited state, $`P_E\mathrm{\Omega }^2/(\mathrm{\Delta }_A^2+\gamma ^2)`$ (valid for low saturation, where $`2\mathrm{\Omega }`$ is the Rabi frequency of the pump, $`\gamma `$ is the polarisation decay rate of the atom, and $`\mathrm{\Delta }_A=\omega _A\omega _P+\mathrm{\Delta }_S`$ is the effective pump detuning from the atomic resonance, $`\omega _A`$, with $`\mathrm{\Delta }_S`$ being the dipole-trap induced AC-Stark shift), the friction is only determined by the cavity parameters ($`g`$ the atom-cavity coupling constant, $`\kappa `$ the field-decay rate of the cavity). Momentum kicks from photon emissions into the resonator lead to a similar force that acts along the cavity axis,
$$𝑭_C=4\mathrm{}g(g𝒗)\frac{\kappa \mathrm{\Delta }_C}{(\mathrm{\Delta }_C^2+\kappa ^2)^2}P_E.$$
(1b)
Photons emitted into the direction of motion are blue detuned due to the Doppler shift. By recoil, these forward emissions also cool the atomic motion. If now the cavity is blue detuned, the emissions into the direction of motion are favored, and hence the atom is cooled along the cavity axis. This simple picture does not consider interference effects leading to a spatial modulation of the cavity field, but it holds true if the force in Eq. (1b) is averaged over a spatial period.
The two above forces are here seen as due to distinct Doppler effects. However, they have a common origin, namely the dependency of the cavity field on the position of the atom. When the atom moves, the field takes a certain time ($`\kappa ^1`$) to adjust to a new steady-state. This time lag gives also rise to a third force along the standing wave axis,
$$𝑭_S^{Cav}=4\mathrm{}\mathrm{\Delta }_S(\mathrm{\Delta }_S𝒗)\frac{\kappa \mathrm{\Delta }_C}{(\mathrm{\Delta }_C^2+\kappa ^2)^2}\frac{g^2P_E}{\mathrm{\Delta }_A^2+\gamma ^2}.$$
(1c)
It acts along the direction of the standing wave and depends in the same way on $`\mathrm{\Delta }_C`$ as the two other friction forces. It follows that the cavity, in combination with the pump and the trap, leads to cooling in three dimensions, with forces that all have same order of magnitude.
Our scheme also profits from a novel cooling force that is directed along the dipole trap. It can be explained by noting that an atom close to a node of the standing wave is not subject to an AC-Stark shift. For a pump frequency resonant to the atomic transition frequency $`\omega _A`$, the atom is resonantly pumped to the excited state. The atom now gains potential energy when moving in the standing wave, which is then lost by spontaneous emission (rate $`2\gamma `$) to the ground state. This is a Sisyphus-like cooling mechanism Dalibard85 ; Taieb94 that uses two different fields for trapping and cooling. The resulting force,
$$𝑭_S^{Sis}=4\mathrm{}\mathrm{\Delta }_S(\mathrm{\Delta }_S𝒗)\frac{\mathrm{\Delta }_A}{2\gamma (\mathrm{\Delta }_A^2+\gamma ^2)}P_E^2,$$
(1d)
is cavity-independent and provides cooling also if $`\mathrm{\Delta }_C<0`$. As shown below, this force alone is sufficient to increase the trapping time. However, for $`\mathrm{\Delta }_C>0`$, cavity forces dominate, and permanent photon scattering into the cavity takes place. Apart from cooling, all forces fluctuate and lead to heating of the atomic motion. The heating rate follows the Lorentzian-shaped cavity resonance (apart from Eq. (1d)). In analogy to free-space laser cooling Cohen92 , a cavity-Doppler limited temperature around $`k_BT\mathrm{}\kappa `$ is expected for $`\mathrm{\Delta }_C\kappa `$. This unique combination of friction forces is unprecedented, and it allows one to cool single dipole-trapped atoms in a cavity along all three directions. We here apply this novel combination of cooling forces to catch and cool a single atom in a high-finesse cavity, as discussed in the following.
In the experiment, we use a dipole-force trap to guide <sup>85</sup>Rb-atoms over a distance of 14 mm from a magneto-optical trap (MOT) into the cavity. This trap is formed by a single horizontally running beam of an Yb:YAG laser, with focus between MOT and cavity. Once the atoms reach the cavity, we switch to a standing-wave dipole trap, formed by a pair of counter-propagating Yb:YAG laser beams (2 Watt, 1030 nm, waist w$`{}_{0}{}^{}=16\mu `$m), tightly focused in the centre of the cavity. The antinodes of the standing wave represent 2.5 mK deep potential wells, i.e. an AC Stark shift of the atomic transition frequency in the centre of the wells of $`\mathrm{\Delta }_S=2\pi \times 100`$MHz.
The 0.5 mm long cavity is formed by two mirrors of 5 cm radius of curvature that have different transmission coefficients (2 ppm and 95 ppm). The relevant atom-cavity parameters are $`(g_0,\kappa ,\gamma )=\mathrm{\hspace{0.17em}2}\pi \times (5,5,3)`$MHz, where $`g_0`$ is the atom-cavity coupling constant in an antinode, averaged over all magnetic sublevels of the $`5^2S_{1/2}(F=3)`$ to $`5^2P_{3/2}(F^{}=4)`$ transition. A Pound-Drever-Hall technique is used to lock the frequency of the TEM<sub>00</sub> mode to the atomic transition frequency. For this purpose, we use a reference laser that is red detuned by eight free spectral ranges (5 nm) from the atomic resonance. This laser acts as an additional standing-wave dipole trap along the cavity axis. Its potential wells are about 30 $`\mu `$K deep and show a good overlap with the resonant mode in the centre of the cavity. Together with the Yb:YAG standing-wave trap, it forms a 2D optical lattice (calculated trap frequencies: $`\nu _{sw}670`$ kHz in direction of the standing-wave trap, $`\nu _{cav}`$ 100 kHz along the cavity axis, and $`\nu _{}10`$ kHz orthogonal to the cavity and the trapping laser).
In addition to these two conservative dipole traps, we continuously drive the $`F=3`$ to $`F^{}=4`$ transition with a pump laser that runs orthogonal to the cavity axis, at an angle of 45 to the standing-wave trap (Fig. 1). Together with a repump laser $`(F=2F^{}=3)`$, the beam has a focus $`(w_0=35\mu `$m$`)`$ at the intersection of cavity mode and dipole trap, and is retroreflected to balance its radiation pressure. To avoid an intensity modulation of the pump, the two counter-propagating beams have orthogonal polarization. Averaged over all magnetic sub-levels, it drives the atoms with a Rabi frequency $`2\mathrm{\Omega }=2\pi \times 30`$MHz (all data, except Fig. 2A). Under these circumstances, the Purcell effect gives rise to photon scattering into the cavity at a rate
$$R_{scat}2\kappa \frac{g^2}{\mathrm{\Delta }_C^2+\kappa ^2}\frac{\mathrm{\Omega }^2}{\mathrm{\Delta }_A^2+\gamma ^2}$$
(2)
for a single atom. For $`\mathrm{\Delta }_C=0`$, this gives a scattering rate of $`R_{scat}1400/`$ms. In the experiment, we must take into account that the atom stops fluorescing whenever it falls into the dark state $`5^2S_{1/2}(F=2)`$, and that repumping to $`F=3`$ takes some time due to the large AC-Stark shift. This leads to blinking, and from the measured count rate, we estimate that photons are scattered at the above rate only 1/5 of the time. Furthermore, due to losses in the imaging system and the limited quantum efficiency of the detector, only 5% of the photons that are scattered into the cavity mode are finally detected behind the cavity mirror of higher transmission.
Once the atoms have been brought into the vicinity of the cavity, they are randomly distributed over the potential wells of the standing wave, with a small probability that an atom actually sits in the cavity. Some atoms have enough kinetic energy to move from well to well, and as soon as such an atom enters the intersection of cavity mode and pump laser $`(\mathrm{\Delta }_C/2\pi =+2`$MHz and $`\omega _P=\omega _A`$ for capturing$`)`$, it scatters photons into the cavity and is therefore cooled. In Fig. 2A, an experimental trace is shown where a single atom suddenly appears in the cavity. First the photon scattering rate is high, as the initially hot atom is poorly localized in the trap and the average experienced Stark shift $`\mathrm{\Delta }_S`$ is small. This changes as the atom gets colder and therefore is better localized in the potential well. After $`\mathrm{\Delta }t100\mu `$s, the atom reaches its final temperature and scatters at a much lower (but constant) rate, since it now experiences a much higher $`\mathrm{\Delta }_S`$ close to the bottom of the potential. Starting from the simple estimation that the total kinetic energy, $`E`$, is lost during $`\mathrm{\Delta }t`$ with $`\dot{E}E/\mathrm{\Delta }t`$, one calculates a mean friction coefficient, $`\beta =\dot{E}/2E1/2\mathrm{\Delta }t`$, between $`5/`$ms (raw data) and $`25/`$ms (assuming blinking as discussed above). This is reasonably close to the expected value, $`\beta =14/`$ms, that one obtains from Eqs. (1c) and (1d) with $`𝑭=\beta m𝒗`$ ($`m`$ the atomic mass).
To prevent further atoms from penetrating into the cavity mode, we apply a filtering procedure. This is accomplished by a 10 ms long interruption of the standing-wave trap. During this time, atoms inside the cavity stay trapped in the shallow intra-cavity dipole trap, but the other atoms are lost. If we turn the trap off and back on adiabatically, the probability for a caught atom to survive this procedure is higher than 50%. As shown in Fig. 2B, we then measure the photon rate to determine the exact number of atoms in the cavity McKeever04:2 . The signal shows only small variations, which indicate a nearly constant atom-cavity coupling, and pronounced individual steps, that occur whenever an atom is lost.
As shown in Fig. 2C, dipole-trapped atoms show an average lifetime of 2.7 s if we switch off the pump laser. However, if an atom is continuously illuminated and coupled to the cavity with $`\mathrm{\Delta }_C0`$, the lifetime increases, reaching values exceeding 20 s. This impressively demonstrates that a strong cooling mechanism is active. For $`\mathrm{\Delta }_C=0`$, in particular, a lifetime of about 17 s is measured. In this case, no cavity cooling is expected, and within our model, the long lifetime comes solely from the Sisyphus-like mechanism, Eq. (1d). Cavity cooling, Eqs. (1a1c), is expected to give much longer trapping times for $`\mathrm{\Delta }_C+\kappa `$. For technical reasons, however, longer trapping times could not be registered in our experiment. We have therefore deliberately modulated the depth of the trapping potential, i.e. the intensity of the Yb:YAG trapping laser, by 30% at a frequency of 7 kHz. This leads to parametric heating and shortens the trapping time, so that systematic lifetime measurements can be made within a reasonable time. Without cooling, the modulation reduces the trapping time to $`(22\pm \mathrm{\hspace{0.17em}5})`$ms. We now have performed measurements such as those in Fig. 2C for several different cavity detunings. The single-atom lifetime as a function of $`\mathrm{\Delta }_C`$ is plotted in Fig. 3. Obviously, the lifetime increases to about 100 ms as soon as the pump laser is present, even if $`\mathrm{\Delta }_C`$ is so large ($`2\pi \times 50`$MHz) that the cavity has no effect. Moreover, a cavity close to resonance significantly extends the lifetime. For a blue detuned cavity, with $`\mathrm{\Delta }_C+\kappa `$, a 20-fold increase of the lifetime to a maximum of $`400`$ms is obtained, whereas a red detuned cavity, with $`\mathrm{\Delta }_C\kappa `$, leads to a minimum lifetime. To compare these results to the theoretical model, the expected friction coefficients are also shown in Fig. 3. Although these values cannot be directly compared with the lifetime, qualitative conclusions can be drawn. First, the cavity-independent Sisyphus cooling, predicted only along the standing-wave axis, accounts for the cavity-independent increase of the lifetime with respect to atoms left in the dark. The variation of the lifetime with $`\mathrm{\Delta }_C`$ finds its correspondence in the predicted course of the friction coefficients due to the cavity cooling for all three dimensions. Obviously, cavity cooling increases the trapping time by a factor of four if compared to Sisyphus cooling alone. Without further limitations, such as background gas collisions, we can therefore expect that the 17 s–long lifetime observed without the trap modulation at $`\mathrm{\Delta }_C=0`$ increases to about one minute in the presence of the cavity cooling at $`\mathrm{\Delta }_C+\kappa `$.
For the same data, we have also analyzed the average count rate per atom using 10 ms-long time bins. Fig. 4 depicts the photon-count histogram. The well-distinct peaks stem from dark counts and traces with one, two or more trapped atoms. From a fit to these data, we can derive the average count rate per atom, $`R_{det}`$, and its statistical spread, $`\sigma (R_{det})`$, that is corrected for shot noise. The results are plotted in the inset of Fig. 4 (with the errorbars indicating $`\sigma (R_{det})`$) as a function of $`\mathrm{\Delta }_C`$. If we assume that all variations of $`R_{det}`$ are caused by variations in the atom-cavity coupling, then we obtain from Eq. (2) $`\mathrm{\Delta }g/g=\sigma (R_{det})/2R_{det}=\pm 8.6`$% (for $`\mathrm{\Delta }_C/2\pi =+4`$MHz). This can only be explained with an atom distribution among the different wells of the standing-wave that is less than $`\pm 9\mu `$m wide. Under the assumption that the distribution in the filter phase is mapped to the standing-wave trap during the adiabatic transfer, this indicates a temperature below $`6\mu `$K during filtering.
To get an estimate of the temperature in the deep 2D-lattice, we analyzed the auto-correlation function of the emitted photon stream. In this signal, a periodic modulation at 2$`\times \nu _{cav}`$ (the trap frequency along the cavity axis) is found, with a visibility of about 5 %. This is indeed expected for an atom that oscillates in the weak intra-cavity dipole trap, without ever reaching the nodes of the cavity mode. The $`k_B\times 30\mu `$K depth of the intra-cavity trap can therefore be seen as an upper limit of the atomic energy. The mean kinetic energy cannot surpass 50 % of this value, which corresponds to a temperature of $`15\mu `$K. Assuming the same temperature in all directions, this limits the mean vibrational quantum number for the motion along the standing-wave dipole trap to $`\overline{n}=0.13`$, i.e. the atom is in the vibrational ground state with at least 88 % probability.
Our experiment demonstrates how to capture single atoms with a three-dimensional cavity cooling scheme. Qualitatively, our results show good agreement with the predictions of a model valid for a two-level atom, although an understanding of the remarkably good localization of the atom calls for a more detailed theoretical analysis, including polarization effects of the incident pump laser Cirac93:3 , the multilevel structure and even the quantum motion of the atom Cirac95:2 ; Zippilli05 . At present, our method allows the preparation of an exactly known number of atoms at low temperatures in the centre of a high-finesse optical cavity. With average trapping times exceeding 15 seconds, single-atom cavity physics is now at a stage where one can start fully controlled atom-photon experiments in the strong atom-cavity coupling regime.
###### Acknowledgements.
This work was supported by the Deutsche Forschungsgemeinschaft (SPP 1078 and SFB 631) and the European Union \[IST (QGATES) and IHP (CONQUEST) programs\]. We are also grateful to Simon Webster for helpful comments.
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# 1 Introduction
## 1 Introduction
The discovery of cosmological acceleration indicates the necessity for modification of the standard cosmological paradigm. One way to explain the accelerated expansion of the Universe is to attribute it to the effect of some unconventional matter such as the cosmological constant or quintessence. In this approach one encounters, however, the difficulty of explaining the small value of the associated energy scale in comparison with other fundamental scales, such as the Planck mass or electroweak symmetry breaking scale. It is reasonable to consider an alternative possibility that the cosmological constant is exactly zero, but its effect is mimicked by modifications of gravitational laws at distance scale of the order of the present cosmological horizon.
A natural framework for modifying gravity at ultra-large distances is provided by models with infinite extra spatial dimensions. A number of brane-world models have been proposed where the four-dimensional graviton is a superposition of bulk KK modes, and thus, is not localized on the brane. Rather, it is quasilocalized, having a non-zero width $`\mathrm{\Gamma }`$ (and, possibly, mass $`m`$) with respect to decay into extra dimensions. In these models gravity is four-dimensional at distances smaller than $`r_c\mathrm{min}\{m^1,\mathrm{\Gamma }^1\}`$ and becomes multi-dimensional at larger scales. The attractive feature of the models with extra-dimensions is that the tiny values of $`m,\mathrm{\Gamma }`$ are generated without strong fine-tuning of parameters.
However, Lorentz invariant extra dimensional models with infrared modifications of gravity face the problems akin to those of four-dimensional massive gravity. Either these theories contain ghosts, fields with the wrong sign of kinetic term, or the propagator of graviton exhibits the van Dam–Veltman–Zakharov (vDVZ) discontinuity originating from a scalar degree of freedom that does not decouple in the massless limit. At the classical level the discontinuity might be cured by nonlinear effects , but at the quantum level the presence of the additional scalar leads to strong coupling at the energy scale at most of order $`(m^2M_{Pl})^{1/3}`$. If $`m`$ is comparable to the present Hubble parameter of the Universe, the theory is in strong coupling regime, and thus looses its predictivity, at distances smaller than 1000 km.
There are indications that these difficulties might be resolved in the Dvali–Gabadadze–Porrati (DGP) model. In particular, it was argued that the scale of strong coupling can be pushed to higher energies in nontrivial classical backgrounds, provided the structure of counterterms in the model is of a special form . However, it is still unclear whether the counterterms in the DGP model actually meet the requirements of Ref. . Probably, the most disappointing fact about the DGP model is that the branch of the cosmological solutions in this theory, which exhibits cosmic acceleration without introduction of the cosmological constant, is plagued by ghost-like instabilities .
In this paper we propose a brane-world model with infrared modification of gravity and violation of 4-dimensional Lorentz-invariance. In our model graviton becomes quasilocalized due to mixing with vector fields which freely propagate in the bulk. This kind of mixing becomes possible once the Lorentz symmetry is spontaneously broken by condensates of the vector fields.
Our approach builds up on the ideas put forward recently in the four-dimensional framework. It was suggested that the problems of massive gravity could be resolved in models incorporating violation of Lorentz invariance . Models with explicit Lorentz violation in the gravitational sector were considered in Refs. and it was shown that they are equivalent to a class of models with gravity coupled to scalar fields with unusual kinetic terms. In the latter formulation the Lorentz symmetry is broken spontaneously by coordinate dependent scalar condensates. Because of nonlinear structure of the scalar kinetic term in this kind of models, they possess a cutoff scale $`\mathrm{\Lambda }_{cutoff}=(mM_{Pl})^{1/2}`$, and can be considered only as low energy effective $`\sigma `$-models. If the graviton mass $`m`$ is of the order of the present Hubble parameter, then $`\mathrm{\Lambda }_{cutoff}(0.01\mathrm{mm})^1`$, which is phenomenologically acceptable, but still unnaturally small in comparison with other fundamental scales. It is unclear whether these $`\sigma `$-models can be extended to complete models making sense above the scale $`\mathrm{\Lambda }_{cutoff}`$.
As a scenario for such a complete model one could envisage a gravitational analog of the Higgs mechanism triggered by spontaneous Lorentz symmetry breaking. In Refs. it was proposed to use coordinate independent vector (in general, tensor) condensates for this purpose. However, in four-dimensional theories, graviton does not acquire mass within this approach: no gap in the dispersion relation of graviton appears . Moreover, the effect of vector condensates on cosmology amounts to nothing but the change of the gravitational constant in the cosmological equations , and no cosmological acceleration is produced.
The brane-world approach opens up a new way to circumvent the problems encountered in four dimensions. In this paper we make the first step along this line: in Sec. 2 we present a model in which graviton becomes quasilocalized due to spontaneous Lorentz symmetry breaking induced by non-zero VEVs of bulk vector fields, and analyze in the subsequent sections the spectrum of linear perturbations about four-dimensionally flat background. In Sec. 3 we analyze the tensor perturbations and find that the characteristic mass $`m_c`$ of KK modes comprising four-dimensional graviton is naturally exponentially small, which allows to take the Lorentz-breaking scale as high as a few tenth of the Planck mass. In Sec. 4 and Sec. 5 we study vector and scalar perturbations and show that the model contains neither tachyons, nor ghosts, and does not exhibit the vDVZ discontinuity. We calculate long distance modification of the Newton law and find quite unexpectedly that attraction of point masses changes to repulsion at distances of order $`1/m_c`$. Thus, the model provides a realization of antigravity at ultra-large distances. This looks promising in view of the possibility to obtain cosmological acceleration at late times. On the other hand, modification of gravitational field of static sources in the model faces severe phenomenological constraints. We discuss possible ways to get around these constraints in the concluding Sec. 6.
A number of important issues are left beyond the scope of the present paper. These are, e.g., cosmology of the proposed model and structure of quantum corrections. We leave them for future investigation.
## 2 The model
We consider a five-dimensional setup with positive tension brane and negative cosmological constant in the bulk — the Randall–Sundrum (RS) model . We add to this model three bulk vector fields $`A_M^a`$, $`a=1,2,3`$, with quartic potential localized on the brane. The action is taken in the following form,
$$\begin{array}{cc}\hfill S=& d^5x\sqrt{g}\left(\frac{R}{16\pi G_5}\mathrm{\Lambda }\frac{1}{4}F_{MN}^aF^{aMN}\right)\hfill \\ & +d^4x\sqrt{\overline{g}}\left(\sigma \frac{\varkappa _1^2}{2}\left(\overline{g}^{\mu \nu }A_\mu ^aA_\nu ^b+v^2\delta ^{ab}\right)^2\frac{\varkappa _2^2}{2}\left(\overline{g}^{\mu \nu }A_\mu ^aA_\nu ^b\frac{1}{3}\delta ^{ab}\overline{g}^{\mu \nu }A_\mu ^cA_\nu ^c\right)^2\right),\hfill \end{array}$$
(1)
where $`\overline{g}_{\mu \nu }`$ is the induced metric on the brane, and
$$F_{MN}^a=_MA_N^a_NA_M^a.$$
The capital Latin indices $`M,N`$ take values $`0,1,2,3,5`$, while the Greek indices $`\mu ,\nu ,\mathrm{}`$ run from $`0`$ to $`3`$. The action (1) is invariant under global $`SO(3)`$ symmetry with vector fields $`A_\mu ^a`$ belonging to fundamental representation. It is straightforward to check that the potential in the second line of Eq. (1) is generic quartic potential invariant under this group. Note that the second term in the potential depends only on the traceless part of the matrix $`\overline{g}^{\mu \nu }A_\mu ^aA_\nu ^b`$. The parameters $`v^2`$, $`\varkappa _1^2`$, $`\varkappa _2^2`$ are assumed to be positive. In addition to the global $`SO(3)`$ symmetry, the bulk part of the action is invariant under $`U(1)\times U(1)\times U(1)`$ gauge transformations,
$$A_M^aA_M^a+_M\alpha ^a.$$
This Abelian gauge symmetry is broken explicitly at the brane.
The model has a static solution with AdS metric in the bulk and constant values of the vector fields,
$`ds^2=dz^2+\mathrm{e}^{2k|z|}\eta _{\mu \nu }dx^\mu dx^\nu ,`$ (2a)
$`A_5^a=A_0^a=0,A_i^a=v\delta _i^a,i=1,2,3,`$ (2b)
where $`k=\sqrt{4\pi G_5\mathrm{\Lambda }/3}`$, and the usual fine-tuning $`\sigma =\mathrm{\Lambda }/k`$ is assumed; the signature of the Minkowski metric $`\eta _{\mu \nu }`$ is $`(+,,,)`$. The VEVs (2b) of the vector fields break $`SO(3)\times `$Lorentz symmetry down to $`SO(3)`$ of spatial rotations accompanied by simultaneous rotations in the internal space. The similar pattern of Lorentz symmetry breaking was considered in the four-dimensional context in Refs. . As we will see, this spontaneous symmetry breaking provides mixing between the graviton zero mode, present in the pure RS case, and the continuum spectrum of KK modes of vector fields. This mixing results in quasilocalization of the graviton.
Let us study the linearized perturbations above the background (2). Imposing the gauge conditions $`g_{5\mu }=0`$, $`g_{55}=1`$ on the metric, one writes the following decomposition,
$`ds^2=dz^2+(\mathrm{e}^{2k|z|}\eta _{\mu \nu }+h_{\mu \nu }(x,z))dx^\mu dx^\nu ,`$
$`A_M^a=v\delta _M^a+a_M^a(x,z).`$
The analysis is simplified by the fact that the energy-momentum tensor of the vector fields in the bulk vanishes to the linear order in perturbations, allowing for fixing the transverse traceless gauge in the bulk ,
$$h_\mu ^\mu =h_{\nu ,\mu }^\mu =0.$$
(3)
Here and throughout the paper indices $`\mu ,\nu ,\mathrm{}`$ are raised (lowered) using the metric $`\eta ^{\mu \nu }`$ ($`\eta _{\mu \nu }`$), and comma denotes derivative. In this gauge the bulk equations for the metric take a fairly simple form,
$$\frac{1}{2}h_{\mu \nu }^{\prime \prime }2k^2h_{\mu \nu }\frac{1}{2[u(z)]^2}h_{\mu \nu ,\lambda }^{,\lambda }=0,$$
(4)
where $`u(z)=\mathrm{e}^{k|z|}`$, and prime denotes derivative with respect to $`z`$. As to the vector fields, the bulk equations read
$$a_{\nu }^{a}{}_{}{}^{\prime \prime }+2ka_{\nu }^{a}{}_{}{}^{}+\frac{1}{u^2}a_{\nu ,\mu }^{a,\mu }=0.$$
(5)
In deriving this equation we made use of the $`[U(1)]^3`$ gauge invariance of the bulk action to impose the conditions
$$a_5^a=0,a_{,\mu }^{a\mu }=0$$
(6)
on the vector fields in the bulk. Imposing the two conditions (6) simultaneously is possible, as they are compatible on-shell. The coordinate frame where the conditions (3) and (6) are satisfied will be referred to below as the bulk frame.
The boundary conditions on the brane are most easily formulated in the Gauss normal (GN) reference frame. In this frame the brane is fixed at $`\overline{z}=0`$ (the quantities with the bar refer to the GN frame), and we assume $`Z_2`$ symmetry across the brane. Then, for general energy-momentum tensor $`T_{\mu \nu }`$ on the brane, the boundary conditions for the metrics at $`\overline{z}=+0`$ read ,
$$\overline{h}_{\mu \nu }^{}+2k\overline{h}_{\mu \nu }=8\pi G_5\left(T_{\mu \nu }\frac{1}{3}\eta _{\mu \nu }T_\lambda ^\lambda \right).$$
(7)
The linearized energy-momentum tensor of the vector fields comes from the Lorentz violating term in the quadratic action arising from the brane potential in Eq. (1). To the quadratic order, the latter term has the form,
$$\begin{array}{cc}\hfill S_{LV}^{(2)}=d^4x& [\frac{\varkappa _1^2v^4}{2}(\overline{h}_{ab}+\frac{1}{v}(\overline{a}_b^a+\overline{a}_a^b))^2\hfill \\ & +\frac{\varkappa _2^2v^4}{2}(\overline{h}_{ab}+\frac{1}{v}(\overline{a}_b^a+\overline{a}_a^b)\frac{1}{3}\delta _{ab}(\overline{h}_{cc}+\frac{2}{v}\overline{a}_c^c))^2].\hfill \end{array}$$
(8)
From this expression one obtains,
$`T_{00}^V=T_{0a}^V=0,`$ (9a)
$`T_{ab}^V=2\varkappa ^2v^4\left(\overline{h}_{ab}+{\displaystyle \frac{1}{v}}(\overline{a}_b^a+\overline{a}_a^b)\right){\displaystyle \frac{2\varkappa _2^2v^4}{3}}\delta _{ab}\left(\overline{h}_{kk}+{\displaystyle \frac{2}{v}}\overline{a}_k^k\right),`$ (9b)
where we introduced $`\varkappa ^2=\varkappa _1^2+\varkappa _2^2`$. Note that the energy-momentum tensor (9) of the vector fields violates the weak energy condition. We will show that this property does not lead to instabilities. At the same time it is crucial for antigravity emerging at ultra-large distances.
For the vector fields, the $`Z_2`$ symmetry and continuity demand that the 5th component vanishes on the brane, $`\overline{a}_5^a(\overline{z}=0)=0`$. The sources for the other components can be read off from (8), resulting in the following boundary conditions at $`\overline{z}=+0`$,
$`\overline{a}_0^a^{}=0,`$ (10a)
$`\overline{a}_i^a^{}=\varkappa ^2v^3\left(\overline{h}_{ai}+{\displaystyle \frac{1}{v}}(\overline{a}_i^a+\overline{a}_a^i)\right){\displaystyle \frac{\varkappa _2^2v^3}{3}}\delta _{ai}\left(\overline{h}_{kk}+{\displaystyle \frac{2}{v}}\overline{a}_k^k\right).`$ (10b)
The metrics in the bulk frame and GN frame are related by a gauge transformation. The form of the latter is restricted by the condition $`\overline{h}_{55}=\overline{h}_{5\mu }=0`$. One obtains
$$\overline{h}_{\mu \nu }=h_{\mu \nu }+u^2(\epsilon _{\mu ,\nu }+\epsilon _{\nu ,\mu })+\frac{1}{k}\epsilon _{,\mu \nu }2ku^2\eta _{\mu \nu }\epsilon ,$$
(11)
where the functions $`\epsilon `$, $`\epsilon _\mu `$ depend only on $`x`$. This transformation corresponds to the following change of coordinates,
$$z=\overline{z}+\epsilon ,x^\mu =\overline{x}^\mu +\frac{1}{2ku^2}\epsilon ^{,\mu }+\epsilon ^\mu .$$
Note that in the bulk frame the brane is displaced from the origin along the $`z`$-direction, its position being described by the brane bending $`\epsilon `$. The relation between the components of the vector fields in the two frames has the form
$`\overline{a}_5^a`$ $`=a_5^a{\displaystyle \frac{v}{u^2}}\epsilon _{,a}+_5\alpha ^a,`$
$`\overline{a}_\mu ^a`$ $`=a_\mu ^a{\displaystyle \frac{v}{2ku^2}}\epsilon _{,a\mu }+v\epsilon _{,\mu }^a+\alpha _{,\mu }^a,`$
where we allowed for the possibility of a $`[U(1)]^3`$ gauge transformation parametrized by the functions $`\alpha ^a(z,x)`$. Imposing the gauge $`\overline{a}_5^a=0`$, which is consistent with the boundary conditions on the brane, we obtain,
$$\overline{a}_\mu ^a=a_\mu ^a+v\epsilon _{,\mu }^a+\beta _{,\mu }^a,$$
(12)
where the functions $`\beta ^a`$ depend only on $`x`$. The set of equations (7), (9)–(11) and (12) results in the following boundary conditions,
$`h_{00}^{}+2kh_{00}+2\epsilon _{,00}=`$ $`{\displaystyle \frac{\lambda _1}{3}}\left(h_{aa}+{\displaystyle \frac{2}{v}}a_a^a+{\displaystyle \frac{1}{k}}\epsilon _{,aa}+6k\epsilon +{\displaystyle \frac{2}{v}}\beta _{,a}^a\right),`$ (13a)
$`h_{0i}^{}+2kh_{0i}+2\epsilon _{,0i}=0`$ $`,`$ (13b)
$`h_{ij}^{}+2kh_{ij}+2\epsilon _{,ij}=\lambda `$ $`\{h_{ij}+{\displaystyle \frac{1}{v}}(a_j^i+a_i^j)+{\displaystyle \frac{1}{k}}\epsilon _{,ij}+{\displaystyle \frac{1}{v}}(\beta _{,j}^i+\beta _{,i}^j)`$
$`{\displaystyle \frac{1}{3}}\delta _{ij}(h_{kk}+{\displaystyle \frac{2}{v}}a_k^k+{\displaystyle \frac{1}{k}}\epsilon _{,kk}+{\displaystyle \frac{2}{v}}\beta _{,k}^k)\},`$ (13c)
$`a_{0}^{a}{}_{}{}^{}=0,`$ (13d)
$`a_{i}^{a}{}_{}{}^{}=\varkappa ^2v^3\{h_{ai}+{\displaystyle \frac{1}{v}}(a_i^a`$ $`+a_a^i)+{\displaystyle \frac{1}{k}}\epsilon _{,ai}+{\displaystyle \frac{1}{v}}(\beta _{,i}^a+\beta _{,a}^i)\}+2\varkappa _1^2v^3k\delta _{ai}\epsilon `$
$`{\displaystyle \frac{\varkappa _2^2v^3}{3}}\delta _{ai}\left(h_{kk}+{\displaystyle \frac{2}{v}}a_k^k+{\displaystyle \frac{1}{k}}\epsilon _{,kk}+{\displaystyle \frac{2}{v}}\beta _{,k}^k\right);`$ (13e)
here we introduced the notations
$$\lambda _1=16\pi G_5\varkappa _1^2v^4,\lambda =16\pi G_5\varkappa ^2v^4.$$
In what follows it is convenient to work in the 4-dimensional Fourier representation,
$$h_{\mu \nu },a_\mu ^a\mathrm{e}^{ip_0x_0+i\mathrm{𝐩𝐱}},$$
and decompose the perturbations into scalar, vector and tensor modes with respect to rotations around the three-momentum $`𝐩`$. To this end one introduces a three-dimensional orthogonal basis consisting of $`𝐩`$ and two unit vectors, $`𝐞^{(\alpha )}`$, $`\alpha =1,2`$. The latter are used to construct a pair of transverse traceless symmetric tensors,
$$d_{ij}^{(1)}=\frac{1}{\sqrt{2}}(e_i^{(1)}e_j^{(2)}+e_j^{(1)}e_i^{(2)}),d_{ij}^{(2)}=\frac{1}{\sqrt{2}}(e_i^{(1)}e_j^{(1)}e_i^{(2)}e_j^{(2)}),$$
and a transverse antisymmetric tensor,
$$f_{ij}=\frac{1}{\sqrt{2}}(e_i^{(1)}e_j^{(2)}e_j^{(1)}e_i^{(2)}).$$
The tensorial decomposition reads,
$`h_{00}=`$ $`\varphi _1,`$ (14a)
$`h_{0i}=`$ $`p_i\varphi _2+e_i^{(\alpha )}\psi _{1\alpha },`$ (14b)
$`h_{ij}=`$ $`\left(\delta _{ij}{\displaystyle \frac{p_ip_j}{𝐩^2}}\right)\varphi _3+p_ip_j\varphi _4+\left(p_ie_j^{(\alpha )}+p_je_i^{(\alpha )}\right)\psi _{2\alpha }+d_{ij}^{(\alpha )}\chi _{1\alpha },`$ (14c)
$`a_0^a=`$ $`p_a\varphi _5+e_a^{(\alpha )}\psi _{3\alpha },`$ (14d)
$`a_i^a=`$ $`\left(\delta _{ia}{\displaystyle \frac{p_ip_a}{𝐩^2}}\right)\varphi _6+p_ip_a\varphi _7+\left(p_ie_a^{(\alpha )}+p_ae_i^{(\alpha )}\right)\psi _{4\alpha }+\left(p_ie_a^{(\alpha )}p_ae_i^{(\alpha )}\right)\psi _{5\alpha }`$
$`+d_{ia}^{(\alpha )}\chi _{2\alpha }+f_{ia}\xi ,`$ (14e)
$`\beta ^a=`$ $`p_a\varphi _8+e_a^{(\alpha )}\psi _{6\alpha },`$ (14f)
with $`𝐩=(p_1,p_2,p_3)`$. There are two fields $`\chi `$ in the symmetric tensor sector; one field $`\xi `$ in the sector of antisymmetric tensors; six fields $`\psi `$ in the vector sector, and eight fields $`\varphi `$ in the scalar sector. Note that all fields in (14), except for $`\varphi _8`$ and $`\psi _6`$, are functions of the fifth coordinate $`z`$.
## 3 Tensor modes
Let us study the spectrum of tensor perturbations. To begin with, one notices that the antisymmetric field $`\xi `$ completely decouples. It satisfies Eq. (5) with the free boundary condition $`\xi ^{}(z=0)=0`$, and forms a continuum spectrum of completely delocalized modes similar to the modes of a free vector field in RS background . We do not consider these modes below.
For the symmetric tensors one obtains from Eqs. (4), (5) the following equations in the bulk,
$`{\displaystyle \frac{1}{2}}\chi _1^{\prime \prime }2k^2\chi _1+{\displaystyle \frac{p^2}{2u^2}}\chi _1=0,`$ (15a)
$`\chi _2^{\prime \prime }+2k\chi _2^{}{\displaystyle \frac{p^2}{u^2}}\chi _2=0,z>0,`$ (15b)
where $`p^2p_0^2𝐩^2`$, and polarization index $`\alpha `$ is omitted to simplify notations. The boundary conditions are read off from Eqs. (13c), (13e):
$`\chi _1^{}+2k\chi _1=\lambda \left(\chi _1+{\displaystyle \frac{2}{v}}\chi _2\right),`$ (16a)
$`\chi _2^{}=\varkappa ^2v^3\left(\chi _1+{\displaystyle \frac{2}{v}}\chi _2\right),z=+0.`$ (16b)
The system (15), (16) can be rewritten in the form of an eigenvalue problem,
$`u^2(\chi _1^{\prime \prime }4k^2\chi _1)\left[(4k2\lambda )\chi _12\mu \widehat{\chi }_2\right]\delta (z)=`$ $`m^2\chi _1,`$ (17a)
$`u^2(\widehat{\chi }_2^{\prime \prime }2k\widehat{\chi }_2^{})+\left[2\mu \chi _1+2\nu \widehat{\chi }_2\right]\delta (z)=`$ $`m^2\widehat{\chi }_2,`$ (17b)
where we denoted $`m^2=p^2`$, $`\nu =2\varkappa ^2v^2`$, $`\mu =\sqrt{\lambda \nu }`$ and introduced $`\widehat{\chi }_2=\sqrt{32\pi G_5}\chi _2`$ to make the operator Hermitian with the scalar product
$$(\eta ,\chi )=_{\mathrm{}}^+\mathrm{}𝑑z\left(\frac{\eta _1^{}\chi _1}{u^2}+\widehat{\eta }_2^{}\widehat{\chi }_2\right).$$
(18)
First, let us make sure that there are no tachyonic modes. The operator (17) is the sum of the diagonal operator appearing in the pure RS case and a positive semi-definite operator
$$\mathrm{\Delta }=\left(\begin{array}{cc}2\lambda & 2\mu \\ 2\mu & 2\nu \end{array}\right)\delta (z).$$
Thus, the absence of tachyons in our case is ensured by their absence in the RS model.
Next, there is no normalizable zero mode. Indeed, if $`m^2=0`$, then $`\chi _1=A_1\mathrm{e}^{2k|z|}`$, $`\widehat{\chi }_2=A_2=const`$ (other solutions to Eqs. (15) are exponentially growing). From Eqs. (16) one obtains that $`2\lambda A_1+2\mu A_2=0`$. Thus, $`A_20`$ and the zero mode is not normalizable with the scalar product (18).
We now turn to the modes with $`m^20`$. These modes belong to continuum spectrum. Introducing the conformal coordinate $`\zeta =\mathrm{e}^{k|z|}/k`$ one writes the general solution of the bulk equations (15) as a combination of Bessel functions,
$`\chi _1=A_1J_2(m\zeta )+B_1N_2(m\zeta ),`$ (19a)
$`\widehat{\chi }_2=m\zeta \left[A_2J_1(m\zeta )+B_2N_1(m\zeta )\right],`$ (19b)
with $`A_1,B_1,A_2,B_2`$ being complex numbers. For each value of the mass $`m`$ there are two eigenvectors $`\chi _m^{(r)}`$, $`r=1,2`$, of Eq. (17), which are normalized as follows,
$$(\chi _m^{(r)},\chi _m^{}^{(s)})=\delta ^{rs}\delta (mm^{}).$$
(20)
The normalization factor is determined by the asymptotics of the modes at large $`\zeta `$. Making use of the asymptotics of the expressions (19) one translates (20) into the relation between the coefficients $`A,B`$,
$$\frac{2k}{m}\left(\left(A_{1m}^{(r)}\right)^{}A_{1m}^{(s)}+\left(B_{1m}^{(r)}\right)^{}B_{1m}^{(s)}\right)+\frac{2m}{k}\left(\left(A_{2m}^{(r)}\right)^{}A_{2m}^{(s)}+\left(B_{2m}^{(r)}\right)^{}B_{2m}^{(s)}\right)=\delta ^{rs}.$$
(21)
Let us turn to the boundary equations (16). In terms of the coefficients $`A,B`$ they take the form,
$`A_1\left(J_1\left({\displaystyle \frac{m}{k}}\right){\displaystyle \frac{\lambda }{m}}J_2\left({\displaystyle \frac{m}{k}}\right)\right)`$ $`+B_1\left(N_1\left({\displaystyle \frac{m}{k}}\right){\displaystyle \frac{\lambda }{m}}N_2\left({\displaystyle \frac{m}{k}}\right)\right)`$
$`A_2{\displaystyle \frac{\mu }{k}}J_1\left({\displaystyle \frac{m}{k}}\right)B_2{\displaystyle \frac{\mu }{k}}N_1\left({\displaystyle \frac{m}{k}}\right)=0,`$ (22a)
$`A_1{\displaystyle \frac{\mu }{m}}J_2\left({\displaystyle \frac{m}{k}}\right)B_1{\displaystyle \frac{\mu }{m}}N_2\left({\displaystyle \frac{m}{k}}\right)`$ $`+A_2\left({\displaystyle \frac{m}{k}}J_0\left({\displaystyle \frac{m}{k}}\right){\displaystyle \frac{\nu }{k}}J_1\left({\displaystyle \frac{m}{k}}\right)\right)`$
$`+B_2\left({\displaystyle \frac{m}{k}}N_0\left({\displaystyle \frac{m}{k}}\right){\displaystyle \frac{\nu }{k}}N_1\left({\displaystyle \frac{m}{k}}\right)\right)=0.`$ (22b)
It is convenient to choose the first eigenvector of the operator (17) to satisfy the relation
$$\lambda \chi _1^{(1)}+\mu \widehat{\chi }_2^{(1)}|_{z=0}=0.$$
Then Eqs. (22) yield
$`A_1^{(1)}=C^{(1)}N_1\left({\displaystyle \frac{m}{k}}\right),`$ $`B_1^{(1)}=C^{(1)}J_1\left({\displaystyle \frac{m}{k}}\right),`$
$`A_2^{(1)}=C^{(1)}{\displaystyle \frac{\lambda k}{\mu m}}N_0\left({\displaystyle \frac{m}{k}}\right),`$ $`B_2^{(1)}=C^{(1)}{\displaystyle \frac{\lambda k}{\mu m}}J_0\left({\displaystyle \frac{m}{k}}\right),`$
where the normalization constant $`C^{(1)}`$ is fixed by Eq. (21). At $`mk`$ it has the form,
$$C^{(1)}=\pi \left(\frac{m}{2k}\right)^{3/2},$$
implying the following amplitude of the graviton field on the brane
$$\chi _1^{(1)}|_{z=0}=\left(\frac{m}{2k}\right)^{1/2}.$$
These modes are analogous to the continuum graviton spectrum in the RS case. They are completely delocalized and provide corrections to the four-dimensional Einstein gravity only at short distances $`r1/k`$.
The orthogonal mode is fixed by Eq. (21). In the general case the corresponding expressions are rather complicated. A considerable simplification occurs, if one assumes the hierarchy between the parameters, $`\lambda \mu \nu kM_5`$, where $`M_5=(16\pi G_5)^{1/3}`$ is the five-dimensional Planck mass. Since the physics at large distances is governed by light modes we are interested in the region of masses $`m^2\lambda k`$. One obtains,
$`A_1^{(2)}={\displaystyle \frac{2\mu }{\nu }}\left({\displaystyle \frac{m}{2k}}\right)^{3/2}\left(ln{\displaystyle \frac{m}{2k}}+𝑪\right),`$ $`B_1^{(2)}={\displaystyle \frac{\pi \mu }{\nu }}\left({\displaystyle \frac{m}{2k}}\right)^{3/2}{\displaystyle \frac{1}{1\frac{2\lambda }{\nu }ln\frac{m}{k}}},`$
$`A_2^{(2)}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{2k}{m}}\right)^{1/2},`$ $`B_2^{(2)}={\displaystyle \frac{\pi \lambda }{2\nu }}\left({\displaystyle \frac{2k}{m}}\right)^{1/2}{\displaystyle \frac{1}{1\frac{2\lambda }{\nu }ln\frac{m}{k}}},`$
where $`𝑪`$ is the Euler constant. An important feature of this mode is the divergence of the amplitudes of the fields on the brane in the small mass limit,
$`\chi _1^{(2)}|_{z=0}={\displaystyle \frac{\mu }{\nu }}\left({\displaystyle \frac{2k}{m}}\right)^{1/2}{\displaystyle \frac{1}{1\frac{2\lambda }{\nu }ln\frac{m}{k}}},`$
$`\widehat{\chi }_2^{(2)}|_{z=0}={\displaystyle \frac{\lambda }{\nu }}\left({\displaystyle \frac{2k}{m}}\right)^{1/2}{\displaystyle \frac{1}{1\frac{2\lambda }{\nu }ln\frac{m}{k}}}.`$
Thus, the modes $`\chi _m^{(2)}`$ give dominant contribution to the Green’s function of the operator (17) at large distances on the brane. Their collection comprises the quasilocalized four-dimensional graviton.
To make the last statement more precise, let us consider production of gravitational waves by an external periodic transverse traceless source on the brane,
$$T_{ij}^{ext}(x,z)=\delta (z)\mathrm{e}^{i\omega t}\frac{d^3𝐩}{(2\pi )^3}\mathrm{e}^{i\mathrm{𝐩𝐱}}d_{ij}^{(\alpha )}(𝐩)T_\alpha ^{ext}(𝐩).$$
The resulting gravitational field is given by the convolution of the source with the Green’s function
$$G(𝐱𝐱^{};\omega )=16\pi G_5𝑑\tau G_{11}(x,x^{};z=z^{}=0)\mathrm{e}^{i\omega \tau },$$
(23)
where $`\tau =tt^{}`$, and $`G_{11}(x,x^{};z,z^{})`$ is the upper left element of the retarded Green’s function of the operator (17),
$$G(x,x^{};z,z^{})=_0^{\mathrm{}}𝑑m\underset{s=1,2}{}\left(\begin{array}{cc}\chi _{1m}^{(s)}(z)\chi _{1m}^{(s)}(z^{})& \chi _{1m}^{(s)}(z)\widehat{\chi }_{2m}^{(s)}(z^{})\\ \widehat{\chi }_{2m}^{(s)}(z)\chi _{1m}^{(s)}(z^{})& \widehat{\chi }_{2m}^{(s)}(z)\widehat{\chi }_{2m}^{(s)}(z^{})\end{array}\right)\frac{d^4p}{(2\pi )^4}\frac{\mathrm{e}^{ip(xx^{})}}{m^2p^2iϵp_0}.$$
(24)
Substituting Eq. (24) into (23) one obtains,
$$G(𝐱𝐱^{};\omega )=\frac{4G_5}{r}_0^{\mathrm{}}𝑑m\underset{s=1,2}{}\left(\chi _{1m}^{(s)}(0)\right)^2\mathrm{e}^{ip_\omega r},$$
(25)
where $`r=|𝐱𝐱^{}|`$, $`p_\omega =\sqrt{\omega ^2m^2}`$ when $`m<\omega `$ and $`p_\omega =i\sqrt{m^2\omega ^2}`$ when $`m>\omega `$. We see that the gravitational field on the brane has the form of a superposition of massive four-dimensional modes. Only modes with $`m<\omega `$ are actually radiated, the other ones exponentially fall off from the source. Thus, as long as we are interested in the gravitational waves, we can integrate in (25) only up to $`m=\omega `$. This means, in particular, that at $`\omega k`$ one can neglect the contribution of the modes $`\chi _m^{(1)}`$ in Eq. (25). Let us introduce
$$m_c=k\mathrm{e}^{\nu /2\lambda },$$
(26)
and study the regime $`m_c\omega k`$. The Green’s function (25) takes the form
$$G(𝐱𝐱^{};\omega )=\frac{4G_N}{r}\mathrm{e}^{i\omega r}\frac{2\lambda }{\nu }_0^\omega \frac{dm}{m}\frac{\mathrm{e}^{i\frac{rm^2}{2\omega }}}{\left(1\frac{2\lambda }{\nu }ln\frac{m}{k}\right)^2},$$
(27)
where $`G_N=G_5k`$ is the four-dimensional Newton constant. At $`r\omega /m_c^2`$ the integral in Eq. (27) is saturated by $`mm_c`$ and we obtain the usual four-dimensional expression for the gravitational wave. We stress that this wave is a superposition of massive modes with masses of order $`m_c`$. In the opposite limit $`r\omega /m_c^2`$ the integral is damped by the rapidly oscillating exponent, so one obtains
$$G(𝐱𝐱^{};\omega )=\frac{4G_N}{r}\frac{\nu }{\lambda ln\frac{k^2r}{2\omega }}\mathrm{e}^{i\omega r}.$$
The gravitational wave gradually dissipates into the fifth dimension.
A few comments are in order. First, the graviton mass scale $`m_c`$ turns out to be independent of the parameters $`\varkappa _1`$, $`\varkappa _2`$ in the action (1):
$$\frac{\nu }{2\lambda }=\frac{1}{16\pi G_5v^2}.$$
This is not very surprising. Indeed, from the very beginning we have been interested in the modes with masses much smaller than the energy scale associated with $`\varkappa _1`$ and $`\varkappa _2`$. Then, the potential term for vector fields is effectively frozen at zero and its parameters do not affect dynamics. Rather, the condition that the potential must be zero, generates mixing between gravitons and vectors which leads to quasilocalization of graviton. The strength of mixing at low energies depends only on $`v`$, and so does the graviton mass.
Second, very small values of $`m_c`$ are generated without fine-tuning. For example, taking $`k,M_5`$ as large as the Planck mass and $`v(M_5/5)^{3/2}`$ one obtains $`m_c`$ as small as $`(10^{28}\mathrm{cm})^1`$, that is the inverse of the present horizon size of the Universe.
Finally, though the Lorentz symmetry is broken in the background (2), the Green’s function (24) for the tensor modes has Lorentz-invariant form.<sup>1</sup><sup>1</sup>1This allows to interpret the additional enhancement of the distance $`r\omega /m_c^2`$, at which the dissipation of the gravitational waves sets in, by the factor $`\omega /m_c`$ as compared to $`m_c^1`$, as a relativistic effect, see Ref. . The place where the Lorentz-breaking is essential is that the very notion of the tensor modes refers to the preferred reference frame.
## 4 Vector modes
The gauge conditions (3), (6) result in the following equations for the vector perturbations,
$`p_0\psi _1+𝐩^2\psi _2=0,`$ (28a)
$`p_0\psi _3+𝐩^2\psi _4+𝐩^2\psi _5=0,`$ (28b)
where we again omitted the polarization index $`\alpha `$. From Eqs. (4), (5) one obtains the bulk equations for the vector modes,
$`{\displaystyle \frac{1}{2}}\psi _I^{\prime \prime }2k^2\psi _I+{\displaystyle \frac{p^2}{2u^2}}\psi _I=0,I=1,2,`$ (29a)
$`\psi _I^{\prime \prime }+2k\psi _I^{}{\displaystyle \frac{p^2}{u^2}}\psi _I=0,I=3,4,5,z>0,`$ (29b)
The junction conditions on the brane read:
$`\psi _1^{}+2k\psi _1=0,`$ (30a)
$`\psi _2^{}+2k\psi _2=\lambda \left(\psi _2+{\displaystyle \frac{2\psi _4}{v}}+{\displaystyle \frac{i\psi _6}{v}}\right),`$ (30b)
$`\psi _3^{}=\psi _5^{}=0,`$ (30c)
$`\psi _4^{}=\varkappa ^2v^3\left(\psi _2+{\displaystyle \frac{2\psi _4}{v}}+{\displaystyle \frac{i\psi _6}{v}}\right).`$ (30d)
Equations (30a), (30b) together with the gauge fixing condition (28a) imply
$$\left(\psi _2+\frac{2\psi _4}{v}\right)|_{z=0}+\frac{i\psi _6}{v}=0.$$
(31)
Thus, the system (30) reduces to a set of homogeneous boundary conditions for the vector part of the gravitational perturbations and for the free vectors in the RS background. In other words, mixing between gravitational and vector fields has no effect in the sector of vector modes. The relations (28), (31) imply three independent vector modes, parametrized, say, by the functions $`\psi _2(z)`$, $`\psi _3(z)`$, $`\psi _4(z)`$.
Let us show the absence of physical zero modes in the vector sector. Assume that there is a normalizable mode with $`p^2=0`$. Then, from Eqs. (29) we find that $`\psi _2=C_2\mathrm{e}^{2k|z|}`$, and that $`\psi _3`$, $`\psi _4`$ do not depend on $`z`$. The latter fact means that the corresponding vector $`a_\mu ^a`$ is independent of $`z`$, see Eq. (14e). Consequently the vector field strength $`F_{\mu \nu }^a`$ is also independent of $`z`$, and, for the mode to be normalizable, $`F_{\mu \nu }^a`$ must be zero. This implies that the vector field $`a_\mu ^a`$ of the zero mode is pure gauge in the bulk. From Eqs. (14b), (14c) we find that the metric of the zero mode is also pure gauge,
$$h_{\mu \nu }(x,z)=(\xi _{\mu ,\nu }+\xi _{\nu ,\mu })\mathrm{e}^{2k|z|},\xi _\mu =iC_{2\alpha }e_\mu ^{(\alpha )}\mathrm{e}^{ipx}.$$
Thus, in the bulk, this mode can be removed by residual gauge transformation left after fixing the transverse-traceless gauge (3), (6). The only degree of freedom we are left with is the vector component $`\psi _6`$ of the field $`\beta ^a`$, see Eq. (14f). This field resides on the brane and describes explicit breaking of the $`[U(1)]^3`$ gauge symmetry by the brane action, cf. Eq. (12). But $`\psi _6`$ vanishes due to Eq. (31). Hence, all components of the mode vanish: there is no normalizable zero mode at all. One concludes that there is only the continuum spectrum in the vector sector with $`p^20`$. The corresponding modes are completely delocalized and have the same form as in the pure RS case.
It is worth mentioning that at $`p^2=0`$ the mode $`\psi _6`$ vanishes due to field equations and boundary conditions. This could be dangerous, if this mode became propagating in a non-trivial background. As we show in Appendix B, this is not the case, which implies also that $`\psi _6`$ remains non-dynamical to higher orders in classical perturbation theory.
## 5 Scalar sector
We now turn to the scalar sector. Having in mind the issue of vDVZ discontinuity, we will consider the gravitational field produced on the brane by an external energy-momentum tensor $`T_{\mu \nu }^{ext}`$ localized on the brane. We do not restrict our analysis to static sources and consider external energy-momentum tensor depending both on time and space. Passing to the 4-dimensional Fourier representation we decompose it as follows,
$$T_{00}^{ext}=t_1,T_{0i}^{ext}=p_it_2+\mathrm{},T_{ij}^{ext}=\left(\delta _{ij}\frac{p_ip_j}{𝐩^2}\right)t_3+p_ip_jt_4+\mathrm{},$$
where we wrote down explicitly only the scalar components, which are of interest to us. The energy-momentum conservation implies
$$p_0t_1+𝐩^2t_2=0,p_0t_2+𝐩^2t_4=0.$$
Below we will consider two gauge-invariant scalar potentials, $`\overline{\varphi }_3`$ and $`\overline{\mathrm{\Phi }}\overline{\varphi }_1+2p_0\overline{\varphi }_2+p_0^2\overline{\varphi }_4`$, of the metric produced by the external source on the brane. The quantities with the bar are defined via decomposition of the metric on the brane in accordance with (14a)–(14c). The potentials $`\overline{\varphi }_3`$, $`\overline{\mathrm{\Phi }}`$ are related to the standard Bardeen variables $`\mathrm{\Phi }`$, $`\mathrm{\Psi }`$ (in the notations of Ref. ) as follows, $`\overline{\varphi }_3=2\mathrm{\Psi }`$, $`\overline{\mathrm{\Phi }}=2\mathrm{\Phi }+\frac{2p_0^2}{𝐩^2}\mathrm{\Psi }`$. Using the relation (11) between the metric induced on the brane and the bulk metric we obtain,
$$\overline{\varphi }_3=\varphi _3|_{z=0}+2k\epsilon ,$$
(32)
$$\overline{\mathrm{\Phi }}=(\varphi _1+2p_0\varphi _2+p_0^2\varphi _4)|_{z=0}+\frac{2kp^2}{𝐩^2}\epsilon .$$
(33)
Our purpose is to compute the potentials $`\overline{\varphi }_3`$ and $`\overline{\mathrm{\Phi }}`$ and to compare them to the expressions obtained in the usual four-dimensional Einstein gravity.
In the bulk, the fields obey the following equations, see Eqs. (4), (5),
$`{\displaystyle \frac{1}{2}}\varphi _I^{\prime \prime }2k^2\varphi _I+{\displaystyle \frac{p^2}{2u^2}}\varphi _I=0,I=1,2,3,4,`$ (34a)
$`\varphi _I^{\prime \prime }+2k\varphi _I^{}{\displaystyle \frac{p^2}{u^2}}\varphi _I=0,I=5,6,7,z>0,`$ (34b)
The gauge conditions (3), (6) take the form,
$$\varphi _12\varphi _3𝐩^2\varphi _4=0,$$
(35a)
$$p_0\varphi _1+𝐩^2\varphi _2=0,$$
(35b)
$$p_0\varphi _2+𝐩^2\varphi _4=0,$$
(35c)
$$p_0\varphi _5+𝐩^2\varphi _7=0.$$
(35d)
As described in Appendix A, the junction conditions in the presence of the external energy-momentum tensor on the brane can be cast into the following form,
$`\varphi _1^{}+2k\varphi _1=2𝐩^2\epsilon +8\pi G_5t_1,`$ (36a)
$`\varphi _2^{}+2k\varphi _2=2p_0\epsilon +8\pi G_5t_2,`$ (36b)
$`\varphi _3^{}+2k\varphi _3=p^2\epsilon 4\pi G_5{\displaystyle \frac{p^2}{𝐩^2}}t_1,`$ (36c)
$`\varphi _4^{}+2k\varphi _4={\displaystyle \frac{2p_0^2}{𝐩^2}}\epsilon +8\pi G_5t_4,`$ (36d)
$`\varphi _5^{}=0,`$ (36e)
$`\varphi _6^{}={\displaystyle \frac{3\nu _1}{2\lambda _1}}vp^2\epsilon 2\pi G_5{\displaystyle \frac{\nu _1}{\lambda _1}}v(t_1+2t_3+𝐩^2t_4),`$ (36f)
$`\varphi _7^{}=0,`$ (36g)
$`p^2\epsilon =\lambda _1\rho \left(\varphi _3+{\displaystyle \frac{2}{v}}\varphi _6+2k\epsilon \right)+{\displaystyle \frac{4\pi G_5}{3}}(t_1+2t_3+𝐩^2t_4),`$ (36h)
where
$$\nu _1=2\varkappa _1^2v^2,\rho =\frac{\lambda _1+\lambda _2}{3\lambda _1+2\lambda _2}.$$
These equations should be supplemented by Eq. (51a) in Appendix A, which determines the longitudinal component $`\varphi _8`$ in terms of the other fields. Note that Eqs. (36c), (36f), (36h) form a closed system. Once the solution of this subsystem is found, the other equations are solved in a straightforward manner.
Before proceeding with the calculation of the gravitational field produced by the source, let us make sure that the scalar sector does not contain instabilities. So, one temporarily sets the external source equal to zero in Eqs. (36). First, we show that there is no tachyonic mode, which would correspond to negative or complex $`p^2`$. For such a mode, let us introduce $`w`$, such that $`k^2w^2=p^2`$ and $`Rew>0`$. Then, for a normalizable mode, one would have
$$\varphi _3=U_3K_2(w\mathrm{e}^{k|z|}),\varphi _6=U_6w\mathrm{e}^{k|z|}K_1(w\mathrm{e}^{k|z|}),$$
where $`K_1`$, $`K_2`$ are modified Bessel functions. From Eqs. (36c), (36f) we obtain
$$U_3=\frac{kw\epsilon }{K_1(w)},U_6=\frac{3\nu _1vk\epsilon }{2\lambda _1K_0(w)}.$$
Inserting these expressions into Eq. (36h) one obtains the following relation,
$$w+\frac{\lambda _1\rho }{k}\left[\frac{K_0(w)}{K_1(w)}+\frac{3\nu _1}{\lambda _1}\frac{K_1(w)}{K_0(w)}\right]=0.$$
(37)
Let us show that the real part of the expression in square brackets is positive in the right half-plane $`Rew>0`$. Indeed, on the boundary of this half-plane one has,
$`Re\left[{\displaystyle \frac{K_1(w)}{K_0(w)}}\right]|_{w=iy}={\displaystyle \frac{2}{\pi |y|(J_0^2(y)+N_0^2(y))}}>0,`$
$`Re\left[{\displaystyle \frac{K_1(w)}{K_0(w)}}\right]|_{|w|\mathrm{},\pi /2\mathrm{Arg}w\pi /2}=1>0.`$
Hence, $`Re\left[\frac{K_1(w)}{K_0(w)}\right]`$, being a harmonic function, is positive everywhere inside the half-plane $`Rew>0`$. This ensures the positivity of the real part of the inverse function, $`Re\left[\frac{K_0(w)}{K_1(w)}\right]>0`$, and thus of the whole expression in square brackets in Eq. (37). One concludes that Eq. (37) has no solutions in the right half-plane, implying the absence of tachyonic modes in the scalar sector.
Second, let us demonstrate that the model is also free of ghosts. The ghost mode, if any, must be localized on the brane. Indeed, the normalization of the modes of continuum spectrum is determined entirely by the bulk action, which by itself is free of ghosts. It is straightforward to see that the modes corresponding to strictly positive $`p^2`$ belong to continuum part of the spectrum. The only dangerous eigenvalue is $`p^2=0`$. But the corresponding mode is unphysical. Indeed, for this mode $`\varphi _2=U_2\mathrm{e}^{2k|z|}`$ and from Eq (36b) one obtains $`\epsilon =0`$. Substituting $`𝐩^2=p_0^2`$ into the gauge fixing conditions (35b), (35c) one finds $`\varphi _1=p_0U_2\mathrm{e}^{2k|z|}`$, $`\varphi _4=\frac{U_2}{p_0}\mathrm{e}^{2k|z|}`$. Then, Eq. (35a) yields $`\varphi _3=0`$. It follows from Eq. (34b) that the scalar component $`\varphi _7`$ of the vector fields does not depend on $`z`$. From Eq. (35d) we find $`\varphi _5=p_0\varphi _7`$. Finally, from Eq. (36h) we obtain $`\varphi _6=0`$. The surviving mode is pure gauge in the bulk, see Eqs. (14),
$`h_{\mu \nu }=(\xi _{\mu ,\nu }+\xi _{\nu ,\mu })\mathrm{e}^{2k|z|},\xi _\mu ={\displaystyle \frac{p_\mu }{2ip_0}}U_2\mathrm{e}^{ipx},`$
$`a_\mu ^a=_\mu \left(ip_a\varphi _7\mathrm{e}^{ipx}\right).`$
The bulk fields can be removed by a residual gauge transformation, leaving the scalar component $`\varphi _8`$ (see Eq. (14f)) of the fields $`\beta ^a`$ on the brane. But the latter is zero according to the field equation (51a). Thus, all the components of the zero mode vanish. There is only continuum spectrum of modes with $`p^2>0`$ in the scalar sector.
We show in Appendix B that the mode $`\varphi _8`$ remains non-dynamical in a non-trivial background as well. Hence, no new propagating degrees of freedom appear both in non-trivial backgrounds and in higher orders of classical perturbation theory.
Having established the absence of instabilities, we now return to the field produced by the source. Our strategy is to express the functions $`\varphi _3(z)`$, $`\varphi _6(z)`$ in terms of the external source and the brane bending $`\epsilon `$ using the bulk equations with the boundary conditions (36c), (36f), and to insert the result into Eq. (36h). This is most easily done by combining the bulk equations with the boundary conditions into the Schrödinger type operators and introducing the Green’s functions of these operators. The latter satisfy the following equations,
$`u^2G_3^{\prime \prime }(z;p)+4k^2u^2G_3(z;p)4k\delta (z)G_3(z;p)p^2G_3(z;p)=\delta (z),`$
$`u^2G_6^{\prime \prime }(z;p)+2ku^2G_6^{}(z;p)p^2G_6(z;p)=\delta (z).`$
Imposing the radiation (outgoing wave) boundary conditions at $`z\pm \mathrm{}`$ one obtains,
$$G_3(z;p)=\frac{H_2^{(1)}(p\mathrm{e}^{k|z|}/k)}{2pH_1^{(1)}(p/k)},$$
(38)
$$G_6(z;p)=\frac{\mathrm{e}^{k|z|}H_1^{(1)}(p\mathrm{e}^{k|z|}/k)}{2pH_0^{(1)}(p/k)},$$
(39)
where $`p=\sqrt{p^2}`$, $`Rep0`$. Making use of these Green’s functions we find,
$`\varphi _3(0)={\displaystyle \frac{pH_2^{(1)}(p/k)}{H_1^{(1)}(p/k)}}\left(\epsilon +4\pi G_5{\displaystyle \frac{t_1}{𝐩^2}}\right),`$ (40)
$`\varphi _6(0)={\displaystyle \frac{H_1^{(1)}(p/k)}{pH_0^{(1)}(p/k)}}\left({\displaystyle \frac{3\nu _1}{2\lambda _1}}vp^2\epsilon +2\pi G_5{\displaystyle \frac{\nu _1}{\lambda _1}}v(t_1+2t_3+𝐩^2t_4)\right).`$ (41)
Substitution of these expressions into Eq. (36h) yields,
$$\begin{array}{cc}& \epsilon \left[p^2+\lambda _1\rho \frac{pH_0^{(1)}(p/k)}{H_1^{(1)}(p/k)}3\nu _1\rho \frac{pH_1^{(1)}(p/k)}{H_0^{(1)}(p/k)}\right]\hfill \\ & \frac{4\pi G_5}{3}T\left[13\nu _1\rho \frac{H_1^{(1)}(p/k)}{pH_0^{(1)}(p/k)}\right]4\pi G_5\lambda _1\rho \frac{pH_2^{(1)}(p/k)}{H_1^{(1)}(p/k)}\frac{t_1}{𝐩^2}=0\hfill \end{array}$$
(42)
where we have used the relation
$$2k\frac{pH_2^{(1)}(p/k)}{H_1^{(1)}(p/k)}=\frac{pH_0^{(1)}(p/k)}{H_1^{(1)}(p/k)}$$
and introduced the notation $`T=t_12t_3𝐩^2t_4`$ for the trace of the energy-momentum tensor.
Further analysis depends on the value of the parameter $`\varkappa _1^2`$. Let us first consider the case $`\varkappa _1^2=0`$, which in Eq. (1) corresponds to the absence of the first term in the potential for vector fields.<sup>2</sup><sup>2</sup>2We do not know whether vanishing of the parameter $`\varkappa _1^2`$ can be ensured by any symmetry requirement. In this case $`\lambda _1=\nu _1=0`$, and Eq. (42) reduces to,
$$\epsilon =\frac{4\pi G_5}{3p^2}T.$$
(43)
The last expression coincides with the result obtained in the pure RS case, cf. . Substitution of (43) into Eqs. (40), (41) leads<sup>3</sup><sup>3</sup>3Note that though $`\lambda _1=\nu _1=0`$, the ratio $`\frac{\nu _1}{\lambda _1}=\frac{1}{8\pi G_5v^2}`$ in (41) is finite. to vanishing of the scalar part of the vector fields, $`\varphi _6=0`$, while the scalar components of the metric are the same as in the pure RS case. Hence, there are only short-distance corrections to the scalar part of the metric in the case $`\varkappa _1^2=0`$. At distances $`r1/k`$, i.e. at small momenta $`pk`$, one arrives at the same expressions as in the four-dimensional Einstein gravity,
$$\overline{\varphi }_3=\frac{8\pi G_N}{𝐩^2}t_1,\overline{\mathrm{\Phi }}=\frac{8\pi G_Np^2}{(𝐩^2)^2}t_1\frac{16\pi G_N}{𝐩^2}t_3,$$
(44)
where $`G_N=G_5k`$ is the four-dimensional Newton constant. Evidently, the gravitational field is free from the vDVZ discontinuity. In particular, the Newton law — the field of a static source — does not get modified at large distances in the case $`\varkappa _1^2=0`$.
Another possibility is that the parameter $`\varkappa _1^2`$ is large. In this case we are interested in the behavior of the fields at distances $`r1/\nu _1`$, i.e., at small values of the momentum, $`pk,\nu _1`$. We also assume as in Sec. 3 the hierarchy $`\lambda _1\nu _1`$. The dominant contributions to the expressions in the square brackets in Eq. (42) come from the terms proportional to $`H_1^{(1)}(p/k)/H_0^{(1)}(p/k)`$. It is worth noting that these terms come from the scalar component $`\varphi _6`$ of the vector fields in Eq. (36h). Keeping only the leading contributions at low momenta, one obtains from Eq. (42),
$$\epsilon =\frac{4\pi G_5}{3p^2}T+\frac{8\pi G_5}{3𝐩^2}t_1\frac{\lambda _1}{\nu _1}ln\frac{p}{k}.$$
(45)
This expression for the brane bending differs from the formula (43) in the pure RS case by the last term, which is small at moderate momenta. However, it becomes important and is logarithmically large relative to the first term in the far infrared. Making use of the expression (45) for the brane bending $`\epsilon `$ and the Green’s function (38) one determines the rest of the gravitational modes, $`\varphi _1`$, $`\varphi _2`$, $`\varphi _4`$. Finally, these expressions, being inserted into Eqs. (32), (33), yield the scalar components of the metric induced on the brane,
$`\overline{\varphi }_3=8\pi G_N{\displaystyle \frac{t_1}{𝐩^2}},`$ (46)
$`\overline{\mathrm{\Phi }}={\displaystyle \frac{8\pi G_Np^2}{(𝐩^2)^2}}t_1{\displaystyle \frac{16\pi G_N}{𝐩^2}}t_3+{\displaystyle \frac{16\pi G_Np^2}{(𝐩^2)^2}}t_1{\displaystyle \frac{\lambda _1}{\nu _1}}ln{\displaystyle \frac{p}{k}}.`$ (47)
The expression (46) coincides with the result of the usual four-dimensional Einstein gravity. Thus, the potential $`\varphi _3`$ does not get modified at large distances at all. As to the second potential, the first two terms in Eq. (47) also coincide with the result of conventional gravity. The last term provides the long-distance correction, which becomes important at the distance
$$r_c=\frac{1}{k}\mathrm{e}^{\nu _1/2\lambda _1}=\frac{1}{k}\mathrm{e}^{1/16\pi G_5v^2}.$$
This distance coincides with $`1/m_c`$, where $`m_c`$ is the graviton mass scale (26).
To make the physical consequences of the formula (47) more clear, let us calculate the gravitational field produced by a point-like static source. We take $`T_{00}=M\delta (𝐱)`$, $`T_{0i}=T_{ij}=0`$. Setting $`t_1=M`$, $`t_3=0`$, $`p^2=𝐩^2`$ in Eq. (47) and performing Fourier transform, one obtains,
$$\overline{\mathrm{\Phi }}(r)=\frac{2G_NM}{r}\left(1\frac{2\lambda _1}{\nu _1}lnkr\right).$$
(48)
while for the other gauge invariant potential we have,
$$\overline{\varphi }_3(r)=\frac{2G_NM}{r}.$$
(49)
These expressions are valid up to corrections at small distances $`r1/k,1/\nu _1`$. The potential (48) is the analog of the Newton potential in our model, it is responsible for gravitational interaction between nonrelativistic massive objects. Remarkably, a contribution to $`\overline{\mathrm{\Phi }}(r)`$ appears, which grows logarithmically with the distance as compared to the standard expression $`(2G_NM/r)`$. Even more strikingly, the sign of this contribution is opposite to that of the standard expression. As a result, the expression (48) describes gradual weakening of gravitational attraction with the distance; at large distances, $`r>1/m_c`$, attraction gets replaced by repulsion. Thus, our model provides an example of a ghost-free theory with antigravity at ultra-large distances.
At first sight it seems surprising that in the model with quasilocalized gravitons we obtain a contribution to the Newton potential which grows at large distances, as compared to the standard four-dimensional expression. Following the conventional line of reasoning one could infer that, as gravitons dissipate at large distances into the fifth dimension, the gravitational interaction should become weaker at large distances, as compared to the four-dimensional case. However, this line of reasoning is incorrect. In theories with Lorentz symmetry breaking, the potential produced by an external source is not directly related to the spectrum of propagating degrees of freedom. This point is illustrated by four-dimensional Lorentz violating massive electrodynamics, considered in Refs. . In that case the electric potential of static sources falls off as $`1/r`$ in spite of the fact that all propagating modes are massive.
The origin of the logarithmically enhanced antigravity in our model can be understood as follows. A point mass gives rise to perturbations of the vector fields which interact with matter via mixing with the metric. The gravitational field is produced by the total energy-momentum tensor composed of $`T_{\mu \nu }^{ext}`$ and the energy-momentum tensor $`T_{\mu \nu }^V`$ of the vectors, given by Eq. (9). The latter tensor falls off slowly from the localized external source, and dominates at large distances. Our results indicate that, insofar as the Newton potential is concerned, the vector fields mimic the effect of negative energy. This is possible because $`T_{\mu \nu }^V`$ violates the weak energy condition.
The gravitational potentials (46), (47) are free from the vDVZ discontinuity. Indeed, in the limit of vanishing graviton mass scale $`m_c`$, which corresponds to $`\lambda _1/\nu _10`$, one recovers the usual four dimensional expressions (44). However, from the phenomenological point of view the difference between the two gravitational potentials (48) and (49) results in a severe phenomenological constraint on the value of the parameter $`\lambda _1/\nu _1`$. Measurements of the light deflection by the gravitational field of the Sun require $`\lambda _1/\nu _1\stackrel{<}{}10^5`$. This value is not unnaturally small: it corresponds to $`v(0.027M_5)^{3/2}`$. However, it pushes the range $`1/m_c`$, where the antigravity sets in, far beyond the present horizon size of the Universe. The mass scale $`m_c`$ in this case is very small and it is not clear whether it can significantly affect the physics within the present horizon. We discuss possible ways of avoiding this phenomenological constraint in the concluding section.
## 6 Discussion
In this paper we presented a Lorentz-violating brane-world model, where gravitons are not completely localized, but are rather quasilocalized on the brane. In other words, the four-dimensional graviton is a collection of KK modes from the continuum spectrum. The characteristic mass $`m_c`$ of these modes is exponentially small when expressed in terms of the parameters of the Lagrangian. We demonstrated that the model is free from tachyonic and ghost-like instabilities. We calculated the metrics produced by an external energy-momentum tensor on the brane, and found that the model does not suffer from the vDVZ discontinuity.
The key observation behind the model is that while the gravitational perturbations in the Randall–Sundrum setup contain a zero mode localized on the brane, the bulk vector fields are completely delocalized. Thus, mixing between the vectors and the would-be zero gravitational mode forces the latter to dissipate into the fifth dimension. This mixing between tensors and vectors becomes possible when the Lorentz symmetry is spontaneously broken by non-zero VEVs of the vector fields. These considerations are rather generic. We believe them to be applicable to a broad class of generalizations of the Randal–Sundrum model with spontaneous Lorentz symmetry breaking by bulk fields. In particular, the vector fields can be replaced by form-fields of higher degrees, which are also completely delocalized in the Randall–Sundrum background.
For general form of the potential term for the vector fields on the brane the model exhibits antigravity at ultra-large distances. Namely, the structure of the Newton potential is such that gravitational attraction between nonrelativistic massive objects gradually weakens as the distance increases, and gets replaced by repulsion at $`r>1/m_c`$. While this property is theoretically appealing, it puts severe phenomenological constraints on the model. The reason is that the second Bardeen potential, characterizing spatial part of the metric of a static source, does not get modified at large distances and thus differs from the Newton potential. Observations of light deflection by the Sun constrain the relative difference between the two gravitational potentials at the level of $`10^4`$. This translates into the constraint $`m_c\stackrel{<}{}M_{Pl}\mathrm{exp}(10^5)`$. This value is negligible compared to the present Hubble parameter of the Universe. It is doubtful whether the model with so tiny graviton mass scale can lead to any interesting phenomenology, in particular, to accelerating expansion of the Universe at the present epoch.
One possibility to avoid this phenomenological problem is to fine-tune the parameter $`\varkappa _1`$ of the vector potential in the action (1) to zero. Then, both Bardeen potentials have at large distances the same form as in the four-dimensional linearized Einstein gravity, while graviton is still quasilocalized and gravity waves dissipate into extra dimensions. The graviton mass scale $`m_c`$ in this case can be comparable to the Hubble parameter or even larger. A drawback of this approach is fine-tuning of $`\varkappa _1`$. We are not aware whether vanishing of this parameter can be imposed by any symmetry.
Another way out is a generalization of the model considered in this paper. Let us sketch a particular example. One introduces a scalar field with dilaton-like coupling to the vector fields on the brane. Namely, one replaces the combination $`\overline{g}^{\mu \nu }A_\mu ^aA_\nu ^b`$ in the brane action in Eq. (1) by $`\overline{g}^{\mu \nu }\mathrm{e}^{\alpha \phi }A_\mu ^aA_\nu ^b`$. The modified model allows for the Lorentz-violating background (2) with $`\phi =0`$. The dilaton does not affect tensor and vector sectors of the linearized perturbation above this background, so the modified model also incorporates quasilocalized gravitons. On the other hand, relative difference between Bardeen potentials depends on the dilaton coupling, and is essentially proportional to $`1/\alpha ^2`$ when $`\alpha `$ is large. Thus, the constraint from light deflection is satisfied once $`\alpha \stackrel{>}{}100`$. We will report more on phenomenology of our model and its generalizations elsewhere.
Two other important issues which we leave for future investigations are the quantum structure of the theory, and cosmology in the model (1) and its generalizations. Of special interest is at what scale the strong coupling at the quantum level sets in, and what kind of late-time cosmological evolution is obtained in these models, in particular, whether they can account for the cosmic acceleration at the present epoch.
### Acknowledgements
We are indebted to S. Dubovsky, M. Libanov and V. Rubakov for many fruitful discussions and helpful suggestions. We are grateful to S. Demidov, D. Krotov, D. Levkov, Kh. Nirov, E. Nugaev for their encouraging interest. This work has been supported in part by the grant of the President of the Russian Federation NS-2184.2003.2, and the Russian Foundation for Basic Research grant 05-02-17363. The work of D.G. has been supported in part by the RFBR grant 04-02-17448, INTAS grant 03-51-5112, the grant of the Russian Science Support Foundation and by the fellowship of the ”Dynasty” foundation (awarded by the Scientific board of ICFPM).
## Appendix A Junction conditions in the scalar sector
In this Appendix we consider the junction conditions for the scalar modes in the presence of an external energy-momentum tensor on the brane. The external source should be added to the r.h.s. of Eqs. (13) according to Eq. (7). As a result one obtains the following junction conditions on the brane,
$`\varphi _1^{}+2k\varphi _12p_0^2\epsilon =`$ $`{\displaystyle \frac{\lambda _1}{3}}\left(2\varphi _3+𝐩^2\varphi _4+{\displaystyle \frac{4}{v}}\varphi _6+{\displaystyle \frac{2𝐩^2}{v}}\varphi _7{\displaystyle \frac{𝐩^2}{k}}\epsilon +6k\epsilon +{\displaystyle \frac{2i𝐩^2}{v}}\varphi _8\right)`$
$`+8\pi G_5\left({\displaystyle \frac{2}{3}}t_1+{\displaystyle \frac{2}{3}}t_3+{\displaystyle \frac{𝐩^2}{3}}t_4\right),`$ (50a)
$`\varphi _2^{}+2k\varphi _2+2p_0\epsilon =`$ $`8\pi G_5t_2,`$ (50b)
$`\varphi _3^{}+2k\varphi _3=\lambda ({\displaystyle \frac{1}{3}}\varphi _3`$ $`{\displaystyle \frac{𝐩^2}{3}}\varphi _4+{\displaystyle \frac{2}{3v}}\varphi _6{\displaystyle \frac{2𝐩^2}{3v}}\varphi _7+{\displaystyle \frac{𝐩^2}{3k}}\epsilon {\displaystyle \frac{2i𝐩^2}{3v}}\varphi _8),`$ (50c)
$`\varphi _4^{}+2k\varphi _42\epsilon =\lambda `$ $`\left({\displaystyle \frac{2}{3𝐩^2}}\varphi _3+{\displaystyle \frac{2}{3}}\varphi _4{\displaystyle \frac{4}{3𝐩^2v}}\varphi _6+{\displaystyle \frac{4}{3v}}\varphi _7{\displaystyle \frac{2}{3k}}\epsilon +{\displaystyle \frac{4i}{3v}}\varphi _8\right)`$
$`+8\pi G_5\left({\displaystyle \frac{1}{3𝐩^2}}t_1{\displaystyle \frac{2}{3𝐩^2}}t_3+{\displaystyle \frac{2}{3}}t_4\right),`$ (50d)
$`\varphi _5^{}=`$ $`0,`$ (50e)
$`\varphi _6^{}=`$ $`\varkappa _1^2v^3\left(\varphi _3+{\displaystyle \frac{2}{v}}\varphi _6+2k\epsilon \right)+\varkappa _2^2v^3\left({\displaystyle \frac{1}{3}}\varphi _3{\displaystyle \frac{𝐩^2}{3}}\varphi _4+{\displaystyle \frac{2}{3v}}\varphi _6{\displaystyle \frac{2𝐩^2}{3v}}\varphi _7+{\displaystyle \frac{𝐩^2}{3k}}\epsilon {\displaystyle \frac{2i𝐩^2}{3v}}\varphi _8\right),`$ (50f)
$`\varphi _7^{}=`$ $`\varkappa _1^2v^3\left(\varphi _4+{\displaystyle \frac{2}{v}}\varphi _7{\displaystyle \frac{1}{k}}\epsilon +{\displaystyle \frac{2k}{𝐩^2}}\epsilon +{\displaystyle \frac{2i}{v}}\varphi _8\right)`$
$`+\varkappa _2^2v^3\left({\displaystyle \frac{2}{3𝐩^2}}\varphi _3+{\displaystyle \frac{2}{3}}\varphi _4{\displaystyle \frac{4}{3𝐩^2v}}\varphi _6+{\displaystyle \frac{4}{3v}}\varphi _7{\displaystyle \frac{2}{3k}}\epsilon +{\displaystyle \frac{4i}{3v}}\varphi _8\right).`$ (50g)
Combining these equations with the gauge fixing conditions (35), we obtain
$`\lambda _1\{(𝐩^2`$ $`\varphi _4{\displaystyle \frac{2𝐩^2}{v}}\varphi _7)|_{z=0}+{\displaystyle \frac{𝐩^2}{k}}\epsilon 2k\epsilon {\displaystyle \frac{2i𝐩^2}{v}}\varphi _8\}`$
$`+\lambda _2\left\{\left({\displaystyle \frac{2}{3}}\varphi _3{\displaystyle \frac{2𝐩^2}{3}}\varphi _4+{\displaystyle \frac{4}{3v}}\varphi _6{\displaystyle \frac{4𝐩^2}{3v}}\varphi _7\right)|_{z=0}+{\displaystyle \frac{2𝐩^2}{3k}}\epsilon {\displaystyle \frac{4i𝐩^2}{3v}}\varphi _8\right\}=0,`$ (51a)
$`2p^2\epsilon =`$ $`{\displaystyle \frac{\lambda _1}{3}}\left\{\left(2\varphi _3+𝐩^2\varphi _4+{\displaystyle \frac{4}{v}}\varphi _6+{\displaystyle \frac{2𝐩^2}{v}}\varphi _7\right)|_{z=0}{\displaystyle \frac{𝐩^2}{k}}\epsilon +6k\epsilon +{\displaystyle \frac{2i𝐩^2}{v}}\varphi _8\right\}`$
$`+8\pi G_5\left({\displaystyle \frac{1}{3}}t_1+{\displaystyle \frac{2}{3}}t_3+{\displaystyle \frac{𝐩^2}{3}}t_4\right).`$ (51b)
The first of these equations plays the role similar to that of Eq. (31) in the vector sector: it determines the field $`\varphi _8`$ in terms of the other fields. Substitution of Eqs. (51) back into Eqs. (50) yields the system (36) considered in the main text.
## Appendix B Absence of ghosts above a non-trivial background
We saw in the main text that the modes $`\varphi _8`$ and $`\psi _{6\alpha }`$, accounting for the explicit $`[U(1)]^3`$ gauge symmetry breaking on the brane, vanish the the linear order in perturbation theory above the background (2). These vanishing modes are generically dangerous. In higher orders of classical perturbation theory (i.e. above non-trivial backgrounds) kinetic terms for these modes may arise, rendering them propagating. Depending on the sign of kinetic term some of these modes may become ghosts. To study this issue we investigate the ultraviolet behavior of the perturbations of the vector fields above a non-trivial vector background close to the static background $`A_i^a=v\delta _i^a`$. For the sake of simplicity we neglect gravity perturbations. This can be done consistently by taking the limit $`G_50`$; we believe that switching on gravity does not spoil our results. The non-trivial background is chosen to be locally space-like $`A_\mu ^a=(0,𝐀^a)`$. In this Appendix we show that the equations of motion for the linear perturbations above the non-trivial vector background imply vanishing of the dangerous modes, in complete analogy to the situation above the static background (2). This demonstrates the ultraviolet stability of the model (1) to higher orders in classical perturbation theory.
Let us consider the brane part of the quadratic action for the fluctuations $`a_\mu ^a`$ of the vector fields above the non-trivial vector background,
$`S^{(2)}=`$ $`{\displaystyle }d^4x[2(\varkappa _1^2+\varkappa _2^2)(A^{\mu b}A_\mu ^ca^{\nu b}a_\nu ^c+A^{\mu b}a_\mu ^ca^{\nu b}A_\nu ^c+A^{\mu b}a_\mu ^ca^{\nu c}A_\nu ^b)`$
$`2\varkappa _1^2v^2a^{\nu b}a_\nu ^b{\displaystyle \frac{2\varkappa _2^2}{3}}(A^{\mu b}A_\mu ^ba^{\nu c}a_\nu ^c+2A^{\mu b}a_\mu ^bA^{\mu c}a_\mu ^c)],`$
where summation over internal indices $`b,c=1,\mathrm{},3`$ is assumed, and the metric on the brane is taken to be flat. The bulk part of the action is irrelevant for us as we are interested in longitudinal modes $`a_\mu ^a=_\mu \beta ^a`$ which are pure gauge in the bulk.
Let us consider equations for the spatial components of the fluctuations, $`a_i^a=_i\beta ^a`$. It is convenient to represent them in the matrix form introducing $`3\times 3`$ matrices $`\widehat{a}=a_i^a`$ and $`\widehat{A}=A_i^a`$. By a suitable transformation, the background matrix $`\widehat{A}`$ can be made symmetric, $`\widehat{A}=\widehat{A}^T`$. Then the equation of motion for the fluctuations reads,
$$4\left(\varkappa _1^2+\varkappa _2^2\right)\left(\widehat{A}^2\widehat{a}+\widehat{A}\widehat{a}^T\widehat{A}+\widehat{a}\widehat{A}^2\right)+4\varkappa _1^2v^2\widehat{a}+\frac{4\varkappa _2^2}{3}\left(2\mathrm{T}\mathrm{r}[\widehat{a}\widehat{A}]\widehat{A}+\mathrm{Tr}[\widehat{A}^2]\widehat{a}\right)=0.$$
(52)
The symmetric part of the fluctuations, $`s_i^a=_i\beta ^a+_a\beta ^i`$, obeys the homogeneous equation,
$$𝒜[\widehat{s}]=0,$$
where the linear operator $`𝒜`$ can be read from Eq. (52). For the static background (2b) the operator $`𝒜`$ is non-degenerate. Thus, it is non-degenerate for non-trivial backgrounds, which are close enough to (2b). One concludes that the matrix $`\widehat{s}`$ vanishes,
$$_i\beta ^b+_b\beta ^i=0.$$
(53)
Differentiating Eq. (53) twice yields the Laplace equation for the divergence, $`\mathrm{\Delta }_a\beta ^a=0`$. Imposing vanishing boundary conditions at spatial infinity we obtain $`_a\beta ^a=0`$. Then, differentiating Eq. (53) once, we find that all the components $`\beta ^a`$ vanish.
Thus the longitudinal modes $`\varphi _8`$ and $`\psi _{6\alpha }`$ do not become propagating in non-trivial background. This result implies the absence of ghosts and/or tachyons in the model (1) to higher orders in classical perturbation theory.
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# Rearrangement inequalities for functionals with monotone integrands
## 1 Introduction
The systematic study of rearrangements begins with the final chapter of “Inequalities” by Hardy, Littlewood, and Pólya . Two inequalities are discussed there at length, the Hardy-Littlewood inequality (Theorems 368-370 and 378 of )
$$_{}u(x)v(x)𝑑x_{}u^{}(x)v^{}(x)𝑑x,$$
(1.1)
and the Riesz rearrangement inequality (, Theorem 370 of )
$$_{}_{}u(x)v(x^{})w(xx^{})𝑑x𝑑x^{}_{}_{}u^{}(x)v^{}(x^{})w^{}(xx^{})𝑑x𝑑x^{}.$$
(1.2)
Here, $`u`$, $`v`$, and $`w`$ are nonnegative measurable functions that vanish at infinity, and $`u^{}`$, $`v^{}`$, and $`w^{}`$ are their symmetric decreasing rearrangements.
The Hardy-Littlewood inequality is a very basic inequality that holds, with suitably defined rearrangements, on arbitrary measure spaces . Its main implication is that rearrangement decreases $`L^2`$-distances . In contrast, the Riesz rearrangement inequality is specific to $``$ and to $`^n`$, where it is closely related with the Brunn-Minkowski inequality of convex geometry. The generalization of Eq. (1.2) from $``$ to $`^n`$ is due to Sobolev , and the inequality is also known as the Riesz-Sobolev inequality. For many applications, the third function in Eq. (1.2) is already radially decreasing, i.e., $`w(xx^{})=K(|xx^{}|)`$ with some nonnegative nonincreasing function $`K`$, such as the heat kernel or the Coulomb kernel (Theorems 371-373 and 380 of ). This special case of the inequality also holds on the standard spheres and hyperbolic spaces , and it still contains the isoperimetric inequality as a limit.
It is a natural question whether these inequalities carry over to more general integral functionals. Under what conditions on $`F`$ do the extended Hardy-Littlewood inequality
$$F(u_1(x),\mathrm{},u_m(x))𝑑xF(u_1^{}(x),\mathrm{},u_m^{}(x))𝑑x$$
(1.3)
and the extended Riesz inequality
$$\begin{array}{c}\mathrm{}F(u_1(x_1),\mathrm{},u_m(x_m))\underset{i<j}{}K_{ij}(d(x_i,x_j))dx_1\mathrm{}dx_m\hfill \\ \mathrm{}F(u_1^{}(x_1),\mathrm{},u_m^{}(x_m))\underset{i<j}{}K_{ij}(d(x_i,x_j))dx_1\mathrm{}dx_m\hfill \end{array}$$
(1.4)
hold for all choices of $`u_1,\mathrm{},u_m`$? In Eq. (1.4) the $`K_{ij}`$ are given nonnegative nonincreasing functions on $`_+`$, and $`d(x,y)`$ denotes the distance between $`x`$ and $`y`$. Eq. (1.3) can be recovered from Eq. (1.4) by choosing $`K_{ij}`$ as a Dirac sequence and passing to the limit. Note that Eq. (1.4) contains only the case of Eq. (1.2) where the third function is a symmetric decreasing kernel. A larger class of integral kernels $`K(x_1,\mathrm{},x_m)`$ was considered in . The full generalization of Riesz’ inequality to products of more than three functions was found by Brascamp-Lieb-Luttinger ; again, one may ask to what class of integrands the Brascamp-Lieb-Luttinger inequality naturally extends.
The main condition on $`F`$ was identified by Lorentz as the second-order monotonicity property
$$F(𝐲+h𝐞_i+k𝐞_j)+F(𝐲)F(𝐲+h𝐞_i)+F(𝐲+k𝐞_j)(ij,h,k>0),$$
(1.5)
where $`𝐲=(y_1,\mathrm{},y_m)`$, and $`𝐞_i`$ denotes the $`i`$-th standard basis vector in $`^m`$. Functions satisfying Eq. (1.5) are called supermodular or $`2`$-increasing in Economics. A smooth function is supermodular, if all its mixed second partial derivatives are nonnegative. Eqs. (1.3) and (1.4) were proved for continuous supermodular integrands depending on $`m=2`$ functions by Crowe-Zweibel-Rosenbloom and Almgren-Lieb (Theorem 2.2 of ). For $`m>2`$, Eq. (1.3) is due to Brock , and Eq. (1.4) is a recent result of Draghici . The purpose of this paper is to dispense with the continuity assumptions on $`F`$ in the theorems of Brock and Draghici, and to characterize the equality cases in some relevant situations. This continues prior work of the second author \[16-19\].
Acknowledgments. We thank Friedemann Brock, Cristina Draghici, and Loren Pitt for useful discussions, and especially Al Baernstein for drawing our attention to Sklar’s theorem. A.B. was partially funded by grants from the National Science Foundation (NSF), the National Sciences and Engineering Research Council of Canada (NSERC), and a University of Toronto Connaught award. H.H. was supported by the Fonds National Suisse de la Recherche Scientifique (FNS).
## 2 Statement of the results
Let $`𝕏`$ denote either the Euclidean space $`^n`$, the sphere $`\mathrm{SS}^n`$, or the hyperbolic space $`^n`$, equipped with the standard distance function $`d(,)`$ and the uniform volume measure $`\lambda `$. Choose a distinguished point $`x^{}𝕏`$ to serve as the origin or the north pole. Consider a nonnegative measurable function $`u`$ on $`𝕏`$. When $`𝕏=^n`$ or $`^n`$, we require $`u`$ to vanish at infinity in the sense that all its positive level sets $`\{x𝕏:u(x)>t\}`$ have finite measure; when $`𝕏=\mathrm{SS}^n`$ this requirement is void. By definition, the symmetric decreasing rearrangement $`u^{}`$ of $`u`$ is the unique upper semicontinuous, nonincreasing function of $`d(x,x^{})`$ that is equimeasurable with $`u`$. Explicitly, if
$$\rho (t)=\lambda \left(\{x𝕏:u(x)>t\}\right)$$
is the distribution function of $`u`$, and $`B_r`$ denotes the open ball of radius $`r`$ centered at $`x^{}`$, then
$$u^{}(x):=sup\{t0:\rho (t)\lambda \left(B_{d(x,x^{})}\right)\}.$$
###### Theorem 1
(Extended Hardy-Littlewood inequality.) Eq. (1.3) holds for all nonnegative measurable functions $`u_1,\mathrm{},u_m`$ that vanish at infinity on $`𝕏=^n,\mathrm{SS}^n`$, or $`^n`$, provided that the integrand $`F`$ is a supermodular Borel measurable function on the closed positive cone $`_+^m`$ with $`F(\mathrm{𝟎})=0`$, and that its negative part satisfies
$$_𝕏F_{}\left(u_i(x)𝐞_i\right)𝑑x<\mathrm{}$$
(2.1)
for $`i=1,\mathrm{},m`$.
Suppose Eq. (1.3) holds with equality, and the integrals are finite. If $`F`$ satisfies Eq. (1.5) with strict inequality for some $`ij`$, all $`𝐲_+^m`$ and all $`h,k>0`$, then
$$\left(u_i(x)u_i(x^{})\right)\left(u_j(x)u_j(x^{})\right)0$$
for almost all $`x,x^{}𝕏`$; in particular, if $`u_i=u_i^{}`$ is strictly radially decreasing, then $`u_j=u_j^{}`$.
The Borel measurability of $`F`$ and the integrability assumption in Eq. (2.1) ensure that the integrals in Eq. (1.3) are well-defined, though they may take the value $`+\mathrm{}`$.
The left hand side of Eq. (1.3) is invariant under volume-preserving diffeomorphisms of $`𝕏`$. More generally, if $`(\mathrm{\Omega },\mu )`$ and $`(\mathrm{\Omega }^{},\mu ^{})`$ are measure spaces and $`\tau :\mathrm{\Omega }\mathrm{\Omega }^{}`$ pushes $`\mu `$ forward to $`\mu ^{}`$ in the sense that $`\mu ^{}(A)=\mu (\tau ^1(A))`$ for all $`\mu ^{}`$-measurable subsets $`A\mathrm{\Omega }^{}`$, then
$$_\mathrm{\Omega }F(u_1(\omega ),\mathrm{},u_m(\omega ))𝑑\mu (\omega )=_\mathrm{\Omega }^{}F(u_1\tau (\omega ^{}),\mathrm{},u_m\tau (\omega ^{}))𝑑\mu ^{}(\omega ^{}).$$
The right hand side of Eq. (1.3) can also be expressed in an invariant form. Define the nonincreasing rearrangement $`u^\mathrm{\#}`$ of $`u`$ as the unique nonincreasing upper semicontinuous function on $`_+`$ that is equimeasurable with $`u`$,
$$u^\mathrm{\#}(\xi ):=sup\{t0:\rho (t)\xi \}.$$
By construction, $`(u\tau )^\mathrm{\#}=u^\mathrm{\#}`$ for any map $`\tau :\mathrm{\Omega }\mathrm{\Omega }^{}`$ that pushes $`\mu `$ forward to $`\mu ^{}`$. On $`𝕏=^n`$, $`\mathrm{SS}^n`$ and $`^n`$, the nonincreasing rearrangement is related with the symmetric decreasing rearrangement by $`u^{}(x)=u^\mathrm{\#}\left(\lambda \left(B_{d(x,x^{})}\right)\right)`$. Theorem 1 implies that
$$_\mathrm{\Omega }F(u_1(\omega ),\mathrm{},u_m(\omega ))𝑑\mu (\omega )_0^{\mu (\mathrm{\Omega })}F(u_1^\mathrm{\#}(\xi ),\mathrm{},u_m^\mathrm{\#}(\xi ))𝑑\xi $$
(2.2)
for all nonnegative measurable functions $`u_1,\mathrm{},u_m`$ on $`\mathrm{\Omega }`$ that vanish at infinity.
When $`\mu `$ is a probability measure, Eq. (2.2) says that the expected value of $`F(Y_1,\mathrm{},Y_m)`$ is maximized among all random variables $`Y_1,\mathrm{},Y_m`$ with given marginal distributions by the perfectly correlated random variables $`Y_1^\mathrm{\#},\mathrm{},Y_m^\mathrm{\#}`$. The joint distribution of the maximizer is uniquely determined, if $`Y_i`$ is continuously distributed for some $`i`$ and Eq. (1.5) is strict for all $`ji`$. In this formulation, the invariance under measure-preserving transformations is evident, since the expected value depends only on the joint distribution of $`Y_1,\mathrm{},Y_m`$. The assumption that $`F`$ is supermodular signifies that each of the random variables enhances the contribution of the others.
###### Theorem 2
(Extended Riesz inequality.) Eq. (1.4) holds for all nonnegative measurable functions $`u_1,\mathrm{},u_m`$ on $`𝕏=^n,\mathrm{SS}^n`$, or $`^n`$ that vanish at infinity, provided that $`F`$ is a supermodular Borel measurable function on $`_+^m`$ with $`F(\mathrm{𝟎})=0`$, each $`K_{ij}`$ is nonincreasing and nonnegative, and the negative part of $`F`$ satisfies
$$_𝕏\mathrm{}_𝕏F_{}\left(u_{\mathrm{}}(x_{\mathrm{}})𝐞_{\mathrm{}}\right)\underset{i<j}{}K_{ij}(d(x_i,x_j))dx_1\mathrm{}dx_m<\mathrm{}$$
(2.3)
for $`\mathrm{}=1,\mathrm{},m`$.
Suppose Eq. (1.4) holds with equality. Assume additionally that the integrals are finite, and that $`K_{ij}(t)>0`$ for all $`i<j`$ and all $`t<\mathrm{diam}𝕏`$. Let $`\mathrm{\Gamma }_0`$ be the graph on the vertex set $`\{1,\mathrm{},m\}`$ which has an edge between $`i`$ and $`j`$ whenever $`K_{ij}`$ is a strictly decreasing function, and let $`ij`$ be from the same component of $`\mathrm{\Gamma }_0`$. If Eq. (1.5) is strict for all $`𝐲_+^m`$ and all $`h,k>0`$, and if $`u_i`$ and $`u_j`$ are non-constant, then $`u_i=u_i^{}\tau `$ and $`u_j=u_j^{}\tau `$ for some translation $`\tau `$ on $`𝕏`$.
## 3 Related work
There are several proofs of the extended Hardy-Littlewood inequality in the literature. For continuous integrands, Lorentz showed by discretization and elementary manipulations of the $`u_i`$ that Eq. (2.2) holds for all measurable functions $`u_1,\mathrm{},u_m`$ on $`\mathrm{\Omega }=(0,1)`$ if and only if $`F`$ is supermodular . By the invariance under measure-preserving transformations, this implies Eq. (1.3), as well as Eq. (2.2) for arbitrary finite measure spaces $`\mathrm{\Omega }`$. However, Lorentz’ paper has had little impact on subsequent developments.
More than thirty years later, Crowe-Zweibel-Rosenbloom proved Eq. (1.3) for $`m=2`$ on $`𝕏=^n`$ . They expressed a given continuous supermodular function $`F`$ on $`_+^2`$ that vanishes on the boundary as the distribution function of a Borel measure $`\mu _F`$,
$$F(y_1,y_2)=\mu _F\left([0,y_1)\times [0,y_2)\right).$$
layer-cake representation
$$F(u_1(x),u_2(x))𝑑x=_{_+^2}\left\{\mathrm{𝟏}_{u_1(x)>y_1}\mathrm{𝟏}_{u_2(x)>y_2}𝑑x\right\}𝑑\mu _F(y_1,y_2),$$
(3.1)
which reduces Eq. (1.3) to the case where $`F`$ is a product of characteristic functions (see Theorem 1.13 in ). Another reduction to products was proposed by Tahraoui . The regularity and boundary conditions on $`F`$ were relaxed by Hajaiej-Stuart, who assumed it to be supermodular, of Carathéodory type (i.e., Borel measurable in the first, continuous in the second variable), and to satisfy some growth and integrability restrictions . Equality statements for their results were obtained by Hajaiej . Using a slightly different layer-cake decomposition, Van Schaftingen-Willem recently established Eq. (2.2) for $`m=2`$, under additional assumptions on $`F`$, for any equimeasurable rearrangement that preserves inclusions .
The drawback of the layer-cake representation is that for $`m>2`$ it requires an $`m`$-th order monotonicity condition on the integrand, which amounts for smooth $`F`$ to the nonnegativity of all (non-repeating) mixed partial derivatives . Brock proved Eq. (1.3) under the much weaker assumption that $`F`$ is continuous and supermodular .
Carlier viewed maximizing the left hand side of Eq. (2.2) for a given right hand side as an optimal transportation problem where the distribution functions of $`u_1,\mathrm{},u_m`$ define mass distributions $`\mu _i`$ on $``$, the joint distribution defines a transportation plan, and the functional represents the cost after multiplying by a minus sign . He showed that the functional achieves its maximum (i.e., the cost is minimized) when the joint distribution is concentrated on a curve in $`^m`$ that is nondecreasing in all coordinate directions, and obtained Eq. (2.2) as a corollary. His proof takes advantage of the dual problem of minimizing
$$\underset{i=1}{\overset{m}{}}_{}f_i(y)𝑑\mu _i(y)$$
over $`f_1,\mathrm{},f_m`$, subject to the constraint that $`f_i(y_i)F(y_1,\mathrm{},y_m)`$ for all $`y_1,\mathrm{},y_m`$.
Theorem 1 can be applied to some integrands that depend explicitly on the radial variable . If $`G`$ is a function on $`_+\times _+^m`$ such that $`F(y_0,\mathrm{},y_m):=G(y_0^1,y_1,\mathrm{},y_m)`$ satisfies the assumptions of Theorem 1, then
$$_^nG(|x|,u_1(x),\mathrm{},u_m(x))𝑑x_^nG(|x|,u_1^{}(x),\mathrm{},u_m^{}(x))𝑑x.$$
(3.2)
Hajaiej-Stuart studied this inequality in connection with the following problem in nonlinear optics . The profiles of stable electromagnetic waves traveling along a planar waveguide are given by the ground states of the energy functional
$$(u)=\frac{1}{2}_{}|u^{}|^2𝑑x_{}G(|x|,u)𝑑x$$
under the constraint $`u_2=c`$. Here, $`x`$ is the position relative to the optical axis, $`G`$ is determined by the index of refraction, and $`c>0`$ is a parameter related to the wave speed . If the index of refraction of the optical media decreases with $`|x|`$, then $`F(r,y)=G(r^1,y)`$ satisfies the assumptions of Theorem 1. Then the first integral shrinks under symmetric decreasing rearrangement by the Pólya-Szegő inequality, the second integral grows by Eq. (3.2), and the $`L^2`$-constraint is conserved. Thus, one may rearrange any minimizing sequence to obtain a minimizing sequence of symmetric decreasing functions. This is a crucial step in the construction of ground states — if $`G`$ violates the monotonicity conditions, then a ground state need not exist . Hajaiej-Stuart worried about restrictive regularity assumptions, because $`G`$ may jump at interfaces between layers of different media.
The Riesz inequality in Eq. (1.4) is non-trivial even when $`F`$ is just a product of two functions. Ahlfors introduced two-point rearrangements to treat this case on $`𝕏=\mathrm{SS}^1`$ , Baernstein-Taylor proved the corresponding result on $`\mathrm{SS}^n`$ , and Beckner noted that the proof remains valid on $`^n`$ and $`^n`$ . When $`F`$ is a product of $`m>2`$ functions, Eq. (1.4) has applications to spectral invariants of heat kernels via the Trotter product formula . This case was settled by Friedberg-Luttinger , Burchard-Schmuckenschläger , and by Morpurgo, who proved Eq. (1.4) more generally for integrands of the form
$$F(y_1,\mathrm{},y_m)=\mathrm{\Phi }\left(\underset{i=1}{\overset{m}{}}y_i\right)$$
(3.3)
with $`\mathrm{\Phi }`$ convex (Theorem 3.13 of ). In the above situations, equality cases have been determined . Almgren-Lieb used the technique of Crowe-Zweibel-Rosenbloom to prove Eq. (1.4) for $`m=2`$ . The special case where $`F(u,v)=\mathrm{\Phi }(|uv|)`$ for some convex function $`\mathrm{\Phi }`$ was identified by Baernstein as a ‘master inequality’ from which many classical geometric inequalities can be derived quickly . Eq. (1.4) for continuous supermodular integrands with $`m>2`$ is due to Draghici .
## 4 Outline of the arguments
In their proofs of Eqs. (1.3) and (1.4), Brock and Draghici showed that the left hand sides increase under two-point rearrangements if $`F`$ is any supermodular Borel integrand . Then they approximated the symmetric decreasing rearrangement with sequences of repeated two-point rearrangements. Baernstein-Taylor had established that such sequences can be made to converge to the symmetric decreasing rearrangement in a space of continuous functions , and Brock-Solynin had proved this convergence in $`L^p`$-spaces . To pass to the desired limits, Brock and Draghici assumed that $`F`$ is continuous and satisfies some boundary and growth conditions.
No new proofs of these inequalities will be given here. Rather, we reduce general supermodular integrands to the known cases of integrands that are also bounded and continuous. This reduction needs more care than the usual density arguments, because pointwise $`a.e.`$ convergence of a sequence of integrands $`F_k`$ does not guarantee pointwise $`a.e.`$ convergence of the compositions $`F_k(u_1,\mathrm{},u_m)`$. Approximation within a class of functions with specified positivity or monotonicity properties can be subtle; for instance, nonnegative functions of $`m`$ variables cannot always be approximated by positive linear combinations of products of nonnegative functions of the individual variables (contrary to Theorem 2.1 and Lemma 4.1 of ).
In Section 5, we prove a variant of Sklar’s theorem which factorizes a given supermodular function on $`_+^m`$ as the composition of a Lipschitz continuous supermodular function on $`_+^m`$ with $`m`$ monotone functions on $`_+`$, and a cutoff lemma that replaces a given supermodular function by a bounded supermodular function. Section 6 is dedicated to the two-point versions of Theorems 1 and 2. Here, we review the proofs of the two-point rearrangement inequalities of Lorentz , Brock , and Draghici and find their equality cases. The main theorems are proved in Section 7 by combining the results from Sections 5 and 6. Adapting Beckner’s argument from , we note that the inequalities in Eq. (1.3) and Eq. (1.4) are strict unless $`u_1,\mathrm{},u_m`$ produce equality in all of the corresponding two-point inequalities, and then apply the results from Section 6. In the final Section 8, we briefly discuss extensions for the Brascamp-Lieb-Luttinger and related inequalities.
## 5 Monotone functions
In this section, we provide two technical results about functions with higher-order monotonicity properties. We begin with an auxiliary lemma for functions of a single variable.
###### Lemma 5.1
(Monotone change of variable.) Let $`\varphi `$ be a nondecreasing real-valued function defined on an interval $`I`$. Then, for every function $`f`$ on $`I`$ satisfying
$$|f(z)f(y)|C(\varphi (z)\varphi (y))$$
(5.1)
for all points $`y<zI`$ with some constant $`C`$, there exists a Lipschitz continuous function $`\stackrel{~}{f}:[inff,supf]`$ such that $`f=\stackrel{~}{f}\varphi `$. Furthermore, if $`f`$ is nondecreasing, then $`\stackrel{~}{f}`$ is nondecreasing.
Proof. If $`t=\varphi (y)`$ we set $`\stackrel{~}{f}(t):=f(y)`$. For $`s<t`$ with $`s=\varphi (y)`$, $`t=\varphi (z)`$, Eq. (5.1) implies that
$$|\stackrel{~}{f}(t)\stackrel{~}{f}(s)|=|f(z)f(y)|C(\varphi (z)\varphi (y))=C(ts).$$
(5.2)
Since $`\stackrel{~}{f}`$ is uniformly continuous on the image of $`\varphi `$, it has a unique continuous extension to the closure of the image. The complement consists of a countable number of open disjoint bounded intervals, each representing a jump of $`\varphi `$, and possibly one or two unbounded intervals. On each of the bounded intervals, we interpolate $`\stackrel{~}{f}`$ linearly between the values that have already been assigned at the endpoints. If $`\varphi `$ is bounded either above or below, we extrapolate $`\stackrel{~}{f}`$ to $`t>sup\varphi `$ and $`t<inf\varphi `$ by constants.
By construction, $`f=\stackrel{~}{f}\varphi `$ and $`\stackrel{~}{f}()=[inff,supf]`$. The continuous extension and the linear interpolation preserve the modulus of continuity of $`\stackrel{~}{f}`$, and hence, by Eq. (5.2),
$$|\stackrel{~}{f}(t)\stackrel{~}{f}(s)|C|ts|$$
(5.3)
for all $`s,t`$. If $`f`$ is nondecreasing, then $`\stackrel{~}{f}`$ is nondecreasing on the image of $`\varphi `$ by definition, and on the complement by continuous extension and linear interpolation.
Lemma 5.1 is related to the elementary fact that a continuous random variable can be made uniform by a monotone change of variables. More generally, if $`\varphi `$ is nondecreasing and right continuous, and its generalized inverse is defined by $`\psi (t)=inf\{y:\varphi (y)t\}`$, then the cumulative distribution functions of two random variables that are related by $`Y=\psi (\stackrel{~}{Y})`$ satisfy
$$F(y)=P(Yy)=P\left(\stackrel{~}{Y}\varphi (y)\right)=\stackrel{~}{F}\left(\varphi (y)\right),$$
i.e., $`F=\stackrel{~}{F}\varphi `$. Choosing $`\varphi =F`$ results in a uniform distribution for $`\stackrel{~}{Y}`$.
The corresponding result for $`m2`$ random variables is known as Sklar’s theorem . The theorem asserts that a collection of random variables $`Y_1,\mathrm{},Y_m`$ with a given joint distribution function $`F`$ can be replaced by random variables $`\stackrel{~}{Y}_1,\mathrm{},\stackrel{~}{Y}_m`$ whose marginals $`\stackrel{~}{Y}_i`$ are uniformly distributed on $`[0,1]`$, and whose joint distribution function $`\stackrel{~}{F}`$ is continuous. The next lemma contains Sklar’s theorem for supermodular functions. Since the lemma follows from the arguments outlined in rather than from the statement of the theorem, we include its proof for the convenience of the reader.
We first introduce some notation. Let $`F`$ be a real-valued function on the closed positive cone $`_+^m`$. For $`i=1,\mathrm{},m`$ and $`h0`$, consider the finite difference operators
$$\mathrm{\Delta }_iF(𝐲;h):=F(𝐲+h𝐞_i)F(𝐲).$$
The operators commute, and higher order difference operators are defined recursively by
$$\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F(𝐲;h_1,\mathrm{},h_{\mathrm{}}):=\mathrm{\Delta }_{i_1\mathrm{}i_\mathrm{}1}\mathrm{\Delta }_i_{\mathrm{}}F((𝐲;h_{\mathrm{}});h_1,\mathrm{},h_\mathrm{}1).$$
If $`F`$ is $`\mathrm{}`$ times continuously differentiable, then
$$\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F(𝐲;h_1,\mathrm{},h_{\mathrm{}})=_0^{h_1}\mathrm{}_0^h_{\mathrm{}}_{i_1}\mathrm{}_i_{\mathrm{}}F\left(𝐲+\underset{i=1}{\overset{\mathrm{}}{}}t_i𝐞_i\right)dt_1\mathrm{}dt_{\mathrm{}}.$$
A function $`F`$ is nondecreasing in each variable if $`\mathrm{\Delta }_iF0`$ for $`i=1,\mathrm{},m`$; it is supermodular, if $`\mathrm{\Delta }_{ij}F0`$ for all $`ij`$. The joint distribution function of $`m`$ random variables satisfies $`\mathrm{\Delta }_{i_1}\mathrm{}\mathrm{\Delta }_i_{\mathrm{}}F0`$ for any choice of distinct indices $`i_1,\mathrm{},i_{\mathrm{}}`$.
###### Lemma 5.2
(Sklar’s theorem.) Assume that $`F`$ is bounded, nondecreasing in each variable, and supermodular on $`_+^m`$. Then there exist bounded nondecreasing functions $`\varphi _1,\mathrm{},\varphi _m`$ on $`_+`$ with $`\varphi _i(0)=0`$ and a Lipschitz continuous function $`\stackrel{~}{F}`$ on $`_+^m`$ such that
$$F(y_1,\mathrm{},y_m)=\stackrel{~}{F}(\varphi _1(y_1),\mathrm{},\varphi _m(y_m)).$$
Furthermore, $`\stackrel{~}{F}`$ is bounded, nondecreasing in each variable, and supermodular. If, in addition, $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F0`$ on $`_+^m\times _+^{\mathrm{}}`$ for some distinct indices $`i_1,\mathrm{},i_{\mathrm{}}`$, then $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}\stackrel{~}{F}0`$.
Proof. Set
$$\varphi _i(y)=\underset{y_j\mathrm{},ji}{lim}\left\{F(y_1,\mathrm{},y_m)|_{y_i=y}F(y_1,\mathrm{},y_m)|_{y_i=0}\right\}.$$
These functions are nonnegative and bounded by $`supFinfF`$. Since $`F`$ is nondecreasing in each variable, they are nonnegative, and since $`F`$ is supermodular, they are nondecreasing and satisfy
$$F(𝐲+h𝐞_i)F(𝐲)\varphi _i(y_i+h)\varphi _i(y_i)$$
(5.4)
for all $`𝐲=(y_1,\mathrm{},y_m)_+^m`$ and all $`h>0`$.
We construct $`\stackrel{~}{F}`$ by changing one variable at a time. For the first variable, we write $`𝐲=(y,\widehat{𝐲})`$ where $`y_+`$ and $`\widehat{𝐲}_+^{m1}`$. By Eq. (5.4), for each $`\widehat{𝐲}_+^{m1}`$, the function $`f(y)=F(y,\widehat{𝐲})`$ satisfies Eq. (5.1) with $`C=1`$ and $`\varphi =\varphi _1`$. By Lemma 5.1, there exists a function $`F_1`$ satisfying
$$F(y,\widehat{𝐲})=F_1(\varphi _1(y),\widehat{𝐲})$$
for all $`(y,\widehat{𝐲})_+^m`$. Furthermore, $`F_1`$ is Lipschitz continuous in the first variable,
$$|F_1(t,\widehat{𝐲})F_1(s,\widehat{𝐲})||ts|.$$
We claim that $`F_1`$ satisfies Eq. (5.4) for all $`j>1`$ with the same function $`\varphi _j`$ as $`F`$. To see this, note that for each $`h>0`$ and every $`\widehat{𝐲}`$,
$$f(y)=\mathrm{\Delta }_jF(y,\widehat{𝐲};h)$$
satisfies the assumptions of Lemma 5.1 with $`C=2`$ and $`\varphi =\varphi _1`$. A moment’s consideration shows that
$$\stackrel{~}{f}(t)=\mathrm{\Delta }_jF_1(t,\widehat{𝐲};h)$$
and the claim follows since $`sup\stackrel{~}{f}=supf\varphi _j(y_j+h)\varphi _j(y_j)`$ by Lemma 5.1.
We next verify that $`F_1`$ has the same monotonicity properties as $`F`$. Suppose that $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F0`$ for some set of $`\mathrm{}1`$ distinct indices $`i_1,\mathrm{},i_{\mathrm{}}`$. If $`1\{i_1,\mathrm{},i_{\mathrm{}}\}`$, we apply Lemma 5.1 to $`f(y)=\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F(y,\widehat{𝐲};h_1,\mathrm{},h_{\mathrm{}})`$, which satisfies Eq. (5.1) with $`C=2^{\mathrm{}}`$ and $`\varphi =\varphi _1`$ for all $`\widehat{𝐲}^{m1}`$ and all $`h_1,\mathrm{},h_{\mathrm{}}0`$. It follows that $`\stackrel{~}{f}(t)=\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F_1(t,\widehat{𝐲};h_1,\mathrm{},h_{\mathrm{}})0`$. On the other hand, if $`i_1=1`$, we apply Lemma 5.1 to $`f(y)=\mathrm{\Delta }_{i_2,\mathrm{},i_{\mathrm{}}}F(y,\widehat{𝐲};h_2,\mathrm{},h_{\mathrm{}})`$. Since $`f(y)`$ is nondecreasing by assumption, $`\stackrel{~}{f}(t)=\mathrm{\Delta }_{i_2,\mathrm{},i_{\mathrm{}}}F_1(t,\widehat{𝐲};h_2,\mathrm{},h_{\mathrm{}})`$ is again nondecreasing, and we conclude that $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F_10`$ also in this case.
Iterating the change of variables for $`i=2,\mathrm{},m`$ gives functions $`F_i`$ satisfying
$$F_{i1}(t_1,\mathrm{},t_{i1},y_i,\mathrm{},y_m)=F_i(t_1,\mathrm{},t_{i1},\varphi _i(y_i),y_{i+1},\mathrm{}y_m),$$
as well as
$$0\mathrm{\Delta }_jF_i(t_1,\mathrm{},t_i,y_{i+1},\mathrm{}y_m;h)\{\begin{array}{cc}h,\hfill & ji\hfill \\ \varphi _j(y_j+h)\varphi _j(y_j),\hfill & j>i.\hfill \end{array}$$
(5.5)
Finally, we set $`\stackrel{~}{F}=F_m`$. It follows from Eq. (5.5) that $`\stackrel{~}{F}`$ satisfies the Lipschitz condition $`|\stackrel{~}{F}(𝐳)\stackrel{~}{F}(𝐲)||z_iy_i|\sqrt{m}|𝐳𝐲|`$ for all $`𝐲,𝐳_+^m`$.
The distribution function of a Borel measure on $`_+^m`$ can be conveniently approximated from below by restricting the measure to a large cube $`[0,L)^m`$. The next lemma constructs the corresponding approximation for functions with weaker monotonicity properties.
###### Lemma 5.3
(Cutoff.) Given a real-valued function $`F`$ in $`_+^m`$, set
$$F^L(y_1,\mathrm{},y_m):=F(\mathrm{min}\{y_1,L\},\mathrm{},\mathrm{min}\{y_m,L\}).$$
If $`F`$ is nondecreasing in each variable, then $`F^LF`$. If $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F0`$ on $`_+^m\times _+^{\mathrm{}}`$ for some distinct indices $`i_1,\mathrm{},i_{\mathrm{}}`$, then $`\mathrm{\Delta }_{i_1,\mathrm{},i_{\mathrm{}}}F^L0`$. In particular, if $`F`$ is supermodular, so is $`F^L`$. If $`F`$ has the property that $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F0`$ on $`_+^m\times _+^{\mathrm{}}`$ for every set of distinct indices $`i_1,\mathrm{},i_{\mathrm{}}`$, then $`FF^L`$ also has this property.
Proof. As in the proof of Lemma 5.2, we modify the variables one at a time. The function $`F^{1,L}(y,\widehat{𝐲}):=F(\mathrm{min}\{y,L\},\widehat{𝐲})`$ has the same monotonicity properties as $`F`$ because $`\mathrm{min}\{y,L\}`$ is nondecreasing in $`y`$.
If $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F0`$ for all collections of distinct indices $`i_1,\mathrm{},i_{\mathrm{}}`$, we write
$$F(y,\widehat{𝐲})F^{1,L}(y,\widehat{𝐲})=\mathrm{\Delta }_1F(y,\widehat{𝐲};[yL]_+),$$
and it follows that $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}(FF^{1,L})0`$ whenever $`1\{i_1,\mathrm{},i_{\mathrm{}}\}`$. For $`i_1=1`$, we write
$$\mathrm{\Delta }_1(F(y,\widehat{𝐲};h)F^{1,L}(y,\widehat{𝐲};h)=\mathrm{\Delta }_1F(\mathrm{max}\{y,L\},\widehat{𝐲};[h[Ly]_+]_+),$$
and conclude that $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}(FF^{1,L})0`$ also in this case.
Repeating the construction for the variables $`y_2,\mathrm{},y_m`$ gives the claims.
## 6 Two-point rearrangements
Let $`𝕏`$ be $`^n`$, $`\mathrm{SS}^n`$, or $`^n`$. A reflection on $`𝕏`$ is an isometry characterized by the properties that (i) $`\sigma ^2x=x`$ for all $`x𝕏`$; (ii) the fixed point set $`H_0`$ of $`\sigma `$ separates $`M`$ into two half-spaces $`H_+`$ and $`H_{}`$ that are interchanged by $`\sigma `$; and (iii) $`d(x,x^{})<d(x,\sigma x^{})`$ for all $`x,x^{}H_+`$. We call $`H_+`$ and $`H_{}`$ the positive and negative half-spaces associated with $`\sigma `$. By convention, we always choose $`H_+`$ to contain the distinguished point $`x^{}`$ of $`𝕏`$ in its closure. The two-point rearrangement, or polarization of a real-valued function $`u`$ with respect to a reflection $`\sigma `$ is defined by
$$u^\sigma (x)=\{\begin{array}{cc}\mathrm{max}\{u(x),u(\sigma x)\},\hfill & xH_+H_0\hfill \\ \mathrm{min}\{u(x),u(\sigma x)\},\hfill & xH_{}.\hfill \end{array}$$
This definition makes sense, and the two-point versions of Eqs. (1.3) and (1.4) hold for any space with a reflection symmetry.
On $`𝕏=^n`$, $`\mathrm{SS}^n`$, and $`^n`$, any pair of points is connected by a unique reflection. The space of reflections forms an $`n`$-dimensional submanifold of the $`n(n+1)/2`$-dimensional space of isometries, and thus has a natural uniform metric. If $`u`$ is measurable, both the composition $`u\sigma `$ and the rearrangement $`u^\sigma `$ depend continuously on $`\sigma `$ in the sense that $`\sigma _k\sigma `$ implies that $`u\sigma _ku\sigma `$ and $`u^{\sigma _k}u^\sigma `$ in measure.
Two-point rearrangements are particularly well-suited for identifying symmetric decreasing functions, because
$$u=u^{}u=u^\sigma \text{for all }\sigma .$$
(6.1)
Functions that are radially decreasing about some point are characterized by
$$u=u^{}\tau \text{for some translation }\tau \text{for all }\sigma \text{, either }u=u^\sigma \text{ or }u=u^\sigma \sigma $$
(6.2)
(see Lemma 2.8 of ).
Integral inequalities for two-point rearrangements typically reduce to elementary combinatorial inequalities for the integrands. The following lemma supplies the elementary inequality for the Hardy-Littlewood and Riesz functionals.
###### Lemma 6.1
(Lorentz two-point inequality.) A real-valued function $`F`$ on $`_+^m`$ is supermodular, if and only if for every pair of points $`𝐳,𝐰_+^m`$.
$$\begin{array}{ccc}F(z_1,\mathrm{},z_m)+F(w_1,\mathrm{},w_m)\hfill & & F(\mathrm{max}\{z_1,w_1\},\mathrm{},\mathrm{max}\{z_m,w_m\})\hfill \\ & & +F(\mathrm{min}\{z_1,w_1\},\mathrm{},\mathrm{min}\{z_m,w_m\}).\hfill \end{array}$$
(6.3)
If $`\mathrm{\Delta }_{ij}F>0`$ for some $`ij`$ then Eq. (6.3) is strict unless $`(z_iw_i)(z_jw_j)0`$.
Proof. Given $`𝐳,𝐰_+^m`$, define $`𝐲,𝐡_+^m`$ by $`y_i=\mathrm{min}\{z_i,w_i\}`$ and $`h_i=|z_iw_i|`$ for $`i=1,\mathrm{}m`$. If $`I\{1,\mathrm{},m\}`$, we use the notation $`𝐡_I=_{iI}h_i𝐞_i`$. Subtracting the left hand side of Eq. (6.3) from the right hand side results in the equivalent statement
$$\mathrm{\Delta }_{IJ}F(𝐲;𝐡_I,𝐡_J):=F(𝐲+𝐡_{IJ})F(𝐲+𝐡_I)F(𝐲+𝐡_J)+F(𝐲)0,$$
(6.4)
where $`I=\{i:z_i<w_i\}`$, and $`J=\{i:z_i>w_i\}`$. If either $`I`$ or $`J`$ is empty, Eq. (6.4) is trivially satisfied. If $`I`$ and $`J`$ each have exactly one element, Eq. (6.4) is equivalent to Eq. (1.5). If one of the sets, say $`I`$, has several elements, then decomposing it into disjoint subsets as $`I=I^{}I^{\prime \prime }`$ gives
$$\mathrm{\Delta }_{IJ}F(𝐲;𝐡_I,𝐡_J)=\mathrm{\Delta }_{I^{}J}F(𝐲+𝐡_{I^{\prime \prime }},𝐡_I^{},𝐡_J)+\mathrm{\Delta }_{I^{\prime \prime }J}F(𝐲,𝐡_I^{},𝐡_J),$$
and Eq. (6.4) follows by recursion. The same recursion implies that if $`\mathrm{\Delta }_{ij}F>0`$ and $`z_iw_i`$ and $`z_jw_j`$ have opposite signs, then the inequality in Eq. (6.4) is strict whenever $`I`$ contains $`i`$, $`J`$ contains $`j`$, and $`h_i,h_j>0`$.
Brock proved that the left hand side of Eq. (1.3) increases under two-point rearrangement :
###### Lemma 6.2
(Hardy-Littlewood two-point inequality.) Let $`F`$ be a supermodular Borel measurable function on $`_+^m`$, and let $`u_1,\mathrm{},u_m`$ be nonnegative measurable functions on $`𝕏`$ satisfying the integrability condition in Eq. (2.1). Then, for any reflection $`\sigma `$ on $`𝕏`$,
$$_𝕏F(u_1(x),\mathrm{},u_m(x))𝑑x_𝕏F(u_1^\sigma (x),\mathrm{},u_m^\sigma (x))𝑑x.$$
(6.5)
Assume furthermore that $`\mathrm{\Delta }_{ij}F>0`$ on $`_+^m\times (0,\mathrm{})^2`$ for some $`ij`$. If Eq. (6.5) holds with equality and the integrals are finite, then
$$\left(u_i(x)u_i(\sigma x)\right)\left(u_j(x)u_j(\sigma x)\right)0a.e..$$
In particular, if $`u_i=u_i^{}`$ is strictly radially decreasing and $`\sigma (x^{})x^{}`$, then $`u_j=u_j^\sigma `$.
Proof. The inequality : The left hand side of Eq. (6.5) can be written as an integral over the positive half-space,
$$(u_1,\mathrm{},u_m):=_{H_+}F(u_1(x),\mathrm{}u_m(x))+F(u_1(\sigma x),\mathrm{}u_m(\sigma x))dx.$$
By Lemma 6.1, with $`z_i=u_i(x)`$ and $`w_i=u_i(\sigma x)`$, the integrand satisfies
$$\begin{array}{c}F(u_1(x),\mathrm{}u_m(x))+F(u_1(\sigma x),\mathrm{}u_m(\sigma x))\hfill \\ F(u_1^\sigma (x),\mathrm{}u_m^\sigma (x))+F(u_1^\sigma (\sigma x),\mathrm{}u_m^\sigma (\sigma x))\hfill \end{array}$$
(6.6)
for all $`xH_+`$. Integrating over $`H_+`$ yields Eq. (6.5).
Equality statement: Assume that $`(u_1,\mathrm{},u_m)=(u_1^\sigma ,\mathrm{},u_m^\sigma )`$ is finite. Then Eq. (6.6) must hold with equality almost everywhere on $`H_+`$. If $`\mathrm{\Delta }_{ij}F>0`$ on $`_+^m\times (0,\mathrm{})^2`$, then Lemma 6.1 implies that $`u_i(x)u_i(\sigma x)`$ and $`u_j(x)u_j(\sigma x)`$ cannot have opposite signs except on a set of zero measure. If moreover $`u_i=u_i^{}`$ is strictly radially decreasing and $`\sigma x^{}x^{}`$, then $`u_i(x)>u_i(\sigma x)`$ for $`a.e.xH_+`$, and Lemma 6.1 implies that $`u_j(x)u_j(\sigma x)`$ for $`a.e.xH_+`$.
Brock completed the proof of Eq. (1.3) by approximating the symmetric decreasing rearrangement with a sequence of two-point rearrangements à la Baernstein-Taylor . We sketch his argument in the simplest case where $`F`$ is a continuous supermodular function that vanishes on the boundary of the positive cone $`_+^m`$, and $`u_1,\mathrm{},u_m`$ are bounded and compactly supported.
By Theorem 6.1 of there exists a sequence of reflections $`\{\sigma _k\}_{k1}`$ such that
$$u_i^{\sigma _1,\mathrm{},\sigma _k}u_i^{}\text{in measure}(k\mathrm{})$$
(6.7)
for $`i=1,\mathrm{},m`$. By Lemma 6.2, the functional increases monotonically along such a sequence. If $`B`$ is a ball centered at $`x^{}`$ that contains the supports of $`u_1,\mathrm{},u_m`$, then the rearranged functions $`u_i^{\sigma _1,\mathrm{},\sigma _k}`$ are also supported on $`B`$, and dominated convergence yields
$$(u_1,\mathrm{},u_m)(u_1^{\sigma _1,\mathrm{},\sigma _k},\mathrm{},u_m^{\sigma _1,\mathrm{},\sigma _k})(u_1^{},\mathrm{},u_m^{})(k\mathrm{}).$$
(6.8)
The corresponding results for Eq. (1.4) are due to Draghici . The two-point inequality is not an immediate consequence of Lemma 6.1, but requires an additional combinatorial argument. This argument was used previously by Morpurgo , and a simpler version appears in .
###### Lemma 6.3
(Riesz two-point inequality.) Assume that $`F`$ is a supermodular Borel measurable function on $`_+^m`$. For each pair of indices $`1i<jm`$, let $`K_{ij}`$ be a nonincreasing function on $`_+`$, and let $`u_1,\mathrm{},u_m`$ be nonnegative measurable functions on $`𝕏`$ satisfying the integrability condition in Eq. (2.3). Then, for any reflection $`\sigma `$,
$$\begin{array}{c}_𝕏\mathrm{}_𝕏F(u_1(x_1),\mathrm{},u_m(x_m))\underset{i<j}{}K_{ij}\left(d(x_i,x_j)\right)dx_1\mathrm{}dx_m\hfill \\ _𝕏\mathrm{}_𝕏F(u_1^\sigma (x_1),\mathrm{},u_m^\sigma (x_m))\underset{i<j}{}K_{ij}\left(d(x_i,x_j)\right)dx_1\mathrm{}dx_m.\hfill \end{array}$$
(6.9)
Assume additionally that that $`K_{ij}(t)>0`$ for all $`i<j`$ and all $`t<\mathrm{diam}𝕏`$. Let $`\mathrm{\Gamma }_0`$ be the graph on $`\{1,\mathrm{},m\}`$ with an edge between $`i`$ and $`j`$ whenever $`K_{ij}`$ is strictly decreasing. If $`\mathrm{\Delta }_{ij}F>0`$ for some $`ij`$ lying in the same connected component of $`\mathrm{\Gamma }_0`$, and that $`u_i`$ and $`u_j`$ are not symmetric under $`\sigma `$. If the integrals in Eq. (6.9) have the same finite value, then either $`u_i=u_i^\sigma `$ and $`u_j=u_j^\sigma `$, or $`u_i=u_i^\sigma \sigma `$ and $`u_j=u_j^\sigma \sigma `$.
Proof. The inequality : The left hand side of Eq. (6.9) can be written as an $`m`$-fold integral over the positive half-space
$`(u_1,\mathrm{},u_m)`$ $`:=`$ $`{\displaystyle _{H_+}}\mathrm{}{\displaystyle _{H_+}}{\displaystyle \underset{\epsilon _i\{0,1\},i=1,\mathrm{},m}{}}\{F(u_1(\sigma ^{\epsilon _1}x_1),\mathrm{},u_m(\sigma ^{\epsilon _m}x_m))\times `$ (6.10)
$`\times {\displaystyle \underset{i<j}{}}K_{ij}\left(d(\sigma ^{\epsilon _i}x_i,\sigma ^{\epsilon _j}x_j)\right)\}dx_1\mathrm{}dx_m.`$
Fix $`x_1,\mathrm{},x_mH_+`$. For each pair of indices $`i<j`$, set $`a_{ij}=K_{ij}\left(d(x_i,\sigma x_j)\right)`$ and $`b_{ij}=K_{ij}\left(d(x_i,x_j)\right)K_{ij}\left(d(x_i,\sigma x_j)\right)`$, so that
$$K_{ij}\left(d(\sigma ^{\epsilon _i}x_i,\sigma ^{\epsilon _j}x_j)\right)=a_{ij}+b_{ij}\mathrm{𝟏}_{\epsilon _i=\epsilon _j}.$$
The product term in Eq. (6.10) expands to
$$\underset{i<j}{}K_{ij}\left(d(\sigma ^{\epsilon _i}x_i,\sigma ^{\epsilon _j}x_j)\right)=\underset{\mathrm{\Gamma }}{}\left(\underset{ijE}{}a_{ij}\right)\left(\underset{ijE}{}b_{ij}\mathrm{𝟏}_{\epsilon _i=\epsilon _j}\right)=:C_\mathrm{\Gamma }\mathrm{𝟏}_{\epsilon _i=\epsilon _j,ijE},$$
where $`\mathrm{\Gamma }`$ runs over all proper graphs on the vertex set $`V=\{1,\mathrm{},m\}`$, and $`E`$ is the set of edges of $`\mathrm{\Gamma }`$. Inserting the expansion into Eq. (6.10) and exchanging the order of summation shows that each graph contributes a nonnegative term
$$C_\mathrm{\Gamma }\underset{\epsilon _i\{0,1\},iV}{}F(u_1(\sigma ^{\epsilon _1}x_1),\mathrm{},u_m(\sigma ^{\epsilon _m}x_m))\mathrm{𝟏}_{\epsilon _i=\epsilon _j,ijE}$$
(6.11)
to the integral in Eq. (6.10). If $`\mathrm{\Gamma }`$ is connected, then
$`{\displaystyle \underset{\epsilon _i\{0,1\},iV}{}}F(u_1(\sigma ^{\epsilon _1}x_1),\mathrm{},u_m(\sigma ^{\epsilon _m}x_m))\mathrm{𝟏}_{\epsilon _i=\epsilon _j,ijE}`$
$`=`$ $`F(u_1(x_1),\mathrm{},u_m(x_m))+F(u_1(\sigma x_1),\mathrm{},u_m(\sigma x_m))`$
$``$ $`F(u_1^\sigma (x_1)),\mathrm{},u_m^\sigma (x_m))+F(u_1^\sigma (\sigma x_1)),\mathrm{},u_m^\sigma (\sigma x_m))`$
$`=`$ $`{\displaystyle \underset{\epsilon _i\{0,1\},iV}{}}F(u_1^\sigma (\sigma ^{\epsilon _1}x_1),\mathrm{},u_m^\sigma (\sigma ^{\epsilon _m}x_m))\mathrm{𝟏}_{\epsilon _i=\epsilon _j,ijE},`$
where the second step follows from Lemma 6.1 with $`z_i=u_i(x_i)`$ and $`w_i=u_i(\sigma x_i)`$.
If $`\mathrm{\Gamma }`$ is not connected, choose a connected component $`\mathrm{\Gamma }^{}`$ and let $`\mathrm{\Gamma }^{\prime \prime }`$ be its complement. Let $`E^{}`$, $`E^{\prime \prime }`$, $`V^{}`$, and $`V^{\prime \prime }`$ be the corresponding edge and vertex sets. The sum in Eq. (6.11) can be decomposed as
$$\underset{\epsilon _i\{0,1\},iV^{\prime \prime }}{}\left\{\underset{\epsilon _i\{0,1\},iV^{}}{}F(u_1(\sigma ^{\epsilon _1}x_1),\mathrm{},u_m(\sigma ^{\epsilon _m}x_m))\mathrm{𝟏}_{\epsilon _i=\epsilon _j,ijE^{}}\right\}\mathrm{𝟏}_{\epsilon _i=\epsilon _j,ijE^{\prime \prime }}.$$
The key observation is that Eq. (6) applies to the term in braces for fixed $`\epsilon _i,iV^{\prime \prime }`$; in other words, the contribution of $`\mathrm{\Gamma }`$ can only increase if $`u_i`$ is replaced by $`u_i^\sigma `$ for all $`iV^{}`$. An induction over the connected components of $`\mathrm{\Gamma }`$ shows that
$$\begin{array}{c}\underset{\epsilon _i\{0,1\},iV}{}F(u_1(\sigma ^{\epsilon _1}x_1),\mathrm{},u_m(\sigma ^{\epsilon _m}x_m))\mathrm{𝟏}_{\epsilon _i=\epsilon _j,ijE}\hfill \\ \underset{\epsilon _i\{0,1\},iV}{}F(u_1^\sigma (\sigma ^{\epsilon _1}x_1),\mathrm{},u_m^\sigma (\sigma ^{\epsilon _m}x_m))\mathrm{𝟏}_{\epsilon _i=\epsilon _j,ijE}\hfill \end{array}$$
for any graph $`\mathrm{\Gamma }=(E,V)`$. Adding the contributions of all graphs shows that the integrand in Eq. (6.10) increases pointwise under two-point rearrangement, and Eq. (6.9) follows.
Equality statement: Let $`\mathrm{\Gamma }_0`$ be the graph defined in the statement of the lemma, and let $`E_0`$ be its edge set. By assumption,
$$C_{\mathrm{\Gamma }_0}=\left(\underset{ijinE_0}{}K_{ij}\left(d(x_i,x_j)\right)K_{ij}\left(d(\sigma x_i,x_j)\right)\right)\left(\underset{ijE_0}{}K_{ij}\left(d(\sigma x_i,x_j)\right)\right)>0$$
for $`a.e.x_1,\mathrm{},x_mH_+`$. If $`\mathrm{\Delta }_{ij}F>0`$, then Lemma 6.1 implies that Eq. (6) is strict unless
$$\left(u_i(x_i)u_i(\sigma x_i)\right)\left(u_j(x_j)u_j(\sigma x_j)\right)0,a.e.x_i,x_jH_+.$$
If $`u_i`$ and $`u_j`$ are not symmetric under $`\sigma `$, the product is not identically zero. Since $`x_i`$ and $`x_j`$ can vary independently, this means that $`u_i(x)u_i(\sigma x)`$ and $`u_j(x)u_j(\sigma x)`$ cannot change sign on $`H_+`$. We conclude that equality in Eq. (6.9) implies that either $`u_i=u_i^\sigma `$ and $`u_j=u_j^\sigma `$, or $`u_i=u_i^\sigma \sigma `$ and $`u_j=u_j^\sigma \sigma `$.
Draghici also used Baernstein-Taylor approximation to obtain Eq. (1.4) from Eq. (6.9). If $`F`$ is bounded and continuous and $`K_{ij}`$ is bounded for $`1i<jm`$, then for bounded functions $`u_1,\mathrm{},u_m`$ that are supported in a common ball $`B`$ the inequality follows from Lemma 6.3 by approximating the symmetric decreasing rearrangement with a sequence of two-point rearrangements, see Eq. (6.7). Dominated convergence applies as in Eq. (6.8), since the integrations extend only over the bounded set $`B^m`$.
## 7 Proof of the main results
Proof of Theorem 1. The inequality for Borel integrands: Let $`F`$ be a supermodular Borel function with $`F(\mathrm{𝟎})=0`$, and let and $`u_1,\mathrm{},u_m`$ be nonnegative measurable functions that vanish at infinity, as in the statement of the theorem. Denote by
$$(u_1,\mathrm{},u_m):=_𝕏F(u_1(x),\mathrm{},u_m(x))𝑑x$$
the left hand side of Eq. (1.3). Replacing $`F(𝐲)`$ by $`F(𝐲)_{i=1}^mF(y_i𝐞_i)`$ and using that $`F(u_i()𝐞_i)`$ and $`F(u_i^{}()𝐞_i)`$ contribute equally to the two sides of Eq. (1.3), we may assume $`F`$ to be nondecreasing in each variable.
Fix $`L>0`$, and replace $`u_i`$ by the bounded function
$$u_i^L(x):=\mathrm{min}\{u_i(x),L\}\mathrm{𝟏}_{\{|x|<L\}}$$
for $`i=1,\mathrm{},m`$. Then
$$F(u_1^L,\mathrm{},u_m^L)=F^L(u_1^L,\mathrm{},u_m^L),$$
(7.1)
where $`F^L`$ is the function defined in Lemma 5.3. By construction, $`F^L`$ is bounded, and by Lemma 5.3 it is nondecreasing and supermodular. By Lemma 5.2, there exist nondecreasing functions $`\varphi _i`$ with $`\varphi _i(0)=0`$ and a continuous supermodular function $`\stackrel{~}{F}^L`$ on $`_+^m`$ such that
$$F^L(y_1,\mathrm{},y_m)=\stackrel{~}{F}^L(\varphi _1(y_1),\mathrm{},\varphi _m(y_m)).$$
(7.2)
Since $`\varphi _i`$ is nondecreasing and vanishes at zero, $`u_i^L`$ is compactly supported, and $`(u_i^L)^{}(u_i^{})^L`$ pointwise by construction, we have
$$(\varphi _iu_i^L)^{}=\varphi _i(u_i^L)^{}\varphi _i(u_i^{})^L$$
(7.3)
for $`i=1,\mathrm{},m`$. By Theorem 1 of ),
$$_𝕏\stackrel{~}{F}^L(\varphi _1u_1^L(x),\mathrm{},\varphi _mu_m^L(x))𝑑x_𝕏\stackrel{~}{F}^L((\varphi _1u_1^L)^{}(x),\mathrm{},(\varphi _mu_m^L)^{}(x))𝑑x.$$
With Eqs. (7.1)-(7.3), this becomes
$$(u_1^L,\mathrm{},u_m^L)((u_1^{})^L,\mathrm{},(u_m^{})^L).$$
Since $`u_i^L(x)=u_i(x)`$ for $`L\mathrm{max}\{u_i(x),|x|\}`$, we see that $`F(u_i^L(x),\mathrm{},u_m^L(x))`$ converges pointwise to $`F(u_1(x),\mathrm{},u_m(x))`$, and Eq. (1.3) follows by monotone convergence.
Equality statement: Combining Eq. (6.5) with Eq. (1.3) and using that $`u_i^\sigma `$ is equimeasurable with $`u_i`$, we see that
$$(u_1,\mathrm{},u_m)(u_1^\sigma ,\mathrm{},u_m^\sigma )(u_1^{},\mathrm{},u_m^{}).$$
Hence equality in Eq. (1.3) implies equality in Eq. (6.5) for every choice of the reflection $`\sigma `$. Given two points $`x,x^{}`$ in $`𝕏`$, choose $`\sigma `$ such that $`\sigma (x)=x^{}`$. If $`\mathrm{\Delta }_{ij}F>0`$ for some $`ij`$, then $`u_i(x)u_i(x^{})`$ and $`u_j(x)u_j(x^{})`$ cannot have opposite signs by Lemma 6.2. If $`u_i=u_i^{}`$ is strictly radially decreasing, then it follows that $`u_j^\sigma =u_j`$ for every reflection $`\sigma `$ that does not fix $`x^{}`$. By Eq. (6.1), $`u_j=u_j^{}`$ as claimed.
Proof of Theorem 2. The inequality for Borel integrands: The proof of Eq. (1.4) proceeds along the same lines as the proof of Eq. (1.3). Let
$$(u_1,\mathrm{},u_m):=_𝕏\mathrm{}_𝕏F(u_1(x_1),\mathrm{},u_m(x_m))\underset{i<j}{}K_{ij}(d(x_i,x_j))dx_1\mathrm{}dx_m$$
be the left hand side of Eq. (1.4). As before, we may assume that $`F`$ is nondecreasing in each variable. We replace $`F`$ with $`\stackrel{~}{F}^L`$, $`u_i`$ with $`\varphi _iu_i^L`$, $`K_{ij}`$ with $`K_{ij}^L=\mathrm{min}\{K_{ij},L\}`$, and set
$$^L(u_1,\mathrm{},u_m):=_𝕏\mathrm{}_𝕏F^L(u_1(x_1),\mathrm{},u_m(x_m))\underset{i<j}{}K_{ij}^L\left(d(x_i,x_j)\right)dx_1\mathrm{}dx_m.$$
Applying Theorem 2.2 of , we obtain with the help of Eqs. (7.1)-(7.3)
$$^L(u_1^L,\mathrm{},u_m^L)^L((u_1^{})^L,\mathrm{},(u_m^{})^L).$$
Eq. (1.4) follows by taking $`L\mathrm{}`$ and using monotone convergence.
Equality statement: Consider the set $`S_i`$ of all reflections $`\sigma `$ of $`𝕏`$ that fix $`u_i`$. If $`u_i`$ is non-constant, then $`S_i`$ is a closed proper subset of the space of all reflections on $`𝕏`$. This subset is nowhere dense, since any open set of reflections generates the entire isometry group of $`𝕏`$. If Eq. (1.4) holds with equality, then the two-point rearrangement inequality in Eq. (6.9) holds with equality for every reflection $`\sigma `$. For $`\sigma S_i`$, Lemma 6.3 implies that either $`u_j=u_j^\sigma `$ or $`u_j=u_j^\sigma \sigma `$. Since $`S_i`$ is nowhere dense, it follows from the continuous dependence of $`u^\sigma `$ on $`\sigma `$ that $`u_j`$ agrees with either $`u_j^\sigma `$ or $`u_j^\sigma \sigma `$ also for $`\sigma S_i`$. By Eq. (6.2), there exists a translation $`\tau `$ such that $`u_j=u_j^{}\tau `$. Lemma 6.3 implies furthermore that $`u_i`$ agrees with $`u_i^\sigma `$ when $`u_j=u_j^\sigma `$, and with $`u_i^\sigma \sigma `$ when $`u_j=u_j^\sigma \sigma `$. We conclude that $`u_i=u_i^{}\tau `$.
## 8 Concluding remarks
In the proof of Eq. (1.4) and its two-point version in Eq. (6.9), the kernels $`K_{ij}`$ played a very different role from the functions $`u_1,\mathrm{},u_m`$ that enter into the integrand. However, the Riesz functional on the left hand side of Eq. (1.2) depends equally on $`u`$, $`v`$, and $`w`$. We will use the connection of Riesz’ inequality with the Brunn-Minkowski inequality to construct examples where the two-point rearrangement fails for Eq. (1.2).
The Brunn-Minkowski inequality says that the measures of two subsets $`A,B^n`$ are related to the measure of their Minkowski sum $`A+B=\{a+b:aA,bB\}`$ by
$$\lambda (A)^{1/n}+\lambda (B)^{1/n}\lambda (A+B)^{1/n}.$$
Recognizing the two sides of the inequality as proportional to the radii of the balls $`A^{}+B^{}`$ and $`(A+B)^{}`$, we rewrite it as the rearrangement inequality
$$\lambda (A^{}+B^{})\lambda (A+B).$$
(8.1)
Eq. (8.1) follows rather directly from Riesz’ inequality in Eq. (1.2), because the support of the convolution of two nonnegative functions is essentially the Minkowski sum of their supports. Conversely, the Brunn-Minkowski inequality enters into the proof of the Brascamp-Lieb-Luttinger inequality , of which Eqs. (1.2) and (1.4) are special cases.
Equality in the Brunn-Minkowski inequality implies that $`A`$ and $`B`$ differ only by sets of measure zero from two independently scaled and translated copies of a convex body . Let $`A=B`$ be an ellipsoid in $`^n`$ with $`n>1`$ that is centered at a point $`c0`$, so that Eq. (8.1) holds with equality. If $`\sigma `$ is the reflection at a hyperplane through $`c`$ that is not a hyperplane of symmetry for $`A`$ and $`B`$, then $`A^\sigma `$ and $`B^\sigma `$ are non-convex, and therefore
$$\lambda (A^\sigma +B^\sigma )>\lambda (A^{}+B^{})=\lambda (A+B).$$
Choosing $`u`$, $`v`$, and $`w`$ as the characteristic functions of $`A`$, $`A+B`$, and $`B`$ provides an example where the Riesz functional strictly decreases under two-point rearrangement. For an example of this phenomenon in one dimension, consider the symmetric decreasing functions
$$u(x)=\mathrm{𝟏}_{|x2|<\epsilon },v(x)=w(x)=\mathrm{𝟏}_{|x1|<\epsilon },$$
and let $`\sigma `$ be the reflection at $`x=1`$. Then
$$u^\sigma (x)=\mathrm{𝟏}_{|x|<\epsilon },v^\sigma (x)=w^\sigma (x)=\mathrm{𝟏}_{|x1|<\epsilon },$$
and if $`0<\epsilon \frac{1}{2}`$, Riesz’ inequality fails for $`\sigma `$,
$$\begin{array}{ccc}_{}_{}u(x)v(x^{})w(xx^{})𝑑x𝑑x^{}\hfill & >& 0\hfill \\ & =& _{}_{}u^\sigma (x)v^\sigma (x)w^\sigma (xx^{})𝑑x𝑑x^{}.\hfill \end{array}$$
While the two-point rearrangement is not useful for Eq. (1.2), the layer-cake representation of Crowe-Zweibel-Rosenbloom shows that
$$\begin{array}{c}_^n_^nF(u(x),v(x^{}),w(xx^{}))𝑑x𝑑x^{}\hfill \\ _^n_^nF(u^{}(x),v^{}(x^{}),w^{}(xx^{}))𝑑x𝑑y\hfill \end{array}$$
(8.2)
for any integrand that can be written as the joint distribution function of a Borel measure $`\mu _F`$ on $`_+^3`$,
$$F(y_1,y_2,y_3)=\mu _F\left([0,y_1)\times [0,y_2)\times [0,y_3)\right).$$
Such integrands are left continuous, vanish at the origin, and satisfy $`\mathrm{\Delta }_{i_1,\mathrm{},i_{\mathrm{}}}F0`$ for every choice of $`\mathrm{}3`$ distinct indices. Lemma 5.2 allows to accommodate integrands in Eq. (8.2) that are only Borel measurable. The main condition is that $`\mathrm{\Delta }_{123}F0`$; the second-order monotonicity conditions can be replaced by integrability assumptions on the negative part $`F_{}`$ similar to Eq. (2.3). To ensure that the functional is finite at least when $`u,v,w`$ are bounded and compactly supported, $`F`$ should vanish on the coordinate axes. For example, Eq. (8.2) holds for
$$F(u,v,w)=\frac{uvw}{(1+u)(1+v)(1+w)}(uv+uw+vw)$$
since $`\mathrm{\Delta }_{123}F>0`$, even though $`\mathrm{\Delta }_{ij}F<0`$ for all $`ij`$.
For Borel integrands satisfying $`\mathrm{\Delta }_{123}F>0`$, equality in Eq. (8.2) implies that every triple of level sets of $`u,v,w`$ produces equality in Eq. (1.2). These equality cases were described in . In particular, if two of the three functions $`u,v,w`$ are known to have continuous distribution functions and the value of the functional is finite, then equality implies that $`u,v,w`$ are equivalent to $`u^{},v^{},w^{}`$ under the symmetries of the functional (see Theorem 2 of ).
By the same line of reasoning, the Brascamp-Lieb-Luttinger inequality implies that
$$(u_1,\mathrm{},u_m):=_^n\mathrm{}_^nF(u_1\left(\underset{j=1}{\overset{k}{}}a_{1j}x_j\right),\mathrm{},u_m\left(\underset{j=1}{\overset{k}{}}a_{mj}x_j\right))𝑑x_1\mathrm{}𝑑x_k$$
increases under symmetric decreasing rearrangement, if $`\mathrm{\Delta }_{i_1,\mathrm{},i_{\mathrm{}}}F0`$ for all choices of distinct indices $`i_1,\mathrm{},i_{\mathrm{}}`$ with $`\mathrm{}m`$. Interesting examples are integrands of the form in Eq. (3.3), where $`\mathrm{\Phi }`$ is completely monotone in the sense that all its distributional derivatives are nonnegative. If $`\mathrm{\Delta }_{i_1\mathrm{}i_{\mathrm{}}}F>0`$ for all choices of $`i_1,\mathrm{},i_{\mathrm{}}`$, then the last statement of Lemma 5.3 can be used to show that the extended Brascamp-Lieb-Luttinger inequality has the same equality cases as the original inequality. The characterization of these equality cases remains an open problem.
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# Charged Lepton Production from Iron Induced by Atmospheric Neutrinos
## I Introduction
The study of neutrino physics with atmospheric neutrinos has a long history with first observations of muons produced by atmospheric muon neutrinos in deep underground laboratories of KGF in India and ERPM in South Africaachar . The indications of some deficit in the atmospheric neutrino flux was known to exist from the early days of these experiments but the evidence was no more than suggestive due to low statistics of the experimental data and anticipated uncertainties in the flux calculationscrouch . The clear evidence of a deficit in the atmospheric muon neutrino flux was confirmed later when data with better statistics were obtained at IMBheines , Kamiokandehirata and Soudanallison experiments. The most likely cause of this deficit is believed to be the phenomena of neutrino oscillationspontecarvo in which the neutrinos produced with muon flavor after passing a certain distance through the atmosphere, manifest themselves as a different flavor. The implication of this phenomena of neutrino oscillation is that neutrinos possess a nonzero mass pointing towards physics beyond the standard model of particle physics. The evidence for neutrino oscillations and a nonzero mass for the neutrinos has also been obtained in the observations made with solarsolar and reactor (anti)neutrinosreactor .
It is well known that, in a two flavor oscillation scenario involving muon neutrino, the probability for a muon neutrino with energy $`E_\nu `$ to remain a muon neutrino after propagating a distance $`L`$ before reaching the detector is given bypontecarvo .
$$P_{\mu \mu }=1sin^22\theta sin^2\left(\frac{1.27\mathrm{\Delta }m^2(eV^2)L(km)}{E_\nu (GeV)}\right)$$
(1)
where $`\mathrm{\Delta }m^2=m_1^2m_2^2`$ is the difference of the squared masses of the two flavor mass eigenstates and $`\theta `$ is the mixing angle between two states. The oscillation parameters $`\mathrm{\Delta }m^2`$ and the mixing angle $`\theta `$ are determined by various observations made in atmospheric neutrino experiments. These include the flavor ratios of muon and electron flavors, angular and $`\frac{L}{E}`$ distributions of muons and electrons produced by atmospheric neutrinos. The first claims of seeing neutrino oscillation in atmospheric neutrinos came from the measurements of ratio of ratios $`R_\nu `$ defined as $`\frac{(\mu /e)_{data}}{(\mu /e)_{MC}}`$ from the observations of fully contained (FC) eventsheines -hirata , but there are now data available from the angular and $`\frac{L}{E}`$ distribution of the atmospheric neutrino induced muon and electron events from SKsk , MACROmacro and Soudansoudan experiments which confirm the phenomena of neutrino oscillations. These experiments are consistent with a value of $`\mathrm{\Delta }m^23.2\times 10^3eV^2`$ and $`sin^22\theta 1`$. The analysis of these data assuming a three flavor neutrino oscillation phenomenology have also been done by many authorsthreefla .
The major sources of uncertainty in the theoretical prediction of the charged leptons of muon and electron flavor produced by the atmospheric neutrinos come from the uncertainties in the calculation of atmospheric neutrino fluxes and neutrino nuclear cross sections. The atmospheric neutrino fluxes at various experimental sites of Kamioka, Soudan and Gransasso have been extensively discussed in literature by many authorsbattistoni -plyaskin . The neutrino nuclear cross sections have also been calculated for various nuclei by many authors using different nuclear modelsdonnelly -singh1 . The aim of the present paper is to study the neutrino nuclear cross section in iron nuclei which are relevant for the atmospheric neutrino experiments performed at Soudansoudan , FREJUSfrejus and NUSEXnusex and planned in future with MINOSminos , MONOLITHmonolith and INOino detectors. The uncertainty in the nuclear production cross section of leptons from iron nuclei by the atmospheric neutrinos are discussed. For our nuclear model, we also discuss the uncertainty due to use of different neutrino fluxes for the sites of Soudan and Gransasso which are relevant to MINOS, MONOLITH and INO detectorsminos -ino .
The momentum and angular distribution of muons and electrons relevant to fully contained events produced by atmospheric neutrinos in iron nuclei are calculated. These leptons of muon and electron flavor characterized by track and shower events include the leptons produced by quasielastic process as well as the inelastic processes induced by charged current interactions. The calculations are done in a model which takes into account nuclear effects like Pauli Blocking, Fermi motion effects and the effect of renormalization of the weak transition strengths in nuclear medium in local density approximation. The model has been successfully applied to describe various electromagnetic and weak processes like photon absorption, electron scattering, muon capture and low energy neutrino reactions in nucleisingh1 ,gil -singh4 . The model can be easily applied to calculate the zenith angle dependence and the $`\frac{L}{E}`$ distribution for stopping and thorough going muon production from iron nuclei which is currently under progress.
The plan of the paper is as follows. In section-II we describe the neutrino(antineutrino) quasielastic inclusive production of leptons $`(e^{},\mu ^{},e^+,\mu ^+)`$ from iron nuclei for various neutrino energies. In section-III we describe the energy dependence of the inelastic production of leptons through $`\mathrm{\Delta }`$-dominance model and highlight the nuclear effects relevant to the energy of fully contained events. In section-IV, we use the atmospheric neutrino flux at Soudan and Gransasso sites as determined by various authors and discuss the flux averaged momentum and angular dependence of leptons corresponding to different flux calculations available at these two sites.
## II Quasielastic Production of Leptons
Quasielastic inclusive production of leptons in nuclei induced by neutrinos has been studied by many authors where nuclear effects have been calculated. Most of these calculations have been done either for $`{}_{}{}^{16}O`$ relevant to IMB and Kamioka experimentsheines -hirata or for $`{}_{}{}^{12}C`$ relevant to LSND and KARMEN experimentlsnd1 . These calculations generally use direct summation method (over many nuclear excited states)donnelly , closure approximationgoulard , Fermi gas modelsmith -smith1 , relativistic mean field approximationkim , continuum random phase approximation (CRPA)engel and local density approximationsingh2 ; singh1 . The calculations for $`{}_{}{}^{56}Fe`$ nucleus have been reported by Bugaev et al.donnelly in a shell model and in Fermi gas model by Gallaghergallagher and Berger et al.berger . In this section we briefly describe the formalism and results of our calculations done for quasielastic inclusive production of leptons for iron nuclei.
### II.1 Formalism
In local density approximation the neutrino nucleus cross section $`\sigma (E_\nu )`$ for a neutrino of energy $`E_\nu `$ scattering from a nucleus $`A(Z,N)`$, is given by
$$\sigma ^A(E_\nu )=2𝑑\stackrel{}{r}\frac{d\stackrel{}{p}}{(2\pi )^3}n_n(\stackrel{}{p},\stackrel{}{r})\sigma ^N(E_\nu )$$
(2)
where $`n_n(\stackrel{}{p},\stackrel{}{r})`$ is the local occupation number of the initial nucleon of momentum $`\stackrel{}{p}`$ (localized at position $`\stackrel{}{r}`$ in the nucleus) and $`\sigma ^N(E_\nu )`$ is cross section for the scattering of neutrino of energy $`E_\nu `$ from a free nucleon given by the expression
$`\sigma ^N(E_\nu )`$ $`=`$ $`{\displaystyle \frac{d^3k^{}}{(2\pi )^3}\frac{m_\nu }{E_\nu }\frac{m_e}{E_e}\frac{M_n}{E_n}\frac{M_p}{E_p}}`$ (3)
$`\times `$ $`\overline{{\displaystyle }}{\displaystyle |T|^2\delta (E_\nu E_l+E_nE_p)}`$
where $`T`$ is the matrix element for the basic process
$$\nu _l(k)+n(p)l^{}(k^{})+p(p^{}),l=e,\mu $$
(4)
written as
$$T=\frac{G_F}{\sqrt{2}}\mathrm{cos}\theta _c\overline{u}(k^{})\gamma _\mu (1\gamma _5)u(k)(J^\mu )^{CC}$$
(5)
$`(J^\mu )^{CC}`$ is the charged current(CC) matrix element of the hadronic current defined as
$`(J^\mu )^{CC}=\overline{u}(p^{})[F_1^V(q^2)\gamma ^\mu +F_2^V(q^2)i\sigma ^{\mu \nu }{\displaystyle \frac{q_\nu }{2M}}`$ (6)
$`+F_A^V(q^2)\gamma ^\mu \gamma ^5]u(p)`$
$`q^2(q=kk^{})`$ is the four momentum transfer square. The form factors $`F_1^V(q^2)`$, $`F_2^V(q^2)`$ and $`F_A^V(q^2)`$ are isovector electroweak form factors written as
$$F_1^V(q^2)=F_1^p(q^2)F_1^n(q^2),F_2^V(q^2)=F_2^p(q^2)F_2^n(q^2),$$
$$F_A^V(q^2)=F_A(q^2)$$
where
$`F_1^{p,n}(q^2)`$ $`=`$ $`{\displaystyle \frac{1}{(1\frac{q^2}{4M^2})}}\left[G_E^{p,n}(q^2){\displaystyle \frac{q^2}{4M^2}}G_M^{p,n}(q^2)\right]`$
$`F_2^{p,n}(q^2)`$ $`=`$ $`{\displaystyle \frac{1}{(1\frac{q^2}{4M^2})}}[G_M^{p,n}(q^2)G_E^{p,n}(q^2)]`$
$$G_E^p(q^2)=\left(1\frac{q^2}{M_v^2}\right)^2$$
(7)
$$G_M^p(q^2)=(1+\mu _p)G_E^p(q^2),G_M^n(q^2)=\mu _nG_E^p(q^2);$$
$$G_E^n(q^2)=(\frac{q^2}{4M^2})\mu _nG_E^p(q^2)\xi _n;\xi _n=\frac{1}{1\lambda _n\frac{q^2}{4M^2}}$$
$$\mu _p=1.79,\mu _n=1.91,M_v=0.84GeV,\text{and}\lambda _n=5.6.$$
The isovector axial vector form factor $`F_A(Q^2)`$ is given by
$$F_A(Q^2)=\frac{F_A(0)}{(1\frac{q^2}{M_A^2})^2}$$
where $`M_A=1.032GeV`$; $`F_A(0)`$=-1.261
In a nuclear process the neutrons and protons are not free and their momenta are constrained by the Pauli principle which is implemented in this model by requiring that for neutrino reactions initial nucleon momentum $`pp_{F_n}`$ and final nucleon momentum $`p^{}=(|\stackrel{}{p}+\stackrel{}{q}|)>p_{F_p}`$ where $`p_{F_{n,p}}=[\frac{3}{2}\pi ^2\rho _{n,p}(r)]^{\frac{1}{3}}`$, are the local Fermi momenta of neutrons and protons at the interaction point in the nucleus defined in terms of their respective nuclear densities $`\rho _{n,p}(r)`$. These constraints are incorporated while performing the integration over the initial nucleon momentum in Eqn.(2) by replacing the energy conserving $`\delta `$-function in Eqn.(3) by $`\frac{1}{\pi }ImU_N(q_0,\stackrel{}{q})`$ where $`U_N(q_0,\stackrel{}{q})`$ is the Lindhard function corresponding to the particle hole(ph) excitations induced by the weak interaction process through W exchange shown in Fig.1(a). In the large mass limit of W boson i.e.$`M_W\mathrm{}`$, this Fig.1(a) is reduced to Fig.1(b) for which the imaginary part of the Lindhard function i.e. $`ImU_N(q_0,\stackrel{}{q})`$ is given by
$$ImU_N(q_0,\stackrel{}{q})=\frac{1}{2\pi }\frac{M_pM_n}{|\stackrel{}{q}|}\left[E_{F_1}A\right]\text{with}$$
(8)
$$q^2<0,E_{F_2}q_0<E_{F_1}\text{and}\frac{q_0+|\stackrel{}{q}|\sqrt{1\frac{4M^2}{q^2}}}{2}<E_{F_1}$$
where $`E_{F_1}`$ and $`E_{F_2}`$ are the local Fermi energy of initial and final nucleons and
$$A=Max[M_n,E_{F_2}q_0,\frac{q_0+|\stackrel{}{q}|\sqrt{1\frac{4M^2}{q^2}}}{2}]$$
The expression for the neutrino nuclear cross section $`\sigma ^A(E_\nu )`$, is then given by:
$`\sigma ^A(E_\nu )={\displaystyle \frac{4}{\pi }}{\displaystyle _{r_{min}}^{r_{max}}}r^2𝑑r{\displaystyle _{p_{l}^{}{}_{}{}^{min}}^{p_{l}^{}{}_{}{}^{max}}}p_{l}^{}{}_{}{}^{2}𝑑p_l{\displaystyle _1^1}d(cos\theta )`$
$`\times {\displaystyle \frac{1}{E_\nu E_l}}\overline{{\displaystyle }}{\displaystyle |T|^2ImU_N[E_\nu E_l,\stackrel{}{q}]}.`$ (9)
Moreover in the nucleus, the $`Q`$ value of the nuclear reaction and the Coulomb distortion of the final lepton in the electromagnetic field of the final nucleus should be taken into account. This is done by modifying the energy conserving $`\delta `$-function $`\delta (E_\nu E_l+E_nE_p)`$ in Eqn.(3) to $`\delta (E_\nu Q(E_l+V_c(r))+E_nE_p)`$ where $`V_c(r)`$ is the Coulomb energy of the produced lepton in the field of final nucleus and is given by
$$V_c(r)=ZZ^{}\alpha 4\pi (\frac{1}{r}_0^r\frac{\rho _p(r^{})}{Z}r_{}^{}{}_{}{}^{2}𝑑r^{}+_r^{\mathrm{}}\frac{\rho _p(r^{})}{Z}r^{}𝑑r^{})$$
(10)
This amounts to evaluation of Lindhard function in Eqn.(8) at $`(q_0(Q+V_c(r)),\stackrel{}{q})`$ instead of $`(q_0,\stackrel{}{q})`$. The implementation of this modification requires a judicious choice of $`Q`$ value for inclusive nuclear reactions in which many nuclear states are excited in iron. We have taken $`Q`$-value of 6.8MeV corresponding to the transition to lowest lying $`1^+`$ state in $`{}_{}{}^{56}Co`$ for $`\nu _l+^{56}Fel^{}+^{56}Co^{}`$ reaction and $`Q`$-value of 4.3MeV corresponding to the transition to lowest lying $`1^+`$ state in $`{}_{}{}^{56}Mn`$ for $`\overline{\nu _l}+^{56}Fel^++^{56}Mn^{}`$ reaction.
The inclusion of $`V_c(r)`$ to modify energy and corresponding momentum of the charged lepton in the Coulomb field of final nucleus in our model is equivalent to the treatment of Coulomb distortion effect in modified effective momentum approximation(MEMA). This approximation has been used in other calculations of charged current neutrino reactionsvolpe and electron scattering at higher energiesguisti .
With these modifications, the final expression for the quasielastic inclusive production from iron nucleus is given by
$`\sigma ^A(E_\nu )={\displaystyle \frac{4}{\pi }}{\displaystyle _{r_{min}}^{r_{max}}}r^2𝑑r{\displaystyle _{p_{l}^{}{}_{}{}^{min}}^{p_{l}^{}{}_{}{}^{max}}}p_{l}^{}{}_{}{}^{2}𝑑p_l{\displaystyle _1^1}d(cos\theta )`$
$`\times {\displaystyle \frac{1}{E_\nu E_l}}\overline{{\displaystyle }}{\displaystyle |T|^2ImU_N[q_0(Q+V_c(r)),\stackrel{}{q}]}.`$ (11)
It is well known that weak transition strengths are modified in the nuclear medium due to presence of strongly interacting nucleons. This modification of the weak transitions strength in the nuclear medium is taken into account by considering the propagation of particle-hole(ph) excitations in the medium. While propagating through the medium, the ph-excitations interact through the nucleon nucleon potential and create other particle hole and $`\mathrm{\Delta }`$-h excitations as shown in Fig.2.
The effect of these excitations are calculated in random phase approximation which is described in Ref.singh2 ; singh4 . The effect of nuclear medium on the renormalization of weak strengths is treated in a non-relativistic frame work. In leading order, the non-relativistic reduction of the weak hadronic current defined in Eqn.6, $`F_2(q^2)`$ term gives a spin-isospin transition operator $`\frac{\stackrel{}{\sigma }\times \stackrel{}{q}}{2M}\stackrel{}{\tau }`$ which is a transverse operator while the $`F_A(q^2)`$ term gives a spin-isospin transition operator $`\stackrel{}{\sigma }\stackrel{}{\tau }`$ which has a longitudinal as well as a transverse part. This representation of the transition operators in the longitudinal and transverse parts is useful when summing the diagrams in Fig.(2) in random phase approximation(RPA) to calculate $`|T|^2`$. While the charge coupling remains unchanged due to nuclear medium effects, the terms proportional to $`F_2^2`$ are affected by the transverse part of the nucleon nucleon potential while the terms proportional to $`F_A^2`$ are affected by transverse as well as longitudinal parts. The effect is to replace the terms like $`F_2^2`$, $`F_A^2`$, $`F_2F_A`$ etc. in the following mannersingh2 ; singh4
$`(F_2^2,F_2F_A)(F_2^2,F_2F_A){\displaystyle \frac{1}{|1U_NV_t|^2}}`$
$`F_A^2\left[{\displaystyle \frac{1}{3}}{\displaystyle \frac{1}{|1U_NV_l|^2}}+{\displaystyle \frac{2}{3}}{\displaystyle \frac{1}{|1U_NV_t|^2}}\right]`$ (12)
where $`V_l`$ and $`V_t`$ are the longitudinal and transverse parts of the nucleon nucleon potential calculated with $`\pi `$ and $`\rho `$ exchanges and modulated by Landau-Migdal parameter $`g^{}`$ to take into account the short range correlation effects and are given by
$`V_l(q)={\displaystyle \frac{f^2}{m_\pi ^2}}\left[{\displaystyle \frac{q^2}{q^2+m_\pi ^2}}\left({\displaystyle \frac{\mathrm{\Lambda }_\pi ^2m_\pi ^2}{\mathrm{\Lambda }_\pi ^2q^2}}\right)^2+g^{}\right],`$
$`V_t(q)={\displaystyle \frac{f^2}{m_\pi ^2}}\left[{\displaystyle \frac{q^2}{q^2+m_\rho ^2}}C_\rho \left({\displaystyle \frac{\mathrm{\Lambda }_{\rho }^{}{}_{}{}^{2}m_\rho ^2}{\mathrm{\Lambda }_{\rho }^{}{}_{}{}^{2}q^2}}\right)^2+g^{}\right]`$ (13)
with $`\mathrm{\Lambda }_\pi =1.3GeV`$, $`C_\rho =2.0`$, $`\mathrm{\Lambda }_\rho =2.5GeV`$, $`m_\pi `$ and $`m_\rho `$ are the pion and rho meson masses and $`g^{}`$ is taken to be $`0.7`$mukh . The effect of $`\mathrm{\Delta }`$h excitations are taken into account by including the Lindhard function $`U_\mathrm{\Delta }`$ for the $`\mathrm{\Delta }`$h excitations and replacing $`U_N`$ by $`U_N+U_\mathrm{\Delta }`$ in Eqn.(12). The complete expression for $`U_N`$ and $`U_\mathrm{\Delta }`$ used in our calculations are taken from oset1 . The different couplings for $`N`$ and $`\mathrm{\Delta }`$ to the nucleon are incorporated in $`U_N`$ and $`U_\mathrm{\Delta }`$ and then the same interaction strengths $`V_l`$ and $`V_t`$ are used for ph and $`\mathrm{\Delta }`$h excitationsgarcia -salcedo .
### II.2 Results
We present the numerical results for the total cross section for the quasielastic processes $`\nu _l(\overline{\nu _l})+^{56}Fel^{}(l^+)+^{56}Co^{}(^{56}Mn^{})`$ as a function of energy for neutrino and anti neutrino reactions on iron in the energy region relevant to the fully contained events of atmospheric neutrinos i.e. $`E_\nu <3GeV`$. The cross sections have been calculated using Eqn.(11) with the nuclear density $`\rho (r)`$ given by a two parameter Fermi densityvries : $`\rho (r)=\frac{\rho (0)}{1.+exp(\frac{rc}{z})}`$ with $`c=3.971fm`$, $`z=0.5935fm`$, $`\rho _n(r)=\frac{(AZ)}{A}\rho (r)`$ and $`\rho _p(r)=\frac{Z}{A}\rho (r)`$.
In Fig.3 we show the numerical results of $`\sigma (E)`$ vs $`E`$, for all flavors of neutrinos i.e. $`\nu _\mu `$, $`\overline{\nu }_\mu `$, $`\nu _e`$ and $`\overline{\nu }_e`$. The reduction due to nuclear effects is large at lower energies but becomes small at higher energies. The energy dependence of the cross sections for muon and electron type neutrinos are similar except for the threshold effects which are seen only at low energies $`(E_\nu <500MeV)`$. This reduction in $`\sigma `$ is due to Pauli blocking as well as due to the weak renormalization of transition strengths which have been separately shown in Fig.4(a) and Fig.4(b) for neutrinos and antineutrinos, where we also show the results in the Fermi gas model given by Llewellyn Smithsmith1 . We plot in Fig.4(a) and Fig.4(b) for neutrinos and antineutrinos, the reduction factor $`R=\frac{\sigma _{nuclear}(E)}{\sigma _{nucleon}(E)}`$ vs E where $`\sigma _{nuclear}(E)`$ is the cross section per neutron(proton) for neutrino(antineutrino) reactions in the nuclear medium. The solid lines show the reduction factor R when only the Pauli suppression is taken into account through the imaginary part of the Lindhard function given in Eqn.8. This is similar to the results of Llewellyn Smithsmith1 in Fermi gas model shown by dashed lines. In this model the total cross section $`\sigma `$ is calculated by using the formula $`\sigma =𝑑q^2R(q^2)(\frac{d\sigma }{dq^2})_{free}`$ where $`R(q^2)`$ describes the reduction in the cross section calculated in the Fermi gas model and includes the effect of Pauli suppression onlysmith1 . However, in our model we get further reduction due to renormalization of weak transition strengths in the nuclear medium when the effects of Fig.2(a) and Fig.2(b) are included. These are shown by dotted lines in Figs.4(a) and 4(b). We find that the reduction at higher energies $`(E>1GeV)`$ is around 20$`\%`$ for neutrinos and 40$`\%`$ for antineutrinos. It is worth noting that the energy dependence of the reduction due to nuclear medium effects is different for neutrinos and antineutrinos. This is due to the different renormalization of various terms like $`F_A^2`$, $`F_2F_A`$ and $`F_2^2`$ in $`|T|^2`$ which enter in different combinations for neutrino and antineutrino reactions. The results for $`\nu _e`$ and $`\overline{\nu }_e`$ cross section are respectively similar to $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ reactions except for the threshold effects and are not shown here.
In Fig.5 and Fig.6 we compare our results for $`\sigma (E)`$ with the results of some earlier experiments which contain nuclear targets like Carbonlsnd , Freonbonnetti -brunner , Freon-Propanepohl and Aluminumbelikov , where the experimental results for the deuteron targetsbaker are not included as they are not subject to the various nuclear effects discussed here. It should be kept in mind that the nuclear targets considered here (except for $`Br`$ in Freon) are lighter than $`Fe`$. Therefore, the reduction in the total cross section due to nuclear effects will be slightly overestimated. For example, for energies $`E_\nu 1GeV`$ the reduction in neutrino(antineutrino) cross section in case of $`{}_{}{}^{56}Fe`$ is $`5\%`$ more than the reduction in case of $`{}_{}{}^{12}C`$athar2 . In comparision to the neutrino(antineutrino) nuclear cross sections as obtained in the Fermi gas model of Llewellyn Smithsmith1 (shown by dashed lines in Fig.5 and Fig.6) we get a smaller result for these cross sections. This reduction in the total cross section leads to an improved agreement with the experimental results as compared to the Fermi gas model results specially for antineutrino reactions(Fig.6). It should be emphasized that the Fermi gas model has no specific mechanism to estimate the renormalization of weak transition strengths in nuclei while in our model this is incorporated by taking into account the RPA correlations.
In Fig.7(a) and Fig.7(b), we show the nuclear medium effects on the momentum and angular distributions i.e. $`\frac{d\sigma }{dp_l}`$ and $`\frac{d\sigma }{dcos\theta _l}`$ of leptons produced in $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ reactions. We find a large suppression in the results specially in the peak region of momentum and angular distributions. Quantitatively similar results are obtained for the case of $`\nu _e`$ and $`\overline{\nu }_e`$ reactions and are not shown here.
We have also studied the effect of Coulomb distortion in the momentum distribution of leptons but find no substantial effect around $`E_\nu =1.0GeV`$. These are found to affect the results only at low energies i.e. $`E_\nu <500MeV`$ where the peak is slightly shifted to lower momentum as shown in Figs.8(a) and 8(b).
## III Inelastic Production of Leptons
The inelastic production process of leptons is the process in which the production of leptons is accompanied by one pion (or more pions). There are many calculations of one pion production by neutrinos from free nucleonsadler but there are only a few calculations which discuss the nuclear effects in these processesadler1 -paschos . In this section we follow the method of Ref.ruso to estimate the nuclear effects and nuclear model dependence of inelastic production cross section of leptons induced by neutrinos from iron nuclei. The calculations are done assuming $`\mathrm{\Delta }`$-dominance of one pion production because the contribution of higher resonances in the energy region of atmospheric neutrinos leading to fully contained events is expected to be small.
### III.1 Formalism
The matrix element for neutrino production reaction of $`\mathrm{\Delta }`$ on proton targets leading to one pion events i.e.
$$\nu _l(k)+p(p)l^{}(k^{})+\mathrm{\Delta }^{++}(p^{})$$
(14)
is given by Eqn.5 where $`J^\mu `$ now defines the matrix element of the transition hadronic current between $`N`$ and $`\mathrm{\Delta }`$ states. The most general form of $`J_{CC}^\mu `$ is written asruso :
$`J_{cc}^\mu `$ $`=`$ $`\overline{\psi }_\alpha (p^{})[({\displaystyle \frac{C_3^V(q^2)}{M}}(g^{\alpha \mu }\mathit{}q^\alpha \gamma ^\mu )`$
$`+`$ $`{\displaystyle \frac{C_4^V(q^2)}{M^2}}(g^{\alpha \mu }qp^{}q^\alpha p^\mu )`$
$`+`$ $`{\displaystyle \frac{C_5^V(q^2)}{M^2}}(g^{\alpha \mu }qpq^\alpha p^\mu )+{\displaystyle \frac{C_6^V(q^2)}{M^2}}q^\alpha q^\mu )\gamma _5`$
$`+`$ $`({\displaystyle \frac{C_3^A(q^2)}{M}}(g^{\alpha \mu }\mathit{}q^\alpha \gamma ^\mu )+{\displaystyle \frac{C_4^A(q^2)}{M^2}}(g^{\alpha \mu }qp^{}q^\alpha p^\mu )`$
$`+`$ $`C_5^A(q^2)g^{\alpha \mu }+{\displaystyle \frac{C_6^A(q^2)}{M^2}}q^\alpha q^\mu )]u(p)`$
where $`\psi _\alpha (p^{})`$ and u(p) are the Rarita Schwinger and Dirac spinors for $`\mathrm{\Delta }`$ and nucleon of momenta $`p^{}`$ and $`p`$ respectively, $`q(=p^{}p=kk^{})`$ is the momentum transfer and $`C_i^V`$(i=3-6) are vector and $`C_i^A`$(i=3-6) are axial vector transition form factors. The vector form factors $`C_i^V`$(i=3-6) are determined by using the conserved vector current(CVC) hypothesis which gives $`C_6^V(q^2)`$=0 and relates $`C_i^V`$(i=3,4,5) to the electromagnetic form factors which are determined from photoproduction and electroproduction of $`\mathrm{\Delta }`$’s. Using the analysis of these experimentspaschos -dufner we take for the vector form factors
$`C_5^V=0,C_4^V={\displaystyle \frac{M}{M_\mathrm{\Delta }}}C_3^V,\text{and}`$
$`C_3^V(q^2)={\displaystyle \frac{2.05}{(1\frac{q^2}{M_V^2})^2}},M_V^2=0.54GeV^2`$ (16)
The axial vector form factor $`C_6^A(q^2)`$ is related to $`C_{5}^{}{}_{}{}^{A}(q^2)`$ using PCAC and is given by
$$C_6^A(q^2)=C_{5}^{}{}_{}{}^{A}(q^2)\frac{M^2}{m_{\pi }^{}{}_{}{}^{2}q^2}$$
(17)
The remaining axial vector form factor $`C_{i=3,4,5}^A(q^2)`$ are taken from the experimental analysis of the neutrino experiments producing $`\mathrm{\Delta }`$’s in proton and deuteron targetskitagaki -barish . These form factors are not uniquely determined but the following parameterizations give a satisfactory fit to the data.
$$C_{i=3,4,5}^A(q^2)=C_i^A(0)\left[1+\frac{a_iq^2}{b_iq^2}\right]\left(1\frac{q^2}{M_{A}^{}{}_{}{}^{2}}\right)^2$$
(18)
with $`C_3^A(0)=0,C_4^A(0)=0.3,C_5^A(0)=1.2,a_4=a_5=1.21,b_4=b_5=2GeV^2`$, $`M_A=1.28GeV`$. Using the hadronic current given in Eqn.(15), the energy spectrum of the outgoing leptons is given by
$$\frac{d^2\sigma }{dE_k^{}d\mathrm{\Omega }_k^{}}=\frac{1}{8\pi ^3}\frac{1}{MM^{}}\frac{k^{}}{E_\nu }\frac{\frac{\mathrm{\Gamma }(W)}{2}}{(WM^{})^2+\frac{\mathrm{\Gamma }^2(W)}{4.}}L_{\mu \nu }J^{\mu \nu }$$
(19)
where $`W=\sqrt{(p+q)^2}`$ and $`M^{}`$ is mass of $`\mathrm{\Delta }`$,
$$L_{\mu \nu }=k_\mu k_\nu ^{}+k_\mu ^{}k_\nu g_{\mu \nu }kk^{}+iϵ_{\mu \nu \alpha \beta }k^\alpha k^\beta ,$$
$$J^{\mu \nu }=\overline{\mathrm{\Sigma }}\mathrm{\Sigma }J^\mu J^\nu $$
and is calculated with the use of spin $`\frac{3}{2}`$ projection operator $`P^{\mu \nu }`$ defined as
$$P^{\mu \nu }=\underset{spins}{}\psi ^\mu \overline{\psi ^\nu }$$
and given by:
$`P^{\mu \nu }={\displaystyle \frac{\mathit{}^{}+M^{}}{2M^{}}}(g^{\mu \nu }{\displaystyle \frac{2}{3}}{\displaystyle \frac{p^\mu p^\nu }{M^2}}`$
$`+{\displaystyle \frac{1}{3}}{\displaystyle \frac{p^\mu \gamma ^\nu p^\nu \gamma _\mu }{M^{}}}{\displaystyle \frac{1}{3}}\gamma ^\mu \gamma ^\nu )`$ (20)
In Eqn.(19), the decay width $`\mathrm{\Gamma }`$ is taken to be an energy dependent P-wave decay width given by
$`\mathrm{\Gamma }(W)={\displaystyle \frac{1}{6\pi }}\left({\displaystyle \frac{f_{\pi N\mathrm{\Delta }}}{m_\pi }}\right)^2{\displaystyle \frac{M}{W}}|𝐪_{cm}|^3\mathrm{\Theta }(WMm_\pi )`$ (21)
where
$$|𝐪_{cm}|=\frac{\sqrt{(W^2m_\pi ^2M^2)^24m_\pi ^2M^2}}{2W}$$
and $`M`$ is the mass of nucleon. The step function $`\mathrm{\Theta }`$ denotes the fact that the width is zero for the invariant masses below the $`N\pi `$ threshold. $`|𝐪_{\mathrm{𝐜𝐦}}|`$ is the pion momentum in the rest frame of the resonance. When the reaction(14) takes place in the nucleus, the neutrino interacts with the nucleon moving inside the nucleus of density $`\rho (r)`$ with its corresponding momentum $`\stackrel{}{p}`$ constrained to be below its Fermi momentum. The produced $`\mathrm{\Delta }`$’s have no such constraints on their momentum. These $`\mathrm{\Delta }^{}s`$ decay through various decay channels in the medium. The most prominent decay mode is $`\mathrm{\Delta }N\pi `$ which produces pions. This decay mode in the nuclear medium is slightly inhibited due to Pauli blocking of the final nucleon momentum modifying the decay width $`\mathrm{\Gamma }`$ used in Eqn.(21). This modification of $`\mathrm{\Gamma }`$ due to Pauli blocking of nucleus has been studied in detail in electromagnetic and strong interactionsoset . The modified $`\mathrm{\Delta }`$ decay width i.e. $`\stackrel{~}{\mathrm{\Gamma }}`$ is written asoset :
$$\stackrel{~}{\mathrm{\Gamma }}=\frac{1}{6\pi }\left(\frac{f_{\pi N\mathrm{\Delta }}}{m_\pi }\right)^2\frac{M|𝐪_{cm}|^3}{W}F(k_F,E_\mathrm{\Delta },k_\mathrm{\Delta })\mathrm{\Theta }(WMm_\pi )$$
(22)
where $`F(k_F,E_\mathrm{\Delta },k_\mathrm{\Delta })`$ is the Pauli correction factor given by
$$F(k_F,E_\mathrm{\Delta },k_\mathrm{\Delta })=\frac{k_\mathrm{\Delta }|𝐪_{cm}|+E_\mathrm{\Delta }E_{p}^{}{}_{cm}{}^{}E_FW}{2k_\mathrm{\Delta }|𝐪_{}^{}{}_{cm}{}^{}|}$$
(23)
where $`k_F`$ is the Fermi momentum, $`E_F=\sqrt{M^2+k_F^2}`$, $`k_\mathrm{\Delta }`$ is the $`\mathrm{\Delta }`$ momentum and $`E_\mathrm{\Delta }=\sqrt{W+k_\mathrm{\Delta }^2}`$.
Moreover, in the nuclear medium there are additional decay channels now open due to two body and three body absorption processes like $`\mathrm{\Delta }NNN`$ and $`\mathrm{\Delta }NNNNN`$ through which $`\mathrm{\Delta }^{}s`$ disappear in the nuclear medium without producing a pion while a two body $`\mathrm{\Delta }`$ absorption process like $`\mathrm{\Delta }N\pi NN`$ gives rise to some more pions. These nuclear medium effects on $`\mathrm{\Delta }`$ propagation are included by using a $`\mathrm{\Delta }`$ propagator in which the $`\mathrm{\Delta }`$ propagator is written in terms of $`\mathrm{\Delta }`$ self energy $`\mathrm{\Sigma }_\mathrm{\Delta }`$. This is done by using a modified mass and width of $`\mathrm{\Delta }`$ in nuclear medium i.e. $`M_\mathrm{\Delta }M_\mathrm{\Delta }+Re\mathrm{\Sigma }_\mathrm{\Delta }`$ and $`\stackrel{~}{\mathrm{\Gamma }}\stackrel{~}{\mathrm{\Gamma }}Im\mathrm{\Sigma }_\mathrm{\Delta }`$. There are many calculations of $`\mathrm{\Delta }`$ self energy $`\mathrm{\Sigma }_\mathrm{\Delta }`$ in the nuclear medium oset -oset3 and we use the results of oset , where the density dependence of real and imaginary parts of $`\mathrm{\Sigma }_\mathrm{\Delta }`$ are parametrized in the following form:
$$Re\mathrm{\Sigma }_\mathrm{\Delta }=40\frac{\rho }{\rho _0}MeV\text{and}$$
$$Im\mathrm{\Sigma }_\mathrm{\Delta }=C_Q\left(\frac{\rho }{\rho _0}\right)^\alpha +C_{A2}\left(\frac{\rho }{\rho _0}\right)^\beta +C_{A3}\left(\frac{\rho }{\rho _0}\right)^\gamma $$
(24)
In Eqn.24, the term with $`C_Q`$ accounts for the $`\mathrm{\Delta }N\pi NN`$ process, the term with $`C_{A2}`$ for two-body absorption process $`\mathrm{\Delta }NNN`$ and the term with $`C_{A3}`$ for three-body absorption process $`\mathrm{\Delta }NNNNN`$. The coefficients $`C_Q`$, $`C_{A2}`$, $`C_{A3}`$, $`\alpha `$, $`\beta `$ and $`\gamma `$ ($`\gamma =2\beta `$) are parametrized in the range $`80<T_\pi <320MeV`$ (where $`T_\pi `$ is the pion kinetic energy) as oset
$`C_i(T_\pi )`$ $`=`$ $`ax^2+bx+c,\text{for}x={\displaystyle \frac{T_\pi }{m_\pi }}`$ (25)
The values of coefficients $`a`$, $`b`$ and $`c`$ are given in Table-I. taken from ref.oset .
With these modifications, which incorporate the various nuclear medium effects on $`\mathrm{\Delta }`$ propagation, the cross section is now written as
$`\sigma `$ $`=`$ $`{\displaystyle \frac{d𝐫}{8\pi ^3}\frac{d𝐤^{}}{E_\nu E_l}\frac{1}{MM^{}}}`$ (26)
$`\times {\displaystyle \frac{\frac{\stackrel{~}{\mathrm{\Gamma }}}{2}Im\mathrm{\Sigma }_\mathrm{\Delta }}{(WM^{}Re\mathrm{\Sigma }_\mathrm{\Delta })^2+(\frac{\stackrel{~}{\mathrm{\Gamma }}}{2.}Im\mathrm{\Sigma }_\mathrm{\Delta })^2}}`$
$`\times \left[\rho _p(𝐫)+{\displaystyle \frac{1}{3}}\rho _n(𝐫)\right]L_{\mu \nu }J^{\mu \nu }`$
The factor $`\frac{1}{3}`$ in front of $`\rho _n`$ comes due to suppression of charged pion production from neutron targets i.e. $`\nu _l+nl^{}+\mathrm{\Delta }^+l^{}+n+\pi ^+`$, as compared to the charged pion production from the proton target i.e.$`\nu _l+pl^{}+\mathrm{\Delta }^{++}l^{}+p+\pi ^+`$, in the nucleus. In case of antineutrino reactions $`\rho _p+\frac{1}{3}\rho _n`$ is replaced by $`\rho _n+\frac{1}{3}\rho _p`$.
### III.2 Results
In this section we present results of inelastic lepton production cross section due to $`\mathrm{\Delta }`$h excitations in iron induced by the charged current neutrino interactions using Eqn.(26). In Fig.9(a) and Fig.9(b), we show the results $`\sigma (E_\nu )E_\nu `$ for $`\nu _e(\overline{\nu }_e)`$ and $`\nu _\mu (\overline{\nu }_\mu )`$ for the inelastic production of lepton accompanied with $`\mathrm{\Delta }`$ and compare this with the cross section for the quasielastic production of leptons discussed in section-II. We see that for $`E_\nu 1.4GeV`$ the lepton production cross section through quasielastic and inelastic production processes are comparable. For energies $`E_\nu 1.4GeV`$ where the atmospheric neutrino energies are important for fully contained events the major contribution comes from the quasielastic events.
The effect of nuclear effects on the $`\mathrm{\Delta }`$ production are shown in Fig.10 for $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ reactions. The results for the $`\nu _e`$ and $`\overline{\nu }_e`$ are similar to Fig.10 and are not shown here. We see that the nuclear medium effects reduce the $`\mathrm{\Delta }`$ production cross section by 5-10$`\%`$.
In Fig.11(a) and Fig.11(b), we show the momentum distribution $`d\sigma /dp`$ and angular distribution $`d\sigma /dcos\theta `$ for muon type neutrinos where we also show the effects of nuclear effects. The nuclear medium effects reduce the cross sections mainly in the peak region of the momentum and angular distributions for muons. The results for the momentum distribution and angular distribution for electrons is similar to the muon distributions. We show our final results on momentum and angular distributions for all charged leptons i.e. $`e^{},e^+,\mu ^{}`$ and $`\mu ^+`$ produced in neutrino and antineutrino reactions with $`\nu _e`$, $`\nu _\mu `$, $`\overline{\nu }_e`$ and $`\overline{\nu }_\mu `$ for $`E_\nu =1GeV`$ in Fig.12(a) and Fig.12(b) in $`{}_{}{}^{56}Fe`$ nuclei.
Here we will like to make some comments about the lepton events produced through the $`\mathrm{\Delta }`$ excitation. The $`\mathrm{\Delta }`$ excitation process gives rise to leptons accompanied by one pion events produced by $`\mathrm{\Delta }N\pi `$ and $`\mathrm{\Delta }N\pi NN`$ processes in the nuclear medium. The pions produced through these processes will undergo secondary nuclear interactions like multiple scattering and possible absorption in iron nuclei while passing through the nucleus and an appropriate model has to be used for their description. Models developed by Salcedo et al.salcedo and Paschos et al.paschos have studied these effects but we do not consider these effects here as they are beyond the scope of this paper. The $`\mathrm{\Delta }`$ excitation process in the nuclear medium also gives rise to quasielastic like events through two body and three body absorption processes like $`\mathrm{\Delta }NNN`$ and $`\mathrm{\Delta }N\mathrm{\Delta }NN`$ where only leptons are present in the final states. The quasielastic-like lepton events are discussed by Kim et al.kim but no quantitative estimates have been made. We have in an earlier paperruso discussed this process only qualitatively but make a quantitative estimate of these events in this paper. We find that around $`E_\nu =1GeV`$ the contribution of these quasielastic like events is about $`1012\%`$ but its effect on the flux averaged production of leptons for atmospheric neutrinos is not very significant. This is discussed in some detail in the next section.
## IV Lepton Production by Atmospheric neutrinos
The energy dependences of the quasielastic and inelastic lepton production cross sections described in sections II and III have been used to calculate the lepton production by atmospheric neutrinos after averaging over the neutrino flux corresponding to the two sites of Soudan and Gransasso, where iron based detectors are being used. There are quite a few calculations of atmospheric neutrino fluxes at these two sites. We use the angle averaged fluxes calculated by Honda et al.honda and Barr et al.barr for the Soudan site and the fluxes of Barr et al.barr and Plyaskinplyaskin for the Gransasso site to calculate the flux averaged cross section $`<\sigma >`$ and also the momentum and the angular distributions $`<\frac{d\sigma }{dp_l}>`$ and $`<\frac{d\sigma }{dcos\theta _l}>`$ for leptons produced by $`\nu _e`$, $`\overline{\nu }_e`$, $`\nu _\mu `$ and $`\overline{\nu }_\mu `$.
### IV.1 Flux averaged Momentum and Angular distributions
In this section we present the numerical results for the flux averaged momentum distributions $`<\frac{d\sigma }{dp_l}>`$ and angular distributions $`<\frac{d\sigma }{dcos\theta _l}>`$ for various leptons produced from $`\nu _e`$, $`\overline{\nu }_e`$, $`\nu _\mu `$, $`\overline{\nu }_\mu `$ at the atmospheric neutrino sites of Soudan and Gransasso. These leptons are produced through the quasielastic as well as inelastic processes. The quasielastic production has been discussed in section-II and the inelastic production has been discussed in section-III. The inelastic production of leptons is accompanied by pions. However, in the nuclear medium where the production of $`\mathrm{\Delta }`$ is followed by $`\mathrm{\Delta }NNN`$, $`\mathrm{\Delta }NNNNN`$ absorption processes, the leptons are produced without pions. These are quasielastic like events. We show the momentum distribution of the leptons produced in iron nuclei corresponding to one pion and quasielastic like events for atmospheric neutrinos relevant to Soudan and Gransasso sites. We show it for $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ for the Soudan site in Figs. 13(a) and 13(b) and for the Gransasso site in Figs.14(a) and 14(b) corresponding to the flux of Barr et al.barr . We see that at both sites, the production cross section of quasielastic-like events is quite small.
We now present our final results for the momentum distribution and angular distribution for various leptons produced by atmospheric neutrinos at Soudan and Gransasso sites in Figs. 15-18. In these figures separate contributions from the quasielastic and inelastic processes to the momentum and angular distributions of charged leptons are shown explicitly. The quasielastic events are those where only a charge lepton is produced in the final state either by a quasielastic process described in section-II or by an inelastic process of $`\mathrm{\Delta }`$-production followed by its subsequent absorption in the nuclear medium as described in section-III. The inelastic events are those events in which a charged lepton in the final state is accompanied by a charged pion as a decay product of deltas excited in the nuclear medium.
We show the results for the momentum distribution $`<d\sigma /dp_l>p_l(l=e^{},e^+,\mu ^+,\mu ^{})`$ for the atmospheric neutrino fluxes of Barr et al.barr at Soudan site in Fig.15 and at Gransasso site in Fig.16. We see that in all cases the major contributions to the charged lepton production comes from the quasielastic processes induced by neutrinos(solid line). The contribution from the quasielastic processes induced by antineutrinos(dotted line) and inelastic processes induced by neutrinos(dashed line) is small and is around $`20\%`$ in the peak region. The contribution due to inelastic processes induced by antineutrinos is very small over the whole region (dashed-dotted lines).
The peak in quasielastic $`\nu _\mu `$ reactions occurs around $`p_l200MeV`$. The peaks in the inelastic $`\nu _\mu `$ and quasielastic $`\overline{\nu _\mu }`$ reactions are slightly shifted towards lower energies. The momentum distribution of the leptons for the quasielastic reaction is peaked more sharply than the momentum distribution of the inelastic reaction.
We have also presented the numerical results for angular distributions of leptons $`<\frac{d\sigma }{dcos\theta _l}>vscos\theta `$ for the atmospheric neutrino fluxes of Barr et al.barr for the Soudan site in Fig.17 and for the Gransasso site in Fig.18. The inelastic lepton production accompanied by pions (dashed line for $`\mu ^{}(e^{})`$ production and dashed-dotted line for $`\mu ^+(e^+)`$ production) and the quasielastic lepton production events (solid line for $`\mu ^{}(e^{})`$ production and dotted line for $`\mu ^+(e^+)`$ production) have been explicitly shown in these figures. We see from these figures that the lepton production cross sections are forward peaked in all cases. The inelastic distribution due to pion production are slightly more forward peaked than the quasielastic distribution. The contribution of the inelastic lepton events are small compared to the quasielastic events and the angular distribution of the flux averaged cross sections are dominated by the quasielastic events.
At Soudan site, we have also studied the momentum and angular distributions for the atmospheric neutrino fluxes of Honda et al.honda . We find that the momentum and angular distributions for the production of muons are similar to the distributions obtained for the flux of Barr et al.barr . In the case of electron production, the use of the flux of Honda et al.honda , gives a slightly smaller value for $`<d\sigma /dp_l>`$ and $`<\frac{d\sigma }{dcos\theta _l}>`$ for electrons as compared to the flux of Barr et al.barr .
Similarly at Gransasso site, the momentum and angular distributions for the atmospheric neutrino fluxes of Plyaskinplyaskin have also been studied. Here, also we find that the momentum and angular distributions for the production of muons are similar to the distributions obtained for the flux of Barr et al.barr . In the case of electron production, the use of the flux of Plyaskinplyaskin , gives a slightly smaller value for $`<d\sigma /dp_l>`$ and $`<\frac{d\sigma }{dcos\theta _l}>`$ for electrons as compared to the flux of Barr et al.barr .
We further find that for the flux of Barr et al.barr , the lepton production cross section at the Soudan site is slightly larger than the Gransasso site for muons as well as for electrons.
### IV.2 Flux averaged Total cross sections and lepton yields
The total cross sections for the production of leptons and its energy dependence have been discussed in section II and section III for quasielastic and inelastic reactions. In this section we calculate the lepton yields $`Y_l`$ for lepton of flavor $`l`$ which we define as
$$Y_l=\mathrm{\Phi }_{\nu _l}\sigma (E_{\nu _l})𝑑E_{\nu _l}$$
where, $`\mathrm{\Phi }_{\nu _l}`$ is the atmospheric neutrino flux of $`\nu _l`$ and $`\sigma (E_{\nu _l})`$ is the total cross section for neutrino $`\nu _l`$ of energy $`E_{\nu _l}`$. We calculate this yield separately for the quasielastic and inelastic events. We define a relative yield of muon over electron type events by $`R`$ as $`R=R_{\mu /e}=\frac{Y_\mu +Y_{\overline{\mu }}}{Y_e+Y_{\overline{e}}}`$ for quasielastic and inelastic events and present our results in Table-II. We study the nuclear model dependence as well as the flux dependence of the relative yield $`R`$.
The results for $`R`$ are presented separately for quasielastic events, inelastic events and the total events in Table-II. For quasielastic events $`\nu _l(\overline{\nu _l})+^{56}Fel^{}(l^+)+X`$, the results are presented for the case of free nucleon by $`R_{FN}`$, for the nuclear case with Fermi gas model description of Llewellyn Smithsmith1 by $`R_{FG}`$ and for the case of nuclear effects within our model by $`R_{NM}`$. We see that there is practically no nuclear model dependence on the value of $`R`$ (compare the values of $`R_{NM}`$, $`R_{FG}`$ and $`R_{FN}`$ for the same fluxes at each site). This is also true for the inelastic production of leptons i.e. $`\nu _l(\overline{\nu _l})+^{56}Fel^{}(l^+)+\pi ^+(\pi ^{})+X`$ for which the ratio for free nucleon case (denoted by $`R_{\mathrm{\Delta }F}`$) and the ratio for the nuclear case in our model (denoted by $`R_{\mathrm{\Delta }N}`$) are presented in Table II. It is, therefore, concluded that there is no appreciable nuclear model dependence on the ratio of total lepton yields for the production of muons and electrons(compare the values of $`R_F`$ and $`R_N`$, where $`R_F`$ shows the ratio of total lepton yields for muon and electron for the case of free nucleon and $`R_N`$ shows the ratio of total yield for muon and electron for the case of nucleon in the nuclear medium).
However, there is some dependence of the ratio R on the atmospheric neutrino fluxes. The flux dependence of $`R`$ can be readily seen from Table-II, for the two sites of Soudan and Gransasso. At the Gransasso site, we see that there is 4-5$`\%`$ difference in the value of $`R_N`$ for the total lepton yields for the fluxes of Barr et al.barr and Plyaskinplyaskin . At Soudan site, the results for the fluxes of Honda et al.honda and Barr et al.barr are within 4-5$`\%`$ but the flux calculation of Plyaskinplyaskin gives a result which is about 10-11$`\%`$ smaller than the results of Honda et al.honda and 7-8$`\%`$ smaller than the results of Barr et al.barr . The flux dependence is mainly due to the quasielastic events. This should be kept in mind while using the flux of Plyaskinplyaskin for making any analysis of the neutrino oscillation experiments.
In Table-III, we present a quantitative estimate of the relative yield of inelastic events $`r_l`$ defined by $`r_l=\frac{Y_{l}^{}{}_{}{}^{\mathrm{\Delta }}}{Y_l}`$, where $`Y_{l}^{}{}_{}{}^{\mathrm{\Delta }}=Y_{l}^{}{}_{}{}^{\mathrm{\Delta }}+Y_{\overline{l}}^{}{}_{}{}^{\mathrm{\Delta }}`$ is the lepton yield due to the inelastic events and $`Y_l`$ is the total lepton yield due to the quasielastic and inelastic events i.e.$`Y_l=Y_{l}^{}{}_{}{}^{q.e.}+Y_{\overline{l}}^{}{}_{}{}^{q.e.}+Y_{l}^{}{}_{}{}^{\mathrm{\Delta }}+Y_{\overline{l}}^{}{}_{}{}^{\mathrm{\Delta }}`$. The relative yield for the case of free nucleon is shown by $`r_l(F)`$ and for the case with the nuclear effects in our model is shown by $`r_l(N)`$. We see that for free nucleons, the relative yield of the inelastic events due to $`\mathrm{\Delta }`$ excitation is in the range of 12-15$`\%`$ for various fluxes at the two sites. The ratio is approximately same for electrons and muons. When the nuclear effects are taken into account this becomes 19-22$`\%`$. This is due to different nature of the effect of nuclear structure on the quasielastic and inelastic production cross scetions which gives a larger reduction in the cross section for the quasielastic case as compared to the inelastic case. This quantitative estimate of $`r_l(N)`$ in iron may be compared with the results in oxygen where the experimental results at Kamiokande give a value of 18$`\%`$kajita11 .
## V Conclusions
We have studied the charged lepton production in iron induced by atmospheric neutrinos at the experimental sites of Soudan and Gransasso. The energy dependence of the total cross sections for the quasielastic and inelastic processes have been calculated in a nuclear model which takes into account the effect of Pauli principle, Fermi motion effects and the renormalization of weak transition strengths in nuclear medium. The inelastic process has been studied through the $`\mathrm{\Delta }`$ dominance model which incorporates the modification of mass and width of $`\mathrm{\Delta }`$ resonance in the nuclear medium. The numerical results for the momentum and angular distributions of the charged lepton production cross section have been presented for muons and electrons. The relative yield of muon to electron production has been studied. In addition to the nuclear model dependence, the flux dependence of the total yields, momentum distribution $`\frac{d\sigma }{dp_l}`$ and $`\frac{d\sigma }{dcos\theta _l}`$ have been also studied. In the following we conclude this paper by summarizing our main results:
1.There is a large reduction due to nuclear effects in the total cross section for quasielastic production cross section specially at lower energies(40-50$`\%`$ around 200MeV) and the reduction becomes smaller at higher energies(15-20$`\%`$ around 500MeV ).
2.For quasielastic reactions we find a larger reduction in the total cross section as compared to the Fermi gas model. The energy dependence of this reduction in cross section at low energies ($`E<500MeV`$) is found to be different for neutrino and antineutrino reactions.
3.The inelastic production of leptons where a charged lepton is accompanied by a pion becomes comparable to the quasielastic production of leptons for $`E_\nu 1.4GeV`$ and increases with increase of energy. The effect of nuclear medium on the inelastic production cross section (in absence of pion re-scattering effects) in not too large ($`10\%`$).
4. For quasielastic reactions the effect of Coulomb distortion of the final lepton in the total cross section is small except at very low energies ($`E<500MeV`$) and becomes negligible when averaged over the flux of atmospheric neutrinos.
5.The flux averaged momentum distribution of leptons produced by atmospheric neutrinos is peaked around the momentum $`p_l200MeV`$ for electrons and muons. The peak for the quasielastic production is sharper than the inelastic production. The inelastic production of leptons contributes about 20$`\%`$ to the total production of leptons.
6.The flux averaged angular distribution of leptons for atmospheric neutrinos for quasielastic as well as inelastic production is sharply peaked in the forward direction.
7.There is a very little flux dependence on the relative yield of muons and electrons at the site of Gransasso. However, at the Soudan site, the atmospheric flux as determined by Plyaskinplyaskin gives a value of relative yield which is smaller than the relative yield using the flux of Honda et al.honda and Barr et al.barr .
8.The nuclear model dependence of the relative yield of muons to electrons is negligible, even though the individual yields for muons and electrons are reduced with the inclusion of nuclear effects specially for the quasielastic production.
## VI Acknowledgment
The work is financially supported by the Department of Science and Technology, Government of India under the grant DST Project No. SP/S2K-07/2000. One of the authors (S.Ahmad) would loke to thank CSIR for financial support.
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# 𝐵→𝜒_{𝑐1}(1𝑃,2𝑃)𝐾 decays in QCD factorization and X(3872)
## Abstract
$`B\chi _{c1}(1P,2P)K`$ decays are studied in QCD factorization by treating charmonia as nonrelativistic bound states. No infrared divergences exist in the vertex corrections, while the logarithmic end-point singularities in the hard spectator corrections can be regularized by a momentum cutoff. Within certain uncertainties we find that the $`B\chi _{c1}(2P)K`$ decay rate can be comparable to $`B\chi _{c1}(1P)K`$, and get $`Br(B^0\chi _{c1}^{}K^0)=Br(B^+\chi _{c1}^{}K^+)2\times 10^4`$. This might imply a possible interpretation for the newly discovered X(3872) that this state has a dominant $`J^{PC}=1^{++}(2P)`$ $`c\overline{c}`$ component but mixed with a substantial $`D^0\overline{D}^0+D^0\overline{D}^0`$ continuum component.
The naively factorizable decay BSW $`B\chi _{c1}K`$ was studiedChao03 in the QCD factorization approach BBNS in which the nonfactorizable vertex and spectator corrections were also estimated, but the numerical results were four times smaller than experimental data. Recently, these decays were also studied in the PQCD approach Li . In both the above approaches, light-cone distribution amplitudes(LCDAs) were used to describe $`\chi _{c1}`$. As argued in Ref. Chao04 , a more appropriate description of charmonium is the nonrelativistic (NR) wave functions which can be expanded in terms of the relative momentum $`q`$ between charm and anticharm quarks. This argument is based on the nonrelativistic nature of heavy quarkonium BBL1 . With careful studies, we find that the two descriptions (i.e.LCDAs and NR) are equivalent for the S-wave charmonium states (see,e.g. Chao02 ), but in the case of P-wave states the light-cone descriptions lose some important contributions in the leading-twist approximation. This is not surprising since $`q`$ can be neglected in S-wave states, but cannot be neglected for P-wave states even in leading order approximation.
On the phenomenological hand, the study of $`B\chi _{c1}(2P)K`$ may help clarify the nature of the recently discovered resonance $`X(3872)`$Belle03 , since the measurements for $`X(3872)`$ favor $`J^{PC}=1^{++}`$Belle05 and hence $`\chi _{c1}(2P)`$ becomes one of the possible assignments for it. On the other hand, aside from the conventional charmoniumBG ; chao1 , a loosely bound S-wave molecule of $`D^0\overline{D}^0+D^0\overline{D}^0`$ has been suggested for X(3872)Tornqvist ; Braaten .
Motivated by the above considerations, in this paper we study the decays $`B\chi _{c1}(1P,2P)K`$ within the framework of QCD factorization by treating the charmonia $`\chi _{c1}(1P,2P)`$ as nonrelativistic bound states with $`m_c/m_b`$ taken to be a fixed value in the heavy $`b`$ quark limit. We will estimate the production rate of $`\chi _{c1}(2P)`$ and argue that the X(3872) may be dominated by the $`\chi _{c1}(2P)`$ charmonium but mixed with some $`D^0\overline{D}^0+D^0\overline{D}^0`$ continuum component.
In the non-relativistic bound-state picture, charmonium can be described by the color-singlet NR wave function. Let $`p`$ be the total momentum of the charmonium and $`2q`$ be the relative momentum between $`c`$ and $`\overline{c}`$ quarks, then $`v^24q^2/p^20.25`$ can be treated as a small expansion parameter BBL1 . For P-wave charmonium $`\chi _{c1}`$, because the wave function at the origin $`_P(0)=0`$, which corresponds to the zeroth order in $`q`$, we must expand the amplitude to first order in $`q`$. Thus we have
$`(B\chi _{c1}K)={\displaystyle \underset{L_z,S_z}{}}1L_z;1S_z|1J_z{\displaystyle \frac{\mathrm{d}^4q}{(2\pi )^3}q_\alpha \delta (q^0)}`$
$`\times \psi _{1M}^{}(q)\mathrm{Tr}[𝒪^\alpha (0)P_{1S_z}(p,0)+𝒪(0)P_{1S_z}^\alpha (p,0)],`$ (1)
where $`𝒪(q)`$ represent the rest of the decay amplitudes and $`P_{1S_z}(p,q)`$ is the spin-triplet projection operator, and $`𝒪^\alpha ,P^\alpha `$ stand for the derivatives of $`𝒪,P`$ with respect to the relative momentum $`q_\alpha `$ Chao04 . The amplitudes $`𝒪(q)`$ can be further factorized as product of $`BK`$ form factors and hard kernel or as the convolution of a hard kernel with light-cone wave functions of B meson and K meson, within QCD factorization approach.
After $`q^0`$ is integrated out, the integral in Eq. ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)) is proportional to the derivative of the P-wave wave function at the origin by
$`{\displaystyle \frac{\mathrm{d}^3q}{(2\pi )^3}q^\alpha \psi _{1M}^{}(q)}=i\epsilon ^\alpha (L_z)\sqrt{{\displaystyle \frac{3}{4\pi }}}_P^{^{}}(0),`$ (2)
where $`\epsilon ^\alpha (L_z)`$ is the polarization vector of an angular momentum-1 system and the value of $`_P^{^{}}(0)`$ for charmonia can be found in, e.g., Ref. Quig .
In contrast to the NR description of $`\chi _{c1}`$, the K-meson is described by LCDAs BBNS :
$`K(p^{})|\overline{s}_\beta (z_2)d_\alpha (z_1)|0`$ $`=`$ (3)
$`{\displaystyle \frac{if_K}{4}}{\displaystyle _0^1}dxe^{i(yp^{}z_2+\overline{y}p^{}z_1)}\{/p^{}\gamma _5\varphi _K(y)\}_{\alpha \beta },`$
where $`y`$ and $`\overline{y}=1y`$ are the momentum fractions of the $`s`$ and $`\overline{d}`$ quarks inside the $`K`$ meson respectively, and $`\varphi _K(x)`$ is the leading twist LCDA of K-meson. The masses of light quarks and $`K`$ meson are neglected in heavy quark limit.
The effective Hamiltonian for $`B\chi _{c1}K`$ reads BBL
$`_{\mathrm{eff}}={\displaystyle \frac{G_F}{\sqrt{2}}}(V_{cb}V_{cs}^{}(C_1𝒪_1+C_2𝒪_2)V_{tb}V_{ts}^{}{\displaystyle \underset{i=3}{\overset{6}{}}}C_i𝒪_i),`$ (4)
where $`G_F`$ is the Fermi constant, $`C_i`$ are the Wilson coefficients and $`V_{q_1q_2}`$ are the CKM matrix elements. The relevant 4-fermion operators $`𝒪_i`$ can be found in Chao04 .
According to BBNS all nonfactorizable corrections are due to Fig.1. These corrections, with operators $`𝒪_i`$ inserted, contribute to the amplitude $`𝒪(q)`$ in Eq. ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)), where the external lines of charm and anti-charm quarks have been truncated. Taking nonfactorizable corrections in Fig.1 into account, the decay amplitude for $`B\chi _{c1}K`$ in QCD factorization is written as
$`i`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}\left[V_{cb}V_{cs}^{}a_2V_{tb}V_{ts}^{}(a_3a_5)\right]`$ (5)
$`\times 12i\sqrt{{\displaystyle \frac{2}{\pi M}}}_P^{^{}}(0)ϵ^{}p_BF_1(M^2),`$
where $`ϵ`$ is polarization vector of $`\chi _{c1}`$ Polar . Here $`F_1`$ is the $`BK`$ form factor and we have used the relation $`F_0(M^2)/F_1(M^2)=1z`$ Cheng , with $`z=M^2/m_B^24m_c^2/m_b^2`$ and $`M`$ is the mass of $`\chi _{c1}`$, to simplify the structure of (5).
The coefficients $`a_i`$ ($`i=2,3,5`$) in the naive dimension regularization(NDR) scheme are given by
$`a_2`$ $`=`$ $`C_2+{\displaystyle \frac{C_1}{N_c}}+{\displaystyle \frac{\alpha _s}{4\pi }}{\displaystyle \frac{C_F}{N_c}}C_1\left(18+12\mathrm{ln}{\displaystyle \frac{m_b}{\mu }}+f_I+f_{II}\right),`$
$`a_3`$ $`=`$ $`C_3+{\displaystyle \frac{C_4}{N_c}}+{\displaystyle \frac{\alpha _s}{4\pi }}{\displaystyle \frac{C_F}{N_c}}C_4\left(18+12\mathrm{ln}{\displaystyle \frac{m_b}{\mu }}+f_I+f_{II}\right),`$
$`a_5`$ $`=`$ $`C_5+{\displaystyle \frac{C_6}{N_c}}{\displaystyle \frac{\alpha _s}{4\pi }}{\displaystyle \frac{C_F}{N_c}}C_6\left(6+12\mathrm{ln}{\displaystyle \frac{m_b}{\mu }}+f_I+f_{II}\right),`$ (6)
where $`C_F=(N_c^21)/(2N_c)`$ and $`\mu `$ is the QCD renormalization scale.
The function $`f_I`$ is calculated from the four vertex correction diagrams (a, b, c, d) in Fig.1 and reads
$`f_I`$ $`=`$ $`{\displaystyle \frac{2z}{2z}}{\displaystyle \frac{4z\mathrm{log}(4)}{2z}}{\displaystyle \frac{4z^2\mathrm{log}(z)}{\left(1z\right)\left(2z\right)}}`$ (7)
$`+{\displaystyle \frac{4\left(32z\right)\left(1z\right)\left(\mathrm{log}(1z)i\pi \right)}{\left(2z\right)^2}},`$
We find that the infrared divergences are canceled between diagrams (a) and (b), (c) and (d) respectively in Fig.1. On the other hand, this function is different from that in Eq. (11) of Ref. Chao03 even when a nonrelativistic limit wave function $`\varphi _{\chi _{c1}}^{NR}(u)=\delta (u1/2)`$ is adopted, as we have mentioned.
For the two spectator correction diagrams (e, f) in Fig.1, the off-shellness of the gluon is natural to be associated with a scale $`\mu _h\sqrt{m_b\mathrm{\Lambda }_{\mathrm{QCD}}}`$, rather than $`\mu _hm_b`$. Following Ref. BBNS , we choose $`\mu =\sqrt{m_b\mathrm{\Lambda }_h}1.4`$ Gev with $`\mathrm{\Lambda }_h=0.5`$ Gev in calculating the hard spectator function $`f_{II}`$ and then, in the leading twist approximation, we get
$`f_{II}`$ $`=`$ $`{\displaystyle \frac{\alpha _\mathrm{s}(\mu _h)C_i(\mu _h)}{\alpha _\mathrm{s}(\mu )C_i(\mu )}}{\displaystyle \frac{8\pi ^2}{N_c}}{\displaystyle \frac{f_Kf_B}{F_1(M^2)m_B^2}}{\displaystyle \frac{1}{1z}}`$ (8)
$`\times {\displaystyle _0^1}d\xi {\displaystyle \frac{\varphi _B(\xi )}{\xi }}{\displaystyle _0^1}dy{\displaystyle \frac{\varphi _K(y)}{y}}[1+{\displaystyle \frac{z}{y(1z)}}],`$
where $`\xi `$ is the momentum fraction of the spectator quark in the $`B`$ meson and $`C_i(\mu _h)`$ ($`i=1,4,6`$) are the NLO Wilson coefficients which can be evaluated by the renormalization group approach BBL .
The spectator contribution depends on the wave function $`\varphi _B`$ through the integral
$`{\displaystyle _0^1}𝑑\xi {\displaystyle \frac{\varphi _B(\xi )}{\xi }}{\displaystyle \frac{m_B}{\lambda _B}}.`$ (9)
Since $`\varphi _B(\xi )`$ is appreciable only for $`\xi `$ of order $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_B`$, $`\lambda _B`$ is of order $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. We will choose $`\lambda _B300`$ MeV in the numerical calculations BBNS .
If we choose the asymptotic form of the K meson twist-2 LCDA , $`\varphi _K(y)=6y(1y)`$, we can find logarithmic end-point singularities in Eq. (8) just like that in Ref. Chao03 , and we parameterize it in a simple way,
$$\frac{dy}{y}=\mathrm{ln}\frac{m_B}{\mathrm{\Lambda }_h}2.4.$$
(10)
The mass of $`\chi _{c1}(1P)`$ $`M_{\chi _{c1}}=3.511`$ Gev is known, but the mass of the missing charmonium $`\chi _{c1}(2P)`$ has to be estimated by, say, potential models. We choose $`M_{\chi _{c1}^{}}=3.953`$ Gev following Ref. Godfrey . Then the form factor $`F_1(M^2)`$ can be determined by light-cone sum rules Ball ,
$`F_1(M_{\chi _{c1}}^2)=0.80,F_1(M_{\chi _{c1}^{}}^2)=1.14.`$ (11)
We also choose $`M_{\chi _{c1}^{}}=3.872`$ Gev and $`F_1(M_{\chi _{c1}^{}}^2)=1.06`$ to study if the $`X(3872)`$ behaves like a $`\chi _{c1}(2P)`$ in their b-production processes.
For numerical analysis, we use the following input parameters :
$`m_b=4.8\text{GeV},m_B=5.28\text{GeV},f_K=160\text{MeV},`$
$`f_B=216\text{MeV}\text{Gray},_{1P}^{^{}}(0)=_{2P}^{^{}}(0)=\sqrt{0.1}\text{GeV}^{5/2},`$
$`C_1(\mu )=1.21(1.082),C_2(\mu )=0.40(0.185),`$
$`C_3(\mu )=0.03(0.014),C_4(\mu )=0.05(0.035),`$
$`C_5(\mu )=0.01(0.009),C_6(\mu )=0.07(0.041),`$
$`\alpha _\mathrm{s}(\mu )=0.35(0.22).`$ (12)
In ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)) the $`\mu `$-dependent quantities at $`\mu _h=1.4`$ Gev ($`\mu =4.4`$ Gev) are shown without (with) parentheses.
Using the above inputs, we get the results of coefficients $`a_i`$ which are listed in Table. 1. With the help of these coefficients $`a_i`$, we calculate the decay branching ratios of decays $`B\chi _{c1}(1P,2P)K`$ with two different choices of $`M_{\chi _{c1}^{}}`$ and get
$`\mathrm{Br}(B^0\chi _{c1}(3511)K^0)`$ $`=`$ $`1.79\times 10^4,`$
$`\mathrm{Br}(B^0\chi _{c1}^{}(3953)K^0)`$ $`=`$ $`1.81\times 10^4,`$
$`\mathrm{Br}(B^0\chi _{c1}^{}(3872)K^0)`$ $`=`$ $`1.78\times 10^4.`$ (13)
Our prediction of $`\mathrm{Br}(B^0\chi _{c1}(3511)K^0)`$ is about 2 times larger than that in Chao03 , although it is still about two times smaller than the recent data BaBar05 . The difference between the theoretical predictions and experimental data may not be as serious as it looks like if we take into account the following uncertainties: (i) We have used a moderate value of $`_{1P}^{^{}}(0)`$ predicted by different potential models Quig in our calculation, and a larger value of $`_{1P}^{^{}}(0)`$ may enhance our prediction in Eq. ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)) significantly. (ii) In evaluation of $`f_{II}`$, we only use the leading twist LCDAs of K-meson, and large uncertainties will arise from the chirally enhanced higher twist effects Cheng . (iii) Since the squared velocity $`v^2`$ of the charm quark in charmonium is about 0.25-0.30, the relativistic corrections may be important for these decays.
Note that although the form factor in (11) and the coefficient $`a_2`$ in Table. 1 increase evidently as the charmonium mass increases, the decreased phase space and kinematic factors in (5) will make a balance, and result in similar decay branching ratios in the charmonium mass region 3.51-3.95 GeV, as shown in ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)). If we neglect the order $`\alpha _\mathrm{s}`$ corrections (i.e., in the naive factorization BSW ), the ratios among these three branching fractions in ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)) would become 1 : 0.74 : 0.69. As estimated in ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)), the branching ratios for $`\chi _{c1}(2P)`$ are
$`\mathrm{Br}(B^0\chi _{c1}^{}K^0)`$ $``$ $`2\times 10^4,`$
$`\mathrm{Br}(B^+\chi _{c1}^{}K^+)`$ $`=`$ $`\mathrm{Br}(B^0\chi _{c1}^{}K^0).`$ (14)
Comparing Eq. ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)) with the measured channel of the $`X(3872)`$ Belle03 :
$`\mathrm{Br}(B^+`$ $``$ $`XK^+)\times _X=(1.3\pm 0.3)\times 10^5,`$ (15)
$`_X`$ $``$ $`\mathrm{Br}(XJ/\psi \pi ^+\pi ^{}),`$
we see that the produced X(3872) looks like the $`\chi _{c1}(2P)`$ if $`_X`$ is sufficient small, say, $`37\%`$. A similar conclusion has recently been obtained in a comprehensive analysis of X(3872) production at the Tevatron and B-factories Bauer . On the other hand, if X(3872) is a loosely bound S-wave molecule of $`D^0\overline{D}^0/D^0\overline{D}^0`$ Tornqvist ; Swanson , a model calculation gives a smaller rate Braaten compared with Eq. ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)):
$`\mathrm{Br}(B^+`$ $``$ $`XK^+)=(0.071)\times 10^4,`$ (16)
which requires a larger $`_X>10\%`$ to be consistent with the experimental data (15). They also predict:
$`Br(B^0X(3872)K^0)<0.1Br(B^+X(3872)K^+).`$ (17)
So the measurement of $`_X`$ and $`Br(B^0X(3872)K^0)`$ is very helpful to identify the nature of X(3872).
Recently, a preliminary result for a new decay mode $`XD^0\overline{D}^0\pi ^0`$ was found by BelleOlsen :
$`\mathrm{Br}(B`$ $``$ $`XK)\times \mathrm{Br}(XD^0\overline{D}^0\pi ^0)`$ (18)
$`=`$ $`(2.2\pm 0.7\pm 0.4)\times 10^4.`$
Eq. (18) implies that $`_X<10\%`$, if it can be confirmed by further measurements. This would disfavor the suggestion that the X(3872) is a loosely bound S-wave molecule of $`D^0\overline{D}^0/D^0\overline{D}^0`$ with predictions of both decaySwanson and productionBraaten .
The above discussions about the X(3872) is based on the assumption that the X(3872) is a pure charmonium $`\chi _{c1}(2P)`$ state. But this cannot be the case due to the coupled channel effects and X(3872) being in extremely close proximity to the $`D^0\overline{D}^0/D^0\overline{D}^0`$ threshold. Perhaps a more realistic model for the X(3872) (for further discussions see chao ) is that the X(3872) has a dominant $`J^{PC}=1^{++}(2P)`$ $`c\overline{c}`$ component which is mixed with a substantial real $`D^0\overline{D}^0/D^0\overline{D}^0`$ continuum component (the $`D^+\overline{D}^{}/D^{}\overline{D}^+`$ continuum component is kinematically forbidden to be mixed in X(3872) and it is the $`ud`$ quark mass difference that causes this isospin violation). Thus X(3872) will have the following features. (1) The production of X(3872) in B meson decays is mainly due to the $`J^{PC}=1^{++}(2P)`$ $`c\overline{c}`$ component, as discussed above. The production of X(3872) at the Tevatron is also due to this $`c\overline{c}`$ component and associated higher Fock states containing the color-octet $`c\overline{c}`$ pair and soft gluons. As was argued chao1 for the prompt charmonium production that cross sections of D-wave charmonia (which were suggested as a tentative candidates for X(3872) in chao1 ) could be as large as $`J/\psi `$ or $`\psi (2S)`$ due to the color-octet mechanism, the P-wave $`(2P)`$ charmonium could also have the comparable production rate to $`J/\psi `$ or $`\psi (2S)`$. But this does not seem to be obvious for a loosely bound S-wave molecule of $`D^0\overline{D}^0/D^0\overline{D}^0`$. (2) On the other hand, the $`D^0\overline{D}^0/D^0\overline{D}^0`$ continuum component in X(3872) will be mainly in charge of the hadronic decays of X(3872) into $`D^0\overline{D}^0/D^0\overline{D}^0`$ or $`D^0\overline{D}^0\pi ^0`$ as well as $`J/\psi \rho ^0`$ and $`J/\psi \omega `$. The latter two decay modes ($`J/\psi \rho ^0`$ and $`J/\psi \omega `$) may come from the first decay mode $`D^0\overline{D}^0/D^0\overline{D}^0`$ and a subsequent rescattering final state interaction and therefore have the same decay amplitudes \[A($`J/\psi \rho ^0`$)=A($`J/\psi \omega `$)\] that are smaller than the first decay mode amplitude. (3) A substantial $`D^0\overline{D}^0/D^0\overline{D}^0`$ continuum component in X(3872) may reduce the production rates in Eq. ( $`B\chi _{c1}(1P,2P)K`$ decays in QCD factorization and X(3872)), and will also reduce the $`X(3872)J/\psi \gamma `$ decay width, which can be as small as $`11`$ KeV BG (note that this 2P-1S E1 transition is sensitive to the model details, see, e.g. Swanson ). This is much smaller then the hadronic decay widths. But a large rate for $`\chi _{c1}(2P)\gamma \psi (2S)`$=60-100 KeV will be expected. These qualitative features are useful in understanding the nature of X(3872) and should be further tested and studied experimentally and theoretically.
In summary, we study the decays $`B\chi _{c1}(1P,2P)K`$ in QCD factorization by treating charmonia as nonrelativistic bound states. We find that there are no infrared divergences in the vertex corrections, and the logarithmic end-point singularities from hard spectator interactions can be regularized by a momentum cutoff. Within certain uncertainties we find the $`B\chi _{c1}(2P)K`$ decay rate can be comparable to $`B\chi _{c1}(1P)K`$, and get $`Br(B^0\chi _{c1}^{}K^0)=Br(B^+\chi _{c1}^{}K^+)2\times 10^4`$. This might imply that the X(3872) has a dominant $`J^{PC}=1^{++}(2P)`$ $`c\overline{c}`$ component but mixed with some $`D^0\overline{D}^0+D^0\overline{D}^0`$ continuum component. The qualitative features of X(3872) are discussed and should be further tested and studied.
$`Note`$. After this work appeared in hep-ph/0506222 we learned some new results from BaBar BaBar05b : $`\mathrm{Br}(B^+X(3872)K^+)<3.2\times 10^4,R=\frac{Br(B^0X(3872)K^0)}{Br(B^+X(3872)K^+)}=0.50\pm 0.30\pm 0.05.`$ We also note that a recent papersuzuki (hep-ph/0508258) obtained similar conclusions to ours for the X(3872).
###### Acknowledgements.
We thank G. Bauer, D. Bernard, S. Olsen, and V. Paradimitriou for stimulating discussions on the experimental status of the X(3872), and E. Swanson for helpful comments. This work was supported in part by the National Natural Science Foundation of China (No 10421003), and the Key Grant Project of Chinese Ministry of Education (No 305001).
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# Probing nucleon strange asymmetry from charm production in neutrino deep inelastic scattering
## 1 Introduction
Nucleon structure is a natural laboratory to understand QCD and is worth to study for its own sake. The nucleon strange quark-antiquark asymmetry is an interesting feature predicted as a natural consequence of the non-perturbative aspect of the nucleon bm96 ; sig87 ; bur92 . Recently, the nucleon strange asymmetry has been suggested oln03 ; kre04 ; dm04 ; alw04 ; dxm04 ; wak04 as a promising mechanism to explain the NuTeV anomaly NT1 ; NT2 within the framework of Standard Model.
While the experimental evidence for such an asymmetry is still inconclusive, there are some approaches such as the global analysis of deep inelastic scattering (DIS) data oln03 ; bar00 , which show a favor for an asymmetric strange sea, in agreement qualitatively with the intrinsic sea theory. On the other hand, the CCFR next-to-leading-order (NLO) analysis of the neutrino induced dimuon production result favors a symmetric strange sea CCFR95 , which is also the case of a recent NuTeV analysis NTmas04 . It seems that more precise and dedicated research is needed to address the problem in a clear way.
The measurement of the strange quark distribution relies on charged current (CC) DIS processes. One method is through parity violating structure functions for isoscalar target in CC DIS: $`F_3^\nu F_3^{\overline{\nu }}=2[s(x)+\overline{s}(x)c(x)\overline{c}(x)]`$, which gives the total distribution of the strange sea. Another way is through the combination of CC parity conserving structure function $`F_2^\nu `$ with the charged lepton DIS structure function $`F_2^\mu `$, for isoscalar target, $`\frac{5}{6}F_2^\nu 3F_2^\mu =x[\frac{4}{3}s(x)\frac{1}{3}\overline{s}(x)c(x)]`$. Such an idea has been applied using high statistic neutrino and charged lepton nucleon DIS data, and the result at low $`x`$ shows a sizable disagreement with the direct measurement from the CCFR dimuon result con98 . The extraction of a small quantity from the difference of two large quantities may suffer from systematic uncertainties, which also seems to be the case for the extraction of the strange asymmetry by CC parity conserving structure functions: $`F_2^\nu F_2^{\overline{\nu }}=2x[s(x)\overline{s}(x)]`$.
A method free from the above drawback is to use the charged current charm production process, which is the main idea of the CCFR and NuTeV dimuon experiments CCFR95 ; CCFR93 ; NTdimu01 , with its leading order (LO) subprocesses being $`\nu _\mu s\mu ^{}c`$ and $`\nu _\mu d\mu ^{}c`$. The latter subprocess is Cabibbo suppressed, thus the charm production in $`\nu `$-induced process is most sensitive to the strange quark distribution in the nucleon. Similarly, the anticharm production in $`\overline{\nu }`$-induced process is sensitive to the antistrange quark distribution, as the corresponding partner subprocesses are $`\overline{\nu }_\mu \overline{s}\mu ^+\overline{c}`$ and $`\overline{\nu }_\mu \overline{d}\mu ^+\overline{c}`$, with the latter subprocess being Cabibbo suppressed.
The oppositely charged dimuon signature is easy to identify and measure in massive detectors, which allow for the collection of high statistics data samples, e.g., the CCFR experiment has a sample of data with 5030 $`\nu _\mu `$ induced events and 1060 $`\overline{\nu }_\mu `$ induced events, and the NuTeV has 5102 $`\nu _\mu `$ induced events and 1458 $`\overline{\nu }_\mu `$ induced events NTdimu01 . However, these two experiments neither show strong support for an asymmetric strange sea, nor can they rule it out bm96 ; kre04 ; bar00 . There are uncertainties in the estimation of the semi-muonic decay of the charmed hadrons CCFR95 ; CHOR02 , e.g., the average semi-leptonic branching ratio for $`\nu `$ and $`\overline{\nu }`$ induced events was only constrained by $`\frac{\overline{B}_c\overline{B}_{\overline{c}}}{\overline{B}_c}=\frac{0.011\pm 0.011}{0.1147}020\%`$ CCFR95 . Besides, the interplay of strange asymmetry and the light quark fragmentation (LQF) effect, as will be discussed in section 4, can only be drawn more clearly in inclusive measurement of charged and neutral charm productions. Thus a direct measurement of charmed hadrons produced in $`\nu `$ and $`\overline{\nu }`$ induced CC DIS will provide more valuable information to probe the s and $`\overline{s}`$ distributions of the nucleon. It is the purpose of this work to show that inclusive charm productions in neutrino and antineutrino induced CC DIS processes will be a promising way to detect the strange quark-antiquark asymmetry.
## 2 Charged current charm production
The differential cross section for charmed hadron $`H^+`$ production in neutrino induced CC DIS can be factorized as
$`{\displaystyle \frac{d^3\sigma _{\nu _\mu N\mu ^{}H^+X}}{d\xi dydz}}={\displaystyle \underset{q}{}}{\displaystyle \frac{d^2\sigma _{\nu _\mu N\mu ^{}qX}}{d\xi dy}}D_q^{H^+}(z),`$ (1)
where the function $`D_q^{H^+}(z)`$ describes the fragmentation of a quark q into the charmed hadron $`H^+`$, with $`z`$ being the momentum fraction of the quark $`q`$ carried by the produced hadron $`H^+`$. For the purpose of this article, the charmed hadron $`H^+`$ is taken to be $`D^+(c\overline{d})`$ or $`D^0(c\overline{u})`$ meson, with $`H^{}`$ denoting its antiparticle $`D^{}(\overline{c}d)`$ or $`\overline{D}^0(\overline{c}u)`$.
It is generally believed that the possibility for light quark fragmentation into charmed hadrons is very small. For example, the Lund string model implemented in some popular Monte Carlo programs predicts a suppression proportional to $`\mathrm{exp}(bm_q^2)`$ for $`q\overline{q}`$ production in the process of hadronization string81 . With a knowledge of the strange suppression $`\lambda 0.3`$ laf95 ; SLD97 , the suppression for charm will be lower than $`10^4`$, which can be safely neglected.
In this case, at leading order, only the $`\nu _\mu sc\mu ^{}`$ and $`\nu _\mu dc\mu ^{}`$ subprocesses contribute to charmed hadron production. For isoscalar target and neglecting target mass effects, the leading order differential cross section for charm production is given by CCFR95 ; CCFR93 :
$`{\displaystyle \frac{d^2\sigma _{\nu _\mu N\mu ^{}cX}}{d\xi dy}}`$ $`=`$ $`{\displaystyle \frac{2G^2ME_\nu }{\pi (1+Q^2/M_W^2)^2}}\left(1{\displaystyle \frac{m_c^2}{2ME_\nu \xi }}\right)`$ (2)
$`\times `$ $`\xi \left[{\displaystyle \frac{d(\xi )+u(\xi )}{2}}|V_{cd}|^2+s(\xi )|V_{cs}|^2\right],`$
where $`\xi `$ is the momentum fraction of the struck quark in the infinite momentum frame. It is introduced with the consideration of a non-negligible charm quark mass, and is related to Bjorken scaling variable $`x`$ through (neglecting light quark mass): $`\xi x(1+m_c^2/Q^2)`$, referred to as slow-rescaling. The term $`(1m_c^2/2ME_\nu \xi )`$ in Eq. (2) is introduced as an energy threshold for charm production and is supported by experiments NOM00 .
## 3 Probing the nucleon strange asymmetry
The differential cross sections for charmed hadrons, namely, $`H^+`$ ($`D^+`$ or $`D^0`$) and $`H^{}`$ ($`D^{}`$ or $`\overline{D}^0`$), produced in $`\nu `$ and $`\overline{\nu }`$ induced CC DIS respectively, are closely related to the $`s`$ and $`\overline{s}`$ distributions of the nucleon, and their difference, as can be seen in the following, is quite sensitive to the nucleon strange asymmetry.
Neglecting the light quark fragmentation effect, and using Eq. (2) and its corresponding partner process for $`\overline{c}`$ production $`\overline{\nu }_\mu N\overline{c}\mu ^+X`$, we can write the difference between $`H^+`$ and $`H^{}`$ production cross sections in CC DIS:
$`f_{H^+}f_H^{}`$ $``$ $`{\displaystyle \frac{d^3\sigma _{\nu _\mu N\mu ^{}H^+X}}{d\xi dydz}}{\displaystyle \frac{d^3\sigma _{\overline{\nu }_\mu N\mu ^+H^{}X}}{d\xi dydz}}`$ (3)
$`=`$ $`{\displaystyle \frac{2G^2ME_\nu }{\pi (1+Q^2/M_W^2)^2}}\left(1{\displaystyle \frac{m_c^2}{2ME_\nu \xi }}\right)`$
$`\times `$ $`\{{\displaystyle \frac{1}{2}}\xi [d_v(\xi )+u_v(\xi )]|V_{cd}|^2`$
$`+\xi [s(\xi )\overline{s}(\xi )]|V_{cs}|^2\}D_c^{H^+}(z),`$
where charge symmetry $`D_c^{H^+}(z)=D_{\overline{c}}^H^{}(z)`$ for fragmentation process is assumed, and $`u_v(\xi )u(\xi )\overline{u}(\xi )`$ and $`d_v(\xi )d(\xi )\overline{d}(\xi )`$ are valence quark distributions of the proton.
From Eq. (3), one sees that two terms, $`\frac{1}{2}\xi [d_v(\xi )+u_v(\xi )]`$ and $`\xi [s(\xi )\overline{s}(\xi )]`$, contribute to the cross section difference $`f_{H^+}f_H^{}`$, with $`|V_{cd}|^20.05`$ and $`|V_{cs}|^20.95`$ PDG04 being their respective weights. The strange asymmetric part of Eq. (3) can be estimated from an integral on variable $`\xi `$, to contribute a fraction
$`P_{\mathrm{SA}}{\displaystyle \frac{2S^{}|V_{cs}|^2}{Q_V|V_{cd}|^2+2S^{}|V_{cs}|^2}},`$ (4)
to the integral of the cross section difference $`𝑑\xi (f_{H^+}f_H^{})`$. Here, $`S^{}`$ and $`Q_V`$ are defined as $`S^{}\xi [s(\xi )\overline{s}(\xi )]𝑑\xi `$ and $`Q_V\xi [d_v(\xi )+u_v(\xi )]𝑑\xi `$.
In Table 1, results of the strange asymmetry from some models accounting for the NuTeV anomaly are listed, together with our estimations of the contributions due to strange asymmetry to the the cross section difference $`f_{H^+}f_H^{}`$, namely, the $`\xi `$ integrated fraction $`P_{\mathrm{SA}}`$.
As shown in Table 1, from the model calculations dm04 ; alw04 ; dxm04 ; wak04 that can explain the NuTeV anomaly, the strange asymmetry contributes a sizable proportion $`(12\%40\%)`$ to the cross section difference. Note that the distribution functions $`\xi [d_v(\xi )+u_v(\xi )]`$ and $`\xi [s(\xi )\overline{s}(\xi )]`$ may evolve with $`Q^2`$, turning flatter and shifting towards smaller $`\xi `$ region as $`Q^2`$ increases. However, their relative feature will remain and the proportion of $`S^{}`$ to $`Q_V`$, will be of the same order in larger $`Q^2`$ and in $`Q_0^2`$. Thus, as to their relative feature, it does not matter much whether the parton distributions are taken at $`Q_0^2`$ or at larger $`Q^2`$. Since the peak of $`\xi [s(\xi )\overline{s}(\xi )]`$ is confined in narrower $`\xi `$ region than $`\xi [d_v(\xi )+u_v(\xi )]`$, its contribution is expected to be more prominent than the integrated one in Table 1. Thus it is promising to measure the strange quark-antiquark asymmetry from $`f_{H^+}f_H^{}`$.
Compared to the sum of the cross sections $`f_{H^+}+f_H^{}`$, the cross section difference $`f_{H^+}f_H^{}`$ is not a very small quantity, as can be seen from the ratio of their integrals,
$`R{\displaystyle \frac{𝑑\xi (f_{H^+}f_H^{})}{𝑑\xi (f_{H^+}+f_H^{})}}{\displaystyle \frac{Q_V|V_{cd}|^2+2S^{}|V_{cs}|^2}{(Q_V+2Q_S)|V_{cd}|^2+2S^+|V_{cs}|^2}},`$ (5)
where $`Q_S\xi [\overline{u}(\xi )+\overline{d}(\xi )]𝑑\xi `$ and $`S^+\xi [s(\xi )+\overline{s}(\xi )]𝑑\xi `$. With a calculation of the $`Q_V`$, $`Q_S`$ and $`S^+`$ from CTEQ5 parametrization at $`Q^2=16`$GeV<sup>2</sup>, together with $`|V_{cd}|^2=0.05`$ and $`|V_{cs}|^2=0.95`$, the ratio $`R`$ is estimated to be about $`20\%`$ ($`25\%`$) for $`2S^{}/Q_V`$ being 0.007 (0.022) from Table 1. Thus the cross section difference $`f_{H^+}f_H^{}`$ is a significant quantity that can be extracted from the semi-inclusive differential cross sections.
Neutrino experiment with emulsion target, like the CHORUS detector, is ideal for the study of charmed hadron production. Compared to dimuon studies, it has a much lower level of background and is free from the uncertainties that exist in charm muonic weak decay processes CHOR02 . And for statistics, CHORUS reported in total about 94000 neutrino CC events located and fully reconstructed, in which about 2000 charm events were observed CHOR04 ; lel04 . This has been compatible with dimuon statistics. If such (or higher) statistics can be achieved with both neutrino and antineutrino beams of high energies in future experiments, the question about strange asymmetry is promising to be settled.
## 4 Light quark fragmentation
The possibility that a light quark fragments into charmed hadrons (associated charm production) can be an interesting effect of non-perturbative QCD, and it has been explored god8489 to explain the unexpected high rate of like-sign dimuons production from many neutrino experiments CDHSW ; CCFR93LSD ; CHARM . Although the field has been inactive for years, and in practice people generally assume the light quark fragmentation (LQF) to be negligible, the physical possibility of a small contribution is not ruled out. In fact, as to our consideration, neutrino experiments can be slightly different from $`e^+e^{}`$ experiments in this respect. The scattered light quark with high momentum can pick up a charm quark or antiquark from nucleon sea to form a $`D`$ meson, and the larger the energy of the light quark, the more the ability that it can pick up a charm from the sea. This energy dependence is apparent in prompt like-sign dimuons production rates in many experiments CDHSW ; CCFR93LSD ; CHARM . Since the scattered quark has most of the energy in the collision, it is much more promising to pick up the charm quark from nucleon sea than other produced quarks in fragmentation process. Thus as to our consideration such fragmentation as $`u\overline{D}^0(\overline{c}u)`$, $`dD^{}(\overline{c}d)`$ can possibly be non-negligible in high energy neutrino experiments.
In case that light quark can fragment into charmed hadrons, the process can manifest itself in a number of observables, such as the prompt like-sign dilepton and trimuon productions in high energy neutrino experiments smi83 , direct observations of two charmed hadrons in nuclear emulsion target CHOR02ONE , and charm production in hadron collisions E791D96 ; dia00 . Among these, the prompt like-sign dimuon productions have been the most seriously studied and we will use some of the data for a quantitative estimate of light quark fragmentation function and its influence on the extraction of strange asymmetry from CC charm production processes.
Prompt like-sigh dimuons ($`\mu ^{}\mu ^{}`$) can be produced through the process: $`\nu +d\mu ^{}+u`$, with the scattered $`u`$ to fragment into $`\overline{D}^0(\overline{c}u)`$ or $`\overline{D}^0(\overline{c}u)`$ meson. Its differential cross section can be expressed as:
$`{\displaystyle \frac{d^3\sigma _{\nu N\mu ^{}\mu ^{}X}}{dxdydz}}={\displaystyle \frac{2G^2ME_\nu |V_{ud}|^2}{\pi (1+Q^2/M_W^2)^2}}x{\displaystyle \frac{u(x)+d(x)}{2}}D_q(z)B_{\overline{D}^0},`$ (6)
where $`D_q(z)`$ is the total fragmentation function for a light quark to fragment into charmed hadrons, defined as $`D_q(z)D_q^D(z)+D_q^D^{}(z)`$, with $`D_q^D(z)D_u^{\overline{D}^0}(z)=D_{\overline{u}}^{D^0}(z)=D_d^D^{}(z)=D_{\overline{d}}^{D^+}(z)`$ and $`D_q^D^{}(z)D_u^{\overline{D}^0}(z)=D_{\overline{u}}^{D^0}(z)=D_d^D^{}(z)=D_{\overline{d}}^{D^+}(z)`$ simply assumed. Note here that $`D_q(z)`$ is energy dependent in analogous to containing an energy suppression factor for charm production. $`B_{\overline{D}^0}`$ is the inclusive muonic decay ratio for $`\overline{D}^0`$ meson decay $`\overline{D}^0\mu ^{}X`$, which is the same for $`\overline{D}^0`$ meson, since all $`\overline{D}^0`$ will decay into $`\overline{D}^0`$ at first.
Similarly, the differential cross section for prompt $`\mu ^+\mu ^+`$ production in $`\overline{\nu }`$ induced DIS on isoscalar target is
$`{\displaystyle \frac{d^3\sigma _{\overline{\nu }N\mu ^+\mu ^+X}}{dxdydz}}={\displaystyle \frac{2G^2ME_\nu |V_{ud}|^2}{\pi (1+Q^2/M_W^2)^2}}x{\displaystyle \frac{\overline{u}(x)+\overline{d}(x)}{2}}D_q(z)B_{D^0}.`$ (7)
Many experimental groups have reported positive results on prompt like-sign dimoun production. Among them the CDHSW CDHSW and CCFR CCFR93LSD data have a high precision and show much consistency with each other. Another high precision experiment CHARM CHARM , which has reported a much higher $`\mu ^{}\mu ^{}`$ production rate, received doubts on their estimate of $`\pi /K`$ decay background CDHSW . Besides, their kinematic cut $`p_\mu >4`$ GeV, which is lower than other experiments ($`p_\mu >9`$ GeV), can permit more $`\mu ^{}\mu ^{}`$ events and thus can produce a higher rate. Since kinematic cut on the second $`\mu `$ reduces event number and thus the rate of like-sign dimuons in CC events: $`\mu ^{}\mu ^{}`$/$`\mu ^{}`$, it will probably underestimate the light quark fragmentation (LQF) effect to use the $`\mu ^{}\mu ^{}`$/$`\mu ^{}`$ data. On the other hand, the ratio of prompt like-sign dimuons to opposite-sign dimuons $`\mu ^{}\mu ^{}`$/$`\mu ^{}\mu ^+`$ is expected to be less influenced by kinematic cut, as both second muons receive the same kinematic cut. Thus, we consider it appropriate to use the $`\mu ^{}\mu ^{}`$/$`\mu ^{}\mu ^+`$ data other than $`\mu ^{}\mu ^{}`$/$`\mu ^{}`$ data to estimate the LQF effect.
The differential cross section for $`\mu ^{}\mu ^+`$ production in $`\nu `$ induced CC DIS on an isoscalar target is given by
$`{\displaystyle \frac{d^3\sigma _{\nu N\mu ^{}\mu ^+X}}{d\xi dydz}}`$ $`=`$ $`{\displaystyle \frac{2G^2ME_\nu }{\pi (1+Q^2/M_W^2)^2}}f_c\overline{B}_c(z)`$ (8)
$`\times `$ $`\left[\xi {\displaystyle \frac{d(\xi )+u(\xi )}{2}}|V_{cd}|^2+\xi s(\xi )|V_{cs}|^2\right]`$
$`+`$ $`\delta \left({\displaystyle \frac{d^3\sigma _{\nu N\mu ^{}\mu ^+X}}{d\xi dydz}}\right)_{\mathrm{LQF}},`$
where $`f_c`$ is the charm suppression factor $`f_c1m_c^2/2ME_\nu \xi `$, and $`\overline{B}_c(z)`$ is the average muonic decay ratio for the charmed hadrons produced in CC DIS, $`\overline{B}_c(z)=_HD_c^H(z)B_H`$, with H being $`D^0`$, $`D^0`$, $`D^+`$, $`D^+\mathrm{}`$. The last term marked with LQF is the light quark fragmentation contribution to $`\mu ^{}\mu ^+`$ production ($`\nu \overline{u}\mu ^{}\overline{d}\mu ^{}D^+(D^+)X^{}\mu ^{}\mu ^+X`$),
$`\delta \left({\displaystyle \frac{d^3\sigma _{\nu N\mu ^{}\mu ^+X}}{d\xi dydz}}\right)_{\mathrm{LQF}}`$ $`=`$ $`{\displaystyle \frac{2G^2ME_\nu |V_{ud}|^2}{\pi (1+Q^2/M_W^2)^2}}D_q(z)\overline{B}_{D^{()+}}`$ (9)
$`\times `$ $`\left[x{\displaystyle \frac{\overline{d}(x)+\overline{u}(x)}{2}}\right](1y)^2,`$
where $`\overline{B}_{D^{()+}}`$ is the average muonic decay ratio of the $`D^+`$ and $`D^+`$ mesons produced from $`\overline{d}`$ quark fragmentation, $`\overline{B}_{D^{()+}}=\frac{1}{D_q}(D_q^DB_{D^+}+D_q^D^{}B_{D^+})`$. Since $`D^+`$ decays to $`D^0`$ with branching ratio $`B67.7\%`$ PDG04 , and to $`D^+`$ with ratio $`1B`$, we have $`\overline{B}_{D^{()+}}=bBB_{D^0}+(1bB)B_{D^+}`$, with $`bD_q^D^{}/D_q`$.
Similarly, the differential cross section for $`\mu ^+\mu ^{}`$ production in $`\overline{\nu }`$ induced CC DIS on isoscalar target reads
$`{\displaystyle \frac{d^3\sigma _{\overline{\nu }N\mu ^+\mu ^{}X}}{d\xi dydz}}`$ $`=`$ $`{\displaystyle \frac{2G^2ME_\nu }{\pi (1+Q^2/M_W^2)^2}}f_c\overline{B}_{\overline{c}}(z)`$ (10)
$`\times `$ $`\left[\xi {\displaystyle \frac{\overline{d}(\xi )+\overline{u}(\xi )}{2}}|V_{cd}|^2+\xi \overline{s}(\xi )|V_{cs}|^2\right]`$
$`+`$ $`\delta \left({\displaystyle \frac{d^3\sigma _{\overline{\nu }N\mu ^+\mu ^{}X}}{d\xi dydz}}\right)_{\mathrm{LQF}},`$
with
$`\delta \left({\displaystyle \frac{d^3\sigma _{\overline{\nu }N\mu ^+\mu ^{}X}}{d\xi dydz}}\right)_{\mathrm{LQF}}`$ $`=`$ $`{\displaystyle \frac{2G^2ME_\nu |V_{ud}|^2}{\pi (1+Q^2/M_W^2)^2}}D_q(z)\overline{B}_{D^{()}}`$ (11)
$`\times `$ $`\left[x{\displaystyle \frac{d(x)+u(x)}{2}}\right](1y)^2.`$
CDHSW has reported prompt dimuon rates $`\sigma _{\mu ^{}\mu ^{}}`$/$`\sigma _{\mu ^{}\mu ^+}`$ and $`\sigma _{\mu ^+\mu ^+}`$/$`\sigma _{\mu ^+\mu ^{}}`$ in $`\nu `$ and $`\overline{\nu }`$ induced DIS. As mentioned previously, these data are less influenced by kinematic cut and thus are better suited for the extraction of the LQF effect. The prompt dimoun rates from CDHSW with visible energy $`E_{\mathrm{vis}}`$ in range $`100200`$ GeV are listed in Table 2.
As can be seen from Table 2, the prompt dimuons rate $`\sigma _{\mu ^{}\mu ^{}}`$/$`\sigma _{\mu ^{}\mu ^+}`$ still show a slight dependence on kinematic cut, though much smaller than the $`\sigma _{\mu ^{}\mu ^{}}`$/$`\sigma _\mu ^{}`$ data do. Thus we can only estimate the order of magnitude for the LQF effect with the reported data.
The rate $`\sigma _{\mu ^{}\mu ^{}}`$/$`\sigma _{\mu ^{}\mu ^+}`$ (without any kinematic cut) can be deduced by Eq. (6) and Eq. (8) with an integral on kinematic variables to approximate
$`{\displaystyle \frac{\sigma _{\mu ^{}\mu ^{}}}{\sigma _{\mu ^{}\mu ^+}}}{\displaystyle \frac{Q_{ud}|V_{ud}|^2}{Q_{ud}|V_{cd}|^2+S|V_{cs}|^2}}{\displaystyle \frac{D_qB_{D^0}}{\overline{f}_c\overline{B}_c}},`$ (12)
where $`Q_{ud}\frac{1}{2}x[u(x)+d(x)]𝑑x`$, $`Sxs(x)𝑑x`$, and $`\overline{f}_c`$ denotes the average of energy suppression factor in Eq. (8). With the measured $`B_{D^0}6.87\%`$ and $`\overline{B}_c8.8\%`$ lel04 , together with $`Q_{ud}`$ and $`S`$ from CTEQ5 at $`Q^2=16`$ GeV<sup>2</sup>, $`D_q/\overline{f}_c`$ is estimated to be:
$`{\displaystyle \frac{D_q}{\overline{f}_c}}0.199{\displaystyle \frac{\sigma _{\mu ^{}\mu ^{}}}{\sigma _{\mu ^{}\mu ^+}}}.`$ (13)
With the experimental data on $`\sigma _{\mu ^{}\mu ^{}}/\sigma _{\mu ^{}\mu ^+}`$ from Table 2, one can easily estimate $`D_q/\overline{f}_c`$ by Eq. (13).
Since we are most interested in the strange quark-antiquark asymmetry here, we will directly address the influence of LQF effect on the extraction of strange asymmetry. Because LQF effect contributes differently for $`\nu `$ induced $`\mu ^{}\mu ^+`$ production and for $`\overline{\nu }`$ induced $`\mu ^+\mu ^{}`$ production, it will give different corrections to $`s`$ and $`\overline{s}`$ distributions, and thus influence the measurement of strange asymmetry from opposite-sign dimuon method. To illustrate this, we will compare the contribution of LQF effect with that of strange asymmetry on the difference between $`\nu `$ and $`\overline{\nu }`$ induced opposite-sign dimuon production cross sections. The latter (strange asymmetry contribution) can be drawn from model predictions in the last column of Table 1, when assuming the average muonic branding ratio of charmed hadrons to be the same for $`\nu `$ and $`\overline{\nu }`$ induced CC DIS $`\overline{B}_c(z)=\overline{B}_{\overline{c}}(z)`$. The former (LQF contribution) can be deduced from Eq. (8)-Eq. (11) with an assumption $`\overline{B}_{D^{()+}}=\overline{B}_{D^{()}}`$, and be compared to the strange asymmetry part with an integral on kinematic variables. The fraction of the LQF contribution is
$`P_{\mathrm{LQF}}`$ $``$ $`{\displaystyle \frac{\delta (\sigma _{\nu N\mu ^{}\mu ^+X}\sigma _{\overline{\nu }N\mu ^+\mu ^{}X})_{\mathrm{LQF}}}{(\sigma _{\nu N\mu ^{}\mu ^+X}\sigma _{\overline{\nu }N\mu ^+\mu ^{}X})_{\mathrm{total}}}}`$ (14)
$``$ $`{\displaystyle \frac{\frac{1}{3}Q_V|V_{ud}|^2}{Q_V|V_{cd}|^2+2S^{}|V_{cs}|^2}}{\displaystyle \frac{D_q\overline{B}_{D^{()+}}}{\overline{f}_c\overline{B}_c}}.`$
To assess $`P_{\mathrm{LQF}}`$, the value of $`\overline{B}_{D^{()+}}`$ is needed. Remember that $`\overline{B}_{D^{()+}}=bBB_{D^0}+(1bB)B_{D^+}`$, with $`bD_q^D^{}/D_q`$. The unknown $`b`$ is the fraction of vector $`D^{}`$ meson in light quark fragmentation. When we set $`b`$ to be $`1/32/3`$, and take $`B_{D^0}6.87\%`$, $`B_{D^+}17.2\%`$ lel04 , we get $`\overline{B}_{D^{()+}}=(13.7\pm 1.2)\%`$. Using Eq. (14) and taking $`2S^{}/Q_V=0.007`$ from Table 1, we get $`P_{\mathrm{LQF}}=(1.73\pm 0.15)\sigma _{\mu ^{}\mu ^{}}/\sigma _{\mu ^{}\mu ^+}`$. Taking $`\sigma _{\mu ^{}\mu ^{}}/\sigma _{\mu ^{}\mu ^+}=(3.5\pm 1.6)\%`$ from Table 2, we get
$`P_{\mathrm{LQF}}=(6.1_{3.1}^{+3.5})\%.`$ (15)
Thus, we get an estimate of the LQF contribution to be a few percent compared to strange asymmetry contribution $`P_{\mathrm{SA}}:12\%40\%`$. However, the constraint of $`P_{\mathrm{LQF}}`$ can also be done with $`\sigma _{\mu ^+\mu ^+}/\sigma _{\mu ^+\mu ^{}}`$ data, and the result is $`P_{\mathrm{LQF}}^{}=(33_{16}^{+19})\%`$, which is very large compared to result from the $`\sigma _{\mu ^{}\mu ^{}}/\sigma _{\mu ^{}\mu ^+}`$ data. This large discrepancy is difficult to explain at present, and may imply an uncertainty in the estimate of the LQF contribution in the opposite-sign dimuon measurements of strange asymmetry.
From the sign and size of $`P_{\mathrm{LQF}}`$, one sees that the LQF effect contributes oppositely to the predicted strange asymmetry contribution on the whole, with a rate that could be non-negligible in opposite-sign dimuon experiments.
The LQF effect also exists in the process of inclusive charm productions that we suggest. For $`D^\pm `$ production, the cross section difference, $`f_{D^+}f_D^{}`$, for $`\nu `$ and $`\overline{\nu }`$ induced CC DIS will include an additional term from light quark fragmentation:
$`\delta (f_{D^+}f_D^{})_{\mathrm{LQF}}`$ $`=`$ $`{\displaystyle \frac{2G^2ME_\nu |V_{ud}|^2}{\pi (1+Q^2/M_W^2)^2}}D_q(z)(1\epsilon )`$ (16)
$`\times `$ $`\left[x{\displaystyle \frac{d_v(x)+u_v(x)}{2}}\right](1y)^2,`$
where $`\epsilon =Bb`$ is introduced with the consideration that part of $`D^+(D^{})`$ will decay into $`D^0(\overline{D}^0)`$ and will not contribute to the cross sections.
For neutral charm production, LQF contributes to $`\overline{D}^0`$ production in $`\nu `$ induced CC DIS ($`\nu +d\mu ^{}+u`$, $`u\overline{D}^0(\overline{c}u)`$), and to $`D^0`$ production in $`\overline{\nu }`$ induced CC DIS. In case that $`D^0`$ and $`\overline{D}^0`$ are not distinguished by emulsion target, the $`\overline{D}^0`$ $`(D^0)`$ production in $`\nu `$ ($`\overline{\nu }`$) induced CC DIS from LQF will be incorporated to $`D^0`$ $`(\overline{D}^0)`$ production in $`\nu `$ ($`\overline{\nu }`$) induced CC DIS. Thus an additional term from LQF will contribute to $`f_{D^0}f_{\overline{D}^0}`$:
$`\delta (f_{D^0}f_{\overline{D}^0})_{\mathrm{LQF}}`$ $`=`$ $`{\displaystyle \frac{2G^2ME_\nu |V_{ud}|^2}{\pi (1+Q^2/M_W^2)^2}}D_q(z)(1\epsilon ^{})`$ (17)
$`\times `$ $`\left[x{\displaystyle \frac{d_v(x)+u_v(x)}{2}}\right],`$
where $`\epsilon ^{}=(1y)^2Bb`$, which is introduced from $`\overline{d}`$ ($`d`$) fragmentation into $`D^+`$ $`(D^{})`$ mesons that then decay into $`D^0`$ $`(\overline{D}^0)`$ and contribute to cross section difference $`f_{D^0}f_{\overline{D}^0}`$.
The proportion of LQF contribution to inclusive charm production cross section difference $`f_{H^+}f_H^{}`$, namely $`P_{\mathrm{LQF}}^{H^\pm }`$, can be estimated similarly to that of dimuon productions. With an integral on kinematic variables of Eq. (3), Eq. (16), Eq. (17), and using charm production fractions $`D_c^{D^+}(z)𝑑z0.26`$ and $`D_c^{D^0}(z)𝑑z0.66`$ for $`E_\nu >80`$GeV lel04 , $`P_{\mathrm{LQF}}^{H^\pm }`$ is estimated (in unite of $`P_{\mathrm{LQF}}`$) to be: $`P_{\mathrm{LQF}}^{D^\pm }1.6P_{\mathrm{LQF}}`$ for $`D^\pm `$ meson productions, and $`P_{\mathrm{LQF}}^{D^0}2.6P_{\mathrm{LQF}}`$ for $`D^0`$, $`\overline{D}^0`$ meson productions.
If the LQF contribution $`P_{\mathrm{LQF}}`$ in opposite-sign dimuons measurement is in the order of a few percent percent and opposite to strange asymmetry contribution $`P_{\mathrm{SA}}:12\%40\%`$, just as we have estimated, the LQF will contribute to inclusive charm production with a larger proportion (in the order of about ten percent or even larger). For inclusive $`D^\pm `$ production, LQF contributes oppositely compared to strange asymmetry when $`x[s(x)\overline{s}(x)]>0`$. On the other hand, for inclusive CC neutral charm $`(D^0,\overline{D}^0)`$ production, LQF contributes positively compared to strange asymmetry when $`x[s(x)\overline{s}(x)]>0`$. A separation of the LQF effect and the strange asymmetry effect can be made from the distinct features of $`f_{D^+}f_D^{}`$ and $`f_{D^0}f_{\overline{D}^0}`$ measured by nuclear emulsion target. Thus, the inclusive measurement of charged and neutral charm production in $`\nu `$ and $`\overline{\nu }`$ induced CC DIS will shed light on both the strange asymmetry and the LQF effect.
Dedicated analysis of charm productions in neutrino experiments and in other processes will be helpful for a more precise estimate and constraint for the light quark fragmentation effect.
## 5 Conclusions
For probing the nucleon strange asymmetry, we analyzed the charged current charm production processes, in particular, the $`\nu _\mu `$ induced $`H^+`$ ($`D^+`$ or $`D^0`$) production and the $`\overline{\nu }_\mu `$ induced $`H^{}`$ ($`D^{}`$ or $`\overline{D}^0`$) production processes. The strange asymmetry from various model calculations that can explain the NuTeV anomaly is shown in general to contribute a sizable proportion ($`12\%40\%`$) to the $`H^\pm `$ differential cross section difference $`f_{H^+}f_H^{}`$. Thus, measurement of these cross sections with high energy neutrino and antineutrino beams on nuclear emulsion target is very promising to detect the strange quark-antiquark asymmetry.
Meanwhile, we analyzed the possible light quark fragmentation (LQF) effect from prompt like-sign dimuon data and studied its influence on the measurement of strange asymmetry. Our result is that the LQF may be an important source that reduces the effect of strange asymmetry from opposite sign dimuon studies. And for inclusive charged current (CC) charm production with emulsion target, since the contributions of LQF are in opposite directions for $`D^\pm `$ and for $`D^0`$ ($`\overline{D}^0`$) productions, a separation of the LQF effect from strange asymmetry effect can be made by the separate measurement of $`D^\pm `$ and neutral charm differential cross sections in CC DIS. Thus the inclusive measurement of charmed hadrons can shed light on both strange asymmetry and the LQF effect. Further analysis and constraint for LQF effect from various experiments will also be helpful for the purpose of measuring the strange asymmetry more reliably.
Acknowledgments: We acknowledge the helpful discussion with Vincenzo Barone. This work is partially supported by National Natural Science Foundation of China (Nos. 10025523, 90103007, and 10421003), by the Key Grant Project of Chinese Ministry of Education (No. 305001), and by the Italian Ministry of Education, University and Research (MIUR).
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# Contents
## 1 A new diagram for spacetime structure
Spacetimes are geometrical objects, independent of the coordinates with which we describe them. However, spacetimes are typically presented and visualized in a specific coordinate system. If the coordinates are poorly chosen, many properties of the spacetime such as horizons, causally connected spacetime points, maximal extensions and null infinity are not readily apparent.
A simplification occurs if a $`D`$ dimensional Lorentzian spacetime has enough fibered directions (like a $`(D2)`$-sphere) or other ignorable directions. One can draw two dimensional diagrams for the remaining directions and such Lorentzian $`+`$ signature spacetime slices can be conformally compactified leading to Penrose diagrams.
Penrose diagrams are quite useful in understanding spacetime geometry and successful especially in understanding causal structure although there are some limitations to this approach. For instance just knowing the Penrose diagram for the subextremal $`Q^2<M^2`$ Reissner-Nordstrøm black hole does not tell us what happens to the spacetime structure in the chargeless or extremal limits. For more complicated spacetimes, Penrose diagrams (which assume symmetry or fibering) can only draw a slice of the spacetime. As a known example, the Penrose diagram for a Kerr black hole does not clearly depict the ring singularity and the possibility of crossing through the interior of the ring into a second universe. In addition, recently analytic continuation has been applied to black hole solutions to yield bubble-type or S-brane solutions. Oftentimes this is done in Boyer-Lindquist type coordinates which are hard to visualize. Again we are not left with a clear picture of the resulting spacetime and the Penrose diagrams are missing important noncompact spatial directions.
It is useful to have an alternative diagram which can also capture important features of a spacetime. For this reason in this paper we expand the notion of drawing spacetimes in Weyl space . Because our diagrams have the appearance of playing cards glued together we will dub them Weyl card diagrams.
To understand the construction of a card diagram we recall that in $`D=4`$ dimensions a Weyl solution in canonical coordinates is written as
$$ds^2=fdt^2+f^1[e^{2\gamma }(d\rho ^2+dz^2)+\rho ^2d\varphi ^2]$$
where $`f`$ and $`\gamma `$ are functions of $`\rho ,z`$. The original Weyl class requires two commuting orthogonal Killing fields $`_t,_\varphi `$ in four dimensions , or $`D2`$ fields for general $`D`$ dimensions .<sup>1</sup><sup>1</sup>1 Non-Weyl, axisymmetric spacetimes in $`D4`$ are discussed in . Sometimes Weyl solutions are called axially-symmetric gravitational solutions although they in fact are more general. We also include the Weyl-Papapetrou class for 2 commuting Killing vectors in $`D=4`$ , and allow charged static solutions in $`D4`$ (see the Appendix to this paper). Furthermore stationary vacuum solutions in $`D4`$ are covered with the recent work of . In four and five dimensions this generalized Weyl class includes spinning charged black holes and rings as well as various arrays of black holes, spacelike-branes, and includes backgrounds like Melvin fluxbranes and spinning ergotubes .
When constructing card diagrams, we will draw only Weyl’s canonical coordinates $`(\rho ,z)`$, or coordinates related to them via a conformal transformation. The Killing coordinates are ignorable and so this diagram is efficient and will show all details of the spacetime. Since there are only two nontrivial coordinates, card diagrams are two dimensional and easy to draw like Penrose diagrams. The difference however is that while Penrose diagrams are truly two-dimensional, card diagrams are drawn as if embedded in three dimensions. When a $`(\rho ,z)`$ region of the spacetime has Euclidean $`++`$ signature, we draw the two coordinates $`(\rho ,z)`$ horizontally; and this makes a horizontal card. For Lorentzian signature $`+`$ regions we use $`(\rho ^{},z)`$ or $`(\tau ,\rho )`$, and draw the timelike coordinate vertically; this makes a vertical card. Causal structure is automatically built into the vertical cards since for example the directions $`(\tau ,\rho )`$ appear conformally only through the combination $`d\tau ^2+d\rho ^2`$. Horizontal cards and vertical cards are attached together at Killing horizons and so card diagrams resemble a gluing-together of a house of playing cards.
In this paper we will present card diagrams for the familiar spacetimes of black holes, as well as expanding bubbles, S-branes, and black rings. Many other spacetimes including the S-dihole, infinite periodic universe, C-metric, and multiple-rod solutions in 4 and 5 dimensions are presented in , and spacetimes derived from 4 and 5 dimensional Kerr geometries will be presented in .
In Section 2 we review the Schwarzschild black hole in the usual coordinates and in Weyl coordinates. By extending through the horizon and properly representing the interior of the black hole we construct the first card diagram. We emphasize the construction of the interior of the black hole as a vertical card comprised of four triangles unfolded across special null lines. We then discuss general card diagram properties such as null lines, list the available card types, and geometrically describe the $`\gamma `$-flip analytic continuation procedure.
In Section 3, we discuss the sub/super/extremal Reissner-Nordstrøm black hole card diagrams, the Kerr black hole, and the black ring/C-metric card diagrams. We then show that a spacetime can have multiple card diagrams and as examples present the elliptic, hyperbolic, and parabolic representations of the charged Witten bubble and charged Spacelike brane which we also call S-Reissner-Nordstrøm. Finally, as a newer example we discuss the (twisted) S-Kerr solution .
We conclude with a discussion in Section 4. We give an appendix on perturbing Witten bubbles and S-branes by introducing Israel-Khan rods, in their hyperbolic or elliptic representations. We also give an appendix on how the higher dimensional vacuum Weyl Ansatz can be extended to include electromagnetic fields.
This paper is a condensed presentation of the card diagrams in .
## 2 Schwarzschild and general card diagrams
In this section we review the Schwarzschild black hole as an example of a Weyl spacetime and then we explain the construction of its associated Weyl card diagram. General features and properties of card diagrams are also developed.
Up to now if a solution such as the Schwarzschild black hole had horizons, then only the regions outside the horizons have been drawn in Weyl coordinates . To go through a nonextremal horizon, the Weyl coordinate $`\rho `$ must be allowed to take imaginary values. We discuss how the horizon can be represented as a junction of four regions which we call four cards. The regions outside the horizon will be drawn as two horizontal cards while the regions between the horizon and the singularity will be drawn as two vertical cards. The interior vertical cards naively are problematic and have fourfold-covered triangles bounded by ‘special null lines’. However the triangles can be unfolded and glued together into a square along the ‘special null lines’ to achieve a singly covered representation of the spacetime in Weyl coordinates by properly choosing branches of a square root in the solution.
Note that card diagrams represent the spaces on which we solve the Laplace equation (horizontal card) or wave equation (vertical card) to find a Weyl metric. For example, the Schwarzschild black hole has a finite length uniform density rod source along the $`z`$-axis generating the potential $`\mathrm{log}f`$, and the remainder of its $`z`$-axis encodes the vanishing of the $`\varphi `$-circle. Thus card diagrams give a full account of the boundary conditions necessary to specify the spacetime.
### 2.1 Schwarzschild black holes
This section will describe the construction of the Schwarzschild black hole card diagram. The Penrose and card diagrams are compared in Fig. 1.
The Schwarzschild metric in spherically symmetric (Schwarzschild) coordinates is
$$ds^2=(12M/r)dt^2+(12M/r)^1dr^2+r^2d\theta ^2+r^2\mathrm{sin}^2\theta d\varphi ^2.$$
(1)
There is a horizon at $`r=2M`$ and a curvature singularity at $`r=0`$. On the other hand, in Weyl’s canonical coordinates the metric ansatz is it is
$$ds^2=fdt^2+f^1(e^{2\gamma }(d\rho ^2+dz^2)+\rho ^2d\varphi ^2)$$
(2)
where $`f`$ and $`\gamma `$ are functions of the coordinates $`\rho `$ and $`z`$:
$`f`$ $`=`$ $`{\displaystyle \frac{(R_++R_{})^24M^2}{(R_++R_{}+2M)^2}},`$ (3)
$`e^{2\gamma }`$ $`=`$ $`{\displaystyle \frac{(R_++R_{})^24M^2}{4R_+R_{}}},`$
$`R_\pm `$ $`=`$ $`\sqrt{\rho ^2+(z\pm M)^2}.`$
Previously attention was restricted to the H1 half-plane $`\rho 0`$, $`\mathrm{}<z<\mathrm{}`$, known as Weyl space, which describes the exterior of the black hole and whose horizon is represented by a ‘rod’ line segment $`\rho =0`$, $`MzM`$; see Fig. 2. Note that the non-Killing 2-metric is conformal to the Euclidean flat space $`d\rho ^2+dz^2`$. The coordinate transformation between Schwarzschild and Weyl coordinates is
$`\rho `$ $`=`$ $`\sqrt{r^22Mr}\mathrm{sin}\theta ,`$ (4)
$`z`$ $`=`$ $`(rM)\mathrm{cos}\theta .`$
Now we wish to ask how Weyl’s coordinates draw the spacetime inside the horizon. The Schwarzschild coordinates (4) tell us that for $`0<r<2M`$, $`\rho `$ is imaginary and so we set $`\rho ^{}=i\rho `$. (In general we must perform an analytic continuation of Weyl coordinates to go through a horizon which are at the zeros of the Weyl functions $`f,e^{2\gamma }`$.) The analytic continuation gives a region with a conformally Minkowskian metric $`d\rho ^2+dz^2`$ and we will draw this region as being vertical and perpendicularly attached to the horizontal card at the horizon $`MzM`$. The vertical direction is always timelike in card diagrams.
Of course the Schwarzschild horizon structure is more complicated than just having two cards joined together. For example we know the Penrose diagram in the $`t,r`$ coordinates is divided into four regions that meet in a $`\times `$-horizon structure. In Weyl coordinates, we already saw that we can extend $`\rho \pm i\rho ^{}`$ or go to negative $`\rho `$. This gives us the four regions of the Schwarzschild black hole in Weyl coordinates. Two regions will be horizontal at real values of $`\rho `$ and two regions will be vertical with imaginary values of $`\rho `$. So in addition to the first horizontal card in front of the horizon, we also have a copy of the horizontal external universe behind the horizon and attach two vertical cards, one above and one mirror copy below the horizontal cards (see Fig. 1). All together, four different regions attach together at the same $`MzM`$ rod horizon.
The four regions labelled H1, H2, V1 and V2 in the Penrose diagram map to the similarly labelled four regions on the card diagram in Fig. 1. Note that the Weyl cards draw the $`r,\theta `$ coordinates which is different from the $`t,r`$ coordinates of the Penrose diagram. However the fact that the radial coordinate $`r`$ describes four distinct regions, two where $`_r`$ is spacelike and two where it is timelike, is still apparent in the Weyl card diagram. So while a Penrose diagram always has a Lorentzian $`+`$ signature, a card diagram will flip from being Euclidean $`++`$ to Lorentzian $`+`$ across a nonextremal Killing horizon.
Let us now examine the construction of the upper vertical card extended in the $`\rho ^{},z`$ directions. Looking at an $`r`$-orbit for $`r[0,2M]`$ on the vertical card, we see that in Weyl coordinates $`0\rho ^{}M\pm z`$. The bounding lines where $`R_\pm =\sqrt{\rho ^2+(z\pm M)^2}=0`$, we will call ‘special null lines,’ and they are a general feature of vertical cards with focal points (the rod endpoints $`z=\pm M`$). Here, special null lines are the envelope of the $`r`$-orbits as we vary $`\theta `$. Inside the horizon the Schwarzschild coordinates apparently fill out a 45-45-90 degree Weyl triangle with hypotenuse length $`2M`$ a total of four times, as shown in Figure 3.
Special null lines play an important role in Weyl card diagrams so let us explain their significance. Keep in mind that we have already broken the manifest spherical symmetry when we have written the Schwarzschild metric in Weyl coordinates, so the existence of preferred special null lines is relative to this chosen axis. Consider the two 3-surfaces
$$R_\pm =rM\pm M\mathrm{cos}\theta =0,$$
(5)
which are drawn in Fig. 4. These surfaces bound the trajectories of light rays that do not move in the Killing directions. The surfaces intersect at $`r=M`$ and partition the black hole interior into four subregions. These regions correspond to the four Weyl triangles.
It is clear from (5) that $`R_\pm `$ is positive outside the horizon and there is no difficulty going to negative values inside the horizon. On the other hand in terms of Weyl coordinates, the functions $`R_\pm =\sqrt{\rho ^2+(z\pm M)^2}`$ are the square root of a positive number when $`\rho ^{}<z\pm M`$ and imaginary if $`\rho ^{}>z\pm M`$. Clearly, instead of dealing with imaginary values of $`R_\pm `$, the way to go ‘beyond’ the special null line $`R_\pm =0`$ is to keep $`\rho ^{}<z\pm M`$ but use the other square root branch for $`R_\pm `$ as it enters in the Weyl functions $`f`$ and $`\gamma `$ by explicitly replacing $`R_\pm R_\pm `$. Since we can pass $`R_+=0`$ and/or pass $`R_{}=0`$, it is clear that the four different branches of the square root functions differentiate the four copies of the Weyl triangle.
From (3), passing each $`R_\pm =0`$ null line changes the sign of $`e^{2\gamma }`$ and hence exchanges the timelike and spacelike nature of $`\rho ^{}`$ and $`z`$. Since the vertical direction on a vertical card always represents time, two of the triangles are drawn turned on their sides, their hypotenuses vertical and so timelike. The four triangles describing the interior of the black hole $`0r2M`$ glue together along the special null lines to fit neatly into a square; see Fig. 5. Because of the unfolding of the triangles, the positive $`z`$-direction on the top triangle (and any attached horizontal cards) points in the opposite direction compared to that on the original horizontal card.
The upper vertical cards is thus a square of length $`2M`$. The bottom of the upper card V1 is the black hole horizon which connects to three other cards in a four card junction. The right and left edges of this vertical card correspond to $`\theta =0,\pi `$ and are the boundaries where $`\rho ^{}=0`$ the $`\varphi `$-circle vanishes. The top edge of the card represents the $`r=0`$ curvature singularity. The second vertical card V2 is built in analogous fashion except the square is built in a downwards fashion towards negative values of $`\rho ^{}`$. Additionally there is a second horizontal card plane, H2, identical to H1 at negative values of the $`\rho `$ coordinate attached to the same horizon along $`[M,M]`$ on the $`z`$-axis.
One typically stops the construction of the Schwarzschild spacetime with the above four regions, and considers the $`r<0`$ part of the metric to be a separate spacetime. However for reasons which become clear when we look at the Reissner-Nordstrøm and Kerr black holes in Sections 3.1.1 and 3.1.4, we continue the card diagram past the singularity and attach two horizontal half-plane corresponding to negative-mass (or $`r<0`$) universes h3 & h4, and further vertical card above the singularity which is identical to V2. The two new horizontal cards each represent negative mass-universes with no horizon and a naked singularity along $`MzM`$. Note that although Penrose diagrams for the $`r>0`$ and $`r<0`$ regions of Schwarzschild cannot attach together since the singularity is spacelike in one region and timelike in another, cards can naturally attach at this singularity. The extended card diagram for Schwarzschild, shown in Fig. 6, is an infinite array of repeated cards representing positive (H1,H2) and negative (h3,h4) mass universes and inside-horizon regions.
### 2.2 General properties of card diagrams
Having described the construction of the card diagram for Schwarzschild we now turn to general remarks about these card diagrams. Card diagrams can be constructed in general D-dimensions starting from a generalized Weyl ansatz including the charged cases discussed in Appendix B. Horizontal cards are conformally Euclidean and represent stationary regions. Vertical cards are always conformally Minkowskian and represent regions with ($`D2`$) spacelike Killing fields. On a vertical card time is always in the vertical direction and a spacetime point’s causal future lies between $`45^{}`$ null lines on the card. Weyl’s coordinates certainly go bad at horizons, so these diagrams are not a full replacement for Penrose diagrams at understanding causal structure or particle trajectories. However, it is clear for example from Fig. 1 that two vertical and two horizontal cards attach together in a orthogonal $`+`$-configuration in precisely the sense of the $`\times `$-horizon structure of the Penrose diagram.
The prototypical horizon is that of the two Rindler and two Milne wedges of flat space; this spacetime has two horizontal half-plane cards and two vertical half plane cards that meet along the horizon being the whole $`z`$-axis. Zooming in on a non-extremal horizon of any card diagram yields this Rindler/Milne picture.
A rod endpoint $`(\rho ,z)=(0,z_i)`$ such as $`z=\pm M`$ for Schwarzschild is a ‘focus’ for the Weyl diagram and often represents the edge of a black hole horizon or the end of an acceleration horizon on a horizontal card. Generally, multi-black hole Weyl spacetimes depend only on distances to foci, such as $`R_\pm `$ in the case of Schwarzschild .
To understand that in general it is natural to change the branch of the distance functions, $`R_i`$, when crossing the special null line that emanates from the foci at $`z=z_i`$, take some Weyl spacetime and imagine moving upwards in a vertical card to meet the special null line where $`R_i=0`$ by increasing time $`\rho ^{}0`$ for fixed spatial $`z`$. Rearranging $`R_i=\sqrt{\rho ^2+(z+z_i)^2}`$ as the semi-ellipse
$$R_i^2+\rho ^2=(z+z_i)^2,$$
we see that a smooth traversal of this semi-ellipse across $`R_i=0`$ requires a change in the sign of $`R_i`$.
In many of our solutions, special null lines are used to reflect vertical triangular cards to create full, rectangular cards. However, in the Bonnor-transformed S-dihole geometry of as well as double Killing rotated extremal geometries and parabolic representations of the bubble and S-Schwarzschild in Sec. 3.2.3, the special null lines will serve as conformal boundaries at null infinity $`^\pm `$.
Boundaries of cards indicate where the metric coefficient along a Killing (circle) direction vanishes. Which circle vanishes is constant over a connected piece of the boundary, even when the boundary turns a right angle onto a vertical card. Furthermore the periodicity to eliminate conical singularities is constant along connected parts of the boundary. For Schwarzschild, the $`\varphi `$-circle vanishes on both connected boundaries and has periodicity $`2\pi `$.
Although not a full replacement for understanding causal structure, it is interesting to consider geodesic trajectories on card diagrams. For example when a light ray is incident from a horizontal card onto a horizon (to enter the upper vertical card), it must turn and meet that horizon perpendicularly. It then appears on the vertical card, again perpendicular to the horizon. Only those light rays which go from the lower vertical card to the upper vertical card directly can meet the horizon rod at a non-right angle; these rays would touch the vertex of the $`\times `$ in a Schwarzschild Penrose diagram. When a light ray on a vertical card hits a boundary where a spacelike circle vanishes, it bounces back at the same angle as drawn on the card relative to the perpendicular.
Spacetimes with a symmetry group larger than the minimal Weyl symmetry can have more than one card diagram representation. Multiple diagrams exist when there is more than one equivalent way to choose $`(D2)`$ Killing congruences on the spacetime manifold. Examples we explicitly discuss in Sec. 3.2 are the 4d Witten bubble and the 4d S-Reissner-Nordstrøm (S-RN) which have three card diagrams corresponding to the three types of Killing congruences on dS<sub>2</sub> and $`𝐇_2`$. These different card diagrams are associated for example with global, patched, and Poincaré coordinates for dS<sub>2</sub> and the different representations will have different applications and reveal different information. On the other hand the S-Kerr solution, whose card diagram is discussed in Sec. 3.1.4, has symmetry group $`U(1)\times 𝐑`$ and its unique card representation looks like the ‘elliptic’ representation of S-RN.
#### 2.2.1 Our deck of cards: The building blocks for Weyl spacetimes
All spacetimes, new and old, in this paper are built from the following card types.
Horizontal cards are always half-planes. They may however have one or more branch cuts which may be taken to run perpendicular to the $`z`$-axis. Undoing one branch cut leads to a horizontal strip with two boundaries; multiple branch cuts lead to some open subset, with boundary, of a Riemann surface.
Vertical cards may be noncompact: Full planes with or without a pair of special null lines; half-planes with vertical or horizontal (horizon) boundary; or quarter planes at any $`45^{}`$ orientation. Vertical cards may also be compact: Squares with a pair of special null lines; or 45-45-90 triangles at any $`45^{}`$ orientation. All horizontal and vertical boundaries represent where a Killing circle vanishes and hence the end of the card. All $`45^{}`$ null boundaries represent instances of $`^\pm `$.
It is satisfying that for a variety of spacetimes including those in , the cards are always of the above rigid types.
There is one basic procedure which can be performed on vertical cards and their corresponding Weyl solutions. It is the analytic continuation $`2\gamma 2\gamma +i\pi `$, which is allowed since $`\gamma `$ is determined by first order PDEs and $`e^{2\gamma }`$ is real of either sign on vertical cards. This continuation is equivalent to multiplying the metric by a minus sign and then analytically continuing the $`D2`$ Killing directions. For charged generalizations of Weyl solutions (see Appendix B), this procedure does not affect the reality of the 1-form gauge field. We call this analytic continuation a $`\gamma `$-flip since the way it acts on a card is to geometrically flip it about a $`45^{}`$ null line (for example, look at the vertical cards in Figs. 19 and 19).
## 3 Card diagrams
In this section we construct the card diagrams for a wide assortment of solutions including black holes, bubbles and S-branes. The card diagrams are shown to be useful in representing continuous changes in the global spacetime structure such as how Reissner-Nordstrøm black holes change as we take their chargeless and extremal limits. For the superextremal black holes we discuss how to deal with branch points and cuts on horizontal cards. The card diagram also clearly represents the Kerr ring singularity and how traversing the interior of the ring leads to a second asymptotic spacetime. The 5d black ring solution, associated C-metric type solutions and twisted S-branes are also discussed.
Furthermore analytic continuation has an interesting interpretation in terms of card diagrams. We will describe the effect of analytic continuation on the card diagrams by examining two known analytic continuations of the Reissner-Nordstrøm black hole, the charged bubble and the S0-brane which we also call S-RN. These time dependent spacetimes each have three card diagram representation and two are obtained via different analytic continuations in Weyl coordinates. The Witten bubble and S-RN are related to each other by what we call a $`\gamma `$-flip which is a geometric realization of analytic continuation.
### 3.1 Black holes
#### 3.1.1 Subextremal $`Q^2<M^2`$ Reissner-Nordstrøm black holes
In the usual coordinates the Reissner-Nordstrøm black hole takes the form
$`ds^2`$ $`=`$ $`\left(1{\displaystyle \frac{2M}{r}}+{\displaystyle \frac{Q^2}{r^2}}\right)dt^2+\left(1{\displaystyle \frac{2M}{r}}+{\displaystyle \frac{Q^2}{r^2}}\right)^1dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)`$ (6)
$`A`$ $`=`$ $`Qdt/r`$
Using the coordinate transformation
$$\rho =\sqrt{r^22Mr+Q^2}\mathrm{sin}\theta ,z=(rM)\mathrm{cos}\theta $$
(7)
we find the Weyl form of Reissner-Nordstrøm black hole
$`ds^2`$ $`=`$ $`fdt^2+f^1(e^{2\gamma }(d\rho ^2+dz^2)+\rho ^2d\varphi ^2)`$ (8)
$`f`$ $`=`$ $`{\displaystyle \frac{(R_++R_{})^24(M^2Q^2)}{(R_++R_{}+2M)^2}}`$
$`e^{2\gamma }`$ $`=`$ $`{\displaystyle \frac{(R_++R_{})^24(M^2Q^2)}{4R_+R_{}}}`$
$`A`$ $`=`$ $`{\displaystyle \frac{2Qdt}{R_++R_{}+2M}}`$
$`R_\pm `$ $`=`$ $`\sqrt{\rho ^2+(z\pm \sqrt{M^2Q^2})^2}=rM\pm \sqrt{M^2Q^2}\mathrm{cos}\theta `$
and the card diagram for $`Q^2<M^2`$ is shown in Fig. 8. The construction of the card diagram proceeds along similar lines to the Schwarzschild card diagram. There are two adjacent horizontal half-planes, H1 and H2, which represent the positive mass asymptotically flat regions. The outer horizon at $`r_+=M+\sqrt{M^2+Q^2}`$ is represented in Weyl space as a rod which lies along the $`z`$-axis for $`\sqrt{M^2Q^2}<z<\sqrt{M^2Q^2}`$. The vertical cards, V1 and V2, are squares of length $`2\sqrt{M^2Q^2}`$ and the diagonal lines connecting opposite corners of the square are special null lines. The top of V2 is the $`r_{}=M\sqrt{M^2Q^2}`$ rod which is a four-card inner horizon. The black hole singularity no longer is on the edge of V2 but instead is on the outer boundary of the horizontal h1 and h2 regions, at $`\rho ^2/Q^2+z^2/M^2=1`$. Now the singularity is timelike and avoidable from the view of an observer on a vertical card. The rest of those horizontal cards, regions h3 and h4, are $`r<0`$ or equivalently $`M<0`$ nakedly singular RN spacetimes.
At each horizon, the card diagram is continued vertically to obtain an infinite tower of cards. In Fig. 8 we show the Penrose diagram for comparison.
Although the chargeless $`Q0`$ limit is hard to understand from Penrose diagrams, it is easy to understand using the card diagram in Fig. 8; the vertical card expands to a $`2M\times 2M`$ square and the singularity degenerates to a line segment coinciding with the inner horizon. Regions h1 and h2 disappear so the singularity is now ‘visible’ from V1 and V2 as well as h3 and h4; the singularity is spacelike relative to vertical cards and timelike for horizontal card observers. This achieves the Schwarzschild infinite array of cards in Fig. 6.
#### 3.1.2 Extremal $`Q^2=M^2`$ Reissner-Nordstrøm black hole
Starting from the above card diagram we now examine the extremal limit $`Q\pm M`$. In this case the vertical cards which represent the regions between the two horizons get smaller and disappear. When $`Q=M`$, the horizontal cards are now only attached at point-like extremal-horizons and only half of the horizontal cards remain connected, see Fig. 9. The region near the point-horizons are anti-de Sitter throats although cards themselves cannot adequately depict the throat region. The throat is a ‘connected’ sequence of points on vertically adjacent horizontal cards.
To understand how to go across the horizon, it is important to remember that there are special null lines in $`r_{}<r<r_+`$. In the extremal case the focal distances are equal ($`R=R_+=R_{}`$) and the null lines $`R_\pm =0`$ degenerate to the origin. So when we pass through the throat to an adjacent $`r<M`$ horizontal card, the sign changes for all occurrences of $`R=\sqrt{\rho ^2+z^2}=rM`$. Half the cards are now absent relative to the subextremal case since we no longer perform analytic continuation to get the vertical card. We can only go from positive real values of $`\rho `$ to negative values of $`\rho `$ alternatively. The singularity appears as a semicircle on the $`r<M`$ cards.
For axisymmetric Majumdar-Papapetrou solutions, this ‘sign change rule’ agrees with that in . Our analysis also applies to non-MP axisymmetric solutions such as the dihole .
#### 3.1.3 Superextremal $`Q^2>M^2`$ Reissner-Nordstrøm black holes
The superextremal $`Q^2>M^2`$ Reissner-Nordstrøm black hole does not have horizons or vertical cards. Its card diagram consists of two horizontal cards, connected along the branch cut $`0\rho \sqrt{Q^2M^2}`$, $`z=0`$. One card has a semi-ellipse singularity passing through the points $`(\rho =0,z=\pm M)`$ and $`(\rho =Q,z=0)`$ (see Fig. 10(a)).
These two horizontal cards are connected in the same sense as a branched Riemann sheet. By choosing Weyl’s canonical coordinates (meaning $`Z=\rho +iz`$ with $`(`$Coef$`dt^2)(`$Coef$`d\varphi ^2)=(`$Re$`Z)^2`$), the solution is no longer accurately represented on the horizontal card. This can be seen by examining the coordinate transformation (8) from Schwarzschild coordinates to Weyl coordinates. For fixed $`r`$ and varying $`\theta `$, the coordinates from $`M<r<\mathrm{}`$ cover the Weyl plane in semi-ellipses which degenerate to the segment $`(0\rho \sqrt{Q^2M^2},z=0)`$, which serves as a branch cut; and $`\rho =\sqrt{Q^2M^2}`$ is the branch point. The range $`\mathrm{}<r<M`$ again covers the half-plane with $`r=0`$ forming an ellipse singularity. Crossing the branch cut means choosing the opposite signs for $`R_\pm =\sqrt{\rho ^2(z\pm i\sqrt{Q^2M^2})^2}`$, and indeed we can think of the superextremal ‘rod’ as being complex-perpendicular to the Weyl $`Z`$-plane.
This double cover of the Weyl plane can be fixed by taking a holomorphic square root; this preserves the conformally Euclidean character of the card diagram. By choosing a new coordinate $`W=\sqrt{Z\sqrt{Q^2M^2}}`$, we map both the positive and negative-mass universes into the region $`(\mathrm{Im}W)^2(\mathrm{Re}W)^2\sqrt{Q^2M^2}`$ (see Fig. 10(b)). The image of the $`z`$-axis boundary is a hyperbola where the $`\varphi `$-circle vanishes. The origin $`W=0`$ is the image of the branch point and the image of the line segment $`(0\rho \sqrt{Q^2+M^2},z=0)`$ is a line connecting the two hyperbolas and intersecting the origin of the $`W`$-plane. The singular nature of $`e^{2\gamma }1/R_+R_{}1/|\mathrm{\Delta }Z|1/|W|^2`$ has been fixed by $`e^{2\gamma }dZd\overline{Z}=4|W|^2e^{2\gamma }dWd\overline{W}`$. Finally the ‘black hole’ singularity is mapped to a curved segment stretching from one hyperbola line to the other, to the left of the branch cut.
Alternatively one can use a Schwarz-Christoffel transformation to map the two universes onto the strip $`|\mathrm{Im}W|W_0`$. This is useful when the horizontal card boundaries are horizons as it allows for multiple horizontal cards to be placed adjacent to each other at the horizons (Fig. 17). This technique of fixing a horizontal card with a branch point will also be used in more complicated geometries such as the hyperbolic representation of S-RN and multi-rod solutions in four and five dimensions .
#### 3.1.4 Kerr black hole
Written in Weyl-Papapetrou form, the Kerr black hole is
$`ds^2`$ $`=`$ $`f(dt\omega d\varphi )^2+f^1(e^{2\gamma }(d\rho ^2+dz^2)+\rho ^2d\varphi ^2),`$ (9)
$`f`$ $`=`$ $`{\displaystyle \frac{(R_++R_{})^24M^2+\frac{a^2}{M^2a^2}(R_+R_{})^2}{(R_++R_{}+2M)^2+\frac{a^2}{M^2a^2}(R_+R_{})^2}},`$
$`e^{2\gamma }`$ $`=`$ $`{\displaystyle \frac{(R_++R_{})^24M^2+\frac{a^2}{M^2a^2}(R_+R_{})^2}{4R_+R_{}}},`$
$`\omega `$ $`=`$ $`{\displaystyle \frac{2aM(M+\frac{R_++R_{}}{2})(1\frac{(R_+R_{})^2}{4(M^2a^2)})}{\frac{1}{4}(R_++R_{})^2M^2+a^2\frac{(R_+R_{})^2}{4(M^2a^2)})}},`$
where $`R_\pm =\sqrt{\rho ^2+(z\pm \sqrt{M^2a^2})^2}=rM\pm \sqrt{M^2a^2}\mathrm{cos}\theta `$. The transformation to Boyer-Lindquist coordinates is $`\rho =\sqrt{r^22Mr+a^2}\mathrm{sin}\theta `$, $`z=(rM)\mathrm{cos}\theta `$.
For $`a^2<M^2`$ the Kerr black hole card diagram (see Figure 11) is similar to Reissner-Nordstrøm except that the singularity is now a point and lies at $`\rho =a`$, $`z=0`$ on each horizontal negative-mass card. The outer and inner ergospheres lie on the positive-and negative-mass cards and are both described by the curve $`z^2=\alpha ^2(\alpha ^2/a^21)\rho ^2\rho ^4/a^2`$ which intersects the rod endpoints at $`z=\pm \alpha =\pm \sqrt{M^2a^2}`$. The boundary of the region with closed timelike curves is also described by a quartic polynomial in Weyl coordinates. Once again the vertical cards have two special null lines where $`R_\pm `$ change sign.
The $`r=0`$ surface in BL coordinates is a semi-ellipse $`\rho ^2/a^2+z^2/M^2=1`$ on the negative-mass card; but it is not a distinguished locus on the card diagram. Attempting to make one loop around the ring in the Kerr-Schild picture clearly does not make a closed loop in Weyl space, whereas two loops in the Kerr-Schild picture will form a single loop around the ring singularity on the card diagram. It is also clear that it is possible to find classical trajectories which avoid the singularity and which safely escape into a second asymptotically flat region.
A card diagram for a charged Kerr-Newman solution can similarly be constructed, with $`\alpha =\sqrt{M^2Q^2a^2}`$.
The extremal Kerr(-Newman) solution has a card diagram like Fig. 9, but the ring singularity is just a point at $`z=0`$ and $`\rho =M`$ on negative-mass cards. Again, $`R=R_+=R_{}`$ and the special null lines degenerate to the origin; crossing the origin (which is a twisted AdS-type throat) entails changing the sign of $`R`$.
The superextremal Kerr(-Newman) solution is similar to the superextremal RN (Fig. 10(b)) except that the curved-segment singularity is replaced by a point ($`z=0`$), and the ergospheres map to an $`\mathrm{}`$-looking locus centered at $`W=0`$.
#### 3.1.5 The black ring
The 5d static black ring solution of is
$`ds^2`$ $`={\displaystyle \frac{F(x)}{F(y)}}dt^2`$
$`+`$ $`{\displaystyle \frac{1}{A^2(xy)^2}}\left[F(x)\left((y^21)d\psi ^2+{\displaystyle \frac{F(y)}{y^21}}dy^2\right)+F(y)^2\left({\displaystyle \frac{dx^2}{1x^2}}+{\displaystyle \frac{1x^2}{F(x)}}d\varphi ^2\right)\right]`$
where $`F(x)=1\mu x`$, $`F(y)=1\mu y`$, and $`0\mu 1`$. The coordinates $`x`$, $`y`$ are 4-focus (including $`\mathrm{}`$) or C-metric coordinates that parametrize a half-plane of Weyl space $`\rho 0`$, $`\mathrm{}<z<\mathrm{}`$:
$`\rho `$ $`=`$ $`{\displaystyle \frac{1}{A(xy)^2}}\sqrt{F(x)F(y)(1x^2)(1y^2)}`$
$`z`$ $`=`$ $`{\displaystyle \frac{(1xy)(F(x)F(y))}{2A(xy)^2}}.`$
The foci are on the $`z`$-axis at $`z=\pm \mu /2A`$ and $`z=1/2A`$. The black ring horizon is also on the $`z`$-axis along $`\mu /2Az\mu /2A`$. The $`\varphi `$-circle vanishes along $`z\mu /2A`$ and $`\mu /2Az1/2A`$, and the $`\psi `$-circle vanishes along $`z1/2A`$. Curves of constant $`y`$ degenerate to the horizon line segment as $`y\mathrm{}`$, and degenerate to the $`(1/2A,\mathrm{})`$ ray (better pictured with a conformally equivalent disk) for $`y1`$. Curves of constant $`x`$ degenerate to the vanishing $`\varphi `$-circle line segment for $`x1`$ and to the ray $`(\mathrm{},\mu /2A)`$ for $`x1`$.
The card diagram is easy to construct and is not much different from the four dimensional Schwarzschild case. Past $`y=\mathrm{}`$ we can go to $`y=+\mathrm{}`$ and hence imaginary $`\rho =i\rho ^{}`$, and move up a $`\mu /A\times \mu /A`$ square with two special null lines. At the top of the square, at $`y=1/\mu `$ we have the curvature singularity. Continuing again to real $`\rho `$ and running $`y`$ down to $`1`$, we map out a (negative-mass) horizontal card. The locus $`y=1`$ is the ray $`z1/2A`$. The space closes off here as the $`\psi `$-circle vanishes, but we formally continue to illustrate how C-metric coordinates run on noncompact vertical cards—this is useful in several applications, such as the Plebański-Demiański solution . Past $`y=1`$, we see that for fixed $`x`$, reducing $`y`$ down to $`x`$ makes a topological half-line in a vertical card with a special null line. Then for $`1<y<x`$, we traverse another vertical card, which we could attach to our original positive-mass horizontal card along $`z>1/2A`$ (see Fig. 12).
Note that when passing through the black ring horizon at $`y=\mathrm{}`$, the Weyl conformal factor
$$e^{2\nu }=\frac{1+\mu }{4A}\frac{Y_{23}}{R_1R_2R_3}\sqrt{\frac{Y_{12}}{Y_{13}}}\sqrt{\frac{R_2\zeta _2}{R_3\zeta _3}},$$
stays real; $`R_3\zeta _3`$ and $`Y_{13}`$ go negative. As we pass the special null lines, explicit appearances of $`R_1`$ and $`R_2`$ in the Weyl functions $`e^{2U_i}`$, $`e^{2\nu }`$ change sign.
The charged ring of is generated by a functional transformation and hence inherits a card diagram structure. In fact, any geometry with $`D2`$ Killing directions written in C-metric coordinates has a card diagram.
### 3.2 Charged Witten bubbles and S-branes
The Schwarzschild black hole can be analytically continued to two different time-dependent geometries, the Witten bubble of nothing with a dS<sub>2</sub> element and S-Schwarzschild with an $`𝐇_2`$ element, and it is instructive to understand how the card diagram changes.
Unlike black hole geometries with their unique card diagrams, these time dependent geometries can have multiple card diagram representations. Both the bubble and S-Schwarzschild have three different card diagram representations corresponding to three different ways to select Killing congruences. These three types of Killing congruences can be understood by representing $`𝐇_2`$ as the unit disk (with its conformal infinity being the unit circle). The orientation-preserving isometries of $`𝐇_2`$ are those Möbius transformations preserving the disk, $`PSL(2,𝐑)`$ . Möbius transformations $`z\frac{az+b}{cz+d}`$ have two complex fixed points, counted according to multiplicity. In the upper half-plane $`z=x+i\sigma `$ representation, $`a`$, $`b`$, $`c`$, and $`d`$ are real, so the fixed points are roots of a real quadratic. Hence they may be (i) distinct on the real boundary (hyperbolic), (ii) degenerate on the real boundary (parabolic), or (iii) nonreal complex conjugate pairs, one interior to the upper half-plane $`𝐇_2`$ (elliptic). Prototypes of Killing fields are (i) $`z(1+ϵ)z`$ for the upper half-plane, (ii) $`zz+ϵ`$ for the upper half-plane; and (iii) $`ze^{iϵ}z`$ for the disk $`|z|<1`$. These are the striped, Poincaré, and azimuthal congruences. In these hyperbolic, parabolic and elliptic representations, the S-Reissner-Nordstrøm (S-RN) and the Witten bubble each have 0, 1, and 2 Weyl foci.
#### 3.2.1 Elliptic representations and extended card diagrams
The bubble of nothing in $`D`$ dimensions has the interpretation as a semi-classical decay mode of the Kaluza-Klein vacuum. A spatial slice is topologically $`S^{D3}\times 𝐑^2`$, where the $`𝐑^2`$ is a cigar with the asymptotic-KK $`S^1`$ closing at some fixed $`r`$, which is an $`S^{D3}`$ bubble. As time passes, the bubble increases in size and ‘destroys’ the spacetime. The solution is obtained as an analytic continuation of a black hole and can be generalized to incorporate gauge fields.
The electrically charged bubble of nothing in its elliptic representation is gotten from (6) by sending $`tix^4`$, $`\varphi i\varphi `$, and to keep the field strength real we need $`QiQ`$. The metric is
$$ds^2=(1\frac{2M}{r}\frac{Q^2}{r^2})(dx^4)^2+(1\frac{2M}{r}\frac{Q^2}{r^2})^1dr^2+r^2(d\theta ^2\mathrm{sin}^2\theta d\varphi ^2)$$
(11)
At $`\theta =0,\pi `$ there are clearly Rindler-type horizons about which we analytically continue $`\theta `$ and obtain the rest of dS<sub>2</sub>, $`d\theta ^2+\mathrm{sinh}^2\theta d\varphi ^2`$. These six patches will precisely correspond to the six cards of Fig. 13.
Let us now turn to the effect of the analytic continuation of the Reissner-Nordstrøm black hole in Weyl coordinates. In Weyl coordinates, the effect of wick rotating $`tix^4`$ turns the horizon of the Schwarzschild card into an $`x^4`$-boundary, while $`\varphi i\varphi `$ turns the boundaries on the horizontal card into noncompact acceleration horizons along the rays $`|z|M`$, $`\rho =0`$. At the two horizons sending $`\theta 0\pm i\theta ,\pi \pm i\theta `$ corresponds to $`\rho \pm i\rho ^{}`$ along the $`|z|\sqrt{M^2+Q^2}`$ rays. We find vertical noncompact one eighth-plane cards with special null lines along $`\rho ^{}=zM`$ for $`zM`$ and $`\rho ^{}=Mz`$ for $`zM`$. Each piece of the plane is a doubly covered triangle and it is necessary to change branches of the function $`R_\pm `$ at the null line, as they appear in (8) to transform the card into a single covered quarter plane card.
The elliptic (and as we will shortly see the hyperbolic) representations of the Witten bubble are simply obtained because their dS<sub>2</sub> Killing congruences are trivially obtained from those on $`S^2`$. Specifically, take the $`S^2`$ embedding into flat space with $`Z=\mathrm{cos}\theta `$, $`X=\mathrm{sin}\theta \mathrm{cos}\varphi `$, $`Y=\mathrm{sin}\theta \mathrm{sin}\varphi `$. Sending $`\varphi i\varphi `$ has the effect $`YiY^{}`$ so the surface becomes $`X^2Y^2+Z^2=1`$ embedded in Minkowski space, or dS<sub>2</sub> with $`\varphi `$ as an elliptic (azimuthal) congruence. On the other hand sending $`\theta \pi /2+i\theta `$, has the effect $`ZiZ^{}`$ giving $`X^2+Y^2Z^2=1`$, which is dS<sub>2</sub> again but with $`\varphi `$ as a hyperbolic (striped) congruence.
Next the S-brane solution of
$$ds^2=(1+\frac{2M}{t}\frac{Q^2}{t^2})(dx^4)^2(1+\frac{2M}{t}\frac{Q^2}{t^2})^1dt^2+t^2(d\theta ^2+\mathrm{sinh}^2\theta d\varphi ^2)$$
(12)
can also be gotten from (6) by taking $`tix^4`$, $`\theta i\theta `$, $`rit`$, and $`MiM`$. From (7) we see that in Weyl’s coordinates this analytic continuation of RN can be implemented by sending $`tix^4`$, $`zi\tau `$, $`MiM`$, up to a real coordinate transformation.
The card diagram for elliptic S-RN (Fig. 14) has the same structure as the Witten bubble (Fig. 13) except that the 6-segment boundary is now $`\theta =0`$ where the $`\varphi `$-circle vanishes. The right and left horizons are at $`t=t_\pm =M\pm \sqrt{M^2+Q^2}`$. The $`t=0`$ singularity is a hyperbola on the $`t_{}tt_+`$ horizontal cards parametrized as $`(\rho ^{},\tau )=(|Q|\mathrm{sinh}\theta ,M\mathrm{cosh}\theta )`$. Any $`Q0`$ gives the same qualitative diagram. The ‘smaller’ connected universe on the card diagram is the negative-mass version of S-RN. Either sign of the mass gives a universe that is cosmologically singular. The Penrose diagram is given in Fig. 15.
In the limit $`Q0`$, the hyperbola singularity degenerates to a straight line covering the horizon at $`t=0`$. The two $`t>2M`$ vertical cards and two $`0<t<2M`$ horizontal cards then form a positive-mass S-Schwarzschild, while each $`t<2M`$ card forms a negative-mass S-Schwarzschild whose singularities begin or end the spacetime.
One can alternatively form the elliptic S-Schwarzschild from the elliptic Witten bubble by performing the $`\gamma `$-flip on any vertical card. This procedure is immediate; the net continuation from Schwarzschild is $`\theta i\theta `$, $`g_{\mu \nu }g_{\mu \nu }`$, and avoids $`rit`$, $`zi\tau `$, and $`MiM`$.
Note how the card representation of the S-brane is quite different from the black hole card diagram while the Penrose diagrams of the two spacetimes are nearly just related by ninety degree rotation. This is because the card diagram shows the compact or noncompact $`\theta `$ direction.
The elliptic form of the card diagrams show that Schwarzschild S-brane, Witten bubble and Schwarzschild solutions have similar structures and in fact they are all related by $`\gamma `$-flips and trivial Killing continuations.<sup>2</sup><sup>2</sup>2Perturbed solutions that can only be obtained from $`zi\tau `$ are considered to be less trivial. Solutions which are related in this manner may be conveniently drawn together in one diagram which simultaneously displays all of their card diagrams. For example in Fig. 16 the S-Schwarzschild solution comprises regions $`1,2,3,4,5`$, the Witten bubble regions $`4,5,6`$, and the Schwarzschild black hole $`6,7,8,9,10`$. Regions $`1,2,10`$ correspond to a singular Witten bubble of negative ‘mass.’ In this diagram we also see that the unification of the special null lines as they extend through all three solutions; we can say $`R_+=0`$ is the long $`\mathrm{}`$-null line and $`R_{}=0`$ is the long $`\mathrm{}`$-null line.
The charged Reissner-Nordstrøm BH/bubble/S-brane solutions cannot be depicted together on such a diagram because $`QiQ`$ changes $`0<r_{}<r_+`$ to $`r_{}<0<r_+`$. Similar diagrams can be found in .
#### 3.2.2 Hyperbolic representations and branch points
The charged Witten bubble
$$ds^2=\left(1\frac{2M}{r}\frac{Q^2}{r^2}\right)(dx^4)^2+\left(1\frac{2M}{r}\frac{Q^2}{r^2}\right)^1dr^2+r^2(d\theta ^2+\mathrm{cosh}^2\theta d\varphi ^2)$$
(13)
can alternatively be obtained from the RN black hole (6) by taking $`\theta \pi /2+i\theta `$ and $`tix^4`$, $`QiQ`$. Here, $`\theta `$ plays the role of time and $`\theta =0`$ is the time where the bubble ‘has minimum size.’ (This statement has meaning if we break $`SO(2,1)`$ symmetry.) To achieve this in Weyl’s coordinates, we put $`zi\tau `$, $`tix^4`$, $`QiQ`$; the resulting space is equivalent to Witten’s bubble by the real coordinate transformation
$$\rho =\sqrt{r^22MrQ^2}\mathrm{cosh}\theta ,\tau =(rM)\mathrm{sinh}\theta .$$
(14)
Thus in Weyl coordinates the only difference between the hyperbolic Witten bubble and the elliptic S-RN is putting $`MiM`$. Witten’s bubble universe is represented in Weyl coordinates as a vertical half-plane card, $`\rho 0`$, $`\mathrm{}<\tau <\mathrm{}`$, where now the $`x^4`$-circle, and not the $`\varphi `$-circle, vanishes at $`\rho =0`$. This boundary is where the bubble of nothing begins. Note that the vertical card now has Minkowski signature and is conformal to $`d\tau ^2+d\rho ^2`$. This vertical card does not have special null lines since the foci are at imaginary values $`\tau =\pm iM`$, and so the spacetime is covered only once by the Schwarzschild coordinates. The bubble does have a rod which is along the imaginary $`\tau `$ axis and which intersects the card at the $`\rho =0`$, $`\tau =0`$ origin. The hyperbolic representation of the charged Witten bubble is therefore just a vertical half-plane. In such a case the card diagrams have just as much causal information as a Penrose diagram and the boundary of the card easily describes the edge of the bubble of nothing.
We obtain a hyperbolic representation for S-RN from (6) by sending $`tix^4`$, $`MiM`$, $`rir`$, $`\theta \pi /2+i\theta `$, and $`\varphi i\varphi `$. The fibered directions are now hyperbolic $`d𝐇_2^2=d\theta ^2+\mathrm{cosh}^2\theta d\varphi ^2`$. In Weyl coordinates, to maintain reality of the solution, we begin with RN (8) and send $`tix^4`$, $`\varphi i\varphi `$, $`MiM`$, and must explicitly change the branch of the square root introducing the minus sign $`R_{}R_{}`$.<sup>3</sup><sup>3</sup>3The naturalness of this sign change is explained in great detail in . This has the interpretation of staying on the same horizontal card and rotating the rod in the complex $`z`$-plane. The foci are then at ($`z=\pm i\sqrt{M^2+Q^2}`$, $`\rho =0`$) with their special null lines intersecting the real half-plane at $`z=0`$, $`\rho =\sqrt{M^2+Q^2}`$. The half-plane is doubly covered, and we will take $`0\rho \sqrt{M^2+Q^2}`$ as the branch cut. The sign change of $`R_{}`$ has effectively reversed the roles of $`r`$ and $`\theta `$ so that, after undoing the branch cut with say a square-root conformal transformation, $`r=r_+>0`$ is one boundary-horizon and $`r=r_{}<0`$ is the other. The hyperbolic angle $`\theta `$ is unbounded $`\mathrm{}<\theta <\mathrm{}`$.
The doubly-covered half-plane is physically cut into two by the $`r=0`$ singularity. At each horizon $`r=r_\pm `$ we have a four-card junction; the double half-plane horizontal card meets another mirror copy as well as two vertical cards. The vertical cards are at $`\rho \pm i\rho ^{}`$ and $`\mathrm{}<z<\mathrm{}`$ and have no special null lines or other features. A full card diagram is shown in Fig. 17. Note that there are no boundaries of this card diagram where a spacelike Killing circle vanishes.
Setting $`R_\pm =\sqrt{\rho ^2+(z\pm i\sqrt{M^2+Q^2})^2}`$ and $`R=R_+`$, the explicit form of hyperbolic S-RN on the horizontal card is
$`ds^2`$ $`=`$ $`f(dx^4)^2+f^1(e^{2\gamma }(d\rho ^2+dz^2)+\rho ^2d\varphi ^2),`$ (15)
$`f`$ $`=`$ $`{\displaystyle \frac{M^2+Q^2(\mathrm{Im}R)^2}{(\mathrm{Im}R+M)^2}},`$
$`e^{2\gamma }`$ $`=`$ $`{\displaystyle \frac{M^2+Q^2(\mathrm{Im}R)^2}{|R|^2}},`$
$`A`$ $`=`$ $`{\displaystyle \frac{Qdx^4}{\mathrm{Im}R+M}}.`$
We can arrive at this spacetime in a simpler way. Take the RN black hole and analytically continue to get the hyperbolic charged Witten bubbles. These universes are nonsingular for $`rr_+`$ or $`rr_{}`$. They have boundaries at $`r=r_\pm `$ where the $`x^4`$-circle vanishes. Now, turn these universes on their sides with the $`\gamma `$-flip. This allows us to decompactify $`x^4`$ and $`r=r_\pm `$ are now Milne horizons—we are looking precisely at the vertical cards of Fig. 17, and they are connected in a card diagram by an $`r_{}rr_+`$ card which is now accessible. We see that generally, vertical half-plane cards parametrized in spherical prolate fashion with no special null lines, when turned on their sides, connect to branched horizontal cards.
The $`Q0`$ limit (hyperbolic S-Schwarzschild) of the card diagram is easy to picture: The singularities of Fig. 17 collapse onto the $`r=r_{}`$ horizon.
One may wonder what happens if we take $`zi\tau `$, on the horizontal card for hyperbolic S-RN, and achieves a vertical card sandwiched between special null lines at $`\rho =\sqrt{M^2+Q^2}+|\tau |`$. This must unfold to give a vertical plane with two intersecting special null lines. This card diagram structure is discussed in .
#### 3.2.3 Parabolic representations
There is a third way to put a Killing congruence on $`𝐇_2`$ or dS<sub>2</sub> using parabolic or Poincaré coordinates. Parametrizing hyperbolic space (Euclideanized AdS<sub>2</sub>) as $`ds^2=\sigma ^2dx^2+\frac{d\sigma ^2}{\sigma ^2}`$, and keeping the Schwarzschild S-brane coordinate $`r`$ and the usual $`x^4`$ we get a Poincaré Weyl representation of S-RN spacetime . It is
$`ds^2`$ $`=`$ $`f(dx^4)^2f^1(e^{2\gamma }(d\rho ^2+dz^2)+\rho ^2dx^2),`$ (16)
$`A`$ $`=`$ $`Qdx^4/r,`$
$`f`$ $`=`$ $`(12M/rQ^2/r^2),`$
$`e^{2\gamma }`$ $`=`$ $`{\displaystyle \frac{r^22MrQ^2}{\sigma ^2(M^2+Q^2)}},`$
$`\rho ^{}`$ $`=`$ $`\sigma \sqrt{r^22MrQ^2},`$
$`z`$ $`=`$ $`\sigma (rM).`$
In this Weyl representation, $`\rho ^{}`$ is timelike on $`rr_+`$ vertical cards which are noncompact $`45^{}`$ wedges, $`0\pm \rho ^{}<z`$. This connects along $`z>0`$ to an $`r_{}rr_+`$ horizontal card; $`rr_{}`$ vertical cards attaches to $`z<0`$. So this is similar to the elliptic representation of S-RN, except the line segment $`\sqrt{M^2+Q^2}<z<\sqrt{M^2+Q^2}`$ has collapsed and the special null lines are now conformal null infinity (Fig. 19). The singularity on the horizontal card is particularly easy to describe in these coordinates. On the first horizontal card it is on a ray $`z/\rho =M/Q`$ and on the second card it is on a ray $`z/\rho =M/Q`$.
If we take the $`rr_+`$ (or $`rr_{}`$) $`45^{}`$ wedges and turn them on their sides via the $`\gamma `$-flip, we get the parabolic version of the $`rr_+`$ (or $`rr_{}`$) charged Witten bubble. The line which used to be the horizon in the S-brane card diagram becomes a boundary which is the ‘minimum volume’ sphere, at $`\rho =0`$. Time is now purely along the $`\tau `$ direction as in the hyperbolic Witten bubble. The special null line is still $`^\pm `$ since $`\rho =|\tau |`$ corresponds to $`r\mathrm{}`$. The vertex of the triangular card is not the end of the spacetime. These wedge cards only represent $`\sigma >0`$ and so the card diagram should be extended to $`\sigma <0`$. The card diagram is an infinite array of $`45^{}`$ wedge cards pointing up and an infinite number pointing down. The vertex of each upward card is attached to its nearest two downward neighbors (one to the left and one to the right), in the dS<sub>2</sub> fashion as shown in Figure 19. One can identify cards so only needs one upward and one downward card with two attachments. Although this card diagram is not the most obvious representation of the Witten bubble, it is useful in understanding the S-Kerr solution of the next subsection as well as the more complicated S-dihole $`𝒰`$, $`𝒰_\pm `$ universes of .
### 3.3 S-Kerr
The twisted S-brane , see also , is also known as S-Kerr, and it is another example of a nonsingular time-dependent solution. Twisted S-branes describe the decay of unstable massive strings. They can be obtained from the Kerr black hole using the following card diagram method. Double Killing continue $`tix^4`$, $`\varphi i\varphi `$ to achieve a $`𝒦_+`$ bubble solution, go to the vertical card via $`\theta i\theta `$, then perform a $`\gamma `$-flip to achieve S-Kerr. S-Kerr has symmetry group $`U(1)\times 𝐑`$ and therefore has a unique card diagram. For the parameter range $`a^2<M^2`$ there are horizons and the card diagram structure is that of the elliptic bubble (Fig. 13). The foci are at $`z=\pm \sqrt{M^2a^2}`$. There is an ergosphere and CTC region on the horizontal card and has the same qualitative shape as it does for the Kerr black hole diagram. In comparison, the Penrose diagram showing (variably twisted) $`x^4`$ and the Boyer-Lindquist coordinate $`r`$ is shown in Fig. 20.
In the extremal limit $`\pm aM`$, $`\theta `$-orbits on the vertical card shift up relative to the special null line, and any fixed ($`r`$, $`\theta `$) point is sent above the null line. The region below the null line disappears in this limit and the horizontal card collapses to a point. Furthermore, those geodesics in the upper-right card can only reach the lower-left card (and the same with upper-left and lower-right), splitting the universe into connected $`45^{}`$ wedges just like the parabolic representation for the Witten bubble (Fig. 19), where the connections are dS<sub>2</sub>-like twisted throats.
The case $`a^2>M^2`$ for S-Kerr does not have horizons and can be represented as a single vertical half-plane card with $`\rho 0`$ and with no special null lines. Just like the hyperbolic Witten bubble, superextremal S-Kerr can be turned on its side with the $`\gamma `$-flip, to yield a new spacetime . In the limit where $`|a|\mathrm{}`$, the $`\theta `$-orbits flatten out on the card, and the solution becomes flat space.
The Penrose diagrams for extremal and superextremal S-Kerr are given in Fig. 21.
## 4 Discussion
In this paper we have examined and extended the utility of the Weyl Ansatz through the construction of an associated Weyl card diagram. Weyl coordinates are an excellent choice for clearly understanding geometric features and so drawing the card diagram is useful. The card diagram conveniently captures most of the interesting properties of a spacetime including its singularities, horizons, and some aspects of its causal structure and null infinity. Card diagrams for families of solutions such as charged and rotating black holes (and bubbles and S-branes) share similar features and can change continuously. They are useful in keeping track of various analytic continuations and mentally partitioning complicated spacetimes into simpler regions. The only technical issues that seem to arise and which we resolved dealt with branches of Weyl distance square-root functions, special null lines on vertical cards, and branched horizontal cards.
Here we summarize the solutions in this paper. The card diagrams correctly capture the different regions of the charged Reissner-Nordstrøm black hole and its various charged and chargeless limits, and its negative mass complement. We also analyzed the Kerr black hole and its singularity structure. In particular the safe passage through the interior of the ring to the second asymptotic universe through $`r=0`$ is clearly depicted.
The Witten bubbles and S-branes each have three card diagrams representations corresponding to the three different choices of Killing congruences on dS<sub>2</sub> or $`𝐇_2`$. The elliptic representations had two foci and six cards. S-RN had a cosmological singularity on the horizontal card splitting the card diagram into two connected universes, whereas the charged bubbles have the same card structure but are nonsingular. The hyperbolic representations have no foci: the bubble is a simple vertical half-plane card and hyperbolic S-RN has a branched horizontal card which we fixed by a conformal mapping of the half-plane. Finally the parabolic representation of the bubble was an infinite array of $`45^{}`$ triangle wedges connected pointwise while S-RN had a 6-card butterfly shape. This parabolic representation showcased special null lines serving as null infinity.
S-Reissner Nordstrøm can be obtained from the bubble in two ways. One may start with the bubble and analytically continue $`MiM`$ in Weyl coordinates. In this case the hyperbolic/elliptic representation of the bubble maps to the elliptic/hyperbolic representation of the S-brane. We also found that these two solutions can be related by what we called the $`\gamma `$-flip, which is conveniently visualized as a flip of the associated cards about a null line. This procedure maintains the number of Weyl foci on the card and so maps the elliptic/parabolic/hyperbolic bubble to the elliptic/parabolic/hyperbolic S-brane. The $`\gamma `$-flip provides a simple and geometric way to relate Schwarzschild with its two analytic continuations, the bubble of nothing and the S-brane. In fact spacetimes related by $`\gamma `$-flips can be simultaneously drawn together in a complexified $`r\theta `$ (or $`\rho ,z`$) spacetime diagram.
Just as we used the $`\gamma `$-flip to turn the hyperbolic Witten bubble on its side and got hyperbolic S-RN, we can take the vertical half-plane card diagrams for the Kerr bubble, dihole wave, superextremal S-Kerr and superextremal S-dihole and apply the $`\gamma `$-flip to yield new spacetimes. These new solutions will be described in .
The card diagram formalism can be further generalized; the recently developed Weyl-Papapetrou formalism for $`D5`$ will yield card diagrams and 5d Kerr-related solutions will appear in . Furthermore card diagrams do not require Weyl’s canonical coordinates. Spacetimes with Weyl-type symmetry and yet where Weyl’s procedure fails algebraically can still admit card diagrams. An example is the inclusion of a nonzero cosmological constant $`\mathrm{\Lambda }`$, where a $`\gamma `$-flip changes the sign of $`\mathrm{\Lambda }`$. Card diagrams for pure (A)dS<sub>D</sub> space for $`D=4,5`$ are presented in . Constant-curvature black holes obtained by quotienting will also have card diagrams. We hope that our methods, and their further generalizations, will have even greater applicability than to the multitude of spacetimes already discussed herein.
## Acknowledgements
We thank D. Jatkar, A. Maloney, W. G. Ritter, A. Strominger, T. Wiseman and X. Yin for valuable discussions and comments. G. C. J. would like to thank the NSF for funding. J. E. W. is supported in part by the National Science Council, the Center for Theoretical Physics at National Taiwan University, the National Center for Theoretical Sciences and would like to thank the organizers of Strings 2004 for support and a wonderful conference where part of this research was conducted.
## Appendix A Perturbed bubbles and S-branes
The chargeless Schwarzschild black hole is easily perturbed as a Weyl solution by adding more rod-horizons, to form Israel-Khan arrays. We can use these solutions to smoothly perturb the Witten bubble in two different ways and also the S-Schwarzschild in two different ways. Addition of charge can be done via Weyl’s electrification method .
The analytic continuation in Weyl space to obtain the hyperbolic Witten bubble is precisely the same as in , except here the Schwarzschild rod crosses $`z=0`$. Even-in-$`z`$ Israel-Khan arrays where no rod crosses $`z=0`$ can be analytically continued to gravitational wave solutions sourced by imaginary black holes, ie rods at imaginary time. We can thus generalize the Witten bubble by adding additional waves by symmetrically placing more rods in addition to the one which crosses $`z=0`$. We dub such an array the ‘hyperbolic-perturbed Witten bubble.’ As these additional rods are made to cover more of the $`z`$-axis and are brought closer and closer to the principal rod, the deformed Witten bubble solution hangs longer with a minimum-radius $`\varphi `$-circle. In the limit where rods occupy the entire $`z`$-axis, we get a static flat solution, which is Minkowski 3-space times a fixed-circumference $`\varphi `$-circle.
The hyperbolic-perturbed Witten bubble can be turned on its side, yielding a hyperbolic-perturbed S-Schwarzschild. It has the card diagram structure of Fig. 17.
We can also perturb S-Schwarzschild by adding rods before analytically continuing $`zi\tau `$, $`MiM`$. We dub this the ‘elliptic-perturbed S-Schwarzschild.’ It is different than hyperbolic-perturbed S-Schwarzschild, and has the card diagram structure of Fig. 14. Turned on its side, it yields an elliptic-perturbed Witten bubble, with the card diagram structure of Fig. 13. It is different than the hyperbolic-perturbed Witten bubble.
In any of the cases, we can choose to analytically continue the mass parameters of the additional rods or not. Additionally, we can displace some rods in the imaginary $`z`$-direction which affects the $`\tau `$-center of their disturbance. If we do everything in an even fashion, i.e. we respect $`\mathrm{Im}\tau \mathrm{Im}\tau `$, the resulting geometry (at real $`\tau `$) will be real. In particular, rotating a rod at $`z>0`$ counterclockwise means rotating its image at $`z<0`$ clockwise. We see in the discussion of the 2-rod example that there may be several choices for branches.
The card diagram techniques allow us to easily construct these two inequivalent families of perturbed Witten bubble and perturbed S-Schwarzschild solutions. These and other multi-rod, S-dihole, and infinite-periodic-universe solutions are described in and cannot be easily described or understood without Weyl coordinates and the construction and language of card diagrams. The nontrivial $`zi\tau `$ continuation is essential.
## Appendix B Appendix: Electrostatic Weyl formalism
The formalism of can be extended for general $`D`$ to include an electrostatic potential. This is somewhat surprising since the electromagnetic energy-momentum tensor
$$T_{\mu \nu }=F_{\mu \rho }F_\nu {}_{}{}^{\rho }\frac{1}{4}g_{\mu \nu }F^2$$
is traceless only in $`D=4`$ and so Einstein’s equations are more complicated. Nevertheless, a cancellation does occur and one may sum the diagonal Killing frame components of the Ricci tensor to achieve a harmonic condition.
Follow the notation of and add a 1-form potential $`A(Z,\overline{Z})dt`$ where $`t=x^1`$ is timelike ($`ϵ_1=1`$) and all other $`x^i`$, $`i=2,\mathrm{},D2`$ are spacelike (with $`ϵ_i=+1`$). The metric takes the form
$$ds^2=e^{2U_1}dt^2+\underset{i=2}{\overset{D2}{}}e^{2U_i}(dx^i)^2+e^{2C}dZd\overline{Z},$$
from which we extract the frame metric
$$g_{\widehat{\mu }\widehat{\nu }}=\mathrm{diag}(1,+1,\mathrm{},+1)\left[\begin{array}{cc}0& 1/2\\ 1/2& 0\end{array}\right].$$
For $`F=dA`$ we have $`F_{\widehat{Z}\widehat{t}}=F_{\widehat{t}\widehat{Z}}=Ae^{U_1C}`$ and $`F_{\widehat{\overline{Z}}\widehat{t}}=F_{\widehat{t}\widehat{\overline{Z}}}=\overline{}Ae^{U_1C}`$, all other components vanishing. We compute $`F^2=8A\overline{}Ae^{2U_12C}`$ and
$`T_{\widehat{t}\widehat{t}}`$ $`=`$ $`2A\overline{}Ae^{2U_12C}`$
$`T_{\widehat{i}\widehat{i}}`$ $`=`$ $`2A\overline{}Ae^{2U_12C}(i1)`$
$`T_{\widehat{Z}\widehat{Z}}`$ $`=`$ $`(A)^2e^{2U_12C}`$
$`T_{\widehat{\overline{Z}}\widehat{\overline{Z}}}`$ $`=`$ $`\overline{T_{\widehat{Z}\widehat{Z}}}`$
$`T_{\widehat{Z}\widehat{\overline{Z}}}`$ $`=`$ $`0.`$
The field equations are $`R_{\widehat{\mu }\widehat{\nu }}\frac{1}{2}g_{\widehat{\mu }\widehat{\nu }}R=T_{\widehat{\mu }\widehat{\nu }}`$; taking the trace, we get
$$R=\frac{4(D4)}{D2}A\overline{}Ae^{2U_12C}$$
and Einstein’s equations are then
$$R_{\widehat{\mu }\widehat{\nu }}=T_{\widehat{\mu }\widehat{\nu }}\frac{2(D4)}{D2}g_{\widehat{\mu }\widehat{\nu }}A\overline{}Ae^{2U_12C}.$$
(17)
Form the sum $`_{i=1}^{D2}R_{\widehat{i}\widehat{i}}ϵ_i`$; the right side of (17) gives
$$(D4)2A\overline{}Ae^{2U_12C}\frac{2(D4)}{D2}(D2)A\overline{}Ae^{2U_12C}=0.$$
Hence (following (2.4)-(2.5) of ) we get
$$\overline{}\mathrm{exp}\left(\underset{i=1}{\overset{D2}{}}U_i\right)=0,$$
the Weyl harmonic condition.
One can add magnetostatic potentials along spatial Killing directions as well. We skip remaining details and give the equations. Let us assume $`x^1`$ is timelike and $`x^i`$ are spacelike for $`i=2,\mathrm{},D2`$, the potential is $`A_\mathit{1}=_{i=1}^{D2}A_idx^i`$, and the metric is $`ds^2=e^{2U_1}(dx^1)^2+_{i=2}^{D2}e^{2U_i}(dx^i)^2+e^{2\nu }(d\rho ^2+dz^2)`$ and $`w=\rho +iz`$, $`_w=\frac{1}{2}(_\rho i_z)`$. Einstein’s equations are
$`\mathrm{\Delta }U_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}({\displaystyle \underset{i=1}{\overset{D2}{}}}(A_i)^2e^{2U_i}+{\displaystyle \frac{D4}{D2}}{\displaystyle \underset{i=1}{\overset{D2}{}}}(A_j)^2e^{2U_i},`$
$`\mathrm{\Delta }U_k`$ $`=`$ $`{\displaystyle \frac{1}{2}}((A_1)^2e^{2U_1}(A_k)^2e^{2U_k}+{\displaystyle \underset{ik,1}{}}(A_i)^2e^{2U_i}`$
$`{\displaystyle \frac{D4}{D2}}{\displaystyle \underset{i1}{}}(A_i)^2e^{2U_i}+{\displaystyle \frac{D4}{D2}}(A_1)^2e^{2U_1}),`$
and
$`_w{\displaystyle \underset{i=1}{\overset{D2}{}}}\nu =2\rho ({\displaystyle \underset{i<j}{}}_wU_i_wU_j+{\displaystyle \frac{(_wA_1)^2e^{2U_1}}{2}}{\displaystyle \underset{i=2}{\overset{D2}{}}}{\displaystyle \frac{(_wA_i)^2}{2}}e^{2U_i}).`$
Maxwell’s equations are
$$\left(A_ie^{2U_i}\right).$$
All Laplacians and divergences are with respect to a flat 3d axisymmetric auxiliary space with coordinates $`\rho ,z`$.
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# Local to Global Normalization Dynamic by Nonlinear Local Interactions
## I Introduction
What is the difference between adaptation and normalization? Are these just two distinct processes, or can they be related? The purpose of this paper is to develop a model whose dynamic smoothly proceeds from local adaptation to global normalization. Mathematical properties of the model are analyzed, and its dynamical properties are evaluated with luminance images. I study the model within the framework of biological vision, where emphasis is laid on understanding the emergence of adaptation within the model’s dynamic. Finally, a method is proposed for freezing the dynamic at the moment when adaptation occurs. But to begin with, I briefly describe how adaptation and normalization contribute to information processing in the brain.
Adaptation refers to the adjustment of a sense organ to the intensity or quality of stimulation Merriam-Webster (2007). There is agreement that adaptation is important for the function of nervous systems, since without corresponding mechanisms any given neuron with its limited dynamic range would stay silent or operate in saturation most of the time Walraven et al. (1990). When considering a population of *cells* (e.g. formal processing units or biological neurons), then adaptation is usually understood as a locally acting process, which can be carried out independently for individual cells or groups of cells, respectively (e.g., individual photoreceptors Carpenter and Grossberg (1981); van Hateren (2005); Keil and Vitrià (2007) vs. groups of photoreceptorsGross and Brajovic (2003); Hong and Grossberg (2004)). Thus, adaptation refers to sensitivity adjustment of output signals as a function of input signals.
Normalization on the other hand usually refers to establishing standardized conditions for one or more qualities. For example, at some stage in the brain, the retinal image may have been normalized with respect to illumination conditions, such that each face or object is represented to have similar illumination patterns, and subsequent recognition stages work in a more robust fashion. Or, once a face image has been detected by an artificial face recognition system, it can be normalized with respect to head tilt or head rotation. In this way a standardized candidate face is obtained, which facilitates matching it to other standardized faces from a database.
Normalization is also used for describing the establishment of standardized conditions for a population of neurons. In this context, normalization processes usually act as gain control mechanisms. For instance, Grossberg Grossberg (1983) proposed “shunting competitive networks” (in his terms) for accurate signal processing in the presence of noise to avoid the noise-saturation dilemma. Because neurons have a fixed input range, weak signals get masked by noise, and neurons’ signal only the noisy fluctuations in the input signal. On the other hand, strong signals cause neurons to saturate, and any variations within the input cannot be distinguished. Shunting networks implement the multiplicative relationship between membrane voltages of neurons and conductance changes that are caused by network input on the one hand and signals on the other. This multiplicative relationship acts as a gain control mechanism that enables these networks to automatically re-tune their sensitivity in response to fluctuating background inputs. As Grossberg demonstrated Grossberg (1983), such networks exhibit a normalization property in the sense that the total (or pooled) activity of all neurons is independent of the number of neurons. Along these lines, Heeger and co-workers proposed a normalization model to account for the observed non-linearities with the cortical simple cell responses, such as response saturation and cross-orientation inhibition Carandini et al. (1997); Carandini and Heeger (1994); Heeger (1992). Similar to Grossberg’s “shunting competitive networks”, in Heeger’s model a neuron’s output activity is adjusted by the pooled activity of a population of many other neurons (“normalization pool”). This “normalization pool” exerts divisive inhibition on the response of a target neuron, and in this way it acts as a gain control mechanism for that cell.
The circuit models proposed by Grossberg and Heeger describe how responses of a group of neurons can be normalized. Both methods rely on the interaction of some target neuron with a number of surrounding neurons. The interaction is brought about by hard-wiring the target neurons with surrounding neurons. In contrast, the normalization scheme introduced in this paper is based on diffusion mechanisms, and thus interactions only take place between adjacent cells. Specifically, within the scope of the present paper, normalization is understood as mapping a set of numbers with finite but in principle arbitrary numerical range onto a fixed target range (below we will see that non-trivial features like contrast enhancement and adaptation phenomena emerge from a network which implements this normalization mechanism).
Whereas in Grossberg’s scheme the normalization process renders the total activity of a group of cells independent of the number of cells (Grossberg (1983)), with my definition of normalization it is clear that in most cases the activity summed over all cells will depend on their number. A further difference concerns the implementation of activity bounds. In Grossberg’s scheme, reversal potentials establish an upper (lower) bound on the activity of each cell which can be reached by excitation (inhibition). However, the highest activity value of the normalized cell population depends on the activity of all other cells (as the total activity is constant). In other words, one cannot rely on the presence of distinguished activity values as it is the case in my approach. In a normalized population of my approach there is always at least one cell which has zero activity, and at least one cell with activity one.
The usual proceeding for normalizing a set of numbers can be subdivided into two successive stages. First, the maximum and the minimum members of the set are determined. These two values are then used in a second stage for re-scaling all set elements such that after re-scaling the elements fall into a pre-defined numerical interval (or numerical range).
If we wish to design a corresponding algorithm for the first stage of the just described process (i.e. finding the maximum and the minimum), we would have to employ two memories for storing the *current* (i.e., a local) maximum and minimum, and compare these values successively with all remaining set elements. After we finished with comparing, the memory would contain the *global* maximum and minimum. Because every set member has to interact explicitly with the memories, the whole process is said to involve global operations. The global nature is mirrored in the connection structure of a correspondingly designed network. Figure 1(a) shows a schematic drawing of such a network, where one distinguished network unit shares connections with all the others. This unit is supposed to represent the maximum (or minimum) activity value of the set of units to which it is connected to. Due to its global connectivity, however, our network seems not to be a very plausible candidate for a “biologically” model, because (biological) neurons are known to interact in a more local fashion. This implausibility can be relaxed by proposing an alternative connectivity pattern (figure 1(b)).
Nevertheless, the two units representing the maximum and the minimum, respectively, need to interact subsequently again with the input units, in order to put into effect the re-scaling operation that implements the gain control mechanism. This means that one would require yet another set of non-local connections, analogously to the pattern shown in figure 1. This led me to the question whether such normalization can be achieved in a more “biological” or local fashion, or even by employing only interactions between adjacent network units. Presuming the existence of corresponding mechanisms, one has to explore in addition whether they could, in principle, be carried out by nerve cells in a biophysically plausible way.
Below I present a network (the *dynamical normalization network*), which achieves normalization by means of lateral propagation of activity between adjacent network cells. To this end, parameterized diffusion operators were developed. In their limit cases, these operators implement non-linear and non-conservative diffusion processes (“pseudo-diffusion”). The dynamic of pseudo-diffusion proceeds from local to global in a continuous fashion, without utilizing any connectivity structure apart from coupling among nearest neighbors.
The dynamic normalization network consists of a total of four layers: an input layer, two diffusion layers, and the normalization or output layer, where all layers interact. Numerical simulations with luminance images revealed that the dynamic of the normalization layer is functionally more rich than just performing a re-scaling of its input. Initially, the dynamic reveals contrast enhancement similar to high-pass filtering.
Furthermore, under certain conditions, an adaptation phenomenon (“dynamic compression”) can be observed in the initial phase of the dynamic. As it is described in detail below (section III.2), the strength of the dynamic compression effect depends on the size of high activity regions in the input, and their relative positions with respect to other local maxima.
## II Formal Description of Nonlinear Diffusion
The dynamic normalization network is based on nonlinear diffusion operators. In order to proof some of their properties, it is necessary that the nonlinear diffusion operators are differentiable. Accordingly, we define at first an operator $`𝒯_\lambda []`$ which is parameterized over $`\lambda `$ as
$$𝒯_\lambda [x]=\frac{\eta x}{1+e^{\lambda x}}$$
(1)
where $`\eta `$ is a normalization constant that is defined as
$$\eta =1+e^{|\lambda |}.$$
(2)
Through the specific choice of $`\lambda `$, we can “steer” the operator $`𝒯_\lambda []`$ continuously from linearity ($`𝒯_0𝒯_{\lambda =0}`$)
$$\lambda =0𝒯_0[x]=x$$
(3)
to half wave rectification (i.e. selection of the maximum between zero and its argument)
$$\underset{\lambda +\mathrm{}}{lim}T_\lambda [x]\mathrm{max}(x,0)$$
(4)
or inverse half wave rectification (i.e. selection of the minimum between zero and its argument)
$$\underset{\lambda \mathrm{}}{lim}T_\lambda [x]\mathrm{min}(0,x).$$
(5)
Notice that the operator satisfies $`T_{\mathrm{}}[x]=T_+\mathrm{}[x]`$.
### II.1 Spatially continuous nonlinear diffusion equation in one dimension
With the operator $`𝒯_\lambda []`$, one can define a general diffusion scheme which contains heat-diffusion as a special case for $`\lambda =0`$. To this end, consider, without loss of generality, the general form of a diffusion equation for a quantity $`f(x,t)`$ (here referred to as “activity”)
$$\frac{f(x)}{t}=\frac{}{x}\left(D(x)\frac{f}{x}\right)$$
(6)
where $`D(x)0`$ is the diffusion coefficient. If $`D(x)`$ depends on $`x`$, then the last equation describes a nonlinear diffusion process, otherwise ordinary heat diffusion. Consequently, by applying the operator $`𝒯_\lambda []`$ on the gradients, the following *pseudo-diffusion* process is obtained (which reduces to heat diffusion for $`\lambda =0`$):
$$\frac{f(x)}{t}=\frac{}{x}\left(D(x)𝒯_\lambda \left[\frac{f}{x}\right]\right).$$
(7)
By defining $`z(x)f(x)/x`$ and differencing we obtain
$$\frac{f(x)}{t}=\frac{D(x)}{x}𝒯_\lambda [z]+D(x)\frac{𝒯_\lambda [z]}{z}\frac{z(x)}{x}.$$
(8)
If $`D(x)=D=\mathrm{const}.`$, the last equation reduces to
$$\frac{f(x)}{t}=D\frac{𝒯_\lambda [z]}{z}\frac{^2f(x)}{x^2}.$$
(9)
The last equation looks in fact like an ordinary diffusion equation if we consider the factor $`D𝒯_\lambda [z]/z`$ as an “effective diffusion coefficient”. But which effect has the derivative $`𝒯_\lambda [z]/z`$? In appendix B it is shown that it approximates a Heaviside (or step) function $`H`$ for $`lim_{\lambda +\mathrm{}}`$, that is
$$\underset{\lambda +\mathrm{}}{lim}\frac{𝒯_\lambda [z]}{z}H(z).$$
(10)
In analogy to the previous case it can be shown that
$$\underset{\lambda \mathrm{}}{lim}\frac{𝒯_\lambda [z]}{z}H(z).$$
(11)
For a given cell, $`\lambda `$ specifies the ratio between negative and positive influx into the cell from its neighbors. Consider the case $`\lambda \mathrm{}`$ for a cell at position $`x`$. If the activity of any adjacent cell is higher, then the gradient $`z(x)f(x)/x`$ will be positive and an influx of activity to cell $`x`$ takes place, because in equation 8 the pseudo-diffusion term $`z(x)/x`$ is multiplied by one as a consequence of equation 10. Equation 10 also implies that any negative gradient at $`x`$ will make the pseudo-diffusion term be multiplied by zero, and thus prevents an influx of negative activity into cell $`x`$. The essence of this mechanism is that activity at $`x`$ can only increase until any gradient has dissipated. As an alternative, this mechanism may be understood as an auto-adaptive diffusion constant which regulates its value according to the current gradient at $`x`$ (figure 2).
For $`\lambda \mathrm{}`$ the mechanism works just vice versa, and the activity for a cell at position $`x`$ may only decrease. The linear (or heat) diffusion equation is obtained for $`\lambda =0`$, where both a positive-valued and a negative-valued influx can enter the cell.
### II.2 Intermediate values of $`\lambda `$ for a two cell system
Intermediate values of $`\lambda `$ attenuate either negative ($`\lambda >0`$) or positive influx ($`\lambda <0`$). The amount of attenuation depends on $`\lambda `$. To illustrate, consider a simplified pseudo-diffusion system which consists only of two cells $`u(t)`$ and $`v(t)`$:
$`{\displaystyle \frac{u}{t}}`$ $`=`$ $`𝒯_\lambda [vu]`$ (12)
$`{\displaystyle \frac{v}{t}}`$ $`=`$ $`𝒯_\lambda [uv]`$
Furthermore, we define the following *surrogate* system
$`{\displaystyle \frac{a}{t}}`$ $`=`$ $`\gamma (ba)`$ (13)
$`{\displaystyle \frac{b}{t}}`$ $`=`$ $`ab`$
with a diffusion coefficient $`\gamma `$. Note that because diffusion coefficients are different for $`a(t)`$ and $`b(t)`$ (that is, $`\gamma `$ and $`1`$, respectively), the last equation implements a nonlinear diffusion system. Without loss of generality, we assume $`\lambda >0`$, and $`u_0v_0>0`$ at $`t=0`$. Furthermore, let both diffusion systems have the same initial conditions $`u_0=a_0`$ and $`b_0=v_0`$. With this configuration of parameters, the influx into cell $`u`$ is negative, and will be attenuated because of $`\lambda >0`$. Dependent on the precise value of $`\lambda `$, the steady-state values of $`u_{\mathrm{}}`$ and $`v_{\mathrm{}}`$ will be situated somewhere between $`(u_0+v_0)/2`$ for $`\lambda =0`$, or $`\mathrm{max}(u_0,v_0)`$ for $`\lambda +\mathrm{}`$. Now, to understand the behavior for $`0<\lambda <\mathrm{}`$, the diffusion coefficient $`\gamma `$ is (numerically) determined such that both diffusion systems (equations 12 and 13) have the same equilibrium state, that is $`u_{\mathrm{}}=a_{\mathrm{}}`$ and $`v_{\mathrm{}}=b_{\mathrm{}}`$ (and also $`u_{\mathrm{}}=v_{\mathrm{}}`$).
With the assumptions $`\lambda >0`$ and $`u_0v_0>0`$, it follows that $`\gamma <1`$, because in order to obtain the same steady-state values for both diffusion systems, the negative influx into cell $`a`$ needs to be attenuated. Figure 2*a* shows that in this case the effective diffusion coefficient $`\gamma `$ and $`\lambda `$ have a sigmoidal relationship. The sigmoid shifts to the left as a function of $`\mathrm{\Delta }u_0v_0`$ (or equivalently $`a_0b_0`$).
Figure 2*b* shows that steady-state values as a function of $`\lambda `$ also follow a sigmoidal relationship. Cell values at convergence smoothly pass from heat diffusion ($`u_{\mathrm{}}[u_0+v_0]/2`$ for $`\lambda =0`$) to implementing a maximum operation ($`u_{\mathrm{}}\mathrm{max}[u_0,v_0]`$ for $`\lambda =\mathrm{}`$). Analogous considerations hold for negative values of $`\lambda `$.
### II.3 Spatially discrete pseudo-diffusion equation in two dimensions
Based on a centered finite difference representation of the Laplacian operator, we define a parameterized diffusion operator acting on a function $`f(x,y)`$ as
$$\begin{array}{cc}\hfill 𝒦_\lambda f(x,y)=& 𝒯_\lambda \left[f(x+1,y)f(x,y)\right]\\ \hfill +& 𝒯_\lambda \left[f(x1,y)f(x,y)\right]\\ \hfill +& 𝒯_\lambda \left[f(x,y+1)f(x,y)\right]\\ \hfill +& 𝒯_\lambda \left[f(x,y1)f(x,y)\right]\end{array}$$
(14)
where a grid spacing of $`\mathrm{\Delta }x=\mathrm{\Delta }y=1`$ is assumed. We will make use of the following compact notation
$$𝒦_+\mathrm{}\underset{\lambda +\mathrm{}}{lim}𝒦_\lambda $$
(15)
and
$$𝒦_{\mathrm{}}\underset{\lambda \mathrm{}}{lim}𝒦_\lambda .$$
(16)
Note that $`𝒦_{\lambda =0}^2`$ from equation 3.
In order to formulate a spatially discrete pseudo-diffusion scheme, we consider a diffusion layer (i.e. a finite grid on which diffusion takes place) with an equal number $`N`$ of rows $`i`$ and columns $`j`$, that is $`1i,jN`$. We use a discrete-in-space and continuous-in-time notation, where $`f_{ij}=f(j,i)`$, $`f_{i,j+1}=f(j+1,i)`$ and so on rem (a). With the above definitions, heat diffusion is described as:
$$\frac{f_{ij}}{t}=D𝒦_0f_{ij}(t)$$
(17)
where $`D=const.`$ is the diffusion coefficient. The process is assumed to start at time $`t=t_0`$ with the initial condition $`f_{ij}(t_0)=s_{ij}`$. From now on we assume that the $`s_{ij}`$ represent an intensity or luminance distribution (i.e., “$`s`$ represents a gray level image).
Since diffusion takes place in a bounded domain (i.e. we have a finite number $`N\times N`$ of grid points), and we also use adiabatic boundary conditions (i.e. there is neither inward flux nor any flux outward over the domain boundary, i.e. $`f_{ij}/t=0`$ for $`(i,j)\{(i,0),(i,N),(0,j),(N,j)\}`$) rem (b), the total activity described by equation 17 does not depend on time (figure 3), that is
$$\underset{i,j}{\overset{N}{}}f_{ij}(t)=\mathrm{const}.$$
(18)
The last expression expresses that diffusion is conservative – activity is neither created nor destroyed. Although the 2-D heat diffusion equation cannot create new activity levels which have not already been present at time $`t_0`$ Koenderink (1984), it can create extrema in activity domains that have a dimension greater than one (cf. Lifshitz and Pizer (1990), p.532). A min-diffusion layer will eventually compute the minimum of all values and is defined as:
$$\frac{a_{ij}}{t}=D𝒦_{\mathrm{}}a_{ij}(t).$$
(19)
A max-diffusion layer will eventually compute the maximum of all values and is defined as:
$$\frac{b_{ij}}{t}=D𝒦_+\mathrm{}b_{ij}(t).$$
(20)
We assume equal initial conditions $`a_{ij}(t_0)=b_{ij}(t_0)=s_{ij}`$ for the min-diffusion layer and the max-diffusion layer at time $`t=t_0`$.
Whereas equation 17 preserves its total activity, the min-diffusion layer and max-diffusion layer, respectively, do not. The total activity of the min-diffusion layer decreases with time and converges to (figure 3)
$$\underset{t\mathrm{}}{lim}\underset{i,j}{\overset{N}{}}a_{ij}(t)=N^2\underset{i,j}{\mathrm{min}}\{s_{ij}\}.$$
(21)
The total activity of the max-diffusion layer increases with time and converges to (figure 3)
$$\underset{t\mathrm{}}{lim}\underset{i,j}{\overset{N}{}}b_{ij}(t)=N^2\underset{i,j}{\mathrm{max}}\{s_{ij}\}.$$
(22)
In other words, all cells $`a_{ij}`$ of the min-diffusion layer will finally contain the global minimum of the input $`s_{ij}`$
$$A:=\underset{i,j}{\mathrm{min}}\{s_{ij}\}=\underset{t\mathrm{}}{lim}a_{ij}(t)i,j$$
(23)
and all cells $`b_{ij}`$ of the max-diffusion layer will end up with the global maximum
$$B:=\underset{i,j}{\mathrm{max}}\{s_{ij}\}=\underset{t\mathrm{}}{lim}b_{ij}(t)i,j.$$
(24)
This can be explained as follows. A cell $`a_{ij}`$ of the min-diffusion layer may only decrease its activity from one time step to the next, until any activity gradient between $`a_{ij}`$ and its nearest neighbors has dissipated. As a consequence, $`a_{ij}`$ adopts the minimum activity value of the neighborhood, including itself. Because the last arguments apply to *all* cells $`a_{ij}`$, eventually all cells will adopt the minimum activity $`\mathrm{min}_{i,j}\{s_{ij}\}`$ at convergence. Convergence occurs if $`a_{ij}=a_{kl}(i,j),(k,l)`$ (i.e. when no more activity gradient exists). The dynamic of the process is illustrated by figure 5.
In an analogous way, in the diffusion process described by the max-diffusion layer, all cells $`b_{ij}`$ could only increase their activity, given the existence of any activity gradient. If any cell has a maximum activity value, then finally all cells will adopt this maximum, since only then all gradients have dissipated.
Hence, both nonlinear diffusion systems are non-conservative, because they do not fulfill requirements analogous to equation 18.
## III Dynamic normalization by next neighbor interactions
Equipped with the pseudo-diffusion operators defined in the last section, we are now ready to define the dynamic normalization network. The network normalizes a given input $`s_{ij}`$ with respect to numerical range, but without taking resort to any global memory for determining the minimum and maximum. Rather, the global minimum and maximum are computed in the min-diffusion layer and the max-diffusion layer, respectively, by only exchanging information between adjacent cells.
We start with the following linear scaling scheme, which is typically used for normalizing a fixed set of numbers (again, see introduction):
$$c_{ij}=\frac{s_{ij}a_{ij}}{b_{ij}a_{ij}}\mathrm{where}1i,jN.$$
(25)
Because of equation 23 and 24 the variable $`c_{ij}`$ will contain (after a sufficiently long time) a normalized representation of $`s_{ij}`$, that is
$$s_{ij}[A,B]c_{ij}[0,1]$$
(26)
($`A`$ and $`B`$ are the global minimum and maximum, respectively, of $`\{s_{ij}\}`$). To arrive at a fully dynamical system, we formally interpret equation 25 as the steady-state solution of
$$\frac{c_{ij}}{t}=b_{ij}(0c_{ij})a_{ij}(1c_{ij})+s_{ij}$$
(27)
which shall be called dynamic normalization. Notice that by using dynamic normalization we naturally avoid the singularity of equation 25 that occurs for $`b_{ij}=a_{ij}`$.
Figure 6 visualizes the state of equation 27 at different time steps. Initially, the dynamic normalization process is similar to high-pass filtering (figure 7), what can be explained as follows. Contrasts are abrupt changes in luminance. Consider a luminance change from dark to bright. Then, the dark side has a local minimum, and the bright side a local maximum, which propagates in the min-diffusion layer and the max-diffusion layer, respectively (figure 5). When the local minimum (maximum) has propagated to the position of the bright (dark) side of the step, then the bright (dark) side will be normalized to one (zero). As the dynamic continues to evolve, local maxima and minima propagate further, thereby “eating” (i.e., annihilating) other smaller local maxima and minima. In figure 6, this annihilation of local maxima and minima, respectively, is visible through a gradual filling-in of image structures from the boundaries. A normalized version of the original image is finally obtained when $`c_{\mathrm{ij}}/t=0`$. Depending on *(i)* how small the $`s_{\mathrm{ij}}`$ are, and *(ii)* the choice of integration step size $`\mathrm{\Delta }t`$, the steady-state of $`c_{\mathrm{ij}}`$ can be reached with delay compared to the steady-states of $`a_{\mathrm{ij}}`$ and $`b_{\mathrm{ij}}`$, respectively. This is now examined in more detail.
### III.1 Time to convergence for dynamic normalization
Figure 8 shows the relationship between the time to convergence and the numerical range of the input $`s_{ij}`$: the smaller the $`s_{ij}`$, the more iterations are necessary to accomplish the mapping expressed by equation 26. Mathematically, this can be seen as follows. Assume that a general solution of equation 27 has the form
$$c_{ij}(t)=C_0e^{t/\tau }+C_1$$
(28)
where $`C_0`$ and $`C_1`$ are constants which are defined by the initial conditions, and $`\tau `$ is a time constant. Plugging the last equation into equation 27 yields
$$C_0e^{t/\tau }\left[b_{ij}a_{ij}\frac{1}{\tau }\right]=C_1\left(b_{ij}a_{ij}\right)+s_{ij}a_{ij}$$
(29)
By identifying
$$\tau =\tau _{ij}(t)\frac{1}{b_{ij}(t)a_{ij}(t)}$$
(30)
we obtain
$$C_1=\frac{s_{ij}a_{ij}}{b_{ij}a_{ij}},$$
(31)
and from the initial condition $`c_{ij}(t=0)=0i,j`$ we furthermore get $`C_0=C_1`$, which finally gives the solution
$$c_{ij}(t)=\frac{s_{ij}a_{ij}}{b_{ij}a_{ij}}\left(1e^{t/\tau _{ij}}\right).$$
(32)
On grounds of the definition of $`\tau `$ (equation 30) we obtain two insights.
First, since the time constant $`\tau `$ of the dynamic normalization process is a function of both $`a_{ij}(t)`$ and $`b_{ij}(t)`$, it is not really a constant, but rather depends on time and space because of equation 19 and 20, respectively. However, $`\tau `$ can be approximated by recalling that $`a_{ij}(t)`$ and $`b_{ij}(t)`$ converge in time and space to the global minimum $`A`$ and global maximum $`B`$, respectively, of the input $`s_{ij}`$ (equation 23 and 24). Thus, $`\tau 1/(BA)`$. This leads to the second insight: the smaller are $`A`$ and $`B`$, the longer it takes dynamic normalization to converge to a steady state. Or, otherwise expressed, the smaller $`\tau `$ is, the faster the system converges.
Notice that when using the steady-state solution (equation 25) of dynamic normalization instead of the full dynamic process (equation 27), no dependency on input contrast is revealed, and the dependence on spatial frequency structure of the input is much weaker.
### III.2 Transient adaptation or dynamic compression
The dynamic normalization layer reveals distinct dynamic phases. In the initial phase, image contrasts are extracted. Contrast enhancement occurs in a subsequent phase. In the final phase, the activity distribution in the dynamic normalization layer is just a re-scaled version of the input. In a second phase between the initial and the final phase, one observes adaptation: image structures with substantially different light intensities in the input are mapped to a smaller range of activities in the dynamic normalization layer. This effect is the *dynamic range compression*. For its illustration an input image was subdivided into four quadrants (“contrast tiles”, figure 9). Each of the tiles has a different range of luminance values. Because the available tonal range for displaying the tiled image is too small to match the range of all tiles, some of the image details in the darker tiles are displayed in black. Nevertheless, a part of these details are rendered visible in the dynamic normalization layer at around 100 iterations (top row in figure 9), implying that cell activities in this layer have less dynamic range than in the input. The compression effect is quantized in figure 10, where each curve represent the mean activity and the maximum activity, respectively, of all cells of one the four tiles. The curves approach each other at around 100 iterations. Thus, the output of the dynamic normalization network can be encoded with a smaller than the original numerical range.
Figure 11 illustrates the mechanism which underlies dynamic compression. A necessary condition for dynamic compression to occur is that the global maximum propagates with finite speed in the max-diffusion layer, and that it is spatially separated from image structures that have less dynamic range ($`=`$ local maxima). When the global maxima has not yet propagated to the local maxima, then image structures are normalized by their “own” local maxima. Since normalization rescales all cell activities to the same target range (all image structures normalize to one), local normalization implies a reduction of the dynamic range. However, local maxima are annihilated as the global maximum propagates, and image structures are now getting normalized by the global maximum. Then, the entire dynamic range of the input image is recovered in the normalization layer, and dynamic range compression is abolished. The recovery of the original dynamic range can be seen when the entropy curves of figure 12(b) reduce to the entropy of the input image at $`1000`$ iterations (dashed horizontal line).
### III.3 Process entropy
Figure 12(b) shows entropy as a function of time computed over the dynamic normalization layer. The entropy reaches a maximum in the time window where dynamic compression occurs. Notice that this maximum in entropy exceeds the entropy of the input image (dashed horizontal line). Because entropy quantifies the degree of flatness of a histogram (or probability distribution), the observed entropy maximum implies that cell activities of the dynamic normalization layer are more homogeneously distributed across the histogram than luminance values of the input. Figure 13 shows how the distribution of activities evolve over time. Initially, cell activities in the dynamic normalization layer are small, and tend to cluster around a single spot in the histogram (the cropped ”hot spot” in the upper left corner of the histogram). Emanating from this “hot spot”, values start to occupy nearly the entire histogram. It is just then when an observer who is monitoring the output of the dynamic normalization network gathers the highest information about the input image.
In the consecutive part of the dynamic, the values are redistributed again in a way that they concentrate around four principal stripes. These stripes correspond to the four contrast tiles. Therefore, dynamic range compression is compatible with adaptation, since adaptation maximizes the transfer of information Wainwright (1999).
### III.4 Adaptation-by-entropy-maximisation
When computing the Shannon entropy of the output of the dynamic normalization network Shannon (1948), one observes an entropy maximum at the dynamic compression effect (figure 12 and 13). Hence, a straightforward algorithm for the adaptation of images is to stop the dynamic normalization process when an entropy maximum is reached (“one feedback loop”). To further enhance the dynamic compression effect, the output at the entropy maximum can be taken again as input to the dynamic normalization network, and once again we can let the dynamic normalization process continue until it reaches a maximum of entropy (“two feedback loops”). The entropy across $`10`$ feedback loops of the just described algorithm is illustrated in figure 14 with the curve designated by “process entropy”. Figure 15 shows the output images obtained for one, two, three and $`20`$ feedback loops. With increasing number of feedback loops, luminance information is suppressed, while contrasts are enhanced. At around $`20`$ loops, one obtains an image which seems to contain only contours, but iterating further enhances also noise and leaves one with an image without any recognizable structures. Figure 14 shows that entropy decreases with increasing number of feedback loops (each data point is the entropy of the output at the indicated number of feedback loops). For the “*tiled*” and the “*4th power*Peppers image, the entropy versus feedback loops has a maximum. Concluding, in terms of entropy, but also by visual inspection, a small number of feedback loops (one or two) seems optimal for the proposed adaptation algorithm. The algorithm should be understood as a “proof-of-concept” rather than a definite tool for image processing, because it occasionally develops artifacts. For example, future versions could address the suppression of the dark zone which emanates from the tiles of the “*tiled*Peppers image.
### III.5 Sensitivity of dynamic normalization for noise
One may argue that an adaptive mechanism designed in a way suggested by dynamic normalization is highly sensitive to noise, because it is based on the computation of minimum and maximum operations. To address this issue, we further distinguish between static noise (i.e. an offset added to the input $`s_{ij}`$ which does not vary with time), and dynamic noise (i.e. an offset added to each layer which varies with time). For the first case we presume the existence of a noise-free input pattern, to which static noise is added. A worst case scenario is on hand if a couple of cells $`s_{ij}`$ have high activities due to noise (“noisy cells”) which lead to an undesired increase of the true dynamic range of the input. If a read-out mechanism for the dynamic normalization layer had only the same dynamic range as the input, then the noise would obscure the relevant information of the input at convergence. Nevertheless, if there were only a few noisy cells in the input layer, then the dynamic compression effect could mitigate the worst case scenario to some extent.
To assess the robustness of dynamic normalization against temporally varying noise, numerical experiments were conducted with additive, normal-distributed noise (“Gaussian noise”), with zero mean and standard deviation $`\sigma `$. Apart, additional simulations were conducted with multiplicative, uniform noise (“white noise”).
#### III.5.1 Additive, normal-distributed noise
Temporally fluctuating normal-distributed noise $`\xi _{ij}(t)`$ was added to the equations 19, 20, 27, and the input $`s_{ij}`$, according to
$$x_{ij}x_{ij}+\sigma \xi _{ij}(t).$$
(33)
In the last equation, $`x_{ij}`$ stands for one of the variables $`a_{ij}`$, $`b_{ij}`$, $`c_{ij}`$, and $`s_{ij}`$, respectively, and “$``$” means that the left hand side is replaced by the right hand side. The noise level is specified by $`\sigma `$ (assuming zero mean), and $`\xi _{ij}(t)`$ is assumed to be not correlated across time and/or spatial positions. A luminance step (32 $`\times `$ 32 pixels) was used as input, with luminance value zero on the dark side (“black patch”, columns 1 to 16), and 1 on the bright side (“white patch”, columns 17 to 32). Thus, the mean activity of the noise free system should approach one at steady-state. We furthermore computed the Michelson contrast $``$ at each position $`(i,j)`$ according to
$$_{ij}=\frac{c_{ij\mathrm{white}}c_{ij\mathrm{black}}}{c_{ij\mathrm{white}}+c_{ij\mathrm{black}}}$$
(34)
where $`j\mathrm{black}`$ means $`1j16`$ and $`j\mathrm{white}`$ means $`17j32`$ (the row index runs over all positions $`1i32`$).
Figure 17(a) shows the temporal evolution of the mean activity of $`<c_{ij}>_{ij\mathrm{white}}`$ (i.e. averaged over white patch positions) for various noise levels $`\sigma `$. Sufficiently high noise levels significantly affect the convergence behavior of dynamic normalization - the response plateau which is seen in the noise-free case is no longer reached. Instead of the plateau, a maximum is approached, the amplitude of which decreases with increasing noise level. Figure 17(b) shows that a similar behavior is also seen for the averaged Michelson-contrast $`<_{ij}>`$: The contrast between the black and the white patch decreases with increasing noise level. This implies that image structures are obscured by noise.
How does noise take influence on dynamic range compression? Three answers exist to this question, and they depend on the noise level. For relatively small noise levels ($`\sigma <0.001`$), no dramatic effect on dynamic compression is observed. For intermediate noise levels ($`\sigma 0.001`$), dynamic compression is enhanced (bottom row in figure 9, and figure 16). Enhancement happens because the net effect of noise is to add an offset, which “lifts” the darker patches of the tiled Peppers image. For higher noise levels, however, the darker patches drown in noise and image details get lost. Consequently, if the goal of dynamic normalization was adaptation, then suitable chosen noise levels would aid to enhance range compression, although this comes at the prize of reduced contrasts in regions with low activities (dark quadrants in the tiled Peppers image).
Notice that additive Gaussian noise can be easily counteracted by proposing additional mechanisms with low-pass characteristics, like spatial or temporal pooling of activity. Then, as long as the noise is not correlated over positions and time, it would simply average out.
#### III.5.2 Multiplicative and normally distributed noise
Multiplicative noise was applied to variables $`a_{ij}`$, $`b_{ij}`$, $`c_{ij}`$, and $`s_{ij}`$, respectively, according to
$$x_{ij}x_{ij}(1+\eta (\mu _{ij}(t)1))$$
(35)
with $`0\mu _{ij}(t)1`$ representing uniformly distributed noise which was uncorrelated across time and/or space. The noise level is specified by $`\eta [0,1]`$. Dynamic normalization is not significantly affected by this type of noise, not even for $`\eta =1`$ (hence results are not shown). Multiplicative noise acts differently on maxima and minima. Maximum activities can only decrease, but never increase beyond their value in the noise free case. Therefore, no spurious maxima are introduced into the max-diffusion layer by the type of multiplicative noise considered here. On the other hand, multiplicative noise can inject spurious minima into the min-diffusion layer, if the lowest luminance value in the input image was bigger than zero. As the minimum luminance values of our images were always zero, they are consequently not affected by the multiplicative noise.
## IV Discussion
### IV.1 Pseudo-diffusion and electrical synapses (gap junctions)
The operator $`𝒦_\lambda `$ models different types of electrical synapses (gap junctions). In its linear version, $`𝒦_{\lambda =0}`$ describes the exchange of both depolarizing (i.e. directed towards a neuron’s firing threshold) and hyperpolarizing (i.e. directed away from a neuron’s firing threshold) currents between adjacent neurons. Networks of electrically coupled neurons are ubiquitous both in the retina (e.g. Raviola and Gilula (1975); Kolb (1977); Nelson et al. (1985); Mills and Massey (1994, 2000)) and the cortex (e.g. Gibson et al. (1999); Galarreta and Hestrin (1999, 2002)). These networks can be modeled by diffusion equations (e.g. Naka and Rushton (1967); Lamb (1976); Winfree (1995); Benda et al. (2001)). Conversely, the operators defined by equation 15 and 16 represent models for rectifying (i.e. voltage sensitive) gap-junctions. Rectifying gap junctions were described in the crayfish (e.g. Edwards et al. (1991, 1998)), and unidirectional and gated gap junctions were reported in the rat (e.g. Bukauskas et al. (2002)) and turtle (e.g. Piccolino et al. (1984); Bukauskas et al. (2002)), respectively.
In organisms, rectifying gap junctions may nevertheless be implemented in a “dirty” fashion. This means that a current flux may not strictly occur in only one direction. Rather, a small amount of current may as well flow in the opposite direction. Such behavior is captured by setting $`\lambda `$ to a finite value $`1|\lambda |<\mathrm{}`$, and was analyzed in figure 2.
### IV.2 Computational aspects
Substitution of two global memories (for the minimum and the maximum activity) by two pseudo-diffusion layers of size $`N\times N`$ leads to a computationally more demanding system, because more memory resources are needed and significantly more computational operations need to be carried out for their simulation. Moreover, because computation of the global maximum or minimum is based on local, diffusion-like interactions, a maximum or a minimum does not propagate from one cell to another from one time step to the next. The diffusion rate cannot be chosen arbitrarily high to guarantee the numerical stability of the process. The time to convergence does not only depend on the pseudo-diffusion layers reaching a steady-state, but is mainly determined by the dynamic normalization layer. The number of iterations that is needed until convergence occurs scales with the numerical range of the input. Thus, for small input values, the number of required iterations can be quite large (see figure 8). Therefore, the dynamic normalization network cannot be seriously considered as an alternative to an ordinary normalization algorithm (i.e., searching the global maximum and minimum, and then rescaling). However, the dynamic normalization network can accomplish different tasks which cannot be accomplished with an ordinary normalization algorithm, for example detection of contrast contours, or compression of the dynamic range of the input.
## V Conclusions
This paper introduced a parameterized diffusion operator (parameter $`\lambda `$) and analyzed some of its properties mathematically and by computer simulations. As a special case, heat diffusion is obtained for $`\lambda =0`$. Diffusion layers which are based on the two limit cases of the operator (for $`\lambda \pm \mathrm{}`$) compute the global maximum and minimum, respectively, of the initial cell activities of the layer. This means that at convergence, all cells of the diffusion layers contain the same activity value – the maximum ($`\lambda \mathrm{}`$) or the minimum ($`\lambda \mathrm{}`$). Based on these operators, a dynamic normalization network was defined (equation 27). Its steady-state solution is functionally equivalent to the ordinary rescaling of a set of numbers (equation 25), but by making the normalization process dynamic, one observes two additional properties: contrast enhancement and dynamic range compression. Both effects occur because at first normalization acts locally, similar to adaptation mechanisms. With increasing time, the normalization process gets continuously more global, until a steady-state is reached. The steady-state corresponds to a rescaling of the input in the dynamic normalization layer.
By exploiting the dynamic compression effect, it should thus be possible to design a powerful adaptation mechanism which maps an input image of an arbitrary numerical range to a smaller target range. To do so, the normalization process has to be “frozen” when dynamic compression occurs. As a first step into that direction, a simple adaptation algorithm based on the maximisation of entropy was proposed (section III.4): the dynamic is frozen as soon as a maximum of entropy is reached, and the output is then fed back as new input to the dynamic normalization network. As a further improvement, the diffusion operators could be modified such that activity exchange between two cells is blocked for sufficiently large activity gradients Keil et al. (2005). Doing so would possibly prevent in figure 9 (first and second row) the global maximum from spreading between tiles, and would normalize each tile independently, such that ideally a dynamic similar to the bottom row in figure 9 is produced.
Systems based on pseudo-diffusion have already turned out to be of utility for a variety of purposes in image processing (for implementing filling-in mechanisms, or winner-takes all inhibition, see Keil et al. (2005)). Pseudo-diffusion systems can generally be used for implementing the max-operation without the need for globally acting pooling units (see for example Yu et al. (2002) and Keil (2003)). The advantage over functionally equivalent but hardwired systems is that the region where normalization takes place can be dynamically adjusted. Furthermore, the maximum operation serves to implement invariance properties in models for object recognition (e.g. Riesenhuber and Poggio (1999, 2000)).
###### Acknowledgements.
This work was supported by the *Juan de la Cierva* program of the Spanish government, and the the MCyT grant SEJ 2006-15095. Further support was provided by the AMOVIP INCO-DC 961646 grant from the European Community. The author likes to thank the anonymous reviewer of *Physica D* for his valuable suggestions which helped to improve the manuscript (the present version is the long version; a shorter version is to be published in *Physica D*).
## Appendix A Material and methods
All simulations were carried out using the Matlab environment (R2006b) on a Linux workstation, where both native Matlab code and mex-files programmed in C++ were used. Diffusion operators were normalized by the number of adjacent cells (normally four, along the domain boundaries three, and in the corners two). Normally, the equations describing the diffusion layers (eqs. 19 and 20), and dynamic normalization (eq. 27) turned out to be numerically stable such that a forward-time-centered-space (FTCS) Euler scheme with step size one is sufficient. (Here, we understand numerical stability such that the solution converges rather than growing in an unbounded fashion). Notice, however, the stability criterion associated with the FTCS-integration of the heat diffusion equation $`D\mathrm{\Delta }t1/2`$ (assuming grid spacing one, see section 19.2 in Press et al. (1997)) where $`D`$ is the diffusion coefficient, and $`\mathrm{\Delta }t`$ is the integration step size. Since we compared pseudo-diffusion with Laplacian or heat diffusion (eq. 17), by default we employed Euler’s method with integration step size $`\mathrm{\Delta }t=0.5`$ and diffusion coefficient $`D=1`$. Exceptions are as follows. Figure 2 was simulated with $`\mathrm{\Delta }t=0.1`$ and $`\mathrm{\Delta }t=0.001`$, respectively. Figure 6, 8, and figures 9 to 16 were integrated with the fourth-order Runge-Kutta method ($`\mathrm{\Delta }t=0.5`$, $`D=1`$). For the compilation of figure 17, again the forth-order Runge-Kutta method was used with $`\mathrm{\Delta }t=0.01`$ and $`D=1/\mathrm{\Delta }t`$, to guarantee numerical stability in the presence of high noise levels.
It should be emphasized that the results presented in this paper do not depend critically on the exact value of neither $`\mathrm{\Delta }t`$ and $`D`$, nor on the specific choice of the integration method. Variation of these parameters leads to a corresponding rescaling of the time axis. Although we exemplified the behavior of the model only by means of two standard images which are commonly used for image processing (Lena and Peppers ), all characteristics of the model can as well be reproduced with other images.
## Appendix B Proof of equation 10 (for $`\lambda \mathrm{}`$ and $`\lambda =0`$)
Consider the derivative of the operator $`𝒯_\lambda []`$ (equation 1),
$$\frac{𝒯_\lambda [z]}{z}=\underset{\mathrm{𝑡𝑒𝑟𝑚}I}{\underset{}{\frac{\eta }{1+e^{\lambda z}}}}+\underset{\mathrm{𝑡𝑒𝑟𝑚}\mathrm{𝐼𝐼}}{\underset{}{\frac{\eta \lambda ze^{\lambda z}}{\left(1+e^{\lambda z}\right)^2}}}$$
(36)
where the following three cases have to be analyzed:
In this case $`\eta =2`$ from equation 2, and
$$\frac{𝒯_0[z]}{z}=\underset{\mathrm{𝑡𝑒𝑟𝑚}I}{\underset{}{\frac{\eta }{1+1}}}+\underset{\mathrm{𝑡𝑒𝑟𝑚}\mathrm{𝐼𝐼}}{\underset{}{0}}=1.$$
(37)
Thus, for $`\lambda =0`$ the derivative is constant one for all $`z`$, and equation 7 reduces to the linear diffusion equation 6.
In this case $`\eta =1`$ from equation 2, and we have to consider three additional cases according to the value of $`z`$. Note that $`z`$ is treated as a constant. Hence,
$$\underset{\lambda +\mathrm{}}{lim}\frac{𝒯_\lambda [z]}{z}=\underset{\mathrm{𝑡𝑒𝑟𝑚}I}{\underset{}{\frac{\eta }{1+1}}}+\underset{\mathrm{𝑡𝑒𝑟𝑚}\mathrm{𝐼𝐼}}{\underset{}{0}}=\frac{1}{2}.$$
(38)
This is to say that if the gradient $`z`$ vanishes, then the derivative is constant with value $`1/2`$.
We start with evaluating term I of equation 36,
$$\underset{\lambda +\mathrm{}}{lim}\frac{\eta }{1+\underset{0}{\underset{}{e^{\lambda |z|}}}}=\eta =1(\mathrm{𝑡𝑒𝑟𝑚}I).$$
(39)
In the numerator of term II appears a product of the kind “$`\mathrm{}0`$”. One may argue that the exponential $`\mathrm{exp}(\lambda |z|)`$ always approaches zero much more faster than the term $`\eta \lambda |z|`$ is able grow (or one may equivalently apply l’Hospital’s rule to this product by applying $`d/d\lambda `$ on each factor),
$$\underset{\lambda +\mathrm{}}{lim}\frac{(\eta \lambda |z|)(e^{\lambda |z|})}{\left(1+e^{\lambda |z|}\right)^2}\underset{\lambda +\mathrm{}}{lim}\frac{\stackrel{\mathrm{const}.}{\stackrel{}{(\eta |z|)}}\stackrel{0}{\stackrel{}{(|z|e^{\lambda |z|})}}}{(1+\underset{0}{\underset{}{e^{\lambda |z|}}})^2}=0(\mathrm{𝑡𝑒𝑟𝑚}\mathrm{𝐼𝐼}).$$
(40)
Thus, for $`lim_{\lambda +\mathrm{}}`$ evaluates equation 36 to 1 for all $`z>0`$.
Evaluating term I,
$$\underset{\lambda +\mathrm{}}{lim}\frac{\eta }{1+\underset{\mathrm{}}{\underset{}{e^{\lambda |z|}}}}=0(\mathrm{𝑡𝑒𝑟𝑚}I).$$
(41)
Evaluating term II (again there is a little more work to do),
$$\underset{\lambda +\mathrm{}}{lim}\frac{\eta \lambda |z|}{2+e^{\lambda |z|}+e^{\lambda |z|}}\underset{\lambda +\mathrm{}}{lim}\frac{\eta }{\underset{0}{\underset{}{e^{\lambda |z|}}}\underset{\mathrm{}}{\underset{}{e^{\lambda |z|}}}}=0(\mathrm{𝑡𝑒𝑟𝑚}\mathrm{𝐼𝐼}).$$
(42)
Hence, for $`lim_{\lambda +\mathrm{}}`$ evaluates equation 36 to 0 for all $`z<0`$.
$`\mathrm{}`$
Summarizing the above we saw that equation 36 behaves approximately rem (c) like a Heaviside function $`H`$ for $`lim_{\lambda +\mathrm{}}`$, thus equation 10 is proofed. The proof of equation 11 (for $`\lambda \mathrm{}`$) proceeds in straight analogy.
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# Dust outflows and inner gaps generated by massive planets in debris disks
## 1 Introduction
Debris disks are disks of dust that surround many main sequence stars. They were discovered by the IRAS satellite in the 1980’s (Aumann et al. 1984; Gillett 1986) and they are preferentially detectable at infrared wavelengths, where the dust re-radiates the light absorbed from the star. Stars harboring debris disks are too old to have remnants of the primordial disk from which the star itself once formed. This is because the dust grain removal processes, such as the Poynting-Robertson (P-R) effect and solar wind drag, act on timescales much shorter than the age of the star, indicating that such “infra-red excess stars” harbor a reservoir of undetected planetesimals producing dust by mutual collisions or by evaporation of comets scattered close to the star (Backman & Paresce 1993). The spectroscopy of systems like $`\beta `$-Pictoris supports this interpretation (e.g. Knacke et al. 1993; Pantin, Lagage & Artymowicz 1997). It seems clear, therefore, that planetesimals are present in these systems. But what about massive planets? High-resolution images of some of these debris disks have revealed the presence of density structure (see Koerner 2001 for a review) and dynamical models have shown that planets can sculpt the dust disks, creating gaps, arcs, rings, warps and clumps of dust (e.g. Roques et al. 1994; Liou & Zook 1999; Mouillet et al. 1997; Wyatt et al. 1999; Moro-Martín & Malhotra 2002; Kuchner & Holman 2003). The combination of both, the very high resolution imaging at long wavelengths and theoretical dynamical models can provide interpretation of the disks’ structure in terms of planetary architectures. This approach has been used in the interpretation of high resolution millimeter interferometry observations of the Vega system (Wilner et al. 2002) and of the submillimeter images of the $`ϵ`$ Eridani system (Ozernoy et al. 2000; Quillen & Thorndike 2002). Recent observations with the $`\mathrm{𝑆𝑝𝑖𝑡𝑧𝑒𝑟}`$ MIPS instrument have confirmed that out of 26 FGK field stars known to have planets by radial velocity studies, 6 show 70$`\mu `$m excess at 3-$`\sigma `$ confidence level, implying the presence of cool material ($`<`$100 K) located beyond 10 AU (Beichman et al. 2005). These stars, with a median age of 4 Gyr, are the first to be identified as having both well-confirmed planetary systems and well-confirmed IR excesses (Beichman et al. 2005). In addition, the first results from the $`\mathrm{𝑆𝑝𝑖𝑡𝑧𝑒𝑟}`$ $`\mathrm{𝐹𝐸𝑃𝑆}`$ Legacy project indicate that inner gaps<sup>1</sup><sup>1</sup>1In this paper a “gap” is an inner depletion zone in the dust disk interior to the planet’s orbit, not an annular depletion zone around the planet’s orbit. appear to be common in cold Kuiper Belt-like disks (Kim et al. 2005). These disks show excesses at 70 $`\mu `$m but not at 24 $`\mu `$m, indicating again the presence of cool dust ($`<`$100 K) located beyond 10 AU. The lack of 24 $`\mu `$m emission yields an upper limit to the amount of warm dust inside 10 AU; this upper limit is 10<sup>-3</sup> to 10<sup>-2</sup> times the lower limit for the mass in the corresponding cold disk. Because the lifetime of the dust particles due to P-R drag is of the order of 1 Myr, it is expected that the density contrast would be erased on this timescale. Kim et al. (2005) suggest that a possible explanation for these inner gaps is that one or more massive planets are dynamically depleting, via gravitational scattering, dust particles generated by an outer belt of planetesimals. All these observations are providing increasing evidence that debris disks and massive planets co-exist around many sun-like stars.
In this paper we report some new results based on numerical modeling regarding the depletion of large dust particles in debris disks by the gravitational perturbations of massive planets. The numerical models used to carry out this study are briefly described in $`\mathrm{\S }`$2. The ejected particles form an “outflow” whose properties (angular confinement, velocity and efficiency of ejection) are characterized in $`\mathrm{\S }`$3.1 as a function of the planet’s mass and orbital elements, and the particle size. The high efficiency of ejection, together with the possible high frequency of debris disks harboring massive planets, suggest that these outflows may be a common phenomenon, whose implications are described in $`\mathrm{\S }`$3.2. The ejection of particles is also responsible for the depletion of dust interior to the orbit of the planet, creating a density contrast that can be measured directly in spatially resolved images or indirectly through the modeling of the spectral energy distribution (SED) of the debris disk. To aid in the interpretation of such observations, in $`\mathrm{\S }`$4 we study the density contrast inside and outside the orbit of the planet, as a function of the planet’s mass and orbital elements and the particle size. Finally, $`\mathrm{\S }`$5 summarizes our results.
## 2 The numerical models
We numerically solve the equations of motion of dust particles generated in a debris disk, analogous to the solar system’s Kuiper Belt. We use a modified version called SKEEL of the multiple time step symplectic method SyMBA (Duncan, Levison & Lee 1998; Moro-Martín & Malhotra 2002). Our models include the combined effects of solar gravity, solar radiation pressure, the P-R effect and solar wind drag, and the gravitational forces of planets. We model the solar system with 7 planets (excluding Mercury and Pluto, and including the mutual perturbations of the planets), and we model hypothetical extra-solar planetary systems with single planets of different masses, semimajor axes and eccentricities (see Tables 1, 2 and 3 for a complete list of models). For some of these systems, the parent bodies of the dust particles are assumed to be distributed in orbits with semimajor axis between 35 and 50 AU, eccentricities such that the perihelion distances are between 35 and 50 AU, and inclinations between 0 and 17, in approximate accord with current estimates of the orbital distribution of the classical Kuiper Belt (Malhotra et al. 2000; Brown 2001). For other systems, the dust-producing planetesimals are randomly distributed in a thinner disk with $`a`$=35–50 AU, $`e`$=0–0.05 and $`i`$=0–0.05 radians. In all our models, the initial values of mean anomaly (M), longitude of ascending node ($`\mathrm{\Omega }`$) and argument of perihelion ($`\omega `$) were randomly distributed between 0 and 2$`\pi `$. We run models for different particle sizes, referred to in terms of their $`\beta `$ value, which is the dimensionless ratio of the radiation pressure force and the gravitational force. For spherical grains,
$$\beta =(3L_{}/16\pi \mathrm{𝐺𝑀}_{}c)(Q_{\mathrm{𝑝𝑟}}/\rho s),$$
(1)
where L and M are the stellar luminosity and mass; for a solar-type star, $`\beta `$=5.7 $`\times `$ 10<sup>-5</sup> Q<sub>pr</sub>/$`\rho `$s, where $`\rho `$ and s are the density and radius of the grain in cgs units (Burns, Lamy & Soter 1979). Q<sub>pr</sub> is the radiation pressure coefficient, a function of the physical properties of the grain and the wavelength of the incoming radiation; the value we use is an average, integrated over the solar spectrum. \[For the correspondence between $`\beta `$ and the particle size see Fig. 5 in Moro-Martín, Wolf & Malhotra (2005).\] The sinks of dust included in our numerical simulations are (1) ejection into unbound orbits, (2) accretion into the planets, and (3) orbital decay to less than 0.5 AU heliocentric distance (0.1 AU for the models with a single planet located at 1 AU). A detailed description of the numerical algorithm used to integrate the equations of motion is given in Moro-Martín & Malhotra (2002).
## 3 Dust outflows from debris disks
Radiation pressure arises from the interception by the dust particles of the momentum carried by the incident stellar photons; it makes the orbits of the dust particles change immediately upon release from their parent bodies (i.e., the meter-to-kilometer size dust-producing planetesimals). For parent bodies in circular orbits, small grains with $`\beta >`$ 0.5 are forced into hyperbolic orbits as soon as they are released. If the parent bodies’ orbits are eccentric, ejection occurs for $`\beta >`$ 0.5(1$``$e) for a particle released at perihelion or aphelion, respectively. In the solar system these particles are known as $`\beta `$-meteoroids (Zook & Berg 1975). These small dust particles leave the system in a “disk wind”, whose angular extent is determined by the inclinations of the parent bodies; this is because radiation pressure is a radial force which does not change the inclinations of the dust particles after their release.
Grains larger than the “blow-out” size, on the other hand, remain on bound orbits upon release, and their orbital evolution is the subject of our study. Their dynamical evolution is affected by the P-R effect, which tends to circularize and shrink their orbits, forcing these particles to slowly drift in toward the central star (Burns, Lamy & Soter 1979). If no planets were present, the final fate of these dust particles would be to drift all the way into the star until they sublimate. Other removal processes may include mutual grain collisions and collisions with interstellar grains, which may comminute the grains to sizes small enough to be blown away by radiation pressure. \[The studies reported here do not include collisional effects; for an estimate of the limitations of our models we refer to Moro-Martín & Malhotra (2002 and 2003).\] When planets are present the story changes: (a) the trapping of particles in mean motion resonances (MMRs) with the planets causes an accumulation of particles at resonant semimajor axes; and (b) sufficiently massive planets can scatter and eject dust particles out of the planetary system. In the case of dust produced in the Kuiper Belt in our solar system, about 80–90% of the dust grains are ejected by close encounters with the giant planets (mainly Jupiter and Saturn), a few percent accrete onto the planets, and the remaining 10–20% drift all the way into the Sun (Liou, Zook & Dermott 1996; Moro-Martín & Malhotra 2003; see also Table 1). Thus, in addition to the afore-mentioned $`\beta `$-meteoroids, an outflow of larger particles produced by gravitational scattering from planets also exists.
### 3.1 Dependence on planetary architecture and particle size
We have explored the characteristics of the large particle outflow and its dependence on planetary architecture and particle size. For the solar system architecture, it is known that the majority of KB dust particles are ejected by Jupiter and Saturn (Liou, Zook & Dermott 1996; Moro-Martín & Malhotra 2003). Motivated by this, we have modeled hypothetical planetary systems consisting of a single planet and a KB-like dust source. These models explore a range of planetary masses (M<sub>p</sub>/M<sub>Jup</sub>=0.03, 0.1, 0.3, 1, 3, and 10), orbital semimajor axis ($`a`$=1, 5.2, 10, 20 and 30 AU), and eccentricities ($`e`$=0, 0.1, 0.2, 0.3, 0.4, 0.5) (see Tables 1, 2 and 3).
Fig. 1 shows examples of the escaping<sup>2</sup><sup>2</sup>2Our definition of “escaping” is that the particles reach a distance 1000 AU from the star (see Fig. 3 to 6); at that point, we stop integrating their orbits. This is not quite equivalent to the precise criterion for ejection, which would be that a particle velocity exceed the escape velocity. However, our numerical studies find that the particles that reach 1000 AU, 30-60% (depending on their $`\beta `$) are in hyperbolic orbits, and more than 90% have orbital eccentricity $`e>`$0.98. This means that even though some of the particles are still bound by the time they reach 1000 AU, it is very likely that they will also be set on hyperbolic orbits within a few orbits, either by subsequent scattering from the planets or due to small additional perturbations not included in our models. particle trajectories for the solar system case, projected in the ecliptic plane (XY; left panel), and in the RZ plane (right panel; where R is the in-plane heliocentric distance and Z is the off-plane out-of-ecliptic distance). These examples are of particles that reach at least 1000 AU and had their last encounter with Jupiter. We see that Jupiter creates a fan-like outflow, mainly confined to the ecliptic, where the trajectories are in the counterclockwise (prograde) direction. The distributions of eccentricity, inclination and perihelion of these Jupiter-ejected particles are presented in Fig. 2. The histograms show that all the particles are either in or very close to hyperbolic orbits; that the scattering rarely changes the inclination of the particles by more than 15 degrees (see also column 9 in Table 1); and that few of the ejected particles leave on orbits of perihelion interior to Jupiter’s orbit.
For the single planet models, Fig. 3–6 show the velocities of the escaping particles at 1000 AU projected in the XY (ecliptic) plane (left) and in the XZ plane (right). At large heliocentric distances the outflow is radial and symmetric, except when the planet is in an eccentric orbit (Fig. 6); the projection in the XZ plane shows that it is largely confined to the ecliptic for Jupiter-mass planets (or smaller), and becomes less confined as the planet mass increases. The angular confinement to the disk can also be seen in Fig. 7 and 8, in the distribution of orbital inclination for the ejected particles, and in column 9 of Tables 1, 2 and 3. This angular confinement is not obvious a priori because the ejection of the particles is due to gravitational scattering, a process that does not necessarily preserve the inclination of the orbits.
In Tables 1–3, we give a list of the single-planet and multiple-planet models that we have simulated (a total of 126 models). Also included in these tables are the statistical results for the fates of the dust particles in each model. We have performed simulations for several solar system models with the same or similar initial conditions of the dust parent bodies and the results indicate that a conservative estimate of the uncertainty in n<sub>1000</sub>, owing to the chaotic dynamics of dust orbital evolution, is $``$10% of the initial number of particles.
Fig. 9 and 10 show the percentage of particles that are gravitationally scattered out from the system, and the velocity at infinity of the ejected particles, as a function of the planet’s mass, semimajor axis and eccentricity, and the particle size. We find the following dependencies (the parentheses show the values explored by our models).
* Particle sizes ($`\beta `$=0.00156, 0.00312, 0.00625, 0.0125, 0.025, 0.044, 0.1, 0.2 and 0.4): It is expected that gravitational scattering is dependent to some extent on the particle size as smaller particles (larger $`\beta `$) migrate past the planet faster, therefore decreasing their probability of ejection. The top panel of Fig. 9 shows that: (1) For a 1M<sub>Jup</sub> planet, the efficiency of ejection decreases as $`\beta `$ increases, reaching a minimum at $`\beta `$$``$0.1–0.2 and increasing thereafter. As mentioned above, the decrease in efficiency is expected because the particle P-R drift velocity is larger for larger $`\beta `$. The increase in efficiency for even larger $`\beta `$ is probably due to the fact that radiation pressure is starting to contribute to the ejection of the particles. (2) The effect described above is more significant for close-in planets (1 AU), i.e. when the particle is deeper in the potential well of the star. (3) Planets $`>`$3M<sub>Jup</sub> in circular orbits between 1 AU and 30 AU eject $`>`$80% of the particles that go past, independently of the particle size. In addition, from the top panel of Fig. 10, we see that there is an increase in $`\overline{v}_{\mathrm{}}`$ as the particle size decreases ($`\beta `$ increases), which is more pronounced when the perturbing planet is closer to the star. The distributions of particle inclinations in Fig. 2 and 7 show that the angular confinement of the ejected particles is similar for all particle sizes. This is not surprising because the inclination perturbation in gravitational scattering is independent of particle size, as particle masses are more than 30 orders of magnitude smaller than the masses of the planets.
* Planet semimajor axis (1, 5.2, 10, 20 and 30 AU): We find that the average dust outflow velocity is larger in the presence of close-in planets than more distant planets of the same mass (see top panels of Fig. 10). This trend is clearly seen in the left panel of Fig. 11; the slope of the line corresponds to approximately $`\overline{v}_{\mathrm{}}`$$``$$`a`$$`{}_{}{}^{0.5}{}_{pl}{}^{}`$, and is consistent with an analytical calculation by Murray, Weingartner & Capobianco (2003).
* Planet mass (0.03, 0.1, 0.3, 1, 3, and 10M<sub>Jup</sub>): The right panel of Fig. 11 shows only a weak dependence on the mass of the planet of the average particle ejection velocity; this is somewhat in contrast with the theoretical prediction, $`v_{\mathrm{}}`$M$`{}_{}{}^{1/4}{}_{pl}{}^{}`$ (Murray, Weingartner & Capobianco 2003). The magnitude of the ejection velocity, $``$3 km $`\mathrm{s}^1`$, in the Jupiter-mass single-planet models (blue line in the top left panel of Fig. 10) is higher than the numerical result in Murray, Weingartner & Capobianco (2003), but agrees better with their analytical estimate. Their analysis, however, assumes that the particle ejection takes place after only a single planetary encounter, whereas our simulations show that typically ejections occur after many planetary encounters. (In our simulations, we track the planetary encounters of dust particles within 3.5 Hill-radius distance from each planet. The number of such encounters that ejected particles suffer is on the order of 10–10<sup>4</sup>, with the lower range being more typical in models with more massive planets, 3–10 M<sub>Jup</sub>). In addition to this complexity, it is important to remember that the effect of the planet’s orbital elements and mass on the outflow parameters (velocity and confinement to the plane) is not only direct, via the close encounters, but also indirect, as the particles encounter the planet with a history of evolution in the MMRs that can change the initial orbital elements of the particles and therefore affect their subsequent dynamical evolution. As an example, the eccentricity distributions of the soon-to-be-ejected particles near the planet show that for the 1 and 3 M<sub>Jup</sub> models, $`e`$0.4–0.5, but for 10 M<sub>Jup</sub>, $`e<`$0.2.
The distribution of inclinations in Fig. 7 shows that for a planet at 1 and 5.2 AU, the angular confinement of the outflow to the disk is affected by the planet’s mass; the more massive the planet the less confinement the outflow has. However, the parameter that is most strongly dependent on the planet’s mass is the number of ejected particles. The bottom left panel of Fig. 9 shows that there is a sharp increase in ejection efficiency when the planet mass increases from 0.3 M<sub>Jup</sub> to 1 M<sub>Jup</sub>: planets $``$0.1 M<sub>Jup</sub> do not eject a significant number of particles, whereas planets $`>`$3 M<sub>Jup</sub> eject $`>`$90% if located between 1–30 AU. A 1 M<sub>Jup</sub> planet at 5–30 AU ejects about 80% of the particles, and about 60% if located at 1 AU.
* Planet eccentricity (0, 0.1, 0.2, 0.3, 0.4 and 0.5): Large planet eccentricities create an asymmetric outflow oriented along the major axis of the planet’s orbit. The number of particles ejected in the apoastron direction exceeds that in the periastron direction by a factor of $``$5 for e=0.5 (see Fig. 6). The asymmetry is due to the fact that the planet spends more time near apoastron and therefore the probability of encounter with a dust particle is higher near apoastron. Fig. 8 and 10 show that the inclinations and the average velocity of the ejected particles at infinity are not affected by the planet’s eccentricity. The efficiency of ejection, however, decreases significantly as the planet’s eccentricity increases: for a 1 M<sub>Jup</sub> planet at 5 AU it decreases from $``$80% to $``$30% when the planet eccentricity is increased from 0 to 0.5 (see bottom right panel in Fig. 9). It is of interest to note that many of the known exo-planets to date have large orbital eccentricities (Marcy et al. 2003); our models predict that the large particle outflow will be asymmetric in these cases.
* Comparison with Solar System: The single-planet analog of the solar system (i.e. only Jupiter in a circular orbit at 5.2 AU) produces a somewhat higher velocity outflow compared with the actual multi-planet solar system. This is mainly due to the effect of Saturn in our solar system: having a larger semimajor axis, Saturn intercepts a fraction of the KB dust grains as they evolve inward due to the P-R drag and ejects them at a somewhat lower velocity, thus depressing the mean velocity of the outflow.
### 3.2 Implications of dust particle outflows
There are several significant implications of this large-particle outflow.
#### 3.2.1 Exo-planetary debris disks and planet formation environment
Stellar surveys show that at least 15% of A-K main sequence stars are surrounded by debris disks, and that the far-infrared excess decreases with stellar age, dropping from about 50% to about 15% after approximately 500 Myr. But these samples are sensitivity-limited, and therefore the occurrence of debris disks could be higher (Lagrange, Backman & Artymowicz 2000 and references therein). Stellar radial velocity surveys indicate that about 7% of the FGK main sequence stars have a Saturn or Jupiter-mass planet within 3 AU (Marcy 2003). Even though the correlation between the presence of planets and debris disks is not known yet, our studies suggest these large-particle dust outflows may be a common phenomena in planetary systems that harbor debris disks. This is of interest because:
(a) These large-particle dust outflows may contribute significantly or even dominate the clearing of circumstellar debris in planetary systems. Hitherto, the main processes that have been considered for such clearing are stellar winds, radiation pressure, sublimation, and collisions. The latter reduce the size of the dust particles until they are small enough to be blown away by radiation pressure. However, as our models indicate, gravitational scattering by giant planets following orbital decay by P-R drag is also significant, and in some cases may be a dominant process, ejecting 50–90% of the dust grain population.
(b) These outflows should be added to the list of processes that link the interplanetary environment to the galactic environment of a star. Planetary systems are prime sites for large particle formation. As such, they can contaminate the immediate vicinity of star-forming regions through this large particle outflow, and thus affect the particle size distribution of their local ISM. It is likely, therefore, that large particle outflows from extra-solar planetary systems may be a source of the large interstellar particles that have been detected in the interplanetary medium.
The presence of an outflow in an exo-planetary system and its detectability will strongly depend on the orbital characteristics of the planet and the orientation of the system. For face-on systems, the expected surface brightness of the dust outflow will be very low, making it very hard to detect astronomically as a radial extension of the debris disk. Additionally, the lack of velocity information from usual infrared measurements will not allow to distinguish between an outflow and a bound disk.
The face-on optical depth of a disk composed of grains of radius $`a`$ and observed at frequency $`\nu `$ is given by (Backman & Paresce 1993): $`\tau `$($`r`$,$`\nu `$)= $`\sigma `$($`r`$)($`\xi `$$`a`$$`\nu `$/c)<sup>q</sup>; where $`\sigma `$($`r`$) cm<sup>2</sup>/cm<sup>2</sup> is the face-on fractional geometric surface density; it is equal to the surface density $`n`$($`r`$), multiplied by the geometric cross section of the grain, $`\sigma `$($`r`$)=$`n`$($`r`$)$`\pi a^\mathit{2}`$. $`\xi `$ is the ratio between the critical wavelength $`\lambda `$<sub>0</sub> up to which the grain absorbs and emits radiation efficiently) and the grain radius $`a`$, and depends on the grain properties (e.g. $`\xi `$$``$$`\lambda `$<sub>0</sub>/$`a`$$``$2$`\pi `$, 1/2$`\pi `$ and 1, for strongly, weakly and moderately absorbing materials; we will use $`\xi `$$``$1). $`q`$ is the power law index of the emissive efficiency $`ϵ`$, such that for $`\lambda `$$`<`$$`\lambda `$<sub>0</sub>, $`ϵ`$$``$1, but for longer wavelengths the emissive efficiency decreases as $`ϵ`$=$`ϵ`$<sub>0</sub>($`\lambda `$<sub>0</sub>/$`\lambda `$)<sup>q</sup>; for the intermediate size regime, where $`a`$ is larger than $`\lambda `$<sub>peak</sub> of the incoming radiation (absorbs efficiently) but smaller than $`\lambda `$<sub>peak</sub> of the grain thermal emission (emits inefficiently), $`q`$=1. And c is the velocity of light.
We can estimate the surface density $`n`$($`r`$) (cm<sup>-2</sup>) at a distance $`r`$ from the central star from mass conservation by equating the mass that is produced in time dt, dN=$`\mathrm{𝑑𝑝𝑟}`$$`f`$<sub>ej</sub>dt, with the mass that crosses the annulus of radius $`r`$ in time dt, dN=$`n`$($`r`$)2$`\pi `$$`r`$$`v`$dt. $`\mathrm{𝑑𝑝𝑟}`$ is the dust production rate in particles per second; $`f`$<sub>ej</sub> is the fraction of particles that are ejected (our numerical studies find $`f`$<sub>ej</sub>$``$50–90%); and $`v`$ is the velocity of the particles at distance $`r`$, for large distances we will take $`v`$$``$$`v`$<sub>esc</sub>=(2$`G`$M/$`r`$)<sup>1/2</sup>. Solving for $`n`$($`r`$) and substituting into $`\sigma `$($`r`$),
$$\tau _{}^{outflow}(r,\nu )=\sigma (r)(\frac{\xi a\nu }{c})^q=\frac{\mathrm{𝑑𝑝𝑟𝑓}_{\mathrm{𝑒𝑗}}}{\mathit{2}\pi r(\mathit{2}GM_{\mathrm{}}/r)^{\mathit{1}/\mathit{2}}}\pi a^\mathit{2}(\frac{\xi a\nu }{c})^q.$$
(2)
We can estimate the optical depth of the solar system’s outflow using the KB dust production rates derived by Landgraf et al. (2002), which are based on Pioneer 10 and 11 measurements and for the Kuiper Belt gives $`\mathrm{𝑑𝑝𝑟}`$$``$2$`\times `$10<sup>14</sup> particles/s (for particles between 0.01 and 6 mm). Because the size distribution is very steep, one can assume that most of the detections are caused by particles just above the detection threshold, i.e. particles with $`a`$$``$5 $`\mathrm{\mu m}`$. For this particle size, $`\beta `$$``$0.05 and $`f`$<sub>ej</sub>$``$0.8, and the optical depth at 60$`\mathrm{\mu m}`$ ($`\nu `$=5$`\times `$10<sup>12</sup>Hz) will then be $`\tau `$<sup>outflow</sup>($`r`$,$`\nu `$)= 2.6$`\times `$10<sup>-14</sup>/$`r`$<sup>1/2</sup> (where $`r`$ is in AU). We can compare this to the optical depth of the Kuiper Belt (bound) disk. From Fig. 11 in Moro-Martín & Malhotra (2002) we can get the surface density that corresponds to a fictitious dust production rate of 100 particles per 1000 years, $`n`$$``$300 particles/AU<sup>2</sup>. Scaling up this density to account for the dust production rate found by Landgraf et al. (2002), we find that $`n`$$``$8.4$`\times `$10<sup>-2</sup> particles/cm<sup>2</sup>, $`\sigma `$$``$6.6$`\times `$10<sup>-8</sup>, so that $`\tau `$<sup>disk</sup>$``$5.5$`\times `$10<sup>-9</sup>. For the solar system, the ratio of the two optical depths is then $``$10<sup>-6</sup>. Other models for the Kuiper Belt dust disk give $`\sigma `$$``$10<sup>-6</sup> (15 times larger than our value; Backman, Dasgupta & Stencel 1995). It is estimated that for a system at 30 pc, the 70 $`\mathrm{\mu m}`$ $`\mathrm{𝑀𝐼𝑃𝑆}`$ array in $`\mathrm{𝑆𝑝𝑖𝑡𝑧𝑒𝑟}`$ will be able to detect a disk with $`\sigma `$$``$3$`\times `$10<sup>-6</sup> (D. Backman, private communication). This means that in order to see the Kuiper Belt dust disk the dust production rate will need to be increased by a factor of $``$3 in Backman’s models, or a factor of $``$45 in our models (using Landgraf’s $`\mathrm{𝑑𝑝𝑟}`$). But in order to see the outflow it will need to be increased by a factor of $``$6$`\times `$10<sup>6</sup> (Backman’s) or 9$`\times `$10<sup>7</sup> (ours). In any case, this increase will make the bound disk be optically thick. In other words, for an optically thin debris disk (where our dynamical models are valid), this outflow is very unlikely to be detected. For younger and more massive edge-on systems, after the giant planets have already formed, it may be possible to detect the outflow out of the plane. In this geometry, the signature of the off-plane outflow will be clearer against the fainter background. However, our dynamical models are not valid in this high-density regime where collisional effects dominate over P-R drag. It is possible that such an outflow may have already been detected with the Advanced Meteor Orbit Radar, which senses plasma signatures produced by extra-terrestrial dust particles ablating in the Earth’s atmosphere: Taylor, Baggaley & Steel (1996) and Baggaley (2000) claim that the main discrete source seems to coincide in direction with $`\beta `$ Pictoris.
#### 3.2.2 Interpretation of in-situ dust detections made by space probes
Recent Ulysses and Galileo dust experiments have led to the surprising discovery of interstellar grains sweeping through the solar system deep within the heliosphere (Grun, Zook & Baguhl 1993). Previously, interstellar grains could only be studied by extinction and polarization measurements of optical starlight, not sensitive to grains larger than 0.3 microns because of their small contribution to the optical cross section, and by infrared emission. These in-situ detections allowed us for the first time to study the mass distribution of interstellar grains within the heliosphere, leading to the surprising discovery of a population of large particles ($`>10^{16}`$ kg, Grun et al. 1994) that are 30 times more massive than the interstellar grains that cause stellar extinction. This finding implies that more mass is locked up in large grains locally than has been estimated from the astronomical measurements. The gas-to-dust ratio derived from astronomical measurements (400–600) is found to be much larger than the value of $``$100 derived from the in-situ detections, implying that the local interstellar cloud exceeds cosmic abundances (Frisch et al. 1999). These very important results rely critically on the correct identification of the origin of the dust grains. This identification is based on a geometrical argument: the direction the grains are coming from, with interstellar grains coinciding with the flow of neutral helium through the solar system; and a dynamical argument: the impact velocity and the expectation that only interstellar grains are on unbound hyperbolic orbits (Grun, Zook & Baguhl 1993). Under the current understanding, the sources of meteoroids in interplanetary space and their orbital properties are assumed as follows: Asteroids: low eccentricity and inclination; Comets: high eccentricity and inclination; Kuiper Belt: low eccentricity and inclination; and Interstellar: hyperbolic, and aligned with the direction of flow of the interstellar gas. However, we have shown in this paper that $``$80–90% of large Kuiper Belt grains ($`\beta <0.5`$) are gravitationally scattered outward by Jupiter and Saturn into hyperbolic orbits; therefore there is the potential of misinterpreting these escaping interplanetary particles as interstellar. In addition, other sources exist such as comets, Asteroid Belt and Trojan asteroids. Due to radiation pressure, some of the dust particles released at those locations will be set on Jupiter crossing orbits, so in principle close encounters with Jupiter could take place resulting on hyperbolic orbits. In the future, we plan to study whether or not these particles may have been detected by Ulysses and Galileo. For the analysis of future in-situ dust detections in the outer solar system, such as with the Cassini Cosmic Dust Analyzer and the Interstellar Probe, it will be important to keep in mind the existence of the large-particle outflow of solar system dust to correctly identify the origin of the massive fast moving particles, whether interplanetary or interstellar. It has been recently announced that the analysis of the ion charge signals in the Cassini dust detector, together with geometric and kinematic considerations, have led to the identification of an interstellar flux at 0.8 AU that is in agreement with the flux measured by Ulysses at 3 AU at the same time (Altobelli et al 2003). But any dust detections by Cassini outside Jupiter’s orbit have not yet been reported.
## 4 On how debris disks with inner gaps signal the presence of massive planets
Recent GTO and $`\mathrm{𝐹𝐸𝑃𝑆}`$ observations with the $`\mathrm{𝑆𝑝𝑖𝑡𝑧𝑒𝑟}`$ MIPS instrument suggest that debris disks and giant planets co-exist and that inner gaps appear to be common in cold Kuiper Belt-like disks (Beichman et al. 2005 and Kim et al. 2005). In view of these observations, it is interesting to study the efficiency of particle ejection ($`\mathrm{\S }`$3.1) and the resulting dust density contrast inside and outside the orbit of the planet, as a function of the planet’s mass and orbital elements and the particle size. It is important to keep in mind, however, that the modeling presented here does not consider the effect of particle collisions, which together with P-R drag could also be responsible for the opening of an inner gap in the dust disk (Wyatt 2005).
If the particles were drifting inward at a constant rate, as set by P-R drag, the ratio n<sub>in</sub>/n<sub>1000</sub> (from Tables 1, 2 and 3) would directly give us an estimate of the density contrast inside and outside the inner boundary of the disk. However, the trapping of particles in MMRs with the planet halts the P-R drift, increasing the number density of particles in that region. The density contrast, therefore, can only be estimated using the radial density profiles that result from the numerical simulations. Fig. 12 shows some of these profiles for a representative set of models. These results, keeping in mind the uncertainties due to the fact that we are modeling the dynamical evolution of a small number of test particles (N$``$100), can help us estimate what planet masses and semimajor axes may be responsible for the inner gaps that are inferred indirectly from the disks’ SEDs, or in few cases, that are seen directly in spatially resolved images. For planets located at 1–30 AU with masses of 1–10M<sub>Jup</sub>, the ratio between the density outside and inside the orbit of the planet is $``$40, whereas for planet masses of 0.03–0.3M<sub>Jup</sub>, this ratio is in the range 3–10. The models show that the radius of the inner depleted region, r<sub>gap</sub>, depends on the mass and the eccentricity of the planet. In Table 4 we show that for the models with planets in circular orbits, r<sub>gap</sub>$``$0.8$`\times `$a<sub>pl</sub> for 1–3M<sub>Jup</sub> and $``$1.2$`\times `$a<sub>pl</sub> for 10M<sub>Jup</sub>. The three bottom planels of Fig. 12 show that for planets with eccentricities in the range 0.3–0.5, the surface density decreases more smoothly and consequently the dust disk would not present a sharp inner edge.
## 5 Conclusions
When a massive planet is located interior to a belt of dust-producing planetesimals, dynamical models have shown that as the dust particles drift inward due to P-R drag, they get trapped in MMRs with the planet, and this well-known effect can sculpt the dust disk creating rings, warps and azimuthal asymmetries. In addition to the trapping in MMRs, gravitational scattering with the planet is responsible for the depletion of dust inside the orbit of the planet. Although this is also a well known effect, to our knowledge it has not been studied in detail in the past. In this paper we have shown that the ejected dust particles form an “outflow”, whose angular confinement, velocity and symmetry depend on the planet’s mass and orbital elements, as well as the particle size. The high efficiency of ejection (for planet masses $``$1M<sub>Jup</sub>), together with the possible high frequency of debris disks harboring massive planets, suggest that these outflows may be a common phenomenon. If this is the case, they may contribute significantly or even dominate the clearing of circumstellar debris in planetary systems, enriching the immediate vicinity of star-forming regions with large dust particles and affecting therefore the particle size distribution of their local ISM. In addition, we have seen how the ejection of particles is responsible for the clearing of dust inside the orbit of the planet, creating a density contrast that can be measured directly in spatially resolved images or indirectly through the modeling of the SED of the debris disk. Indeed, recent $`\mathrm{𝑆𝑝𝑖𝑡𝑧𝑒𝑟}`$ observations suggest that debris disks and giant planets co-exist and that inner gaps appear to be common in cold Kuiper Belt-like disks (Beichman et al. 2005, Kim et al. 2005). To aid in the interpretation of these observations, we have studied the efficiency of particle ejection and the resulting dust density contrast inside and outside the orbit of the planet, as a function of the planet’s mass and orbital elements and the particle size.
Acknowledgments
We thank Hal Levison for providing the SKEEL computer code, Alberto Noriega-Crespo, Dana Backman, George Rieke, Re’em Sari and Mark Sykes for useful discussion, and the anonymous referee for very helpful comments on how to improve the manuscript. This work is part of the Spitzer FEPS Legacy project (http://feps.as.arizona.edu). We acknowledge NASA for research support (contract 1224768 administered by JPL and grants NAG5-10343 and NAG5-11661), and IPAC, the $`\mathrm{𝑆𝑝𝑖𝑡𝑧𝑒𝑟}`$ Science Center and the Max-Plank-Institute in Heidelberg for providing access to their facilities. AMM is also supported by the Lyman Spitzer Fellowship at Princeton University.
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# Sets of rigged paths with Virasoro characters
## 1. Introduction
Consider a representation $`M`$ of a vertex operator algebra $`𝒱`$. Let $`\varphi ^{(i)}(z)𝒱`$ be a collection of fields, and let $`\varphi _n^{(i)}`$ be the corresponding Fourier coefficients.
Question. Find a set $`I`$ of sequences of pairs $`(i_j,n_j)`$ and a vector $`vM`$ such that vectors $`\{\varphi _{n_L}^{(i_L)}\mathrm{}\varphi _{n_1}^{(i_1)}\varphi _{n_0}^{(i_0)}v|(i_j,n_j)I\}`$ form a basis of $`M`$.
Vectors of such form are usually called “monomials”. The question of finding a monomial basis, i.e., a basis consisting of monomials, is an important, well-known and old problem, solved in many interesting non-trivial cases. Examples include: integrable $`\widehat{𝔰𝔩}_n`$ modules in terms of $`e_{1i}(z)`$ currents \[LP\] and \[Pr\]; “big” and “small” coinvariants of $`\widehat{𝔰𝔩}_2`$ integrable modules in terms of the currents $`e(z),f(z),h(z)`$ \[FKLMM1\] and \[FKLMM2\]; $`(2,2n+1)`$-Virasoro minimal series in terms of the Virasoro current \[FF\]; and $`(3,3n\pm 1)`$-Virasoro minimal modules tensored with Fock spaces in terms of an abelian current \[FJM\].
Following the same philosophy, in \[FJMMT\] we constructed a basis of the $`(p,p^{})`$-Virasoro minimal series representations $`M_{r,s}^{(p,p^{})}`$ ($`1rp1`$, $`1sp^{}1`$) with $`1<p^{}/p<2`$ and $`p3`$, in terms of the $`(2,1)`$-primary field $`\varphi (z)`$. The basis has the form
(1.1)
$$\varphi _{n_L}^{(r_L,r_{L1})}\mathrm{}\varphi _{n_1}^{(r_1,r_0)}|b(s),s$$
where $`r_0=b(s),r_L=r,r_i=r_{i+1}\pm 1`$, $`\varphi _n^{(r,r^{})}:M_{r^{},s}^{(p,p^{})}M_{r,s}^{(p,p^{})}`$ are the Fourier coefficients of $`\varphi (z)`$, and $`n_i\mathrm{\Delta }_{r_i,s}\mathrm{\Delta }_{r_{i1},s}+`$. For each fixed $`s`$, $`1b(s)p1`$ is so chosen that the conformal dimension $`\mathrm{\Delta }_{r,s}`$ of the space $`M_{r,s}^{(p,p^{})}`$ attains the minimum at $`r=b(s)`$. We have also the condition
(1.2)
$$n_{i+1}n_iw(r_{i+1},r_i,r_{i1})(1iL1)$$
where $`w(r^{\prime \prime },r^{},r)`$ $`(r=r^{}\pm 1,r^{}=r^{\prime \prime }\pm 1)`$ are certain rational numbers (see (2.2)–(2.5) below). It was shown that the condition (1.2) is a consequence of the quadratic relations in the algebra generated by Fourier coefficients of the primary field $`\varphi (z)`$ .
In this paper we consider the case $`p^{}/p>2`$. In this case, in addition to the quadratic relations, the algebra has cubic relations. To our surprise, we found that in many cases (if not all) the cubic relations result in a very simple exclusion rule in addition to (1.2):
(1.3) $`n_{i+2}n_i1.`$
The condition (1.3) is void in the case $`1<p^{}/p<2`$ since it follows from (1.2). We conjecture that in the general case $`p<p^{}`$ the monomials satisfying (1.2) and (1.3) form a basis (see Conjecture 2.3 for the precise statement).
The purpose of this paper is to give some evidences for this conjecture. Our results are two-fold. First, we show that the combinatorial set defined by conditions (1.2), (1.3) (which we call the set of ‘rigged paths’) has the same graded character as that of the Virasoro module for any $`(p,p^{})`$ minimal theory with $`s=1`$ (see Theorem 2.2). Second, we prove the conjecture for the $`(3,3n\pm 1)`$ minimal theories with $`s=1`$ (see Theorem 5.5).
The first statement is purely combinatorial. We prove it by showing that the characters of Virasoro modules and the characters of the combinatorial sets enjoy the same recurrence when the parameter changes from $`p^{}p`$ to $`p^{}`$. (see Proposition 3.2 and Proposition 3.3).
The second statement is based on the work \[FJM\], in which the representation space
$`V^{(3,p^{})}=\left(M_{1,1}^{(3,p^{})}(_{n:\mathrm{even}}_{n\beta })\right)\left(M_{2,1}^{(3,p^{})}(_{n:\mathrm{odd}}_{n\beta })\right)`$
of the abelian current $`a(z)=\varphi (z)\mathrm{\Phi }_\beta (z)`$ is studied. Here $`\mathrm{\Phi }_\beta (z):_{n\beta }_{(n+1)\beta }`$ is a certain bosonic vertex operator acting on the bosonic Fock space $`_{n\beta }`$. It was shown that the abelian current $`a(z)`$ satisfies cubic relations, and by exploiting the cubic relations a monomial basis of the space $`W^{(3,p^{})}`$ generated from the vector $`|1,1|0`$, where $`|0_0`$ is the highest weight vector, was constructed. From the construction for $`W^{(3,p^{}3)}`$, we deduce a monomial basis of the space
$`M^{(3,p^{})}=M_{1,1}^{(3,p^{})}M_{2,1}^{(3,p^{})}.`$
The story is as follows. We construct a filtration of the space $`M^{(3,p^{})}`$ by using the operator $`\varphi (z)`$. This filtration induces a current $`\stackrel{~}{a}(z)`$ acting on the corresponding graded space, from the operator $`a(z)`$ acting on $`V^{(3,p^{})}`$. The correlation functions for the operator $`\stackrel{~}{\varphi }(z)`$ are equal to those for the operator $`\stackrel{~}{a}(z)`$ up to simple factors. The latter belong to the space of correlation functions of the operator $`a(z)`$ with $`p^{}`$ replaced by $`p^{}3`$ up to simple factors. The spanning property of the monomials (1.1) with (1.2), (1.3) is deduced by using these identifications of correlation functions and the monomial basis of the space $`W^{(3,p^{}3)}`$ constructed in\[FJM\].
The plan of our paper is as follows. In Section 2 we define the combinatorial set of rigged paths and formulate Conjecture 2.3. In Section 3 we show that Virasoro characters satisfy a recurrence relation when the parameter changes from $`p^{}p`$ to $`p^{}`$. In Section 4 we show that the combinatorial sets of rigged paths satisfy the same recurrence relation. Section 5 is devoted to the proof of the conjecture in the case $`(p,p^{})=(3,3n\pm 1)`$.
In the paper by P. Jacob and P. Mathieu \[JM\], the problem of constructing monomial basis is studied, and combinatorial conditions similar to ours are proposed. Their study is restricted to the $`(3,p)`$ case, and in this case the paths $`(r_L,\mathrm{},r_1)`$ do not appear in combinatorics. We thank one of the referees for attracting our attentions to this paper.
We also thank another referee for providing us with the simple proof of Proposition 3.4 as given in the below.
## 2. Rigged paths
### 2.1. Minimal series
We recall some definitions about minimal conformal field theory. For the details we refer to \[DMS\]. Let Vir be the Virasoro algebra with the standard $``$-basis $`\{L_n\}_n`$ and $`c`$ satisfying
$`[L_m,L_n]=(mn)L_{m+n}+{\displaystyle \frac{c}{12}}m(m^21)\delta _{m+n,0},[c,L_n]=0.`$
Fix a pair $`(p,p^{})`$ of relatively prime positive integers. We set
$`t={\displaystyle \frac{p^{}}{p}}.`$
We assume $`p3`$ so that the $`(2,1)`$ primary field exists (see below). We consider the minimal series of representations $`M_{r,s}^{(p,p^{})}`$ ($`1rp1,1sp^{}1`$) of Vir. The module $`M_{r,s}^{(p,p^{})}`$ is generated by a vector $`|r,s`$ called the highest weight vector. It is an irreducible module. The central element $`c`$ acts as the scalar
$`c_{p,p^{}}=136\left(t+{\displaystyle \frac{1}{t}}\right).`$
The highest weight vector satisfies
$`L_n|r,s=0\text{ if }n>0,L_0|r,s=\mathrm{\Delta }_{r,s}^{(t)}|r,s,`$
where
$$\mathrm{\Delta }_{r,s}^{(t)}=\frac{(rts)^2(t1)^2}{4t}.$$
We fix $`s`$. The $`(2,1)`$ primary field
$`\varphi ^{(r\pm 1,r)}(z)={\displaystyle \underset{n\mathrm{\Delta }_{r,s}^{(t)}+\mathrm{\Delta }_{r\pm 1,s}^{(t)}}{}}\varphi _n^{(r\pm 1,r)}z^{n\mathrm{\Delta }_{2,1}^{(t)}}`$
is a generating series of linear operators $`\varphi _n^{(r\pm 1,r)}`$ acting as
$`\varphi _n^{(r\pm 1,r)}:M_{r,s}^{(p,p^{})}M_{r\pm 1,s}^{(p,p^{})}.`$
Up to a scalar multiple, they are characterized by the commutation relations with the Virasoro generators:
$`[L_n,\varphi ^{(r\pm 1,r)}(z)]=z^n\left(z+(n+1)\mathrm{\Delta }_{2,1}^{(t)}\right)\varphi ^{(r\pm 1,r)}(z).`$
The module $`M_{r,s}^{(p,p^{})}`$ is graded by eigenvalues of the operator $`L_0`$:
$`M_{r,s}^{(p,p^{})}`$ $`=`$ $`_d\left(M_{r,s}^{(p,p^{})}\right)_d,`$
$`(M_{r,s}^{(p,p^{})})_d`$ $`=`$ $`\{|vM_{r,s}^{(p,p^{})}L_0|v=d|v\}.`$
The character $`\chi _{r,s}^{(p,p^{})}(q)`$ is defined by the formula
$`\chi _{r,s}^{(p,p^{})}(q)=\mathrm{tr}_{M_{r,s}^{(p,p^{})}}q^{L_0}.`$
The primary field preserves the grading:
$`\varphi _n^{(r^{},r)}(M_{r,s}^{(p,p^{})})_d(M_{r^{},s}^{(p,p^{})})_{d+n}.`$
### 2.2. Quadratic relations
In \[FJMMT\] we derived a set of quadratic relations for the Fourier components of the $`(2,1)`$ primary field of the following form:
(2.1)
$$\underset{\genfrac{}{}{0pt}{}{r^{}}{n+n^{}=d}}{}c_{n,n^{}}^{(r,r^{},r^{\prime \prime })}\varphi _n^{(r,r^{})}\varphi _n^{}^{(r^{},r^{\prime \prime })}=0,$$
where $`r,r^{\prime \prime }`$ and $`d`$ are fixed, and the coefficients $`c_{n,n^{}}^{(r,r^{},r^{\prime \prime })}`$ are such that if $`n^{}>N`$ for some $`N`$ they are zero except for finitely many of them. Let us describe a consequence of the quadratic relations in a little bit different way than in \[FJMMT\]. Define weights $`w^{(t)}(r,r^{},r^{\prime \prime })=w^{(t)}(r^{\prime \prime },r^{},r)=w^{(t)}(pr,pr^{},pr^{\prime \prime })`$ where $`rr^{},r^{}r^{\prime \prime }\{1,1\}`$ by
(2.2) $`w^{(t)}(r,r\pm 1,r\pm 2)`$ $`=`$ $`{\displaystyle \frac{t}{2}},\text{if }1r,r\pm 1,r\pm 2p1,`$
(2.3) $`w^{(t)}(r,r+1,r)`$ $`=`$ $`2{\displaystyle \frac{t}{2}}+[rt]rt,\text{if }2r,r+1p1,`$
(2.4) $`w^{(t)}(r,r1,r)`$ $`=`$ $`1{\displaystyle \frac{t}{2}}[rt]+rt,\text{if }1r1,rp2,`$
(2.5) $`w^{(t)}(1,2,1)`$ $`=`$ $`w^{(t)}(p1,p2,p1)=3{\displaystyle \frac{3t}{2}}.`$
Here $`[x]`$ is the integer part of $`x`$.
###### Proposition 2.1.
Any monomial of the form
(2.6)
$$\varphi _{n_L}^{(r_L,r_{L1})}\mathrm{}\varphi _{n_1}^{(r_1,r_0)}$$
can be written as an infinite linear combination of monomials satisfying
(2.7)
$$n_{i+1}n_iw^{(t)}(r_{i+1},r_i,r_{i1})0.$$
###### Proof.
Essentially, the proof is given in \[FJMMT\] (see Proposition 3.3 in that paper). Note, however, that in \[FJMMT\] the weights $`w^{(t)}(1,2,1)`$ and $`w^{(t)}(p1,p2,p1)`$ are given as special cases of (2.3) and (2.4), respectively. This is because the range of $`t`$ was restricted to $`1<t<2`$, wherein the two expressions coincide. In \[FJMMT\], the quadratic relations were derived without the restriction on $`t`$, and the statement that any monomial of the form (2.6) can be rewritten into an infinite linear combination of those which satisfy (2.7) was, in effect, proved without the restriction $`t<2`$ by using (2.5), instead of using the special cases of (2.3) and (2.4). ∎
We give a few more remarks.
As stated above, if we rewrite a monomial by using the relations of the form (2.1), we obtain an infinite linear combination. Note, however, that if we fix the degree $`n_1+\mathrm{}+n_L`$ and restrict $`n_1N`$ for some $`N`$, there exists only finitely many monomials satisfying the condition (2.7). Therefore, such infinite sums are meaningful after completion.
Let us consider vectors of the form
(2.8)
$$\varphi _{n_L}^{(r_L,r_{L1})}\mathrm{}\varphi _{n_1}^{(r_1,r_0)}|r_0,s,$$
where $`|r_0,s`$ is the highest weight vector of $`M_{r_0,s}^{(p,p^{})}`$. By abuse of terminology we call such vectors monomials.
Proposition 2.1 implies that any vector of the above form belongs to the linear span (in the sense of finite linear combinations) of those satisfying (2.7) and the highest weight condition
$`n_1+\mathrm{\Delta }_{r_0,s}^{(t)}\mathrm{\Delta }_{r_1,s}^{(t)}0.`$
### 2.3. Statement of the results
Now we restrict to the case
$`t>1\text{ and }(r_0,s)=(1,1).`$
We say a monomial
(2.9)
$$\varphi _{n_L}^{(r_L,r_{L1})}\mathrm{}\varphi _{n_1}^{(r_1,r_0)}|1,1$$
is admissible if and only if (2.7), and
(2.10)
$$n_1\mathrm{\Delta }_{2,1}^{(t)}=\frac{3t2}{4},$$
and
(2.11)
$$n_{i+2}n_i1$$
when $`r_{i1},r_i,r_{i+1}r_{i+2}\{r,r+1\}`$ for some $`1rp2`$, $`1iL2`$, are satisfied.
Set
(2.12) $`v^{(t)}(r)`$ $`=`$ $`1w^{(t)}(r,r+1,r)w^{(t)}(r+1,r,r+1).`$
We have
$`v^{(t)}(r)`$ $`=`$ $`\{\begin{array}{cc}p^{}5\hfill & \text{if }p=3;\hfill \\ [2t3]\hfill & \text{if }p>3\text{ and }r=1,p2;\hfill \\ [(r+1)t][rt]2\hfill & \text{if }1<r<p2.\hfill \end{array}`$
Although $`w^{(t)}`$ is not necessarily an integer, $`v^{(t)}`$ are integers. Note also that $`v^{(t)}(r)=v^{(t)}(p1r)`$.
A rigged path of length $`L`$ is a table of integers of the form
$`P=\left(\begin{array}{ccccc}r_L& r_{L1}& \mathrm{}& r_1& r_0\\ & \sigma _{L1}& \mathrm{}& \sigma _1& \sigma _0\end{array}\right)`$
where $`r_0=1`$, $`1r_ip1`$ $`(0iL)`$ and
(2.13)
$$r_{i+1}r_i\{1,1\}(0iL1).$$
In the usual terminology, a sequence of integers $`(r_i)`$ satisfying (2.13) is called a path. A rigged path is decorated by the rigging $`(\sigma _i)`$. For brevity, we often call a rigged path simply a path.
A path is called admissible at level $`t`$ if $`\sigma _i0(0iL1)`$, and
(2.14)
$$\sigma _i+\sigma _{i+1}v^{(t)}(r)$$
when $`r_{i1},r_i,r_{i+1},r_{i+2}\{r,r+1\}`$ for some $`1rp2`$, $`1iL2`$.
Note that these conditions correspond to (2.10), (2.7) and (2.11), respectively, if we set
$`\sigma _0=n_1\mathrm{\Delta }_{2,1}^{(t)},\sigma _i=n_{i+1}n_iw^{(t)}(r_{i+1},r_i,r_{i1})(1iL1).`$
If $`t<2`$, we have $`v^{(t)}(r)0`$ and the condition (2.14) follows from the positivity of $`\sigma _i`$’s. If $`t>2`$, we have
$`v^{(t)}(r)0(1rp2),`$
and moreover
(2.15) $`v^{(t)}(1)=v^{(t)}(p2)1.`$
We denote by $`C_L^{(t)}`$ the set of rigged paths of length $`L`$ which are admissible at level $`t`$, and by $`C_{L,r}^{(t)}`$ the subset of $`C_L^{(t)}`$ consisting of paths such that $`r_L=r`$. The subset $`C_{L,r}^{(t)}`$ is empty unless $`rL+1mod2`$.
We define the degree of $`PC_L^{(t)}`$ by
$`d(P)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{L}{}}}n_i`$
$`=`$ $`L\mathrm{\Delta }_{2,1}^{(t)}+{\displaystyle \underset{i=1}{\overset{L1}{}}}(Li)w^{(t)}(r_{i+1},r_i,r_{i1})+{\displaystyle \underset{i=0}{\overset{L1}{}}}(Li)\sigma _i,`$
and the character of $`C_{L,r}^{(t)}`$ by
$`\mathrm{ch}_qC_{L,r}^{(t)}={\displaystyle \underset{PC_{L,r}^{(t)}}{}}q^{d(P)}.`$
Our main result is the following identity:
###### Theorem 2.2.
We have
$`\chi _{r,1}^{(p,p^{})}(q)={\displaystyle \underset{L0}{}}\mathrm{ch}_qC_{L,r}^{(t)}.`$
Theorem 2.2 will be proved at the end of Section 4. This result motivates us to make the following conjecture:
###### Conjecture 2.3.
For $`1rp1`$, the set of admissible monomials of the form (2.8), where $`r_L=r,r_0=1`$ and $`s=1`$, is a basis of $`M_{r,1}^{(p,p^{})}`$.
The case $`1<t<2`$ of the conjecture has been proved in \[FJMMT\]. We prove Conjecture 2.3 in the case $`p=3`$ in Section 5.
## 3. Recurrence structure
### 3.1. Recurrence relation for Virasoro characters
Recall the formula for $`\chi _{r,s}^{(p,p^{})}(q)`$ \[RC\]:
(3.1) $`\chi _{r,s}^{(p,p^{})}(q)={\displaystyle \frac{q^{\mathrm{\Delta }_{r,s}^{(t)}}}{(q)_{\mathrm{}}}}\left({\displaystyle \underset{n}{}}q^{pp^{}n^2+(p^{}rps)n}{\displaystyle \underset{n}{}}q^{pp^{}n^2+(p^{}r+ps)n+rs}\right).`$
Here $`(q)_{\mathrm{}}=_{j=1}^{\mathrm{}}(1q^j)`$. In the case $`1<t<2`$, $`\chi _{r,1}^{(p,p^{})}(q)`$ was written in the following form, \[FJMMT, W\]:
(3.2) $`\chi _{r,1}^{(p,p^{})}(q)=q^{\mathrm{\Delta }_{r,1}^{(t)}}{\displaystyle \underset{mr1\mathrm{mod}\mathrm{\hspace{0.17em}2}}{}}{\displaystyle \frac{1}{(q)_m}}K_{m,r}^{(p,p^{}p)}(q),`$
where
(3.3) $`K_{m,r}^{(p,\overline{p^{}})}(q)=q^{\frac{m^2(r1)^2}{4}}{\displaystyle \underset{n}{}}q^{p\overline{p^{}}n^2+\overline{p^{}}nr}\left(\left[{\displaystyle \genfrac{}{}{0pt}{}{m}{\frac{mr+1}{2}pn}}\right]\left[{\displaystyle \genfrac{}{}{0pt}{}{m}{\frac{m+r+1}{2}+pn}}\right]\right).`$
Here we have used the notation
$`(q)_n={\displaystyle \underset{j=1}{\overset{n}{}}}(1q^j),\left[{\displaystyle \genfrac{}{}{0pt}{}{m}{n}}\right]={\displaystyle \frac{(q)_m}{(q)_n(q)_{mn}}}.`$
The identity (3.2) can be generalized to the case $`t>2`$:
###### Proposition 3.1.
Set $`k=[t]`$. Then the following equality holds:
(3.4) $`\chi _{r,1}^{(p,p^{})}(q)=q^{\mathrm{\Delta }_{r,1}^{(t)}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{m_0,\mathrm{},m_{k1}0}{m_0r1\mathrm{mod}\mathrm{\hspace{0.17em}\hspace{0.17em}2}}}{}}{\displaystyle \frac{q^{Q^{(k)}(m_0,\mathrm{},m_{k1})\frac{k1}{4}(r^21)}}{(q)_{m_0}\mathrm{}(q)_{m_{k1}}}}K_{m_0,r}^{(p,p^{}kp)}(q).`$
Here
$`Q^{(k)}(m_0,\mathrm{},m_{k1})`$ $`=`$ $`{\displaystyle \frac{k1}{4}}m_0^2+{\displaystyle \underset{j=1}{\overset{k1}{}}}(kj)m_j^2`$
$`+`$ $`{\displaystyle \underset{j=1}{\overset{k1}{}}}(kj)m_0m_j+2{\displaystyle \underset{1j<j^{}k1}{}}(kj^{})m_jm_j^{}`$
$`+`$ $`{\displaystyle \frac{k1}{2}}m_0+{\displaystyle \underset{j=1}{\overset{k1}{}}}(kj)m_j.`$
Proposition 3.1 will be proved in Section 3.2.
For $`L_0`$ denote by $`\chi _{r,1:L}^{(p,p^{})}(q)`$ the right hand side of (3.4) with the sum replaced by the partial sum over $`m_0,\mathrm{},m_{k1}`$ satisfying $`m_0+2(m_1+\mathrm{}+m_{k1})=L`$. Then we have
$`\chi _{r,1}^{(p,p^{})}(q)={\displaystyle \underset{L0}{}}\chi _{r,1;L}^{(p,p^{})}(q).`$
It follows from Proposition 3.1 that
###### Proposition 3.2.
We have
(3.5) $`\chi _{r,1;L}^{(p,p^{})}(q)={\displaystyle \underset{m0}{}}{\displaystyle \frac{q^{\frac{L^2}{4}+\frac{L}{2}}}{(q)_m}}\chi _{r,1;L2m}^{(p,p^{}p)}(q).`$
In Section 4 we prove (see Theorem 4.13)
###### Proposition 3.3.
(3.6) $`\mathrm{ch}_qC_{L,r}^{(t)}={\displaystyle \underset{m0}{}}{\displaystyle \frac{q^{\frac{L^2}{4}+\frac{L}{2}}}{(q)_m}}\mathrm{ch}_qC_{L2m,r}^{(t1)}.`$
Theorem 2.2 follows from these identities.
### 3.2. Proof of Proposition 3.1
In the following we set $`1/(q)_n=0`$ for $`n<0`$. To prove Proposition 3.1 we use the following formula (the proof given here is provided by one of the referees).:
###### Proposition 3.4.
For $`l_0`$ and $`\mu `$, we have
(3.7) $`{\displaystyle \frac{1}{(q)_{\mathrm{}}}}={\displaystyle \underset{N_0,\mathrm{},N_l0}{}}{\displaystyle \frac{q^{_{j=0}^lN_j^2+\mu _{j=0}^lN_j}}{(q)_{N_0+\mu }(q)_{N_0}(q)_{N_1N_0}\mathrm{}(q)_{N_lN_{l1}}}}.`$
###### Proof.
Consider the following functions $`g_l(z)(l=0,1,\mathrm{})`$:
$`g_l(z)={\displaystyle \underset{N_0,\mathrm{},N_l0}{}}{\displaystyle \frac{z^{_{j=0}^lN_j}q^{_{j=0}^lN_j^2}}{(zq)_{N_0}(q)_{N_0}(q)_{N_1N_0}\mathrm{}(q)_{N_lN_{l1}}}},`$
where $`(z)_n:=_{j=0}^{n1}(1zq^j)`$. Let us prove that $`g_l(z)=1/(zq)_{\mathrm{}}`$. Then we obtain (3.7) by setting $`z=q^\mu `$.
First note that $`g_0(z)`$ can be rewritten in terms of the basic hypergeometric series $`{}_{2}{}^{}\varphi _{1}^{}`$ as follows:
$`g_0(z)=\underset{a,b\mathrm{}}{lim}{}_{2}{}^{}\varphi _{1}^{}(a,b;zq;q,zq/ab),`$
where
$`{}_{2}{}^{}\varphi _{1}^{}(a,b;c;q,z):={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(a)_n(b)_n}{(c)_n(q)_n}}z^n.`$
On the other hand we have the $`q`$-Gauss sum identity (see, e.g., \[GR\]):
$`{}_{2}{}^{}\varphi _{1}^{}(a,b;c;q,c/ab)={\displaystyle \frac{(c/a)_{\mathrm{}}(c/b)_{\mathrm{}}}{(c)_{\mathrm{}}(c/ab)_{\mathrm{}}}}.`$
This implies that $`g_0(z)=1/(zq)_{\mathrm{}}`$.
Next consider the case where $`l>0`$. Rewrite the definition of $`g_l(z)`$ as
$`g_l(z)={\displaystyle \underset{N_1\mathrm{},N_l0}{}}{\displaystyle \frac{z^{_{j=1}^lN_j}q^{_{j=1}^lN_j^2}}{(q)_{N_1}(q)_{N_1N_0}\mathrm{}(q)_{N_lN_{l1}}}}f_{N_1}(z),`$
where
$`f_N(z)=(q)_N{\displaystyle \underset{j=0}{\overset{N}{}}}{\displaystyle \frac{z^jq^{j^2}}{(zq)_j(q)_j(q)_{Nj}}}.`$
It is easy to see that
$`f_N(z)=\underset{a\mathrm{}}{lim}{}_{2}{}^{}\varphi _{1}^{}(a,q^N;zq;q,zq^{N+1}/a).`$
Using the $`q`$-Gauss sum formula again, we see that $`f_N(z)=1/(zq)_N`$. This implies $`g_l(z)=g_{l1}(z)`$ and therefore we get $`g_l(z)=g_0(z)=1/(zq)_{\mathrm{}}`$. ∎
Proof of Proposition 3.1.
Substituting (3.3) to the right hand side of (3.4), we obtain two sums, which we refer to as I and II. First consider I:
(3.8) $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{m_0,\mathrm{},m_{k1}0}{m_0r1\mathrm{mod}\mathrm{\hspace{0.17em}\hspace{0.17em}2}}}{}}{\displaystyle \frac{q^{Q^{(k)}(m_0,\mathrm{},m_{k1})\frac{k1}{4}(r^21)}}{(q)_{m_0}\mathrm{}(q)_{m_{k1}}}}`$
$`\times q^{\frac{m_0^2(r1)^2}{4}}{\displaystyle \underset{n}{}}q^{p(p^{}kp)n^2+(p^{}kp)nr}\left[{\displaystyle \genfrac{}{}{0pt}{}{m_0}{\frac{m_0r+1}{2}pn}}\right].`$
Set
$`N_0={\displaystyle \frac{m_0r+1}{2}}pn,N_jN_{j1}=m_j(j=1,\mathrm{},k1)`$
and rewrite (3.8) as a summation over $`N_0,\mathrm{},N_{k1}0`$ and $`n`$. Then (3.8) becomes
$`{\displaystyle \underset{n}{}}q^{pp^{}n^2+(p^{}rp)n}F_k(r1+2pn),`$
where
$`F_k(\mu )={\displaystyle \underset{N_0,\mathrm{},N_{k1}0}{}}{\displaystyle \frac{q^{_{j=0}^{k1}N_j^2+\mu N_0+(\mu +1)_{j=1}^{k1}N_j}}{(q)_{N_0+\mu }(q)_{N_0}(q)_{N_1N_0}\mathrm{}(q)_{N_{k1}N_{k2}}}}.`$
In the same way rewrite II by setting
$`N_0={\displaystyle \frac{m_0+r+1}{2}}+pn,N_jN_{j1}=m_j(j=1,\mathrm{},k1).`$
The result is
$`{\displaystyle \underset{n}{}}q^{pp^{}n^2+(p^{}r+p)n+r}F_k(r12pn).`$
Thus we find that the right hand side of (3.4) is equal to
(3.9) $`q^{\mathrm{\Delta }_{r,1}^{(t)}}{\displaystyle \underset{n}{}}q^{pp^{}n^2+(p^{}rp)n}\left(F_k(r1+2pn)q^{r+2pn}F_k(r12pn)\right).`$
Set $`\mu =r+2pn`$. The last part of (3.9) reads $`F_k(\mu 1)q^\mu F_k(\mu 1)`$. In the sum $`F_k(\mu 1)`$, we replace $`N_j`$ by $`N_j+\mu `$. Then, we have
$`F_k(\mu 1)q^\mu F_k(\mu 1)`$
$`=(1q^\mu ){\displaystyle \underset{N_0,\mathrm{},N_{k1}0}{}}{\displaystyle \frac{q^{_{j=0}^{k1}N_j^2+\mu _{j=0}^{k1}N_j}}{(q)_{N_0+\mu }(q)_{N_0}(q)_{N_1N_0}\mathrm{}(q)_{N_{k1}N_{k2}}}}={\displaystyle \frac{1q^\mu }{(q)_{\mathrm{}}}}.`$
Here we used (3.7) in the last equality. Hence (3.9) is equal to the right hand side of the formula (3.1). This completes the proof. ∎
## 4. The bijection
### 4.1. Particles
In each path $`PC_L^{(t)}`$, we locate “particles”.
Recall, that a path $`PC_L^{(t)}`$ is a table of integers
(4.1) $`P=\left(\begin{array}{ccccc}r_L& r_{L1}& \mathrm{}& r_1& r_0\\ & \sigma _{L1}& \mathrm{}& \sigma _1& \sigma _0\end{array}\right).`$
The integers $`r_x`$ satisfy the conditions $`r_0=1`$, $`1r_xp1`$ $`(1xL)`$ and $`|r_xr_{x1}|=1`$ $`(1xL)`$. The integers $`\sigma _x`$ satisfy the conditions $`\sigma _x0`$ $`(0xL1)`$ and
(4.2) $`\sigma _x+\sigma _{x1}`$ $``$ $`v^{(t)}(r)(2xL1)`$
if
(4.3) $`r_{x+1},r_x,r_{x1},r_{x2}\{r,r+1\}\text{ for some }1rp2.`$
It is convenient to define $`\sigma _L=\mathrm{}`$, and use the convention that $`\mathrm{}\pm 1=\mathrm{}`$.
Our aim is to relate the sets $`C_L^{(t1)}`$ ($`L=0,1,2,\mathrm{}`$) to the sets $`C_L^{(t)}`$ ($`L=0,1,2,\mathrm{}`$). There is an injective mapping from $`C_L^{(t1)}`$ to $`C_L^{(t)}`$ such that the image of the mapping is equal to the subset of $`C_L^{(t)}`$ consisting of paths $`P`$ for which the inequality (4.2) is strict, and $`\sigma _x1`$ ($`1xL1`$), and $`\sigma _x2`$ if $`r_{x+1}=r_{x1}=1`$ or $`p1`$. We say the path $`P`$ has no particle in such case. When these conditions are violated, we observe the appearance of “particles” in $`P`$. We define the number of particles in $`P`$. Then, we construct a bijection from $`C_{L2m}^{(t1)}\times \pi _m`$ to the subset of paths in $`C_L^{(t)}`$ with $`m`$ particles. Here we denote by $`\pi _m`$ the set of partitions $`(\lambda _1,\mathrm{},\lambda _m)`$ of length $`m`$, i.e.,
$`\lambda _1\lambda _2\mathrm{}\lambda _m0.`$
We now give a precise definition of particles.
Given a path $`PC_L^{(t)}`$, we define an equivalence relation in the set $`\{1,\mathrm{},L1\}`$. We say a neighboring pair of integers $`x,x1\{1,\mathrm{},L1\}`$ is connected if and only if (4.3) is valid and
(4.4) $`\sigma _x+\sigma _{x1}`$ $`=`$ $`v^{(t)}(r)(2xL1).`$
We say $`xy`$, where $`x,y\{1,\mathrm{},L1\}`$, if and only if all neighboring pairs of integers in the interval between $`x`$ and $`y`$ are connected.
We call an equivalence class $`B`$ for this equivalence relation a block of particles if one of the following is satisfied:
(4.5) $`|B|2,`$
(4.6) $`B=\{x\}\text{ and }\sigma _x=1,(r_{x+1},r_x,r_{x1})=(1,2,1)\text{ or }(p1,p2,p1),`$
(4.7) $`B=\{x\}\text{ and }\sigma _x=0,r_{x+1}=r_{x1}.`$
We denote by $`(P)`$ the set of blocks for the path $`P`$. We call a block $`B`$ an isolated particle if $`B`$ consists of one element, $`|B|=1`$.
For each block $`B(P)`$, let $`\mathrm{max}(B),\mathrm{min}(B)`$ be the largest and the smallest integer in $`B`$. The blocks in $`P`$ are naturally ordered: $`B>B^{}`$ if and only if $`\mathrm{min}(B)>\mathrm{max}(B^{})`$.
If $`x`$ belongs to a block $`B`$, we have $`r_{x+1}=r_{x1}`$. We denote this number by $`r_x^{}`$.
###### Definition 4.1.
Let $`PC_L^{(t)}`$. We define a map $`m:(P)_0`$. Namely, for $`B(P)`$, we set
(4.8) $`m(B)=\{\begin{array}{cc}\left[\frac{|B|}{2}\right]\hfill & \text{ if }\sigma _x2\text{ or if }\sigma _x=1\text{ and }r_x^{}1,p1;\hfill \\ \left[\frac{|B|+1}{2}\right]\hfill & \text{ if }\sigma _x=0\text{ or if }\sigma _x=1\text{ and }r_x^{}=1,p1.\hfill \end{array}`$
where $`x=\mathrm{max}(B)`$. The number $`m(B)`$ is called the number of particles in the block $`B`$. We also set
(4.9) $`m(P)={\displaystyle \underset{B(P)}{}}m(B).`$
The number $`m(P)`$ is called the number of particles in the path $`P`$.
If $`|B|`$ is odd, then we have $`\sigma _{\mathrm{min}(B)}=\sigma _{\mathrm{max}(B)}`$. Therefore, one can use $`x=\mathrm{min}(B)`$ in Definition 4.1. Note also that if $`m(P)[L/2]`$.
### 4.2. Propagation of particles
Roughly speaking, blocks in a path $`P`$ are the location of particles in the sequence of integers
$`L1,L2,\mathrm{},2,1.`$
We move the particles in $`P`$ by changing blocks locally in this sequence without changing the total number of particles. We number the particles from the left to the right. We define the left (resp., right) move of the $`j`$-th particle, $`M_j^+`$ (resp., $`M_j^{}`$). These operations change a path $`PC_L^{(t)}`$ to a path $`M_j^\pm PC_L^{(t)}`$. For some paths $`PC_L^{(t)}`$, the path $`M_j^\pm PC_L^{(t)}`$ is not defined. If a block contains more than one particles, i.e., $`m(B)2`$, the left move is defined only for the leftmost particle in the block, and the right move is only for the rightmost one. In other words, particles in a block are located without gaps and they cannot pass each other. We start with the definition of the domain of the operations $`M_j^\pm `$.
We set
(4.10) $`D^+(P)`$ $`=`$ $`\{\mathrm{max}(B)|B(P)\},`$
(4.11) $`D^{}(P)`$ $`=`$ $`\{\begin{array}{cc}\{\mathrm{min}(B)|B(P)\}\backslash \{1\}\hfill & \text{ if }\sigma _0=\sigma _1=0;\hfill \\ \{\mathrm{min}(B)|B(P)\}\hfill & \text{ otherwise}.\hfill \end{array}`$
We define maps $`m_P^\pm :D^\pm (P)\{1,\mathrm{},m(P)\}`$:
(4.12) $`m_P^+(x)=1+{\displaystyle \underset{\genfrac{}{}{0pt}{}{B(P)}{\mathrm{max}(B)>x}}{}}m(B),`$
(4.13) $`m_P^{}(x)={\displaystyle \underset{\genfrac{}{}{0pt}{}{B(P)}{\mathrm{min}(B)x}}{}}m(B).`$
Set
(4.14) $`I^\pm (P)=\mathrm{Im}(m_P^\pm )\{1,\mathrm{},m(P)\}.`$
The map $`m_P^\pm `$ is injective, and therefore, on the image we can define the inverse mappings
(4.15) $`x_P^\pm =(m_P^\pm )^1:I^\pm (P)`$ $``$ $`\{1,\mathrm{},L1\},`$
(4.16) $`j`$ $``$ $`x_P^\pm (j).`$
The left move $`M_j^+P`$ is defined if and only if $`jI^+(P)`$, and the right move $`M_j^{}P`$ is defined if and only if $`jI^{}(P)`$. In such cases, we define the position of the $`j`$-th particle by $`x_P^+(j)`$ or $`x_P^{}(j)`$. If $`jI^+(P)I^{}(P)`$, then we have $`x_P^+(j)=x_P^{}(j)`$. In this case, the corresponding particle is isolated.
In general, for each $`1jm(P)`$ there exists a unique block $`B_j(P)`$ such that
(4.17) $`1+{\displaystyle \underset{\genfrac{}{}{0pt}{}{B(P)}{B>B_j}}{}}m(B)j{\displaystyle \underset{\genfrac{}{}{0pt}{}{B(P)}{BB_j}}{}}m(B).`$
We say that the $`j`$th particle is located in the block $`B_j`$. However, we do not specify its position except for the leftmost one or the rightmost one.
###### Definition 4.2.
Let $`PC_L^{(t)}`$. Recall our convention $`\sigma _L=\mathrm{}`$. Fix $`jI^\pm (P)`$ and set $`x=x_P^\pm (j)`$. We define
(4.18) $`M_j^\pm P=\left(\begin{array}{ccccc}r_L^\pm & r_{L1}^\pm & \mathrm{}& r_1^\pm & r_0^\pm \\ & \sigma _{L1}^\pm & \mathrm{}& \sigma _1^\pm & \sigma _0^\pm \end{array}\right),`$
by setting $`r_y^\pm =r_y`$ $`(0yL)`$ and $`\sigma _y^\pm =\sigma _y`$ $`(0yL1)`$ except for the following.
Case 1 $`\sigma _x0:`$
(4.19) $`\sigma _x^\pm `$ $`=`$ $`\sigma _x1,`$
(4.20) $`\sigma _{x1}^\pm `$ $`=`$ $`\sigma _{x1}+1.`$
Case 2 $`\sigma _x=0`$ and $`r_x^{}=1,p2:`$
(4.21) $`\sigma _x^\pm `$ $`=`$ $`1,`$
(4.22) $`\sigma _{x\pm 1}^\pm `$ $`=`$ $`\sigma _{x\pm 1}1.`$
Case 3 $`\sigma _x=0`$ and $`r_x^{}1,p2:`$
(4.23) $`r_x^\pm `$ $`=`$ $`2r_x^{}r_x,`$
(4.24) $`\sigma _{x\pm 1}^\pm `$ $`=`$ $`\{\begin{array}{cc}\sigma _{x\pm 1}v^{(t)}(r)1\hfill & \text{ if }r_{x\pm 2}=r_x;\hfill \\ \sigma _{x\pm 1}+v^{(t)}(r+\epsilon )\hfill & \text{ if }r_{x\pm 2}=r_x+2\epsilon ,\hfill \end{array}`$
(4.25) $`\sigma _{x1}^\pm `$ $`=`$ $`\{\begin{array}{cc}\sigma _{x1}v^{(t)}(r)\hfill & \text{ if }r_{x2}=r_x;\hfill \\ \sigma _{x1}+v^{(t)}(r+\epsilon )+1\hfill & \text{ if }r_{x2}=r_x+2\epsilon ,\hfill \end{array}`$
where $`r=\mathrm{min}\{r_x,r_x^{}\}`$ and $`\epsilon =\pm 1`$.
The moves $`M_i^\pm `$ change the degree of a path by $`\pm 1`$ and do not change the number of particles:
###### Lemma 4.3.
Let $`PC_{L,r}^{(t)}`$ and $`jI^\pm (P)`$. We have $`M_j^\pm PC_{L,r}^{(t)}`$, $`m(M_j^\pm P)=m(P)`$ and $`d(M_j^\pm P)=d(P)\pm 1`$.
The proof is only case-checking.
### 4.3. Properties of $`M_j^\pm `$
We list properties of moves $`M_j^\pm `$. In most cases, we omit proofs because they are only case-checkings. However, the following remark might help understanding.
If $`jI^+(P)`$ and the $`j`$-th particle belongs to a block $`B`$ with more than one particles, after the move $`M_j^+`$ this particle quits the block and the number of particles in $`B`$ decreases. The $`j`$-th particle either becomes an isolated particle or joins in another block to increase the number of particles in that block. The consideration of $`M_j^{}`$ is similar.
The moves $`M_j^+`$ and $`M_j^{}`$ are the inverse to each other:
###### Lemma 4.4.
Let $`jI^\pm (P)`$. Then, we have $`jI^{}(M_j^\pm P)`$ and $`M_j^{}M_j^\pm P=P`$.
The moves $`M_i^\epsilon `$ and $`M_j^\epsilon `$ are commutative as far as they are defined:
###### Lemma 4.5.
Suppose that $`ij+1`$. If $`M_i^\pm M_j^\pm P`$ is defined, then $`M_j^\pm M_i^\pm P`$ is also defined and $`M_j^\pm M_i^\pm P=M_i^\pm M_j^\pm P`$.
The moves $`M_i^+`$ and $`M_j^{}`$ are commutative as far as they are defined:
###### Lemma 4.6.
Suppose that $`ij`$. If $`M_i^\pm M_j^{}P`$ is defined, then $`M_j^{}M_i^\pm P`$ is also defined and $`M_j^{}M_i^\pm P=M_i^\pm M_j^{}P`$.
We are particularly interested in the relation between the moves $`M_{j+1}^+`$ and $`M_j^{}`$:
###### Lemma 4.7.
The move $`M_{j+1}^+P`$ is defined if and only if the move $`M_j^{}P`$ is defined.
The move $`M_{j+1}^+M_j^{}P`$ is defined if and only if the move $`M_j^{}M_{j+1}^+P`$ is defined, and in such a case we have $`M_{j+1}^+M_j^{}P=M_j^{}M_{j+1}^+P`$.
###### Corollary 4.8.
Let $`PC_L^{(t)}`$ and let $`l`$ be a positive integer. The move $`(M_{j+1}^+)^lP`$ is defined if and only if the move $`(M_j^{})^lP`$ is defined.
###### Proof.
We use induction on $`l`$. The case $`l=1`$ is proved in Lemma 4.7. Suppose that we have proved the statement for $`l1`$. By Lemma 4.7 $`(M_j^{})^lP=M_j^{}(M_j^{})^{l1}P`$ is defined if and only if $`M_{j+1}^+(M_j^{})^{l1}P`$ is defined. Again by Lemma 4.7 $`M_{j+1}^+(M_j^{})^{l1}P`$ is defined if and only if $`M_j^{l1}M_{j+1}^+P`$ is defined. By induction hypothesis $`(M_j^{})^{l1}M_{j+1}^+P`$ is defined if and only if $`(M_{j+1}^+)^{l1}M_{j+1}^+P`$ is defined. Thus, we have proved that $`(M_{j+1}^+)^lP`$ is defined if and only if the path $`(M_j^{})^lP`$ is defined. ∎
### 4.4. Rigging of particles
We define the rigging of particles $`\lambda _j(P)`$ $`(j=1,\mathrm{},m(P))`$. Let $`\pi _m`$ be the set of partitions of length $`m`$. The rigging $`(\lambda _1(P),\lambda _2(P),\mathrm{},\lambda _{m(P)}(P))`$ is by definition an element of $`\pi _{m(P)}`$.
Let $`PC_L^{(t)}`$. We set formally $`\lambda _{m(P)+1}(P)=0`$. Starting from $`j=m(P)`$ we define $`\lambda _j(P)`$ inductively by requiring that $`\lambda _j(P)\lambda _{j+1}(P)`$ is equal to the maximal number $`l`$ such that the path $`(M_j^{})^lP`$ is defined. Thus we obtain a partition. We call $`\lambda _j(P)`$ the rigging of the $`j`$-th particle in the path $`P`$.
Note that by Corollary 4.8 if $`jm(P)`$ the number $`\lambda _j(P)\lambda _{j+1}(P)`$ is also equal to the maximal number $`l`$ such that path $`(M_{j+1}^+)^lP`$ is defined. Therefore, using Lemma 4.4, we have
###### Proposition 4.9.
Let $`PC_L^{(t)}`$.
If $`M_j^+P`$ is defined if and only if $`\lambda _j(P)<\lambda _{j1}(P)`$ and in such case $`\lambda _i(M_j^+P)=\lambda _i(P)+\delta _{i,j}`$.
If $`M_j^{}P`$ is defined if $`\lambda _j(P)>\lambda _{j+1}(P)`$ and in such case $`\lambda _i(M_j^{}P)=\lambda _i(P)\delta _{i,j}`$.
Now we describe paths with zero rigging. Fix an integer $`m1`$. The following lemma follows from the definitions.
###### Lemma 4.10.
There is a bijection from the set of rigged paths in $`C_{L,r}^{(t)}`$ which have $`m`$ particles and zero rigging $`\lambda _1=\mathrm{}=\lambda _m=0`$ to the set of rigged paths in $`C_{L2m,r}^{(t)}`$ which have no particles. The bijection maps the path
$`P_{(m)}=\left(\begin{array}{cccccccccc}r_L& \mathrm{}& r_{2m+1}& 1& 2& \mathrm{}& 1& 2& 1& \mathrm{2\; 1}\\ & \mathrm{}& \sigma _{2m+1}& \sigma _{2m}& 0& \mathrm{}& v^{(t)}(1)& 0& v^{(t)}(1)& \mathrm{0\; 0}\end{array}\right)C_{L,r}^{(t)}`$
to the path
$`P_{(0)}=\left(\begin{array}{cccc}r_L& \mathrm{}& r_{2m+1}& 1\\ & \mathrm{}& \sigma _{2m+1}& \sigma _{2m}v^{(t)}(1)\end{array}\right)C_{L2m,r}^{(t)}.`$
Examples. We set $`v=v^{(t)}(1)`$.
$`L=2,m(P)=1`$:
$`\left(\begin{array}{ccc}1& 2& 1\\ & 0& 0\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccc}1& 2& 1\\ & 1& 0\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccc}1& 2& 1\\ & 0& 1\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccc}1& 2& 1\\ & 1& 1\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
$`L=3,m(P)=1,a0`$:
$`\left(\begin{array}{cccc}2& 1& 2& 1\\ & v+a& 0& 0\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{cccc}2& 1& 2& 1\\ & v+a1& 1& 0\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{cccc}2& 1& 2& 1\\ & v+a1& 0& 1\end{array}\right)`$
$`\stackrel{M_1^+}{}`$ $`\left(\begin{array}{cccc}2& 1& 2& 1\\ & v& 0& a\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{cccc}2& 1& 2& 1\\ & v1& 1& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}\stackrel{M_1^+}{}\left(\begin{array}{cccc}2& 1& 2& 1\\ & 0& v& a\end{array}\right)`$
$`\stackrel{M_1^+}{}`$ $`\left(\begin{array}{cccc}2& 3& 2& 1\\ & 0& 0& a\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{cccc}2& 1& 2& 1\\ & 0& v+1& a\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{cccc}2& 3& 2& 1\\ & 0& 1& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
$`L=4,m(P)=1,a0,2bv`$:
$`\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & b& v+a& 0& 0\end{array}\right)\stackrel{M_1^+}{}\mathrm{}\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & b& v& 0& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & b& vb& b& a\end{array}\right)`$
$`\stackrel{M_1^+}{}\mathrm{}\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 0& v& b& a\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 1& v& b& a\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 0& v+1& b& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
$`L=4,m(P)=1,a0,b>v`$:
$`\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & b& v+a& 0& 0\end{array}\right)\stackrel{M_1^+}{}\mathrm{}\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & b& v& 0& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & b& 0& v& a\end{array}\right)`$
$`\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 3& 2& 1\\ & bv1& 0& 0& a\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & b1& 0& v+1& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
$`\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & v& 0& b& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
$`L=4,m(P)=1,a,b0`$: Set $`v^{}=v^{(t)}(2)`$.
$`\left(\begin{array}{ccccc}3& 2& 1& 2& 1\\ & b& v+a& 0& 0\end{array}\right)\stackrel{M_1^+}{}\mathrm{}\stackrel{M_1^+}{}\left(\begin{array}{ccccc}3& 2& 1& 2& 1\\ & b& 0& v& a\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}3& 2& 3& 2& 1\\ & v^{}+b& 0& 0& a\end{array}\right)`$
$`\stackrel{M_1^+}{}\left(\begin{array}{ccccc}3& 2& 1& 2& 1\\ & b1& 0& v+1& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}\stackrel{M_1^+}{}\left(\begin{array}{ccccc}3& 2& 1& 2& 1\\ & 0& 0& v+b& a\end{array}\right)`$
$`\stackrel{M_1^+}{}\left(\begin{array}{ccccc}3& 2& 3& 2& 1\\ & v^{}& 0& b& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}\stackrel{M_1^+}{}\left(\begin{array}{ccccc}3& 2& 3& 2& 1\\ & 0& v^{}& b& a\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}3& 4& 3& 2& 1\\ & 0& 0& b& a\end{array}\right)`$
$`\stackrel{M_1^+}{}\left(\begin{array}{ccccc}3& 2& 3& 2& 1\\ & 0& v^{}+1& b& a\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}3& 4& 3& 2& 1\\ & 0& 1& b& a\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
$`L=4,m(P)=2`$:
$`\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 0& v& 0& 0\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 1& v& 0& 0\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 0& v+1& 0& 0\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 1& v+1& 0& 0\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
$`M_2^+M_2^+M_2^+`$
$`\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 1& v1& 1& 0\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 0& v& 1& 0\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 1& v& 1& 0\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
$`M_2^+M_2^+`$
$`\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 0& v& 0& 1\end{array}\right)\stackrel{M_1^+}{}\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 1& v& 0& 1\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
$`M_2^+`$
$`\left(\begin{array}{ccccc}1& 2& 1& 2& 1\\ & 1& v1& 1& 1\end{array}\right)\stackrel{M_1^+}{}\mathrm{}`$
### 4.5. Bijection
We are in a position to construct the bijection which leads to Proposition 3.6.
Let $`\overline{P}C_L^{(t1)}`$
$`\overline{P}=\left(\begin{array}{cccc}\overline{r}_L& \overline{r}_{L1}& \mathrm{}& \overline{r}_0\\ & \overline{\sigma }_{L1}& \mathrm{}& \overline{\sigma }_0\end{array}\right).`$
Define the numbers $`\sigma _0,\mathrm{},\sigma _{L1}`$ and the path $`P_{(0)}C_L^{(t)}`$ by the formulas
$`\sigma _0`$ $`=`$ $`\overline{\sigma }_0,`$
$`\sigma _i`$ $`=`$ $`\overline{\sigma }_i+1/2+w^{(t1)}(\overline{r}_{i+1},\overline{r}_i,\overline{r}_{i1})w^{(t)}(\overline{r}_{i+1},\overline{r}_i,\overline{r}_{i1}),(1iL1),`$
$`P_{(0)}`$ $`=`$ $`\left(\begin{array}{cccc}\overline{r}_L& \overline{r}_{L1}& \mathrm{}& \overline{r}_0\\ & \sigma _{L1}& \mathrm{}& \sigma _0\end{array}\right).`$
###### Lemma 4.11.
We have an inclusion
$`\iota _0:C_{L,r}^{(t1)}`$ $``$ $`C_{L,r}^{(t)},`$
$`\overline{P}`$ $``$ $`P_{(0)}.`$
The image of $`\iota _0`$ coincides with the subset of paths in $`C_{L,r}^{(t)}`$ with no particles.
###### Proof.
For $`i>1`$ by a simple computation we have
$`\sigma _i=\overline{\sigma }_i`$ if $`r_{i1}r_{i+1}=\pm 2`$,
$`\sigma _i=\overline{\sigma }_i+1`$ if $`r_{i1}r_{i+1}=0`$ and $`r_{i1}1,p1`$,
$`\sigma _i=\overline{\sigma }_i+2`$ if $`r_{i1}=r_{i+1}=1`$ or $`r_{i1}=r_{i+1}=p1`$.
The statement of the lemma follows straightforwardly. ∎
Let $`\overline{P}C_{L2m,r}^{(t1)}`$, and let $`\lambda =(\lambda _1,\mathrm{},\lambda _m)\pi _m`$ be a partition. The path $`P_{(0)}=\iota _0(\overline{P})C_{L2m,r}^{(t)}`$ is constructed in Lemma 4.11. The path $`P_{(m)}C_{L,r}^{(t)}`$ is constructed from $`P_{(0)}`$ in Lemma 4.10.
###### Lemma 4.12.
Notation being as above, for any non-negative integer $`m`$ we have an inclusion
$`\iota _m:C_{L2m,r}^{(t1)}\times \pi _m`$ $``$ $`C_{L,r}^{(t)},`$
$`(\overline{P},\lambda )`$ $``$ $`(M_m^+)^{\lambda _m}\mathrm{}(M_2^+)^{\lambda _2}(M_1^+)^{\lambda _1}P_{(m)}.`$
The image of $`\iota _m`$ coincides with the subset of paths in $`C_{L,r}^{(t)}`$ with $`m`$ particles. We have
(4.26) $`d(\iota _m(\overline{P},\lambda ))=d(\overline{P})+{\displaystyle \underset{i=1}{\overset{m}{}}}\lambda _i+L^2/4+L/2.`$
###### Proof.
The degree of a path is defined in (2.3). The relation (4.26) follows by a straightforward computation. The rest follows from Lemma 4.9. ∎
We obtain the main result of this section.
###### Theorem 4.13.
The map $`\iota :=_{m=0}^{\mathrm{}}\iota _m`$ defines a bijection
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}(C_{L2m,r}^{(t1)}\times \pi _m)C_{L,r}^{(t)}`$
with the property $`(\text{4.26})`$.
Proposition 3.3 is a direct consequence of Theorem 4.13.
## 5. The case $`p=3`$
The aim of this section is to deduce Conjecture 2.3 for $`p=3`$ from the work \[FJM\].
Throughout this section, we fix an integer $`p^{}>3`$ coprime to $`3`$. Omitting the upper index we write the $`(2,1)`$-field $`\varphi ^{(r^{},r)}(z)`$ as $`\varphi (z)`$, since for $`p=3`$ the choice of $`r^{}=r\pm 1`$ is uniquely determined from $`r`$. We set
$`M^{(3,p^{})}=M_{1,1}^{(3,p^{})}M_{2,1}^{(3,p^{})}.`$
In the following we deal with bi-graded $``$-vector spaces of the form $`X=_{d,L}X_{d,L}`$ with $`dimX_{d,L}<\mathrm{}`$. We will refer to the index $`d`$ and $`L`$ as degree and weight, respectively. We set $`X_L=_dX_{d,L}`$ and define the restricted dual space by $`X_L^{}=_dHom_{}(X_{d,L},)`$.
### 5.1. Extended modules and monomial basis
First we review the results of \[FJM\] which are relevant to us.
Consider the Heisenberg algebra with generators $`\{h_n\}_n`$ satisfying $`[h_m,h_n]=m\delta _{m+n,0}`$. Denote by
$`_\gamma =[h_1,h_2,\mathrm{}]|\gamma `$
the Fock space with highest weight vector $`|\gamma `$ ($`\gamma `$), where $`h_n|\gamma =0`$ ($`n>0`$) and $`h_0|\gamma =\gamma |\gamma `$. We set
$`\beta =\sqrt{{\displaystyle \frac{p^{}2}{2}}}.`$
We use the vertex operator which acts on $`=_n_{n\beta }`$,
$`\mathrm{\Phi }_\beta (z)=:\mathrm{exp}(\beta {\displaystyle \underset{n0}{}}{\displaystyle \frac{h_n}{n}}z^n)e^{\beta Q}z^{\beta h_0}:.`$
Here $`::`$ stands for the normal ordering symbol, and $`e^{\beta Q}:_\gamma \stackrel{}{}_{\gamma +\beta }`$ is an isomorphism of vector spaces such that $`[h_n,e^{\beta Q}]=0`$ ($`n0`$), $`e^{\beta Q}|\gamma =|\gamma +\beta `$. Set
$`V^{(3,p^{})}`$ $`=`$ $`_LV_L^{(3,p^{})}\text{ where }V_L^{(3,p^{})}=\{\begin{array}{cc}M_{1,1}^{(3,p^{})}_{L\beta }\hfill & \text{ if }L\text{ is even};\hfill \\ M_{2,1}^{(3,p^{})}_{L\beta }\hfill & \text{ if }L\text{ is odd},\hfill \end{array}`$
$`|\mathrm{vac}`$ $`=`$ $`|1,1|0V_0^{(3,p^{})}.`$
Introduce a field with coefficients in $`End(V^{(3,p^{})})`$,
$`a(z)=\varphi (z)\mathrm{\Phi }_\beta (z).`$
The parameter $`\beta `$ is so chosen that $`a(z)`$ has the expansion
$`a(z)={\displaystyle \underset{n}{}}a_nz^{n1},`$
and that the coefficients $`a_n`$’s are mutually commutative,
(5.1) $`[a_m,a_n]=0(m,n).`$
We have also $`a_n|\mathrm{vac}=0`$ for $`n0`$. Later we will use the relation
(5.2) $`a(z)^2=k\mathrm{id}:\mathrm{\Phi }_\beta (z)^2:,`$
where $`k`$ is a nonzero constant. We have
(5.3) $`[h_n,a(z)]=\beta z^na(z).`$
Let $`A`$ be a subalgebra of $`\mathrm{End}(V^{(3,p^{})})`$ generated by $`a_n`$’s. Our main concern is the following subspace of $`V^{(3,p^{})}`$ generated from $`|\mathrm{vac}`$ by acting with $`a_n`$’s:
$`W^{(3,p^{})}=A|\mathrm{vac}V^{(3,p^{})}.`$
We introduce a bi-grading to $`W^{(3,p^{})}`$ by assigning the degree and weight as $`\mathrm{deg}a_n=n`$, $`\mathrm{wt}a_n=1`$ and $`\mathrm{deg}|\mathrm{vac}=0`$, $`\mathrm{wt}|\mathrm{vac}=0`$.
Let $`\mathrm{\Lambda }_L=[x_1,\mathrm{},x_L]^{𝔖_L}`$ denote the space of symmetric polynomials in the variables $`x_1,\mathrm{},x_L`$. The restricted dual space $`W_L^{(3,p^{})}`$ is identified with a subspace of $`\mathrm{\Lambda }_L`$ by
(5.4) $`W_L^{(3,p^{})}\mathrm{\Lambda }_L,v|v|a(x_L)\mathrm{}a(x_1)|\mathrm{vac}.`$
If $`n>0`$, $`h_n|\mathrm{vac}=0`$ and from (5.3) we see that $`W_L^{(3,p^{})}`$ is invariant by $`h_n`$. The right action of $`h_n`$ ($`n>0`$) on $`W_L^{(3,p^{})}`$ corresponds to the multiplication by $`\beta _{j=1}^Lx_j^n`$ on $`\mathrm{\Lambda }_L`$. Therefore, the image $`I_L^{(3,p^{})}`$ of (5.4) is an ideal of $`\mathrm{\Lambda }_L`$.
The linear isomorphism $`\tau =\mathrm{id}e^{2\beta Q}:V^{(3,p^{})}V^{(3,p^{})}`$ has an effect of shifting the indices
(5.5) $`\tau |\mathrm{vac}_{2l}=|\mathrm{vac}_{2l+2},\tau a_n\tau ^1=a_{np^{}+2},`$
where $`|\mathrm{vac}_{2l}=|1,1|2l\beta `$. We have an increasing filtration (\[FJM\], Proposition 2.3)
(5.6) $`W^{(3,p^{})}\tau ^1(W^{(3,p^{})})\tau ^2(W^{(3,p^{})})\mathrm{},`$
(5.7) $`V^{(3,p^{})}={\displaystyle \underset{N0}{}}\tau ^N(W^{(3,p^{})}).`$
Let $`\widehat{\mathrm{\Lambda }}_L=[x_1^{\pm 1},\mathrm{},x_L^{\pm 1}]^{𝔖_L}`$ be the space of symmetric Laurent polynomials, and let $`\widehat{I}_L^{(3,p^{})}`$ be the ideal of $`\widehat{\mathrm{\Lambda }}_L`$ generated by $`I_L^{(3,p^{})}`$. By (5.1), (5.5) and (5.7), we see that the subspace of $`\widehat{\mathrm{\Lambda }}_L`$ spanned by all matrix elements
(5.8) $`v|a(x_L)\mathrm{}a(x_1)|u(v|(V^{(3,p^{})})^{},|uV^{(3,p^{})})`$
coincides with $`\widehat{I}_L^{(3,p^{})}`$. Note also that $`f\widehat{I}_L^{(3,p^{})}`$ if and only if $`x_1\mathrm{}x_Lf\widehat{I}_L^{(3,p^{})}`$.
The following fact is proved in \[FJM\], Theorem 2.8.
###### Proposition 5.1.
For each $`L0`$, the space $`W_L^{(3,p^{})}`$ has a basis consisting of elements
$`a_{\lambda _L}\mathrm{}a_{\lambda _1}|\mathrm{vac},`$
where $`\lambda =(\lambda _L,\mathrm{},\lambda _1)`$ runs over all $`L`$-tuples of positive integers satisfying
(5.9) $`\lambda _L\mathrm{}\lambda _1,`$
(5.10) $`\lambda _{i+2}\lambda _ip^{}2(1iL2).`$
Let us reformulate Proposition 5.1 in a form suitable for later use.
###### Proposition 5.2.
For any $`\lambda =(\lambda _L,\mathrm{},\lambda _1)^L`$ with $`\lambda _L\mathrm{}\lambda _1`$, there exist unique $`c_{\lambda \mu }`$ such that the following identity holds as operators on $`V^{(3,p^{})}:`$
(5.11) $`a_{\lambda _L}\mathrm{}a_{\lambda _1}={\displaystyle \underset{\mu }{}}c_{\lambda \mu }a_{\mu _L}\mathrm{}a_{\mu _1},`$
where the sum in the right hand side is taken over all $`\mu =(\mu _L,\mathrm{},\mu _1)^L`$ satisfying (5.9), (5.10) and $`_{i=1}^L\mu _i=_{i=1}^L\lambda _i`$. The coefficients $`c_{\lambda ,\mu }`$ are shift invariant $`:`$
$`c_{\lambda +(1,\mathrm{},1),\mu +(1,\mathrm{},1)}=c_{\lambda ,\mu }.`$
###### Proof.
The sum in the right hand side of (5.11) is a well defined element of $`\mathrm{End}(V^{(3,p^{})})`$ because for a given vector in $`V^{(3,p^{})}`$ only finitely many monomials $`a_{\mu _L}\mathrm{}a_{\mu _1}`$ acts non-trivially. For $`N`$ and $`\mu =(\mu _L,\mathrm{},\mu _1)^L`$, let us write $`\mu ^{(N)}=(\mu _L^{(N)},\mathrm{},\mu _1^{(N)})`$, $`\mu _i^{(N)}=\mu _i+N(p^{}2)`$. Choosing $`N`$ so that $`\lambda _1^{(N)}>0`$, we apply Proposition 5.1 to $`\lambda ^{(N)}`$. There exist unique constants $`c_{\lambda \mu }^{(N)}`$ such that
$`\left(a_{\lambda _L^{(N)}}\mathrm{}a_{\lambda _1^{(N)}}{\displaystyle \underset{\mu }{}}c_{\lambda \mu }^{(N)}a_{\mu _L^{(N)}}\mathrm{}a_{\mu _1^{(N)}}\right)|\mathrm{vac}=0,`$
where the sum is taken over $`\mu `$ satisfying (5.9), (5.10) and $`\mu _1>N(p^{}2)`$. Since $`a_n`$’s commute and $`W^{(3,p^{})}`$ is cyclic over $`[a_1,a_2,\mathrm{}]`$, the operator inside the parentheses vanishes on $`W^{(3,p^{})}`$. Applying $`\tau ^N`$ on both sides, we deduce that
$`a_{\lambda _L}\mathrm{}a_{\lambda _1}={\displaystyle \underset{\mu }{}}c_{\lambda \mu }^{(N)}a_{\mu _L}\mathrm{}a_{\mu _1}\text{on }\tau ^N(W^{(3,p^{})}).`$
If $`N^{}>N`$, then $`\tau ^N(W^{(3,p^{})})\tau ^N^{}(W^{(3,p^{})})`$. Therefore we have
$`{\displaystyle \underset{\mu }{}}(c_{\lambda \mu }^{(N)}c_{\lambda \mu }^{(N^{})})a_{\mu _L}\mathrm{}a_{\mu _1}=0\text{on }\tau ^N(W^{(3,p^{})}),`$
or equivalently
$`{\displaystyle \underset{\mu }{}}(c_{\lambda \mu }^{(N)}c_{\lambda \mu }^{(N^{})})a_{\mu _LN(p^{}2)}\mathrm{}a_{\mu _1N(p^{}2)}=0\text{on }W^{(3,p^{})}.`$
From Proposition 5.1 we have $`c_{\lambda \mu }^{(N^{})}=c_{\lambda \mu }^{(N)}`$ for all $`\mu `$ with $`\mu _1>N(p^{}2)`$. Setting $`c_{\lambda \mu }=lim_N\mathrm{}c_{\lambda \mu }^{(N)}`$ and noting (5.7), we obtain (5.11). The uniqueness of $`c_{\lambda \mu }`$ is clear.
The last statement follows from the commutation relation (5.3). ∎
###### Corollary 5.3.
A symmetric Laurent polynomial $`f=_{\lambda _1,\mathrm{},\lambda _L}f_{\lambda _L,\mathrm{},\lambda _1}x_L^{\lambda _L}\mathrm{}x_1^{\lambda _1}`$ belongs to $`\widehat{I}_L^{(3,p^{})}`$ if and only if
$`f_{\lambda _L,\mathrm{},\lambda _1}={\displaystyle \underset{\mu }{}}c_{\lambda \mu }f_{\mu _L,\mathrm{},\mu _1}`$
holds for all $`\lambda =(\lambda _L,\mathrm{},\lambda _1)^L`$ with $`\lambda _L\mathrm{}\lambda _1`$. In particular, if $`f_{\lambda _L,\mathrm{},\lambda _1}=0`$ holds for all $`\lambda `$ satisfying (5.9), (5.10), then $`f=0`$.
We quote one more fact which follows immediately from \[FJM\], Propositions 4.2 and 4.4.
###### Proposition 5.4.
Let $`\widehat{D}_L`$ be the subspace of $`\widehat{\mathrm{\Lambda }}_L`$ consisting of Laurent polynomials vanishing on the diagonal $`x_ix_j=0`$ $`(ij)`$. Then
$`\widehat{I}_L^{(3,p^{})}\widehat{D}_L=\widehat{I}_L^{(3,p^{}3)}\times {\displaystyle \underset{1i<jL}{}}(x_ix_j)^2.`$
When $`p^{}=4,5`$, we define $`\widehat{I}_L^{(3,p^{}3)}`$ in the right hand side to be $`\widehat{\mathrm{\Lambda }}_L`$.
### 5.2. Monomial basis for $`M^{(3,p^{})}`$
For $`p=3`$, the admissibility of the rigged paths is equivalent to the following conditions for the $`n_i`$’s (see Section 2.3).
(5.12) $`n_1_0+\mathrm{\Delta }_{2,1},`$
(5.13) $`n_{i+1}n_i{\displaystyle \frac{p^{}}{2}}+3(1iL1),`$
(5.14) $`n_{i+2}n_i1(1iL2).`$
Our goal in this section is to prove the following result.
###### Theorem 5.5.
The set of monomials
$`\varphi _{n_L}\mathrm{}\varphi _{n_1}|1,1`$
satisfying the conditions (5.12), (5.13) and (5.14) is a basis of $`M^{(3,p^{})}`$.
We have shown in Theorem 2.2 that the character of this set matches with that of $`M^{(3,p^{})}`$. For the proof of Theorem 5.5, it is therefore sufficient to show that this set spans $`M^{(3,p^{})}`$. For that purpose, we consider a filtration of $`M^{(3,p^{})}`$
$`\{0\}=F_1F_0\mathrm{}F_i\mathrm{}M^{(3,p^{})},`$
defined by
(5.15) $`F_0=|1,1,`$
(5.16) $`F_i=F_{i1}+{\displaystyle \underset{n+(1)^{i1}\mathrm{\Delta }_{2,1}}{}}\varphi _nF_{i1},`$
where $`\mathrm{\Delta }_{2,1}=(p^{}2)/4`$. In the right hand side of (5.16), we have
$`\varphi _n=\{\begin{array}{cc}\varphi _n^{(2,1)}\hfill & \text{if }i\text{ is odd};\hfill \\ \varphi _n^{(1,2)}\hfill & \text{if }i\text{ is even}.\hfill \end{array}`$
If $`rr^{\prime \prime }`$ the action of $`\varphi _n^{(r^{},r)}`$ is zero on $`M_{r^{\prime \prime },1}^{(3,p^{})}`$ by definition. Since the operator product expansion of $`\varphi (z)`$ contains the energy-momentum tensor $`T(z)`$, we have $`_{i0}F_i=M^{(3,p^{})}`$.
Let
$`\stackrel{~}{\varphi }_n:F_i/F_{i1}F_{i+1}/F_i`$
denote the operator induced from $`\varphi _n`$ on the associated graded space $`\stackrel{~}{M}^{(3,p^{})}=_{i0}F_i/F_{i1}`$. From the remark made above, the proof of Theorem 5.5 reduces to the following statement.
###### Lemma 5.6.
Let $`\stackrel{~}{B}_L^{(3,p^{})}`$ denote the set of vectors
(5.17) $`\stackrel{~}{\varphi }_{n_L}\mathrm{}\stackrel{~}{\varphi }_{n_1}|1,1(n_i(1)^i\mathrm{\Delta }_{2,1})`$
satisfying the conditions (5.12), (5.13) and (5.14). Then $`\stackrel{~}{M}_L^{(3,p^{})}=\mathrm{span}_{}\stackrel{~}{B}_L^{(3,p^{})}`$.
The rest of this subsection is devoted to the proof of Lemma 5.6.
We set
$`\stackrel{~}{\varphi }(z)={\displaystyle \underset{n+(1)^i\mathrm{\Delta }_{2,1}}{}}\stackrel{~}{\varphi }_nz^{n\mathrm{\Delta }_{2,1}}\text{on }F_i/F_{i1}.`$
The field $`\stackrel{~}{\varphi }(z)`$ satisfies a quadratic relation following from that of $`\varphi (z)`$ (see e.g.\[FJMMT\]),
$`{\displaystyle \underset{k0}{}}c_k\stackrel{~}{\varphi }_{n_2k}\stackrel{~}{\varphi }_{n_1+k}+{\displaystyle \underset{k0}{}}c_k\stackrel{~}{\varphi }_{n_1+p^{}/22k}\stackrel{~}{\varphi }_{n_2p^{}/2+2+k}=0,`$
where $`_{k0}c_kz^k=(1z)^{p^{}/22}`$. This can be rewritten as a relation of the form
(5.18) $`\stackrel{~}{\varphi }_{n_2}\stackrel{~}{\varphi }_{n_1}={\displaystyle \underset{l1}{}}C_{n_1,n_2,l}\stackrel{~}{\varphi }_{n_2l}\stackrel{~}{\varphi }_{n_1+l}(n_2n_1{\displaystyle \frac{p^{}}{2}}+2)`$
with some $`C_{n_1,n_2,l}`$. Given any monomial (5.17), we apply (5.18) to rewrite it as a linear combination of monomials satisfying the condition (5.13). In each step of rewriting, $`_{i=1}^Ln_i`$ is invariant and $`_{i=1}^Lin_i`$ strictly increases. Since $`_{j=1}^in_j>0`$ for each $`i`$, $`_{i=1}^Lin_i`$ is bounded from above. Therefore, the process terminates after a finite number of steps. To prove Lemma 5.6, we must also fulfill the second condition (5.14). To this end, we make use of the knowledge in the previous subsection.
Consider the filtration of $`V^{(3,p^{})}`$ induced from that of $`M^{(3,p^{})}`$.
$`V_i^{(3,p^{})}`$ $`=`$ $`_L(V_L^{(3,p^{})})_i,`$
$`(V_L^{(3,p^{})})_i`$ $`=`$ $`\{\begin{array}{cc}(F_iM_{1,1}^{(3,p^{})})_{L\beta }\hfill & \text{ if }L\text{ is even};\hfill \\ (F_iM_{2,1}^{(3,p^{})})_{L\beta }\hfill & \text{ if }L\text{ is odd}.\hfill \end{array}`$
Define
$`\stackrel{~}{a}_n:V_i^{(3,p^{})}/V_{i1}^{(3,p^{})}V_{i+1}^{(3,p^{})}/V_i^{(3,p^{})}`$
to be the operator induced from $`a_n:V_i^{(3,p^{})}V_{i+1}^{(3,p^{})}`$. Setting $`\stackrel{~}{a}(z)=_n\stackrel{~}{a}_nz^{n1}`$, we have clearly the relation
$`\stackrel{~}{a}(z)=\stackrel{~}{\varphi }(z)\mathrm{\Phi }_\beta (z).`$
Since the relations (5.11) are homogeneous, they remain valid for $`\stackrel{~}{a}_n`$’s. From the relation (5.2), we have also $`\stackrel{~}{a}(z)^2=0`$.
Now for $`\stackrel{~}{v}|\stackrel{~}{M}^{(3,p^{})}`$ and $`|\stackrel{~}{u}\stackrel{~}{M}^{(3,p^{})}`$, we consider a matrix element
$`g(x_1,\mathrm{},x_L)=\stackrel{~}{v}|\stackrel{~}{\varphi }(x_L)\mathrm{}\stackrel{~}{\varphi }(x_1)|\stackrel{~}{u}`$
which is a formal series in $`x_1,\mathrm{},x_L`$.
###### Lemma 5.7.
There exists an element $`f\widehat{I}_L^{(3,p^{}3)}`$ such that
$`g(x_1,\mathrm{},x_L)=f(x_1,\mathrm{},x_L)\times {\displaystyle \underset{1i<jL}{}}(x_jx_i)^{\beta ^2+2}.`$
Here the right hand side means its expansion in the domain $`|x_L|>\mathrm{}>|x_1|`$.
###### Proof.
Set $`\stackrel{~}{v}_L|=\stackrel{~}{v}|L\beta |`$, $`|\stackrel{~}{u}_0=|\stackrel{~}{u}|0`$, and consider
$`F(x_1,\mathrm{},x_L)`$ $`=`$ $`\stackrel{~}{v}_L|\stackrel{~}{a}(x_L)\mathrm{}\stackrel{~}{a}(x_1)|\stackrel{~}{u}_0`$
$`=`$ $`g(x_1,\mathrm{},x_L)\times {\displaystyle \underset{1i<jL}{}}(x_jx_i)^{\beta ^2}.`$
We have $`F\widehat{I}_L^{(3,p^{})}`$ because the defining relations (5.11) for the operator algebra dual to the ideal $`F\widehat{I}_L^{(3,p^{})}`$ are equally valid for $`\stackrel{~}{a}(z)`$. Moreover the relation $`\stackrel{~}{a}(z)^2=0`$ implies $`F\widehat{D}_L`$. Hence from Proposition 5.4, we can write
$`F(x_1,\mathrm{},x_L)=f(x_1,\mathrm{},x_L){\displaystyle \underset{1i<jL}{}}(x_jx_i)^2`$
with some $`f\widehat{I}_L^{(3,p^{}3)}`$. This proves the assertion. ∎
###### Lemma 5.8.
Let $`|\stackrel{~}{u}\stackrel{~}{M}^{(3,p^{})}`$ be a homogeneous element and let $`m=(m_3,m_2,m_1)`$ be a triple of integers such that the relations (5.13), (5.14) for $`L=3`$ are not satisfied. Then we have a relation
(5.19) $`\stackrel{~}{\varphi }_{m_3}\stackrel{~}{\varphi }_{m_2}\stackrel{~}{\varphi }_{m_1}|\stackrel{~}{u}={\displaystyle \underset{n}{}}d_{m,n}\stackrel{~}{\varphi }_{n_3}\stackrel{~}{\varphi }_{n_2}\stackrel{~}{\varphi }_{n_1}|\stackrel{~}{u}`$
with some $`d_{m,n}`$, where the sum is taken over $`n=(n_3,n_2,n_1)`$ satisfying the relations (5.13), (5.14) and $`_{i=1}^3n_i=_{i=1}^3m_i`$, $`_{i=1}^3in_i>_{i=1}^3im_i`$.
###### Proof.
Given $`|\stackrel{~}{u}`$, let $`X`$ be the linear span of all monomials $`\stackrel{~}{\varphi }_{m_3}\stackrel{~}{\varphi }_{m_2}\stackrel{~}{\varphi }_{m_1}|\stackrel{~}{u}`$, and let $`X_0`$ be the subspace spanned by those whose indices $`m`$ satisfy (5.13) and (5.14). We show $`X=X_0`$.
Let $`\stackrel{~}{v}|\stackrel{~}{M}^{(3,p^{})}`$ be an element orthogonal to the space $`X_0`$. Consider the matrix element $`g(x_1,x_2,x_3)=\stackrel{~}{v}|\stackrel{~}{\varphi }(x_3)\stackrel{~}{\varphi }(x_2)\stackrel{~}{\varphi }(x_1)|\stackrel{~}{u}`$. Let $`f`$ be as in Lemma 5.7 with $`L=3`$, and expand them as
$`f(x_1,x_2,x_3)={\displaystyle \underset{\lambda _1,\lambda _2,\lambda _3}{}}f_{\lambda _3,\lambda _2,\lambda _1}x_3^{\lambda _3}x_2^{\lambda _2}x_1^{\lambda _1},`$
$`g(x_1,x_2,x_3)={\displaystyle \underset{n_1,n_2,n_3}{}}g_{n_3,n_2,n_1}x_3^{n_3\mathrm{\Delta }_{2,1}}x_2^{n_2\mathrm{\Delta }_{2,1}}x_1^{n_1\mathrm{\Delta }_{2,1}},`$
where $`g_{n_3,n_2,n_1}=\stackrel{~}{v}|\stackrel{~}{\varphi }_{n_3}\stackrel{~}{\varphi }_{n_2}\stackrel{~}{\varphi }_{n_1}|\stackrel{~}{u}`$.
Suppose $`\lambda _3\lambda _2\lambda _1`$ and $`\lambda _3\lambda _1p^{}5`$. The coefficient $`f_{\lambda _3,\lambda _2,\lambda _1}`$ is a linear combination of $`g_{n_3,n_2,n_1}`$ such that
$`\lambda _i=n_i\mathrm{\Delta }_{2,1}+(i1)(\beta ^22)\alpha _{i1}+\alpha _i,`$
with some $`\alpha _1,\alpha _20`$ and $`\alpha _0=\alpha _3=0`$. In particular
$`n_3n_1=\lambda _3\lambda _1(p^{}6)+\alpha _1+\alpha _21.`$
Using the quadratic relation (5.18), we can further rewrite $`g_{n_3,n_2,n_1}`$ as linear combinations of those satisfying (5.13). Since $`n_3n_1=_{i=1}^3in_i2_{i=1}^3n_i`$ does not decrease, the resulting terms all satisfy (5.14) as well. By the choice of $`\stackrel{~}{v}|`$, $`g_{n_3,n_2,n_1}=0`$ holds for all such $`n=(n_3,n_2,n_1)`$. It follows that $`f_{\lambda _3,\lambda _2,\lambda _1}=0`$ holds for all $`\lambda =(\lambda _3,\lambda _2,\lambda _1)`$ with $`\lambda _3\lambda _2\lambda _1`$ and $`\lambda _3\lambda _1p^{}5`$. Applying Corollary 5.3 we conclude that $`f=0=g`$. This shows $`X=X_0`$.
It remains to verify that if $`m`$ violates (5.13) or (5.14), then $`m_3m_1<n_3n_1`$ for all terms in the right hand side of (5.19). This is evident if (5.14) is violated. Otherwise we can use the quadratic relation (5.18) to rewrite it as a linear combination of those satisfying (5.13). In the process $`m_3m_1`$ strictly increases, and the verification reduces to the first case. ∎
Proof of Lemma 5.6. Given a monomial (5.17), suppose the conditions (5.13)–(5.14) are not valid for a triple $`(m_{i+1},m_i,m_{i1})`$. We apply Lemma 5.8 to reduce it. In the process, $`_{i=1}^Ln_i`$ does not change and $`_{i=1}^Lin_i`$ strictly increases. Therefore the process terminates after a finite number of steps. Proof of Lemma 5.6 is now complete. ∎
Acknowledgments. BF is partially supported by grants RFBR-02-01-01015, RFHR-01-01-00906, INTAS-00-00055. MJ is partially supported by the Grant-in-Aid for Scientific Research B2, no. 16340033, and TM is partially supported by (A1) no.13304010, Japan Society for the Promotion of Science. EM is partially supported by the National Science Foundation (NSF) grant DMS-0140460. YT is partially supported by Grant-in-Aid for Young Scientists (B) No. 17740089.
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# BOSON GAUGE COUPLINGS AT LEP
## 1 Introduction
The LEP $`e^+e^{}`$collider has been the most important experimental environment to study the electroweak interactions and precisely test the Standard Model. Since 1989 each of the four LEP experiments collected data corresponding to an integrated luminosity larger than 1000 pb$`{}_{}{}^{}1`$, of which about 200 pb$`{}_{}{}^{}1`$ at a centre-of-mass energy of the $`Z`$mass peak (from 1989 to 1995). After 1996 a second phase started, LEP2, with collected lumonosity of 750 pb$`{}_{}{}^{}1`$ per experiment and energy above the WW production threshold (centre-of-mass energy above 161 GeV), allowing to produce for the first time two real (on shell) massive bosons WW, ZZ, decaying in a four fermions final state. These new channels opened at LEP2, allow the precise measurement of the mass of the W but are also sensitive to the trilinear gauge couplings (TGC) predicted in the Standard Model $`^\mathrm{?}`$. LEP ended collecting data in 2000 but only recently all the most important analyses in the W sector have been completed and the results of the four LEP experiments combined. This proceeding is a review of the measurements of the boson gauge couplings at LEP. All the results here presented have been published by ALEPH $`^\mathrm{?}`$, L3 $`^\mathrm{?}`$, OPAL $`^\mathrm{?}`$ and DELPHI $`^\mathrm{?}`$ collaborations.
## 2 Triple gauge boson couplings
The triple gauge boson couplings (TGC) can be classified into charged couplings (cTGC) when the triple boson vertex is WW$`\gamma `$ or WWZ and neutral couplings (nTGC) for the case of ZZZ, ZZ$`\gamma `$ or Z$`\gamma `$$`\gamma `$. The first are forseen in the Standard Model and are linked to the non-abelian structure of the $`\mathrm{SU}(2)_\mathrm{L}\times \mathrm{U}(1)_\mathrm{Y}`$ gauge simmetry. The neutral nTGC at the contrary, do not exist in the Standard Model at tree level.
### 2.1 Charged triple gauge couplings
In order to measure the charged triple gauge couplings (cTGC) at LEP is necessary to have a vertex with two W and a neutral boson (Z,$`\gamma `$). Within the Standard Model, this vertex can be achieved through the production of two massive W with the $`e^+e^{}W^+W^{}`$ process available for the first time at LEP2 or, with lower sensitivity, through the single W production with the process $`e^+e^{}W^{}e^\pm \nu (\overline{\nu })`$ and through the single photon production $`e^+e^{}\gamma \overline{\nu _e}\nu _e`$ .
The first proof of the existence of the charged TGC is in the measurement of the WW cross section. Three diagrams are responsible for the WW production as shown in fig. 1: the t-channel with the neutrino exchange (fig. 1,a), the s-channel with a Z exchange (fig. 1,b) and the s-channel with a virtual photon exchange (fig. 1,c). Both the s-channel diagrams have a TGC vertex.
All the three diagrams are needed to make the WW cross section converging to finite values at high energies. If the $`\mathrm{SU}(2)_\mathrm{L}\times \mathrm{U}(1)_\mathrm{Y}`$ was an abelian gauge simmetry, the TGC verteces would be forbidden and the WW cross section would diverge when increasing the centre-of-mass energy. The LEP2 data fit precisely the Standard Model predictions and confirm the presence of the TGCs and the non-abelian structure of the $`\mathrm{SU}(2)_\mathrm{L}\times \mathrm{U}(1)_\mathrm{Y}`$ gauge simmetry, as can be viewed in fig. 2.
The most general form for an effective charged TGC Lagrangian consistent with Lorentz invariance involves 14 complex couplings, 7 for the WW$`\gamma `$ and 7 for the WWZ vertex. Most of these couplings are C- or P-violating while in the Standard Model C- and P-conservation are predicted in the TGCs sector. Assuming C and P conservation the 14 complex couplings are reduced to 6 real couplings: $`g_1^\gamma `$, $`g_1^Z`$, $`k_\gamma `$, $`k_Z`$, $`\lambda _\gamma `$ and $`\lambda _Z`$. In the Standard model $`g_1^\gamma =1`$, $`g_1^Z=1`$, $`k_\gamma =1`$, $`k_Z=1`$, $`\lambda _\gamma =0`$ and $`\lambda _Z=0`$. The couplings can be related to physical properties of the gauge bosons, like the electric dipole, quadrupole and magnetic moment. For instance the W anomalous magnetic moment can be written as
$$\mu _W=\frac{e}{2M_W}(1+k_\gamma +\lambda _\gamma )$$
(1)
The requirement of local $`\mathrm{SU}(2)_\mathrm{L}\times \mathrm{U}(1)_\mathrm{Y}`$ gauge invariance introduce the further constraints
$`\mathrm{\Delta }k_Z`$ $`=`$ $`\mathrm{\Delta }k_\gamma \mathrm{tan}^2\theta _W+\mathrm{\Delta }g_1^Z`$ (2)
$`\lambda _\gamma `$ $`=`$ $`\lambda _Z`$
$`g_1^\gamma `$ $`=`$ $`1`$
with $`\mathrm{\Delta }`$ indicating the deviation from the Standard Model predictions and $`\theta _W`$ the electroweak mixing angle, leaving 3 independent real couplings $`^\mathrm{?}`$: $`g_1^Z`$, $`k_\gamma `$ and $`\lambda _\gamma `$.
These couplings have been experimentally tested with the $`e^+e^{}W^+W^{}`$ sample collected at LEP2. The angular distribution of the $`e^+e^{}W^+W^{}`$ cross section are more sensitive to TGCs than the inclusive measurement. A precise study of the TGC parameters is achieved with the analysis of the differential WW cross sections respect to 5 angles:
* the angle $`\theta _W`$ between the $`W^{}`$ and the initial $`e^{}`$ in the WW rest frame
* the polar and azimuthal angles of the fermion in the decay $`W^{}f\overline{f}`$ calculated in the rest frame of the $`W^{}`$
* the corresponding polar and azimuthal angles of the fermion in the decay of the $`W^+`$.
All the possible WW decays, shown in tab. 1 together with the branching ratios and the typical efficiencies, are taken into account for this analysis.
In the semileptonic $`W^+W^{}qql\nu `$ decays the charge of the lepton identifies without ambiguity the W<sup>-</sup> and the missing momentum provides information on the direction and the energy of the undetected neutrino. Since quark and anti-quark are not identified in the hadronic W decays, there are two ambiguous solutions that normally enter both with a weight 0.5 in the angular distributions used for the measurements. Ambiguities are more important in the fully hadronic and fully leptonic channels. In the first case jets are paired using a likelihood fit to the W mass and the correct charge is extracted from an estimator. In the second case the consistency between the W mass and the mass of the $`l\nu `$ system helps in reconstructing the neutrino momenta. The quadratic nature of the constraints, however, always yield to a two-fold ambiguity: solutions obtained by flipping both neutrinos are equally valid. Again, both solutions enter in the experimental distribution with equal weight.
The charged TGCs are measured by the LEP experiments using the also the $`e^+e^{}W^{}e^\pm \nu (\overline{\nu })`$
and $`e^+e^{}\gamma \overline{\nu _e}\nu _e`$ events that are sensitive to the WW$`\gamma `$ vertex. The $`e^+e^{}W^{}e^\pm \nu (\overline{\nu })`$ events are reconstructed both in the leptonic $`Wl\nu `$and in the hadronic $`Wq\overline{q}`$decay channels. The leptonic channel selection requires a single good track identified as an electron or a muon with a missing energy and momentum. The hadronic selection requires two hadronic jets and a missing energy an momentum with an angular distribution not compatible with a WW event. The $`e^+e^{}\gamma \overline{\nu _e}\nu _e`$ events are obtained through two different processes shown in fig.( 3, a): the radiative return to the Z and the WW fusion. The two processes are well separated in the distribution of the recoil mass from the photon, in fig.( 3, b). While the first kind of events allow an alternative way to measure the number of neutrino family, by counting the number of events under the Z mass peak in fig. 3, the WW fusion processes, above the the Z mass peak, are sensitive to the TGC.
Fits to the triple gauge couplings are performed with methods where only one parameter is allowed to vary and the other two are fixed to their Standard Model prediction to improve the precision of measurement. The constraints obtained on the three triple gauge couplings from the combination of the LEP experiments $`^\mathrm{?}`$ are shown in fig. 4.
To study the correlations in the measurements of the TGC couplings also fits in 3-dimension allowing the 3 parameter free to vary or in 2-dimension, constraining one parameter to its Standard Model value and fit the other two, are implemented. The two dimension fit results are plotted in in fig. 5.
All the three cuplings are consistent with the Standard Model expectation, with a few percent precision.
ALEPH also performed a fit to all the 14 complex couplings, relaxing all the constraints on C- and P-conservation and $`\mathrm{SU}(2)_\mathrm{L}\times \mathrm{U}(1)_\mathrm{Y}`$ gauge invariance $`^\mathrm{?}`$. Out of all the 28 real parameters, one at the time is allowed to vary and the other are fixed to the Standard model predictions. The results of this test are:
$`Re(g_1^\gamma )`$ $`=`$ $`1.123\pm 0.091`$ (3)
$`Re(g_1^Z)`$ $`=`$ $`1.066\pm 0.073`$
$`Re(k_\gamma )`$ $`=`$ $`1.071\pm 0.062`$
$`Re(k_Z)`$ $`=`$ $`1.065\pm 0.061`$
All the other 24 parameters are consistent with zero ($`<`$ 0.05 $`÷`$ 0.20 at 95 $`\%`$ C.L.), in perfect agreement with the Standard Model.
### 2.2 Neutral triple gauge couplings
Neutral TGC, as ZZZ, ZZ$`\gamma `$ or Z$`\gamma `$$`\gamma `$ do not exist in the Standard Model at tree level and are negligible at the LEP energies. The presence of anomalous neutral couplings would affect the $`e^+e^{}Z(\gamma ^{})Z(\gamma ^{})`$ processes. The couplings describing the nTGC $`^\mathrm{?}`$ in the on shell ZZ production are 4, namely: $`f_4^\gamma `$, $`f_5^\gamma `$, $`f_4^Z`$ and $`f_5^Z`$, while 8 parameters are needed to describe the on shell production of Z$`\gamma `$: $`h_1^\gamma `$, $`h_2^\gamma `$, $`h_3^\gamma `$, $`h_4^\gamma `$, $`h_1^Z`$, $`h_2^Z`$, $`h_3^Z`$, $`h_4^Z`$. Similarly to the WW case, the differential cross section of the ZZ events respect to the Z polar angle and the angular distribution distributions of the fermions from the Z decay are taken into account, to estimate the nTGC couplings. Fit results are shown in tab. 2.
All these couplings are compatible with zero therefore no departure from the cross sections predicted by the Standard Model for these processes is observed.
## 3 Quartic Gauge Couplings
The Standard Model predicts the existence of verteces with four gauge bosons, like WW$`\gamma \gamma `$, WWZ$`\gamma `$, even if their rates are negligible at LEP energies. On the other hand, as in the case of the TGCs, the completely neutral verteces with four bosons, like ZZ$`\gamma \gamma `$, are not allowed in the Standard Model. The processes sensitive to the charged QGC are $`e^+e^{}W^+W^{}\gamma `$ , with a final state identical to that of a real photon radiative correction to the $`e^+e^{}W^+W^{}`$ process; and the WW fusion with two photon in the final state $`e^+e^{}\overline{\nu }\nu \gamma \gamma `$. The neutral QGC could be obtained through a process like $`e^+e^{}\gamma \gamma Z`$ .
The quartic couplings QGC are not dimensionless and are always referred to a parameter $`\mathrm{\Lambda }`$ which has the dimension of a mass and is commonly set to the W mass: $`\mathrm{\Lambda }=M_W`$. Three parameters are used to describe the charged QGC, $`a_0^W`$, $`a_c^W`$ and $`a_n^W`$ while only two parameters are necessary for the neutral QGC: $`a_0^Z`$ and $`a_c^Z`$. The constraints on these parameters, combining the four LEP experiments fit results, are showed intab. 3.
In both neutral and charged QGC no evidence of these anomalous couplings and therefore no deviation from the Standard Model has been observed.
## 4 Conclusions
The trilinear couplings TGC have been measured at LEP2 with a precision of few percent. The proof of the non-abelian stucture of the $`\mathrm{SU}(2)_\mathrm{L}\times \mathrm{U}(1)_\mathrm{Y}`$ gauge simmetry is obtained for the first time and this is one of the most important results of LEP2. No deviation from the Standard Model prediction is observed. The search of possible anomalous couplings for the trilinear neutral gauge couplings and the quartic gauge couplings is also presented. All the combined results are in good agreement with the Standard Model and no anomalous coupling has found.
## Acknowledgments
I would like to thank Paolo Azzurri, Timothy Barklow, Stephane Jezequel, Roberto Tenchini and Andrea Venturi for their help in preparing this talk and the LEP Electroweak Working Group for the combination results and the plots here presented.
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# Untitled Document
Parallel Bundles in Planar Map Geometries
Linfan Mao
Institute of Systems Science of Academy of Mathematics and Systems
Chinese Academy of Sciences, Beijing 100080, P.R.China
E-mail: [email protected]
Abstract: Parallel lines are very important objects in Euclid plane geometry and its behaviors can be gotten by one’s intuition. But in a planar map geometry, a kind of the Smarandache geometries, the situation is complex since it may contains elliptic or hyperbolic points. This paper concentrates on the behaviors of parallel bundles in planar map geometries, a generalization of parallel lines in plane geometry and obtains characteristics for parallel bundles.
Key Words: parallel bundle, planar map, Smarandache geometry, map geometry, classification.
AMS(2000): 05C15, 20H15, 51D99, 51M05
$`1`$ Introduction
A map is a connected topological graph cellularly embedded in a surface. On the past century, many works are concentrated on to find the combinatorial properties of maps, such as to determine whether exists a particularly embedding on a surface ($`[7][11]`$) or to enumerate a family of maps ($`[6]`$). All these works are on the side of algebra, not the object itself, i.e., geometry. For the later, more attentions are given to its element’s behaviors, such as, the line, angle, area, curvature, $`\mathrm{}`$, see also $`[12]`$ and $`[14]`$. For returning to its original face, the conception of map geometries is introduced in $`[10]`$. It is proved in $`[10]`$ that the map geometries are nice model of the Smarandache geometries. They are also a new kind of intrinsic geometry of surfaces ($`[1]`$). The main purpose of this paper is to determine the behaviors of parallel bundles in planar geometries, a generalization of parallel lines in the Euclid plane geometry.
An axiom is said Smarandachely denied if the axiom behaves in at least two different ways within the same space, i.e., validated and invalided, or only invalided but in multiple distinct ways.
A Smarandache geometry is a geometry which has at least one Smarandachely denied axiom($`1969`$)($`[5][13]`$).
In $`[3][4]`$, Iseri presented a nice model of the Smarandache geometries, called $`s`$-manifolds by using equilateral triangles, which is defined as follows($`[3][5][9]`$):
An $`s`$-manifold is any collection $`𝒞(T,n)`$ of these equilateral triangular disks $`T_i,1in`$ satisfying the following conditions:
$`(i)`$ Each edge $`e`$ is the identification of at most two edges $`e_i,e_j`$ in two distinct triangular disks $`T_i,T_j,1i,jn`$ and $`ij`$;
$`(ii)`$ Each vertex $`v`$ is the identification of one vertex in each of five, six or seven distinct triangular disks.
The conception of map geometries without boundary is defined as follows ($`[10]`$).
Definition $`1.1`$ For a given combinatorial map $`M`$, associates a real number $`\mu (u),0<\mu (u)<\pi `$, to each vertex $`u,uV(M)`$. Call $`(M,\mu )`$ a map geometry without boundary, $`\mu (u)`$ the angle factor of the vertex $`u`$ and to be orientablle or non-orientable if $`M`$ is orientable or not.
In $`[10]`$, it has proved that map geometries are the Smarandache geometries. The realization of each vertex $`u,uV(M)`$ in $`R^3`$ space is shown in the Fig.$`1`$ for each case of $`\rho (u)\mu (u)>2\pi `$, $`=2\pi `$ or $`<2\pi `$, call elliptic point, euclidean point and hyperbolic point, respectively.
$`\rho (u)\mu (u)<2\pi `$ $`\rho (u)\mu (u)=2\pi `$ $`\rho (u)\mu (u)>2\pi `$
Fig.$`1`$
Therefore, a line passes through an elliptic vertex, an euclidean vertex or a hyperbolic vertex $`u`$ has angle $`\frac{\rho (u)\mu (u)}{2}`$ at the vertex $`u`$. It is not $`180^{}`$ if the vertex $`u`$ is elliptic or hyperbolic. Then what is the angle of a line passes through a point on an edge of a map? It is $`180^{}`$? Since we wish the change of angles on an edge is smooth, the answer is not. For the Smarandache geometries, the parallel lines in them are need to be given more attention. We have the following definition.
Definition $`1.2`$ A family $``$ of infinite lines not intersecting each other in a planar geometry is called a parallel bundle.
In the Fig.$`2`$, we present all cases of parallel bundles passing through an edge in planar geometries, where, (a) is the case of points $`u,v`$ are same type with $`\rho (u)\mu (u)=\rho (v)\mu (v)`$, (b) and (c) the cases of same types with $`\rho (u)\mu (u)>\rho (v)\mu (v)`$ and (d) the case of $`u`$ is elliptic and $`v`$ hyperbolic.
Fig.$`2`$
Here, we assume the angle at the intersection point is in clockwise, that is, a line passing through an elliptic point will bend up and a hyperbolic point will bend down, such as the cases (b),(c) in the Fig.$`2`$. For a vector $`\stackrel{}{O}`$ on the Euclid plane, call it an orientation. We classify parallel bundles in planar map geometries along an orientation $`\stackrel{}{O}`$.
$`2.`$ A condition for parallel bundles
We investigate the behaviors of parallel bundles in the planar map geometries. For this object, we define a function $`f(x)`$ of angles on an edge of a planar map as follows.
Definition $`2.1`$ Denote by $`f(x)`$ the angle function of a line $`L`$ passing through an edge $`uv`$ at the point of distance $`x`$ to $`u`$ on the edge $`uv`$.
Then we get the following result.
Proposition $`2.1`$ A family $``$ of parallel lines passing through an edge $`uv`$ is a parallel bundle iff
$$\frac{df}{dx}|_+0.$$
Proof If $``$ is a parallel bundle, then any two lines $`L_1,L_2`$ will not intersect after them passing through the edge $`uv`$. Therefore, if $`\theta _1,\theta _2`$ are the angles of $`L_1,L_2`$ at the intersect points of $`L_1,L_2`$ with $`uv`$ and $`L_2`$ is far from $`u`$ than $`L_2`$, then we know that $`\theta _2\theta _1`$. Whence, for any point with $`x`$ distance from $`u`$ and $`\mathrm{\Delta }x>0`$, we have that
$$f(x+\mathrm{\Delta }x)f(x)0.$$
Therefore, we get that
$$\frac{df}{dx}|_+=\underset{\mathrm{\Delta }x+0}{lim}\frac{f(x+\mathrm{\Delta }x)f(x)}{\mathrm{\Delta }x}0.$$
As the cases in the Fig.$`1`$.
Now if $`\frac{df}{dx}|_+0`$, then $`f(y)f(x)`$ if $`yx`$. Since $``$ is a family of parallel lines before meeting $`uv`$, whence, any two lines in $``$ will not intersect each other after them passing through $`uv`$. Therefore, $``$ is a parallel bundle. $`\mathrm{}`$
A general condition for a family of parallel lines passing through a cut of a planar map being a parallel bundle is the following.
Proposition $`2.2`$ Let $`(M,\mu )`$ be a planar map geometry, $`C=\{u_1v_1,u_2v_2,\mathrm{},u_lv_l\}`$ a cut of the map $`M`$ with order $`u_1v_1,u_2v_2,\mathrm{},u_lv_l`$ from the left to the right, $`l1`$ and the angle functions on them are $`f_1,f_2,\mathrm{},f_l`$, respectively, also see the Fig.$`3`$.
Fig.$`3`$
Then a family $``$ of parallel lines passing through $`C`$ is a parallel bundle iff for any $`x,x0`$,
$`f_1^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+f_{3+}^{}(x)0`$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+\mathrm{}+f_{l+}^{}(x)0.`$
Proof According to the Proposition $`2.1`$, see the following Fig.$`4`$,
Fig.$`4`$
we know that any lines will not intersect after them passing through $`u_1v_1`$ and $`u_2v_2`$ iff for $`\mathrm{\Delta }x>0`$ and $`x0`$,
$$f_2(x+\mathrm{\Delta }x)+f_{1+}^{}(x)\mathrm{\Delta }xf_2(x).$$
That is,
$$f_{1+}^{}(x)+f_{2+}^{}(x)0.$$
Similarly, any lines will not intersect after them passing through $`u_1v_1,u_2v_2`$ and $`u_3v_3`$ iff for $`\mathrm{\Delta }x>0`$ and $`x0`$,
$$f_3(x+\mathrm{\Delta }x)+f_{2+}^{}(x)\mathrm{\Delta }x+f_{1+}^{}(x)\mathrm{\Delta }xf_3(x).$$
That is,
$$f_{1+}^{}(x)+f_{2+}^{}(x)+f_{3+}^{}(x)0.$$
Generally, any lines will not intersect after them passing through $`u_1v_1,u_2v_2,\mathrm{},u_{l1}v_{l1}`$ and $`u_lv_l`$ iff for $`\mathrm{\Delta }x>0`$ and $`x0`$,
$$f_l(x+\mathrm{\Delta }x)+f_{l1+}^{}(x)\mathrm{\Delta }x+\mathrm{}+f_{1+}^{}(x)\mathrm{\Delta }xf_l(x).$$
Whence, we get that
$$f_{1+}^{}(x)+f_{2+}^{}(x)+\mathrm{}+f_{l+}^{}(x)0.$$
Therefore, a family $``$ of parallel lines passing through $`C`$ is a parallel bundle iff for any $`x,x0`$, we have that
$`f_1^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+f_{3+}^{}(x)0`$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+\mathrm{}+f_{l+}^{}(x)0.`$
This completes the proof. $`\mathrm{}`$.
Corollary $`2.1`$ Let $`(M,\mu )`$ be a planar map geometry, $`C=\{u_1v_1,u_2v_2,\mathrm{},u_lv_l\}`$ a cut of the map $`M`$ with order $`u_1v_1,u_2v_2,\mathrm{},u_lv_l`$ from the left to the right, $`l1`$ and the angle functions on them are $`f_1,f_2,\mathrm{},f_l`$. Then a family $``$ of parallel lines passing through $`C`$ is still parallel lines after them leaving $`C`$ iff for any $`x,x0`$,
$`f_1^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+f_{3+}^{}(x)0`$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+\mathrm{}+f_{l1+}^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+\mathrm{}+f_{l+}^{}(x)=0.`$
Proof According to the Proposition $`2.2`$, we know the condition is a necessary and sufficient condition for $``$ being a parallel bundle. Now since lines in $``$ are parallel lines after them leaving $`C`$ iff for any $`x0`$ and $`\mathrm{\Delta }x0`$, there must be that
$$f_l(x+\mathrm{\Delta }x)+f_{l1+}^{}(x)\mathrm{\Delta }x+\mathrm{}+f_{1+}^{}(x)\mathrm{\Delta }x=f_l(x).$$
Therefore, we get that
$$f_{1+}^{}(x)+f_{2+}^{}(x)+\mathrm{}+f_{l+}^{}(x)=0\mathrm{}$$
When do the parallel lines parallel the initial parallel lines after them passing through a cut $`C`$ in a planar map geometry? The answer is in the following result.
Proposition $`2.3`$ Let $`(M,\mu )`$ be a planar map geometry, $`C=\{u_1v_1,u_2v_2,\mathrm{},u_lv_l\}`$ a cut of the map $`M`$ with order $`u_1v_1,u_2v_2,\mathrm{},u_lv_l`$ from the left to the right, $`l1`$ and the angle functions on them are $`f_1,f_2,\mathrm{},f_l`$. Then the parallel lines parallel the initial parallel lines after them passing through $`C`$ iff for $`x0`$,
$`f_1^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+f_{3+}^{}(x)0`$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+\mathrm{}+f_{l1+}^{}(x)0`$
and
$$f_1(x)+f_2(x)+\mathrm{}+f_l(x)=l\pi .$$
Proof According to the Proposition $`2.2`$ and Corollary $`2.1`$, we know the parallel lines passing through $`C`$ is a parallel bundle.
We calculate the angle $`\alpha (i,x)`$ of a line $`L`$ passing through an edge $`u_iv_i,1il`$ with the line before it meeting $`C`$ at the intersection of $`L`$ with the edge $`u_iv_i`$, where $`x`$ is the distance of the intersection point to $`u_1`$ on $`u_1v_1`$, see also the Fig.$`4`$. By the definition, we know the angle $`\alpha (1,x)=f(x)`$ and $`\alpha (2,x)=f_2(x)(\pi f_1(x))=f_1(x)+f_2(x)\pi `$.
Now if $`\alpha (i,x)=f_1(x)+f_2(x)+\mathrm{}+f_i(x)(i1)\pi `$, then similar to the case $`i=2`$, we know that $`\alpha (i+1,x)=f_{i+1}(x)(\pi \alpha (i,x))=f_{i+1}(x)+\alpha (i,x)\pi `$. Whence, we get that
$$\alpha (i+1,x)=f_1(x)+f_2(x)+\mathrm{}+f_{i+1}(x)i\pi .$$
Notice that a line $`L`$ parallel the initial parallel line after it passing through $`C`$ iff $`\alpha (l,x)=\pi `$, i.e.,
$$f_1(x)+f_2(x)+\mathrm{}+f_l(x)=l\pi .$$
This completes the proof. $`\mathrm{}`$
$`3.`$ Linear condition and combinatorial realization for parallel bundles
For the simplicity, we can assume the function $`f(x)`$ is linear and denoted it by $`f_l(x)`$. We can calculate $`f_l(x)`$ as follows.
Proposition $`3.1`$ The angle function $`f_l(x)`$ of a line $`L`$ passing through an edge $`uv`$ at the point with distance $`x`$ to u is
$$f_l(x)=(1\frac{x}{d(uv)})\frac{\rho (u)\mu (v)}{2}+\frac{x}{d(uv)}\frac{\rho (v)\mu (v)}{2},$$
where, $`d(uv)`$ is the length of the edge $`uv`$.
Proof Since $`f_l(x)`$ is linear, we know that $`f_l(x)`$ satisfies the following equation.
$$\frac{f_l(x)\frac{\rho (u)\mu (u)}{2}}{\frac{\rho (v)\mu (v)}{2}\frac{\rho (u)\mu (u)}{2}}=\frac{x}{d(uv)},$$
Calculation shows that
$$f_l(x)=(1\frac{x}{d(uv)})\frac{\rho (u)\mu (v)}{2}+\frac{x}{d(uv)}\frac{\rho (v)\mu (v)}{2}.\mathrm{}$$
Corollary $`3.1`$ Under the linear assumption, a family $``$ of parallel lines passing through an edge $`uv`$ is a parallel bundle iff
$$\frac{\rho (u)}{\rho (v)}\frac{\mu (v)}{\mu (u)}.$$
Proof According to the Proposition $`2.1`$, a family of parallel lines passing through an edge $`uv`$ is a parallel bundle iff for $`x,x0`$, $`f^{}(x)0`$, i.e.,
$$\frac{\rho (v)\mu (v)}{2d(uv)}\frac{\rho (u)\mu (u)}{2d(uv)}0.$$
Therefore, a family $``$ of parallel lines passing through an edge $`uv`$ is a parallel bundle iff
$$\rho (v)\mu (v)\rho (u)\mu (u).$$
Whence,
$$\frac{\rho (u)}{\rho (v)}\frac{\mu (v)}{\mu (u)}.\mathrm{}$$
For a family of parallel lines pass through a cut, we have the following condition for it being a parallel bundle.
Proposition $`3.2`$ Let $`(M,\mu )`$ be a planar map geometry, $`C=\{u_1v_1,u_2v_2,\mathrm{},u_lv_l\}`$ a cut of the map $`M`$ with order $`u_1v_1,u_2v_2,\mathrm{},u_lv_l`$ from the left to the right, $`l1`$. Then under the linear assumption, a family $`L`$ of parallel lines passing through $`C`$ is a parallel bundle iff the angle factor $`\mu `$ satisfies the following linear inequality system
$$\rho (v_1)\mu (v_1)\rho (u_1)\mu (u_1)$$
$$\frac{\rho (v_1)\mu (v_1)}{d(u_1v_1)}+\frac{\rho (v_2)\mu (v_2)}{d(u_2v_2)}\frac{\rho (u_1)\mu (u_1)}{d(u_1v_1)}+\frac{\rho (u_2)\mu (u_2)}{d(u_2v_2)}$$
$$\mathrm{}\mathrm{}\mathrm{}\mathrm{}$$
$`{\displaystyle \frac{\rho (v_1)\mu (v_1)}{d(u_1v_1)}}`$ $`+`$ $`{\displaystyle \frac{\rho (v_2)\mu (v_2)}{d(u_2v_2)}}+\mathrm{}+{\displaystyle \frac{\rho (v_l)\mu (v_l)}{d(u_lv_l)}}`$
$``$ $`{\displaystyle \frac{\rho (u_1)\mu (u_1)}{d(u_1,v_1)}}+{\displaystyle \frac{\rho (u_2)\mu (u_2)}{d(u_2,v_2)}}+\mathrm{}+{\displaystyle \frac{\rho (u_l)\mu (u_l)}{d(u_l,v_l)}}.`$
Proof Under the linear assumption, for any integer $`i,i1`$, we know that
$$f_{i+}^{}(x)=\frac{\rho (v_i)\mu (v_i)\rho (u_i)\mu (u_i)}{2d(u_iv_i)}$$
by the Proposition $`3.1`$. Whence, according to the Proposition $`2.2`$, we get that a family $`L`$ of parallel lines passing through $`C`$ is a parallel bundle iff the angle factor $`\mu `$ satisfies the following linear inequality system
$$\rho (v_1)\mu (v_1)\rho (u_1)\mu (u_1)$$
$$\frac{\rho (v_1)\mu (v_1)}{d(u_1v_1)}+\frac{\rho (v_2)\mu (v_2)}{d(u_2v_2)}\frac{\rho (u_1)\mu (u_1)}{d(u_1v_1)}+\frac{\rho (u_2)\mu (u_2)}{d(u_2v_2)}$$
$$\mathrm{}\mathrm{}\mathrm{}\mathrm{}$$
$`{\displaystyle \frac{\rho (v_1)\mu (v_1)}{d(u_1v_1)}}`$ $`+`$ $`{\displaystyle \frac{\rho (v_2)\mu (v_2)}{d(u_2v_2)}}+\mathrm{}+{\displaystyle \frac{\rho (v_l)\mu (v_l)}{d(u_lv_l)}}`$
$``$ $`{\displaystyle \frac{\rho (u_1)\mu (u_1)}{d(u_1,v_1)}}+{\displaystyle \frac{\rho (u_2)\mu (u_2)}{d(u_2,v_2)}}+\mathrm{}+{\displaystyle \frac{\rho (u_l)\mu (u_l)}{d(u_l,v_l)}}.`$
This completes the proof. $`\mathrm{}`$
For planar maps underlying a regular graph, we have the following interesting results for parallel bundles.
Corollary $`3.2`$ Let $`(M,\mu )`$ be a planar map geometry with $`M`$ underlying a regular graph, $`C=\{u_1v_1,u_2v_2,\mathrm{},u_lv_l\}`$ a cut of the map $`M`$ with order $`u_1v_1,u_2v_2,\mathrm{},u_lv_l`$ from the left to the right, $`l1`$. Then under the linear assumption, a family $`L`$ of parallel lines passing through $`C`$ is a parallel bundle iff the angle factor $`\mu `$ satisfies the following linear inequality system
$$\mu (v_1)\mu (u_1)$$
$$\frac{\mu (v_1)}{d(u_1v_1)}+\frac{\mu (v_2)}{d(u_2v_2)}\frac{\mu (u_1)}{d(u_1v_1)}+\frac{\mu (u_2)}{d(u_2v_2)}$$
$$\mathrm{}\mathrm{}\mathrm{}\mathrm{}$$
$$\frac{\mu (v_1)}{d(u_1v_1)}+\frac{\mu (v_2)}{d(u_2v_2)}+\mathrm{}+\frac{\mu (v_l)}{d(u_lv_l)}\frac{\mu (u_1)}{d(u_1v_1)}+\frac{\mu (u_2)}{d(u_2v_2)}+\mathrm{}+\frac{\mu (u_l)}{d(u_lv_l)}$$
and particularly, if assume that all the lengths of edges in $`C`$ are the same, then
$`\mu (v_1)`$ $``$ $`\mu (u_1)`$
$`\mu (v_1)+\mu (v_2)`$ $``$ $`\mu (u_1)+\mu (u_2)`$
$`\mathrm{}\mathrm{}`$ $`\mathrm{}`$ $`\mathrm{}\mathrm{}`$
$`\mu (v_1)+\mu (v_2)+\mathrm{}+\mu (v_l)`$ $``$ $`\mu (u_1)+\mu (u_2)+\mathrm{}+\mu (u_l).`$
Certainly, by choosing different angle factors, we can also get combinatorial conditions for existing parallel bundles under the linear assumption.
Proposition $`3.3`$ Let $`(M,\mu )`$ be a planar map geometry, $`C=\{u_1v_1,u_2v_2,\mathrm{},u_lv_l\}`$ a cut of the map $`M`$ with order $`u_1v_1,u_2v_2,\mathrm{},u_lv_l`$ from the left to the right, $`l1`$. If for any integer $`i,i1`$,
$$\frac{\rho (u_i)}{\rho (v_i)}\frac{\mu (v_i)}{\mu (u_i)},$$
then under the linear assumption, a family $`L`$ of parallel lines passing through $`C`$ is a parallel bundle.
Proof Notice that under the linear assumption, for any integer $`i,i1`$, we know that
$$f_{i+}^{}(x)=\frac{\rho (v_i)\mu (v_i)\rho (u_i)\mu (u_i)}{2d(u_iv_i)}$$
by the Proposition $`3.1`$. Whence, $`f_{i+}^{}(x)0`$ for $`i=1,2,\mathrm{},l`$. Therefore, we get that
$`f_1^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)0`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+f_{3+}^{}(x)0`$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`f_{1+}^{}(x)+f_{2+}^{}(x)+\mathrm{}+f_{l+}^{}(x)0.`$
By the Proposition $`2.2`$, we know that a family $`L`$ of parallel lines passing through $`C`$ is a parallel bundle. $`\mathrm{}`$
$`4.`$ Classification of parallel bundles
For a cut $`C`$ in a planar map geometry and $`eC`$, denote by $`f_e(x)`$ the angle function on the edge $`e`$, $`f(C,x)=\underset{eC}{}f_e(x)`$. If $`f(C,x)`$ is independent on $`x`$, then we abbreviate it to $`f(C)`$. According to the results in the Section $`2`$ and $`3`$, we can classify the parallel bundles with a given orientation $`\stackrel{}{O}`$ in planar map geometries into the following $`15`$ classes, where, each class is labelled by a $`4`$-tuple $`0,1`$ code.
Classification of parallel bundles
$`(1)`$ $`𝒞_{1000}`$: for any cut $`C`$ along $`\stackrel{}{O}`$, $`f(C)=|C|\pi `$;
$`(2)`$ $`𝒞_{0100}`$: for any cut $`C`$ along $`\stackrel{}{O}`$, $`f(C)<|C|\pi `$;
$`(3)`$ $`𝒞_{0010}`$: for any cut $`C`$ along $`\stackrel{}{O}`$, $`f(C)>|C|\pi `$ ;
$`(4)`$ $`𝒞_{0001}`$: for any cut $`C`$ along $`\stackrel{}{O}`$, $`f_+^{}(C,x)>0`$ for $`x,x0;`$
$`(5)`$ $`𝒞_{1100}`$: There exist cuts $`C_1,C_2`$ along $`\stackrel{}{O}`$, such that $`f(C_1)=|C_1|\pi `$ and $`f(C_2)=c<|C_2|\pi `$;
$`(6)`$ $`𝒞_{1010}`$: there exist cuts $`C_1,C_2`$ along $`\stackrel{}{O}`$, such that $`f(C_1)=|C_1|\pi `$ and $`f(C_2)>|C_2|\pi `$;
$`(7)`$ $`𝒞_{1001}`$: there exist cuts $`C_1,C_2`$ along $`\stackrel{}{O}`$, such that $`f(C_1)=|C_1|\pi `$ and $`f_+^{}(C_2,x)>0`$ for $`x,x0;`$
$`(8)`$ $`𝒞_{0110}`$: there exist cuts $`C_1,C_2`$ along $`\stackrel{}{O}`$, such that $`f(C_1)<|C_1|\pi `$ and $`f(C_2)>|C_2|\pi `$;
$`(9)`$ $`𝒞_{0101}`$: there exist cuts $`C_1,C_2`$ along $`\stackrel{}{O}`$, such that $`f(C_1)<|C_1|\pi `$ and $`f_+^{}(C_2,x)>0`$ for $`x,x0;`$
$`(10)`$ $`𝒞_{0011}`$: there exist cuts $`C_1,C_2`$ along $`\stackrel{}{O}`$, such that $`f(C_1)>|C_1|\pi `$ and $`f_+^{}(C_2,x)>0`$ for $`x,x0;`$
$`(11)`$ $`𝒞_{1110}`$: there exist cuts $`C_1,C_2`$ and $`C_3`$ along $`\stackrel{}{O}`$, such that $`f(C_1)=|C_1|\pi `$, $`f(C_2)<|C_2|\pi `$ and $`f(C_3)>|C_3|\pi `$;
$`(12)`$ $`𝒞_{1101}`$: there exist cuts $`C_1,C_2`$ and $`C_3`$ along $`\stackrel{}{O}`$, such that $`f(C_1)=|C_1|\pi `$, $`f(C_2)<|C_2|\pi `$ and $`f_+^{}(C_3,x)>0`$ for $`x,x0;`$
$`(13)`$ $`𝒞_{1011}`$: there exist cuts $`C_1,C_2`$ and $`C_3`$ along $`\stackrel{}{O}`$, such that $`f(C_1)=|C_1|\pi `$, $`f(C_2)>|C_2|\pi `$ and $`f_+^{}(C_1,x)>0`$ for $`x,x0;`$
$`(14)`$ $`𝒞_{0111}`$: there exist cuts $`C_1,C_2`$ and $`C_3`$ along $`\stackrel{}{O}`$, such that $`f(C_1)<|C_1|\pi `$, $`f(C_2)>|C_2|\pi `$ and $`f_+^{}(C_1,x)>0`$ for $`x,x0;`$
$`(15)`$ $`𝒞_{1111}`$: there exist cuts $`C_1,C_2,C_3`$ and $`C_4`$ along $`\stackrel{}{O}`$, such that $`f(C_1)=|C_1|\pi `$, $`f(C_2)<|C_2|\pi `$, $`f(C_3)>|C_3|\pi `$ and $`f_+^{}(C_4,x)>0`$ for $`x,x0.`$
Notice that only the first three classes may be parallel lines after them passing through the cut $`C`$. All of the other classes are only parallel bundles, not parallel lines in the usual meaning.
Proposition $`4.1`$ For an orientation $`\stackrel{}{O}`$, the $`15`$ classes $`𝒞_{1000}𝒞_{1111}`$ are all the parallel bundles in planar map geometries.
Proof Not loss of generality, we assume $`C_1,C_2,\mathrm{},C_m,m1`$, are all the cuts along $`\stackrel{}{O}`$ in a planar map geometry $`(M,\mu )`$ from the upon side of $`\stackrel{}{O}`$ to its down side. We find their structural characters for each case in the following discussion.
$`𝒞_{1000}`$: By the Proposition $`2.3`$, a family $``$ of parallel lines parallel their initial lines before meeting $`M`$ after the passing through $`M`$.
$`𝒞_{0100}`$: By the definition, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if
$$f(C_1)f(C_2)\mathrm{}f(C_m)<\pi .$$
Otherwise, some lines in $``$ will intersect. According to the Corollary $`2.1`$, they parallel each other after they passing through $`M`$ only if
$$f(C_1)=f(C_2)=\mathrm{}=f(C_m)<\pi .$$
$`𝒞_{0010}`$: Similar to the case $`𝒞_{0100}`$, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if
$$\pi <f(C_1)f(C_2)\mathrm{}f(C_m)$$
and parallel each other after they passing through $`M`$ only if
$$\pi <f(C_1)=f(C_2)=\mathrm{}=f(C_m).$$
$`𝒞_{0001}`$: Notice that by the proof of the Proposition $`2.3`$, a line has angle $`f(C,x)(|C|1)\pi `$ after it passing through $`C`$ with the initial line before meeting $`C`$. In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if for $`x_i,x_i0,1im`$,
$$f(C_1,x_1)f(C_2,x_2)\mathrm{}f(C_m,x_m).$$
Otherwise, they will intersect.
$`𝒞_{1100}`$: In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is an integer $`k,2km`$, such that
$$f(C_1)f(C_2)\mathrm{}f(C_{k1})<f(C_k)=f(C_{k+1})=\mathrm{}=f(C_m)=\pi .$$
Otherwise, they will intersect.
$`𝒞_{1010}`$: Similar to the case $`𝒞_{1100}`$, in this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is an integer $`k,2km`$, such that
$$\pi =f(C_1)=f(C_2)=\mathrm{}=f(C_k)<f(C_{k+1})\mathrm{}f(C_m).$$
Otherwise, they will intersect.
$`𝒞_{1001}`$: In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is an integer $`k,l,1k<lm`$, such that for $`x_i,x_i0,1ik`$ or $`lim`$,
$`f(C_1,x_1)`$ $``$ $`f(C_2,x_2)\mathrm{}f(C_k,x_k)<f(C_{k+1})`$
$`=`$ $`f(C_{k+2})=\mathrm{}=f(C_{l1})=\pi <f(C_l,x_l)\mathrm{}f(C_m,x_m).`$
Otherwise, they will intersect.
$`𝒞_{0110}`$: In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is integers $`k,1k<m`$, such that
$$f(C_1)f(C_2)\mathrm{}f(C_k)<\pi <f(C_{k+1})\mathrm{}f(C_m).$$
Otherwise, they will intersect.
$`𝒞_{0101}`$: In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is integers $`k,1km`$, such that for $`x_i,x_i0,1im`$,
$$f(C_1,x_1)f(C_2,x_2)\mathrm{}f(C_k,x_k)<\pi f(C_{k+1},x_{k+1})\mathrm{}f(C_m,x_m),$$
and there must be a constant in $`f(C_1,x_1),f(C_2,x_2),\mathrm{},f(C_k,x_k)`$.
$`𝒞_{0011}`$: In this case, the situation is similar to the case $`𝒞_{0101}`$ and there must be a constant in $`f(C_{k+1},x_{k+1}),f(C_{k+2},x_{k+2}),\mathrm{},f(C_m,x_m)`$.
$`𝒞_{1110}`$: In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is an integer $`k,l,1k<lm`$, such that
$`f(C_1)`$ $``$ $`f(C_2)\mathrm{}f(C_k)<f(C_{k+1})`$
$`=`$ $`\mathrm{}=f(C_{l1})=\pi <f(C_l)\mathrm{}f(C_m).`$
Otherwise, they will intersect.
$`𝒞_{1101}`$: In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is an integer $`k,l,1k<lm`$, such that for $`x_i,x_i0,1ik`$ or $`lim`$,
$`f(C_1,x_1)`$ $``$ $`f(C_2,x_2)\mathrm{}f(C_k,x_k)<f(C_{k+1})`$
$`=`$ $`\mathrm{}=f(C_{l1})=\pi <f(C_l,x_l)\mathrm{}f(C_m,x_m)`$
and there must be a constant in $`f(C_1,x_1),f(C_2,x_2),\mathrm{},f(C_k,x_k)`$. Otherwise, they will intersect.
$`𝒞_{1011}`$: In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is an integer $`k,l,1k<lm`$, such that for $`x_i,x_i0,1ik`$ or $`lim`$,
$`f(C_1,x_1)`$ $``$ $`f(C_2,x_2)\mathrm{}f(C_k,x_k)<f(C_{k+1})`$
$`=`$ $`\mathrm{}=f(C_{l1})=\pi <f(C_l,x_l)\mathrm{}f(C_m,x_m)`$
and there must be a constant in $`f(C_l,x_l),f(C_{l+1},x_{l+1}),\mathrm{},f(C_m,x_m)`$. Otherwise, they will intersect.
$`𝒞_{0111}`$: In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is an integer $`k,1km`$, such that for $`x_i,x_i0`$,
$$f(C_1,x_1)f(C_2,x_2)\mathrm{}f(C_k,x_k)<\pi <f(C_l,x_l)\mathrm{}f(C_m,x_m)$$
and there must be a constant in $`f(C_1,x_1),f(C_2,x_2),\mathrm{},f(C_k,x_k)`$ and a constant in $`f(C_l,x_l),f(C_{l+1},x_{l+1}),\mathrm{},f(C_m,x_m)`$. Otherwise, they will intersect.
$`𝒞_{1111}`$: In this case, a family $``$ of parallel lines is a parallel bundle along $`\stackrel{}{O}`$ only if there is an integer $`k,l,1k<lm`$, such that for $`x_i,x_i0,1ik`$ or $`lim`$,
$`f(C_1,x_1)`$ $``$ $`f(C_2,x_2)\mathrm{}f(C_k,x_k)<f(C_{k+1})`$
$`=`$ $`\mathrm{}=f(C_{l1})=\pi <f(C_l,x_l)\mathrm{}f(C_m,x_m)`$
and there must be a constant in $`f(C_1,x_1),f(C_2,x_2),\mathrm{},f(C_k,x_k)`$ and a constant in $`f(C_l,x_l),f(C_{l+1},x_{l+1}),\mathrm{},f(C_m,x_m)`$. Otherwise, they will intersect.
Following the structural characters of the classes $`𝒞_{1000}𝒞_{1111}`$, by the Proposition $`2.2`$, $`2.3`$ and Proposition $`3.1`$, we know that any parallel bundle is in one of the classes $`𝒞_{1000}𝒞_{1111}`$ and each class in $`𝒞_{1000}𝒞_{1111}`$ is non-empty. This completes the proof. $`\mathrm{}`$
A example of parallel bundle in a planar map geometry is shown in the Fig.$`5`$, in where the number on a vertex $`u`$ denotes the number $`\rho (u)\mu (u)`$.
Fig.$`5`$
$`5.`$ Generalization
All the planar map geometries considered in this paper are without boundary. For planar map geometries with boundary, i.e., some faces are deleted ($`[10]`$), which are correspondence with the maps with boundary ($`[2]`$). We know that they are the Smarandache non-geometries, satisfying one or more of the following conditions:
($`A1^{}`$)It is not always possible to draw a line from an arbitrary point to another arbitrary point.
($`A2^{}`$)It is not always possible to extend by continuity a finite line to an infinite line.
($`A3^{}`$)It is not always possible to draw a circle from an arbitrary point and of an arbitrary interval.
($`A4^{}`$)not all the right angles are congruent.
($`A5^{}`$)if a line, cutting two other lines, forms the interior angles of the same side of it strictly less than two right angle, then not always the two lines extended towards infinite cut each other in the side where the angles are strictly less than two right angle.
Notice that for an one face planar map geometry $`(M,\mu )^1`$ with boundary, if we choose all points being euclidean, then $`(M,\mu )^1`$ is just the Poincaré’s model for the hyperbolic geometry.
Using the neutrosophic logic idea, we can also define the conception of neutrosophic surface as follow, comparing also with the surfaces in $`[8]`$ and $`[14]`$.
Definition $`5.1`$ A neutrosophic surface is a Hausdorff, connected, topological space $`S`$ such that every point $`v`$ is elleptic, euclidean, or hyperbolic.
For this kind of surface, we present the following problem for the further researching.
Problem $`5.1`$ To determine the behaviors of elements, such as, the line, angle, area, $`\mathrm{}`$, in neutrosophic surfaces.
Notice that results in this paper are just the behaviors of line bundles in a neutrosophic plane.
References
A.D.Aleksandrov and V.A.Zalgaller, Intrinsic geometry of surfaces, American Mathematical Society, 1967.
R.P.Bryant and D.Singerman, Foundations of the theory of maps on surfaces with boundary,Quart.J.Math.Oxford(2),36(1985), 17-41.
H.Iseri, Smarandache manifolds, American Research Press, Rehoboth, NM,2002.
H.Iseri, Partially Paradoxist Smarandache Geometries, http://www.gallup.unm. edu/s̃marandache/Howard-Iseri-paper.htm.
L.Kuciuk and M.Antholy, An Introduction to Smarandache Geometries, Mathematics Magazine, Aurora, Canada, Vol.12(2003)
Y.P.Liu, Enumerative Theory of Maps, Kluwer Academic Publisher, Dordrecht / Boston / London (1999).
Y.P.Liu, Embeddability in Graphs, Kluwer Academic Publisher, Dordrecht / Boston / London (1995).
Mantredo P.de Carmao, Differential Geometry of Curves and Surfaces, Pearson Education asia Ltd (2004).
L.F.Mao, Automorphism groups of maps, surfaces and Smarandache geometries, American Research Press, Rehoboth, NM,2005.
L.F.Mao, A new view of combinatorial maps by Smarandache’s notion, arXiv: Math.GM/0506232, will also appear in Smarandache Notion J.
B.Mohar and C.Thomassen, Graphs on Surfaces, The Johns Hopkins University Press, London, 2001.
V.V.Nikulin and I.R.Shafarevlch, Geometries and Groups, Springer-Verlag Berlin Heidelberg (1987)
F. Smarandache, Mixed noneuclidean geometries, eprint arXiv: math/0010119, 10/2000.
J.Stillwell, Classical topology and combinatorial group theory, Springer-Verlag New York Inc., (1980).
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# Oscillations of high energy neutrinos in matter: Precise formalism and parametric resonance
## Abstract
We present a formalism for precise description of oscillation phenomena in matter at high energies or high densities, $`V>\mathrm{\Delta }m^2/2E`$, where $`V`$ is the matter-induced potential of neutrinos. The accuracy of the approximation is determined by the quantity $`\mathrm{sin}^22\theta _m\mathrm{\Delta }V/2\pi V`$, where $`\theta _m`$ is the mixing angle in matter and $`\mathrm{\Delta }V`$ is a typical change of the potential over the oscillation length ($`l2\pi /V`$). We derive simple and physically transparent formulas for the oscillation probabilities, which are valid for arbitrary matter density profiles. They can be applied to oscillations of high energy ($`E>10`$ GeV) accelerator, atmospheric and cosmic neutrinos in the matter of the Earth, substantially simplifying numerical calculations and providing an insight into the physics of neutrino oscillations in matter. The effect of parametric enhancement of the oscillations of high energy neutrinos is considered. Future high statistics experiments can provide an unambiguous evidence for this effect.
preprint: YITP-SB-05-16
1. Introduction. Neutrino physics enters a new phase now, where the objectives are precision measurements of the parameters, studies of subleading oscillation effects and searches for new physics beyond the already standard picture, which includes non-zero neutrino masses and mixing. Detection of neutrinos from new sources, in particular, of cosmic neutrinos, is in the agenda.
Substantial new information is expected from the studies of high energy ($`E>1`$ GeV) neutrinos. This includes investigations of atmospheric neutrinos with new large volume detectors fatm , long baseline accelerator experiments LBL and detection of cosmic neutrinos from galactic and extragalactic sources, as well as of neutrinos produced in interactions of cosmic rays with the solar atmosphere solatm . Another possible source of neutrinos is annihilation of hypothetical WIMP’s in the center of the Earth and the Sun WIMPS . In all these cases beams of high energy neutrinos can propagate significant distances in the matter of the Earth (or of the Sun) and therefore undergo oscillations/conversions in matter.
Increased accuracy and reach of neutrino experiments put forward new and more challenging demands to the theoretical description of neutrino oscillations. In the present Letter we derive very simple and accurate analytic formulas describing oscillations of neutrinos in matter. Our primary goal is to study oscillations of high energy neutrinos note1 , but the formulas we obtain are actually applicable in a wide range of neutrino energies. They simplify substantially numerical calculations and allow a deep insight into the physics of neutrino conversions in matter. In particular, they provide a useful tool for studying parametric resonance enhancement of neutrino oscillations.
The parametric resonance can occur in oscillating systems with varying parameters due to a specific correlation between the rate of the change of these parameters and the values of the parameters themselves. In the case of neutrino oscillations in matter, the parametric enhancement is realized when the variation of the matter density along the neutrino trajectory is in a certain way correlated with the change of the oscillation phase ETC ; Akh1 . This effect is different from the MSW effect W ; MS , and yet can result in a strong amplification of the oscillations. We apply the formalism developed here to study the parametric resonance in oscillations of high energy neutrinos in the Earth.
2. Formalism. We consider oscillations in the 3-flavor neutrino system ($`\nu _e,\nu _\mu ,\nu _\tau `$), with the mass squared differences $`\mathrm{\Delta }m_{31}^2`$ and $`\mathrm{\Delta }m_{21}^2`$ responsible for the oscillations of atmospheric and solar neutrinos, respectively. In certain situations, the full three-flavor problem is approximately reduced to the effective two-flavor ones, in which the electron neutrino $`\nu _e`$ oscillates into a combination of $`\nu _\mu `$ and $`\nu _\tau `$. We shall be mainly interested in oscillations of neutrinos with energies $`E>\mathrm{\Delta }m_{31}^2/2V`$, where the matter-induced potential of neutrinos $`V(x)\sqrt{2}G_FN_e(x)`$, with $`N_e(x)`$ the electron number density in matter and $`G_F`$ the Fermi constant. Numerically, this corresponds to $`E>810`$ GeV for the matter of the Earth. In this case the 1-2 mixing is strongly suppressed by matter, and the problem is reduced to an effective two-flavor one, described by the potential $`V`$, mass squared difference $`\mathrm{\Delta }m^2\mathrm{\Delta }m_{31}^2`$ and vacuum mixing angle $`\theta \theta _{13}`$ (which is assumed to be non-zero) ADLS . In particular, the oscillations of electron neutrinos are determined by the transition probability $`P_2P(\nu _e\nu _a)`$, where $`\nu _a=\mathrm{sin}\theta _{23}\nu _\mu +\mathrm{cos}\theta _{23}\nu _\tau `$ is the state which mixes with $`\nu _e`$ in the third mass eigenstate. In terms of $`P_2`$ the flavor transition probabilities are $`P(\nu _e\nu _\mu )=P(\nu _\mu \nu _e)=\mathrm{sin}^2\theta _{23}P_2`$, $`P(\nu _e\nu _\tau )=P(\nu _\tau \nu _e)=\mathrm{cos}^2\theta _{23}P_2`$ ADLS . The third state (orthogonal to $`\nu _a`$) decouples from the rest of the system and evolves independently.
In the ($`\nu _e,\nu _a)`$ basis, the evolution matrix $`S(x)`$ describing neutrino oscillations satisfies the equation
$$i\frac{dS}{dx}=H(x)S,$$
(1)
with the Hamiltonian
$$H(x)=\frac{V}{2}(\begin{array}{cc}1& 0\\ 0& 1\end{array})+\delta (\begin{array}{cc}\mathrm{cos}2\theta & \mathrm{sin}2\theta \\ \mathrm{sin}2\theta & \mathrm{cos}2\theta \end{array}).$$
(2)
Here
$$\delta \mathrm{\Delta }m^2/4E,$$
(3)
and the first (potential) term dominates in the high energy limit. However, in most situations of interest the neutrino path length in matter $`L`$ satisfies $`\delta L1`$, therefore we cannot consider the whole second term as a small perturbation, and the effect of $`\delta `$ on the neutrino energy level splitting should be taken into account. For this reason we split the Hamiltonian as
$$H=H_0+H_I,$$
(4)
with
$$H_0=\omega (\begin{array}{cc}1& 0\\ 0& 1\end{array}),H_I=\mathrm{sin}2\theta \delta (\begin{array}{cc}ϵ& 1\\ 1& ϵ\end{array}).$$
(5)
Here
$$\omega (x)\sqrt{(V/2\delta \mathrm{cos}2\theta )^2+\delta ^2\mathrm{sin}^22\theta },$$
(6)
$`2\omega `$ being the difference of the eigenvalues of $`H(x)`$;
$$ϵ\frac{\mathrm{cos}2\theta \delta V/2+\omega }{\mathrm{sin}2\theta \delta }\frac{\delta }{V}\mathrm{sin}2\theta 1.$$
(7)
The ratio of the second and the first terms in the Hamiltonian in (5) is determined by the mixing angle in matter $`\theta _m`$:
$$\frac{\mathrm{sin}2\theta \delta }{\omega }=\mathrm{sin}2\theta _m.$$
(8)
Therefore for $`\mathrm{sin}2\theta _m1`$ the term $`H_I`$ can be considered as a perturbation. Furthermore, according to (7), $`ϵ\mathrm{sin}2\theta _m`$, so that the diagonal terms in $`H_I`$ can be neglected in the lowest approximation.
We seek the solution of eq. (1) in the form
$$S=S_0S_I,$$
(9)
where $`S_0`$ is the solution of the evolution equation with $`H`$ replaced by $`H_0`$. From (5) we find
$$S_0(x)=\left(\begin{array}{cc}e^{i\varphi (x)}& 0\\ 0& e^{i\varphi (x)}\end{array}\right),$$
(10)
where
$$\varphi (x)_0^x𝑑x^{}\omega (x^{})$$
(11)
is the adiabatic phase. Then, according to (1), (4) and (9), the matrix $`S_I`$ satisfies the equation
$$i\frac{dS_I}{dx}=S_0^1H_IS_0S_I=\stackrel{~}{H}_IS_I,$$
(12)
where $`\stackrel{~}{H}_IS_0^1H_IS_0`$ is the perturbation Hamiltonian in the “interaction” representation. Eq. (12) can be solved by iterations: $`S_I=\mathrm{𝟙}+S_I^{(1)}+\mathrm{}`$, which leads to the standard perturbation series for the $`S`$ matrix. For neutrino propagation between $`x=0`$ and $`x=L`$ we have
$$S_I(L)=\mathrm{𝟙}i_0^L𝑑x\stackrel{~}{H}_I(x)_0^L𝑑x\stackrel{~}{H}_I(x)_0^x𝑑x^{}\stackrel{~}{H}_I(x^{})\mathrm{}$$
(13)
Taking $`S_I`$ to the first order, we obtain from (9) the evolution matrix
$$S(L)=S_0(L)\left[\mathrm{𝟙}i\delta \mathrm{sin}2\theta _0^L𝑑x\left(\begin{array}{cc}0& e^{i2\varphi (x)}\\ e^{i2\varphi (x)}& 0\end{array}\right)\right].$$
(14)
The $`\nu _e\nu _a`$ transition probability $`P_2(L)`$ is given by the squared modulus of the off-diagonal element $`[S(L)]_{ae}`$:
$$P_2=\delta ^2\mathrm{sin}^22\theta \left|_0^L𝑑xe^{i2\varphi (x)}\right|^2.$$
(15)
For density profiles that are symmetric with respect to the center of the neutrino trajectory, $`V(x)=V(Lx)`$, eq. (15) gives
$$P_2=4\left(\frac{\mathrm{\Delta }m^2}{4E}\right)^2\mathrm{sin}^22\theta \left[_0^{L/2}𝑑z\mathrm{cos}2\varphi (z)\right]^2,$$
(16)
where $`z=xL/2`$ is the distance from the midpoint of the trajectory and $`\varphi (z)`$ is the phase acquired between this midpoint and the point $`z`$.
The transition probability $`P_2`$ scales with neutrino energy essentially as $`E^2`$. The accuracy of eq. (15) also improves with energy as $`E^2`$. This is illustrated by panels (a) and (b) of fig. 1, which show $`P_2`$ as a function of neutrino energy for several trajectories inside the Earth and as a function of the zenith angle of the neutrino trajectory $`\mathrm{\Theta }`$ for several neutrino energies. One can see that already for $`E8`$ MeV the accuracy of our analytic formula is extremely good. Note that when neutrinos do not cross the Earth’s core ($`\mathrm{cos}\mathrm{\Theta }>0.837`$) and so experience a slowly changing potential $`V(x)`$, the accuracy of the approximation (15) is very good even in the MSW resonance region $`E`$5 – 8 GeV. The accuracy of eq. (15) is also good for energies below $`2`$ GeV (not shown in the figure); however, in this region the domain of the applicability of (15) is relatively narrow, since for $`E0.5`$ GeV the oscillations driven by the “solar” parameters ($`\mathrm{\Delta }m_{21}^2,\theta _{12}`$) can no longer be neglected.
To understand the remarkable accuracy of eq. (15), we find the correction $`\mathrm{\Delta }P_2`$ to the transition probability in (15) emerging in the next nontrivial order in $`H_I`$. Note that from the above considerations one can expect $`\mathrm{\Delta }P_2/P_2`$ to be proportional to $`\mathrm{sin}2\theta _m`$. Furthermore, as can be easily seen, for uniform matter eq. (15) reproduces the exact transition probability; therefore one expects $`\mathrm{\Delta }P_2/P_2V^{}`$. A straightforward calculation indeed gives
$$\frac{\mathrm{\Delta }P_2}{P_2}\mathrm{sin}^22\theta _m\frac{\mathrm{\Delta }V}{4\pi \omega }\mathrm{sin}^22\theta _m\frac{\mathrm{\Delta }V}{2\pi V},$$
(17)
where $`\mathrm{\Delta }V`$ is the change of the potential over the oscillation length $`\pi /\omega `$, and the last equality holds in the high energy regime. For slowly changing density this is equivalent to
$$\frac{\mathrm{\Delta }P_2}{P_2}\mathrm{sin}^22\theta _m\frac{V^{}}{4\omega ^2}.$$
(18)
Introducing the adiabaticity parameter $`\gamma =4\pi \omega /(\mathrm{sin}2\theta _m\mathrm{\Delta }V)`$, we find that $`\mathrm{\Delta }P_2/P_2\mathrm{sin}2\theta _m\gamma ^1`$, and therefore for small mixing in matter our approximation is better than the adiabatic one. At the same time, for $`\mathrm{\Delta }V/4\pi \omega <1`$ it is better than the simple expansion in powers of $`\mathrm{sin}^22\theta _m`$.
The matter density profile of the Earth satisfies $`V^{}/V^20.5`$, and therefore for neutrino oscillations in the Earth our approximation is expected to work well when $`\mathrm{sin}^22\theta _m1`$. This is fulfilled in the high energy (or, equivalently, high density) limit $`EV/\mathrm{\Delta }m^2\mathrm{cos}2\theta `$, i.e. above the MSW resonance. If the vacuum mixing angle is small (i.e. $`\theta =\theta _{13}`$), our expansion parameter is also small below the resonance. The above formalism applies in this low energy case as well, with only minor modifications: the sign of $`H_0`$ in (5) has to be flipped, and correspondingly one has to replace $`\omega \omega `$ in eq. (7). These modifications are necessary because of the interchange of the eigenvalues of the Hamiltonian upon crossing the MSW resonance. Expressions for the transition probability in eqs. (15), (16) remain unchanged. Thus, our results are in general valid outside the MSW resonance region, which for small $`\theta `$ is very narrow. For the non-resonant channels ($`\overline{\nu }`$ channels for $`\mathrm{\Delta }m^2>0`$ or $`\nu `$ channels for $`\mathrm{\Delta }m^2<0`$) and small vacuum mixing our formulas are valid in the whole diapason of energies because $`\mathrm{sin}2\theta _m`$ is always small.
If $`\theta _{13}`$ is very small or vanishes, $`\nu _e\nu _{\mu ,\tau }`$ oscillations are driven by $`\mathrm{\Delta }m^2=\mathrm{\Delta }m_{21}^2`$ and the large mixing angle $`\theta =\theta _{12}`$. The oscillation probabilities can then be expressed through another effective two-flavor probability, $`\stackrel{~}{P}_2\stackrel{~}{P}_2(\mathrm{\Delta }m_{21}^2,\theta _{12},V(x))`$, in terms of which the flavor transition probabilities are $`P(\nu _e\nu _\mu )=\mathrm{cos}^2\theta _{23}\stackrel{~}{P}_2`$, $`P(\nu _e\nu _\tau )=\mathrm{sin}^2\theta _{23}\stackrel{~}{P}_2`$ PerSm . For $`EV/\mathrm{\Delta }m_{21}^2\mathrm{cos}2\theta _{12}`$ (which for the typical densities inside the Earth corresponds to $`E0.5`$ GeV), the probability $`\stackrel{~}{P}_2`$ is very well approximated by eq. (15).
Let us consider the case of symmetric matter. Integrating (16) by parts, one finds
$$P_2=\mathrm{sin}^22\theta _m^0\left[\mathrm{sin}\varphi _L+\omega _0_0^{L/2}𝑑z\frac{d\omega }{dz}\frac{1}{\omega ^2}\mathrm{sin}2\varphi (z)\right]^2,$$
(19)
where $`\theta _m^0\theta _m(V_0)`$, $`\omega _0\omega (V_0)`$, $`V_0`$ being the potential at the initial and final points of the neutrino trajectory, and $`\varphi _L`$ is the adiabatic phase acquired along the entire neutrino path. If the potential changes slowly with distance, so that $`\omega ^2d\omega /dz1`$, the second term in (19) can be neglected, and $`P_2`$ reduces to the usual adiabatic probability in symmetric matter: $`P_{\mathrm{adiab}}=\mathrm{sin}^22\theta _m^0\mathrm{sin}^2\varphi _L`$. The second term in (19) describes the effects of violation of adiabaticity.
Let us apply the above results to neutrino beams (atmospheric, accelerator, cosmic neutrinos) crossing the Earth. According to the PREM model PREM , the Earth density profile can be described as several spherical shells of radii $`R_i`$ with rather smooth density change within the shells and sharp change at the borders between them. Then, along a direction from the center of the Earth outwards, $`\omega (z)`$ decreases abruptly from $`\omega _i^+`$ to $`\omega _i^{}`$ in very narrow regions around $`R_i`$. Therefore $`d\omega /dz`$ is large in these narrow regions and small outside them. The integration in (19) can then be easily done, leading to
$$P_2=\mathrm{sin}^22\theta _m^0\left[\mathrm{sin}\varphi _L\omega _0\mathrm{\Sigma }_i\frac{\omega _i^+\omega _i^{}}{\omega _i^+\omega _i^{}}\mathrm{sin}2\varphi _i\right]^2.$$
(20)
Here $`\varphi _i`$ is the adiabatic phase acquired by neutrinos between the points $`z=0`$ and $`z=R_i`$. We will use eq. (20) for interpreting the results of our calculations.
3. Parametric enhancement of oscillations. In the PREM model of the Earth density all the density jumps between different shells except those between the mantle and core are relatively small PREM . Therefore the density profile felt by neutrinos crossing the core of the Earth can be approximated by three layers (mantle \- core - mantle) with slow density changes within the layers. Eq. (20) then gives
$$P_2=\mathrm{sin}^22\theta _m^0\left[\mathrm{sin}(\varphi _c+2\varphi _m)\frac{\omega _0}{\omega _m}\left(1\frac{\omega _m}{\omega _c}\right)\mathrm{sin}\varphi _c\right]^2,$$
(21)
where $`\omega _m`$ and $`\omega _c`$ are the values of $`\omega (x)`$ in the mantle and core on the respective sides of their border, and $`\varphi _m`$ and $`\varphi _c`$ are the phases acquired in the mantle (one layer) and core. Eq. (21) corresponds to the adiabatic neutrino propagation inside the core and mantle and violation of the adiabaticity at the borders between them. In the approximation of constant effective densities in the mantle and core we have $`\omega _0=\omega _m^{\mathrm{eff}}`$, and the non-adiabatic term is proportional to $`(1\omega _c^{\mathrm{eff}}/\omega _m^{\mathrm{eff}})`$.
For neutrino trajectories that cross the mantle only ($`\varphi _c=0`$), eq. (21) reduces to the adiabatic probability. The passage of neutrinos through the core can lead to an enhancement of the oscillations. As follows from (21), the maximum enhancement of the probability can be achieved when $`\mathrm{sin}(\varphi _c+2\varphi _m)`$ and $`\mathrm{sin}\varphi _c`$ are of opposite sign and maximal amplitude:
$$\mathrm{sin}\varphi _c=\mathrm{sin}(\varphi _c+2\varphi _m)=\pm 1,$$
(22)
that is, when
$$\varphi _c=\pm \frac{\pi }{2}+2\pi n,\varphi _m=\pm \frac{\pi }{2}+2\pi k.$$
(23)
Here $`n`$ and $`k`$ are integers, and the signs in front of $`\pi /2`$ are correlated. In this case the enhancement factor is
$$\frac{P_2^{max}}{\mathrm{sin}^22\theta _m}=\left[1+\frac{\omega _0}{\omega _m}\left(1\frac{\omega _m}{\omega _c}\right)\right]^2\left(2\frac{V_m}{V_c}\right)^22.5,$$
(24)
where $`\mathrm{sin}^22\theta _m`$ in the denominator corresponds to the maximum possible transition probability for neutrinos crossing only the mantle, and we have taken into account that at high energies $`\omega _m/\omega _cV_m/V_c`$.
The condition in eq. (23) and the enhancement described by eq. (24) are the particular cases of the parametric resonance condition and the parametric enhancement of neutrino oscillations ETC ; Akh1 . In Akh2 it was shown that in the case of matter consisting of alternating layers of two different constant densities and (in general) different widths the parametric resonance condition is
$$X_3(\mathrm{sin}\varphi _m\mathrm{cos}\varphi _c\mathrm{cos}2\theta _m+\mathrm{cos}\varphi _m\mathrm{sin}\varphi _c\mathrm{cos}2\theta _c)=0,$$
(25)
where $`\theta _{m,c}`$ and $`\varphi _{m,c}`$ are the mixing angles and the acquired oscillations phases in the layers $`m`$ and $`c`$. This condition can also be used as an approximate one when matter density inside the layers is not constant but varies sufficiently slowly. For neutrino oscillations in the Earth we identify the layers $`m`$ and $`c`$ with the mantle and core. Since in the energy region $`\mathrm{sin}2\theta \delta V`$ one has $`\theta _m\theta _c\pi /2`$, condition (25) reduces to
$$X_3\mathrm{sin}(\varphi _m+\varphi _c)=0.$$
(26)
Eq. (23) is a particular realization of this condition. In the high energy limit the parametric resonance condition (23) was previously considered in the framework of active-sterile atmospheric neutrino oscillations in Liu . It should be noted that, while the realization (23) of condition (25) corresponds to the maximal possible parametric enhancement of oscillations of high energy neutrinos in the Earth, a sizable amplification is also possible if the equality $`X_3=0`$ is realized differently, i.e. when the two terms in (25) do not separately vanish but cancel each other note2 .
The parametric resonance conditions in eq. (25) or (23) require a subtle correlation between the matter density of the Earth, the distances that neutrinos travel in the Earth’s mantle and core (which are not independent), and also in general neutrino energy and oscillation parameters. Therefore it is far from obvious that these conditions can actually be satisfied. Amazingly, this is indeed the case Liu ,Pet ,Akh2 .
Our present analysis shows that, for $`\nu _e\nu _{\mu ,\tau }`$ oscillations of high energy neutrinos in the Earth, there are two regions of the zenith angles of neutrino trajectories in which (25) is satisfied: $`\mathrm{cos}\mathrm{\Theta }=1÷0.93`$ and $`\mathrm{cos}\mathrm{\Theta }=0.88÷0.84`$. This is illustrated in panel (b) of fig. 1. Since at high energies matter suppresses neutrino mixing, one could expect that for the trajectories crossing the core ($`\mathrm{cos}\mathrm{\Theta }<0.837`$), where the densities are higher, the probability would be suppressed or at least would not change. Instead, we see two prominent peaks there, exceeding maximal allowed by the MSW effect value of probability $`\mathrm{sin}^22\theta _m`$ by up to a factor of 2. This is the result of the parametric enhancement of neutrino oscillations. As can be seen from the figure, for core-crossing trajectories the positions of the peaks of $`P_2`$ essentially coincide with zeros of $`X_3`$.
For neutrino energies $`E10`$ – 15 GeV, the oscillation phases corresponding to the peak with $`\mathrm{cos}\mathrm{\Theta }0.93`$ (the inner peak) are $`\varphi _m\pi /4`$, $`\varphi _c7\pi /4`$, while for the peak with $`\mathrm{cos}\mathrm{\Theta }=0.88÷0.84`$ (the outer peak), they are $`\varphi _m0.35\pi `$, $`\varphi _c0.65\pi `$. In both peaks to a good accuracy $`\varphi _m+\varphi _c=n\pi `$, so that eq. (26) is satisfied. The phases in the outer peak are closer to the realization (23) of the parametric resonance condition, and therefore in this peak the parametric enhancement of oscillations is closer to the maximal possible one. From panel (c) of fig. 1 one can see that at $`E21`$ GeV the maximum enhancement condition (23) can be exactly realized in the outer peak. For neutrinos of very high energies ($`E100`$ GeV), it can be realized nearly exactly in the inner peak, whereas the outer peak becomes slimmer and lower, and at energies $`E300`$ GeV virtually disappears.
4. In conclusion – We have derived very simple and accurate integral formulas describing neutrino oscillations in matter with arbitrary density profiles. They can be applied to all possible situations where the effective mixing in matter is small. In particular, our results can be applied to atmospheric, accelerator and cosmic neutrinos crossing the Earth as well as to cosmic neutrinos crossing the Sun. They can be used for probing the effects of small jumps in the density profile of the Earth as well as evaluating the effects of uncertainties in this profile on the interpretation of the results of future accelerator experiments. They can also be employed for studying the effects of proper averaging over the energy spectrum of the neutrino beam. We used the obtained formulas to study the parametric enhancement of oscillations of high energy neutrinos in the Earth and identified two peaks in the zenith angle distribution of core-crossing neutrinos, which are due to the parametric effects. Observation of these peaks in future high statistics experiments with high energy neutrinos would provide an unambiguous evidence for the parametric resonance effects in neutrino oscillations in matter.
This work was supported in part by SFB-375 für Astro-Teilchenphysik der Deutschen Forschungsgemeinschaft (EA), the National Science Foundation grant PHY0354776 (MM) and by the Alexander von Humboldt Foundation (AS).
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# The Average Shape of Transport-Limited Aggregates
## Abstract
We study the relation between stochastic and continuous transport-limited growth models, which generalize conformal-mapping formulations of diffusion-limited aggregation (DLA) and viscous fingering, respectively. We derive a nonlinear integro-differential equation for the asymptotic shape (average conformal map) of stochastic aggregates, whose mean-field approximation is the corresponding continuous equation, where the interface moves at its local expected velocity. Our equation accurately describes advection-diffusion-limited aggregation (ADLA), and, due to nonlinear averaging over fluctuations, the average ADLA cluster is similar, but not identical, to an exact solution of the mean-field dynamics. Similar results should apply to all models in our class, thus explaining the known discrepancies between average DLA clusters and viscous fingers in a channel geometry.
Developing effective mean field approximations to nonlinear stochastic equations constitutes a major challenge in various active fields of statistical physics, e.g. hydrodynamic turbulence Frisch and self organized criticality Bak . Straightforward derivation of such approximate theories typically involves uncontrolled assumptions required to ”close” an infinite hierarchy of equations for moments of the probability distribution of the stochastic field. An alternative approach consists of deriving asymptotic solutions to a deterministic version of the original stochastic dynamics, assuming that such solutions capture the behavior of ensemble average of the original stochastic field Bak . Such approach, however, might turn out to be unreliable as well, since it underestimates the possible effects of noise on the asymptotic evolution of a stochastic field.
A nontrivial example in which such approach has been advanced over the last two decades is the fractal morphology of patterns observed in computer simulations of the celebrated diffusion limited aggregation (DLA) model witten81 . Since the relation between the mathematical formulations of the stochastic DLA process and the continuous viscous fingering dynamics was established Paterson84 , the striking similarity between patterns observed in both processes has triggered various attempts to interpret viscous fingering dynamics as a mean field for DLA stanley91 ; arneodo89 ; barra02 ; swinney03 .
In this Letter, we study the connection between a broad class of stochastic transport-limited aggregation processes and their continuous counterparts bazant03 . In our models, growth is fuelled by nonlinear, non-Laplacian transport processes, such as advection-diffusion and electrochemical conduction, which satisfy conformally invariant equations bazant04 . Stochastic and continuous dynamics are defined by generalizing conformal-mapping formulations of DLA hastings98 and viscous fingering polubarinova45 ; shraiman84 , respectively. We show that the continuous dynamics is a self-consistent mean-field approximation of the stochastic dynamics, which, nevertheless, does not accurately predict the average shape of a random ensemble of aggregates.
We consider a set of two-dimensional scalar fields, $`𝝋=\{\phi _1,\phi _2,\mathrm{},\phi _M\}`$, whose gradients produce quasi-static, conserved “flux densities”,
$$𝐅_i=\underset{j=1}{\overset{M}{}}C_{ij}(𝝋)\mathbf{}\phi _j,\mathbf{}𝐅_i=0$$
(1)
in $`\mathrm{\Omega }_z(t)`$, the exterior of a singly connected domain that represents a growing aggregate at time $`t`$. (The coefficients, $`\{C_{ij}\}`$ may be nonlinear functions of the fields.) The crucial property of the nonlinear system (1) is its conformal invariance bazant04 : If $`𝝋(w,\overline{w})`$, not necessarily harmonic, is a solution in a domain $`\mathrm{\Omega }_w`$ and $`w=f(z)`$ is a conformal map from $`\mathrm{\Omega }_z`$ to $`\mathrm{\Omega }_w`$, then $`𝝋(f(z),\overline{f(z)})`$ is a solution in $`\mathrm{\Omega }_z`$. Using this fact, the evolving domain, $`\mathrm{\Omega }_z(t)`$, can be described by the conformal map, $`z=g(w,t)`$, from the exterior of (say) the unit disk, $`\mathrm{\Omega }_w`$.
Growth is driven by a combination of flux densities, $`𝐐=_{i=1}^NB_i(𝝋)𝐅_i`$, on the boundary with a local growth rate, $`\sigma =\widehat{𝒏}𝐐`$, where $`\widehat{𝒏}`$ is the unit normal vector at $`z\mathrm{\Omega }_z(t)`$. For continuous, deterministic growth, each boundary point $`z`$ moves with a velocity, $`𝒗(z)=\alpha \sigma (z)\widehat{𝒏}(z)`$, where $`\alpha `$ is a constant. For discrete, stochastic growth, the initial seed, $`\mathrm{\Omega }_z(t_0=0)`$, is iteratively advanced by elementary “bump” maps representing particles of area $`\lambda _0`$ at times $`t_1,\mathrm{},t_n`$. The waiting time $`t_nt_{n1}`$ is an exponential random variable with a mean set by the total integrated flux. The probability density to add the $`n`$th particle in a boundary element $`(z,z+dz)\mathrm{\Omega }_z(t_{n1})`$ is proportional to $`\sigma (z)|dz|`$.
The classical models are recovered in the simplest case of one field ($`M=1`$). DLA corresponds to stochastic growth by diffusion, $`𝐅=𝐐=D\mathbf{}c`$, from a distant source ($`c\mathrm{log}|z|`$ as $`|z|\mathrm{}`$) to an absorbing cluster ($`c=0`$ for $`z\mathrm{\Omega }_z(t)`$), where $`c`$ is the particle concentration and $`D`$ the diffusivity. Viscous fingering corresponds to continuous growth by the same process, where $`c`$ becomes the fluid pressure and $`D`$ the permeability.
The simplest, nontrivial models with multiple fields ($`M=2`$) involve diffusion in a fluid flow. The stochastic case is advection-diffusion-limited aggregation (ADLA) bazant03 , illustrated in Fig. 1a. Particles are deposited around a circular seed of radius, $`L_o`$, from potential flow, $`𝒗=\phi `$, of uniform velocity $`U`$ far from the aggregate. The dimensionless transport problem is
$$\mathrm{Pe}_o\mathbf{}\phi \mathbf{}c=\mathbf{}^2c,\mathbf{}^2\phi =0,z\mathrm{\Omega }_z(t)$$
(2)
$$c=0,\widehat{𝒏}\mathbf{}\phi =0,\sigma =\widehat{𝒏}\mathbf{}c,z\mathrm{\Omega }_z(t)$$
(3)
$$c1,\mathbf{}\phi \widehat{𝒙},|z|\mathrm{},$$
(4)
where $`c`$ is the concentration of particles. Here $`x`$, $`\phi `$, $`c`$, and $`\sigma `$ are in units of $`L_o`$, $`UL_o`$, $`C`$, and $`DC/L_o`$, respectively, and $`\mathrm{Pe}_o=UL_o/D`$ is the initial Péclet number. Numerical solutions and asymptotic approximations are discussed in Ref. choi04 .
The transport problem is conformally invariant, except for the boundary condition, Eq. (4), which alters the flow speed upon conformal mapping. Instead, we choose to fix the mapped background flow and replace $`\mathrm{Pe}_o`$ with the renormalized Péclet number, $`\mathrm{Pe}(t)=A_1(t)\mathrm{Pe}_o`$, when Eq. (2) is transformed from $`\mathrm{\Omega }_z(t)`$ to $`\mathrm{\Omega }_w`$. The “conformal radius”, $`A_1(t)`$, is the coefficient of $`w`$ in the Laurent expansion of $`g(w,t)`$ and scales with the radius of the growing object hastings98 ; davidovitch99 . Since $`A_1(t)\mathrm{}`$ for any initial condition, the flux approaches a self-similar form,
$$\sigma (\theta ;\mathrm{Pe})2\sqrt{\mathrm{Pe}/\pi }\mathrm{sin}(\theta /2)\text{as}\mathrm{Pe}\mathrm{}.$$
(5)
More generally, there is a universal crossover from DLA ($`\sigma =`$ constant) to this stable fixed point, where $`\mathrm{Pe}(t)=A_1(t)\mathrm{Pe}_o`$ is the appropriate scaling variable choi05 .
The continuous analog of ADLA is a simple model for solidification from a flowing melt kornev94 . More generally, continuous dynamics in our class is described by a nonlinear equation,
$$\mathrm{Re}(\overline{wg^{}}g_t)=\alpha \sigma (w;\mathrm{Pe}(t))\text{for}|w|=1.$$
(6)
which generalizes the Polubarinova-Galin equation for Laplacian growth polubarinova45 ; shraiman84 ($`\sigma =1`$). In the case of advection-diffusion kornev94 , only low-$`\mathrm{Pe}`$ approximations are known, but we have found an exact high-$`\mathrm{Pe}`$ solution of the form, $`g(w,t)=A_1(t)G_c(w)`$, where
$$A_1(t)=t^{2/3},G_c(w)=w\sqrt{11/w}.$$
(7)
This similarity solution to Eq. (6) with $`\alpha \sigma (\theta ,t)=\sqrt{A_1(t)}\mathrm{sin}(\theta /2)`$ describes the long-time limit, according to Eq. (5). (We do not know the uniqueness or stability of this solution or whether it can be approached without singularities from general initial conditions.) Just as the Saffman-Taylor finger solution (for $`\sigma =1`$) has been compared to DLA in a channel geometry somfai\_ball02 , this analytical result begs comparison with ADLA.
As in Ref. bazant03 , we grow ADLA clusters by a modified Hastings-Levitov algorithm hastings98 . The random attachment of the $`n`$th particle to the cluster is described by perturbing the boundary $`\mathrm{\Omega }_z(t_{n1})\mathrm{\Omega }_z(t_n)`$ by a “bump” of characteristic area $`\lambda _0`$. This leads to the recursive dynamics
$$g_n(w)=g_{n1}\varphi _{\lambda _n,\theta _n}(w),g_n(w)=g(w,t_n)$$
(8)
where $`\varphi _{\lambda ,\theta }(w)`$ is a specific map, conformal in $`\mathrm{\Omega }_w`$, that slightly distorts the unit circle by a bump of area $`\lambda `$ around the angle $`\theta `$. The parameter, $`\lambda _n=\lambda _0|g_{n1}^{}(e^{i\theta _n})|^2`$ is the area of the pre-image of such bump under the inverse map $`g^1`$. The angle $`\theta _n`$ is chosen with a probability density $`p(\theta ;\mathrm{Pe}(t_n))\sigma (e^{i\theta };\mathrm{Pe}(t_n))`$, so the expected growth rate is the same as in the continuous dynamics.
For an ensemble of $`n`$-particle aggregates, a natural definition of average cluster shape is the conformal map, $`G_n(w)`$, defined by averaging at a point $`w\mathrm{\Omega }_w`$ all the maps, $`G_n(w)=g_n(w)/A_1(t_n)`$, rescaled to have a unit conformal radius. We then ask: *What is the limiting average cluster shape, $`G_{\mathrm{}}(w)=lim_n\mathrm{}G_n(w)`$, and how does it compare to the similarity solution, $`G_c(w)`$, of the continuous growth equation (6)?* The same questions apply to any of our transport-limited growth models, but here we focus on ADLA as a representative case.
To provide numerical evidence, we grow 2000 ADLA clusters of size $`n=10^5`$ using the semi-circular bump function in Ref. davidovitch99 (with $`a=1/2`$). To reduce fluctuations, we aggregate small particles, $`\lambda _0=1/16`$, on a large initial seed ($`g_0(w)=w,|w|=1`$). To reach at the asymptotic limit faster and also match the assumption of $`G_c(w)`$, we fix the angular probability measure, $`p_{\mathrm{}}(\theta )=\mathrm{sin}(\theta /2)/4`$ for $`\mathrm{Pe}=\mathrm{}`$, throughout the growth. In Fig. 2a, we plot the average contour of the ensemble, $`G_n(e^{i\theta })`$ at $`n=10^5`$ along with that of the continuous solution, $`G_c(e^{i\theta })`$. To give a sense of fluctuations, we also plot a “cloud” of points, $`G_n(e^{i\theta })`$ over uniformly sampled values of $`\theta `$. Fig. 2b is the zoom-in of the boxed region in Fig. 2a, where we also show $`G_n(e^{i\theta })`$ at $`n=10^3`$ and $`n=10^4`$. Although the convergence of $`G_n(w)`$ is easily extrapolated, the $`n=10^5`$ line has not reached at the asymptotic limit yet. The branch point at $`w=1`$ seems to be related to the slow convergence. Ignoring the unconverged area, $`G_n(w)`$ and $`G_c(w)`$ are quite similar, and yet clearly not same. The average, $`G_n(w)`$, better captures the ensemble morphology reflected by the cloud pattern than $`G_c(w)`$ and the opening angles at the branch point of the two curves are also different choi05 .
Next we derive an equation for $`G_n(w)`$ in the asymptotic regime. For growing aggregates $`\lambda _n0`$ as $`n\mathrm{}`$ davidovitch99 . Following Hastings hastings97 , we use Eq. (8) to derive a linearized recursive equation for $`G_{n+1}(w)`$ for $`|we^{i\theta _{n+1}}|\sqrt{\lambda _{n+1}}`$. Letting $`(\lambda ,\theta )`$ denote the parameters of the $`(n+1)`$th bump, we obtain:
$`G_{n+1}`$ $`(w)(1a\lambda )(G_n(w)+a\lambda H_\theta (w)G_n^{}(w))`$ (9)
$`G_n(w)+a\lambda (H_\theta (w)G_n^{}(w)G_n(w)).`$
where $`H_\theta (w)=w(w+e^{i\theta })/(we^{i\theta })`$ and we use $`A_1(t_{n+1})=(1+\lambda )^aA_1(t_n)`$.
Stationarity of the ensemble of rescaled clusters implies:
$$G_n(w)=_0^{2\pi }𝑑\theta p_{\mathrm{}}(\theta )G_{n+1}(w).$$
(10)
Our analysis applies for conformally invariant transport-limited growth from an isolated seed with general angular probability distributions, although we will focus on the case of ADLA, $`p_{\mathrm{}}(\theta )=\mathrm{sin}(\theta /2)/4`$.
Using Eq. (9), we get the fixed-point condition:
$$_0^{2\pi }𝑑\theta p_{\mathrm{}}(\theta )\lambda G_{\mathrm{}}(w)_0^{2\pi }𝑑\theta p_{\mathrm{}}(\theta )\lambda G_{\mathrm{}}^{}(w)H_\theta (w).$$
(11)
To facilitate further analysis, we approximate the left hand side of Eq. (11) as
$$_0^{2\pi }𝑑\theta p_{\mathrm{}}\lambda G_{\mathrm{}}_0^{2\pi }𝑑\theta p_{\mathrm{}}\lambda G_{\mathrm{}}$$
(12)
and the right hand side similarly. We checked the validity of this assumption by numerical evaluation of these two quantities for increasing values of $`n`$, finding less than 1% discrepancy for the largest clusters ($`n=10^5`$). Although the stronger assumption, $`\lambda G_{\mathrm{}}\lambda G_{\mathrm{}}`$, is not valid, particularly near the bump center $`w=e^{i\theta }`$, the correlation seems to be canceled out in the integration along the angle.
With these assumptions, we arrive at a nonlinear integro-differential equation for $`G_{\mathrm{}}(w)`$, the limiting average cluster shape:
$`{\displaystyle \frac{G_{\mathrm{}}(w)}{G_{\mathrm{}}(w)^{}}}`$ $`={\displaystyle \frac{_0^{2\pi }𝑑\theta p_{\mathrm{}}(\theta )|G_{\mathrm{}}^{}(e^{i\theta })|^2H_\theta (w)}{_0^{2\pi }𝑑\theta p_{\mathrm{}}(\theta )|G_{\mathrm{}}^{}(e^{i\theta })|^2}}`$ (13)
$`={\displaystyle _0^{2\pi }}𝑑\theta p_{\mathrm{}}(\theta )\mathrm{\Lambda }_{\mathrm{}}(\theta )H_\theta (w)`$ (14)
where we introduce a conditional probability density,
$$\mathrm{\Lambda }_n(\theta )=\frac{|G_n^{}(e^{i\theta })|^2}{_0^{2\pi }𝑑\theta ^{}p_{\mathrm{}}(\theta ^{})|G_n^{}(e^{i\theta ^{}})|^2},$$
(15)
proportional to the average local size of a bump’s pre-image (Jacobian factor), $`\lambda _n`$, at angle $`\theta `$. In deriving Eq. (13) we assume that $`A_1(t_n)|G_n^{}|^2A_1(t_n)|G_n^{}|^2`$ as $`n\mathrm{}`$ choi05 . The stochastic nature of the aggregates is manifested through the two different averages in Eqs. (14)–(15).
To check the validity of Eq. (14), we obtain $`\mathrm{\Lambda }_{\mathrm{}}`$ from simulations and solve for $`G_{\mathrm{}}(w)`$. As shown in Fig. 3, the measured curves for $`\mathrm{\Lambda }_n(\theta )`$ for $`n=10^3,\mathrm{\hspace{0.17em}10}^4`$ and $`10^5`$ are nearly identical, so we conclude that $`\mathrm{\Lambda }_{10^5}(\theta )`$ is a good approximation of $`\mathrm{\Lambda }_{\mathrm{}}(\theta )`$. Now we solve Eq. (14) by expanding $`G_{\mathrm{}}(w)`$ by a Laurent series and finding recurrence relations for the coefficients, which involve integrals of $`\mathrm{\Lambda }_{\mathrm{}}(\theta )`$. We calculate 200 first coefficients and reconstruct $`G_{\mathrm{}}(w)`$. The image of the unit circle under $`G_{\mathrm{}}(w)`$, shown in Fig.2a (thick gray), is in excellent agreement with the converging pattern of $`G_n(w)`$.
The surprising difference between the convergence rates of the average Jacobian $`\mathrm{\Lambda }_n(\theta )`$ and the average map itself $`G_n(w)`$ is intimately related to the multifractal nature of the distribution of the stretching factor (harmonic measure) $`|G_n^{}(e^{i\theta })|`$. Since this factor is very large around the cusp at $`\theta =0`$, which is dominant during the growth process, fluctuations at the cusp do not contribute to negative moments of the distribution of $`|G_n^{}(e^{i\theta })|`$ and thus negative moments converge much faster than positive ones, and faster than the average map itself. Since $`\mathrm{\Lambda }_n(\theta )`$ comes from averaging $`|G_n^{}(e^{i\theta })|^2`$, this observation explains its fast convergence. This argument illustrates how the two averages interact in Eqs. (14)–(15) and suggests that the faster convergence of $`\mathrm{\Lambda }_n(\theta )`$ dominates the morphology.
With the validity of Eq. (14) established, we may consider its mean-field version, where the ensemble average is replaced by a single conformal map, given by:
$$\frac{G_c(w)}{G_c^{}(w)}=_0^{2\pi }𝑑\theta p_{\mathrm{}}(\theta )\mathrm{\Lambda }_c(\theta )H_\theta (w),$$
(16)
$$\mathrm{\Lambda }_c(\theta )=\frac{|G_c^{}(e^{i\theta })|^2}{_0^{2\pi }𝑑\stackrel{~}{\theta }p_{\mathrm{}}(\stackrel{~}{\theta })|G_c^{}(e^{i\stackrel{~}{\theta }})|^2}.$$
(17)
Not surprisingly, the similarity solution for continuous growth, Eq. (7), is an exact solution of Eqs. (16)–(17). In fact, it is possible to derive Eqs. (16)–(17) from a different representation of Eq. (6), which has been done for the case of DLA, $`p_{\mathrm{}}(\theta )=1/2\pi `$, in Ref. hastings97 . Elsewhere choi05 , we obtain an analytical form for $`\mathrm{\Lambda }_c(\theta )`$ for ADLA, which is plotted in Fig. 3 (thick gray). A small, but significant, difference between $`\mathrm{\Lambda }_c(\theta )`$ and $`\mathrm{\Lambda }_{\mathrm{}}(\theta )`$ is apparent, especially at $`\theta =\pi /4`$ and $`7\pi /4`$.
The solution in Eq. (7) can be interpreted as a self-consistent mean-field approximation for the average conformal map, $`G_{\mathrm{}}(w)`$. However, fluctuations in the ensemble manifest themselves through the different averages in Eqs. (14)–(15). As long as $`|G_{\mathrm{}}^{}(w)|^2`$ is different from $`|G_{\mathrm{}}(w)^{}|^2`$, $`\mathrm{\Lambda }_{\mathrm{}}(\theta )\mathrm{\Lambda }_c(\theta )`$, and thus the deviation of $`G_{\mathrm{}}(w)`$ from $`G_c(w)`$ is inevitable.
We believe that the assumptions leading to Eq. (14) are quite general, and not specific to ADLA, so the continuous dynamics should be a mean field theory (in this sense) for any stochastic aggregation, driven by conformally invariant transport processes, Eq. (1). We conclude, therefore, that the solution to the continuous dynamics, although similar, is not identical to the ensemble-averaged cluster shape. An exceptional case is DLA in radial geometry, where isotropy implies the trivial solution, $`G_{\mathrm{}}(w)=w`$ and $`\mathrm{\Lambda }_c(\theta )=1`$. Clearly, $`G_c(w)=w`$ and $`\mathrm{\Lambda }_c(\theta )=1`$ solves Eq. (14) with $`p_{\mathrm{}}(\theta )=1/2\pi `$.
We expect, however, that this identity between the mean-field approximation and the average shape of stochastic clusters will be removed with any symmetry breaking, either in the model equations (such as ADLA) or in the BCs (e.g. DLA in a channel). This result is consistent with recent simulations of DLA in a channel geometry somfai\_ball02 , which show that the average cluster shape, $`G_n(w)`$, is similar, but not identical, to any of the Saffman-Taylor “fingers”, which solve the continuous dynamics. We expect that an analogous equation to Eq. (14), relating $`G_{\mathrm{}}(w)`$ and $`|G_{\mathrm{}}^{}(e^{i\theta })|^2`$, will hold in a channel geometry, and Saffman-Taylor fingers should be exact solutions to the mean-field approximation of that equation.
We conclude by emphasizing that, although Eq. (14) is a necessary condition for the average shape of transport-limited aggregates in the class, Eq. (1), it does not provide a basis for complete statistical theory. Such a theory would likely consist of an infinite set of independent equations connecting a hierarchy of moments of the multifractal distributions of maps $`\{G_{\mathrm{}}(w)\}`$, and derivatives $`\{G_{\mathrm{}}^{}(w)\}`$. Multifractality may speed up convergence, as for Eq. (14), or slow down convergence of other equations in this set. The mean-field approximation, Eq. (16), which corresponds to the continuous growth process, can be considered as leading a hierarchy of closure approximations to this set.
This work was supported in part by the MRSEC program of the National Science Foundation under Award No. DMR 02-13282 (M.Z.B), and by Harvard MRSEC (B.D).
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# Skew-product for group-valued edge labellings of Bratteli diagrams
## 1.1. Definition.
A Bratteli diagram is an infinite directed graph $`(V,E)`$, where $`V`$ is the vertex set and $`E`$ is the edge set. Both $`V`$ and $`E`$ are partitioned into non-empty disjoint finite sets
$$V=V_0V_1V_2\mathrm{}\mathrm{and}E=E_1E_2\mathrm{}$$
There are two maps $`r,s:EV`$ the range and source maps. The following properties hold:
1. $`V_0=\{v_0\}`$ consists of a single point, referred to as the ‘top vertex’ of the Bratteli diagram
2. $`r(E_n)V_n,s(E_n)V_{n1},n=1,2,\mathrm{}`$. Also $`s^1(v)\varphi `$ $`vV`$ and $`r^1(v)\varphi `$ for all $`vV_1,V_2,\mathrm{}`$.
Maps between Bratteli diagrams are assumed to preserve gradings and intertwine the range and source maps. If $`vV_n`$ and $`wV_m`$, where $`m>n`$, then a path from $`v`$ to $`w`$ is a sequence of edges $`(e_{n+1},\mathrm{},e_m)`$ such that $`s(e_{n+1})=v,r(e_m)=w`$ and $`s(e_{j+1})=r(e_j)`$. Infinite paths from $`v_0V_0`$ are defined similarly. The Bratteli diagram is called simple if for any $`n=0,1,2,\mathrm{}`$ there exists $`m>n`$ such that every vertex of $`V_n`$ can be joined to every vertex of $`V_m`$ by a path.
## 1.2. Definition.
An ordered Bratteli diagram $`(V,E,)`$ is a Bratteli diagram $`(V,E)`$ together with a linear order on $`r^1(v),vV\{v_0\}=V_1V_2V_3\mathrm{}`$. We say that an edge $`eE_n`$ is a maximal edge (resp. minimal edge) if $`e`$ is maximal (resp. minimal) with respect to the linear order in $`r^1(r(e))`$.
Given $`vV_n`$, it is easy to see that there exists a unique path $`(e_1,e_2,\mathrm{},e_n)`$ from $`v_0`$ to $`v`$ such that each $`e_i`$ is maximal (resp. minimal).
Note that if $`m>n`$, then for any $`wV_m`$, the set of paths starting from $`V_n`$ and ending at $`w`$ obtains an induced (lexicographic) linear order:
$$(e_{n+1},e_{n+2},\mathrm{},e_m)>(f_{n+1},f_{n+2},\mathrm{},f_m)$$
if for some $`i`$ with $`n+1im,e_j=f_j`$ for $`1<jm`$ and $`e_i>f_i`$.
## 1.3. Definition.
A properly ordered Bratteli diagram is a simple ordered Bratteli diagram $`(V,E)`$ which possesses a unique infinite path $`x_{\mathrm{max}}=(e_1,e_2,\mathrm{})`$ such that each $`e_i`$ is a maximal edge and a unique infinite path $`x_{\mathrm{min}}=(f_1,f_2,\mathrm{})`$ such that each $`f_i`$ is a minimal edge.
Given a properly ordered Bratteli diagram $`B=(V,E,)`$ we denote by $`X_B`$ its infinite path space. So
$$X_B=\{(e_1,e_2,\mathrm{})e_iE_i,r(e_i)=s(e_{i+1}),i=1,2,\mathrm{}\}$$
For an initial segment $`(e_1,e_2,\mathrm{},e_n)`$ we define the cylinder sets
$$U(e_1,e_2,\mathrm{}e_n)=\{(f_1,f_2,\mathrm{})X_Bf_i=e_i,1in\}.$$
By taking cylinder sets to be a basis for open sets $`X_B`$ becomes a topological space. We exclude trivial cases (where $`X_B`$ is finite, or has isolated points). Thus, $`X_B`$ is a Cantor set. $`X_B`$ is a metric space, where for two paths $`x,y`$ whose initial segments to level $`m`$ agree but not to level $`m+1,d(x,y)=1/m+1`$.
## 1.4. Definition.
(Vershik map for a properly ordered Bratteli diagram). If $`x=(e_1,e_2,\mathrm{}e_n,\mathrm{})X_B`$ and if at least one $`e_i`$ is not maximal define
$$V_B(x)=y=(f_1,f_2,\mathrm{},f_j,e_{j+1},e_{j+2},\mathrm{})X_B$$
where $`e_1,e_2,\mathrm{},e_{j1}`$ are maximal, $`e_j`$ is not maximal and has $`f_j`$ as successor in the linearly ordered set $`r^1(r(e_j))`$ and $`(f_1,f_2,\mathrm{},f_{j1})`$ is the minimal path from $`v_0`$ to $`s(f_j)`$. Extend the above $`V_B`$ to all of $`X_B`$ by setting $`V_B(x_{\mathrm{max}})=x_{\mathrm{min}}`$. Then $`(X_B,V_B)`$ is a Cantor minimal dynamical system.
Next, we describe the construction of a dynamical system associated to a non-properly ordered Bratteli diagram. The Bratteli diagram need not be simple. To motivate this construction, it is perhaps worthwhile to begin by indicating how it works in the case of an ordered Bratteli diagram associated to a nested sequence of Kakutani-Rohlin partitions of a Cantor dynamical system $`(X,T)`$.
## 1.5. Definition.
A Kakutani-Rohlin partition of the Cantor minimal system $`(X,T)`$ is a clopen partition $`𝒫`$ of the kind
$$𝒫=\{T^jZ_kkA\mathrm{and}0j<h_k\}$$
where $`A`$ is a finite set and $`h_k`$ is a positive integer. The $`k^{th}`$ tower $`𝒮_k`$ of $`𝒫`$ is $`\{T^jZ_k0j<h_k\}`$ ; its floors are $`T^jZ_k,(0j<h_k)`$. The base of $`𝒫`$ is the set $`Z=_{kA}Z_k`$.
Let $`\{𝒫_n\},(n)`$ be a sequence of Kakutani-Rohlin partitions
$$𝒫_n=\{T^jZ_{n,k}kA_n,\text{and }0j<h_{n,k}\},$$
with $`𝒫_0=\{X\}`$ and with base $`Z_n=_{kA_n}Z_{n,k}`$. We say that this sequence is nested if, for each $`n`$,
1. $`Z_{n+1}Z_n`$
2. $`𝒫_{n+1}`$ refines the partition $`𝒫_n`$.
For the Bratteli-Vershik system $`(X_B,V_B)`$ of sections 1.3-1.4, one obtains a Kakutani-Rohlin partition $`𝒫_n`$ for each $`n`$ by taking the sets in the partition to be the cylinder sets $`U(e_1,e_2,\mathrm{}e_n)`$ of section 1.3 and taking as the base of the partition the union $`U(e_1,e_2,\mathrm{}e_n)`$ over minimal paths (i.e., each $`e_i`$ is a minimal edge). This is a nested sequence.
1.6. To any nested sequence $`\{𝒫_n\},(n)`$ of Kakutani-Rohlin partitions we associate an ordered Bratteli diagram $`B=(V,E,)`$ as follows (see \[DHS, section 2.3\]): the $`A_n`$ towers in $`𝒫_n`$ are in $`11`$ correspondence with $`V_n`$, the set of vertices at level $`n`$. Let $`v_{n,k}V_n`$ correspond to the tower $`𝒮_{n,k}=\{T^jZ_{n,k}0j<h_{n,k}\}`$ in $`𝒫_n`$. We refer to $`T^jZ_{n,k},0j<h_{n,k}`$ as floors of the tower $`𝒮_{n,k}`$ and to $`h_{n,k}`$ as the height of the tower. We will exclude nested sequences of K-R partitions where the infimum (over $`k`$ for fixed $`n`$) of the height $`h_{n,k}`$ does not go to infinity with $`n`$. Let us view the tower $`𝒮_{n,k}`$ against the partition $`𝒫_{n1}=\{T^jZ_{n1,k}kA_{n1}`$, and $`0j<h_{n1,k}\}`$. As the floors of $`𝒮_{n,k}`$ rise from level $`j=0`$ to level $`j=h_{n,k}1`$, $`𝒮_{n,k}`$ will start traversing a tower $`𝒮_{n1,i_1}`$ from the bottom to the top floor, then another tower $`𝒮_{n1,i_2}`$ from the bottom to the top floor, then another tower $`𝒮_{n1,i_3}`$ likewise and so on till a final segment of $`𝒮_{n,k}`$ traverses a tower $`𝒮_{n1,i_m}`$ from the bottom to the top. Note that in this final step the top floor $`T^jZ_{n,k}`$ for $`j=h_{n,k}1`$ of $`𝒮_{n,k}`$ reaches the top floor $`T^qZ_{n1,i_m}`$ for $`q=h_{n1,i_m}1`$ of $`𝒮_{n1,i_m}`$ as a consequence of the assumption $`Z_nZ_{n1}`$ and the fact that $`T^1`$ (union of bottom floors) = union of top floors. Bearing in mind this order in which $`𝒮_{n,k}`$ traverses $`𝒮_{n1,i_1},𝒮_{n1,i_2},\mathrm{},𝒮_{n1,i_m}`$ we associate $`m`$ edges, ordered as $`e_{1,k}<e_{2,k}<\mathrm{}<e_{m,k}`$ and we set the range and source maps for edges by $`r(e_{j,k})=v_{n,k}`$ and $`s(e_{j,k})=v_{n1,i_j}`$. Note that $`m`$ depends on the index $`kA_n`$ (and that by convention the indexing sets $`A_n`$ are disjoint). $`E_n`$ is the disjoint union over $`kA_n`$ of the edges having range in $`V_n`$.
1.7. For $`xX`$, we define $`x_n𝒫_n^{},n`$ as follows: $`x_n=(x_{n,i})_i`$, where $`x_{n,i}𝒫_n`$ is the unique floor in $`𝒫_n`$ to which $`T^i(x)`$ belongs. If $`m>n`$, let $`j_{m,n}:𝒫_m𝒫_n`$ be the unique map defined by $`j_{m,n}(F)=F^{}`$ if $`FF^{}`$. (By abuse of notation, we use the same symbol $`F`$ to denote a point of the finite set $`𝒫_m`$ and also to denote the subset of $`X`$, in the partition $`𝒫_n`$, which $`F`$ represents). An important property of the map
$$X\underset{n}{}(𝒫_n^{}),x(x_1,x_2,\mathrm{}),x_n=(x_{n,i})_i,$$
defined above is the following:
1.8. If $`F`$ and $`TF`$ are two successive floors of a $`𝒫_n`$-tower and if $`x_{n,i}=F`$ then $`x_{n,i+1}=TF`$. If $`x_{n,i}`$ is the top floor of a $`𝒫_n`$-tower, then $`x_{n,i+1}`$ is the bottom floor of a $`𝒫_n`$-tower. More importantly, given integers $`K`$ and $`n`$, there exist $`m>n`$ and a single tower $`𝒮_{m,k}`$ of level $`m`$ such that the finite sequence $`(x_{n,i})_{KiK}`$ is an interval segment contained in
$$\{j_{m,n}(T^{\mathrm{}}(Z_{m,k}))0\mathrm{}<h_{m,k}\}.$$
This is a consequence of the assumption that the infimum of the heights of level-$`n`$ towers goes to infinity. It is true that $`x_{n,i}=j_{m,n}(x_{m,i})`$, but the sequence $`(x_{m,i})_{KiK}`$ need not be an interval segment of $`\{T^{\mathrm{}}(Z_{m,k})0\mathrm{}<h_{m,k}\}`$.
The foregoing observations in the case of an ordered Bratteli diagram associated to a nested sequence of Kakutani-Rohlin partitions gives us the hint to define a dynamical system $`(X_B,T_B)`$ of a non properly ordered Bratteli diagram $`B=(V,E,)`$ as follows:
## 1.9. Definition.
For each $`n`$ define $`\varpi _n=`$ the set of paths from $`V_0`$ to $`V_n`$. There is an obvious truncation map $`j_{m,n}:\varpi _m\varpi _n`$ which truncates paths from $`V_0`$ to $`V_m`$ to the initial segment ending in $`V_n`$. For each $`vV_n`$, the set $`\varpi (v)`$ of paths from $`\{\}V_0`$ ending at $`v`$ will be called a ‘$`\varpi _n`$-tower parametrised by $`v`$’; the cardinality $`|\varpi (v)|`$ will be referred to as the height of this tower. Each tower is a linearly ordered set (whose elements may be referred to as floors of the tower) since paths from $`v_0`$ to $`v`$ acquire a linear order (cf. 1.2). We will exclude unusual examples of ordered Bratteli diagram where the infimum of the height of level-$`n`$ towers does not go to infinity with $`n`$, (for example like \[HPS, Example 3.2\]). Now, we define
1.10. Definition. $`X_B=\{x=(x_1,x_2,\mathrm{},x_n,\mathrm{})\}`$ where
1. $`x_n=(x_{n,i})_i\varpi _n^{}`$,
2. $`j_{m,n}(x_{m,i})=x_{n,i}`$ for $`m>n`$ and $`i`$ and
3. given $`n`$ and $`K`$ there exists $`m`$ such that $`m>n`$ and a vertex $`vV_m`$, such that the interval segment $`x_n[K,K]:=(x_{n,K},x_{n,K+1},\mathrm{},x_{n,K})`$ is obtained by applying $`j_{m,n}`$ to an interval segment of the linearly ordered set of paths from $`v_0`$ to $`v`$.
The condition (iii) is the crucial part of the definition. Without it what one gets is an inverse system.
The condition (iii) implies that a property similar to (1.8) holds. Since each $`\varpi _n`$ is a finite set $`\varpi _n^{}`$ has a product topology which makes it a compact set - in fact a Cantor set. Likewise, $`_n(\varpi _n^{})`$ is again a Cantor set. Thus, $`X_B_n(\varpi _n^{})`$ has an induced topology. The lemma below and the following proposition are analogous to corresponding facts for the Vershik model associated to properly ordered Bratteli diagrams.
## 1.11. Lemma.
The topological space $`X_B`$ is compact.
## Proof.
We show that $`X_B`$ is closed in $`_n\varpi _n^{}`$. Let $`z=(z_1,z_2,\mathrm{},z_n,\mathrm{})`$, where $`z_n\varpi _n^{}`$. Assume that $`z=lim_m\mathrm{}w_m`$, where $`w_m_n\varpi _n^{}`$ and moreover $`w_mX_B`$. Let $`K_1,K_2,\mathrm{},`$ be a strictly increasing sequence of positive integers. Define neighbourhoods $`U_1,U_2,\mathrm{},U_m,\mathrm{}`$ of $`z`$ in $`_n\varpi _n^{}`$ shrinking to $`z`$ by $`U_m=\left\{z^{}=(z_1^{},z_2^{},\mathrm{},z_m^{},\mathrm{})_n\varpi _n^{}\right\}`$ where $`z_k^{}`$ has the same coordinates as $`z_k`$ in the range $`[K_m,K_m]`$, i.e. $`z_k^{}[K_m,K_m]=z_k[K_m,K_m]`$ for $`1km`$. Since $`z=limw_j`$, for any given $`m`$, $`J(m)`$, such that $`\mathrm{`}\mathrm{`}j>J(m)^{\prime \prime }w_jU_m`$. But, since $`w_jX_B`$, we conclude that for $`1km`$, the interval segments $`z_k[K_m,K_m]`$ are obtained by applying $`j_{M,k}`$ to an interval segment of the sequence of floors of a single $`\varpi _M`$tower (for some $`M)`$. This shows that $`zX_B`$. $`\mathrm{}`$
Denote by $`T_B`$ the restriction of the shift operator to $`X_B`$. So, if $`x=(x_1,x_2,\mathrm{},x_n,\mathrm{}),`$ where $`x_n=(x_{n,i})_i\varpi _n^{}`$, then $`T_B(x)=(x_1^{},x_2^{},\mathrm{},x_n^{},\mathrm{})`$, where $`x_n^{}=(x_{n,i}^{})_i\varpi _n^{}`$ and $`x_{n,i}^{}=x_{n,i+1}`$.
$`(X_B,T_B)`$ will be called the dynamical system associated to $`B=(V,E,)`$.
## 1.12. Proposition.
If $`B=(V,E,)`$ is a simple ordered Bratteli diagram, then $`(X_B,T_B)`$ is a Cantor minimal dynamical system.
## Proof.
Let $`x,yX_B`$, where $`x=(x_1,\mathrm{},x_m,\mathrm{})`$ and $`y=(y_1,\mathrm{},y_m,\mathrm{})`$ satisfy the conditions i), ii) and iii) of definition (1.10). We will show that $`x`$ belongs to the closure of the orbit $`\{T_B^i(y)i\}`$. For integers $`n`$ and $`K`$ define the neighbourhood $`U(n,K,x)`$ of $`x`$ to be the set
$`\{z=(z_1,z_2,\mathrm{},z_m,\mathrm{})X_Bz_m[K,K]=x_m[K,K]`$ for $`1mn\}`$.
Given any neighbourhood $`U`$ of $`x`$, choose $`n`$ and $`K`$ such that $`UU(n,K,x)`$. Since $`xX_B`$, we can choose $`J`$ and a $`𝒫_J`$-tower $`𝒮_v`$ of paths from $`v_0`$ ending at a fixed $`vV_J`$, such that $`x_n[K,K]`$ is obtained by applying $`j_{J,n}`$ to an interval segment of $`𝒮_v`$. Since the Bratteli diagram $`(V,E,)`$ is simple, we can choose $`L>J`$ such that every point of $`V_J`$ is connected to every point of $`V_L`$ by a path. This implies that $`x_n[K,K]`$ occurs as an interval segment of the sequence obtained by applying $`j_{L,n}`$ to the sequence of floors of any $`𝒫_L`$-tower. Now, choose a $`<b`$ such that $`y_L[a,b]`$ is the sequence of all floors of a $`𝒫_L`$-tower. From the preceding observation, we can choose $`d`$ such that $`ad<d+2Kb`$ and $`x_n[K,K]=j_{L,n}(y_L[d,d+2K])`$. But since $`yX_B`$, $`y`$ satisfies condition (ii) in definition 1.10. So $`x_n[K,K]=y_n[d,d+2K]`$ and we conclude $`T_B^{d+K}(y)U(n,K,x)`$. $`\mathrm{}`$
In (1.7), given a nested sequence of Kakutani-Rohlin partitions of $`(X,T)`$, we defined a map from $`(X,T)`$ to the dynamical system $`(X_B,T_B)`$ of the associated ordered Bratteli diagram. It follows that if $`(X,T)`$ is minimal, and if the Bratteli diagram of the nested sequence of K-R partitions is a simple Bratteli diagram, then $`(X,T)(X_B,T_B)`$ is onto. If the topology of $`(X,T)`$ is spanned by the collection of the clopen sets belonging to the K-R partitions then clearly the map $`(X,T)(X_B,T_B)`$ is injective. In particular, if the Bratteli diagram is properly ordered then the Bratteli-Vershik system is naturally isomorphic to the system given by our construction in 1.10.
Note that the same term ‘towers’ has been used to denote two separate but related objects \[in (1.5) and (1.9)\]. For $`vV_n`$, let $`y`$ be a path from $`\{\}`$ to $`v`$ in $`(V,E,)`$. So, $`y`$ is a ‘floor’ (consisting of the single element $`y`$) belonging to the $`\varpi _n`$\- tower $`\varpi (v)`$ (a finite set) parametrized by $`vV_n`$ \- all in the sense of $`(1.9)`$. Here, $`\varpi (v)`$= all paths from $`\{\}`$ to $`v`$. Put $`_y=\{x=(x_1,x_2,\mathrm{},x_n,\mathrm{})X_Bx_{n,0}=y\}`$. $`_y`$ is a clopen set of the Cantor set $`X_B`$. Put $`𝒫_n=\{_yy\varpi (v),vV_n\}.`$ Then, in the sense of $`(1.5)`$ $`𝒫_n`$ is a K-R partition of $`X_B`$ whose base is the union of $`_y,(y\text{ minimal }\varpi (v),vV_n)`$. Its towers $`S_v`$ are parametrized by $`vV_n`$: $`S_v=\{_yy\varpi (v)\}`$. $`_y,(y\varpi (v))`$ are the floors of the tower $`S_v`$. (We encountered this K-R partition earlier in the case of the Bratteli-Vershik system at the end of 1.5).It is easy to see that the ordered Bratteli diagram obtained from $`\{_yy\varpi (v),vV_n\}`$ is $`(V,E,)`$. (compare the last two lines in the proof of \[HPS, Theorem 4.5\] and also the opening observation in the proof of \[DHS, Proposition 16\]).
2. Skew-product dynamical systems
2.1. We recall the notion of ‘skew-product dynamical systems’ which was defined by Matui in \[M\]. Let $`(X,T)`$ be a Cantor dynamical system. Let $`G`$ be a finite group and $`c:XG`$ a continuous function. We set $`Y=X\times G`$ and define a continuous map $`S:YY`$ by $`S(x,g)=(T(x),gc(T(x)))`$. The dynamical system $`(Y,S)`$ is called the skew-product extension associated to the $`G`$-valued cocycle $`c`$. For $`gG`$, define $`\gamma _g:YY`$ by $`\gamma _g(x,h)=(x,gh)`$. Clearly, $`\gamma _gS=S\gamma _g,g`$. Furthermore, the map $`X\times GX`$ given by projection to the first factor displays $`(Y,S)`$ as an extension of $`(X,T)`$, (or, $`(X,T)`$ as a factor of $`(Y,S))`$. In a different context, Kumjian and Pask defined in \[KP\] the $`C^{}`$-algebra associated to a graph and further when the edges of a graph are labelled with values in a finite group $`G`$, they considered a ‘skew-product graph’ and ‘skew-product $`C^{}`$-algebra’ associated to the labelling.
Motivated by these two works, we employ our construction of subsection $`1.10`$ to associate a ‘skew-product’ dynamical system when the edges of an ordered Bratteli diagram are labelled with values in a finite group $`G`$.
2.2. Let $`(V,E,)`$ be an ordered Bratteli diagram. Let $`G`$ be a finite group and $`\lambda :EG`$ a map defined on the edges with values in $`G`$ ($`\lambda \stackrel{defn.}{=}`$ the ‘labels’ on edges). We define a new Bratteli diagram $`B(\lambda )\stackrel{defn.}{=}(V_\lambda ,E_\lambda )`$ as follows: $`V_{0,\lambda }=\{\},V_{n,\lambda }=V_n\times G,(n1)`$ and $`E_{n,\lambda }=E_n\times G,(n1)`$. The source and range maps on $`E_\lambda `$ are defined by $`r(e,g)=(r(e),g)`$, and $`s(e,g)=(s(e),g\lambda (e))`$. If $`eE_1`$, we define $`s(e,g)=\{\}`$. We define a map $`\pi :(V_\lambda ,E_\lambda )(V,E)`$ by $`\pi (v_{\lambda ,0})=v_0,\pi (v,g)=v`$ if $`vV_n,n1`$ and $`\pi (e,g)=e`$. One sees that $`\pi `$ commutes with the range and source map. It is easy to see that $`\pi |_{r^1(v,g)}`$ is a bijection onto $`r^1(v),vV\{v_0\}`$ and $`gG`$. Thus if $`(V,E,)`$ is an ordered Bratteli diagram there is a unique order in $`E_\lambda `$ such that $`\pi :(V_\lambda ,E_\lambda ,)(V,E,)`$ is order preserving. In the sequel we assume that $`(V_\lambda ,E_\lambda )`$ is equipped with this order. The dynamical system $`(X_\lambda ,T_\lambda )\stackrel{defn.}{=}`$ the system constructed in $`1.10`$ for this ordered Bratteli diagram $`B(\lambda )=(V_\lambda ,E_\lambda ,)`$ will be called the skew-product system for the edge labelling $`\lambda `$. We also remark that
$$\pi :(V_\lambda ,E_\lambda )(V,E)\text{ has the ‘}\text{unique path lifting}\text{’ property }$$
in the following sense. If $`m>n1`$, and $`(e_n,e_{n+1},\mathrm{},e_m)`$ is a path in $`(V,E)`$ from $`V_{n1}`$ to $`V_m`$ with $`r(e_m)=v`$ then for any $`gG`$, there is a unique path $`(\stackrel{~}{e}_n,\stackrel{~}{e}_{n+1},\mathrm{},\stackrel{~}{e}_m)`$ in $`(V_\lambda ,E_\lambda )`$ which maps onto $`(e_n,e_{n+1},\mathrm{},e_m)`$ under $`\pi `$ and such that $`r(\stackrel{~}{e}_m)=(v,g)`$.
\[CAUTION :- Even in the presence of unique path lifting property two different edges on the left with the same source may map into the same edge on the right. See example (2.8) below. What the property asserts is that two different edges on the left with the same range cannot map to the same edge on the right.\]
The group $`G`$ acts on $`(V_\lambda ,E_\lambda ,)`$ by $`\gamma _g(v,h)=(v,gh)`$ and $`\gamma _g(e,h)=(e,gh)`$ for $`vV\{v_0\}`$ and of course $`\gamma _g(v_{0,\lambda })=v_{0,\lambda }`$.
Given the labelling $`\lambda `$ we can extend the labelling to paths from $`V_{n1}`$ to $`V_m`$, for $`m>n`$. With notation as above, we define $`\lambda (e_n,e_{n+1},\mathrm{},e_m)=\lambda (e_m)\lambda (e_{m1})\mathrm{}\lambda (e_n)`$. Let $`\{n_k\}_{k=0}^{\mathrm{}}`$ be a subsequence of $`\{0,1,2,\mathrm{}\}`$ where we assume $`n_0=0`$. A Bratteli diagram $`(V^{},E^{})`$ is called a ‘telescoping’ of $`(V,E)`$ if $`V_k^{}=V_{n_k}`$ and $`E_k^{}`$ consists of paths $`(e_{n_{k1}+1},\mathrm{},e_{n_k})`$ from $`V_{n_{k1}}`$ to $`V_{n_k}`$ in $`(V,E)`$, the range and source maps being the obvious ones. Thus, for every telescoping the labelling $`\lambda `$ on the edges of $`(V,E)`$ gives rise to a labelling on the edges of $`(V^{},E^{})`$.
2.3. Remark. Suppose there is a telescoping $`(V^{},E^{})`$ of $`(V,E)`$ such that the induced labelling $`\lambda `$ on $`(V^{},E^{})`$ has the following property: given $`k`$, there exists $`v^{}V_k^{}`$ and $`w^{}V_{k+1}^{}`$, such that $`\lambda (s^1(v^{})r^1(w^{}))=G`$. If in addition $`(V,E)`$ is simple, then $`(V_\lambda ,E_\lambda )`$ is simple.
2.4. Stationary Bratteli diagrams. A Bratteli diagram is stationary if the diagram repeats itself after level $`1`$. (One may relax by allowing a period from some level onwards; but, a telescoping will be stationary in the above restricted sense.) If $`(V,E,)`$ is an ordered Bratteli diagram and the diagram together with the order repeats itself after level $`1`$, then $`(V,E,)`$ will be called a stationary ordered Bratteli diagram. We refer the reader to \[DHS, section (3.3)\] for the usual definition of a substitutional system and how they give rise to stationary Bratteli diagrams. Some details are recalled below. Let $`(V,E,)`$ be as above and suppose moreover that it is a simple Bratteli diagram and that $`\lambda `$ is a labelling of the edges with values in a finite group $`G`$. Assume that the labelling is stationary: so, we have
1. an enumeration $`\{v_{n,1},v_{n,2},\mathrm{},v_{n,L}\}`$ of $`V_n,n1`$,
2. for $`n>1`$ and $`1jL`$ an enumeration $`\{e_{n,j,1},e_{n,j,2},\mathrm{},e_{n,j,a_j}\}`$ of $`r^1(e_{n,j})`$ which is assumed to be listed in the linear order in $`r^1(v_{n,j})`$,
3. in the enumerations above, $`L`$ does not depend on $`n`$ and $`a_j`$ depends only on $`j`$ and not on $`n`$. Moreover, if $`n,m>1`$, if $`1jL,1kL,1ia_j`$, then $`\mathrm{`}\mathrm{`}s(e_{n,j,i})=v_{n1,k}`$$`\mathrm{`}\mathrm{`}s(e_{m,j,i})=v_{m1,k}`$”,
4. with notation as in 3) above, $`\lambda (e_{n,j,i})=\lambda (e_{m,j,i})`$.
If $`S`$ is a set of generators for $`G`$ and if in addition to the above, we also assume that $`\lambda (\{r^1(v_{n,j})\}\{s^1(v_{n1,k})\})S\{e\},j,k`$ between 1 and $`L`$ then $`(V_\lambda ,E_\lambda )`$ is simple.
2.5. Substitutional systems. Let $`A`$ be an alphabet set. Write $`A^+`$ for the set of words of finite length in the alphabets of $`A`$. Let $`\sigma :AA^+`$ be a substitution, written, $`\sigma (a)=\alpha \beta \gamma \mathrm{}`$ The stationary ordered Bratteli diagram $`B=(V,E,)`$ associated to $`(A,\sigma )`$ (cf. \[DHS, section 3.3\] can be described as
$$V_n=A,\text{ }n1,V_0=\{\}$$
$$E_n=\{(a,k,b)a,bA,k,\text{ }a\text{ is the }k^{th}\text{ alphabet in the word }\sigma (b)\}.$$
(The reader who prefers a more carefully evolved notation can consider intoducing an extra factor ‘$`\times \{n\}`$’ so that vertices and edges at different levels are seen to be disjoint). The source and range maps $`s`$ and $`r`$ are defined by $`s(a,k,b)=a,r(a,k,b)=b`$. In the linear order in $`r^1(b)`$, $`(a,k,b)`$ is the $`k^{th}`$ edge. If $`\lambda `$ is a stationary labelling on edges with values in a finite group $`G`$, then define a new ‘skew-product’ substitutional system $`(A_\lambda ,\sigma _\lambda )`$ as follows:=
1. $`A_\lambda =A\times G`$
2. $`E_{n,\lambda }=\{[(a,g),k,(b,h)]\text{ (i) }(a,k,b)E_n,\text{ (ii) }g=h\lambda (a,k,b)\}`$.
Define $`\sigma _\lambda :A_\lambda A_\lambda ^+`$ by the rule that
$$(a,g)\text{ is the }k^{\mathrm{𝑡ℎ}}\text{ alphabet in the word }\sigma _\lambda (b,h)$$
if
$$[(a,g),k,(b,h)]E_{n,\lambda }.$$
Set $`V_\lambda =\{\}\{_nV_{n,\lambda }\}`$ where each $`V_{n,\lambda }`$ is a (disjoint) copy of $`A_\lambda `$ and likewise, $`E_\lambda =_nE_{n,\lambda }`$. The source and range maps are defined by $`s([(a,g),k,(b,h)])=(a,g)`$ and $`r([(a,g),k,(b,h)])=(b,h)`$. Then $`(V_\lambda ,E_\lambda ,)`$ is the stationary ordered Bratteli diagram arising from the substitutional system $`(A_\lambda ,\sigma _\lambda )`$. The group $`G`$ acts on $`(V_\lambda ,E_\lambda ,)`$ by
$$g(a,g_1)=(a,gg_1)$$
$$g[(a,g_1),k,(b,g_2)]=[(a,gg_1),k,(b,gg_2)].$$
The ordered Bratteli diagram $`(V_\lambda ,E_\lambda ,)`$ thus obtained is the same as the skew-product of (2.2) if one starts with the $`(V,E,)`$ in the beginning of (2.5).
To the stationary ordered Bratteli diagram $`B`$ of $`(A,\sigma )`$ (which may not be properly ordered unless $`\sigma `$ is a primitive, aperiodic, proper substitution, – see \[DHS, section 3\]) we can associate a dynamical system $`X_B`$ following the construction of 1.10; we will see that this is naturally isomorphic to the substitutional dynamical system $`X_\sigma `$ associated to $`(A,\sigma )`$ defined for example in \[DHS, section 3.3.1\]. Let us use the nested sequence $`(𝒫_n)`$ of Kakutani-Rohlin partitions in $`X_\sigma `$ defined in \[DHS, corollary 13\]; the ordered Bratteli diagram associated by section 1.6 to this nested sequence $`(𝒫_n)`$ is the same as the stationary ordered Bratteli diagram $`B`$ in the beginning of this section. (this is the opening observation in the proof of \[DHS, Proposition 16\], but, this observation also remains valid for improper substitutions). Next apply the map $`X_\sigma X_B`$ that was defined in 1.7. As remarked after the proof of Proposition 1.12, this map is onto. That this map is also injective can be seen as follows: let $`xX_\sigma `$. So, $`x=(a_i)_iA^{}`$. Let the image of $`x`$ in $`X_B`$ by the map of 1.7 be denoted by $`(x_1,x_2,\mathrm{})`$ where $`x_n=(x_{n,i})_i𝒫_n^{}`$. The sequence $`x=(a_i)_iA^{}`$ can be read off from the first co-ordinate $`x_1=(x_{1,i})_i𝒫_1^{}`$. Choose any element $`zx_{1,i}`$. Writing, uniquely, as in \[DHS, Corollary 12(ii)\] $`z=T_\sigma ^k(\sigma (y))`$ where $`yX_\sigma `$ and $`0k<|\sigma (y_0)|,`$ it is at once seen that $`a_i\text{ is the }k^{th}\text{ alphabet in }\sigma (y_0)`$ and is independent of $`zx_{1,i}`$. In fact it is a part of the definition \[DHS, Corollary 13\] of the Kakutani-Rohlin partition $`𝒫_n`$ that $`y_0\text{ and }k`$ are independent of $`zx_{1,i}`$
2.6. Example. Consider the (proper) substitutional (dynamical) system in the alphabet set $`A=\{X,Y\}`$ given by $`\sigma (X)=XXY`$, $`\sigma (Y)=XYY`$. The corresponding (stationary) Bratteli diagram $`(V,E,)`$ can be represented by
Diagram 1
2.7. Example $`(2.6)`$ continued: For our substitutional system $`A=\{X,Y\}`$, $`\sigma :AA^+`$ given by $`\sigma (X)=XXY,\sigma (Y)=XYY`$, the associated Bratteli diagram is stationary. Each $`V_n`$ can be identified with $`\{X,Y\}`$. There are six edges connecting $`V_n`$ and $`V_{n+1}`$ noted $`e_{XX}^1,e_{YX}^2,e_{YX}^3,e_{XY}^1,e_{XY}^1,e_{YY}^2,e_{YY}^3`$ with source and range maps defined by
$$\begin{array}{ccc}s(e_{XX}^1)=X\hfill & ,\hfill & r(e_{XX}^1)=X\hfill \\ s(e_{XX}^2)=X\hfill & ,\hfill & r(e_{XX}^2)=X\hfill \\ s(e_{YX}^3)=Y\hfill & ,\hfill & r(e_{YX}^3)=X\hfill \\ s(e_{XY}^1)=X\hfill & ,\hfill & r(e_{XY}^1)=Y\hfill \\ s(e_{YY}^2)=Y\hfill & ,\hfill & r(e_{YY}^2)=Y\hfill \\ s(e_{YY}^3)=Y\hfill & ,\hfill & r(e_{YY}^3)=Y\hfill \end{array}$$
The edges with range $`X`$ are linearly ordered as $`\{e_{XX}^1,e_{XX}^2,e_{YX}^3\}`$. The edges with range $`Y`$ are linearly ordered as $`\{e_{XY}^1,e_{YY}^2,e_{YY}^3\}`$. Consider the stationary labelling with values in the group $`G=_2=\{\overline{0},\overline{1}\}`$ given by $`\lambda (e_{YY}^2)=\overline{1}`$ and $`\lambda (f)=\overline{0}`$ for the remaining edges.
We now describe $`(V_\lambda ,E_\lambda )`$. Recall $`V_{n,\lambda }=V_n\times G`$. Thus each $`V_{n,\lambda }`$ can be identified with $`\{X,Y\}\times _2`$. We identify $`\{X,Y\}\times _2`$ with the set of alphabets $`\{x,x^{},y,y^{}\}`$ by $`x=(X,\overline{0}),x^{}=(X,\overline{1}),y=(Y,\overline{0}),y^{}=(Y,\overline{1})`$. Then the edges $`V_{n,\lambda }V_{n+1,\lambda }`$ can be diagrammatically represented by
Diagram 2
One concludes that $`(V_\lambda ,E_\lambda ,)`$ is the stationary ordered Bratteli diagram associated to the substitutional system on the alphabet set $`A_\lambda =\{x,x^{},y,y^{}\}`$ where the substitution $`\sigma _\lambda :A_\lambda A_\lambda ^+`$ is given by
$$\begin{array}{ccc}x\hfill & \hfill & xxy\hfill \\ x^{}\hfill & \hfill & x^{}x^{}y^{}\hfill \\ y\hfill & \hfill & xy^{}y\hfill \\ y^{}\hfill & \hfill & x^{}yy^{}\hfill \end{array}$$
The Bratteli diagram $`(V_\lambda ,E_\lambda ,)`$ is simple:- any vertex at level $`n`$ can be connected to any vertex at level $`n+3`$ by a path.
2.8. Example where $`G`$ is infinite. Let $`G`$ be a residually finite group. Let $`GN_0N_1\mathrm{}N_i\mathrm{}`$ be a decreasing sequence of cofinite normal subgroups of $`G`$ such that $`_i(N_i)=\{e\}`$. Let $`\phi _i:GG/N_i`$ be the canonical projection. We use freely the notation for a substitutional system $`(A,\sigma )`$ set up in $`(2.5)`$. Let $`\lambda `$ be a stationary labelling on $`E=\{(a,k,b)a\text{ is the }k^{th}\text{ alphabet in }\sigma (b)\}`$ with values in $`G`$. The skew-product system in this example is defined as follows:-
$$V_{n,\lambda }=A\times (G/N_n),n$$
$$E_{n,\lambda }=\{[(a,\phi _{n1}(g)),k,(b,\phi _n(h))]\}$$
where $`a,bA,(a,k,b)E`$ and $`\phi _{n1}(g)=\phi _{n1}(h\lambda (a,k,b))`$.
In the statements below, we do not assume that $`(V_\lambda ,E_\lambda ,)`$ arises from a stationary labelling of a stationary ordered Bratteli diagram. Nor do we assume that $`(V_\lambda ,E_\lambda ,)`$ is simple. We only assume that
1. $`(V,E,)`$ is a simple ordered Bratteli diagram and
2. $`\lambda `$ is a labelling with values in $`G`$.
2.9. Let $`(X,T)`$ and $`(X_\lambda ,T_\lambda )`$ be the Cantor dynamical systems associated respectively to $`(V,E,)`$ and $`(V_\lambda ,E_\lambda ,)`$ as in subsection (1.10). The map $`\pi :(V_\lambda ,E_\lambda ,)(V,E,)`$ sends paths from $`\{\}`$ to $`(v,g)`$ in $`(V_\lambda ,E_\lambda )`$ to paths from $`\{\}`$ to $`v`$ in $`(V,E)`$ and respects truncation (cf. 1.9). The unique path lifting property implies that $`\pi `$ maps a $`\varpi _{n,\lambda }`$\- tower in $`(V_\lambda ,E_\lambda ,)`$ parametrised by $`(v,g)V_{n,\lambda }`$ bijectively onto the $`\varpi _n`$-tower in $`(V,E,)`$ parametrised by $`vV_n`$ (and of course, respects the linear order of the floors). Thus $`\pi `$ induces a map $`\pi :(X_\lambda ,T_\lambda )(X,T)`$ between the two dynamical systems. The action of $`G`$ on $`(V_\lambda ,E_\lambda ,)`$ given by $`\{\gamma _g\}_{gG}`$ gives rise to a free action $`\{\gamma _g\}_{gG}`$ of $`G`$ on $`(X_\lambda ,T_\lambda )`$ such that $`\pi \gamma _g=\pi `$.
2.10. Remark. As observed by Matui \[M\], when $`G`$ acts freely on a Cantor dynamical system $`(Y,S)`$, the latter is isomorphic to the skew-product dynamical system arising from a $`G`$-valued cocycle on the quotient of the action. If one starts with a cocycle on $`(X,T)`$ then for the associated skew-product system $`(Y,S)`$ where $`Y=X\times G`$, the quotient is naturally isomorphic to $`(X,T)`$. But, in our situation $`(X_\lambda ,T_\lambda )`$ arises from a $`G`$-valued labelling of the edges of $`(V,E,)`$, the dynamical system associated to $`(V,E,)`$ is $`(X,T)`$, but the quotient of $`(X_\lambda ,T_\lambda )`$ by the $`G`$-action $`\{\gamma _g\}_{gG}`$ may still be a nontrivial extension of $`(X,T)`$. We will make explicit the direct link between the labelling and the cocycle on the quotient. In the next few sections, we describe an ordered Bratteli diagram for the quotient of $`(X_\lambda ,T_\lambda )`$ by the $`G`$-action $`\{\gamma _g\}_{gG}`$. The Bratteli diagram thus constructed naturally maps to $`(V,E,)`$ with unique path lifting property.
Recall that by definition $`X_\lambda =\{x=(x_1,x_2,\mathrm{},x_n,\mathrm{})\}`$ where
1. $`x_n=(x_{n,i})_i(\varpi _{n,\lambda })^{}`$
2. $`j_{m,n}(x_{m,i})=x_{n,i}`$ for $`m>n`$ and $`i`$
3. Given $`n`$ and $`K`$, $`\text{ }m>n`$ and a vertex $`vV_{m,\lambda }`$ such that the interval segment $`x_n[K,K]:=(x_{n,K},x_{n,K+1},\mathrm{},x_{n,K})`$ is obtained by applying $`j_{m,n}`$ to an interval segment of the linearly ordered set of paths from $`v_0`$ to $`v`$.
2.11. Definition: “Loops lift to loops”.
If $`(V,E,)`$ and $`(V^{},E^{},)`$ are two ordered Bratteli diagrams and $`\pi :(V,E,)(V^{},E^{},)`$ is a morphism with unique path-lifting property we say that $`\pi `$ “lifts loops to loops” (or, that $`\pi `$ has the “loops lifting to loops” property) if the following condition is satisfied:-
(“loops on the right”) Let $`u^{}V_m^{},v^{}V_n^{}`$. Let $`km`$ and $`kn`$. Suppose $`\alpha ^{}`$ and $`\beta ^{}`$ are paths from $`V_k^{}`$ ranging at $`u^{}`$. Suppose $`\gamma ^{}`$ and $`\delta ^{}`$ are paths from $`V_k^{}`$ ranging at $`v^{}`$. In the lexicographic order induced on paths from $`V_k^{}`$ to $`u^{}`$ assume that $`\alpha ^{}`$ is the successor of $`\beta ^{}`$. Similarly for paths from $`V_k^{}`$ to $`v^{}`$ assume that $`\gamma ^{}`$ is the successor of $`\delta ^{}`$. Assume that $`\alpha ^{}`$ and $`\gamma ^{}`$ have the same source in $`V_k^{}`$ and that, likewise, $`\beta ^{}`$ and $`\delta ^{}`$ have the same source in $`V_k^{}`$. (Thus, one has a loop:- going from the range of $`\alpha ^{}`$ to the source of $`\beta ^{}`$ via $`\beta ^{}`$, then to the range of $`\delta ^{}`$ via $`\delta ^{}`$, then to the source of $`\alpha ^{}`$ via $`\gamma ^{}`$ and then to the range of $`\alpha ^{}`$ via $`\alpha ^{}`$.)
(‘pull-back of the above loop to the left’) Now, let $`uV_m,vV_n`$ and assume $`\pi (u)=u^{},\pi (v)=v^{}`$. Let $`\alpha `$ (resp. $`\beta )`$ be the unique path from $`V_k`$ to $`V_m`$ lying above $`\alpha ^{}`$, (resp. $`\beta ^{})`$ and ranging at $`u`$. Similarly, let $`\gamma `$ (resp. $`\delta )`$ be the unique path from $`V_k`$ to $`V_n`$ lying above $`\gamma ^{}`$ (resp. $`\delta ^{})`$ and ranging at $`v`$. If $`\pi `$ has the property that under the above conditions,
$$\text{“{source of }\beta =\text{ source of }\delta \}\mathrm{"}\mathrm{`}\mathrm{`}\{\text{ source of }\alpha =\text{ source of }\gamma \}\mathrm{"}$$
then we say that $`\pi `$ has the ‘loops lifting to loops’ property.
Diagram 3
Finally, we say that the labelling $`\lambda :EG`$ has the “loops lifting to loops” property if $`(V_\lambda ,E_\lambda ,)(V,E,)`$ has the above property.
2.11.1. Example. Let $`\beta :VG`$ be a function. The labelling $`\lambda `$ defined on $`E`$ by $`\lambda (e)=\beta (r(e))^1\beta (s(e))`$ has the “loops lifting to loops” property.
2.11.2. Example. Let $`s,tG`$. Let $`\lambda `$ be a $`G`$-valued stationary labelling on the $`2adic`$ odometer system(cf. \[GJ,pp.1691,1695\]) whose image consists of $`\{s,t\}`$.
The Bratteli diagram of the odometer has one vertex at each level and two edges, say, $`min`$ and $`max`$ from level $`k`$ to $`k+1`$ for $`k1`$, with $`\lambda (min)=s`$ and $`\lambda (max)=t`$. In 2.11, we take
$`\alpha ^{}\text{ = }max`$ from level 1 to 2,
$`\beta ^{}\text{ = }min`$ from level 1 to 2,
$`\gamma ^{}\text{ = }(min,max)`$ from level 1 to 3,
$`\delta ^{}\text{ = }(max,min)`$ from level 1 to 3.
Denote by $`w^{},u^{},v^{}`$ the vertices at levels 1,2,3. Following the description of vertices and edges for $`(V_\lambda ,E_\lambda )`$ as in 2.5,
first (resp.second) edge ranging at $`(u^{},g)`$ has source $`(w^{},gs)`$ (resp.$`(w^{},gt)`$.
first (resp.second) edge ranging at $`(v^{},h)`$ has source $`(u^{},hs)`$ (resp.$`(u^{},ht)`$.
source $`(w^{},gt)`$ to range $`(u^{},g)`$ is a lift $`\alpha `$ of $`\alpha ^{}`$.
source $`(w^{},gs)`$ to range $`(u^{},g)`$ is a lift $`\beta `$ of $`\beta ^{}`$.
$`(w^{},hts)`$ to $`(u^{},ht)`$ to $`(v^{},h)`$ is a lift $`\gamma `$ of $`\gamma ^{}`$.
$`(w^{},hst)`$ to $`(u^{},hs)`$ to $`(v^{},h)`$ is a lift $`\delta `$ of $`\delta ^{}`$.
If the “loops lifting to loops” property holds then for the above choice of $`\alpha ^{},\beta ^{},\gamma ^{},\delta ^{}`$ the condition that
$`\text{ “{source of }\beta =\text{ source of }\delta \}\mathrm{"}\mathrm{`}\mathrm{`}\{\text{ source of }\alpha =\text{ source of }\gamma \}\mathrm{"}`$
yields the property in the group
$`\text{ “}\{(w^{},gs)=(w^{},hst)\}\mathrm{"}\mathrm{`}\mathrm{`}\{(w^{},gt)=(w^{},hts)\}\mathrm{"}.`$
This reduces to the condition that $`st^1s^1=ts^1t^1`$. One can easily see that analogous conditions with other choices of $`\alpha ^{},\beta ^{},\gamma ^{},\delta ^{}`$ are all consequences of the property $`st^1s^1=ts^1t^1`$.
Assume that $`G`$ is generated by $`s,t`$ and that $`st^1s^1=ts^1t^1`$. Then either $`s=t`$, (so $`G`$ is cyclic), or, $`G`$ has the following simple description:
Put $`x=ts^1`$. Then, $`x^2=ts^1ts^1=ts^1ts^1t^1t=ts^1st^1s^1t=s^1t`$. So, $`t^1xt=s^1t=x^2`$. Let $`H`$ be the subgroup of $`G`$ generated by $`x`$. Let $`J`$ be the subgroup of $`G`$ generated by $`t`$. Note that elements of $`J`$ normalize $`H`$. Also, $`JH=\{e\}`$ the unit element. In fact since $`x`$ and $`x^2`$ are conjugate, the order of $`x`$ has to be $`2k+1`$ an odd number. If $`x^{\mathrm{}}=t^mJH`$, where, $`0\mathrm{}<2k+1`$ then $`x^{\mathrm{}}=t^m=t^1x^{\mathrm{}}t=x^2\mathrm{}`$, which implies that $`x^{\mathrm{}}=\{e\}`$ which is not possible (unless $`\mathrm{}=0`$) since $`\mathrm{}<`$ order of $`x`$.
Hence, $`G`$ is the semi-direct product of $`J`$ and $`H`$ for the action given by $`t^1xt=x^2`$. The skew-product Bratteli diagram $`(V_\lambda ,E_\lambda ,)`$ has the following connectivity structure. Telescoping to levels $`\{0,n_1=1,n_2,n_3,\mathrm{},n_i,\mathrm{}\}`$ where $`n_i=1+(i1)\times \{\text{order of }G\}`$, a vertex $`(v,g)`$ in level $`n_i`$ is connected to a vertex $`(w,g^{})`$ in level $`n_{i+1}`$ by an edge if and only if $`g\text{ and }g^{}`$ belong to the same $`H`$-coset.
Our main result about the skew-product dynamical systems $`(X_\lambda ,T_\lambda )`$ which were constructed in $`2.2`$ are contained in the following.
2.12.Theorem. (I) Suppose that the map $`\pi :(V_\lambda ,E_\lambda ,)(V,E,)`$ has the property ‘loops lift to loops’. Then the quotient of $`(X_\lambda ,T_\lambda )`$ by the action of $`G`$ is canonically isomorphic to $`(X,T)`$.
(II) In general, there is a commutative diagram
$$\begin{array}{ccc}(V_\lambda ,E_\lambda ,)& \stackrel{\pi }{}& (V,E,)\\ \rho ^{}& & \rho \\ (\stackrel{~}{V}_\mu ,\stackrel{~}{E}_\mu ,)& \stackrel{\stackrel{~}{\pi }}{}& (\stackrel{~}{V},\stackrel{~}{E},)\end{array}$$
and the induced diagram
$$\begin{array}{ccc}(X_\lambda ,T_\lambda )& \stackrel{\pi _B}{}& (X,T)\\ \rho _B^{}& & \rho _B\\ (\stackrel{~}{X}_\mu ,\stackrel{~}{T}_\mu )& \stackrel{\stackrel{~}{\pi }_B}{}& (\stackrel{~}{X},\stackrel{~}{T})\end{array}$$
between the corresponding dynamical systems, where
(i) All the horizontal and vertical arrows in the first diagram have the ‘unique path lifting property’,
(ii) $`(\stackrel{~}{V}_\mu ,\stackrel{~}{E}_\mu ,)`$ is the skew-product Bratteli diagram associated to the $`G`$-valued labelling $`\mu `$ on $`(\stackrel{~}{V},\stackrel{~}{E},)`$ defined by $`\mu =\lambda \rho `$
(iii) $`\rho ^{}`$ commutes with the $`G`$-action and $`\rho _B^{}`$ is an isomorphism and
(iv) the map $`\stackrel{~}{\pi }:(\stackrel{~}{V}_\mu ,\stackrel{~}{E}_\mu ,)(\stackrel{~}{V},\stackrel{~}{E},)`$ has the property “loops lift to loops”. (Consequently, by Part (I), the quotient of $`(X_\lambda ,T_\lambda )`$ by the action of $`G`$ is isomorphic to $`(\stackrel{~}{X},\stackrel{~}{T})`$.)
Remark. Note that Part (II) is almost a converse to Part (I). If we are given that the quotient of $`(X_\lambda ,T_\lambda )`$ by the action of $`G`$ is isomorphic to $`(X,T)`$, we are not deducing from this that $`\pi :(V_\lambda ,E_\lambda ,)(V,E,)`$ has the property ‘loops lift to loops’. Instead, Part (II) gives us a closely related morphism with the property ‘loops lift to loops’ which can be nicely fit into a square with the given morphism $`\pi `$ as one of its sides and in which all the arrows have the unique path lifting property; the vertical arrows are compatible with the given labelling and its pull-back and induce isomorphisms in the corresponding dynamical systems.
Poof of part (I). By hypothesis $`\pi :(V_\lambda ,E_\lambda ,)(V,E,)`$ has the property ‘loops lifts to loops’. Suppose that $`\pi _B:(X_\lambda ,E_\lambda ,)(X,E,)`$ is not an isomorphism, modulo the $`G`$-action on $`X_\lambda `$. Let $`y`$ and $`z`$ be two different $`G`$-orbits in $`X_\lambda `$, which have the same image $`x`$ in $`X`$. We employ the notation of $`(2.9)`$ and $`(2.10)`$. Write
$`y=(y_n)_n,y_n=(y_{n,i})_i,y_{n,i}\varpi _{n,\lambda }`$
$`z=(z_n)_n,z_n=(z_{n,i})_i,z_{n,i}\varpi _{n,\lambda }`$
$`x=(x_n)_n,x_n=(x_{n,i})_i,x_{n,i}\varpi _n`$.
By the hypothesis, $`y`$ and $`z`$ have the image $`x`$ in $`X`$ and also the $`G`$-orbits of $`y`$ and $`z`$ are different. By shifting $`y`$ and $`z`$ by a suitable power of the shift operator $`T_\lambda `$ and replacing $`z`$ if necessary by another element within the orbit $`\{\gamma _gzgG\}`$, we can assume that, for some $`k`$, $`y_{k,0}=z_{k,0}`$ and $`y_{k,1}z_{k,1}`$. In view of (1.8) it follows that
(2.12.1.) $`y_{k,0}`$ is the top floor of a $`\varpi _{k,\lambda }`$-tower $`\varpi (u,g_1)`$
and
(2.12.2.) $`y_{k,1}`$ is the lowest floor of a $`\varpi _{k,\lambda }`$-tower $`\varpi (w,g_2)`$.
In the same way,
(2.12.3.) $`z_{k,0}`$ is the top floor of the $`\varpi _{k,\lambda }`$-tower $`\varpi (u,g_1)`$
and
(2.12.4.) $`z_{k,1}`$ the lowest floor of a $`\varpi _{k,\lambda }`$-tower $`\varpi (w,g_3)`$.
By property ((1.10), (iii)) there exists $`mk`$, a vertex $`(a,h)V_{m,\lambda }`$ and two successive floors $`F^0,F^1`$ of the $`\varpi _{m,\lambda }`$ tower $`\varpi (a,h)`$ such that $`j_{m,k}(F^0)=y_{k,0}`$ and $`j_{m,k}(F^1)=y_{k,1}`$. Similarly, there exists $`nk`$, a vertex $`(b,h^\stackrel{~}{})V_{n,\lambda }`$ and two successive floors $`\stackrel{~}{F}^0,\stackrel{~}{F}^1`$ of the $`\varpi _{n,\lambda }`$-tower $`\varpi (b,h^\stackrel{~}{})`$ such that $`j_{m,k}(\stackrel{~}{F}^0)=y_{k,0}`$ and $`j_{m,k}(\stackrel{~}{F}^1)=y_{k,1}`$. The floors $`F^0`$ and $`F^1`$ correspond to two successive paths from $`\{\}V_{0,\lambda }`$ to $`(a,h)V_{m,\lambda }`$ in the linearly ordered set of paths from $`\{\}`$ to $`(a,h)`$. The path $`F^0`$ is traced by first tracing a path $`F^0[0,k]`$ from level $`0`$ to level $`k`$ and following it by a path from level $`k`$ to level $`m`$. The part $`F^0[0,k]`$ being the truncation $`j_{m,k}(F^0)`$ represents the floor $`y_{k,0}`$. By $`(\mathrm{2.12.1})`$, $`y_{k,0}`$ is the unique maximal path from $`\{\}V_{0,\lambda }`$ to $`(u,g_1)V_{k,\lambda }`$. Similarly, the path $`F^1`$ from level $`0`$ to level $`m`$ initially traces the truncation $`F^1[0,k]=j_{m,k}(F^1)`$ which represents the floor $`y_{k,1}`$ followed by $`F^1[k,m]`$ from level $`k`$ to level $`m`$. By $`(\mathrm{2.12.2})`$, $`y_{k,1}`$ is the unique minimal path (lowest floor) from $`\{\}V_{0,\lambda }`$ to $`(w,g_2)V_{k,\lambda }`$. Since $`F^1`$ is the successor of $`F^0`$ in the lexicographic order of paths from $`\{\}`$ to $`(a,h)V_{m,\lambda }`$, we conclude from the above observation that
(2.12.5) $`F^1[k,m]`$ is the successor of $`F^0[k,m]`$ in the lexicographic order of paths from $`V_{k,\lambda }`$ to $`(a,h)V_{m,\lambda }`$.
Similarly, if $`\stackrel{~}{F}^0[k,n]`$ and $`\stackrel{~}{F}^1[k,n]`$ represent the truncation of $`\stackrel{~}{F}^0`$ and $`\stackrel{~}{F}^1`$ respectively from level $`k`$ to level $`n`$ then
(2.12.6) $`\stackrel{~}{F}^1[k,n]`$ is the successor of $`\stackrel{~}{F}^0[k,n]`$ in the lexicographic order of paths from $`V_{k,\lambda }`$ to $`(a,h^\stackrel{~}{})V_{n,\lambda }`$.
We show how this leads to a contradiction of the property ‘loops lift to loops’. Denote the paths $`F^0[k,m]`$, $`F^1[k,m]`$, $`\stackrel{~}{F}^0[k,n]`$ and $`\stackrel{~}{F}^1[k,n]`$ by $`\beta ,\alpha ,\delta `$ and $`\gamma `$ respectively. Let their images under $`\pi `$ be denoted by $`\beta ^{},\alpha ^{},\delta ^{}`$ and $`\gamma ^{}`$ respectively.
We will observe that with the above data we have something on the left which is almost a loop, but not actually a loop eventhough its image on the right is a loop (cf. (2.11)).
(object on the right which is a loop) $`\alpha ^{}`$ and $`\beta ^{}`$ are paths from $`V_k`$ ranging at $`a`$. Also $`\gamma ^{}`$ and $`\delta ^{}`$ are paths from $`V_k`$ ranging at $`b`$. $`\alpha ^{}`$ is the successor of $`\beta ^{}`$. For paths from $`V_k`$ to $`V_n,\gamma ^{}`$ is the successor of $`\delta ^{}`$. The paths $`\alpha ^{}`$ and $`\gamma ^{}`$ have the same source in $`V_k`$, $`\beta ^{}`$ and $`\delta ^{}`$ have the same source in $`V_k`$. (Thus, one has a loop:- going from the range of $`\alpha ^{}`$ to the source of $`\beta ^{}`$ via $`\beta ^{}`$, then to the range of $`\delta ^{}`$ via $`\delta ^{}`$, then to the source of $`\alpha ^{}`$ via $`\gamma ^{}`$ and then to the range of $`\alpha ^{}`$ via $`\alpha ^{}`$.)
(object on the left lying over the above loop) $`(a,h)V_{m,\lambda },(b,h^\stackrel{~}{})V_{n,\lambda }`$, $`\pi (a,h)=a,\pi (b,h^\stackrel{~}{})=b`$. $`\alpha `$ (resp. $`\beta )`$ is the unique path from $`V_{k,\lambda }`$ to $`V_{m,\lambda }`$ lying above $`\alpha ^{}`$, (resp. $`\beta ^{})`$ and ranging at $`(a,h)`$. Similarly, $`\gamma `$ (resp. $`\delta )`$ is the unique path from $`V_{k,\lambda }`$ to $`V_{n,\lambda }`$, lying above $`\gamma ^{}`$ (resp. $`\delta ^{})`$ and ranging at $`(b,h^\stackrel{~}{})`$. $`\alpha `$ is the successor of $`\beta `$. $`\gamma `$ is the successor of $`\delta `$.
(object on the left is almost a loop) Moreover source of $`\beta =`$ source of $`\delta `$.
But source of $`\alpha `$ source of $`\gamma `$, which contradicts the ‘loops lifting to loops’ property. $`\mathrm{}`$
Now we begin the proof of the part II in the statement of the theorem.
2.13. We will now define two nested sequences of K-R partitions of $`X_\lambda `$, both acted on by the group $`G`$. Recall $`V_{n,\lambda }=V_n\times G`$ (for $`n1)`$. For $`(v,g)V_n\times G`$, let $`y`$ be a path from $`\{\}`$ to $`(v,g)`$ in $`(V_\lambda ,E_\lambda ,)`$. So, $`y`$ is a ‘floor’ belonging to the $`\varpi _{n,\lambda }`$\- tower $`\varpi (v,g)`$ parametrized by $`(v,g)`$. Put $`_y=\{x=(x_1,x_2,\mathrm{},x_n,\mathrm{})X_\lambda x_{n,0}=y\}`$
$$𝒫_{n,\lambda }=\{_yy\varpi (v,g),(v,g)V_{n,\lambda }\}.$$
Then $`\{𝒫_{n,\lambda }\}_n`$ is a nested sequence of K-R partitions of $`X_\lambda `$ acted on by $`G`$. But, the topology of $`X_\lambda `$ need not be spanned by the collection of clopen sets $`\{_y\},(y\varpi (v,g),(v,g)V_{n,\lambda },n)`$. In contrast, the topology of $`X_\lambda `$ is indeed spanned by the collection of clopen sets in another nested sequence $`\{𝒬_{n,\lambda }\}_n`$ of K-R partitions, defined below. Let $`\varpi =\varpi (u,g_1),\varpi ^{}=\varpi (v,g_2),\varpi ^{\prime \prime }=\varpi (w,g_3)`$ be three $`\varpi _{n,\lambda }`$-towers and $`y`$ a floor of $`\varpi ^{}`$. For any $`xX_\lambda `$ and for any $`n`$ if $`x_{n,i}`$ is a floor of a $`\varpi _{n,\lambda }`$-tower $`\overline{\varpi }`$, then for some $`a,b`$ such that $`aib`$, the segment $`x_n[a,b]`$ is just the sequence of floors in $`\overline{\varpi }`$. We define
$`(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)=`$ the clopen subset of $`_y`$ consisting of the elements $`x=(x_1,x_2,\mathrm{},x_n,\mathrm{})`$ with the property that for some $`a_1<a_20<a_3<a_4`$, the segment $`x_n[a_1,a_21]`$ is the sequence of floors of $`\varpi `$, the segment $`x_n[a_2,a_31]`$ is the sequence of floors of $`\varpi ^{}`$ and the segment $`x_n[a_3,a_4]`$ is the sequence of floors of $`\varpi ^{\prime \prime }`$. Some of the sets $`(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)`$ may be empty, but the non-empty sets $`(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)`$ form a K-R partition which we denote by $`𝒬_{n,\lambda }`$. For fixed $`\varpi ,\varpi ^{},\varpi ^{\prime \prime }`$ the subcollection $`\{(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)\}`$ as $`y`$ varies through the floors of $`\varpi ^{}`$, is a $`𝒬_{n,\lambda }`$-tower parametrized by $`[(u,g_1),(v,g_2),(w,g_3)]`$. We denote this $`𝒬_{n,\lambda }`$-tower by $`𝒮_{(\varpi ,\varpi ^{},\varpi ^{\prime \prime })}`$. The floors of the tower $`𝒮_{(\varpi ,\varpi ^{},\varpi ^{\prime \prime })}`$ are $`\{(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)\}`$ as $`y`$ runs through the sequence of floors of $`\varpi ^{}`$.
2.14. Lemma. $`\{𝒬_{n,\lambda }\}_n`$ is a nested sequence of K-R partitions of $`X_\lambda `$ acted on by $`G`$. The topology of $`X_\lambda `$ is spanned by the clopen sets in this sequence of partitions.
## Proof.
Each $`𝒬_{n,\lambda }=\{(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)\}`$ (for a fixed $`n)`$ is a K-R partition. This is evident on going through the definitions. Since the infimum of the height of $`\varpi _{n,\lambda }`$-towers $`(\varpi ,\varpi ^{}`$ etc.) goes to infinity as $`n\mathrm{}`$, it is evident from the definition of the topology of $`X_\lambda `$ that the clopen sets $`(\varpi ,\varpi ^{},\varpi ^{\prime \prime };y)`$, (for various $`n,\varpi _{n,\lambda }`$\- towers $`\varpi ,\varpi ^{},\varpi \mathrm{"}`$, floors $`y`$ of $`\varpi ^{})`$ span the topology of $`X_\lambda `$. $`\mathrm{}`$
2.15. Remark. The above construction involving ‘triples’ can be introduced starting with any simple ordered Bratteli diagram, (not necessarily the diagram $`(V_\lambda ,E_\lambda ,)`$ which had a $`G`$-action). For this, let $`(V,E,)`$ be an arbitrary simple, ordered Bratteli diagram. Define $`(V^𝒬,E^𝒬,)`$ as follows: $`V_0^𝒬=\{\}`$, a single point.
$`V_n^𝒬`$ consists of triples $`(u,v,w)V_n\times V_n\times V_n`$ such that for some $`yV_m`$ where $`m>n`$, the level-$`m`$ tower $`\varpi _y`$ passes successively through the level-$`n`$ tower $`\varpi _u`$, then $`\varpi _v`$ and then $`\varpi _w`$. An edge $`\stackrel{~}{e}E_n^𝒬`$ is a triple $`(u,e,w)`$ such that $`e`$ is an edge of $`(V,E)`$ and $`(u,r(e),w)V_n^𝒬`$. Let
1. $`\{e_1,e_2,\mathrm{},e_k\}`$ be all the edges in $`r^1(r(e))`$,
2. $`\{f_1,f_2,\mathrm{},f_{\mathrm{}}\}`$ be all the edges in $`r^1(u)`$ and
3. $`\{g_1,g_2,\mathrm{},g_m\}`$ be all the edges in $`r^1(w)`$.
The sources of $`(u,e_1,w),(u,e_2,w),\mathrm{},(u,e_k,w)`$ are defined to be
$`(s(f_{\mathrm{}}),s(e_1),s(e_2)),(s(e_1),s(e_2),s(e_3)),\mathrm{},(s(e_{k1}),s(e_k),s(g_1))`$
respectively. The range of $`(u,e,w)`$ is of course $`(u,r(e),w)`$.
The map $`(u,v,w)v,(u,e,w)e`$ from $`(V^𝒬,E^𝒬)`$ to $`(V,E)`$ has unique path lifting property; in particular it gives rise to the ordered Bratteli diagram $`(V^𝒬,E^𝒬,)`$.
(Example (2.6) continued). We illustrate the tripling construction described above for the substitutional system $`A=\{X,Y\},\sigma (X)=XXY,\sigma (Y)=XYY`$. One can easily verify that the tripling yields a stationary ordered Bratteli diagram corresponding to substitutional system with alphabet set $`A^𝒬A\times A\times A`$ described by
$$A^𝒬=\{[X,X,Y],[Y,X,X],[X,Y,X],[X,Y,Y],[Y,X,Y],[Y,Y,X]\}$$
and substitution $`\sigma ^𝒬`$ described by
$$\begin{array}{ccc}\sigma ^𝒬([X,X,Y])\hfill & =\hfill & [Y,X,X][X,X,Y][X,Y,X]\hfill \\ \sigma ^𝒬([Y,X,X])\hfill & =\hfill & [Y,X,X][X,X,Y][X,Y,X]\hfill \\ \sigma ^𝒬([X,Y,X])\hfill & =\hfill & [Y,X,Y][X,Y,Y][Y,Y,X]\hfill \\ \sigma ^𝒬([X,Y,Y])\hfill & =\hfill & [Y,X,Y][X,Y,Y][Y,Y,X]\hfill \\ \sigma ^𝒬([Y,X,Y])\hfill & =\hfill & [Y,X,X][X,X,Y][X,Y,X]\hfill \\ \sigma ^𝒬([Y,Y,X])\hfill & =\hfill & [Y,X,Y][X,Y,Y][Y,Y,X]\hfill \end{array}$$
In this example, since $`\sigma `$ is a proper substitution, $`\sigma ^𝒬`$ depends only on the middle coordinate. Later we will illustrate with a more interesting example.
2.16. The ordered Bratteli diagrams $`B^𝒬(\lambda )`$ and $`\overline{B}^𝒬(\lambda )`$. The construction $`(1.6)`$ of an ordered Bratteli diagram applied to the nested sequence of K-R partitions $`\{𝒬_{n,\lambda }\}_n`$ gives rise to an ordered Bratteli diagram. The vertices $`V_{n,\lambda }^𝒬`$ are in 1-1 correspondence with the towers $`𝒮_{[\varpi (u,g_1),\varpi (v,g_2),\varpi (w,g_3)]}`$ of $`𝒬_{n,\lambda }`$. Note that not all choices $`[(u,g_1),(v,g_2),(w,g_3)]`$ may give rise to a non-empty set $`𝒮_{[\varpi (u,g_1),\varpi (v,g_2),\varpi (w,g_3)]}`$. We may denote the vertex corresponding to (non-empty) $`𝒮_{[\varpi (u,g_1),\varpi (v,g_2),\varpi (w,g_3)]}`$ by $`[(u,g_1),(v,g_2),(w,g_3)]`$. If $`[(u,g_1),(v,g_2),(w,g_3)]V_{n,\lambda }^𝒬`$, then $`[(u,gg_1),(v,gg_2),(w,gg_3)]V_{n,\lambda }^𝒬,gG`$. So $`G`$ acts on $`V_{n,\lambda }^Q`$ freely if $`n1`$. As in $`(1.6)`$ the edges connecting $`V_{n,\lambda }^𝒬`$ and $`V_{n+1,\lambda }^𝒬`$ and the linear order between edges with the same range simply reflects the sequence of $`𝒬_{n,\lambda }`$-towers traversed by a given $`𝒬_{n+1,\lambda }`$\- tower. It is then clear that $`G`$-acts on the ordered Bratteli diagram $`(V_\lambda ^𝒬,E_\lambda ^𝒬,)`$. We denote by $`(\overline{V}_\lambda ^𝒬,\overline{E}_\lambda ^𝒬,)`$ the ordered Bratteli diagram which is the quotient by $`G`$-action on $`(V_\lambda ^𝒬,E_\lambda ^𝒬,)`$. We write $`B^𝒬(\lambda )`$ and $`\overline{B}^𝒬(\lambda )`$ for the Bratteli diagrams $`(V_\lambda ^𝒬,E_\lambda ^𝒬,)`$ and $`(\overline{V}_\lambda ^𝒬,\overline{E}_\lambda ^𝒬,)`$ respectively.
2.17. Examples 2.4 and 2.7 continued. For the ordered Bratteli diagram $`(V_\lambda ,E_\lambda ,)`$ which arose in (2.7) from the labelling $`\lambda `$ defined there, the ‘tripling’ construction $`(2.15)`$ yields a Bratteli diagram $`(V_\lambda ^𝒬,E_\lambda ^𝒬,)`$, which is associated to the primitive substitutional system below. The alphabet set $`A_\lambda ^𝒬`$ consists of the twenty alphabets
1. \[x,x,y\], \[x,y,x\], \[x,y,x’\], \[y,x,x\], \[y’,x,x\]
2. \[y,x,y’\], \[y’,x,y’\],\[y’,y,x’\], \[y’,y,x\], \[x’,y,y’\]
3. \[x’,x’,y’\], \[x’,y’,x’\], \[x’,y’,x\], \[y’,x’,x’\], \[y,x’,x’\]
4. \[y’,x’,y\], \[y,x’,y\],\[y,y’,x\], \[y,y’,x’\], \[x,y’,y\]
The substitution $`\sigma _\lambda ^𝒬:A_\lambda ^𝒬(A_\lambda ^𝒬)^+`$ is given by
$`\sigma _\lambda ^𝒬([x,x,y])=[y,x,x][x,x,y][x,y,x]`$
$`\sigma _\lambda ^𝒬([x,y,x])=[y,x,y^{}][x,y^{},y][y^{},y,x]`$
$`\sigma _\lambda ^𝒬([x,y,x^{}])=[y,x,y^{}][x,y^{},y][y^{},y,x]`$
$`\sigma _\lambda ^𝒬([y,x,x])=[y,x,x][x,x,y][x,y,x]`$
$`\sigma _\lambda ^𝒬([y^{},x,x])=[y^{},x,x][x,x,y][x,y,x]`$
$`\mathrm{}`$ etc..
The general rule for the substitution $`\sigma _\lambda ^𝒬`$ is expalined in the more general situation in $`(2.18)`$, below. The group $`_2`$ acts on $`(A_\lambda ^𝒬,\sigma _\lambda ^𝒬)`$ and on $`(V_\lambda ^𝒬,E_\lambda ^𝒬,)`$ via the transposition $`xx^{},yy^{}`$. The quotient substitutional system $`(\overline{A}_\lambda ^𝒬,\overline{\sigma }_\lambda ^𝒬)`$ with ten alphabets maps to $`(A,\sigma )`$ of example (2.6). The induced map between the dynamical systems is not an isomorphism.
2.18. Tripling for a substitutional system $`(A,\sigma )`$. If $`(A,\sigma )`$ is a substitutional system, define $`A^𝒬`$ to be the unique smallest subset of $`A\times A\times A`$ with the property that $`(a,b,c)A^𝒬`$ if and only if the word $`abc`$ occurs as a subword of $`\sigma ^n(d)`$ for some $`dA`$ and some $`n`$. Define
$$\sigma ^𝒬:A^𝒬(A^𝒬)^+$$
by $`\sigma ^𝒬[(a,b,c)]=(a_m,b_1,b_2)(b_1,b_2,b_3)\mathrm{}(b_{n2},b_{n1},b_n)(b_{n1},b_n,c_1)`$, where $`\sigma (b)=b_1b_2\mathrm{}b_n`$, and $`a_m`$ is the last alphabet in $`\sigma (a)`$, while $`c_1`$ is the first alphabet in $`\sigma (c)`$. Suppose moreover that the ordered Bratteli diagram associated to $`(A,\sigma )`$ is equipped with a stationary labelling $`\lambda `$. Then following the construction of $`(2.5)`$ we get a skew-product substitutional system $`(A_\lambda ,\sigma _\lambda )`$. The tripling described above applied to this $`(A_\lambda ,\sigma _\lambda )`$ gives a substitutional system $`(A_\lambda ^𝒬,\sigma _\lambda ^𝒬)`$.
2.19. Even if we start with a primitive (cf. \[DHS,3.3.1\]) and proper substitutional system and a stationary labelling on the edges with the additional property that the maximal and minimal edges are both labelled by the trivial element of the group, the property “loops lift to loops” may generally fail; the quotient of the dynamical system associated to $`(A_\lambda ^𝒬,\sigma _\lambda ^𝒬)`$ by the $`G`$-action will in general be a nontrivial extension of the dynamical system of $`(A,\sigma )`$. This is an important difference between our construction of skew-product systems arising from $`G`$-valued labellings and Matui’s construction of skew-product systems arising from $`G`$-valued cocycles.
2.20. It now remains to finish the proof of part (II) of Theorem (2.12). Our construction of $`\overline{B}^𝒬(\lambda )=(\overline{V}_\lambda ^𝒬,\overline{E}_\lambda ^𝒬,)`$ in $`(2.16)`$ was motivated precisely to serve as a candidate for the $`\stackrel{~}{B}`$ in the statement of part II of the theorem. Thus, we set $`\stackrel{~}{V}=\overline{V}_\lambda ^𝒬`$, $`\stackrel{~}{E}=\overline{E}_\lambda ^𝒬`$, $`\stackrel{~}{B}=(\stackrel{~}{V},\stackrel{~}{E},)=(\overline{V}_\lambda ^𝒬,\overline{E}_\lambda ^𝒬,)=\overline{B}^𝒬(\lambda )`$. A vertex of $`\overline{V}_{n,\lambda }^𝒬`$ is represented by the $`G`$-orbit of a triple $`[(u,g_1),(v,g_2),(w,g_3)]`$. Define $`\rho :\overline{V}_\lambda ^𝒬V`$ by sending the above vertex to $`vV_n`$. The set $`\overline{E}_{n,\lambda }^𝒬`$ of edges of $`\overline{B}^𝒬(\lambda )`$ from $`\overline{V}_{n1,\lambda }^𝒬`$ to $`\overline{V}_{n,\lambda }^𝒬`$ ranging at $`[(u,g_1),(v,g_2),(w,g_3)]`$ is represented by the triple $`[(u,g_1),(e,g_2),(w,g_3)]`$ where $`e`$ is an edge from $`V_{n1}`$ to $`V_n`$ ranging at $`v`$. Define $`\rho :\overline{E}_\lambda ^𝒬E`$ by sending the above edge (namely the $`G`$-orbit of $`[(u,g_1),(e,g_2),(w,g_3)]`$) to $`e`$.
Define a labelling $`\mu `$ on $`\overline{E}_\lambda ^𝒬`$ by $`\mu =\lambda \rho `$. We show that the corresponding skew-product of $`(\overline{V}_\lambda ^𝒬,\overline{E}_\lambda ^𝒬,)`$ by $`\mu `$ is precisely $`(V_\lambda ^𝒬,E_\lambda ^𝒬,)`$. For this we define a map
$$\mathrm{\Phi }:(V_\lambda ^𝒬,E_\lambda ^𝒬,)(\overline{V}_\lambda ^𝒬\times G,\overline{E}_\lambda ^𝒬\times G,)$$
as follows: let $`x=[(u,g_1),(v,g_2),(w,g_3)]V_\lambda ^𝒬`$. Let the $`G`$-orbit of $`x`$ be denoted by $`\overline{x}\overline{V}_\lambda ^𝒬`$. Then, define $`\mathrm{\Phi }(x)=(\overline{x},g_2)`$. Similarly, let $`\eta =[(u,g_1),(e,g_2),(w,g_3)]`$ be an edge of $`E_\lambda ^𝒬`$. Let the $`G`$-orbit of $`\eta `$ be denoted by $`\overline{\eta }\overline{E}_\lambda ^𝒬`$. Then, define $`\mathrm{\Phi }(\eta )=(\overline{\eta },g_2)`$. Then $`\mathrm{\Phi }`$ is an isomorphism between the two ordered Bratteli diagrams.
Referring to the requirements in the statement of Part (II) of Theorem $`(2.12)`$, the map $`\rho ^{}:(V_\lambda ^𝒬,E_\lambda ^𝒬,)(V_\lambda ,E_\lambda ,)`$ is easy to define:
$$\rho ^{}[(u,g_1),(v,g_2),(w,g_3)]=(v,g_2)\text{and}\text{ }\rho ^{}[(u,g_1),(e,g_2),(w,g_3)]=(e,g_2).$$
The map $`\rho ^{}:B^𝒬(\lambda )=(V_\lambda ^𝒬,E_\lambda ^𝒬,)(V_\lambda ,E_\lambda ,)=B(\lambda )`$ induces a map $`\rho _B^{}`$ between the corresponding dynamical systems $`X_{B^𝒬(\lambda )}`$ and $`X_{B(\lambda )}`$. The inverse $`\beta `$ of this map is defined as follows:
Let $`zX_{B(\lambda )}`$. So $`z=(z_1,\mathrm{},z_n,\mathrm{})`$, where $`z_n=(z_{n,i})_i(\varpi _{n,\lambda })^{}`$. Fix $`n`$ and $`i`$. So $`z_{n,i}`$ is a floor of a $`\varpi _{n,\lambda }`$-tower $`𝒮_{(v,g_2)}`$. Choose $`a,b,c,d`$ such that $`a<bi<c<d`$, the interval segment $`z_n[b,c1]`$ is the sequence of all the floors of the tower $`𝒮_{(v,g_2)}`$, $`z_n[a,b1]`$ is the sequence of all the floors of some $`\varpi _{n,\lambda }`$-tower $`𝒮_{(u,g_1)}`$ and $`z_n[c,d1]`$ is the sequnce of all the floors of some $`\varpi _{n,\lambda }`$-tower $`𝒮_{(w,g_3)}`$. Define
$$y=(y_1,\mathrm{},y_n,\mathrm{})X_{B^𝒬(\lambda )}\text{where }y_n=(y_{n,i})_i\left(\varpi _{n,\lambda }^𝒬\right)^{}$$
by $`y_{n,i}=((u,g_1),(v,g_2),(w,g_3);z_{n,i})`$. This map defines an inverse of $`\rho _B^{}`$.
Finally, we prove the property ‘loops lift to loops’ for the map
$$(V_\lambda ^𝒬,E_\lambda ^𝒬,)(\overline{V}_\lambda ^𝒬,\overline{E}_\lambda ^𝒬,).$$
To verify the property ‘loops lift to loops’ for the above map, we start with some data on the left-side which is ‘almost’ a loop, we make a further assumption that the image is actually a loop and then we have to conclude that the data we started with on the left-side must in fact be a loop.
2.21. The data on the left which is almost a loop. Let $`\widehat{u}V_{m,\lambda }^𝒬`$ and $`\widehat{v}V_{n,\lambda }^𝒬`$. Let $`km`$ and $`kn`$. Let $`\alpha `$ and $`\beta `$ be paths from $`V_{k,\lambda }^𝒬`$ ranging at $`\widehat{u}`$. Let $`\gamma ,\delta `$ be paths from $`V_{k,\lambda }^𝒬`$ ranging at $`\widehat{v}`$. Assume that $`\beta `$ and $`\delta `$ have the same source in $`V_{k,\lambda }^𝒬`$.
2.22. The data ‘image on the right is actually a loop’. Denote by $`\alpha ^{},\beta ^{},\gamma ^{},\delta ^{}`$ the paths in $`(\overline{V}_\lambda ^𝒬,\overline{E}_\lambda ^𝒬,)`$ which are the images of $`\alpha ,\beta ,\gamma ,\delta `$. Suppose that $`\alpha ^{}`$ is the successor of $`\beta ^{}`$ and that $`\gamma ^{}`$ is the successor of $`\delta ^{}`$. Assume that $`\alpha ^{}`$ and $`\gamma ^{}`$ have the same source in $`\overline{V}_{k,\lambda }^𝒬`$ and that $`\beta ^{}`$ and $`\delta ^{}`$ have the same source in $`\overline{V}_{k,\lambda }^𝒬`$.
To conclude that the data on the left must in fact be a loop we have to show that $`\alpha `$ and $`\gamma `$ have the same source in $`V_{k,\lambda }^𝒬`$.
Let source of $`\beta =`$ source of $`\delta =[(a,g_1),(b,g_2),(c,g_3)]V_{m,\lambda }`$. Since $`\alpha `$ is the succssor of $`\beta `$, necessarily the source of $`\alpha `$ has the form $`[(b,g_2),(c,g_3),(d,g_4)]`$ for some $`dV_k`$ and $`g_4G`$. Likewise, since $`\gamma `$ is the successor of $`\delta `$ the source of $`\gamma `$ must be of the form $`[(b,g_2),(c,g_3),(d^{},g_5)]`$, for some $`d^{}V_k`$ and $`g_5G`$. But it is a consequence of the assumption “source of $`\alpha ^{}`$ = source of $`\gamma ^{}`$” that $`[(b,g_2),(c,g_3),(d,g_4)]`$ and $`[(b,g_2),(c,g_3),(d^{},g_5)]`$ must lie in the same $`G`$-orbit, i.e., for some $`hG`$, $`[(b,g_2),(c,g_3),(d,g_4)]`$ = $`[(b,hg_2),(c,hg_3),(d^{},hg_5)]`$.
Therefore, we conclude that $`d=d^{}`$ and $`g_4=g_5`$. Thus, the property “loops lift to loops” is verified.
This ends the proof of the theorem. $`\mathrm{}`$
2.23. We already gave an example (cf. 2.17) where the map $`\rho _B`$ of Theorem 2.12 is not an isomorphism. To give another example let $`G=𝔖_3`$ be the group consisting of the six permutations of the symbols $`\{1,2,3\}`$. Let $`s`$ and $`t`$ be the elements of $`𝔖_3`$ corresponding to the 2-cycle $`(1,2)`$ and the 3-cycle $`(1,2,3)`$ respectively. Let $`A`$ be any alphabet set $`\{a_\tau \tau 𝔖_3\}`$. Define a substitution $`\sigma :AA`$ by $`\sigma (a_\tau )=a_{\tau t}a_{\tau s}`$. The group $`𝔖_3`$ acts on $`(A,\sigma )`$ by $`\gamma _g(a_\tau )=a_{g\tau }`$. This is a skew-product system (cf. 2.5) associated to a stationary labelling $`\lambda `$ on the edges of the $`2adic`$ odometer system whose image consists of the two elements $`\{s,t\}𝔖_3`$. In this example the map $`\rho _B`$ of the theorem from the quotient of the dynamical system for $`(A,\sigma )`$ by $`𝔖_3`$ to the $`2adic`$ odometer system is not an isomorphism. This follows from an observation of Matui \[M, see para before lemma 2.2\] and from the fact that $`𝔖_3`$ is not a cyclic group.
2.24. Finally it follows from the assertion part (II) (i) of the theorem and from \[M, Lemma 2.2\] that there must exist a continuous function $`c:\overline{X}_{B(\lambda )}^𝒬G`$ such that $`(X_{B(\lambda )}^𝒬,T_{B(\lambda )}^𝒬)(\overline{X}_{B(\lambda )}^𝒬\times G,\mathrm{\Psi })`$ where $`\mathrm{\Psi }(y,g)=(\overline{T}_{B(\lambda )}^𝒬(y),gc(\overline{T}_{B(\lambda )}^𝒬(y))`$.
A function $`c`$ with this property is easy to describe. Let $`y`$ belong to level $`n`$ tower of $`\overline{X}_{B(\lambda )}^𝒬`$ parametrized by the $`G`$-orbit of $`[(u,g_1),(v,g_2),(w,g_3)]\overline{V}_{n,\lambda }^𝒬`$. Define (for $`n1`$) $`c_n(y)=g_1^1g_2`$ if $`y`$ belongs to the bottom floor of the above tower ; otherwise, define $`c_n(y)=1G`$. All the cocycles $`c_n`$ are cohomologous ($`\mathrm{𝑐𝑓}.`$\[M, definition 2.1\]).
2.25. We also remark that Matui’s (cf. \[M\]) dynamical systems $`(X\times G,\psi )`$ of the form $`\psi (x,g)=(T(x),gc(T(x))`$ for a $`G`$-valued cocycle $`c`$ on $`X`$ arise as the skew-product system of this article associated to a particular $`G`$-valued labelling of the edges of a Bratteli diagram for $`(X,T)`$. In fact, let us use Matui’s notation in \[M, p.140\]. There he considers a nested sequence of K-R partitions $`(𝒫_n)_n`$ in $`X`$ such that $`c`$ is constant on the sets of the partition and a nested sequence of K-R partitions $`(𝒬_n)_n`$ in $`X\times G`$. They are related as follows:- denoting, as Matui does (cf. \[M,p.138\]), the floors of $`𝒫_n`$ in $`X`$ by $`X(n,v,k)`$, where $`X(n,v,k)`$ is the $`k^{th}`$ floor of a level-$`n`$ tower $`𝒮_v`$, the sets of the partition $`𝒬_n`$ in $`X\times G`$ are $`X(n,v,k)\times \{g\},(gG)`$,(\[loc.cit ,p.140\]). Let $`(V,E,)`$ be the Bratteli diagram of $`(𝒫_n)_n`$. Let $`wV_{n+1}`$. Let $`(e_1,e_2,\mathrm{},e_{\mathrm{}})`$ be the edges in $`r^1(w)`$ enumerated in the linear order in $`r^1(w)`$. Let $`h(v)`$ denote the height of the K-R tower $`𝒮_v`$ for a vertex $`vV`$. For $`i=1,2,\mathrm{},\mathrm{}`$, choose $`(k_i)_i`$, such that
1. $`1k_1<k_2<k_3<\mathrm{}<k_{\mathrm{}}h(w)`$ and
2. $`X(n+1,w,k_i)X(n,v_i,h(v_i))`$ for some $`v_iV_n`$.
There is a unique level-$`(n+1)`$ tower for the K-R partition $`𝒬_n`$ of $`X\times G`$ whose top floor is $`X(n+1,w,h(w))\times \{e\}`$ where $`eG`$ is the identity element of $`G`$. Let the floors of this tower be $`X(n+1,w,1)\times \{g_1\},X(n+1,w,2)\times \{g_2\},X(n+1,w,3)\times \{g_3\},\mathrm{},X(n+1,w,h(w))\times \{g_{h(w)}\}`$. Then the labelling of $`e_i`$ is $`g_{k_i}`$.
We end with a definition and a comment and indicate scope for further progress.
2.26. Definition. Let $`\lambda `$ and $`\mu `$ be two $`G`$-valued edge labellings of $`(V,E,)`$. We say that $`\lambda `$ and $`\mu `$ are cohomologous if there exists a function $`\beta :VG`$ such that
$$\beta (r(e))\mu (e)=\lambda (e)\beta (s(e)),\text{for }eE_n,(n1).$$
Such a map $`\beta `$ gives rise to an isomorphism
$$\mathrm{\Phi }_\beta :(V_\lambda ,E_\lambda ,)(V_\mu ,E_\mu ,)$$
where $`\mathrm{\Phi }_\beta (v,g)=(v,g\beta (v))\text{ and }\mathrm{\Phi }_\beta (e,g)=(e,g\beta (r(e)))`$. This notion is related to Matui’s definition as follows: recall $`(\mathrm{𝑐𝑓}.[M,2.1])`$ that Matui calls two $`G`$-valued cocycles $`c,c^{}`$ cohomologous if $``$ a continuous function $`b:XG`$ such that $`c(T(x))b(T(x))=b(x)c^{}(T(x))`$. Let $`\lambda \text{ and }\mu `$ be the labellings associated to $`c\text{ and }c^{}`$, respectively as in (2.25). We show that $`\lambda `$ and $`\mu `$ are cohomologous. Let us use the notation of (2.24). By doing a telescoping of $`(V,E,)`$ if necessary, we can suppose that $`c,c^{}`$ and $`b`$ are all constant on the sets of the K-R partition $`𝒫_1`$. Then we just have to take
$$\beta (v)=b(x),\text{ where }x\text{ the top floor of the tower }𝒮_v$$
for a vertex $`vV`$.
We do not know of any example of a substitutional system $`(A_\lambda ,\sigma _\lambda )`$ (see $`2.5`$) arising from a stationary labelling $`\lambda `$ with values in a non-cyclic group $`G`$ , such that $`(A,\sigma )`$ and $`(A_\lambda ,\sigma _\lambda )`$ are both Toeplitz flows defined below as usual (cf. eg.\[GJ,p.1695\]).
2.27. Definition. A Toeplitz sequence is a non-periodic sequence $`\eta =(\eta _n)_n`$ in $`A^{}`$, where $`A`$ is any finite alphabet set, so that for each $`m`$ there exists $`n`$ so that $`\eta _m=\eta _{m+kn}`$ for all $`k`$.
The following gives such an example with $`G`$ a cyclic group of order $`k`$. Let $`(A,\sigma )`$ be a primitive aperiodic proper substituion of constant length $`n`$. By \[GJ, Corollary 9,p.1698\] this gives rise to a Toeplitz flow. Assume that $`n`$ is congruent to 1 mod $`k`$. Let $`z`$ denote a generator for $`G`$. We take the stationary labelling $`\lambda `$ which assigns $`z^i`$ to the $`i^{th}`$ edge ranging at any vertex of the stationary Bratteli diagram for $`(A,\sigma )`$ that we described in subsection 2.5. It is not hard to see that if $`\eta =(a_i)_i`$ is a Toeplitz sequence for $`(A,\sigma )`$ then $`\stackrel{~}{\eta }\text{ }\stackrel{def.}{=}\{(a_i,z^i)\}_i`$ is a Toeplitz sequence for $`(A_\lambda ,\sigma _\lambda )`$.
2.28. Further observations. After the first version of the article was submitted for publication we noticed why, for non-cyclic groups $`G`$, one cannot have Toeplitz sequences occurring in primitive substitutional systems $`(A_\lambda ,\sigma _\lambda )`$ arising from $`G`$-valued stationary edge labellings $`\lambda `$ of another substitutional system $`(A,\sigma )`$. Here’s a proof. Suppose $`\stackrel{~}{\eta }\stackrel{def.}{=}(a_i,g_i)_i`$ is a Toeplitz sequence in $`(A_\lambda ,\sigma _\lambda )`$ (which is assumed to be primitive). Then $`gG,\stackrel{~}{\eta }_g\stackrel{def.}{=}(a_i,gg_i)_i`$ are also Toeplitz sequences. Further, they all admit the same period structure $`p=(p_0,p_1,p_2,\mathrm{})`$ (cf. \[DKL, p.220, GJ,p.1695\]). More importantly, all these Toeplitz sequences lie in the same (minimal) subshift dynamical system $`X_\lambda `$ in the alphabets $`A_\lambda `$ and the $``$\- orbit of any one of these Toeplitz sequences $`\stackrel{~}{\eta }_g`$ is dense in $`X_\lambda `$. Let $`(G_p,1)`$ be the maximal uniformly continuous factor of $`(X_\lambda ,T_\lambda )`$ (cf.\[GJ, p.1696\]) and let $`\pi :(X_\lambda ,T_\lambda )(G_p,1)`$ denote the corresponding factor map. For a Toeplitz sequence $`\omega ,\pi ^1(\pi (\omega ))=\{\omega \}`$.(cf.\[W\], \[GJ,p.1696\], \[DL,Theorem 6,p.167\]). Define elements $`h_g,(gG)`$ of the monothetic group $`G_p`$ by $`h_g=\pi (\stackrel{~}{\eta }_g)`$. They are pairwise distinct. For $`gG`$, denote by $`\tau _g:(X_\lambda ,T_\lambda )(X_\lambda ,T_\lambda )`$ the unique isomorphism which replaces an alphabet $`(a,\theta )`$ by $`(a,g\theta )`$, for $`\theta G`$. Thus, for example, $`\tau _g(\stackrel{~}{\eta }_\theta )=(\stackrel{~}{\eta }_{g\theta })`$.
As observed in \[DKL, sec.2, paragraph 1\], we have the following commutative diagram
$$\begin{array}{ccc}(X_\lambda ,T_\lambda )& \stackrel{\tau _g}{}& (X_\lambda ,T_\lambda )\\ \pi & & \pi \\ (G_p,1)& \stackrel{h_g}{}& (G_p,1)\end{array}$$
where the bottom map denotes group operation by $`h_g`$. Composition of two such diagrams for $`g_1,g_2G`$ yields the diagram for $`g_1g_2`$. Thus, the group $`G`$ sits faithfully as a finite subgroup of $`G_p`$. Hence, it has to be isomorphic to a subgroup of the finite cyclic group $`_{p_i}`$ for one of the essential periods in $`p_0,p_1,p_2,\mathrm{},p_i,\mathrm{}`$. $`\mathrm{}`$
On the other hand, if we drop the ‘Toeplitz’ requirement, the subshift dynamical system of 2.23 in the alphabet set $`A=\{a_\tau \tau 𝔖_3\}`$ provides examples of sequences $`\eta =(a_{\tau _i})_i`$ such that all the sequences $`\eta _g\stackrel{def.}{=}(a_{g\tau _i})_i,(\tau 𝔖_3)`$ lie in the same (minimal) subshift dynamical system in the alphabets $`A`$.
References
1. \] T. Downarowiz, J. Kwiatkowski and Y. Lacroix. A criterion for Toeplitz flows to be topologically isomorphic and applications, Coll. Math.68, (1995) 219-228.
2. \]T. Downarowiz and Y. Lacroix. Almost 1-1 extensions of Furstenberg-Weiss type and applications to Toeplitz flows, Studia Math. 130, (1998) 149-170.
3. \] F. Durand, B. Host and C. Skau. Substitutional dynamical systems, Bratteli diagrams and dimension groups, Ergodic Th. and Dynam. Sys. 19, (1999) 953-993.
4. \] R. Gjerde and Ø. Johansen. Bratteli-Vershik models for Cantor minimal systems: applications to Toeplitz flows, Ergodic Th. and Dynam. Sys. 20, (2000) 1687-1710.
5. \] R.H. Herman, I.F. Putnam and C.F. Skau. Ordered Bratteli diagrams, dimension groups and topological dynamics, Internat. J. Math. 3, (1992) 827-864.
6. \] A. Kumjian and D. Pask. $`C^{}`$-algebras of directed graphs and group actions, Ergodic Th. and Dynam. Sys. 19, (1999) 1503-1519.
7. \] H. Matui. Finite order automorphisms and dimension groups of Cantor minimal systems, J. Math. Soc. Japan 54, (2002) 135-160.
8. \] S. Williams. Toeplitz minimal flows which are not uniquely ergodic, Z. Wahrs. verw. Gebiete. 67, (1984) 95-107.
Lamath, Le Mont Houy
Université de Valenciennes
59313 Valenciennes Cedex 9 (France)
E-mail address : [email protected]
School of Mathematics
Tata Institute of Fundamental Research
Homi Bhabha Road, Colaba
400 005 Mumbai (India)
E-mail address : [email protected]
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# A critical assessment of the pairing symmetry in NaxCoO₂⋅𝑦H2O
## Abstract
We examine each of the symmetry-allowed pairing states of Na<sub>x</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O and compare their properties to what is experimentally and theoretically established about the compound. In this way, we can eliminate the vast majority of states that are technically allowed and narrow the field to two, both of $`f`$-wave type states. We discuss the expected features of these states and suggest experiments that can distinguish between them. We also discuss odd-frequency gap pairing and how it relates to available experimental evidence
Na<sub>x</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O is a novel superconductor which, despite a relatively low superconducting temperature KTHS+03 of only $``$ 5K, has recently attracted substantial experimental and theoretical attention. Much of the interest is driven by an as-of-yet unresolved pairing state that is presumed to be highly unusual and possibly (as we will argue in this Letter, likely) even more unconventional than the $`d`$-wave superconductivity of the cuprates or $`p`$-wave superconductivity of the ruthenates. A survey of the current literature for Na<sub>x</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O reveals that in the two years since its discovery, various groups have proposed $`s`$-waveYKMY ; MYHW+03 ; KKYT03 , $`p`$-wave $`𝐳(x+iy)`$ MKCM+04 ; YYMM05 , $`d`$-wave $`x^2y^2+2ixy`$DSMS04 ; OIM03 ; TWHY+04 ; YYMM05 and various versions of $`f`$ KKYT03 ; YYMM05 ; WHKO+03 and $`i`$-wave states, as well as an odd-frequency triplet $`s`$ state MDJ+b . To our knowledge, Na<sub>x</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O holds the record for the greatest number of different superconducting symmetries proposed for one compound.
Synthesis of single crystal Na<sub>x</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O is difficult and polycrystalline samples often exhibit inhomogeneities in Na distribution and H<sub>2</sub>O accumulation DPC+04 ; BGU+04 . The compound is furthermore chemically unstable at ambient temperature and humidity MLF+03 , making it difficult to handle and characterize. For these reasons, well-reproducible and reliable experimental results that could be expected to unravel the precise superconducting state have been slow to emerge. Still, there are several experimental facts that are rather well established, reproducible and which bear immediate relevance to the superconducting symmetry. Other facts that follow from confirmed knowledge about the crystal and electronic structure allow the exclusion of at least a few of the symmetry-allowed pairing states. An absence of knowledge about which states conform to symmetry requirements, which are excluded by experiment and which are physically unreasonable has frequently lead both theorists and experimentalists to concentrate on those pairing symmetries which are compatible with a specific data set or a specific theory of pairing to the detriment of a broader, consistent picture.
In this Letter, we list all the different symmetry representations that are compatible with a hexagonal crystal structure, according to the seminal work of Sigrist and UedaMSKU91 . Based on what is currently known experimentally about Na<sub>x</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O, we discuss which states can be eliminated from consideration with a reasonable degree of confidence. We will show that surprisingly few candidates survive this scrutiny, and that all of these are highly unconventional and, in a sense, more exotic than either the high-T<sub>c</sub> cuprates or the $`p`$-wave ruthenates. The list of allowed symmetries MSKU91 up to L=3 (i.e. up to the $`f`$-states) is given in Table LABEL:sym.
There are 25 states in this table, excluding the last one which will be discussed separately later. We will show that all but two of them are incompatible with the experimental data.
First, we decide which facts are to be considered as firmly established. Some potentially very important probes, such as the temperature dependence of the Knight shift, are still controversial in the sense that different authors report contradictory results YKMY ; WHKO+03 ; kobayashi ; MKCM+04 . We have therefore singled out three pieces of evidence on which all or practically all publications agree. These are:
Two-dimensionality. Electronic structure calculations for the hydrated compound show an anisotropy in the Fermi velocity of at least an order of magnitude MDJ04 , which is supported by an experimentally measured resistive anisotropy RJBCS+03 ; FCC+03 of 10<sup>3</sup> \- 10<sup>4</sup>, corresponding to a Fermi velocity anisotropy of 30 to 100 (the resistivity anisotropy of the unhydrated, high Na content compound, Na<sub>0.75</sub>CoO$`_2,`$ which should be substantially lower than that of Na<sub>0.3</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O, was foundBCS+04 to be as high as 500), indicating that the transport along $`c`$ is probably incoherent. This is firm evidence that the electronic structure is very strongly 2D.
As found experimentallySPB+04 ; LMH+04 and explained theoretically MDJ05 , the magnetic anisotropy of the unhydrated high-Na compound is very small, primarily because each Co couples with 7 Co atoms in neighboring layers. While there are no data on the magnetic coupling at $`x=0.3,`$ nor for the hydrated compound, one can estimate the reduction in magnetic coupling from the ratio of the squared Fermi velocities, which is about 20. Thus in the hydrated compound, magnetic interaction should also be 2D.
Finally, in an interesting difference from both the cuprates and ruthenates, the Co and O phonons should also be 2D in this system. Of course, water vibration need not be such, but the absence of the hydrogen isotope effect RJBCS+03 clearly indicates their irrelevance for the superconducting pairing.
Therefore, we conclude that the superconducting order parameter in Na <sub>0.3</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O should be 2D.
Absence of superconductivity-induced spontaneous magnetic moments below $`T_c`$. Some of the superconducting states listed in Table LABEL:sym (#9,12) are nonunitary and have a spontaneous magnetization in the superconducting state. Others (#6,18,21) break the time-reversal symmetry for a Cooper pair by virtue of a nonzero pair orbital moment. In both cases, the resulting nonzero local magnetic moments are supposed to be detectableMSKU91 . Note that net magnetization is not present in the latter case, due to domain formation and internal Meissner screening, but crystallographic defects and grain or domain boundaries should still host nonzero local moments. One of the main arguments in favor of the axial $`(x+iy)\widehat{𝐳}`$ state in Sr<sub>2</sub>RuO<sub>4</sub>APM03 was the fact that muon spectroscopy revealed the appearance of disordered static magnetic moments below $`T_c.`$ The accepted interpretation of this finding is that the pairing symmetry has a nonzero orbital moment. Muon spin rotation experiments for Na<sub>0.3</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub> have been reported WHKO+03 ; YJU+04 and no indications of static moments below $`T_c`$ have been found. To our knowledge, there are no other works reporting detection or non-detection of static magnetic moments in this compound. We nevertheless feel confident to include this fact in our compendium for two reasons: i) one of the reports comes from a group WHKO+03 which has performed similar measurements on other superconductors and was previously able to detect local moments in PrOs<sub>4</sub>Sb<sub>12</sub>YAAT+03 , and ii) this is a rare example of an experiment in which poor sample quality makes the effect more pronounced rather than obscuring it.
Therefore, we conclude that neither nonunitary nor $`L0`$ states are possible in Na<sub>0.3</sub>CoO$`{}_{2}{}^{}y`$H <sub>2</sub>O
Absence of a finite superconducting gap. Several experimental groups have reported experiments indirectly probing the density of states (DOS) in the superconducting phase. Such experiments, primarily calorimetry, were instrumental in clarifying the symmetry of pairing in such novel superconductors as SrRuO<sub>4</sub>APM03 and MgB<sub>2</sub>MgB2rev . These experiments measure the temperature dependence of either specific heatHDY+05 ; NORAF+ or relaxation rates (NMRYKMY , NQRKIYI+ ; TFGZ+ or $`\mu `$SRWHKO+03 ; AKAK+03 ). In all work that we are aware of, the authors agree that the low temperature behavior of the DOS is not exponentialnote . As of yet, no group has reported measurements at temperatures low enough to allow for a reasonably confident determination of the exact temperature dependence, but most authors suggest a $`T^3`$ (line nodes) behavior for $`T2`$ K and finite-DOS linear behavior at lower $`T`$.
These results exclude states with a fully developed sizeable gap on all Fermi surfaces.
Armed with these three facts, let us now test the 25 states listed above against them. First, we eliminate all states that have, by symmetry, strong $`z`$-dependence of the order parameter. These are the ten states #3, 7, 8, 9, 15, 19, 20, 21, 24, 25. Of the remaining 15 states, #1, 2, 10, 11, 12, 13, 14, and 18 have no symmetry restriction that would require them to have node lines or points. In general, there exists the possibility of accidental rather than symmetry-induced nodes, or regions with a finite but extremely small order parameter, similar to the so-called “extended $`s`$ ” state, earlier considered for superconducting cuprates. In the case of Na<sub>0.3</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O,, such accidental nodes seem highly unlikely due to its specific fermiology. The Fermi surface of Na<sub>0.3</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O consists of one relatively small nearly circular cylinder, and, possibly (predicted by the theory, but not yet confirmed experimentally) six tiny pockets surrounding the first Fermi surface. Whether the latter actually exist or not, a pairing interaction that would be consistent with the hexagonal symmetry and at the same time enforce sign change of the order parameter on these small Fermi surfaces should itself change sign with a variation of the wave vector on the order of 0.2–0.25 of the Brillouin zone dimension. It is hardly possible to imagine a physically meaningful pairing interaction of this sort.
Finally, the $`\mu SR`$ experiment allows us to exclude the non-unitary state #12. Note that several states that we have already excluded because they are allowed by symmetry to have a full gap, are additionally excluded as having nonzero orbital moment. This leaves us with the six states: #4,5, 16, 17, 22 and 23. There is good reason to believe that the first four are not realized. As an example, we consider the first pair, $`x\widehat{𝐳}`$ and $`y\widehat{𝐳}.`$ In the linear approximation (that is, expanding the free energy to second order in the order parameter), they are degenerate with each other and with state #6, ($`x+iy)\widehat{𝐳}`$ (see Ref. MSKU91 ). Strong coupling effects, the spin-orbit interaction and other effects can, formally, tilt the energy balance in favor of these states, but only in the 4th order in the order parameter (the distinction between the states themselves appears only in the 6th order). However, all three states belong to the same symmetry representation and the first two have node lines, while ($`x+iy)\widehat{𝐳}`$ has full gap. This should lead to a considerable difference in pairing energy that favors the fully gapped state for the same amplitude of the order parameter and makes stabilization of the other states highly questionable.
Excluding states #4,5,16, and 17 on the basis of this argument, we are left with only two possible states: $`x(x^23y^2)\widehat{𝐳}`$ and $`y(y^23x^2)\widehat{𝐳}.`$ Assuming only one Fermi surface, the $`a_{1g}`$ one around the $`\mathrm{\Gamma }`$ point, we cannot make a further distinction between the two. However, if one accepts the Fermi surface from band structure calculations, the latter state has a node line on all $`e_g^{}`$ Fermi surfaces (See Fig. 1). Given the small size of these pockets, the line nodes cause a near loss of pairing for 2/3 of all electrons at the Fermi level, which is energetically unlikely. Therefore, we conclude that if the $`e_g^{}`$ derived Fermi surface pockets actually exist, the most likely superconducting symmetry among all possible superconducting states with an even-frequency order parameter is the $`f`$ state $`x(x^23y^2)\widehat{𝐳}`$.
This is the main result of our paper. However, before concluding, we would like to remark that an odd-frequency triplet $`s\widehat{𝐳}`$ state, proposed earlier by us MDJ+b , is also compatible with the criteria introduced above. It has no orbital moment, it is unitary \[as opposed to another triplet $`s`$ state, $`s(\widehat{𝐱}\pm i\widehat{𝐲})],`$ it is 2D, and, despite being isotropic, it has finite DOS at zero energy, giving rise to the observed nonexponential specific heat and other DOS-sensitive quanitities. As opposed to the $`x\widehat{𝐳}`$ and $`y\widehat{𝐳}`$ case above, $`s\widehat{𝐱}`$ and $`s\widehat{𝐲}`$ are also isotropic and thus do not have any additional disadvantage in terms of the pairing energy. The ground state in this case is defined either by the spin-orbit induced magnetic anisotropy (if the easy magnetization axis is in the plane, $`s\widehat{𝐱}`$ or $`s\widehat{𝐲}`$ is favored, otherwise $`s\widehat{𝐳}`$), or by the spin-orbit induced anisotropy of the pairing interaction.
Finally, we would like comment on the (still controversial) Knight shift experiments. The absence of a Knight shift decay below $`T_c`$ has been taken as a decisive argument in favor of the $`(x\pm iy)\widehat{𝐳}`$ state in Sr<sub>2</sub>RuO<sub>4</sub>APM03 . The $`f`$ state that has emerged from our discussion also corresponds to Cooper pairs with spins in the $`xy`$ plane and thus to a constant Knight shift. The same is true for the $`s\widehat{𝐳}`$ state. On the other hand, $`s\widehat{𝐱}`$ or $`s\widehat{𝐲}`$ would show a reduced, though not exponentially reduced, Knight shift below $`T_c.`$ Of the states with an even frequency gap (#1-25), the states which, formally, should not show a Knight shift reduction below $`T_c`$ along some directions are #4, 5, 6, 15, 22 and 23.
In conclusion, we have shown that, based on established experimental evidence and a knowledge of the electronic structure of Na<sub>0.3</sub>CoO$`{}_{2}{}^{}y`$H<sub>2</sub>O, many superconducting states that are allowed by symmetry can be eliminated from consideration. Of those with an even-frequency gap, only the $`f`$ states, $`x(x^23y^2)\widehat{𝐳}`$, and $`y(y^23x^2)\widehat{𝐳}`$ are fully compatible with what is known about this compound. The former has no line nodes along the $`e_g^{}`$ hole pockets and is therefore energetically favorable, given the existence of $`e_g^{}`$ states at the Fermi level. In terms of odd-freqency gap states, s$`\widehat{𝐱}`$, s$`\widehat{𝐲}`$, and s$`\widehat{𝐳}`$ are all consistent with experimental reports.
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# 1 Introduction
## 1 Introduction
The purpose of the MAGIC Telescope is the observation of high energy gamma radiation from celestial objects. When the gamma quanta hit the earth atmosphere they initiate a cascade of photons, electrons and positrons. The latter radiate short flashes of Cherenkov light which can be recorded by a Cherenkov telescope. The FWHM of the pulses is about 2 ns.
In order to sample this pulse shape with the 300 MSamples/s FADC system , the original pulse is folded with a stretching function leading to a FWHM greater than 6 ns. To increase the dynamic range of the MAGIC FADCs the signals are split into two branches with gains differing by a factor 10. Figure 1a) shows a typical average of identical signals.
In order to discriminate the small signals from showers in the energy range below 100 GeV against the light of the night sky (LONS) the highest possible signal to noise ratio, signal reconstruction resolution and a small bias are important.
Monte Carlo (MC) based simulations predict different time structures for gamma and hadron induced shower images as well as for images of single muons . An accurate arrival time determination may therefore improve the separation power of gamma events from the background events. Moreover, the timing information may be used in the image cleaning to discriminate between pixels whose signal belongs to the shower and pixels which are dominated by randomly timed background noise.
## 2 Digital Filter
The goal of the digital filtering method is to optimally reconstruct from FADC samples the amplitude and arrival time of a signal whose shape is known. Thereby, the noise contributions to the amplitude and arrival time reconstruction are minimized.
For the digital filtering method, three assumptions have to be made:
* The normalized signal shape has to be always constant.
* The noise properties must be constant, especially independent of the signal amplitude.
* The normalized noise auto-correlation has to be constant.
Due to the artificial pulse stretching by about 6 ns on the receiver board all three assumptions are fullfilled to a good approximation. For a more detailed discussion see .
Let $`g(t)`$ be the normalized signal shape, $`E`$ the signal amplitude and $`\tau `$ the time shift between the physical signal and the predicted signal shape. Then the time dependence of the signal, $`y(t)`$, is given by $`y(t)=Eg(t\tau )+b(t),`$ where $`b(t)`$ is the time-dependent noise contribution. For small time shifts $`\tau `$ (usually smaller than one FADC slice width), the time dependence can be linearized. Discrete measurements $`y_i`$ of the signal at times $`t_i(i=1,\mathrm{},n)`$ have the form $`y_i=Eg_iE\tau \dot{g}_i+O(\tau ^2)+b_i`$, where $`\dot{g}(t)`$ is the time derivative of the signal shape.
The correlation of the noise contributions at times $`t_i`$ and $`t_j`$ can be expressed in the noise autocorrelation matrix $`𝑩`$: $`B_{ij}=b_ib_jb_ib_j`$. Figure 2 shows the noise autocorrelation matrix for an open camera. It is dominated by LONS pulses shaped to 6 ns.
The signal amplitude $`E`$, and the product $`E\tau `$ of amplitude and time shift, can be estimated from the given FADC measurements $`𝒚=(y_1,\mathrm{},y_n)`$ by minimizing the deviation of the measured FADC slice contents from the known pulse shape with respect to the known noise auto-correlation: $`\chi ^2(E,E\tau )=(𝒚E𝒈E\tau \dot{𝒈})^T𝑩^1(𝒚E𝒈E\tau \dot{𝒈})+O(\tau ^2)`$ (in matrix form). This leads to the following solution for $`\overline{E}`$ and $`\overline{E\tau }`$:
$$\overline{E}=𝒘_{\text{amp}}^T(t_{\text{rel}})𝒚+O(\tau ^2)\mathrm{with}𝒘_{\text{amp}}(t_{\text{rel}})=\frac{(\dot{𝒈}^T𝑩^1\dot{𝒈})𝑩^1𝒈(𝒈^T𝑩^1\dot{𝒈})𝑩^1\dot{𝒈}}{(𝒈^T𝑩^1𝒈)(\dot{𝒈}^T𝑩^1\dot{𝒈})(\dot{𝒈}^T𝑩^1𝒈)^2},$$
(1)
$$\overline{E\tau }=𝒘_{\text{time}}^T(t_{\text{rel}})𝒚+O(\tau ^2)\mathrm{with}𝒘_{\text{time}}(t_{\text{rel}})=\frac{(𝒈^T𝑩^1𝒈)𝑩^1\dot{𝒈}(𝒈^T𝑩^1\dot{𝒈})𝑩^1𝒈}{(𝒈^T𝑩^1𝒈)(\dot{𝒈}^T𝑩^1\dot{𝒈})(\dot{𝒈}^T𝑩^1𝒈)^2},$$
(2)
where $`t_{\text{rel}}`$ is the relative phase between $`g(t)`$ and the FADC clock. Thus $`\overline{E}`$ and $`\overline{E\tau }`$ are given by a weighted sum of the discrete measurements $`y_i`$ with the weights for the amplitude, $`w_{\text{amp}}(t_{\text{rel}})`$, and time shift, $`w_{\text{time}}(t_{\text{rel}})`$, plus $`O(\tau ^2)`$. To reduce $`O(\tau ^2)`$ the fit can be iterated using $`g(t_1=t\tau )`$ and the weights $`w_{\mathrm{amp}/\mathrm{time}}(t_{\text{rel}}+\tau )`$ . Figure 2 a) shows the amplitude and timing weights for the MC pulse shape. The first weight $`w_{\mathrm{amp}/\mathrm{time}}(t_0)`$ is plotted as a function of $`t_{\text{rel}}`$ in the range $`[0.5,0.5[T_{\text{ADC}}`$, the second weight in the range $`[0.5,1.5[T_{\text{ADC}}`$ and so on.
The expected contributions of the noise to the error of the estimated amplitude and timing only depend on the the shape $`g(t)`$, and the noise auto-correlation $`B`$. Analytic expressions can be found in references .
## 3 Performance and Discussion
Figure 2b) shows the measured timing resolution for different calibration LED pulses as a function of the mean reconstructed pulse charge. For signals of 10 photo-electrons the timing resolution is as good as 700 ps, for very large signals a timing resolution of about 200 ps can be achieved.
Figure 3 shows the charge and arrival time resolution as a function of the input pulse height for MC simulations (no PMT time spread and no gain fluctuations) assuming an extra-galactic background for different signal extraction algorithms. The digital filter yields the best charge and timing resolution of the studied algorithms .
For known constant signal shapes and noise auto-correlations the digital filter yields the best theoretically achievable signal and timing resolution. Due to the pulse shaping of the Cherenkov signals the algorithm can be applied to reconstruct their charge and arrival time, although there are some fluctuations of the pulse shape and noise behavior. The digital filter reduces the noise contribution to the error of the reconstructed signal. Thus it is possible to lower the image cleaning levels and the analysis energy threshold . The timing resolutions is as good as a few hundred ps for large signals.
## Acknowledgements
The authors thank F. Goebel, Th. Schweizer and W. Wittek for discussions and suggestions.
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# 1 Introduction
## 1 Introduction
The most popular classical noncommutative field theory (see, e.g. review ) can be realized on ordinary smooth field functions $`f(x),g(x)`$ on $`R^4`$ using the following pseudolocal representation of the $``$-product :
$`fg=fe^Pg=fg+\frac{i}{2}\vartheta _{mn}_mf_ng\frac{1}{8}\vartheta _{mn}\vartheta _{rs}_m_rf_n_sg+\mathrm{},`$
$`fPg=\frac{i}{2}\vartheta _{mn}_mf_ng.`$ (1.1)
where $`x_m`$ are the coordinates of $`R^4`$, $`_m=/x_m`$, and $`\vartheta _{mn}`$ are some constants ( $`m,n=1,2,3,4`$). All products of the functions and their derivatives in the right-hand side are commutative. It is evident that nonlinear interactions in these noncommutative (nonlocal) field theories are not invariant with respect to the standard Lorentz transformations of local fields.
The quantum group structures in this noncommutative algebra of functions were found and analyzed in -. The basic point of this interpretation is connected with the twist operator acting on tensor products of functions
$$=\mathrm{exp}(𝒫),𝒫=\frac{i}{2}\vartheta ^{mn}P_mP_n$$
(1.2)
where $`P_mf=_mf`$. The strict definition of the noncommutative product is
$$fg=\mu fg,\mu fg=fg$$
(1.3)
where $`\mu `$ is the multiplication map in the commutative algebra. Thus, this twist operator is the quantum-group analog of the pseudolocal operator $`\mathrm{exp}(P)`$ (1.1).
Let us consider generators of the Poincaré group $`P_m`$ and $`M_{mn}`$. By definition, the twist-deformed Poincaré group $`U_t(P_m,M_{mn})`$ has the undeformed Lie algebra of generators; however, its coproduct is deformed
$`\mathrm{\Delta }_t(P_m)=P_m1+1P_m,`$
$`\mathrm{\Delta }_t(M_{mn})=\mathrm{exp}(𝒫)(M_{mn}1+1M_{mn})\mathrm{exp}(𝒫).`$ (1.4)
The exact constructions of maps between differential operators on commutative and noncommutative algebras of functions were formulated in recent papers of the Munchen group . It was shown that the $``$-product (1.3) transforms covariantly in $`U_t(P_m,M_{mn})`$, but the Leibniz rule for deformed transformations is changed according to Eq.(1.4). 4D-space integrals of the covariant $``$-products of fields are invariant with respect to $`U_t(P_m,M_{mn})`$. The quadratic free interactions possess also the standard Poincaré invariance.
We shall consider the quantum group interpretation of the non-anticommutative deformations in the Euclidean supersymmetric theories - <sup>1</sup><sup>1</sup>1Note that the alternative quantum-group deformations of supersymmetries with a more complex supersymmetric geometry were considered earlier , however, we shall not discuss these models here.. The basic $``$-product of these models is realized on the standard local superfields, and supersymmetry generators can be presented as the 1-st order differential operators on the undeformed superspace. The left-handed Grassmann coordinates of these superspaces do not anticommute with respect to the $``$-product, but the basic chiral bosonic coordinates commute with all superspace coordinates. The non-anticommutative superspace is defined exactly as the $``$-product algebra on ordinary functions of the superspace coordinates. The twist elements for the nilpotent deformations can be constructed in terms of the left supersymmetry generators.
Section 2 is devoted to the analysis of the twist deformation of the $`N=(\frac{1}{2},\frac{1}{2})`$ supersymmetry . We derive the unusual Leibniz rules for the deformed transformations on the products of superfields or the products of component fields. Twist deformation of the Euclidean $`N=(1,1)`$ supersymmetry in the chiral and harmonic superspaces is considered in Sect. 3. The corresponding nilpotent operator $`𝒫`$ is analogous to the basic bi-differential operator of the $`N=(1,1)`$ deformations in Refs..
The twist interpretation allow us to understand correctly transformation properties of $``$-products of superfields by analogy with the deformed transformations of the $``$-products of fields (1.1) in the noncommutative field theory . At the level of the pseudolocal superfield formalism, the $`t`$-supersymmetry is equivalent to the $``$-covariance principle for the noncommutative algebra of superfields which means the similarity of transformations of local superfields and their $``$-products. The covariance principle allow us to obtain a simplified field-theoretical derivation of the unusual Leibniz rules for the deformed supersymmetry transformations on $``$-products of superfields. Known deformed supersymmetric actions in the non-anticommutative superspaces are manifestly invariant with respect to the corresponding twist-deformed supersymmetry, and this invariance explains naturally all selection rules of these theories which could seem formal earlier.
## 2 Twist-deformed $`N=(\frac{1}{2},\frac{1}{2})`$ supersymmetry
We use the chiral coordinates $`z^M=(y_m,\theta ^\alpha ,\overline{\theta }^{\dot{\alpha }})`$ in the Euclidean superspace R$`(4|2,2)`$, where $`m=1,2,3,4,\alpha =1,2,`$ and $`\dot{\alpha }=\dot{1},\dot{2}`$. The central and antichiral 4D coordinates are, respectively,
$$x_m=y_mi\theta \sigma _m\overline{\theta },\overline{y}_m=y_m2i\theta \sigma _m\overline{\theta }$$
(2.1)
and $`(\sigma _m)_{\alpha \dot{\alpha }}`$ are the SO(4) Weyl matrices. Note that these coordinates are pseudoreal with respect to the special conjugation
$$(y_m)^{}=y_m,(\theta ^\alpha )^{}=\epsilon _{\alpha \beta }\theta ^\beta ,(\overline{\theta }^{\dot{\alpha }})^{}=\epsilon _{\dot{\alpha }\dot{\beta }}\overline{\theta }^{\dot{\beta }},$$
(2.2)
so one can use the reality condition for the even Euclidean chiral superfield $`\varphi (y,\theta )`$. The generators of the Euclidean $`N=(\frac{1}{2},\frac{1}{2})`$ supersymmetry SUSY$`(\frac{1}{2},\frac{1}{2})`$ have the following form:
$`L_\alpha ^\beta =L_\alpha ^\beta (y)+L_\alpha ^\beta (\theta )=\frac{1}{4}(\sigma _m\overline{\sigma }_n)_\alpha ^\beta (y_n_my_m_n)+\theta ^\beta _\alpha \frac{1}{2}\delta _\alpha ^\beta \theta ^\gamma _\gamma ,`$
$`R_{\dot{\alpha }}^{\dot{\beta }}=R_{\dot{\alpha }}^{\dot{\beta }}(y)+R_{\dot{\alpha }}^{\dot{\beta }}(\overline{\theta })=\frac{1}{4}(\overline{\sigma }_m\sigma _n)_{\dot{\alpha }}^{\dot{\beta }}(y_m_ny_n_m)+\overline{\theta }^{\dot{\beta }}\overline{}_{\dot{\alpha }}\frac{1}{2}\delta _{\dot{\alpha }}^{\dot{\beta }}\overline{\theta }^{\dot{\gamma }}\overline{}_{\dot{\gamma }},`$
$`O=\theta ^\alpha _\alpha \overline{\theta }^{\dot{\alpha }}\overline{}_{\dot{\alpha }},Q_\alpha =_\alpha ,\overline{Q}_{\dot{\alpha }}=\overline{}_{\dot{\alpha }}2i\theta ^\alpha _{\alpha \dot{\alpha }},P_m=_m,`$ (2.3)
where $`(\overline{\sigma }_m)^{\dot{\alpha }\alpha }=\epsilon ^{\alpha \beta }\epsilon ^{\dot{\alpha }\dot{\beta }}(\sigma _m)_{\beta \dot{\beta }}`$, and $`_M=(_m,_\alpha ,\overline{}_{\dot{\alpha }})`$ are partial derivatives in the chiral coordinates. Generators $`L_\alpha ^\beta ,R_{\dot{\alpha }}^{\dot{\beta }}`$ and $`O`$ correspond to the automorphism group SU(2)$`{}_{L}{}^{}\times `$SU(2)$`{}_{R}{}^{}\times `$O(1,1). The SUSY$`(\frac{1}{2},\frac{1}{2})`$ transformations can be separated as follows
$`\delta A=(g+G)A,g=P_c+R_\rho +Q_ϵ,G=L_\lambda +aO+\overline{Q}_{\overline{ϵ}},`$ (2.4)
$`P_c=c_mP_m,L_\lambda =\lambda _\beta ^\alpha L_\alpha ^\beta ,R_\rho =\rho _{\dot{\beta }}^{\dot{\alpha }}R_{\dot{\alpha }}^{\dot{\beta }},Q_ϵ=ϵ^\alpha Q_\alpha ,\overline{Q}_{\overline{ϵ}}=\overline{ϵ}^{\dot{\alpha }}\overline{Q}_{\dot{\alpha }},`$
where the corresponding combinations of operators and transformation parameters are introduced. These definitions will be convenient in the deformed supersymmetry.
We shall use the notation S$`(4|2,2)`$ or C$`(4|2,0)`$ for the supercommutative algebras of general or chiral superfields. The bilinear multiplication map $`\mu `$ connects the tensor product of superfields with the local supercommutative product in S$`(4|2,2)`$
$$\mu AB=AB=(1)^{p(A)p(B)}BA.$$
(2.5)
The standard coproduct map is defined on the generators of SUSY$`(\frac{1}{2},\frac{1}{2})`$ (2.4)
$$\mathrm{\Delta }(g)=g1+1g,\mathrm{\Delta }(G)=G1+1G.$$
It determines the action of these generators on the tensor product of superfields and yields the standard Leibniz rule for supersymmetry transformations on the local product of superfields $`\delta (AB)=(\delta A)B+A\delta B`$.
The non-anticommutative deformation $`\widehat{z}=(y_m,\widehat{\theta }^\alpha ,\overline{\theta }^{\dot{\alpha }})`$ of the coordinates of the Euclidean $`N=(\frac{1}{2},\frac{1}{2})`$ superspace was considered in . The basic operator relation of the non-anticommutative superspace is
$$T^{\alpha \beta }(\widehat{\theta })=\widehat{\theta }^\alpha \widehat{\theta }^\beta +\widehat{\theta }^\beta \widehat{\theta }^\alpha C^{\alpha \beta }=0,$$
(2.6)
where $`C^{\alpha \beta }`$ are some constants. The operator superfields $`\widehat{A}(y,\widehat{\theta },\overline{\theta })`$ and $`\widehat{B}(y,\widehat{\theta },\overline{\theta })`$ with the antisymmetric ordering of the $`\widehat{\theta }^\alpha `$ decomposition contain the highest terms $`\epsilon _{\alpha \beta }\widehat{\theta }^\alpha \widehat{\theta }^\beta `$. In the pseudolocal representation, we consider the usual superfields $`A(z)`$ and $`B(z)`$ as the supercommutative images of these operator superfields. The corresponding $``$-product of superfields $`A(z)`$ and $`B(z)`$ is defined via the generators of the left $`N=(\frac{1}{2},0)`$ supersymmetry
$`AB=Ae^PB=AB\frac{1}{2}(1)^{p(A)}C^{\alpha \beta }Q_\alpha AQ_\beta B\frac{1}{32}C^{\alpha \beta }C_{\alpha \beta }Q^2AQ^2B,`$
$`APB=\frac{1}{2}(1)^{p(A)}C^{\alpha \beta }Q_\alpha AQ_\beta B,P^3=0,`$ (2.7)
where $`p(A)`$ is the $`Z_2`$ grading, and $`P`$ is the nilpotent bi-differential operator. The deformed algebras S$`{}_{}{}^{}(4|2,2)`$ and C$`{}_{}{}^{}(4|2,0)`$ use this noncommutative product for general or chiral superfields, respectively.
The twist operator in this supersymmetry was introduced in (see also discussions in )
$`=\mathrm{exp}(𝒫),𝒫=\frac{1}{2}C^{\alpha \beta }Q_\alpha Q_\beta ,`$
$`(AB)=AB\frac{1}{2}(1)^{p(A)}C^{\alpha \beta }Q_\alpha AQ_\beta B\frac{1}{32}C^{\alpha \beta }C_{\alpha \beta }Q^2AQ^2B.`$ (2.8)
The bilinear map $`\mu _{}`$ in S$`{}_{}{}^{}(4|2,2)`$ can be defined via this twist operator
$`AB\mu _{}AB=\mu \mathrm{exp}(𝒫)AB,`$ (2.9)
so $``$ is the quantum group analog of the pseudolocal operator $`e^P`$ (2.7).
By analogy with the map between differential operators on commutative and noncommutative algebras of functions , one can easily define the corresponding differential operator $`\widehat{X}_D`$ on S$`{}_{}{}^{}(4|2,2)`$ for any differential operator $`D`$ on the supercommutative algebra. In the case of the 1-st order operator $`D_1=\xi ^M(z)_M`$, the image $`\widehat{X}_{D_1}`$ contains, in general, terms with higher derivatives on S$`{}_{}{}^{}(4|2,2)`$
$`(\widehat{}_MA)=_MA,(\widehat{}_Mz^N)=\delta _M^N,`$
$`(\widehat{X}_{D_1}A)=(D_1A)=\mu _{}\mathrm{exp}(𝒫)(\xi ^M(z)_M)(1A)`$
$`=\xi ^M(z)_MA+\frac{1}{2}(1)^{p(D_1)}C^{\alpha \beta }Q_\alpha \xi ^M(z)_MQ_\beta A`$
$`\frac{1}{32}C^{\alpha \beta }C_{\alpha \beta }Q^2\xi ^M(z)_MQ^2A`$ (2.10)
where $`p(D_1)`$ is the $`Z_2`$ grading of $`D_1`$. For instance, the deformed images of generators $`\overline{Q}_{\dot{\alpha }}`$ and $`L_a^\beta `$ (2.3) are the following second-order differential operators on S$`{}_{}{}^{}(4|2,2)`$:
$`(\widehat{\overline{Q}}_{\dot{\alpha }}A)=(\overline{}_{\dot{\alpha }}2i\theta ^\alpha _{\alpha \dot{\alpha }}+iC^{\alpha \beta }_{\alpha \dot{\alpha }}Q_\beta )A=\overline{Q}_{\dot{\alpha }}A,`$
$`(\widehat{L}_\alpha ^\beta A)=(L_\alpha ^\beta \frac{1}{2}C^{\beta \gamma }Q_\gamma Q_\alpha )A=L_\alpha ^\beta A,`$ (2.11)
while $`\widehat{g}=g`$ and $`\widehat{O}=O`$. The deformed operators can be used in the operator representation of S$`{}_{}{}^{}(4|2,2)`$ to check directly the covariance of the basic operator $`T^{\alpha \beta }(\widehat{\theta })`$ (2.6) with respect to the transformations of the deformed supersymmetry
$`\widehat{\overline{Q}}_{\dot{\alpha }}T^{\alpha \beta }=0,\widehat{L}_\sigma ^\rho T^{\alpha \beta }=\delta _\sigma ^\alpha T^{\rho \beta }+\delta _\sigma ^\beta T^{\alpha \rho }\delta _\sigma ^\rho T^{\alpha \beta }.`$ (2.12)
The Lie superalgebra of the deformed generators with hats is isomorphic to the Lie superalgebra of the undeformed supersymmetry generators (2.3).
The coproduct $`\mathrm{\Delta }_t(G)=e^𝒫\mathrm{\Delta }(G)e^𝒫`$ in SUSY$`{}_{t}{}^{}(\frac{1}{2},\frac{1}{2})`$ is deformed on $`G=\overline{Q}_{\overline{ϵ}}+L_\lambda +aO`$, in particular,
$`\mathrm{\Delta }_t(\overline{Q}_{\overline{ϵ}})=\overline{Q}_{\overline{ϵ}}1+1\overline{Q}_{\overline{ϵ}}+i\overline{ϵ}^{\dot{\alpha }}C^{\alpha \beta }(_{\alpha \dot{\alpha }}Q_\beta Q_\alpha _{\beta \dot{\alpha }}),`$
$`\mathrm{\Delta }_t(L_\lambda )=L_\lambda 1+1L_\lambda +\frac{1}{2}C^{\rho \sigma }(\lambda _\rho ^\alpha Q_\alpha Q_\sigma +\lambda _\sigma ^\alpha Q_\rho Q_\alpha ),`$
$`\mathrm{\Delta }_t(O)=O1+1OC^{\alpha \beta }Q_\alpha Q_\beta ,`$ (2.13)
while $`e^𝒫\mathrm{\Delta }(g)e^𝒫=\mathrm{\Delta }(g)=g1+1g`$.
Acting by the composition of $`\mu _{}`$ and coproduct $`\mathrm{\Delta }_t(\overline{Q}_{\overline{ϵ}})`$ on the tensor product of superfields, one can obtain the following relation:
$`\widehat{\delta }_{\overline{ϵ}}(AB)\mu _{}\mathrm{\Delta }_t(\overline{Q}_{\overline{ϵ}})AB=(\overline{Q}_{\overline{ϵ}}A)BA\overline{Q}_{\overline{ϵ}}B`$
$`i\overline{ϵ}^{\dot{\alpha }}C^{\alpha \beta }[(1)^{p(A)}_{\alpha \dot{\alpha }}AQ_\beta BQ_\alpha A_{\beta \dot{\alpha }}B]=\overline{Q}_{\overline{ϵ}}(AB).`$ (2.14)
The last formula can be derived from the pseudolocal definition of $`AB`$ (2.7)
$$\widehat{\delta }_{\overline{ϵ}}(AB)=[\overline{Q}_{\overline{ϵ}},(Ae^PB)]=(\overline{Q}_{\overline{ϵ}}A)e^PB+A[\overline{Q}_{\overline{ϵ}},e^P]B+Ae^P(\overline{Q}_{\overline{ϵ}}B).$$
(2.15)
Note that Eq.(2.14) can be treated as the deformed Leibniz rule for $`\widehat{\delta }_{\overline{ϵ}}`$.
It is important to formulate the principle of covariance of the algebra S$`{}_{}{}^{}(4|2,2)`$ (or C$`{}_{}{}^{}(4|2,0)`$) with respect to all deformed transformations of SUSY$`{}_{t}{}^{}(\frac{1}{2},\frac{1}{2})`$: the primary superfields $`A,B`$ and their $``$-product $`(AB)`$ transform analogously
$`\widehat{\delta }_GA=GA=\widehat{G}A,`$
$`\widehat{\delta }_G(AB)=G(AB)=\widehat{G}(AB).`$ (2.16)
where $`\widehat{\delta }_G=(\widehat{\delta }_{\overline{ϵ}}+\widehat{\delta }_\lambda +\widehat{\delta }_a)`$. In the pseudolocal formalism of the noncommutative superfield theory, these relations are sufficient to derive the deformed Leibniz rules.
It is evident that the algebra S$`{}_{}{}^{}(4|2,2)`$ is not covariant with respect to the undeformed supersymmetry, e.g., $`\delta _{\overline{ϵ}}(AB)=(\delta _{\overline{ϵ}}AB)+(A\delta _{\overline{ϵ}}B)\overline{ϵ}^{\dot{\alpha }}\overline{Q}_{\dot{\alpha }}(AB).`$
The deformed transformations act noncovariantly on the supercommutative product of superfields $`AB`$. For instance, it is not difficult to consider the following transformations of the ordinary product of the even chiral superfields:
$`\widehat{\delta }_{\overline{ϵ}}(\varphi _1\varphi _2)\mu \mathrm{\Delta }_t(\overline{Q}_{\overline{ϵ}})\varphi _1\varphi _2=\overline{Q}_{\overline{ϵ}}(\varphi _1\varphi _2)i\overline{ϵ}^{\dot{\alpha }}C^{\alpha \beta }(_{\alpha \dot{\alpha }}\varphi _1Q_\beta \varphi _2`$
$`Q_\alpha \varphi _1_{\beta \dot{\alpha }}\varphi _2)=\widehat{\delta }_{\overline{ϵ}}(a_1a_2)+O(\theta ),`$
$`\widehat{\delta }_\lambda (\varphi _1\varphi _2)\mu \mathrm{\Delta }_t(L_\lambda )\varphi _1\varphi _2=\lambda _\beta ^\alpha L_\alpha ^\beta (\varphi _1\varphi _2)`$
$`\frac{1}{2}C^{\rho \sigma }(\lambda _\rho ^\alpha Q_\alpha \varphi _1Q_\sigma \varphi _2+\lambda _\sigma ^\alpha Q_\rho \varphi _1Q_\alpha \varphi _2)=\widehat{\delta }_\lambda (a_1a_2)+O(\theta ).`$ (2.17)
The first terms coincide with the transformations of the undeformed supersymmetry. Using the $`\theta `$-decomposition of these superfield formulas one can obtain the deformed transformations of the products of component fields, for instance,
$`\widehat{\delta }_{\overline{ϵ}}(a_1a_2)=i\overline{ϵ}^{\dot{\alpha }}C^{\alpha \beta }(_{\alpha \dot{\alpha }}a_1\psi _{\beta 2}\psi _{\alpha 1}_{\beta \dot{\alpha }}a_2),`$
$`\widehat{\delta }_\lambda (a_1a_2)=\lambda _\beta ^\alpha L_\alpha ^\beta (y)(a_1a_2)\frac{1}{2}C^{\rho \sigma }(\lambda _\rho ^\alpha \psi _{\alpha 1}\psi _{\sigma 2}+\lambda _\sigma ^\alpha \psi _{\rho 1}\psi _{\alpha 2}).`$ (2.18)
Let us consider two even chiral superfields $`\varphi _1`$ and $`\varphi _2`$ in the chiral basis
$`\varphi _i=a_i+\theta ^\alpha \psi _{i\alpha }+\theta ^2f_i,`$ (2.19)
$`Q_\alpha \varphi _i=\psi _{i\alpha }+2\theta _\alpha f_i,Q^2\varphi _i=4f_i.`$
The $`\theta `$-decomposition of the $``$-product of two chiral superfields depends on these components and constants $`C^{\alpha \beta }`$
$`\mathrm{\Phi }_{12}=\varphi _1\varphi _2=+\theta ^\alpha \mathrm{\Psi }_\alpha +\theta ^2F,`$ (2.20)
$`=a_1a_2\frac{1}{2}C^{\alpha \beta }\psi _{1\alpha }\psi _{2\beta }\frac{1}{2}C^{\alpha \beta }C_{\alpha \beta }f_1f_2,`$
$`\mathrm{\Psi }_\alpha =a_1\psi _{2\alpha }+a_2\psi _{1\alpha }C_{\alpha \beta }(f_1\psi _2^\beta f_2\psi _1^\beta ),`$
$`F=a_1f_2+a_2f_1\frac{1}{2}\psi _1^\alpha \psi _{2\alpha }.`$ (2.21)
These relations can be treated as a deformed tensor calculus for chiral component multiplets. The $`t`$-supersymmetry transformations of the composite components (2.20) are completely analogous to the transformations of the basic components $`a_i,\psi _{\alpha i},f_i`$
$`\widehat{\delta }_{\overline{ϵ}}=0,\widehat{\delta }_{\overline{ϵ}}\mathrm{\Psi }_\alpha =2i\overline{ϵ}^{\dot{\alpha }}_{\alpha \dot{\alpha }},\widehat{\delta }_{\overline{ϵ}}F=i\overline{ϵ}^{\dot{\alpha }}_{\alpha \dot{\alpha }}\mathrm{\Psi }^\alpha ,`$
$`\widehat{\delta }_\lambda =\lambda _\beta ^\alpha L_\alpha ^\beta (y),\widehat{\delta }_\lambda \mathrm{\Psi }_\gamma =\lambda _\gamma ^\alpha \mathrm{\Psi }_\alpha \lambda _\beta ^\alpha L_\alpha ^\beta (y)\mathrm{\Psi }_\gamma ,\widehat{\delta }_\lambda F=\lambda _\beta ^\alpha L_\alpha ^\beta (y)F.`$ (2.22)
These transformations are compatible with the noncovariant component transformations (2.18).
The non-anticommutative deformation of the Euclidean model for an arbitrary number of the chiral and antichiral superfields $`\varphi _a`$ and $`\overline{\varphi }_a`$ is based on the superfield action $`S_{}(\varphi _a,\overline{\varphi }_a)`$ . Each term of the $``$-polynomial decomposition of this action is separately invariant with respect to SUSY$`{}_{t}{}^{}(\frac{1}{2},\frac{1}{2})`$, while the quadratic terms like $`d^8z\varphi _a\overline{\varphi }_a=d^8z\varphi _a\overline{\varphi }_a`$ possess also the undeformed supersymmetry.
## 3 Twist-deformed $`N=(1,1)`$ supersymmetry
The nilpotent deformations of the Euclidean $`N=(1,1)`$ supersymmetry were considered in the framework of the harmonic-superspace approach . Harmonic-superspace coordinates contain the SU(2)/U(1) harmonics $`u_i^\pm `$ and the chiral superspace coordinates
$$z^M=(y_m,\theta _k^\alpha ,\overline{\theta }^{\dot{\alpha }k}),y_m=x_m+i\theta _k\sigma _m\overline{\theta }^k$$
(3.1)
where $`x_m`$ are the central 4D coordinates. The spinor derivatives $`D_\alpha ^k`$ and $`\overline{D}_{\dot{\alpha }k}`$ in these coordinates are
$$D_\alpha ^k=_\alpha ^k+2i\overline{\theta }^{\dot{\alpha }k}_{\alpha \dot{\alpha }},\overline{D}_{\dot{\alpha }k}=\overline{}_{\dot{\alpha }k}.$$
(3.2)
Using harmonic projections of the Grassmann coordinates $`\theta ^{\pm \alpha }=u_k^\pm \theta ^{\alpha k}`$ and $`\overline{\theta }^{\pm \dot{\alpha }}=u_k^\pm \overline{\theta }^{\dot{\alpha }k}`$ one can define the analytic coordinates $`(x_A,\theta ^\pm ,\overline{\theta }^\pm )`$. The corresponding representation of spinor derivatives and supersymmetry generators can be found in . The bosonic part of the harmonic superspaces is $`R^4\times S^2`$, and the left-right Grassmann dimensions of general, chiral or analytic superspaces are (4,4), (4,0) or (2,2), respectively. We shall use the notation S$`(4,2|4,4)`$ for the supercommutative algebra and S$`{}_{}{}^{}(4,2|4,4)`$ for the non-anticommutative algebra of general harmonic superfields, respectively.
It is convenient to use the following differential representation of the SUSY(1,1) generators on S$`(4,2|4,4)`$:
$`T_l^k=\theta _l^\alpha _\alpha ^k+\frac{1}{2}\delta _l^k\theta _j^\alpha _\alpha ^j+\overline{\theta }^{\dot{\alpha }k}\overline{}_{\dot{\alpha }l}\frac{1}{2}\delta _l^k\overline{\theta }^{\dot{\alpha }j}\overline{}_{\dot{\alpha }j}u_l^\pm ^k+\frac{1}{2}\delta _l^ku_j^\pm ^j,`$
$`L_\alpha ^\beta =L_\alpha ^\beta (y)+\theta _k^\beta _\alpha ^k\frac{1}{2}\delta _\alpha ^\beta \theta _k^\gamma _\gamma ^k,R_{\dot{\alpha }}^{\dot{\beta }}=R_{\dot{\alpha }}^{\dot{\beta }}(y)+\overline{\theta }^{\dot{\beta }k}\overline{}_{\dot{\alpha }k}\frac{1}{2}\delta _{\dot{\alpha }}^{\dot{\beta }}\overline{\theta }^{\dot{\gamma }k}\overline{}_{\dot{\gamma }k},`$
$`O=\theta _k^\alpha _\alpha ^k\overline{\theta }^{\dot{\alpha }k}\overline{}_{\dot{\alpha }k},Q_\alpha ^k=_\alpha ^k,\overline{Q}_{\dot{\alpha }k}=\overline{}_{\dot{\alpha }k}2i\theta _k^\alpha _{\alpha \dot{\alpha }},P_m=_m`$ (3.3)
where $`L_\alpha ^\beta (y)`$ and $`R_{\dot{\alpha }}^{\dot{\beta }}(y)`$ are defined above (2.3), and partial harmonic derivatives act as follows $`^lu_k^\pm =\delta _k^l`$. For our purposes, it is convenient to consider the following combinations of SUSY(1,1) generators and corresponding parameters:
$`g=P_c+R_\rho +Q_ϵ,G=T_u+L_\lambda +\overline{Q}_{\overline{ϵ}}+aO,`$ (3.4)
$`P_c=c_mP_m,T_u=u_l^kT_k^l,L_\lambda =\lambda _\beta ^\alpha L_\alpha ^\beta ,R_\rho =\rho _{\dot{\beta }}^{\dot{\alpha }}R_{\dot{\alpha }}^{\dot{\beta }},Q_ϵ=ϵ_k^\alpha Q_\alpha ^k,\overline{Q}_{\overline{ϵ}}=\overline{ϵ}^{\dot{\alpha }k}\overline{Q}_{\dot{\alpha }k}.`$
The $`N=(1,1)`$ twist operator $`=\mathrm{exp}(𝒫)`$ contains the nilpotent operator
$$𝒫=\frac{1}{2}C_{kl}^{\alpha \beta }Q_\alpha ^kQ_\beta ^l,𝒫^5=0$$
(3.5)
where $`C_{kl}^{\alpha \beta }`$ are some constants. The non-anticommutative product in the corresponding deformed algebra S$`{}_{}{}^{}(4,2|4,4)`$ can be defined by equivalent formulas
$`AB=A\mathrm{exp}(P)B=\mu \mathrm{exp}(𝒫)AB=\mu _{}AB`$ (3.6)
where $`\mu `$ and $`\mu _{}`$ are product maps for S$`(4,2|4,4)`$ and S$`{}_{}{}^{}(4,2|4,4)`$ and $`P`$ is the basic operator from
$$APB=\frac{1}{2}(1)^{p(A)}C_{kl}^{\alpha \beta }Q_\alpha ^kAQ_\beta ^lB=\mu 𝒫AB.$$
(3.7)
The non-anticommutative algebras of the $`N=(1,1)`$ chiral or G-analytic superfields can be defined analogously.
By analogy with Eq.(2.10), one can construct the image operator $`\widehat{X}_{D_1}`$ on S$`{}_{}{}^{}(4,2|4,4)`$ for any 1-st order differential operator $`D_1=\xi ^M(z)_M`$ on S$`(4,2|4,4)`$. The twisted coproduct of SUSY<sub>t</sub>(1,1) $`\mathrm{\Delta }_t(G)=e^𝒫\mathrm{\Delta }(G)e^𝒫`$ is deformed on $`G=(\overline{Q}_{\overline{ϵ}}+T_u+L_\lambda +aO)`$, in particular,
$`\mathrm{\Delta }_t(\overline{Q}_{\overline{ϵ}})=\overline{Q}_{\overline{ϵ}}1+1\overline{Q}_{\overline{ϵ}}+i\overline{ϵ}^{\dot{\alpha }k}C_{kj}^{\alpha \beta }_{\alpha \dot{\alpha }}Q_\beta ^ji\overline{ϵ}^{\dot{\alpha }k}C_{ik}^{\alpha \beta }Q_\alpha ^i_{\beta \dot{\alpha }},`$
$`\mathrm{\Delta }_t(T_u)=T_u1+1T_u\frac{1}{2}u_k^lC_{lj}^{\alpha \beta }Q_\alpha ^kQ_\beta ^j\frac{1}{2}u_k^lC_{jl}^{\alpha \beta }Q_\alpha ^jQ_\beta ^k,`$
$`\mathrm{\Delta }_t(L_\lambda )=L_\lambda 1+1L_\lambda +\frac{1}{2}\lambda _\beta ^\alpha C_{kl}^{\beta \gamma }Q_\alpha ^kQ_\gamma ^l+\frac{1}{2}\lambda _\beta ^\alpha C_{kl}^{\rho \beta }Q_\rho ^kQ_\alpha ^l,`$
$`\mathrm{\Delta }_t(O)=O1+1OC_{kl}^{\alpha \beta }Q_\alpha ^kQ_\beta ^l,`$ (3.8)
while $`e^𝒫\mathrm{\Delta }(g)e^𝒫=\mathrm{\Delta }(g)`$ for $`g=Q_ϵ+P_c+R_\rho `$.
The $``$-products of arbitrary $`N=(1,1)`$ superfields preserve covariance with respect to all deformed transformations of SUSY<sub>t</sub>(1,1)
$`\widehat{\delta }_G(AB)=\mu _{}\mathrm{\Delta }_t(G)AB=G(AB).`$ (3.9)
The deformed Leibniz rules can be derived directly from these covariant relations.
The twisted supersymmetry acts noncovariantly on the supercommutative product of superfields, for instance,
$`\widehat{\delta }_{\overline{ϵ}}(AB)\mu \mathrm{\Delta }_t(\overline{Q}_{\overline{ϵ}})AB`$
$`=\overline{Q}_{\overline{ϵ}}(AB)i\overline{ϵ}^{\dot{\alpha }k}C_{kj}^{\alpha \beta }(1)^{p(A)}_{\alpha \dot{\alpha }}AQ_\beta ^jB+i\overline{ϵ}^{\dot{\alpha }k}C_{ik}^{\alpha \beta }Q_\alpha ^iA_{\beta \dot{\alpha }}B.`$ (3.10)
It is easy to define the $`t`$-deformed transformations on the products of the $`N=(1,1)`$ component fields using the corresponding Grassmann decompositions.
In the special case of the singlet deformation , the twist operator is defined by the parameter $`I`$ and the SU(2)$`\times `$SU(2)<sub>L</sub> invariant constant tensor
$$C_{kl}^{\alpha \beta }=2I\epsilon ^{\alpha \beta }\epsilon _{kl}𝒫_s=IQ_\alpha ^iQ_i^\alpha .$$
(3.11)
The $`𝒫_s`$-twist deformation vanishes for SU(2) and SU(2)<sub>L</sub> transformations.
The Leibniz rules for differential operators $`D=(_m,D_\alpha ^k,\overline{D}_{\dot{\alpha }k},/u_k^\pm )`$ are standard for the general $`Q`$-deformation (3.5)
$$D(AB)=(DAB)+(1)^{p(D)p(A)}(ADB).$$
(3.12)
The $``$-product preserves differential constraints of chirality, antichirality and Grassmann analyticity . All superfield actions using $``$-products in the non-anticommutative $`N=(1,1)`$ harmonic superspace are invariant with respect to the quantum group SUSY<sub>t</sub>(1,1), and this invariance is a natural basic principle of these deformed theories. Free quadratic parts of these actions possess also the undeformed $`N=(1,1)`$ supersymmetry. The simple examples of the superfield-density terms for the analytic hypermultiplet $`q^+,\stackrel{~}{q}^+`$ and the U(1) gauge potential $`V^{++}`$ in the deformed theory are
$$\stackrel{~}{q}^+(D^{++}q^++[V^{++},q^+]_{})+\lambda q^+q^+\stackrel{~}{q}^+\stackrel{~}{q}^+.$$
(3.13)
The $`t`$-supersymmetry transformations of any term $`L_{}^{+4}`$ in this density are
$`\widehat{\delta }_GL_{}^{+4}=(\overline{ϵ}^{\dot{\alpha }k}\overline{Q}_{\dot{\alpha }k}+u_k^lT_l^k+l_\beta ^\alpha L_a^\beta +aO)L_{}^{+4},`$ (3.14)
so the analytic-superspace integrals of these variations vanish. We hope that the manifest SUSY<sub>t</sub>(1,1) covariance could help to prove the nonrenormalization theorems in $`t`$-deformed harmonic-superfield theories by analogy with the corresponding undeformed theories.
## 4 Conclusions
We analyzed the twist deformations of the Euclidean $`N=(\frac{1}{2},\frac{1}{2})`$ and $`N=(1,1)`$ supersymmetries. By analogy with the formalism of the deformed Minkowski space , we construct explicitly the map between differential operators on ordinary and deformed superspaces (2.10). This map connects the standard representation for the supersymmetry generators with the corresponding operator representation on the deformed superspace. It is shown that the noncommutative $``$-products of primary superfields transform covariantly in these $`t`$-deformed supersymmetries. This covariance is a basic principle of the superfield formalism of the deformed theories. The Grassmann-coordinate decompositions of the $``$-product superfields define the deformed tensor calculus for the components of primary superfields. The ordinary supercommutative products of primary superfields or component fields are not covariant with respect to the deformed supersymmetries.
Any polynomial terms of the superfield actions in the non-anticommutative $`N=(\frac{1}{2},\frac{1}{2})`$ and $`N=(1,1)`$ superspaces are manifestly invariant with respect to the corresponding $`t`$-deformed supersymmetries. The bilinear free parts of these actions are also invariant under the standard supersymmetry transformations. The deformation constants of the non-anticommutative superfield theories break some undeformed (super)symmetries, however, these parameters can be treated as ‘coupling constants’ compatible with the deformed supersymmetries. We hope that $`t`$-supersymmetries would help to analyze nonrenormalization theorems using the superfield effective actions in these theories.
I am grateful to P.P. Kulish for stimulating discussions of quantum-group symmetries in the noncommutative field theory. This work was partially supported by DFG grant 436 RUS 113/669-2 , by RFBR grants 03-02-17440 and 04-02-04002, by NATO grant PST.GLG.980302 and by grants of the Heisenberg-Landau and Votruba-Blokhintsev programs.
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# Collapse and Fragmentation of Rotating Magnetized Clouds. I. Magnetic Flux - Spin Relation
## 1 Introduction
It has long been recognized that magnetic field and rotation affect collapse of a molecular cloud, and accordingly, star formation. The magnetic and centrifugal forces, as well as the pressure force, oppose the self-gravity of the cloud and delay star formation. Magnetic field and rotation are coupled. Magnetic field is twisted and amplified by rotation. The twisted magnetic field brakes cloud rotation and launches outflows.
In spite of its importance, only a limited number of numerical simulations have been performed for the coupling of magnetic field and rotation in a collapsing molecular cloud. The first numerical simulation of self-gravitating rotating magnetized clouds were performed by Dorfi (1982). He found formation of bar-like structure for a cloud rotating perpendicular to the magnetic field and that of ring-like structure for an aligned rotator. However, the grid resolution was limited so that the simulation was stopped when the density increased by 200 times from the initial value. The spatial resolution was limited also in other simulations in 1980’s by Phillips & Monaghan (1985) and Dorfi (1989), who studied the cloud with toroidal magnetic field and that with oblique magnetic field, respectively. The spatial resolution was improved greatly by Tomisaka (1998, 2002). He considered an initially filamentary cloud of $`n_{\mathrm{max}}=\mathrm{\hspace{0.17em}10}^4\mathrm{cm}^3`$ and followed the evolution up to the emergence of magnetically driven outflows from the first core of $`n_{\mathrm{max}}>\mathrm{\hspace{0.17em}10}^{11}\mathrm{cm}^3`$, where $`n_{\mathrm{max}}`$ denotes the maximum density. However, his computation was two dimensional and could not take account of asymmetry around the rotation axis. Basu & Mouschovias (1994, 1995a, 1995b) and Nakamura & Li (2002, 2003); Li & Nakamura (2002) have got rid of the symmetry around the axis but introduced the thin disk approximation. The magnetic braking could not be taken into account in these simulations because of the thin disk approximation. Although Boss (2002) has performed three-dimensional simulations, he has employed approximate magnetohydrodynamical equations. The approximation neglects torsion of the magnetic field and replaces magnetic tension with the dilution of the gravity. A fully three-dimensional numerical simulation has just been initiated by Machida, Tomisaka & Matsumoto (2004a), Hosking & Whitworth (2004), and Matsumoto &Tomisaka (2004).
The recent fully three-dimensional simulations have demonstrated that fragmentation of the cloud depends on the magnetic field strength. When the magnetic field is weak, a rotating cloud fragments after the central density exceeds the critical density, $`5\times 10^{10}\mathrm{cm}^3`$, i.e., after the formation of Larson’s first core (Larson, 1969). The magnetic field changes its direction and strength during the collapse of the cloud. Thus it is important to study how strong a magnetic field the first core has.
In this and subsequent papers, we show 144 models in which a filamentary cloud collapses to form a magnetized rotating first core. All the models are constructed using the fully three-dimensional numerical simulation code used in Machida, Tomisaka & Matsumoto (2004a, hereafter MTM04). This paper shows the evolution by the first core formation stage, i.e., the stages before the maximum density reaches the critical density, $`5\times 10^{10}\mathrm{cm}^3`$. The later stages, i.e., fragmentation of the first core and emergence of outflows, are shown in the subsequent paper (Machida, Matsumoto, Hanawa & Tomisaka, 2004c, hereafter PaperII).
From analysis of 144 models, we find two variables which characterize the evolution of magnetic field and rotation. The first one is the ratio of the angular velocity to the magnetic field. This remains nearly constant while the maximum density increases from $`5\times 10^2\mathrm{cm}^3`$ to $`5\times 10^{10}\mathrm{cm}^3`$. The second characteristic variable is the sum of the ratio of the magnetic pressure to the gas pressure and the square of the angular velocity in units of the freefall timescale. This variable converges to a certain value. We refer to the convergence as the magnetic flux - spin ($`B\mathrm{\Omega }`$) relation in the following. We discuss the evolution of the magnetic flux density and angular velocity by means of these two characteristic variables.
This paper is organized as follows: Section 2 denotes the framework of our models and the assumptions employed. Section 3 describes methods of numerical simulations. Section 4 presents typical models in the first four subsections and compares various models in the last subsection. Section 5 discusses implications of the magnetic flux - spin relation and some applications of our models to observations.
## 2 Model
We consider formation of protostars through fragmentation of a filamentary molecular cloud by taking account of its magnetic field and self-gravity. The magnetic field is assumed to be coupled with the gas for simplicity although the molecular gas is only partially ionized. Then the dynamics of the cloud are described by the ideal magnetohydrodynamical (MHD) equations,
$`{\displaystyle \frac{\rho }{t}}+(\rho 𝒗)=0,`$ (1)
$`\rho {\displaystyle \frac{𝒗}{t}}+\rho (\text{v})𝒗=P{\displaystyle \frac{1}{4\pi }}𝑩\times (\times 𝑩)\rho \varphi ,`$ (2)
$`{\displaystyle \frac{𝑩}{t}}=\times (𝒗\times 𝑩),`$ (3)
$`^2\varphi =4\pi \mathrm{G}\rho ,`$ (4)
where $`\rho `$, $`𝒗`$, $`P`$, $`𝑩`$, and $`\varphi `$ denote the density, velocity, pressure, magnetic flux density and gravitational potential, respectively. The ideal MHD approximation is fairly good as long as the gas density is lower than $`10^{11}`$ cm<sup>-3</sup> (Nakano, 1988; Nakano, Nishi & Umebayashi, 2002). The gas pressure is assumed to be
$$P=c_s^2\rho \left[1+\left(\frac{n}{n_{\mathrm{cri}}}\right)^{2/5}\right],$$
(5)
where $`n`$ denotes the number density and is related to the mass density $`\rho `$ by
$$\rho =2.3\times 1.67\times 10^{24}\times n.$$
(6)
The critical number density is set to be $`n_{\mathrm{cri}}=5\times 10^{10}\mathrm{cm}^3`$ (Masunaga & Inutsuka, 2000) and the sound speed is assumed to be $`c_s=0.19`$ km s<sup>-1</sup>. Thus, this equation of state means that the gas is isothermal at $`T=10`$K for $`nn_{\mathrm{cri}}`$ and adiabatic for $`nn_{\mathrm{cri}}`$.
Our initial model is the same as that of Tomisaka (2002) except for the azimuthal perturbation. It is expressed as
$`\rho `$ $`=`$ $`\rho _{\mathrm{c},0}\left[1+(r^2/8H^2)\right]^2\left[1+\delta \rho _z(z)\right]\left[1+\delta \rho _\phi (r,\phi )\right],`$ (7)
$`𝒗`$ $`=`$ $`r\mathrm{\Omega }_{\mathrm{c},0}\left[1+(r^2/8H^2)\right]^{1/2}𝒆_\phi ,`$ (8)
$`𝑩`$ $`=`$ $`B_{\mathrm{c},0}\left[1+(r^2/8H^2)\right]^1\left[1+\delta B_z(r,\phi )\right]𝒆_z,`$ (9)
where
$`H^2={\displaystyle \frac{c_s^2+B_c^2/8\pi \rho _c}{4\pi G\rho _{\mathrm{c},0}2\mathrm{\Omega }_c^2}}.`$ (10)
in the cylindrical coordinates, $`(r,\phi ,z)`$. This initial model denotes a magnetohydrodynamical equilibrium (Stodółkiewicz, 1963) when $`\delta \rho _z(z)`$, $`\delta \rho _\phi (r,\phi )`$ and $`\delta B_z(r,\phi )`$ are not taken into account. The initial density is $`n_{\mathrm{c},0}=5\times 10^2`$ cm<sup>-3</sup> on the axis ($`r=\mathrm{\hspace{0.17em}0}`$). The filamentary cloud is supported in part by the magnetic field and rotation. This equilibrium is unstable against fragmentation in the $`z`$-direction. The perturbation in the $`z`$-direction is assumed to be
$`\delta \rho _z`$ $`=`$ $`A_z\mathrm{cos}\left(2\pi z/\lambda _{\mathrm{max}}\right),`$ (11)
where
$`\lambda _{\mathrm{max}}\left[{\displaystyle \frac{c_s}{(4\pi G\rho _{\mathrm{c},0})^{1/2}}}\right]{\displaystyle \frac{2\pi (1+\alpha /2+\beta )^{1/2}}{0.72\left[(1+\alpha /2+\beta )^{1/3}0.6\right]}},`$ (12)
and
$$\beta 2\omega _c^2H^2/c_s^2.$$
(13)
The symbol, $`\lambda _{\mathrm{max}}`$, denotes the wavelength of the fastest growing perturbation (Matsumoto et al., 1994).
The azimuthal perturbation is assumed to be
$`\delta \rho _\phi ,\delta B_\phi `$ $`=`$ $`\{\begin{array}{cc}A_\phi (r/H)^m\mathrm{cos}(m\phi ),\mathrm{for}r_cH,\hfill & \\ A_\phi \mathrm{cos}(m\phi ),\mathrm{for}r_c>H,\hfill & \end{array}`$ (16)
where the azimuthal wavenumber is assumed to be $`m`$ = 2. The radial dependence is chosen so that the density perturbation remains regular at the origin ($`r=\mathrm{\hspace{0.17em}0}`$) at one time step after the initial stage. The ratio of density to the magnetic flux density is constant in the $`\phi `$-direction for a given $`r`$ and $`z`$ \[see equations (7) and (9) \].
The initial model is characterized by four nondimensional parameters: twice the magnetic-to-thermal pressure ratio,
$$\alpha =B_{\mathrm{zc},0}^2/(4\pi \rho _{\mathrm{c},0}c_{\mathrm{s},0}^2),$$
(17)
the angular velocity normalized by the free-fall timescale,
$$\omega =\mathrm{\Omega }_{\mathrm{c},0}/\sqrt{4\pi \mathrm{G}\rho _{\mathrm{c},0}},$$
(18)
the amplitude of the perturbation in the $`z`$-direction, $`A_z`$, and that of the non-axisymmetric perturbation, $`A_\phi `$. The former two specify the equilibrium model, while the later two do the perturbations. We made 144 models by combining values listed in Table 1. The results depend little on the values of $`A_z`$, thus $`A_z`$ is fixed to be 0.1 in most models.
## 3 Numerical Method
We employed the same 3D MHD nested grid code as that used in MTM04. It incorporates the 3D nested grid code of Matsumoto & Hanawa (2003b) for a hydrodynamical simulation and the approximate Riemann solver for the MHD equation (Fukuda & Hanawa, 1999). This 3D MHD nested grid code integrates equations (1) through (5) numerically, with a finite difference scheme on the Cartesian coordinates. The solution is a second order accurate, both in space and in time by virtue of the Monotone Upstream Scheme for Conservation Law (see e.g., Hirsh 1990). The Poisson equation is solved by the multigrid iteration (Matsumoto & Hanawa, 2003a). We have used Fujitsu VPP 5000, vector-parallel supercomputers, for 40 hours to make a typical model shown in this paper.
The nested grid consists of concentric hierarchical rectangular subgrids to gain high spatial resolution near the origin. Each rectangular grid has the same cell number ($`=128\times 128\times 32`$) but a different cell width, $`h(\mathrm{})\mathrm{\hspace{0.17em}2}^\mathrm{}5\lambda _{\mathrm{max}}`$, where $`\mathrm{}`$ denotes the level of the grid and ranges from 1 to $`\mathrm{}_{\mathrm{max}}`$. Thus the coarsest rectangular grid of $`\mathrm{}=\mathrm{\hspace{0.17em}1}`$ covers the whole computation region of $`\lambda _{\mathrm{max}}x\lambda _{\mathrm{max}}`$, $`\lambda _{\mathrm{max}}y\lambda _{\mathrm{max}}`$, and $`0z\lambda _{\mathrm{max}}/2`$. The solution in $`z<\mathrm{\hspace{0.17em}0}`$ is constructed from that in $`z\mathrm{\hspace{0.17em}0}`$ by the mirror symmetry with respect to $`z=\mathrm{\hspace{0.17em}0}`$. The maximum level number is set at $`\mathrm{}_{\mathrm{max}}=\mathrm{\hspace{0.17em}3}`$ at the initial stage ($`t=\mathrm{\hspace{0.17em}0}`$). A new finer subgrid is generated whenever the minimum local Jeans length $`\lambda _\mathrm{J}`$ becomes smaller than $`h(\mathrm{}_{\mathrm{max}})/8`$. Since the density is highest always in the finest subgrid, the generation of the new subgrid ensures the Jeans condition with a margin of a factor of 2 (c.f. Truelove et al., 1997). We have adopted the hyperbolic divergence cleaning method of Dedner et al. (2002) to obtain the magnetic field of $`\mathbf{}𝑩`$ free.
## 4 Results
We have followed all the models shown in this paper until the central density exceeds $`n_\mathrm{c}\mathrm{\hspace{0.17em}10}^{15}\mathrm{cm}^3`$. This paper describes the first half of the evolution for each model, i.e., the stages of $`n_\mathrm{c}n_{\mathrm{cr}}=\mathrm{\hspace{0.17em}5}\times 10^{10}\mathrm{cm}^3`$. The second half is described in the subsequent paper (Paper II).
Our models are characterized mainly by the strength of the magnetic field ($`\alpha `$) and the angular velocity ($`\omega `$). They are classified into four groups: (A) models having small $`\alpha (<0.1)`$ and small $`\omega (<0.1)`$, (B) those having large $`\alpha (0.1)`$ and small $`\omega (<0.1)`$, (C) those having small $`\alpha (<0.1)`$ and large $`\omega (0.1)`$, and (D) those having large $`\alpha (0.1)`$ and large $`\omega (0.1)`$. In other words, the model cloud has a weak magnetic field and rotates slowly in group A, while it has a relatively strong magnetic field and relatively large angular momentum in group D. Each group is described separately in the following subsections, in each of which two typical models of $`A_\phi =\mathrm{\hspace{0.17em}0.01}`$ (S) and 0.2 (L) are shown. The typical models are named after the group (A, B, C, or D) and $`A_\phi `$ (S or L). Model AS has $`\alpha =\mathrm{\hspace{0.17em}0.01}`$, $`\omega `$ = 0.01, and $`A_\phi `$ = 0.01, for example. Table 2 shows the values of $`\alpha `$, $`\omega `$, $`A_z`$, and $`A_\phi `$ for the 8 typical models shown in the following subsections. It also shows the initial magnetic field ($`B_{zc,0}`$), the initial angular velocity ($`\mathrm{\Omega }_{c0}`$), the wavelength of the perturbation in the $`z`$-direction ($`\lambda _{\mathrm{max}}`$), the mass ($`M`$) of the gas contained in the region of $`|z|\lambda _{\mathrm{max}}/2`$, and the epoch at which the density becomes infinity ($`t_\mathrm{f}`$).
### 4.1 Weak Magnetized and Slowly Rotating Cloud
This subsection displays model AS as a typical model having a weak magnetic field and slow rotation. Model AS has parameters $`\alpha =0.01`$, $`\omega =0.01`$ and $`A_\phi =0.01`$. Fig. 1 shows the cloud evolution in model AS by a series of cross sections.
As shown in Fig. 1, a gas cloud is transformed from a prolate one to an oblate one on the $`y=0`$ plane (see lower panels), while it maintains a round shape on the $`z=0`$ plane (upper panels) in the period of $`5.5\times 10^2\mathrm{cm}^3<n_c<7.6\times 10^{10}\mathrm{cm}^3`$. The velocity field is almost spherically symmetric while the central density increases from $`5\times 10^5\mathrm{cm}^3`$ to $`2\times 10^9\mathrm{cm}^3`$. The dense cloud is prolate and elongated in the $`z`$-direction in the lower panel of Fig. 1 (b), while it is nearly spherical in the lower panel of Fig. 1 (c). In this early collapse phase, the cloud contracts along the major axis (i.e. $`z`$-axis), regardless of the magnetic field and rotation as discussed in Bonnell et al. (1996). An oblate core is seen in panel Fig. 1 (e) and a thin disk is seen in the lower panel of Fig. 1 (f). The collapse is dynamical at the stages shown in Fig. 1 (c) through Fig. 1 (f). The radial infall velocity reaches $`v_r`$ = $`0.52`$ km s<sup>-1</sup> on the $`z=\mathrm{\hspace{0.17em}0}`$ plane in Fig. 1 (e), while the vertical infall velocity does $`v_z=\pm 0.58`$ km s<sup>-1</sup> on the $`z`$-axis. The rotation velocity is $`v_\phi `$ = 0.047 km s<sup>-1</sup> at maximum and much smaller than the infall velocities. This means gas contracts spherically in this phase. The difference between the radial and vertical infall velocities is still small ($`|v_{r,\mathrm{max}}|=\mathrm{\hspace{0.17em}0.61}`$ km s<sup>-1</sup> and $`|v_{z,\mathrm{max}}|=\mathrm{\hspace{0.17em}0.8}`$ km s<sup>-1</sup>) in Fig. 1 (f) although a high density disk is formed.
The density increase is well approximated by $`\rho _\mathrm{c}=\mathrm{\hspace{0.17em}2.2}/[4\pi G(tt_\mathrm{f})^2]`$ in the period of $`5\times 10^4\mathrm{cm}^3n\mathrm{\hspace{0.17em}10}^9\mathrm{cm}^3`$ as shown by the thick solid curve in Fig. 2. The offset is taken to be $`t_\mathrm{f}=\mathrm{\hspace{0.17em}5.96}\times 10^6`$yr so that the central density increases in proportion to the inverse square of the time in the widest span in $`\mathrm{log}\rho _\mathrm{c}`$, as shown in Larson (1969). Remember that the similarity solution of Larson (1969) and Penston (1969) gives $`\rho _\mathrm{c}=\mathrm{\hspace{0.17em}1.667}/[4\pi G(tt_\mathrm{f})^2]`$ for spherical collapse of a non-magnetized non-rotating isothermal cloud. We have checked that the density increase is well approximated by $`\rho _\mathrm{c}\mathrm{\hspace{0.17em}1.65}/[4\pi G(tt_\mathrm{f})^2]`$ in a non-magnetized and non-rotating cloud of our test calculation. The density increase is 15 % slower in model AS than in the similarity solution, since $`(2.2/1.667)^{1/2}\mathrm{\hspace{0.17em}1.15}`$. This small difference is due to the rotation and magnetic field.
To evaluate the change in the core shape shown in Fig. 1, we measure the moment of inertia for the high density gas of $`\rho \mathrm{\hspace{0.17em}0.1}\rho _\mathrm{c}`$. We derive the major axis ($`h_l`$), minor axis ($`h_s`$), and $`z`$-axis ($`h_z`$) from the moment of inertia according to Matsumoto & Hanawa (1999). The oblateness is defined as $`\epsilon _{\mathrm{ob}}(h_lh_s)^{1/2}/h_z`$ and the axis ratio is defined as $`\epsilon _{\mathrm{ar}}=h_l/h_s\mathrm{\hspace{0.17em}1}`$.
The oblateness is denoted by the thick solid curve as a function of time in Fig. 3 (a) and as a function of the central density in Fig. 3 (b). The oblateness is nearly constant at $`\epsilon _{\mathrm{ob}}=0.27`$ in the period of $`t4\times 10^6`$yr (or $`5\times 10^2\mathrm{cm}^3n_c5\times 10^3\mathrm{cm}^3`$). It increases and reaches $`\epsilon _{\mathrm{ob}}=1`$ at the stage of $`n_c=2\times 10^6\mathrm{cm}^3`$, which is shown in the lower panel of Fig. 1 (c). The oblateness reaches $`\epsilon _{\mathrm{ob}}=2.9`$, when the disk-like structure is formed at $`n_c=7.6\times 10^{10}\mathrm{cm}^3`$ as shown in the lower panel of Fig. 1 (f). The increase in $`\epsilon _{\mathrm{ob}}`$ is monotonic in the period of $`n_\mathrm{c}>5\times 10^3\mathrm{cm}^3`$.
The axis ratio is denoted by the thick solid curve as a function of time in Fig. 3 (c) and as a function of the central density in Fig. 3 (d). The axis ratio decreases from $`\epsilon _{\mathrm{ar}}`$ = 0.01 to $`7\times 10^4`$ after oscillating once over the period of $`t5\times 10^6`$ yr ( or $`n_\mathrm{c}<5\times 10^3\mathrm{cm}^3`$). It increases in proportion to $`n_\mathrm{c}^{1/6}`$ over the period of $`n_\mathrm{c}>5\times 10^3\mathrm{cm}^3`$. The growth rate of $`\epsilon _{\mathrm{ar}}`$ coincides with that of the bar mode growing in the spherical runaway collapse (Hanawa & Matsumoto 1999). The axis ratio grows up to $`3.5\times 10^3`$ in the isothermal collapse phase as shown in Fig. 3 (d). In order to examine dependence on the axis ratio, we compare models AS and AL, of which model parameters are the same except for the amplitude of the non-axisymmetric perturbation, $`A_\phi `$. The value of $`\epsilon _{\mathrm{ar}}`$ is 20 times larger in model AL than in model AS at a given stage and reaches $`\epsilon _{\mathrm{ar}}=6.8\times 10^2`$ at $`n_\mathrm{c}=n_{\mathrm{cri}}`$. The axis ratio is proportional to $`A_\phi `$. The oblateness is nearly the same in models AS and AL. The non-axisymmetric perturbation grows linearly in models AS and AL.
The magnetic flux density increases as the density increases. The left panel of Fig. 4 shows the square root of the ratio of the magnetic pressure to the gas pressure, $`B_{\mathrm{zc}}/(8\pi \rho _\mathrm{c}c_\mathrm{s}^2)^{1/2}`$, as a function of $`n_\mathrm{c}`$. Note that the ordinate is normalized by the initial value. It increases in proportion to one sixth the power of the density, i.e., $`B_{\mathrm{zc}}/(8\pi \rho _\mathrm{c}c_\mathrm{s}^2)^{1/2}n_\mathrm{c}^{1/6}`$, in the period of $`10^6\mathrm{cm}^3n_\mathrm{c}10^9\mathrm{cm}^3`$. This means that the magnetic field increases in proportion to $`B_{\mathrm{zc}}n_\mathrm{c}^{2/3}`$. This increase in $`B_{\mathrm{zc}}`$ is consistent with the spherical collapse of the core. When the collapse is spherically symmetric, the density and magnetic field increase inversely proportional to the cubic and square of the radius, respectively, since the magnetic field is frozen in the gas. Hence, the magnetic field is proportional to two thirds the power of the density, $`B_{\mathrm{zc}}\rho _\mathrm{c}^{2/3}`$.
After the central density exceeds $`10^9\mathrm{cm}^3`$, the growth of the magnetic field slows down. This slowdown coincides with the change in the core shape. The core is significantly oblate in the period of $`n_\mathrm{c}>10^9\mathrm{cm}^3`$. Remember that the magnetic field is proportional to the square root of the density ($`B_{\mathrm{zc}}n_c^{1/2}`$) when a magnetized disklike gas cloud collapses (Scott & Black 1980). This is because the disk is nearly in a hydrostatic equilibrium in the $`z`$-direction and the isothermal disk has the relation $`n\mathrm{\Sigma }^2`$. We use the terminology, the “disk collapse”, for this radial collapse of a disklike gas cloud. In the disk collapse, the magnetic flux density increases in proportion to the surface density ($`B_{\mathrm{zc}}\mathrm{\Sigma }_\mathrm{c}`$) since the gas is frozen in a magnetic flux tube. The relations, $`n_\mathrm{c}\mathrm{\Sigma }_\mathrm{c}^2`$ and $`B_{\mathrm{zc}}\mathrm{\Sigma }_\mathrm{c}`$, yield $`B_{\mathrm{zc}}n_\mathrm{c}^{1/2}`$. In the period of $`n_\mathrm{c}>10^9\mathrm{cm}^3`$, the growth rate of the magnetic field is intermediate between those for the spherical collapse and for the disk collapse. This is consistent with the density change over the same period.
As well as the magnetic flux density, the angular velocity of the core increases as the density increases. The right panel of Fig. 4 shows the ratio of the angular velocity to the magnetic field ($`\mathrm{\Omega }_\mathrm{c}/B_{\mathrm{zc}}`$) normalized by the initial value ($`\mathrm{\Omega }_{\mathrm{c},0}/B_{zc,0}`$) as a function of $`n_\mathrm{c}`$. The ratio is nearly constant at the initial value. This is because both the specific angular momentum ($`j`$) and the magnetic flux ($`\mathrm{\Phi }`$) are conserved for a central magnetic flux tube. Both the angular velocity and magnetic field increase proportionally to the inverse square of the tube radius. Hence the ratio is constant in both the spherical and the disk collapse. The conservation of the specific angular momentum implies that none of the magnetic torque, gravitational torque, and $`\phi `$-component of the pressure force are significant.
Since $`\mathrm{\Omega }_\mathrm{c}/B_{\mathrm{zc}}`$ is nearly constant, the angular velocity increases in proportion to $`n_\mathrm{c}^{2/3}`$ in the period of $`10^6\mathrm{cm}^3n_\mathrm{c}10^9\mathrm{cm}^3`$ and the growth of $`\mathrm{\Omega }_\mathrm{c}`$ slows down in the period of $`n_\mathrm{c}>10^9\mathrm{cm}^3`$. When measured in units of the free fall timescale, the angular velocity increases in proportion to $`\mathrm{\Omega }_\mathrm{c}(4\pi G\rho _\mathrm{c})^{1/2}n_\mathrm{c}^{1/6}`$ in the former period. The angular velocity in units of the free fall timescale denotes the square root of the ratio of the centrifugal force to the gravitational force. The magnetic field and rotation strengthen in the same manner during the spherical collapse, since both $`B_{\mathrm{zc}}(8\pi \rho _\mathrm{c}c_\mathrm{s}^2)^{1/2}`$ and $`\mathrm{\Omega }_\mathrm{c}(4\pi G\rho _\mathrm{c})^{1/2}`$ increase in proportion to $`\rho _\mathrm{c}^{1/6}`$.
Model AS is similar to model B of Matsumoto et al. (1997), although our model AS includes a very weak magnetic field. The magnetic field influences little the cloud collapse. Fig. 5 shows that the magnetic field is not twisted but pinched at the stages of $`n_\mathrm{c}\mathrm{\hspace{0.17em}10}^6`$ cm<sup>-3</sup> as shown in panels (d) - (f) of Fig. 5. Each panel denotes the magnetic field lines for the corresponding stage shown in each panel of Fig. 1. This weak magnetic field has no significant effect. When $`\alpha <\mathrm{\hspace{0.17em}0.1}`$ and $`\omega <0.1`$, the effects of magnetic field and rotation are very small.
### 4.2 Strongly Magnetized and Slowly Rotating Cloud
Model BL is shown as a typical example of models in this subsection having large $`\alpha `$ and small $`\omega `$. Model BL has parameters $`\alpha =0.1`$, $`\omega =0.01`$ and $`A_\phi =0.2`$. The parameters of model BL are the same as those of model AL except for $`\alpha `$, which is 0.01 for model AL and 0.1 for BL (Table 2). When $`\alpha \mathrm{\hspace{0.17em}0.1}`$, the magnetic pressure becomes comparable to the gas pressure in the course of cloud collapse and decelerates the radial collapse significantly. The magnetic braking is also effective in models BL and BS.
Also in model BL the high density core changes its form from prolate to oblate as the central density increases, as shown in Fig. 6, which is the same as Fig. 1 but for model BL. The change in the core shape is due to the magnetic field, which is amplified during the spherical collapse. The ratio of the magnetic pressure to the gas pressure is 0.11 at the stage of $`n_\mathrm{c}=2\times 10^6`$, while it is only 0.05 at the initial stage. Each panel of Fig. 6 denotes the density and velocity distribution at the stage of (a) $`n_\mathrm{c}=\mathrm{\hspace{0.17em}8.2}\times 10^3\mathrm{cm}^3`$, (b) $`5.6\times 10^4\mathrm{cm}^3`$, (c) $`7.9\times 10^6\mathrm{cm}^3`$, and (d) $`6.0\times 10^{10}\mathrm{cm}^3`$. At the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}8.2}\times 10^3\mathrm{cm}^3`$, the oblateness is $`\epsilon _{\mathrm{ob}}`$ = 0.58 in model BL \[Fig. 6(a)\] while $`\epsilon _{\mathrm{ob}}`$ = 0.45 in model AL \[Fig. 1(b)\] . The core is more oblate in model BL than in models AL and AS when compared at a given stage with the same central density \[Fig. 3(b)\]. The oblateness is $`\epsilon _{\mathrm{ob}}`$ = 5.3 at the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5}\times 10^{10}\mathrm{cm}^3`$ in model BL, while $`\epsilon _{\mathrm{ob}}=2.9`$ in model AS (see Fig. 3). The oblateness increases slowly over the period of $`3\times 10^3\mathrm{cm}^3n_\mathrm{c}10^{10}\mathrm{cm}^3`$ in model BL and is saturated around $`\epsilon _{\mathrm{ob}}5`$ in the period of $`n_\mathrm{c}10^{10}\mathrm{cm}^3`$.
As well as in models AS and AL, the axis ratio decreases from the initial value of $`\epsilon _{\mathrm{ar}}=0.2`$ to 0.015 in the period of $`n_\mathrm{c}\mathrm{\hspace{0.17em}10}^4\mathrm{cm}^3`$ in model BL. Then it switches to growing in proportion to $`n_\mathrm{c}^{1/6}`$ in the period of $`n_\mathrm{c}\mathrm{\hspace{0.17em}10}^4\mathrm{cm}^3`$. The core is elliptic on the $`xy`$ plane and the axis ratio is $`\epsilon _{\mathrm{ar}}`$ = 0.23 at the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5}\times 10^{10}\mathrm{cm}^3`$, as shown in the upper panel of Fig. 6 (d). The amplitude of the non-axisymmetic perturbation is linearly proportional to the initial amplitude. The axis ratio is always smaller by a factor of 20 in model BS than in model BL when compared at the stage of a given central density. Models BS and BL have the same model parameters except for $`A_\phi `$.
Since the core is appreciably oblate, the infall velocity is higher in the vertical direction than in the radial direction. At the stage of $`n_\mathrm{c}=5\times 10^{10}\mathrm{cm}^3`$ the radial infall velocity is $`v_r=0.46\mathrm{km}\mathrm{s}^1`$ at maximum in $`z=\mathrm{\hspace{0.17em}0}`$ plane while the vertical infall velocity is $`v_z=0.7\mathrm{km}\mathrm{s}^1`$ at maximum on the $`z`$-axis. The radial infall velocity is smaller and the vertical infall velocity is larger than in model AS. This asymmetry is due to the magnetic field. The rotation velocity is quite small ($`v_\phi =0.03\mathrm{km}\mathrm{s}^1`$) and the centrifugal force is negligible. The density increase due to collapse is slower in model BL than in model AS. The growth of the central density is well approximated by $`\rho _\mathrm{c}=\mathrm{\hspace{0.17em}2.5}/[4\pi G(tt_\mathrm{f})^2]`$. The growth rate is 7 % smaller than that of model AS at a given central density.
Also in model BL, the magnetic field strengthens as the density increases. The ratio of the magnetic pressure to the gas pressure increases very slowly in the period of $`n_\mathrm{c}10^6\mathrm{cm}^3`$ (see Fig. 4). It is saturated around $`(B_{\mathrm{zc}}/B_{\mathrm{zc},0})^2/(\rho _c/\rho _{c,0})1.3`$ in the period of $`n_\mathrm{c}10^6\mathrm{cm}^3`$. The magnetic field decelerates the radial collapse appreciably as shown earlier.
Fig. 7 shows the magnetic field lines at the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}6}\times 10^{10}\mathrm{cm}^3`$. The magnetic field lines break at the levels of $`z15`$ AU and 30 AU near the disk surface. The latter break corresponds to a fast-mode MHD shock, which is essentially the same as the shock waves seen in Norman, Wilson & Barton (1980), Matsumoto et al. (1997), and Nakamura et al. (1999). They are squeezed and vertical to the midplane below the shock front while they are open above the shock front. The disk formation is due to the magnetic field.
The magnetic field extracts angular momentum from the core. As shown in Fig. 4 (b), the ratio of angular velocity to the magnetic field decreases by 30 % in the period of $`5\times 10^2\mathrm{cm}^3<n_c5\times 10^{10}\mathrm{cm}^3`$ in model BL, although it remains constant in model AL. The decrease is due to the magnetic braking. The twisted magnetic field transfers the angular momentum of the core outwards. The specific angular momentum of the core is 70 % of the initial value at the stage of $`n_\mathrm{c}=5\times 10^{10}\mathrm{cm}^3`$. The angular velocity normalized by free-fall timescale \[$`\mathrm{\Omega }_c/(4\pi G\rho _\mathrm{c})^{1/2}`$\] increases slightly from 0.01 to 0.015 in the period of $`5\times 10^2<n_c5\times 10^{10}\mathrm{cm}^3`$ in model BL, while it spins up from 0.01 to 0.06 in model AL. We discuss this difference again in §4.5 in which we compare the increase in $`\mathrm{\Omega }_c`$ for various models.
The efficiency of the magnetic braking is qualitatively similar in models BS and BL.
Models BL and BS are similar to model C of Tomisaka (1995) and model B1 of Nakamura et al. (1999), although the earlier models include neither rotation nor non-axisymmetric perturbation. The rotation and non-axisymmetric perturbation have little effect on the cloud collapse in our models BL and BS. When the initial magnetic pressure is larger than a tenth of the gas pressure ($`\alpha 0.1`$), initially weak magnetic field is amplified during the collapse and affects the evolution of the core. The magnetic pressure decelerates the radial collapse and leads to disk formation. Also the magnetic braking is appreciable.
### 4.3 Weakly Magnetized and Rapidly Rotating Cloud
This subsection describes model CS as a typical example for models having small $`\alpha `$ and large $`\omega `$. Model CS has parameters of $`\alpha =\mathrm{\hspace{0.17em}0.01}`$, $`\omega =\mathrm{\hspace{0.17em}0.5}`$, and $`A_\phi =\mathrm{\hspace{0.17em}0.01}`$. When $`\omega \mathrm{\hspace{0.17em}0.1}`$, rotation affects the collapse of the cloud significantly.
In model CS, a rotating disk forms at an early stage of low central density. Each panel of Fig. 8 denotes the density and velocity distribution at the stages of (a) $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5.2}\times 10^3\mathrm{cm}^3`$, (b) $`6.5\times 10^4\mathrm{cm}^3`$, (c) $`5.7\times 10^6\mathrm{cm}^3`$, and (d) $`8.3\times 10^{10}\mathrm{cm}^3`$. The rotating disk is clearly seen at the stage of $`6.5\times 10^4\mathrm{cm}^3`$. The oblateness reaches $`\epsilon _{\mathrm{ob}}=3.0`$ at the stage of $`4.0\times 10^4\mathrm{cm}^3`$ and is saturated around $`\epsilon _{\mathrm{ob}}\mathrm{\hspace{0.17em}3.5}`$ in the period of $`5\times 10^5\mathrm{cm}^3n_\mathrm{c}<5\times 10^{10}\mathrm{cm}^3`$ as shown in Fig. 3 (b).
The axis ratio increases up to $`\epsilon _{\mathrm{ar}}=\mathrm{\hspace{0.17em}0.2}`$ by the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5}\times 10^{10}\mathrm{cm}^3`$ in model CS (see Fig. 3). At the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5}\times 10^{10}\mathrm{cm}^3`$, the axis ratio is larger in model CS than in model AS, while it is the same at the initial stage. The difference arises in the period of $`n_\mathrm{c}\mathrm{\hspace{0.17em}5}\times 10^4\mathrm{cm}^3`$. The axis ratio remains around 0.01 in model CS, while it decreases to $`7\times 10^4`$ in model AS. The axis ratio grows roughly in proportion to $`n_\mathrm{c}^{1/6}`$ in the period of $`n\mathrm{\hspace{0.17em}5}\times 10^4\mathrm{cm}^3`$ both in model AS and in model BS. We have confirmed that the non-axisymmetric perturbation is proportional to the initial perturbation by comparing with model CL of which initial parameters are the same as those of model CS except for $`A_\phi `$. The axis ratio is 20 times larger in model CL than in model CS in the period of $`n_\mathrm{c}\mathrm{\hspace{0.17em}1.0}\times 10^9\mathrm{cm}^3`$. The axis ratio reaches $`\epsilon _{\mathrm{ar}}=10.2`$ and the high density core has a bar shape at the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5}\times 10^{10}\mathrm{cm}^3`$ in model CL.
The increase in the central density is approximated by $`\rho _\mathrm{c}=\mathrm{\hspace{0.17em}6.2}/[4\pi G(tt_\mathrm{f})^2]`$. The rate of the increase is appreciably smaller than those of models AS and BS. It is 1.93 times smaller than that of the spherical collapse at a given central density. The relatively slow collapse is due to fast rotation.
In the period of $`n_\mathrm{c}5\times 10^5\mathrm{cm}^3`$ the cloud collapses mainly in the vertical direction along the magnetic field. Accordingly the magnetic field increases a little and the ratio of the magnetic pressure to the gas pressure decreases in this period (see Fig. 4). Note that the square root of the ratio of the magnetic pressure to the gas pressure decreases in proportion to $`\rho _\mathrm{c}^{1/2}`$ when the collapse is purely vertical along the magnetic field. Also, the angular velocity increases a little and decreases in proportion to $`\rho _\mathrm{c}^{1/2}`$ when measured in the freefall timescale.
In the period of $`5\times 10^2\mathrm{cm}^3<n_c5\times 10^{10}\mathrm{cm}^3`$, the magnetic field ($`B_{\mathrm{zc}}`$) strengthens and the angular velocity of the core ($`\mathrm{\Omega }_\mathrm{c}`$) continue to increases. However, the ratio of the magnetic pressure to the gas pressure remains nearly constant around $`(B_{\mathrm{zc}}/B_{\mathrm{zc},0})^2/(\rho _c/\rho _{c,0})=0.5`$ for $`5\times 10^5\mathrm{cm}^3<n_c5\times 10^{10}\mathrm{cm}^3`$ as shown in Fig. 4. The angular velocity measured in the freefall timescale is also nearly constant around $`\mathrm{\Omega }_c(4\pi G\rho _\mathrm{c})^{1/2}=0.2`$. In other words, both the magnetic field and angular velocity increase in proportion to $`\rho _\mathrm{c}^{1/2}`$. These dependences of $`B_{zc}`$ and $`\mathrm{\Omega }_\mathrm{c}`$ on $`\rho _\mathrm{c}`$ indicate that the core collapses in the radial direction while maintaining a disk shape. They are the same as those in the similarity solution for a self-gravitationally collapsing gas disk (Tomisaka 1995, 2002; Nakamura et al. 1995; Matsumoto et al. 1997; Saigo & Hanawa 1998).
The ratio of the angular velocity to the magnetic field is constant in the period of $`n_\mathrm{c}<10^7\mathrm{cm}^3`$. It increases up to 1.2 by the stage of $`n_c=5\times 10^{10}\mathrm{cm}^3`$. This increase is due to the torsional Alfvén wave. The magnetic braking is not significant in model CS. The ratio of the angular velocity to the magnetic field is constant during the collapse as shown in Fig. 4 (b). This confirms that the specific angular momentum is conserved. The magnetic field is twisted by fast rotation as shown in Fig. 9. It is also bent at the shock front as well as in model BS. The twisted magnetic field is too weak to have any appreciable dynamical effects.
The infall velocity is higher vertically than radially. The maximum infall velocity is $`v_{r,\mathrm{max}}=0.37\mathrm{km}\mathrm{s}^1`$ radially and $`v_{z,\mathrm{max}}=\pm 0.59\mathrm{km}\mathrm{s}^1`$ at the stage of $`n_c=5\times 10^{10}\mathrm{cm}^3`$. The maximum rotation velocity is $`v_\phi =0.36\mathrm{km}\mathrm{s}^1`$ and exceeds the sound speed at the same stage. Thus both the infall and rotation are supersonic. This dynamically infalling gas disk is similar to infalling envelopes observed in HL Tau (Hayashi et al., 1993) and L1551 IRS5 (Ohashi et al., 1996; Saito et al., 1996) in the sense that the radial infall velocity is comparable with the rotation velocity. The vertical inflow along the $`z`$-axis forms shock waves twice, once at the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5.7}\times 10^6\mathrm{cm}^3`$ \[see Fig. 8 (c)\] and at that of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}8.3}\times 10^{10}\mathrm{cm}^3`$ \[see Fig. 8 (d)\]. The former forms at $`z=\pm 4\times 10^3`$ AU, and the latter at $`z=\pm \mathrm{\hspace{0.17em}40}`$ AU. These shock waves are essentially the same as those seen in Matsumoto et al. (1997). The oblateness has a temporal maximum value at the stages of the shock formation.
### 4.4 Strongly Magnetized and Rapidly Rotating Cloud
This subsection describes model DL as a typical example of models having large $`\alpha `$ and large $`\omega `$. Model DL has parameters $`\alpha =1.0`$, $`\omega =0.5`$ and $`A_\phi =0.2`$. When $`\alpha 0.1`$ and $`\omega \mathrm{\hspace{0.17em}0.1}`$, both magnetic field and rotation affect the collapse of the cloud significantly. The magnetic braking is also effective.
Fig. 10 shows formation of a magnetized rotating disk that deforms to an elongated high density bar in model DL. Each panel denotes the density and velocity distribution at the stage of (a) $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5.6}\times 10^3\mathrm{cm}^3`$, (b) $`6.8\times 10^4\mathrm{cm}^3`$, (c) $`5.1\times 10^6\mathrm{cm}^3`$, and (d) $`5.3\times 10^{10}\mathrm{cm}^3`$. The high density gas has an oblateness of $`\epsilon _{\mathrm{ob}}=4.2`$ at the stage of $`n_c=6.8\times 10^4\mathrm{cm}^3`$. The oblateness reaches its maximum at $`n_\mathrm{c}4\times 10^5\mathrm{cm}^3`$ and oscillates around $`\epsilon _{\mathrm{ob}}5`$ in the period of $`4\times 10^5\mathrm{cm}^3n_\mathrm{c}\mathrm{\hspace{0.17em}5}\times 10^{10}\mathrm{cm}^3`$ (see Fig. 3). As a result of the strong magnetic field and fast rotation, the disk forms at an earlier stage in model DL than in the other models shown in the previous subsections.
The increase in the central density is approximated by $`n_\mathrm{c}=4.9/|4\pi G(tt_\mathrm{f})^2|`$. This is slower than in models AS and BS, although faster than in model CS. The density increase is faster in a model having a larger $`\alpha `$ for a given $`\omega \mathrm{\hspace{0.17em}0.1}`$. This is because a stronger magnetic field brakes the rotating core more effectively and the centrifugal force is reduced more. Remember that the density increase is slower in a model having a larger $`\alpha `$, when $`\omega 0.1`$. When the angular momentum of the cloud is very small, the centrifugal force is negligible and its reduction due to the magnetic braking is unimportant. A stronger magnetic field decelerates the collapse through higher magnetic pressure and tension. The magnetic field plays two roles: acceleration of the collapse through magnetic braking, and deceleration of the collapse through magnetic pressure and tension. The former dominates for $`\omega 0.1`$ while the latter dominates for $`\omega <0.1`$.
As shown in Fig. 10 (d) upper panel, the disk is elongated into a bar of $`\epsilon _{\mathrm{ar}}=15`$ at the stage of $`5.3\times 10^{10}\mathrm{cm}^3`$. As well as in model CL, the axis ratio remains nearly constant at the beginning and increases in proportion to $`n_\mathrm{c}^{1/6}`$ from an early stage of $`n_\mathrm{c}<5\times 10^3\mathrm{cm}^3`$ in model DL (see Fig. 3). An elongated bar forms in the models in which the non-axisymmetric perturbation is relatively large ($`A_\phi >0.2`$) and does not diminish in the early phase. Remember that the axis ratio decreases in the period of $`n_\mathrm{c}1.0\times 10^4\mathrm{cm}^3`$ in models AL and BL. The axis ratio increases in proportion to $`n_\mathrm{c}^{1/6}`$ in all models while the core collapses dynamically. The final axis ratio depends on the initial value and the amount of damping in the early phase. The initial damping is smaller when either the initial magnetic field or rotation is larger.
A similar bar structure is also seen in Durisen et al. (1986) and Bate (1998). The bar structure develops as a result of the bar mode instability, when the ratio of rotational to gravitational energy ($`\beta `$) of the core exceeds $`\beta =0.274`$. This $`\beta `$ is related to $`\omega `$ by $`\beta =\omega ^2`$, when the cloud is spherical and has constant density and angular velocity. Thus, the criterion for the ‘bar mode instability’, $`\beta >0.274`$, corresponds to $`\omega >0.523`$. The value of $`\omega `$ continues to decrease until it converges to $`\omega 0.2`$ as denoted in the following subsection. Therefore, the condition $`\omega >0.523`$ is never fulfilled, and the bar should form in model D by another mechanism. While the bar mode instability of Hanawa & Matsumoto (1999) works only in a dynamically collapsing cloud, the bar mode instability of Durisen et al. (1986) and Bate (1998) does in a cloud in hydrostatic equilibrium.
There is another evidence that the bar is formed not by fast rotation in our models. The bar forms also in the model of ($`\alpha `$, $`\omega `$, $`A_\phi `$) = (3, 0, 0.2), which is listed as model 56 in Table 2 of Paper II. The bar formation can not be due to rotation since the cloud does not rotate in this model. (See Table 2 of Paper II for the list of the models in which the bar forms at the end of the isothermal phase.)
As well as in models CS and CL, the vertical infall dominates over the radial infall in the period of $`n_\mathrm{c}7\times 10^4\mathrm{cm}^3`$ in models DS and DL. The ratio of the magnetic pressure to the gas pressure normalized by its initial value decreases to 0.5 in the period as shown in Fig. 4. Then it oscillates around 0.5 in the period of $`10^5n_c5\times 10^{10}\mathrm{cm}^3`$. The epoch of disk formation coincides with that at which the ratio of the magnetic pressure to the gas pressure reaches its first local minimum value.
The vertical inflow also forms shock waves twice in model DL. Fig. 10(c) shows the outer shock located at $`z=\pm 7000`$ AU. The flow is nearly vertical above the front while it is horizontal below. The epoch of shock formation coincides with that of the temporarily maximum oblateness as in model CL.
The magnetic braking slows the spin of the collapsing disk in model DL. The initial central angular velocity is $`\mathrm{\Omega }_c/(4\pi G\rho _\mathrm{c})^{1/2}=0.5`$ in both models CL and DL. The central angular velocity decreases to $`\mathrm{\Omega }_c/(4\pi G\rho _\mathrm{c})^{1/2}=0.21`$ by the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5}\times 10^4\mathrm{cm}^3`$ in model DL, while it decrease to 0.24 in model CL. The ratio of the angular velocity to the magnetic field ($`\mathrm{\Omega }_c/B_{zc}`$) normalized by the initial value ($`\mathrm{\Omega }_{\mathrm{c},0}/B_{\mathrm{zc},0}`$) decreases to 70 % of the initial value in model DL (see Fig. 4). The magnetic braking is effective in the period of $`n_\mathrm{c}7\times 10^5\mathrm{cm}^3`$ as in model BL. The ratio of the angular velocity to the magnetic field increases in the period of $`7\times 10^5\mathrm{cm}^3<n_\mathrm{c}<5\times 10^8\mathrm{cm}^3`$. This spin is due to the magnetic torque. The twist of the magnetic field is bounded by the shock front and the torsional Alfvén wave is reflected there. Thus the angular momentum is not released from the core in model DL.
The ratio of the magnetic pressure to gas pressure decreases from 0.5 to 0.1 in the period of $`5\times 10^2\mathrm{cm}^3n_\mathrm{c}7\times 10^5\mathrm{cm}^3`$, and remains around 0.1 in the period of $`n_\mathrm{c}7\times 10^5\mathrm{cm}^3`$. Thus, importance of the magnetic force relative to the centrifugal force decreases in models DS and DL. Fig. 11 illustrates the magnetic field for the stages shown in Fig. 10. The magnetic field lines are twisted but less pinched than in model CS. They are twisted at a higher $`z`$ in model DL than in model CS. As shown in Fig. 11 (d), the magnetic field is squeezed to stem vertically from the bar and the magnetic flux density is large in the bar. In Fig. 11 (d), magnetic field lines are bent at $`z40`$AU, which corresponds to the shock front. Inside the shock front, the magnetic field lines are ran vertically and hardly twisted, while twisted moderately outside of the shock front.
Models DL and DS have the same initial model parameters except for $`A_\phi `$, which is smaller by a factor of 20 in model DS. As a result, only the axis ratio differs appreciably between models DL and DS. A high density disk is seen at the stage of $`n_\mathrm{c}=5\times 10^{10}\mathrm{cm}^3`$ in model DS while the elongated bar is seen in model DL. The axis ratio is smaller by a factor 20 in model DL throughout the evolution.
### 4.5 Magnetic Flux - Spin Relation
The filamentary cloud fragments to form a high density core of $`n_\mathrm{c}>5\times 10^{10}`$ in all the models computed. The formation of the core is dynamical and the central density increases in proportion to the inverse square of the time. As the central density increases, the magnetic field increases roughly in proportion to a power of $`\rho _\mathrm{c}`$. The power index, however, differs and depends on the geometry of the collapse as shown in the previous subsections. Also the angular velocity increases in proportion to a power of $`\rho `$ and the power index depends on the geometry of the collapse and the strength of the magnetic field. To summarize the increase in $`B_{\mathrm{zc}}`$ and $`\mathrm{\Omega }_c`$ we have plotted the evolutionary locus of the core in Fig. 12. The abscissa denotes the square root of the ratio of the magnetic pressure to the gas pressure, $`B_{\mathrm{zc}}(8\pi c_\mathrm{s}^2\rho _\mathrm{c})^{1/2}`$, in the logarithmic scale. The ordinate denotes the angular velocity normalized by the freefall timescale, $`\mathrm{\Omega }_c(4\pi G\rho _\mathrm{c})^{1/2}`$, on the logarithmic scale. Each curve denotes the evolutionary locus for a model. The asterisks denote the initial stages. The circles, squares and triangles denote the stages of $`n_\mathrm{c}=5\times 10^4\mathrm{cm}^3`$, $`5\times 10^6\mathrm{cm}^3`$, and $`5\times 10^8\mathrm{cm}^3`$, respectively. Models without magnetic field are shown inside the upper left box. Models without rotation are shown inside the lower right box.
The evolutionary loci are systematically ordered in Fig. 12. They are aligned to evolve toward the upper right in the lower left region. Models AL and AS belong to this region of weak magnetic field and slow rotation. On the other hand, the loci are aligned to evolve toward the lower left in the upper region (fast rotation) and in the right region (strong magnetic field). Models CL, CS, DL, and DS belong to these regions. Models BL and BS appear in the middle of the panel. Their loci are nearly horizontal and the angular velocity measured on the freefall timescale does not increase as a result of the magnetic braking.
We can deduce several rules for the collapse of a magnetized rotating gas cloud from Fig. 12. First all the loci seem to converge on the curve,
$$\frac{\mathrm{\Omega }_c^2}{(0.2)^2\mathrm{\hspace{0.33em}4}\pi G\rho _\mathrm{c}}+\frac{B_{\mathrm{zc}}^2}{(0.36)^2\mathrm{\hspace{0.33em}8}\pi c_s^2\rho _c}=1.$$
(19)
Equation (19) is denoted by the thick solid curve in Fig. 12. We call this curve the magnetic flux - spin relation or $`B\mathrm{\Omega }`$ relation in the following. The first term of equation (19) is proportional to the square of the angular velocity normalized by the freefall timescale and accordingly is proportional to the ratio of the centrifugal force to the gravity. The second term is proportional to the ratio of the magnetic pressure to the gas pressure. The numerators are proportional to the anisotropic forces which suppress only the radial infall. Convergence to equation (19) indicates that the sum of the centrifugal and magnetic forces are regulated to be at a certain value. This rule involves the rules found by Matsumoto et al. (1997) and Nakamura et al. (1999) as a special case. The former showed that the ratio of the specific angular momentum to the core mass converges to a half of the critical value $`j=0.5(2\pi GM/c_\mathrm{s})`$ for models having no magnetic field. The latter showed that the ratio of the magnetic field to the surface density tends to be a half of the critical one, i.e., $`B_{\mathrm{zc}}=0.5(2\pi G)^{1/2}\mathrm{\Sigma }`$, for collapse of a non-rotating cloud. See the models shown inside the upper left box and those inside the lower right box to confirm that they also converge to equation (19).
The magnetic flux - spin relation is related to formation of the shock waves. The first shock wave forms exactly when the evolutionary locus reaches the $`B\mathrm{\Omega }`$ relation. After the shock formation, the evolutionary locus leaves it temporarily and reaches it again at the formation of the second shock wave. The shock strength is also related to the distance between the initial stage and the $`B\mathrm{\Omega }`$ relation on the $`B\mathrm{\Omega }`$ diagram. The shock wave is stronger in a model starting from a more distant place from equation (19) on the diagram. No shock wave forms in a model of which the initial stage is close to equation (19) on the diagram (see, e.g., the model of $`\alpha =0.1`$ and $`\omega =0.1`$ shown in Fig. 12).
Second, the magnetic braking is appreciable only in models having $`\alpha 0.1`$. The effect of the magnetic braking is evaluated from the slope on the diagram, $`d\mathrm{log}\mathrm{\Omega }_\mathrm{c}/d\mathrm{log}B_{\mathrm{zc}}`$. When the specific angular momentum is conserved, the slope is $`d\mathrm{log}\mathrm{\Omega }_\mathrm{c}/d\mathrm{log}B_{\mathrm{zc}}=1`$ as discussed in subsection 4.3. The slope differs appreciably from unity near the lower part of the $`B\mathrm{\Omega }`$ relation in Fig. 12. It is appreciably smaller than unity on the left-hand side of the $`B\mathrm{\Omega }`$ relation, while it is appreciably larger on the right-hand side. When the initial magnetic field is very weak, its magnetic torque is negligible. When the initial magnetic field is strong, the vertical collapse dominates. The magnetic braking reduces the specific angular momentum by $`3040\%`$ by the stage of $`n_c=\mathrm{\hspace{0.17em}10}^6\mathrm{cm}^3`$. However, it does not operate effectively beyond the stage. We will discuss the implication of equation (19) in the next section.
## 5 Discussion
As shown in the previous section, the magnetic flux density and angular velocity converge on equation (19) in Fig. 12, the $`B\mathrm{\Omega }`$ diagram. Thus, we can evaluate the magnetic flux density and angular velocity of the first core to be
$$\left(\frac{\mathrm{\Omega }_c}{2.57\times 10^3\mathrm{yr}^1}\right)^2+\left(\frac{B_{\mathrm{zc}}}{1.50\times 10^4\mu \mathrm{G}}\right)^2=1,$$
(20)
by substituting $`\rho _\mathrm{c}=\mathrm{\hspace{0.17em}1.92}\times 10^{13}`$ g$`\mathrm{cm}^3`$ (equivalent to $`n_\mathrm{c}=5\times 10^{10}\mathrm{cm}^3`$) and $`c_\mathrm{s}=\mathrm{\hspace{0.17em}0.19}\mathrm{km}\mathrm{s}^1`$, into equation (19). equation (20) implies that the first core has either the ‘standard’ magnetic flux density (15 mG) or the ‘standard’ angular velocity ($`2.57\times 10^3`$ yr<sup>-1</sup>), unless the initial magnetic field is very weak and the rotation is very slow. When both the magnetic flux density and angular velocity are negligible, the cloud collapses almost spherically and hence both $`B_{\mathrm{zc}}`$ and $`\mathrm{\Omega }_\mathrm{c}`$ increase in proportion to $`n_\mathrm{c}^{2/3}`$. Thus, either the magnetic flux density or the angular velocity reach the standard value at $`n_\mathrm{c}=5\times 10^{10}\mathrm{cm}^3`$ when either $`B1.50\mu \mathrm{G}`$ or $`\mathrm{\Omega }_\mathrm{c}2.57\times 10^7\mathrm{yr}^1`$ at the stage of $`n_\mathrm{c}=\mathrm{\hspace{0.17em}5}\times 10^4\mathrm{cm}^3`$. The latter condition is equivalent to $`|\mathbf{}\times 𝒗|\mathrm{\hspace{0.17em}2.49}\times 10^3\mathrm{km}\mathrm{s}^1\mathrm{pc}^1`$.
The standard magnetic flux density is approximately a half of the critical one, as mentioned in subsection 4.4. The latter is evaluated to be
$`B_{\mathrm{cr}}`$ $`=`$ $`2\pi \sqrt{G}\mathrm{\Sigma }`$ (21)
$`=`$ $`\sqrt{8\pi \rho c_\mathrm{s}^2}`$ (22)
at the limit of the geometrically thin self-gravitationally bound gas disk (Nakano & Nakamura, 1978). Also, the standard angular velocity is approximately a half of the critical one. The critical angular velocity is defined so that the centrifugal force balances with the gravity. Then it is evaluated to be
$$\mathrm{\Omega }_{\mathrm{cr}}=\sqrt{\frac{4\pi G\rho }{3}}$$
(23)
for a uniform gas sphere. Thus, equation (19) implies that either the magnetic flux density or the angular velocity is regulated to be a half of the critical value.
Equation (19) also predicts anti-correlation between the magnetic flux density and angular velocity of the first core. In other words, only one of the magnetic flux density or the angular velocity is close to the standard value. Then we can make a new index, the ratio of angular velocity to the magnetic flux density, for identifying whether the magnetic field dominates over rotation during the cloud collapse. If it is larger than the ratio of the standard values,
$`{\displaystyle \frac{\mathrm{\Omega }_{\mathrm{st}}}{B_{\mathrm{st}}}}`$ $`=`$ $`0.39\sqrt{G}c_{\mathrm{s}}^{}{}_{}{}^{1}`$ (24)
$`=`$ $`1.69\times 10^7\left({\displaystyle \frac{c_\mathrm{s}}{0.19\mathrm{km}\mathrm{s}^1}}\right)^1\mathrm{yr}^1\mu \mathrm{G}^1,`$ (25)
the centrifugal force dominates over the magnetic force. Otherwise the magnetic force dominates over the centrifugal force. This analysis suggests that there exist two types of first core: magnetic first core and spinning first core. We discuss the difference between them in Paper II.
We shall apply the above discussion to L1544, the prestellar core, of which rotation and magnetic field have been measured. The rotation velocity is evaluated to be $`0.09\mathrm{km}\mathrm{s}^1`$ at $`r`$ = 15000 AU by Ohashi et al. (1999) and $`0.14\mathrm{km}\mathrm{s}^1`$ at $`r`$ = 7000 AU by Williams et al. (1999). These velocity gradients correspond to $`1.26\times 10^6`$ yr<sup>-1</sup> and $`4.21\times 10^6`$ yr<sup>-1</sup>. On the other hand, the line-of-sight magnetic field is evaluated to be $`+11\pm 2\mu \mathrm{G}`$ by Crutcher & Troland (2000). Combining these values, we obtain $`\mathrm{\Omega }/B`$ = $`1.1\times 10^7`$ yr<sup>-1</sup> $`\mu \mathrm{G}^1`$ and $`3.8\times 10^7`$ yr<sup>-1</sup> $`\mu \mathrm{G}^1`$. If we take account of uncertainty of the observed values, the magnetic force dominates over the centrifugal force only marginally.
It should be noted that Crutcher et al. (2004) derived a much stronger magnetic field ($`140\mu `$G) for L1544 from linear polarization of the dust emission. They derived the value under the assumption that the randomness of the magnetic field can be ascribed to turbulent motion. If the magnetic field is as strong as 140 $`\mu `$G at the distance of 10000 AU, the magnetic force should dominate over the centrifugal force. However, their method gives a magnetic field an order of magnitude stronger compared with the values derived by the Zeeman effect. Possible systematic errors should be examined.
Next, we discuss the speed of dynamical collapse in the molecular cloud core. Aikawa et al. (2001) discussed the possibility of deriving the collapse speed from the chemical abundance in the prestellar core L1544. They computed chemical evolution in a molecular cloud core, assuming that the density evolution is the same as that of the Larson-Penston similarity solution, or by a factor $`f`$ slower. They concluded that the observed chemical anomaly in L1544 is consistent with the model based on the Larson-Penston similarity solution from comparison with the slow collapse models of $`f`$ = 3 and 10. The model of $`f`$ = 3 is supposed to mimic a molecular cloud core of which collapse is slowed down owing to rotation, magnetic field, or turbulence. Our simulation has shown that the slowing by magnetic field and rotation is appreciably smaller. The slow-down factor is evaluated to be $`f`$ = 1.22 for model BS and 1.93 for model CL. The small slow-down factor makes the chemical diagnosis harder.
Finally, we discuss the effect of the ambipolar diffusion. The evolution of magnetically subcritical cloud including the ambipolar diffusion has been investigated by Basu & Mouschovious (1994, 1995a, 1995b) under the disk approximation. They showed that the magnetically supercritical core is formed in the subcritical cloud for ambipolar diffusion after 10-20 freefall time passed. Once the supercritical core is formed, the magnetic field is hardly extracted from the core, because the ambipolar diffusion is much slower than the freefall (Basu & Mouschovias, 1994). Thus, our ideal MHD approximation is valid since our model cloud is supercritical from the initial stage (see Table 2). The ambipolar diffusion may have an important role in the initially subcritical cloud and after the formation of the dense core ($`n_c>10^{11}\mathrm{cm}^3`$).
## Acknowledgments
We have greatly benefited from discussion with Prof. M. Y. Fujimoto, Prof. A. Habe and Dr. K. Saigo. Numerical calculations were carried out with a Fujitsu VPP5000 at the Astronomical Data Analysis Center, the National Astronomical Observatory of Japan. This work was supported partially by the Grants-in-Aid from MEXT (15340062, 14540233 \[KT\], 16740115 \[TM\]).
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# Raising/Lowering Maps and Modules for the Quantum Affine Algebra 𝑈_𝑞((𝔰𝔩)̂₂) Keywords: Quantum group, quantum affine algebra, affine Lie algebra (𝔰𝔩)̂₂, raising and lowering maps, tridiagonal pair. 2000 Mathematics Subject Classification. Primary: 17B37. Secondary: 16W35, 20G42, 81R50
## 1 Introduction
A quantum affine algebra is a $`q`$-analogue of the universal enveloping algebra of a Kac-Moody Lie algebra of affine type. Quantum affine algebras were first introduced and studied by V.G. Drinfeld and M. Jimbo in relation to the Yang-Baxter equation of mathematical physics. Since then quantum affine algebras have played an important role in various areas of mathematics and physics, for example see , , , . In this paper we will be concerned with the quantum affine algebra $`U_q(\widehat{𝔰𝔩}_2)`$.
We study the finite dimensional modules for $`U_q(\widehat{𝔰𝔩}_2)`$. In V. Chari and A. Pressley classified the finite dimensional irreducible $`U_q(\widehat{𝔰𝔩}_2)`$-modules up to isomorphism. These modules were further studied in , . However, a complete classification of all finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-modules is still unknown. In the present paper we obtain a result that may help in this classification. We now introduce some notation and recall the definition of $`U_q(\widehat{𝔰𝔩}_2)`$.
Throughout this paper $`𝕂`$ will denote an algebraically closed field. We fix a nonzero scalar $`q𝕂`$ that is not a root of unity. We will use the following notation:
$`[n]={\displaystyle \frac{q^nq^n}{qq^1}},n=0,1,\mathrm{}.`$
###### Definition 1.1
\[4, Definition 2.2\] The quantum affine algebra $`U_q(\widehat{𝔰𝔩}_2)`$ is the unital associative $`𝕂`$-algebra with generators $`e_i^\pm `$, $`K_i^{\pm 1}`$, $`i\{0,1\}`$, which satisfy the following relations:
$`K_iK_i^1`$ $`=`$ $`K_i^1K_i=1,`$ (1)
$`K_0K_1`$ $`=`$ $`K_1K_0,`$ (2)
$`K_ie_i^\pm K_i^1`$ $`=`$ $`q^{\pm 2}e_i^\pm ,`$ (3)
$`K_ie_j^\pm K_i^1`$ $`=`$ $`q^2e_j^\pm ,ij,`$ (4)
$`e_i^+e_i^{}e_i^{}e_i^+`$ $`=`$ $`{\displaystyle \frac{K_iK_i^1}{qq^1}},`$ (5)
$`e_0^\pm e_1^{}`$ $`=`$ $`e_1^{}e_0^\pm ,`$ (6)
$`(e_i^\pm )^3e_j^\pm [3](e_i^\pm )^2e_j^\pm e_i^\pm +[3]e_i^\pm e_j^\pm (e_i^\pm )^2e_j^\pm (e_i^\pm )^3=0,ij.`$ (7)
We call $`e_i^\pm `$, $`K_i^{\pm 1}`$ the Chevalley generators for $`U_q(\widehat{𝔰𝔩}_2)`$ and refer to (7) as the cubic $`q`$-Serre relations.
We now give two definitions, state our main result, and then make some comments concerning its significance.
###### Definition 1.2
Let $`V`$ denote a vector space over $`𝕂`$ with finite positive dimension. By a
decomposition of $`V`$ we mean a sequence $`U_0,U_1,\mathrm{},U_d`$ consisting of nonzero subspaces of $`V`$ such that $`V=_{i=0}^dU_i`$ (direct sum). For notational convenience we set $`U_1:=0,U_{d+1}:=0`$.
###### Definition 1.3
Let $`V`$ denote a vector space over $`𝕂`$ with finite positive dimension. Let $`U_0,U_1,\mathrm{},U_d`$ denote a decomposition of $`V`$. Let $`K:VV`$ denote the linear transformation such that, for $`0id`$, $`U_i`$ is an eigenspace for $`K`$ with eigenvalue $`q^{2id}`$. We refer to $`K`$ as the linear transformation that corresponds to the decomposition $`U_0,U_1,\mathrm{},U_d`$.
###### Note 1.4
With reference to Definition 1.3, we note that $`K`$ is invertible. Moreover, for $`0id`$, $`U_i`$ is the eigenspace for $`K^1`$ with eigenvalue $`q^{d2i}`$. We observe that $`K^1`$ is the linear transformation that corresponds to the decomposition $`U_d,U_{d1},\mathrm{},U_0`$.
We will be concerned with the following situation.
###### Assumption 1.5
Let $`V`$ denote a vector space over $`𝕂`$ with finite positive dimension. Let $`U_0,U_1,\mathrm{},U_d`$ denote a decomposition of $`V`$. Let $`K`$ denote the linear transformation that corresponds to $`U_0,U_1,\mathrm{},U_d`$ as in Definition 1.3. Let $`R:VV`$ and $`L:VV`$ be linear transformations such that
1. $`RU_iU_{i+1}(0id1),RU_d=0`$,
2. $`LU_iU_{i1}(1id),LU_0=0`$,
3. for $`0id/2`$ the restriction $`R^{d2i}|_{U_i}:U_iU_{di}`$ is a bijection,
4. for $`0id/2`$ the restriction $`L^{d2i}|_{U_{di}}:U_{di}U_i`$ is a bijection,
5. $`R^3L[3]R^2LR+[3]RLR^2LR^3=0`$,
6. $`L^3R[3]L^2RL+[3]LRL^2RL^3=0`$.
We now state our main result.
###### Theorem 1.6
Adopt Assumption 1.5. Then there exists a unique $`U_q(\widehat{𝔰𝔩}_2)`$-module structure on $`V`$ such that $`(Re_1^{})V=0`$, $`(Le_0^{})V=0`$, $`(KK_0)V=0`$, $`(K^1K_1)V=0`$, where $`e_1^{}`$, $`e_0^{}`$, $`K_0`$, $`K_1`$ are Chevalley generators for $`U_q(\widehat{𝔰𝔩}_2)`$ as in Definition 1.1.
The proof of Theorem 1.6 will take up most of the paper until Section 9. In Sections 10 and 11 we determine which $`U_q(\widehat{𝔰𝔩}_2)`$-modules arise from the construction of Theorem 1.6.
###### Remark 1.7
Not all finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-modules arise from the construction of Theorem 1.6. However as we will see, every finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-module is a direct sum of submodules, each of which arises from Theorem 1.6 up to a routine normalization. Thus Theorem 1.6 can be viewed as a step towards the classification of the finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-modules.
###### Remark 1.8
The proof of Theorem 1.6 involves modifying a construction used in and . The construction originally arose from the study of tridiagonal pairs. According to \[8, Definition 1.1\] a tridiagonal pair is an ordered pair $`(A,A^{})`$ of diagonalizable linear transformatons on a finite dimensional vector space $`V`$ such that (i) the eigenspaces of $`A`$ (resp. $`A^{}`$) can be ordered as $`V_0,V_1,\mathrm{},V_d`$ (resp. $`V_0^{},V_1^{},\mathrm{},V_d^{}`$) with $`A^{}V_iV_{i1}+V_i+V_{i+1}`$ (resp. $`AV_i^{}V_{i1}^{}+V_i^{}+V_{i+1}^{}`$) for $`0id`$; (ii) there are no nonzero proper subspaces of $`V`$ which are invariant under both $`A`$ and $`A^{}`$. See , , for connections between tridiagonal pairs and $`U_q(\widehat{𝔰𝔩}_2)`$.
###### Remark 1.9
In G. Benkart and P. Terwilliger determine the finite dimensional irreducible modules for the standard Borel subalgebra of $`U_q(\widehat{𝔰𝔩}_2)`$. The authors adopt Assumption 1.5(i),(ii),(v),(vi) but replace Assumption 1.5(iii),(iv) with the assumption that $`V`$ is irreducible as a $`(K,R,L)`$-module. From this assumption they obtain a $`U_q(\widehat{𝔰𝔩}_2)`$-module structure on $`V`$ as in Theorem 1.6. The $`U_q(\widehat{𝔰𝔩}_2)`$-module structure that they obtain is irreducible while the $`U_q(\widehat{𝔰𝔩}_2)`$-module structure given by our Theorem 1.6 is not necessarily irreducible. As far as we know Theorem 1.6 does not imply the result in nor does the result in imply Theorem 1.6.
###### Remark 1.10
In the proof of Theorem 1.6 we will need a number of lemmas that are similiar to lemmas appearing in . The assumptions in Theorem 1.6 and in are similiar except that in the authors assume $`V`$ is irreducible as a $`(K,R,L)`$-module and we are not making this assumption. However, in this irreducibility condition is not used in the proofs of the lemmas we require. In such cases we simply cite without proof.
## 2 Preliminaries
In this section we make a few observations about Assumption 1.5.
###### Note 2.1
With reference to Assumption 1.5, if we replace $`U_i`$ by $`U_{di}`$ for $`0id`$, and replace $`(K,R,L)`$ by $`(K^1,L,R)`$, then the assumption is still satisfied.
###### Lemma 2.2
With reference to Assumption 1.5, for $`0id/2`$ and $`0jd2i`$, the restriction $`R^j|_{U_i}:U_iU_{i+j}`$ is an injection.
Proof: Immediate from Assumption 1.5(iii). $`\mathrm{}`$
###### Lemma 2.3
With reference to Assumption 1.5, the following (i),(ii) hold.
1. $`KR=q^2RK`$,
2. $`KL=q^2LK`$.
Proof: Immediate from Assumption 1.5 (i),(ii) and Definition 1.3. $`\mathrm{}`$
## 3 An outline of the proof for Theorem 1.6
Our proof of Theorem 1.6 will consume most of the paper from Section 4 to Section 9. Here we sketch an overview of the argument.
We begin by adopting Assumption 1.5. To start the construction of the $`U_q(\widehat{𝔰𝔩}_2)`$-action on $`V`$ we require that the linear transformations $`Re_1^{}`$, $`Le_0^{}`$, $`K^{\pm 1}K_0^{\pm 1}`$, and $`K^{\pm 1}K_1^1`$ vanish on $`V`$. This gives the actions of the elements $`e_1^{},e_0^{},K_0^{\pm 1},K_1^{\pm 1}`$ on $`V`$. We define the actions of $`e_0^+,e_1^+`$ on $`V`$ as follows. First we prove that $`K+R`$ is diagonalizable on $`V`$. Then we show that the set of distinct eigenvalues of $`K+R`$ on $`V`$ is $`\{q^{2id}|\mathrm{\hspace{0.17em}0}id\}`$. For $`0id`$ we let $`V_i`$ denote the eigenspace of $`K+R`$ on $`V`$ associated with the eigenvalue $`q^{2id}`$. So $`V_0,V_1,\mathrm{},V_d`$ is a decomposition of $`V`$. Next we define the subspaces $`W_i`$ as follows.
$`W_i=(U_0+\mathrm{}+U_i)(V_0+\mathrm{}+V_{di})(0id).`$
We show that $`W_0,W_1,\mathrm{},W_d`$ is a decomposition of $`V`$. Then we apply Note 2.1 to the above argument to obtain the following results. $`K^1+L`$ is diagonalizable on $`V`$. The set of distinct eigenvalues of $`K^1+L`$ on $`V`$ is $`\{q^{d2i}|\mathrm{\hspace{0.17em}0}id\}`$. For $`0id`$ we let $`V_i^{}`$ denote the eigenspace of $`K^1+L`$ on $`V`$ associated with the eigenvalue $`q^{d2i}`$. So $`V_0^{},V_1^{},\mathrm{},V_d^{}`$ is a decomposition of $`V`$. Next we define the subspaces $`W_i^{}`$ as follows.
$`W_i^{}=(U_i+\mathrm{}+U_d)(V_{di}^{}+\mathrm{}+V_d^{})(0id).`$
Then $`W_0^{},W_1^{},\mathrm{},W_d^{}`$ is a decomposition of $`V`$. Next we define the linear transformation $`B:VV`$ (resp. $`B^{}:VV`$) such that for $`0id`$, $`W_i`$ (resp. $`W_i^{}`$) is an eigenspace for $`B`$ (resp. $`B^{}`$) with eigenvalue $`q^{2id}`$ (resp. $`q^{d2i}`$). We let $`e_1^+`$ act on $`V`$ as $`IK^1B`$ times $`q^1(qq^1)^2`$. We let $`e_0^+`$ act on $`V`$ as $`IKB^{}`$ times $`q^1(qq^1)^2`$. We display some relations that are satisfied by $`B,B^{},L,R,K,K^1`$. Using these relations we argue that the above actions of $`e_0^\pm ,e_1^\pm ,K_0^{\pm 1},K_1^{\pm 1}`$ satisfy the defining relations for $`U_q(\widehat{𝔰𝔩}_2)`$. In this way we obtain the required action of $`U_q(\widehat{𝔰𝔩}_2)`$ on $`V`$.
## 4 The linear transformation $`A`$
In this section we define and discuss a linear transformation that will be useful.
###### Definition 4.1
With reference to Assumption 1.5, let $`A:VV`$ denote the following linear transformation:
$`A=K+R.`$
###### Lemma 4.2
With reference to Definition 4.1 and Assumption 1.5, the following (i),(ii) hold.
1. For $`0id`$ the action of $`Aq^{2id}I`$ on $`U_i`$ coincides with the action of $`R`$ on $`U_i`$,
2. $`(Aq^{2id}I)U_iU_{i+1}`$, $`0id`$.
Proof: Immediate from Definition 4.1, Definition 1.3, and Assumption 1.5(i). $`\mathrm{}`$
###### Lemma 4.3
\[3, Lemma 4.13\] With reference to Definition 4.1 and Assumption 1.5, the following holds. $`A`$ is diagonalizable on $`V`$ and the set of distinct eigenvalues of $`A`$ is
$`\{q^{2id}|\mathrm{\hspace{0.17em}0}id\}`$. Moreover, for $`0id`$, the dimension of the eigenspace for $`A`$ associated with $`q^{2id}`$ is equal to the dimension of $`U_i`$.
###### Definition 4.4
With reference to Definition 4.1 and Lemma 4.3, for $`0id`$ we let $`V_i`$ denote the eigenspace for $`A`$ with eigenvalue $`q^{2id}`$. For notational convenience we set $`V_1:=0,V_{d+1}:=0`$. We observe that $`V_0,V_1,\mathrm{},V_d`$ is a decomposition of $`V`$.
###### Lemma 4.5
Let the decomposition $`U_0,U_1,\mathrm{},U_d`$ be as in Assumption 1.5 and let the decomposition $`V_0,V_1,\mathrm{},V_d`$ be as in Definition 4.4. Then for $`0id`$ the spaces $`U_i`$, $`U_{di}`$, $`V_i`$, $`V_{di}`$ all have the same dimension.
Proof: Immediate from Lemma 4.3 and Assumption 1.5 (iii). $`\mathrm{}`$
###### Definition 4.6
With reference to Lemma 4.5, for $`0id`$ we let $`\rho _i`$ denote the common dimension of $`U_i`$, $`U_{di}`$, $`V_i`$, $`V_{di}`$.
###### Lemma 4.7
With reference to Definition 4.6, the following (i)–(iii) hold.
1. $`\rho _i0`$, $`0id`$,
2. $`dim(V)=_{i=0}^d\rho _i`$,
3. $`\rho _i=\rho _{di}`$, $`0id`$.
Proof: Immediate by Definition 4.6 and since $`U_0,U_1,\mathrm{},U_d`$ is a decomposition of $`V`$. $`\mathrm{}`$
###### Lemma 4.8
\[3, Lemma 5.2\] With reference to Assumption 1.5 and Definition 4.4,
$`U_i+\mathrm{}+U_d=V_i+\mathrm{}+V_d,0id.`$
###### Lemma 4.9
\[3, Lemma 5.3\] With reference to Assumption 1.5 and Definition 4.4,
$`(K^1q^{d2i})V_iV_{i+1},0id.`$ (8)
## 5 The subspaces $`W_i`$
###### Definition 5.1
With reference to Assumption 1.5 and Definition 4.4, we define
$`W_i=(U_0+\mathrm{}+U_i)(V_0+\mathrm{}+V_{di}),0id.`$
For notational convenience we set $`W_1:=0,W_{d+1}:=0`$.
The goal of this section is to prove the following theorem.
###### Theorem 5.2
With reference to Definition 5.1, the sequence $`W_0,W_1,\mathrm{},W_d`$ is a decomposition of $`V`$.
We prove Theorem 5.2 in three steps. First, we show the sum $`_{i=0}^dW_i`$ is direct. Second, we show $`V=_{i=0}^dW_i`$. Finally, we show $`W_i0`$ for $`0id`$.
The following definition and the next few lemmas will be useful in proving the sum $`_{i=0}^dW_i`$ is direct.
###### Definition 5.3
With reference to Assumption 1.5 and Definition 4.4, we define
$`W(i,j)=\left({\displaystyle \underset{h=0}{\overset{i}{}}}U_h\right)\left({\displaystyle \underset{h=0}{\overset{j}{}}}V_h\right),1i,jd+1.`$
With reference to Definition 5.1, note that $`W(i,di)=W_i`$ for $`0id`$.
###### Lemma 5.4
With reference to Assumption 1.5, Definition 4.1, and Definition 5.3, the following (i),(ii) hold.
1. $`(Aq^{2jd}I)W(i,j)W(i+1,j1)`$, $`0i,jd`$,
2. $`(K^1q^{d2i}I)W(i,j)W(i1,j+1)`$, $`0i,jd`$.
Proof: (i) Using Definition 5.3 and Definition 4.4, we have $`(Aq^{2jd}I)W(i,j)_{h=0}^{j1}V_h`$. Using Definition 5.3 and Lemma 4.2(ii) we have $`(Aq^{2jd}I)W(i,j)_{h=0}^{i+1}U_h`$. Combining these facts we obtain the desired result.
(ii) Using Definition 5.3 and Lemma 4.9, we have $`(K^1q^{d2i}I)W(i,j)_{h=0}^{j+1}V_h`$. Using Definition 5.3 and Note 1.4 we have $`(K^1q^{d2i}I)W(i,j)_{h=0}^{i1}U_h`$. Combining these facts we obtain the desired result. $`\mathrm{}`$
###### Lemma 5.5
With reference to Definition 5.3,
$`W(i,d1i)=0,0id1.`$
Proof: Define $`T=_{i=0}^{d1}W(i,d1i)`$. It suffices to show $`T=0`$. By Lemma 5.4(ii) we find $`K^1TT`$. Recall that $`K^1`$ is diagonalizable on $`V`$ and so $`K^1`$ is diagonalizable on $`T`$. Also, $`U_jT`$ $`(0jd)`$ are the eigenspaces for $`K^1|_T`$. Thus $`T=_{j=0}^d(U_jT)`$ (direct sum). Suppose, towards a contradiction, $`T0`$. Then there exists $`j`$ $`(0jd)`$ such that $`U_jT0`$. Define $`t:=\text{min}\{j|\mathrm{\hspace{0.17em}0}jd,U_jT0\}`$ and $`r:=\text{max}\{j|\mathrm{\hspace{0.17em}0}jd,U_jT0\}`$. Of course $`tr`$. We will now show
$`r+td.`$ (9)
If $`d/2<t`$ then (9) holds since $`tr`$. So now assume $`0td/2`$. Let $`xU_tT`$ such that $`x0`$. By Assumption 1.5(i), we find $`R^{d2t}xU_{dt}`$. Also, by Assumption 1.5(iii), we find $`R^{d2t}x0`$. By Lemma 5.4(i), we find $`ATT`$. Using this and Lemma 4.2(i), we find $`R^{d2t}xT`$. So $`R^{d2t}xU_{dt}T`$. Combining these facts we find $`U_{dt}T0`$. This shows (9). Define $`y:=\text{max}\{j|\mathrm{\hspace{0.17em}0}jd1,W(j,d1j)0\}`$. By the definition of $`T`$ we find $`TU_0+\mathrm{}+U_y`$ and so
$`yr.`$ (10)
We will now show
$`dyt+1.`$ (11)
By the definition of $`y`$ we have $`W(y,d1y)0`$. By Definition 5.3 we have $`W(y,d1y)V_0+\mathrm{}+V_{d1y}`$. Using these facts and since $`V_0,V_1,\mathrm{},V_d`$ is a decomposition of $`V`$ we find $`W(y,d1y)V_{dy}+\mathrm{}+V_d`$. Therefore $`TV_{dy}+\mathrm{}+V_d`$. Using this and Lemma 4.8 we find $`TU_{dy}+\mathrm{}+U_d`$ and (11) follows. Adding (9), (10), and (11) we find $`01`$ for a contradiction. The result follows. $`\mathrm{}`$.
###### Lemma 5.6
With reference to Definition 5.1, the sum $`_{i=0}^dW_i`$ is direct.
Proof: It suffices to show $`(W_0+\mathrm{}+W_{i1})W_i=0`$ for $`1id`$. Let $`i`$ be given. By Definition 5.1, $`W_0+\mathrm{}+W_{i1}U_0+\mathrm{}+U_{i1}`$. Also by Definition 5.1, $`W_iV_0+\mathrm{}+V_{di}`$. By this and Definition 5.3 we find $`(W_0+\mathrm{}+W_{i1})W_iW(i1,di)`$. But $`W(i1,di)=0`$ by Lemma 5.5 and so $`(W_0+\mathrm{}+W_{i1})W_i=0`$. $`\mathrm{}`$
The following definition and the next few lemmas will be useful in proving $`V=_{i=0}^dW_i`$.
###### Definition 5.7
With reference to Assumption 1.5, we define
$`H_i=\{vU_i|R^{d2i+1}v=0\},0id/2.`$
###### Lemma 5.8
With reference to Assumption 1.5, Definition 4.4, and Definition 5.7,
$`H_i=(V_i+\mathrm{}+V_{di})U_i,0id/2.`$
Proof: Immediate from Definition 4.4, Definition 5.7, and Lemma 4.2(i). $`\mathrm{}`$
###### Lemma 5.9
With reference to Assumption 1.5 and Definition 5.7,
$`U_i={\displaystyle \underset{j=0}{\overset{min(i,di)}{}}}R^{ij}H_j(directsum),0id.`$
Proof: Case 1: $`0id/2`$. The proof is by induction on $`i`$. Observe the result holds for $`i=0`$ since $`U_0=H_0`$ by Definition 5.7 and Assumption 1.5(i). Next assume $`i1`$. By induction and Lemma 2.2 we find
$`RU_{i1}={\displaystyle \underset{j=0}{\overset{i1}{}}}R^{ij}H_j(directsum).`$ (12)
We now show
$`U_i=RU_{i1}+H_i(directsum).`$ (13)
Using Assumption 1.5(i) and Definition 5.7, we have $`RU_{i1}+H_iU_i`$. We now show $`U_iRU_{i1}+H_i`$. Let $`xU_i`$. By Assumption 1.5(i),(iii) there exists $`yU_{i1}`$ such that $`R^{d2i+2}y=R^{d2i+1}x`$. Using this we find $`xRyH_i`$. So $`xRU_{i1}+H_i`$. We have now shown equality in (13). It remains to show that the sum in (13) is direct. To do this we show $`RU_{i1}H_i=0`$. Let $`xRU_{i1}H_i`$. By Definition 5.7 we have $`R^{d2i+1}x=0`$. Also, there exists $`yU_{i1}`$ such that $`x=Ry`$. Combining these facts with Assumption 1.5(iii) we find $`y=0`$ and then $`x=0`$. We have now shown the sum in (13) is direct and this completes the proof of (13). Combining (12) and (13) we find
$`U_i={\displaystyle \underset{j=0}{\overset{i}{}}}R^{ij}H_j(directsum).`$
Case 2: $`d/2<id`$. This case follows immediately from Case 1 and Assumption 1.5(iii). $`\mathrm{}`$
Recall that $`\text{End}(V)`$ is the $`𝕂`$-algebra consisting of all linear transformations from $`V`$ to $`V`$.
###### Definition 5.10
With reference to Definition 4.1, let $`𝒟`$ denote the $`𝕂`$-subalgebra of $`\text{End}(V)`$ generated by $`A`$.
We will be concerned with the following subspace of $`V`$. With reference to Definition 5.7 and Definition 5.10, for $`0id/2`$ we define $`𝒟H_i=\mathrm{span}\{Xh|X𝒟,hH_i\}`$.
###### Lemma 5.11
With reference to Assumption 1.5, Definition 5.7, and Definition 5.10,
$`𝒟H_i={\displaystyle \underset{j=0}{\overset{d2i}{}}}R^jH_i(directsum),0id/2.`$ (14)
Proof: Let $`i`$ be given. Define $`\mathrm{\Delta }=_{j=0}^{d2i}R^jH_i`$. We first show $`𝒟H_i=\mathrm{\Delta }`$. Recall $`H_iU_i`$ by Definition 5.7. By this and Lemma 4.2(i) we find $`\mathrm{\Delta }𝒟H_i`$. We now show $`𝒟H_i\mathrm{\Delta }`$. Since $`𝒟`$ is generated by $`A`$ and since $`\mathrm{\Delta }`$ contians $`H_i`$ it suffices to show that $`\mathrm{\Delta }`$ is $`A`$-invariant. We now show $`\mathrm{\Delta }`$ is $`A`$-invariant. For $`0jd2i`$ and $`hH_i`$ we show $`AR^jh\mathrm{\Delta }`$. Using Assumption 1.5(i) and Lemma 4.2(i) we find $`AR^jhR^jH_i+R^{j+1}H_i`$. Recall $`R^{d2i+1}H_i=0`$ by Definition 5.7. By these comments we find $`AR^jh\mathrm{\Delta }`$. This completes the proof that $`\mathrm{\Delta }`$ is $`A`$-invariant and it follows $`𝒟H_i\mathrm{\Delta }`$. We have now shown $`𝒟H_i=\mathrm{\Delta }`$. It remains to show that the sum $`_{j=0}^{d2i}R^jH_i`$ is direct. This follows since $`U_0,U_1,\mathrm{},U_d`$ is a decomposition of $`V`$ and since $`R^jH_iU_{i+j}`$ $`(0jd2i)`$ by Assumption 1.5(i). $`\mathrm{}`$
###### Lemma 5.12
With reference to Assumption 1.5 and Definition 5.10,
$`V={\displaystyle \underset{i=0}{\overset{d/2}{}}}𝒟H_i(directsum).`$
Proof: By Lemma 5.9 and since $`U_0,U_1,\mathrm{},U_d`$ is a decomposition of $`V`$,
$`V={\displaystyle \underset{i=0}{\overset{d}{}}}{\displaystyle \underset{j=0}{\overset{min(i,di)}{}}}R^{ij}H_j(directsum).`$
In this sum we interchange the order of summation and find
$`V={\displaystyle \underset{i=0}{\overset{d/2}{}}}{\displaystyle \underset{j=0}{\overset{d2i}{}}}R^jH_i(directsum).`$
The result now follows by Lemma 5.11. $`\mathrm{}`$
###### Lemma 5.13
With reference to Definition 5.1,
$`V={\displaystyle \underset{i=0}{\overset{d}{}}}W_i.`$
Proof: Define $`V^{}=_{i=0}^dW_i`$. We show $`V=V^{}`$. By construction $`V^{}V`$. We now show $`VV^{}`$. For $`0id`$ we set $`j=di`$ in Lemma 5.4(i) and find $`(Aq^{d2i}I)W_iW_{i+1}`$. Using this we find $`AV^{}V^{}`$. By this and Definition 5.10 we find $`𝒟V^{}V^{}`$. By Definition 5.1 and Lemma 5.8 we find $`H_jW_j`$ for $`0jd/2`$. Therefore $`H_jV^{}`$ for $`0jd/2`$. By these comments we find $`𝒟H_jV^{}`$ for $`0jd/2`$. Now $`VV^{}`$ in view of Lemma 5.12. We have now shown $`V=V^{}`$. $`\mathrm{}`$
###### Corollary 5.14
With reference to Assumption 1.5, Definition 4.4, and Definition 5.1, the following (i)–(iii) hold.
1. $`W_0+\mathrm{}+W_i=U_0+\mathrm{}+U_i`$, $`0id`$,
2. $`W_i+\mathrm{}+W_d=V_0+\mathrm{}+V_{di}`$, $`0id`$,
3. $`dim(W_i)=\rho _i`$, $`0id`$.
Proof: (i) Let $`i`$ be given. Define $`\mathrm{\Delta }=W_0+\mathrm{}+W_i`$ and $`\mathrm{\Gamma }=U_0+\mathrm{}+U_i`$. We show $`\mathrm{\Delta }=\mathrm{\Gamma }`$. By Definition 5.1 we have $`\mathrm{\Delta }\mathrm{\Gamma }`$. Thus it suffices to show dim$`(\mathrm{\Delta })`$ $`=`$ dim$`(\mathrm{\Gamma })`$. By construction dim$`(\mathrm{\Delta })`$ $``$ dim$`(\mathrm{\Gamma })`$. Suppose, towards a contradiction, that dim$`(\mathrm{\Delta })`$ $`<`$ dim$`(\mathrm{\Gamma })`$. Using Definition 4.6, Lemma 5.6 and since $`U_0,U_1,\mathrm{},U_d`$ is a decomposition of $`V`$ we find
$`{\displaystyle \underset{h=0}{\overset{i}{}}}\text{dim}(W_h)<{\displaystyle \underset{h=0}{\overset{i}{}}}\rho _h.`$ (15)
By Definition 5.1 we have $`_{h=i+1}^dW_h_{h=0}^{di1}V_h`$. Using Definition 4.6, Lemma 4.7(iii), Lemma 5.6 and since $`V_0,V_1,\mathrm{},V_d`$ is a decomposition of $`V`$ we find
$`{\displaystyle \underset{h=i+1}{\overset{d}{}}}\text{dim}(W_h){\displaystyle \underset{h=i+1}{\overset{d}{}}}\rho _h.`$ (16)
By Lemma 5.6 and Lemma 5.13 we find
$`\text{dim}(V)={\displaystyle \underset{h=0}{\overset{d}{}}}\text{dim}(W_h).`$ (17)
Adding (15)–(17) we find $`\text{dim}(V)`$ $`<`$ $`_{h=0}^d\rho _h`$. This contradicts Lemma 4.7(ii) and the result follows.
(ii) Similar to (i).
(iii) By (i), Lemma 5.6, and since $`U_0,U_1,\mathrm{},U_d`$ is a decomposition of $`V`$ we find $`_{h=0}^i\text{dim}(W_h)=_{h=0}^i\rho _h`$ for $`0id`$. The result follows. $`\mathrm{}`$
###### Corollary 5.15
With reference to Definition 5.1,
$`W_i0,0id.`$
Proof: Immediate from Lemma 4.7(i) and Corollary 5.14(iii). $`\mathrm{}`$
Combining Lemma 5.6, Lemma 5.13, and Corollary 5.15 we obtain Theorem 5.2.
## 6 Interchanging $`R`$ and $`L`$
In this section we use Note 2.1 to obtain results that are analogous to the results in Sections 5 and 6.
###### Definition 6.1
With reference to Assumption 1.5, let $`A^{}:VV`$ denote the following linear transformation:
$`A^{}=K^1+L.`$
###### Lemma 6.2
With reference to Definition 6.1 and Assumption 1.5, the following holds. $`A^{}`$ is diagonalizable on $`V`$ and the set of distinct eigenvalues of $`A^{}`$ is $`\{q^{2id}|\mathrm{\hspace{0.17em}0}id\}`$. Moreover, for $`0id`$, the dimension of the eigenspace for $`A^{}`$ associated with $`q^{d2i}`$ is equal to the dimension of $`U_i`$.
Proof: Apply Note 2.1 to Lemma 4.3. $`\mathrm{}`$
###### Definition 6.3
With reference to Definition 6.1 and Lemma 6.2, for $`0id`$ we let $`V_i^{}`$ denote the eigenspace for $`A^{}`$ with eigenvalue $`q^{d2i}`$. For notational convenience we set $`V_1^{}:=0,V_{d+1}^{}:=0`$. We observe that $`V_0^{},V_1^{},\mathrm{},V_d^{}`$ is a decomposition of $`V`$.
###### Definition 6.4
With reference to Assumption 1.5 and Definition 6.3, we define
$`W_i^{}=(U_i+\mathrm{}+U_d)(V_{di}^{}+\mathrm{}+V_d^{}),0id.`$
For notational convenience we set $`W_1^{}:=0,W_{d+1}^{}:=0`$.
###### Theorem 6.5
With reference to Definition 6.4, the sequence $`W_0^{},W_1^{},\mathrm{},W_d^{}`$ is a decomposition of $`V`$.
Proof: Apply Note 2.1 to Theorem 5.2. $`\mathrm{}`$
## 7 The linear transformations $`B`$ and $`B^{}`$
In this section we introduce the linear transformations $`B`$, $`B^{}`$ and present a number of relations involving $`A`$, $`A^{}`$, $`B`$, $`B^{}`$, $`K`$, $`K^1`$.
###### Definition 7.1
With reference to Definition 5.1 and Definition 6.4, we define the following linear transformations.
1. Let $`B:VV`$ be the linear transformation such that for $`0id`$, $`W_i`$ is an eigenspace for $`B`$ with eigenvalue $`q^{2id}`$.
2. Let $`B^{}:VV`$ be the linear transformation such that for $`0id`$, $`W_i^{}`$ is an eigenspace for $`B^{}`$ with eigenvalue $`q^{d2i}`$.
###### Lemma 7.2
\[3, Lemma 7.2\] With reference to Definition 4.1, Definition 6.1, and Definition 7.1,
$`{\displaystyle \frac{qABq^1BA}{qq^1}}=I,`$ (18)
$`{\displaystyle \frac{qA^{}B^{}q^1B^{}A^{}}{qq^1}}=I,`$ (19)
$`{\displaystyle \frac{qBA^{}q^1A^{}B}{qq^1}}=I,`$ (20)
$`{\displaystyle \frac{qB^{}Aq^1AB^{}}{qq^1}}=I.`$ (21)
###### Lemma 7.3
\[3, Lemma 9.1\] With reference to Assumption 1.5 and Definition 7.1,
$`{\displaystyle \frac{qBK^1q^1K^1B}{qq^1}}=I,`$ (22)
$`{\displaystyle \frac{qB^{}Kq^1KB^{}}{qq^1}}=I.`$ (23)
###### Lemma 7.4
\[3, Lemma 10.1\] With reference to Definition 7.1, the following (i),(ii) hold.
1. $`B^3B^{}[3]B^2B^{}B+[3]BB^{}B^2B^{}B^3=0`$,
2. $`B^3B[3]B^2BB^{}+[3]B^{}BB^2BB^3=0`$.
## 8 The proof of Theorem 1.6 (existence)
In this section we prove the existence part of Theorem 1.6.
###### Definition 8.1
With reference to Assumption 1.5 and Definition 7.1, let $`r:VV`$ and $`l:VV`$ denote the following linear transformations:
$`r={\displaystyle \frac{IKB^{}}{q(qq^1)^2}},l={\displaystyle \frac{IK^1B}{q(qq^1)^2}}.`$
###### Lemma 8.2
With reference to Definition 8.1, the following (i),(ii) hold.
1. $`B=Kq(qq^1)^2Kl`$,
2. $`B^{}=K^1q(qq^1)^2K^1r`$.
Proof: Immediate from Definition 8.1. $`\mathrm{}`$
###### Theorem 8.3
With reference to Assumption 1.5 and Definition 8.1, the following (i)–(ix) hold.
1. $`KK^1=K^1K=I`$,
2. $`KR=q^2RK`$, $`KL=q^2LK`$,
3. $`Kr=q^2rK`$, $`Kl=q^2lK`$,
4. $`rR=Rr`$, $`lL=Ll`$,
5. $`lRRl=\frac{K^1K}{qq^1}`$, $`rLLr=\frac{KK^1}{qq^1}`$,
6. $`R^3L[3]R^2LR+[3]RLR^2LR^3=0`$,
7. $`L^3R[3]L^2RL+[3]LRL^2RL^3=0`$,
8. $`r^3l[3]r^2lr+[3]rlr^2lr^3=0`$,
9. $`l^3r[3]l^2rl+[3]lrl^2rl^3=0`$.
Proof: (i) Immediate from Note 1.4.
(ii) These equations hold by Lemma 2.3.
(iii) Evaluate Lemma 7.3 using Lemma 8.2.
(iv) Evaluate (20),(21) using Definition 4.1, Definition 6.1, and Lemma 8.2 and simplify the result using (ii),(iii) above.
(v) Evaluate (18),(19) using Definition 4.1, Definition 6.1, and Lemma 8.2 and simplify the result using (ii),(iii) above.
(vi),(vii) These relations hold by Assumption 1.5(v),(vi).
(viii), (ix) Evaluate Lemma 7.4 using Lemma 8.2 and simplify the result using (iii) above. $`\mathrm{}`$
###### Theorem 8.4
With reference to Assumption 1.5 and Definition 8.1, $`V`$ supports a $`U_q(\widehat{𝔰𝔩}_2)`$-module structure for which the Chevalley generators act as follows:
generator $`e_0^{}`$ $`e_1^{}`$ $`e_0^+`$ $`e_1^+`$ $`K_0`$ $`K_1`$ $`K_0^1`$ $`K_1^1`$ action on $`V`$ $`L`$ $`R`$ $`r`$ $`l`$ $`K`$ $`K^1`$ $`K^1`$ $`K`$
Proof: To see that the above action on $`V`$ determines a $`U_q(\widehat{𝔰𝔩}_2)`$-module, compare the equations in Theorem 8.3 with the defining relations for $`U_q(\widehat{𝔰𝔩}_2)`$ in Definition 1.1. $`\mathrm{}`$
Proof of Theorem 1.6 (existence): The existence part of Theorem 1.6 is immediate from Theorem 8.4. $`\mathrm{}`$
Note that the $`U_q(\widehat{𝔰𝔩}_2)`$-module structure given by Theorem 1.6 is not necessarily irreducible.
## 9 The proof of Theorem 1.6 (uniqueness)
In this section we prove the uniqueness part of Theorem 1.6.
The quantum algebra $`U_q(sl_2)`$ and its finite dimensional modules will be useful in proving uniqueness. We begin by recalling the definition of $`U_q(sl_2)`$.
###### Definition 9.1
\[11, Definition 1.1\] The quantum algebra $`U_q(sl_2)`$ is the unital associative $`𝕂`$-algebra with generators $`k,k^1,e,f`$ which satisfy the following relations:
$`kk^1=k^1k=1,`$
$`ke=q^2ek,`$
$`kf=q^2fk,`$
$`effe={\displaystyle \frac{kk^1}{qq^1}}.`$
We now recall the finite dimensional irreducible $`U_q(sl_2)`$-modules.
###### Lemma 9.2
\[11, Theorem 2.6\] With reference to Definition 9.1, there exists a family
$`V_{ϵ,d},ϵ\{1,1\},d=0,1,2,\mathrm{}`$
of finite dimensional irreducible $`U_q(sl_2)`$-modules with the following properties. The module $`V_{ϵ,d}`$ has a basis $`u_0,u_1,\mathrm{},u_d`$ satisfying:
$`ku_i=ϵq^{d2i}u_i,0id,`$ (24)
$`fu_i=[i+1]u_{i+1},0id1,fu_d=0,`$ (25)
$`eu_i=ϵ[di+1]u_{i1},1id,eu_0=0.`$ (26)
Moreover, every finite dimensional irreducible $`U_q(sl_2)`$-module is isomorphic to exactly one of the modules $`V_{ϵ,d}`$.
###### Remark 9.3
If the characteristic Char$`(𝕂)=2`$ then in Lemma 9.2 we view $`\{1,1\}`$ as having a single element.
###### Lemma 9.4
\[11, Proposition 2.3\] Let $`V`$ denote a finite dimensional $`U_q(sl_2)`$-module. If Char$`(𝕂)2`$ then the action of $`k`$ on $`V`$ is diagonalizable.
###### Remark 9.5
\[11, p. 19\] Assume Char$`(𝕂)=2`$. We display a finite dimensional $`U_q(sl_2)`$-module on which the action of $`k`$ is not diagonalizable. Let $`k,k^1,e,f`$ denote the generators for $`U_q(sl_2)`$ as in Definition 9.1. Let $`k`$ and $`k^1`$ act on the vector space $`𝕂^2`$ as $`\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$ and let $`e`$ and $`f`$ act on $`𝕂^2`$ as $`0`$. Then $`𝕂^2`$ is a finite dimensional $`U_q(sl_2)`$-module and the action of $`k`$ on $`𝕂^2`$ is not diagonalizable.
###### Lemma 9.6
\[11, Theorem 2.9\] Let $`V`$ denote a finite dimensional $`U_q(sl_2)`$-module. If the action of $`k`$ on $`V`$ is diagonalizable then $`V`$ is the direct sum of irreducible $`U_q(sl_2)`$-submodules.
###### Lemma 9.7
Let $`V`$ be a finite dimensional $`U_q(sl_2)`$-module. Assume the action of $`k`$ on $`V`$ is diagonalizable. For $`ϵ\{1,1\}`$ and for an integer $`d0`$ let $`vV`$ denote an eigenvector for $`k`$ with eigenvalue $`ϵq^d`$. Then $`ev=0`$ if and only if $`f^{d+1}v=0`$ .
Proof: Immediate from Lemma 9.2 and Lemma 9.6. $`\mathrm{}`$
###### Lemma 9.8
Let $`V`$ denote a finite dimensional vector space over $`𝕂`$. Suppose there are two $`U_q(sl_2)`$-module structures on $`V`$. Assume the actions of $`k`$ on $`V`$ given by the two module structures agree and are diagonalizable. Assume the actions of $`f`$ on $`V`$ given by the two module structures agree. Then the actions of $`e`$ on $`V`$ given by the two module structures agree.
Proof: Let $`E:VV`$ (resp. $`E^{}:VV`$) denote the action of $`e`$ on $`V`$ given by the first (resp. second) module structure. We show $`(EE^{})V=0`$. Using Lemma 9.6 and refering to the first module structure we find $`V`$ is the direct sum of irreducible $`U_q(sl_2)`$-submodules. Let $`W`$ be one the irreducible submodules in this sum. It suffices to show $`(EE^{})W=0`$. By Lemma 9.2, there exists a nonnegative integer $`d`$ and $`ϵ\{1,1\}`$ such that $`W`$ is isomorphic to $`V_{ϵ,d}`$. Therefore, the eigenvalues for $`k`$ on $`W`$ are $`ϵq^{d2i}`$ $`(0id)`$, and dim$`(W)=d+1`$. Let $`uW`$ be an eigenvector for $`k`$ with eigenvalue $`ϵq^d`$. By Lemma 9.2, the vectors $`u,fu,\mathrm{},f^du`$ are a basis for $`W`$. We show $`(EE^{})f^iu=0`$ for $`0id`$. First assume $`i=0`$. Using Lemma 9.2 we find $`Eu=0`$. Also by Lemma 9.2 we find $`f^{d+1}u=0`$ and so $`E^{}u=0`$ by Lemma 9.7. We have now shown $`(EE^{})u=0`$. Next let $`i1`$. By induction on $`i`$ we may assume
$`(EE^{})f^{i1}u=0.`$ (27)
Using Definition 9.1 and since the actions of $`k`$ (resp. $`f`$) on $`V`$ given by the two module structures agree we find $`EffE=E^{}ffE^{}`$. Hence
$`f(EE^{})=(EE^{})f.`$ (28)
Applying $`f`$ to (27) and using (28) we find $`(EE^{})f^iu=0`$. This shows $`(EE^{})W=0`$ and the result follows. $`\mathrm{}`$
The following lemma relates $`U_q(\widehat{𝔰𝔩}_2)`$-modules to $`U_q(sl_2)`$-modules.
###### Lemma 9.9
Let $`V`$ denote a finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-module. Then for $`i\{0,1\}`$, $`V`$ supports a $`U_q(sl_2)`$-module structure such that $`K_ik`$, $`e_i^+e`$, and $`e_i^{}f`$ vanish on $`V`$.
Proof: Immediate from Definition 1.1 and Definition 9.1. $`\mathrm{}`$
Proof of Theorem 1.6 (uniqueness): Suppose there exist two $`U_q(\widehat{𝔰𝔩}_2)`$-module structures on $`V`$ satisfying the conditions of Theorem 1.6. We show that these two module structures agree. By construction the actions of $`e_0^{}`$ (resp. $`e_1^{},K_0,K_1`$) on $`V`$ given by the two $`U_q(\widehat{𝔰𝔩}_2)`$-module structures agree. We now show that the actions of $`e_0^+`$ on $`V`$ given by the two $`U_q(\widehat{𝔰𝔩}_2)`$-module structures agree. Note that the actions of $`K_0`$ on $`V`$ given by the two $`U_q(\widehat{𝔰𝔩}_2)`$-module structures are diagonalizable. Using Lemma 9.9, $`V`$ supports two $`U_q(sl_2)`$-module structures given by the two actions of $`K_0,e_0^+,e_0^{}`$ on $`V`$. Using Lemma 9.8, we find that the actions of $`e_0^+`$ on $`V`$ given by the two $`U_q(\widehat{𝔰𝔩}_2)`$-module structures agree. Similarly, the actions of $`e_1^+`$ on $`V`$ given by the two $`U_q(\widehat{𝔰𝔩}_2)`$-module structures agree. We have now shown the two $`U_q(\widehat{𝔰𝔩}_2)`$-module structures of $`V`$ agree. $`\mathrm{}`$
## 10 Which $`U_q(\widehat{𝔰𝔩}_2)`$-modules arise from Theorem 1.6?
Theorem 1.6 gives a way to construct finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-modules. Not all finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-modules arise from this construction; in this section we determine which ones do.
###### Definition 10.1
Let $`V`$ denote a finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-module. Let $`d`$ denote a nonnegative integer. We say $`V`$ is basic of diameter $`d`$ whenever there exists a decomposition $`U_0,U_1,\mathrm{},U_d`$ of $`V`$ and linear transformations $`R:VV`$ and $`L:VV`$ satisfying Assumption 1.5 such that the given $`U_q(\widehat{𝔰𝔩}_2)`$-module structure on $`V`$ agrees with the $`U_q(\widehat{𝔰𝔩}_2)`$-module structure on $`V`$ given by Theorem 1.6.
Our goal for this section is to determine which $`U_q(\widehat{𝔰𝔩}_2)`$-modules are basic. We begin with a lemma.
###### Lemma 10.2
Let $`V`$ denote a finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-module. If Char$`(𝕂)2`$ then the actions of $`K_0`$ and $`K_1`$ on $`V`$ are diagonalizable.
Proof: For $`i\{0,1\}`$, view $`V`$ as a $`U_q(sl_2)`$-module under the action of $`K_i,e_i^+,e_i^{}`$ as in Lemma 9.9. The result now follows immediately by Lemma 9.4. $`\mathrm{}`$
###### Remark 10.3
Assume Char$`(𝕂)=2`$. We display a finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-module on which the actions of $`K_0`$ and $`K_1`$ are not diagonalizable. Let $`e_i^\pm ,K_i^{\pm 1}`$, $`i\{0,1\}`$, denote the Chevalley generators for $`U_q(\widehat{𝔰𝔩}_2)`$ as in Definition 1.1. Let $`K_0^{\pm 1}`$ and $`K_1^{\pm 1}`$ act on the vector space $`𝕂^2`$ as $`\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$ and let $`e_0^\pm `$ and $`e_1^\pm `$ act on $`𝕂^2`$ as $`0`$. Then $`𝕂^2`$ is a finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-module and the actions of $`K_0`$ and $`K_1`$ on $`𝕂^2`$ are not diagonalizable.
###### Theorem 10.4
Let $`d`$ denote a nonnegative integer and let $`V`$ denote a finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-module. With reference to Definition 10.1, the following are equivalent.
1. $`V`$ is basic of diameter $`d`$.
2. $`(K_0K_1I)V=0`$, the action of $`K_0`$ on $`V`$ is diagonalizable, and the set of distinct eigenvalues for $`K_0`$ on $`V`$ is $`\{q^{2id}|\mathrm{\hspace{0.17em}0}id\}`$.
Proof: (i) $``$ (ii): Let $`U_0,U_1,\mathrm{},U_d`$ be the decomposition of $`V`$ from Definition 10.1. Let $`K:VV`$ be the linear transformation that corresponds to this decomposition as in Definition 1.3. By Definition 10.1 and Theorem 1.6, we find $`KK_0`$ and $`K^1K_1`$ vanish on $`V`$. The result follows.
(ii) $``$ (i): For $`0id`$ let $`U_iV`$ be the eigenspace for the action of $`K_0`$ on $`V`$ with eigenvalue $`q^{2id}`$. Note that the sequence $`U_0,U_1,\mathrm{},U_d`$ is a decomposition of $`V`$. Define linear transformations $`R:VV`$ and $`L:VV`$ by $`(Re_1^{})V=0`$ and $`(Le_0^{})V=0`$. We now check that $`R`$ and $`L`$ satisfy Assumption 1.5(i)–(vi). Using (3) and (4), it is routine to check that $`R`$ and $`L`$ satisfy Assumption 1.5(i),(ii). Next we verify that $`R`$ satisfies Assumption 1.5(iii). View $`V`$ as a $`U_q(sl_2)`$-module under the action of $`K_1,e_1^+,e_1^{}`$ as in Lemma 9.9. Since $`(K_0K_1I)V=0`$ and the action of $`K_0`$ on $`V`$ is diagonalizable we find that the action of $`K_1`$ on $`V`$ is diagonalizable. Thus, by Lemma 9.6, $`V`$ is a direct sum of irreducible $`U_q(sl_2)`$-submodules. Let $`W`$ be one of the irreducible submodules in this sum. Note that $`RWW`$. Using Lemma 9.2, we find that for $`0id/2`$ the restriction $`R^{d2i}|_{WU_i}:WU_iWU_{di}`$ is a bijection. It follows that for $`0id/2`$ the restriction $`R^{d2i}|_{U_i}:U_iU_{di}`$ is a bijection. We have now shown $`R`$ satisfies Assumption 1.5(iii). The proof that $`L`$ satisfies Assumption 1.5(iv) is similar. By (7) we find $`R`$ and $`L`$ satisfy Assumptions 1.5(v),(vi). We have now shown that $`R`$ and $`L`$ satisfy Assumption 1.5(i)–(vi) and so Theorem 1.6 applies. It remains to show that the given $`U_q(\widehat{𝔰𝔩}_2)`$-module structure on $`V`$ agrees with the $`U_q(\widehat{𝔰𝔩}_2)`$-module structure on $`V`$ given by Theorem 1.6. Using the uniqueness statement in Theorem 1.6 it suffices to show each of $`Re_1^{}`$, $`Le_0^{}`$, $`KK_0`$, $`K^1K_1`$ vanish on $`V`$, where $`K:VV`$ is the linear transformation that corresponds to the decomposition $`U_0,U_1,\mathrm{},U_d`$ as in Definition 1.3. The first two equations were mentioned earlier. Using Definition 1.3 we find $`(KK_0)V=0`$. Since $`(K_0K_1I)V=0`$ we find $`(K^1K_1)V=0`$. We have now shown that the given $`U_q(\widehat{𝔰𝔩}_2)`$-module structure on $`V`$ agrees with the $`U_q(\widehat{𝔰𝔩}_2)`$-module structure on $`V`$ given by Theorem 1.6. $`\mathrm{}`$
## 11 The relationship between general $`U_q(\widehat{𝔰𝔩}_2)`$-modules and basic $`U_q(\widehat{𝔰𝔩}_2)`$-modules
Throughout this section $`V`$ will denote a nonzero finite dimensional $`U_q(\widehat{𝔰𝔩}_2)`$-module (not necessarily irreducible) on which the actions of $`K_0`$ and $`K_1`$ are diagonalizable (see Lemma 10.2).
In this section we will show, roughly speaking, that $`V`$ is made up of basic $`U_q(\widehat{𝔰𝔩}_2)`$-modules. We will use the following definition.
###### Definition 11.1
For $`ϵ_0,ϵ_1\{1,1\}`$ we define
$`V_{even}^{(ϵ_0,ϵ_1)}=\text{span}\{vV|K_0v=ϵ_0q^iv,K_1v=ϵ_1q^iv,i,ieven\},`$
$`V_{odd}^{(ϵ_0,ϵ_1)}=\text{span}\{vV|K_0v=ϵ_0q^iv,K_1v=ϵ_1q^iv,i,iodd\}.`$
###### Theorem 11.2
With reference to Definition 11.1,
$`V={\displaystyle \underset{(ϵ_0,ϵ_1)}{}}{\displaystyle \underset{\sigma }{}}V_\sigma ^{(ϵ_0,ϵ_1)}(directsumofU_q(\widehat{𝔰𝔩}_2)modules),`$ (29)
where the first sum is over all ordered pairs $`(ϵ_0,ϵ_1)`$ with $`ϵ_0,ϵ_1\{1,1\}`$, and the second sum is over all $`\sigma \{even,odd\}`$.
Proof: Using (2) and Definition 11.1 we find $`V_{even}^{(ϵ_0,ϵ_1)}`$ and $`V_{odd}^{(ϵ_0,ϵ_1)}`$ are invariant under the action of $`K_0^{\pm 1}`$ and $`K_1^{\pm 1}`$. Using (3) and (4) we find $`V_{even}^{(ϵ_0,ϵ_1)}`$ and $`V_{odd}^{(ϵ_0,ϵ_1)}`$ are invariant under the action of $`e_0^\pm ,e_1^\pm `$. Thus the subspaces on the right hand side of (29) are $`U_q(\widehat{𝔰𝔩}_2)`$-submodules of $`V`$. We now show the sum on the right hand side of (29) equals $`V`$. Recall that the actions of $`K_0`$ and $`K_1`$ on $`V`$ are both diagonalizable. Using this and (2) we find that the actions of $`K_0`$ and $`K_1`$ on $`V`$ are simultaneously diagonalizable. Thus $`V`$ is the direct sum of common eigenspaces for the actions of $`K_0`$ and $`K_1`$ on $`V`$. It remains to show that any common eigenvector for the actions of $`K_0`$ and $`K_1`$ on $`V`$ is in one of the subspaces on the right hand side of (29). Let $`v`$ denote a common eigenvector for the action of $`K_0`$ and $`K_1`$ on $`V`$. By construction there exists $`\alpha ,\beta 𝕂`$ such that $`K_0v=\alpha v`$ and $`K_1v=\beta v`$. Using Lemma 9.9, Lemma 9.6, and Lemma 9.2 we find that there exists an $`ϵ_0\{1,1\}`$ and an integer $`i`$ such that $`\alpha =ϵ_0q^i`$. For every $`m`$ define $`T_m:=\{xV|K_0x=ϵ_0q^{i+2m}x`$ and $`K_1x=\beta q^{2m}x\}`$, and define $`T:=_mT_m`$. Observe $`K_0K_1ϵ_0q^i\beta I`$ vanishes on $`T`$. Using (3) and (4) we find $`T`$ is a $`U_q(\widehat{𝔰𝔩}_2)`$-module. Also, $`0vT`$ and so $`T`$ is not the zero module. Let $`W`$ denote an irreducible $`U_q(\widehat{𝔰𝔩}_2)`$-module contained in $`T`$. By \[4, Proposition 3.2\], there exists an $`ϵ\{1,1\}`$ such that $`K_0K_1ϵI`$ vanishes on $`W`$. Also, $`K_0K_1ϵ_0q^i\beta I`$ vanishes on $`W`$. So we find $`ϵ=ϵ_0q^i\beta `$. Define $`ϵ_1=ϵϵ_0^1`$. Then $`ϵ_1\{1,1\}`$ and $`\beta =ϵ_1q^i`$. We have now shown $`K_0v=ϵ_0q^iv`$ and $`K_1v=ϵ_1q^iv`$. Therefore $`vV_{even}^{(ϵ_0,ϵ_1)}`$ if $`i`$ is even or $`vV_{odd}^{(ϵ_0,ϵ_1)}`$ if $`i`$ is odd. This shows that the sum on the right hand side of (29) equals $`V`$. By Definition 11.1 the sum in (29) is direct. $`\mathrm{}`$
###### Lemma 11.3
With reference to Definition 10.1, Definition 11.1, and Theorem 11.2, the following are equivalent.
1. $`V=V_{even}^{(1,1)}`$.
2. $`V`$ is basic of even diameter.
3. The spaces $`V_{even}^{(1,1)}`$, $`V_{even}^{(1,1)}`$, $`V_{even}^{(1,1)}`$, $`V_{odd}^{(1,1)}`$, $`V_{odd}^{(1,1)}`$, $`V_{odd}^{(1,1)}`$, $`V_{odd}^{(1,1)}`$ are all zero.
Proof: (i)$``$(ii): We use Theorem 10.4. By Definition 11.1, we find $`K_0K_1I`$ vanishes on $`V`$. Recall that the action of $`K_0`$ on $`V`$ is diagonalizable. Using Lemma 9.9, Lemma 9.6 and Lemma 9.2 we find that there exists a nonnegative integer $`d`$ such that the set of distinct eigenvalues for the action of $`K_0`$ on $`V`$ is $`\{q^{2id}|\mathrm{\hspace{0.17em}0}id\}`$. So by Theorem 10.4 $`V`$ is basic of diameter $`d`$. By Definition 11.1 $`d`$ is even.
(ii)$``$(i): Immediate from Theorem 10.4 and Definition 11.1.
(i)$``$(iii): Immediate from Theorem 11.2. $`\mathrm{}`$
###### Lemma 11.4
With reference to Definition 10.1, Definition 11.1, and Theorem 11.2, the following are equivalent.
1. $`V=V_{odd}^{(1,1)}`$.
2. $`V`$ is basic of odd diameter.
3. The spaces $`V_{odd}^{(1,1)}`$, $`V_{odd}^{(1,1)}`$, $`V_{odd}^{(1,1)}`$, $`V_{even}^{(1,1)}`$, $`V_{even}^{(1,1)}`$, $`V_{even}^{(1,1)}`$, $`V_{even}^{(1,1)}`$ are all zero.
Proof: Similar to the proof of Lemma 11.3. $`\mathrm{}`$
Refering to (29), even though the six submodules $`V_{even}^{(1,1)}`$, $`V_{even}^{(1,1)}`$, $`V_{even}^{(1,1)}`$, $`V_{odd}^{(1,1)}`$, $`V_{odd}^{(1,1)}`$, $`V_{odd}^{(1,1)}`$ are not basic they become basic after a routine normalization. This is explained in the following lemma, definition, and remark.
###### Lemma 11.5
\[4, Prop. 3.3\] For any choice of scalars $`ϵ_0,ϵ_1`$ from $`\{1,1\}`$, there exists a $`𝕂`$-algebra automorphism of $`U_q(\widehat{𝔰𝔩}_2)`$ such that
$`K_iϵ_iK_i,e_i^+ϵ_ie_i^+,e_i^{}e_i^{},`$
for $`i\{0,1\}`$. We refer to the above automorphism as $`\tau (ϵ_0,ϵ_1)`$.
###### Definition 11.6
Let $`W`$ denote a $`U_q(\widehat{𝔰𝔩}_2)`$-module. Let $`\tau `$ be an automorphism of $`U_q(\widehat{𝔰𝔩}_2)`$. We define a new $`U_q(\widehat{𝔰𝔩}_2)`$-module structure on $`W`$ as follows. For $`xU_q(\widehat{𝔰𝔩}_2)`$ and $`wW`$ define $`x.w(\mathrm{new}\mathrm{action})=\tau (x).w(\mathrm{original}\mathrm{action})`$. We refer to this new $`U_q(\widehat{𝔰𝔩}_2)`$-module structure as $`W`$ twisted via $`\tau `$.
###### Remark 11.7
With reference to (29) each submodule $`V_\sigma ^{(ϵ_0,ϵ_1)}`$ becomes basic upon twisting via $`\tau (ϵ_0,ϵ_1)`$.
## 12 Acknowledgment
I would like to express my gratitude to my thesis advisor Paul Terwilliger for introducing me to this subject and for his many useful suggestions.
Darren Funk-Neubauer Department of Mathematics University of Wisconsin-Madison 480 Lincoln Drive Madison, WI 53706-1388 USA
email: [email protected]
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# Can the rotation of the dark matter halo of our galaxy be detected through its effect on the cosmic microwave background polarisation?
## I Introduction
Galaxies and clusters of galaxies are believed to be embedded in dark matter halos whose masses are significantly larger than those of the visible baryonic components, namely, stars and gas. Determining the nature of the dark matter is currently one of the most challenging problems in astrophysics dark . The dark matter halo makes its presence felt only through its gravitational field which, therefore, merits particular attention. Several proposals on modeling the dark matter halo can be found in halogravf ; partdark .
The dark matter halos are predicted to have angular momentum. In the hierarchical structure formation scenario the halos acquire angular mementum through tidal interactions of the halos with their surroundings peebles . The angular momentum is usually quantified using the dimensionless spin parameter $`\lambda =JE^{1/2}/GM^{5/2}`$ where $`E`$ is the total energy of the halo and $`M`$ is its mass. The spin parameter $`\lambda `$ of a halo is essentially the ratio of its angular momentum to that needed for rotational support. Simulations show $`\lambda `$ to be in the range $`0.020.1`$ vitvitska .
In Newtonian gravity the gravitational field of an object does not depend on whether the object is rotating or not, and the angular momentum of a dark matter halo does not manifest itself in its gravitational field. Interestingly, this is not true in Einstein’s theory of gravity where the gravitational field of a rotating object is different from that of a static one. The angular momentum manifests itself itself through the gravomagnetic effect gravomag , which is analogous to the magnetic field produced by a rotating charge.
The gravomagnetic effect is possibly the only signature of the angular momentum of a dark matter halo, and any detectable consequence of this effect holds the only hope of measuring the angular momentum of dark matter halos. In this paper, for a rotating dark matter halo, we calculate the gravomagnetic field both inside and outside the halo. The gravitational field inside the halos of galaxies and clusters of galaxies is weak, and we adopt the weak field limit of Einstein’s theory in our calculation.
The Faraday effect is a well known manifestation of magnetic fields on astronomical scales, and the gravo-magnetic field produced by rotating dark matter halos will produce an analogous gravitational Faraday effect gravofarad . This will cause the plane of polarization of linearly polarized electro-magnetic radiation to rotate as it propagates through the gravitational field of a rotating dark-matter halo. The gravitational Faraday rotation of the dark matter halo of our Galaxy will be present in all astronomical observations carried out from our location inside the halo. In particular, the gravitational Faraday effect will effect the observed polarisation pattern of the Cosmic Microwave Background Radiation (CMBR). In this paper, we calculate the expected signature of the gravitational Faraday effect of the dark matter halo of our Galaxy and investigate if this will be detectable in the large angular scale CMBR polarisation pattern cmbrpol .
A brief outline of the paper follows. In Section 2 we calculate the metric for a rotating dark matter halo and use this to determine the gravo-magnetic field both inside and outside the halo. In Section 3 we calculate the Faraday effect that will be produced by the gravo-magnetic field of the halo and estimate its effect on the CMBR polarisation pattern as seen by an observer inside the halo. In Section 4 we present results and conclusions.
It should be noted that throughout this paper we have restricted our analysis to the dark matter halo of a spiral galaxy. A similar analysis can also be applied to clusters of galaxies cluster without major modifications, but this has not been considered here.
## II The Gravomagnetic Field
The gravitational field associated with galaxy dark matter halos is weak and it is justified to describe it using a metric
$$g_{\mu \nu }=\eta _{\mu \nu }+h_{\mu \nu }$$
(1)
which deviates only slightly from the flat space-time Minkowski metric $`\eta _{\mu \nu }`$ ie. $`h_\nu ^\mu 1`$. It may be noted that we adopt the signature $`(,+,+,+)`$ for the metric and the weak field (linearised) Einstein equations in the Lorentz gauge shutz are given by
$$\mathrm{}H^{\alpha \beta }=\frac{16\pi G}{c^4}T^{\alpha \beta }$$
(2)
where
$$H^{\alpha \beta }=h^{\alpha \beta }\frac{h}{2}\eta ^{\alpha \beta }.$$
(3)
To begin with, we ignore the rotation and consider a static, spherically symmetric halo. The most general static, spherically symmetric spacetime has metric coefficients of the form $`h_{00}=2\mathrm{\Phi }(r)`$, $`h_{i0}=0`$ and $`h_{ij}=2[\mathrm{\Phi }(r)+\mathrm{\Psi }(r)]\delta _{ij}`$ where $`\mathrm{\Phi }(r)`$ and $`\mathrm{\Psi }(r)`$ are two arbitrary functions of the radial distance $`r`$. Further, if we assume that the dark matter does not have relativistic pressure (ie. $`T_{00}=\rho c^2`$ is the only non-zero component of the energy momentum tensor) it then follows from eq. (2) that $`\mathrm{\Psi }(r)=0`$. The metric now contains only one unknown function $`\mathrm{\Phi }(r)`$ which is the usual Newtonian gravitational potential divided by $`c^2`$. For spiral galaxies $`\mathrm{\Phi }(r)`$ can be determined from HI rotation curves. Assuming the HI to be rotating in circular orbits, the geodesic equation reduces to
$$\frac{d\mathrm{\Phi }}{dr}=\frac{1}{r}\frac{v_c^2(r)}{c^2}$$
(4)
where $`v_c(r)`$ is the velocity of the HI cloud at a distance $`r`$ fromt the center of the galaxy. The outer parts of spiral galaxies usually show a flat rotation curve ie. $`v_c`$ is a constant rc . For this, the potential, up to an arbitrary constant of integration, is
$$\mathrm{\Phi }(r)=\frac{v^2}{c^2}\mathrm{ln}r.$$
(5)
Using this in the Einstein’s equation (2) gives the total density in the outer parts of spiral galaxies to be $`\rho (r)=v^2/4\pi Gr^2`$. This density is usually significantly larger than the density of the visible matter and this is the dark matter problem. The density profile in the inner parts of galaxies is not very well determined and it is still an issue of debate. For the purposes of this paper we assume the density profile
$$\rho (r)=\frac{v^2}{4\pi G(r^2+r_0^2)}$$
(6)
for the dark matter halo throughout the galaxy. Here $`r_0`$ is the core radius and it is assumed that the Dark Matter halo is of size $`R`$ ie. $`\rho (r)=0`$ for $`r>R`$.
We now incorporate the rotation of the halo. Retaining the spherical density profile given by eq. (6) we assume that the halo is made up of spherical shells, each shell executing rigid rotation with angular velocity $`\stackrel{}{\omega }(r)`$ whose magnitude and direction can vary from shell to shell. It is also assumed that the rotational velocities $`\stackrel{}{v}=\stackrel{}{\omega }\times 𝐫`$ are non-relativistic ($`\stackrel{}{v}/c1`$). The change introduced by the rotation is that we now have a non-zero time-space component of the energy-momentum tensor $`T^{0i}=\rho cv^i`$. This results in non-zero values for the metric coefficients $`𝐡=h_{oi}`$ whose values have to be determined through
$$^2𝐡(𝐫)=\frac{16\pi G}{c^3}\rho (r)(\stackrel{}{\omega }\times 𝐫).$$
(7)
The solution for $`𝐡`$ is given by
$$𝐡(𝐬)=\frac{4G}{c^3}\frac{\rho (r)(\stackrel{}{\omega }(𝐫)\times 𝐫)}{|𝐬𝐫|}𝑑\tau $$
(8)
The integral in eq. (8) can be simplified and written as a sum of two parts as
$$𝐡(𝐬)=\frac{16\pi G}{3c^3}\left(\frac{1}{s^2}_0^sr^4\rho (r)\stackrel{}{\omega }(r)𝑑r+s_s^Rr\rho (r)\stackrel{}{\omega }(r)\right)\times \frac{𝐬}{s}$$
(9)
where the first integral is the contributionfrom the region interior to the point $`𝐬`$ where $`𝐡`$ is being calculated and the second integral is for the exterior region.
The angular momentum of a shell of radius $`r`$ and thickness $`dr`$ is
$$d𝐋(r)=\frac{8\pi }{3}\stackrel{}{\omega }(r)\rho (r)r^4dr$$
(10)
Using this, we can write eq. (8) as
$$𝐡(𝐬)=\frac{2G}{c^3}\left(\frac{𝐋_s}{s^2}+s_s^R\frac{d𝐋(r)}{r^3}\right)\times \frac{𝐬}{s}$$
(11)
Here $`𝐋_s`$ denotes the total angular momentum of all the shells interior to $`𝐬`$. The first term in eq. (11) for $`𝐡(𝐬)`$ is the contribution from the total angular momentum interior to $`𝐬`$ and the term involving the integral is the contribution from the angular momentum of the shells exterior to $`𝐬`$.
It is worth noting that only the first term contributes outside the halo and relates $`𝐡(𝐬)`$ to the angular momentum of the whole halo
$$𝐡(𝐬)=2\frac{𝐋\times 𝐬}{s^3}$$
(12)
Calculating the geodesic equation for a non-relativistic test particle moving in the gravitation field of the halo we have
$$\frac{d𝐯}{dt}=(c^2\mathrm{\Phi })+𝐯\times (\times 𝐡c)$$
(13)
where we see that we can identify $`c𝐡`$ as the equivalent of the vector potential $`𝐀`$ in electro-magnetism, and we can identify $`𝐁_g=(\times 𝐡c)`$ as the gravi-magnetic field. This will affect all motion inside the dark matter halo and it is possible that there may be a detectable effect of this force. We do not consider this possibility here but instead we focus on the gravitational Faraday effect which is a different manifestation of the gravi-magnetic field.
## III Gravitational Faraday Effect
Having obtained the gravi–magnetic field let us now see if we can predict any observable effect. It is known that the plane of polarisation of a light ray passing through a region with a gravimagnetic field may get rotated due to the gravitational Faraday effect. This is similar to the magnetooptic Faraday effect gravofarad except that the gravitational analogue is achromatic (wavelength independent). The rotation of the plane of polarisation is a consequence of the parallel transport of the polarisation vector along the path of the light ray. The resultant rotation of the plane of polarisation is given by (see paper by Sereno in gravofarad ) the expression:
$$\mathrm{\Omega }=\frac{1}{2}𝐁_g𝑑𝐥$$
(14)
where the integral is along the trajectory of the light ray. For the CMBR being observed from inside our Galaxy we consider a light ray propagating along the line of sight starting from infinity (actually the last scattering surface) and ending at the observer.
The quantity of interest from the observational point of view is the difference in the rotation angle for CMBR along different lines of sight. There will be no detectable effect if the rotation angle were the same along all lines of sight. We consider CMBR photons arriving along two different lines of sight say $`𝐧_1`$ and $`𝐧_2`$. The light rays reaching us along the two different lines of sight will traverse two different paths and the rotation angles for their planes of polarisation will be given by
$$\mathrm{\Omega }_1=\frac{1}{2}𝐁_𝐠𝑑𝐥_1$$
(15)
and
$$\mathrm{\Omega }_2=\frac{1}{2}𝐁_g𝑑𝐥_2$$
(16)
where the integrals are still from infinity to the observer but along the two different paths. Now consider a path joining the two points at infinity. This third path lies entirely at infinity where the gravimagnetic field vanishes. Thus the integral over this third path vanishes. Thus if we consider the difference of the two rotation angles $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ we can safely add to it the integral along the third path at infinity since this does affect the net result. Thus the difference of the two angles yields a closed integral of the form
$$\mathrm{\Omega }_1\mathrm{\Omega }_2=\frac{1}{2}𝐁_g𝑑𝐥$$
(17)
Applying Stokes’ theorem we obtain
$$\mathrm{\Omega }_1\mathrm{\Omega }_2=\frac{1}{2}\times 𝐁_g𝑑𝐬=\frac{c}{2}\times (\times 𝐡)𝑑𝐬$$
(18)
where the integral is over the surface enclosed by the two referred paths. Note that we are working in the Lorentz gauge where the divergence of the vector $`𝐡`$ vanishes. Also the Laplacian of this vector is given as
$$^2𝐡=\frac{16\pi G}{c^3}\rho (r)\left(\omega \times 𝐫\right)$$
(19)
Thus the expression for the difference in the two rotation angles becomes
$$\mathrm{\Omega }_1\mathrm{\Omega }_2=\frac{8\pi G}{c^2}\rho (r)(\omega \times 𝐫)𝑑𝐬=\frac{8\pi G}{c^2}𝐣𝑑𝐬$$
(20)
where the quantity $`𝐣=\rho 𝐯`$ is the matter current density.
Thus we see that the angle difference is proportional to nothing but the matter flux across the plane enclosed by the two referred paths. Only a finite portion of this surface lying within the spherical halo contributes to this.
For simplicity we assume that all the shells in the halo have a common rotation axis. Relative to the center of the halo, we shall refer to one direction along the rotation axis as the North pole and the opposite direction as the South pole, and the plane perpendicular to the rotation axis as the equator (Figure 1) . From the symmetry of the problem we expect $`𝐁_g`$ to be perpendicular to the light ray for a lines of sight along the equator. Hence we do not expect any gravitational Faraday effect along the equator. Further the effect is expected to be maximum at the poles, the direction of the rotation being opposite in the two hemispheres and the magnitude of the rotation increasing with elevation away from the equatorial plane.
We next estimate the rotation of the polarisation angle along one of the poles by calculating the matter flux through a surface shown in Figure 1. Using the density profile in eq. (6) and assuming that the all shells of the halo rotate with the same velocity as the HI in the spiral galaxy embedded in the halo we have
$$\mathrm{\Omega }_N=\pi \frac{v_c^3}{c^3}_{r_0}^R\frac{rdr}{r^2+r_0^2}$$
(21)
where $`\mathrm{\Omega }_N`$ is the rotation angle along the North pole. Calculating this we have
$$\mathrm{\Omega }_N=\pi \frac{v_c^3}{c^3}\mathrm{ln}\left(\frac{R}{r_0}\right)$$
(22)
where we have assumed $`r_0R`$.
## IV Result and Conclusion
Galaxy rotational velocities are trypically $`200\mathrm{k}\mathrm{m}/\mathrm{s}`$ for which $`(v_c/c)^33\times 10^{10}`$. The only hope for a substantial rotation is through the factor $`\mathrm{ln}(R/r_0)`$. The size of the core radius and the extent of the halo are both largely not well determined. Assuming that the core radius is of the order of a few $`\mathrm{pc}`$ and that the halo extends to a few $`\mathrm{Mpc}`$ we see that the factor $`\mathrm{ln}(R/r_0)`$ will have a value in the range $`1015`$. It then follows that the rotation will be less than $`1^{^{\prime \prime }}`$ which is too small to have a detectable signature in the CMBR.
It should be noted that the gravitational Faraday effect is of the order of $`(v_c/c)^3`$ and hence it is small for most non-relativistic situations.
The analysis carried out here depends crucially on the fact that the density profile falls as $`r^2`$ at large $`r`$, and hence the matter flux (and the rotation angle) increases only logarithmically with the size of the halo. It is possible that our conclusions would be different (and the rotation detectable) if the density profile were to be shallower.
A worthwhile future investigation will be to study gravitational lensing in the line element we have obtained for the rotating galaxy halo. The existent literature sereno seems to suggest that the deflection angle will be significantly different because of gravomagnetic effects. It remains to be seen what lies in store for us if the deflection is worked out for the metric we have derived here.
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# Untitled Document
Une généralisation du théorème de
Kobayashi-Ochiai
Frédéric Campana et Mihai Păun
0. Introduction
Soit $`X`$ une variété kählérienne compacte $`X`$, de dimension $`n`$, munie d’une métrique $`\omega `$, et une application non dégénérée $`\phi :\mathrm{}^nX`$.
Dans cet article, nous établissons des relations entre la croissance de $`\phi `$ (mesu-rée par le degré moyen ou la fonction caractéristique) et la positivité du fibré canonique $`K_X`$ (mesurée par sa pseudo-effectivité, et sa dimension numérique). Le principe général étant que la croissance de $`\phi `$ augmente avec la positivité de $`K_X`$.
Afin d’introduire notre premier résultat, on rappelle la notion de degré moyen (voir par exemple le travail de K. Kodaira ) si $`\phi :\mathrm{\Delta }_rX`$ est une application holomorphe, son dégré moyen est par définition
$$\mathrm{deg}(\phi |\mathrm{\Delta }_r):=_{\mathrm{\Delta }_r}\phi ^{}\omega ^n/_X\omega ^n.$$
Dans ce contexte, on a le théorème suivant (essentiellement optimal, par l’exemple des tores complexes):
Théorème 1. Soit $`(X,\omega )`$ une variété complexe compacte de dimension $`n`$. S’il existe une application holomorphe $`\phi :\mathrm{}^nX`$ non dégénérée, et telle que
$$\underset{r\mathrm{}}{lim}1/r^{2n}\mathrm{deg}(\phi |\mathrm{\Delta }_r)=0,$$
$`(1)`$
alors le fibré canonique de $`X`$ n’est pas pseudo-effectif.
Il convient de rappeller ici que sous les mêmes hypothèses, Kodaira montre dans que le tous les plurigenres de $`X`$ sont nuls; le théorème 1 peut être vu comme une amélioration de son résultat.
En particulier, si la variété $`X`$ est projective, cette amélioration, combinée avec le théorème principal de résolvent la conjecture d’uniréglage dans cette situation:
Corollaire. Soit $`X`$ une variété projective $`n`$-dimensionelle, et $`\phi :\mathrm{}^nX`$ une application holomorphe non dégénérée, et qui vérifie la condition (1). Alors $`X`$ est uniréglée.
On remarque que, dès que $`dim(X)2`$, il existe des exemples de variétés $`X`$ et d’applications holomorphes $`\phi :\mathrm{}^nX`$ non-dégénérées de volume fini ((Fatou-Bieberbach) dont l’image a un complémentaire contenant un ouvert non-vide.
La démonstration du théorème 1 ne fournit aucun rapport entre les courbes rationnelles (produites par la théorie de Mori) et les images des droites complexes par l’application $`\phi `$. Intuitivement, on s’attend à ce que “beaucoup” parmi les images de droites par $`\phi `$ se compactifient, mais ceci est faux en général: voir un contre-exemple à la fin du chapitre 1, dans lequel aucune courbe algébrique de $`\mathrm{}^n`$ n’a une image algébrique par $`\phi `$). Nous remercions TC Dinh et N. Sibony pour nous avoir indiqué la référence , dans laquelle on trouve une réponse positive à la question sur les domaines de Fatou-Bieberbach que nous nous posions.
Dans la théorie de Nevanlinna équidimensionelle, developpée entre autres par Griffiths, King, Stoll…, on définit la fonction caractéristique $`T_\omega (\phi ,.)`$ d’une application holomorphe (substitut du degré dans le cadre compact) par l’égalité:
$$T_\omega (\phi ,r):=_0^r\frac{dt}{t^{2n1}}_{|z|<t}\phi ^{}\omega \omega _{euc}^{n1},$$
$`\omega `$ est une métrique hermitienne sur $`X`$, et $`\omega _{euc}:=i\overline{}|z|^2`$ est la métrique canonique sur $`\mathrm{}^n`$ (voir par exemple pour une interprétation de cette quantité, et certaines de ses propriétés).
Considérons les données suivantes: $`X`$ est une variété kählérienne compacte de fibré canonique $`K_X`$ est pseudo-effectif, et $`\phi :\mathrm{}^nX`$ est une application holomorphe non-dégénérée. Le résultat suivant montre que la fonction caractéristique de $`\phi `$ croit d’autant plus vite que la dimension numérique de $`K_X`$ est grande.
Théorème 2. Soit $`(X,\omega )`$ une variété kählérienne compacte, de dimension $`n`$, de fibré canonique $`K_X`$ pseudo-effectif. Considérons une application holomorphe $`\phi :\mathrm{}^nX`$ non dégénérée. Alors il existe une constante positive $`C>0`$ telle que
$$T(\phi ,r)^{1\nu /n}Cr^2,$$
pour tout $`r>0`$, où $`\nu =\nu (K_X)`$ est la dimension numérique de $`K_X`$.
Le théorème précédent généralise ainsi le résultat classique suivant, dû à Kodaira, Griffiths, Kobayashi–Ochiai. Théorème (, , ) Soit $`X`$ une variété de dimension $`n`$, de type général. Alors pour toute application $`\phi :\mathrm{}^nX`$, on a: $`\phi ^{}(\omega ^n)=0`$, et $`\phi `$ est donc dégénérée.
Si le fibré canonique $`K_X`$ est seulement situé sur le bord du cône big, il semble beaucoup plus délicat d’analyser les conséquences du fait que la pseudo-forme volume de Kobayashi de $`X`$ soit dégénérée (évidemment, on ne peut pas espérer un énoncé analogue au théorème précedent, comme le montre le cas des tores). Néanmoins, on voudrait proposer le problème suivant. Conjecture 1. Soit $`(X,\omega )`$ une variété kählérienne compacte telle que $`K_X`$ est pseudo-effectif. Supposons que $`X`$ est revêtue par $`\mathrm{}^n`$.
Alors la dimension numérique de $`K_X`$ est zero.(Lorsque $`X`$ est projective, la condition de pseudo-effectivité est superflue, puisque $`K_X`$ est nef).
Cette conjecture généralise Kobayashi-Ochiai et implique une célèbre conjecture d’Iitaka, par des arguments standard.
Pour les courbes entières $`\phi :\mathrm{}X`$, on dispose du lemme de “reparamétrisation” de Brody: étant donnée $`\phi `$, non-constante, il existe une application $`\widehat{\phi }:\mathrm{}X`$, telle que $`sup_{\mathrm{}}|\widehat{\phi }^{}(t)|=1`$. Le procédé de Brody ne donne que des résultats partiels en plusieurs variables: si $`\phi :\mathrm{}^mX`$ est une application holomorphe, alors il existe une suite $`\phi _k:\mathrm{\Delta }_kX`$ (ici, $`\mathrm{\Delta }_k`$ désigne le polydisque de rayon égal à $`k`$) telle que $`\phi _k^{}\omega ^m{\displaystyle \frac{d\lambda }{_j(1|t_j/k|^2)^2}}`$, avec égalité à l’origine. On peut montrer que si, de plus, les images des applications $`\phi _k`$ ne s’applatissent pas lorsque $`k\mathrm{}`$, alors la dimension numérique de $`K_X`$ vaut zéro (la preuve de cette affirmation est très voisine de celle qui sera donnée pour le théorème 2, donc on la laisse au soin du lecteur).
Le dernier paragraphe de cet article concerne les applications quasi-conformes. Par analogie avec la pseudo-forme volume de Kobayashi, on introduit au troisième chapitre la notion de pseudo-forme volume quasi-conforme. Comme conséquence du théorème 2 on obtient: Corollaire. Soit $`X`$ une variété kählérienne compacte, telle que $`K_X`$ soit nef. Si la pseudo-forme volume quasi-conforme est non-dégénérée, alors $`\nu (K_X)1`$.
Ce corollaire nous a été suggeré par Y.-T. Siu, que nous remercions vivement.
Dans le même esprit, une application $`\phi :\mathrm{}^nX`$ est dite quasi-conforme en moyenne s’il existe une constante positive $`\delta `$ telle que
$$_0^r\frac{dt}{t^{2n1}}_{B(t)}|J(\phi )|_\omega ^{2/n}𝑑\lambda \delta T_\omega (\phi ,r)$$
pour tout $`r0`$; on note $`J(\phi )`$ le jacobien de l’application $`\phi `$. Pour une telle application, nous allons montrer le théorème suivant:
Théorème 3. Soit $`(X,\omega )`$ une variété kählérienne compacte de dimension $`n`$, telle que $`K_X`$ est pseudo-effectif. S’il existe une application holomorphe $`\phi :\mathrm{}^nX`$, quasi-conforme en moyenne, alors $`\nu (K_X)=0`$.
On remarquera que les théorèmes 2 et 3 etablissent des cas particuliers de la conjecture proposée ci-dessus.
1. Applications de degré moyen à croissance lente
Considérons une application holomorphe $`\phi :\mathrm{\Delta }_rX`$, où $`X`$ est une variété complexe compacte de dimension $`n`$, et $`\mathrm{\Delta }_r`$ désigne le polydisque de rayon $`r`$ dans $`\mathrm{}^n`$. On munit $`X`$ d’une métrique hermitienne $`\omega `$; la quantité:
$$\mathrm{deg}(\phi |\mathrm{\Delta }_r):=_{\mathrm{\Delta }_r}\phi ^{}\omega ^n/_X\omega ^n$$
peut être interprétée comme le degré moyen de $`\phi `$ sur ce polydisque.
Dans ce contexte, notre résultat est le suivant: Théorème 1. Soit $`(X,\omega )`$ une variété complexe compacte de dimension $`n`$. S’il existe une application holomorphe $`\phi :\mathrm{}^nX`$ non dégénérée, et telle que:
$$\underset{r\mathrm{}}{lim}1/r^{2n}\mathrm{deg}(\phi |\mathrm{\Delta }_r)=0,$$
$`(1)`$
alors le fibré canonique de $`X`$ n’est pas pseudo-effectif.
Remarquons que le produit d’un tore de dimension $`n1`$ par $`\mathrm{}^1`$ montre que l’exposant $`2n`$ est optimal comme entier pair.
Comme on l’a déjà rappelé dans l’introduction , Kodaira montre dans que sous les hypothèses du théorème 1, on a $`H^0(X,K_X^m)=0`$, pour tout $`m1`$. Maintenant si $`L`$ est un fibré en droites sur $`X`$, tel que pour un certain $`m0`$ on ait $`H^0(X,L^m)0`$, alors $`L`$ est pseudo-effectif. Ainsi, le théorème 1 peut-etre vu comme généralisation du théorème de Kodaira. D’autre part, bien qu’en général la pseudo-effectivité d’un fibré ne soit pas equivalente à l’existence de sections non-nulles pour une de ses puissances, dans le cas du fibré canonique cela devrait être vrai, d’après la conjecture d’abondance.
Preuve. Supposons que $`K_X`$ soit pseudo-effectif; alors il existe une fonction $`fL^1(X)`$, semi-continue supérieurement, telle que:
$$\mathrm{\Theta }_h(K_X)+i\overline{}f0$$
$`(2)`$
au sens des courants sur la variété $`X`$, où $`h`$ est la métrique duale de $`det(\omega )`$ sur le fibré canonique (voir , par exemple).
Prenons l’image inverse de (2) sur $`\mathrm{}^n`$ via l’application $`\phi `$; on a
$$i\overline{}(\mathrm{log}(J_\omega (\phi )^2e^{f\phi }))0$$
$`(3)`$
et alors la fonction $`\tau :\mathrm{}^n\mathrm{}_+`$ définie par $`\tau (z):=J(\phi ,z)_\omega ^2e^{f\phi (z)}`$ sera psh sur $`\mathrm{}^n`$. On rappelle maintenant la propriété de convexité suivante des fonctions psh (voir , par exemple).
Lemme. La fonction
$$M_\tau (t_1,\mathrm{},t_n):=_{[0,2\pi ]^n}\tau (t_1e^{i\theta _1},\mathrm{},t_ne^{i\theta _n})𝑑\theta $$
est croissante en chaque variable, et convexe en $`\mathrm{log}(t_j)`$.
Grâce a ce lemme, on a la suite d’inégalités
$$\begin{array}{cc}\hfill _{\mathrm{\Delta }_r}\tau (z)𝑑\lambda (z)=& _{[0,r]^n}t_1\mathrm{}t_nM_\tau (t_1,\mathrm{},t_n)𝑑t_1\mathrm{}𝑑t_n\hfill \\ \hfill & _{[r_0,r]^n}t_1\mathrm{}t_nM_\tau (t_1,\mathrm{},t_n)𝑑t_1\mathrm{}𝑑t_n\hfill \\ \hfill & M_\tau (r_0,\mathrm{},r_0)_{[r_0,r]^n}t_1\mathrm{}t_n𝑑t_1\mathrm{}𝑑t_n=\hfill \\ \hfill =& \left(\frac{r^2r_0^2}{2}\right)^nM_\tau (r_0,\mathrm{},r_0).\hfill \end{array}$$
Par ailleurs, la fonction $`f`$ est bornée supérieurement sur $`X`$, et donc
$$_{\mathrm{\Delta }_r}\tau (z)𝑑\lambda (z)C_{\mathrm{\Delta }_r}J_\omega (\phi )_z^2𝑑\lambda =_{\mathrm{\Delta }_r}\phi ^{}\omega ^n$$
La condition de croissance imposée par hypothèse sur le degré moyen de l’application $`\phi `$ montre que $`M_\tau (r_0,\mathrm{},r_0)=0`$, et ceci est vrai pour tout rayon $`r_0>0`$. Notre hypothèse permet également de changer l’origine de $`\mathrm{}^n`$, quitte à faire une translation, donc en resumé $`\tau 0`$. Mais ceci est clairement impossible, car l’image de l’application $`\phi `$ contient un ouvert de la variété $`X`$. La contradiction ainsi obtenue montre que $`K_X`$ n’est pas pseudo-effectif, et notre théorème est démontré.
Supposons à présent que $`X`$ est une variété projective, de dimension $`n`$. D’après les résultats de , on sait que si le fibré canonique n’est pas pseudo-effectif, alors $`X`$ est unireglée. Le théorème 3 admet donc le corollaire suivant. Corollaire Soit $`X`$ une variété projective de dimension $`n`$. S’il existe une application holomorphe $`\phi :\mathrm{}^nX`$ non dégénérée, et telle que:
$$\underset{r\mathrm{}}{lim}1/r^{2n}\mathrm{deg}(\phi |\mathrm{\Delta }_r)=0,$$
alors $`X`$ est uniréglée.
Conjecturalement, toute compactification de $`\mathrm{}^n`$ devrait être rationnelle. Donc on peut voir le corollaire précédent comme un premier pas vers cette conjecture. Malheureusement, la façon dont les courbes rationnelles apparaissent dans notre résultat (i.e. via la théorie de Mori) n’est pas du tout explicite.
Il est d’ailleurs possible qu’aucune courbe algébrique de $`\mathrm{}^n`$ n’ait pour image une courbe algébrique de $`X`$, même si $`\mathrm{deg}(\phi |\mathrm{\Delta }_r)`$ est uniformément bornée: soit $`\mathrm{\Omega }\mathrm{}^n\mathrm{}^n\mathrm{}^n=X`$ un domaine de Fatou-Bieberbach obtenu comme bassin d’attraction d’un automorphisme régulier (voir , 2.2, p. 124, et p.125 pour des exemples) de $`\mathrm{}^n,n3`$. Alors $`\mathrm{\Omega }`$ ne contient aucune courbe algébrique (par , Théorème 2.4.4 et remarque 2.4.5 (1) , p. 136), et on est dans la situation évoquée ci-dessus.
La méthode de démonstration du théorème 1 précédent devrait permettre d’ établir la généralisation suivante (illustrée par le produit $`X=\mathrm{}^{np+1}\times A`$, où $`A`$ est une variété abélienne de dimension $`p1`$):
Question: Soit $`X`$ une variété projective de dimension $`n`$. S’il existe une application holomorphe $`\phi :\mathrm{}^nX`$ non dégénérée, et telle que:
$$\underset{r\mathrm{}}{lim}1/r^{2p}\mathrm{deg}(\phi |\mathrm{\Delta }_r)=0,$$
et si $`f:XY`$ est une application méromorphe surjective telle que $`dim(Y)=p`$, alors: $`Y`$ est-elle uniréglée?
Une réponse affirmative à cette question montrerait que le quotient rationnel de $`X`$ (voir et ) est de dimension au plus $`(p1)`$. En particulier, si $`\mathrm{deg}(\phi |\mathrm{\Delta }_r)`$ est borné, $`X`$ (et donc toute compactification de $`\mathrm{}^n`$) serait du moins rationnellement connexe.
2. Preuve du théorème 2
2.1 On commence par rappeler la notion de dimension numérique d’une classe de cohomologie pseudo-effective, telle qu’elle a été introduite par Boucksom, Tsuji dans , , .
Soit $`\{\alpha \}H^{1,1}(X,\mathrm{})`$ une classe de cohomologie pseudo-effective (i.e., il existe un courant positif fermé $`T\{\alpha \}`$, voir pour une présentation plus complète de cette notion). Pour chaque nombre réel $`\epsilon >0`$, la classe de cohomologie $`\{\alpha +\epsilon \omega \}`$ est “big”, et il existe un courant positif fermé $`T_\epsilon \{\alpha +\epsilon \omega \}`$, tel que ses singularités sont concentrées le long d’un ensemble analytique $`Y_\epsilon `$. On définit le nombre d’intersection mobile (voir ) comme suit
$$(\alpha ^k\omega ^{nk}):=\underset{\epsilon 0}{lim}\underset{T_\epsilon \{\alpha +\epsilon \omega \}}{sup}_{XY_\epsilon }T_\epsilon ^k\omega ^{nk}.$$
Définition () Soit $`\{\alpha \}H^{1,1}(X,\mathrm{})`$ une classe de cohomologie pseudo-effective sur $`X`$. La dimension numérique de $`\alpha `$ est définie par
$$\nu (\alpha ):=\mathrm{max}\{k\mathrm{}/(\alpha ^k\omega ^{nk}0\}.$$
Il a été demontré par S.Boucksom (voir ) qu’une classe $`\{\alpha \}`$ est de dimension numérique maximale si et seulement si elle contient un réprésentant strictement positif. Concernant l’autre cas extrême $`\nu (\alpha )=0`$, la situation est malheureusement beaucoup moins bien comprise (cf. pour quelques résultats dans cette direction).
2.2 Supposons à présent qu’on ait $`\nu :=\nu (K_X)1`$. Alors pour tout $`\epsilon >0`$, on a un courant positif $`T_\epsilon c_1(K_X)+\epsilon \{\omega \}`$, dont les singularités sont concentrées le long d’un ensemble analytique $`Y_\epsilon `$, tel que
$$_{XY_\epsilon }T_\epsilon ^\nu \omega ^{n\nu }\delta _0>0$$
uniformément par rapport à $`\epsilon `$. Autrement dit, il existe une famille de modifications $`\mu _\epsilon :X_\epsilon X`$ telle que $`\mu _\epsilon ^{}T_\epsilon =[E_\epsilon ]+\stackrel{~}{\alpha }_\epsilon `$, où $`E_\epsilon `$ est un $`\mathrm{}`$-diviseur effectif, (dont la partie $`\mu _\epsilon `$-exceptionnelle provient des singularités de $`T_\epsilon `$ en codimension 2), et $`\stackrel{~}{\alpha }_\epsilon `$ est une $`(1,1)`$–forme semi-positive sur $`X_\epsilon `$, tels que
$$_{X_\epsilon }\stackrel{~}{\alpha }_\epsilon ^\nu \mu _\epsilon ^{}\omega ^{n\nu }\delta _0>0.$$
2.3 On rappelle maintenant que si $`\mu :X_1X`$ est un éclatement de centre lisse $`YX`$, alors il existe une métrique $`h`$ sur le fibré associé au diviseur exceptionnel $`E`$ telle que pour tout $`0<\delta 1`$, la forme différentielle $`\stackrel{~}{\omega }_\epsilon :=\mu ^{}\omega \delta \mathrm{\Theta }_h(E)`$ est positive définie sur $`X_1`$. Ainsi, pour chaque $`\epsilon >0`$, on construit $`\stackrel{~}{\omega }_\epsilon `$ une métrique kählérienne sur $`X_\epsilon `$, telle que:
(i)$`\stackrel{~}{\omega }_\epsilon \mu _\epsilon ^{}\omega `$ soit un multiple (négatif) du diviseur exceptionnel.
(ii)en chaque point de $`X_\epsilon `$, on a $`\stackrel{~}{\omega }_\epsilon 1/2\mu _\epsilon ^{}\omega `$.
(iii)le volume de $`(X_\epsilon ,\stackrel{~}{\omega }_\epsilon )`$ soit majoré par le volume de $`(X,2^{1/n}\omega )`$.
En effet, les conditions (i)-(iii) sont clairement satisfaites dans le cas d’un seul éclatement de centre lisse (quitte à choisir le paramètre $`\delta `$ ci-dessus assez petit). Le cas général est déduit du fait qu’on peut supposer que la modification $`\mu _\epsilon `$ est une composée d’éclatements de centres lisses.
2.4 Nous utilisons maintenant le théorème de Yau , afin d’obtenir dans la classe $`\{\stackrel{~}{\alpha }_\epsilon +\epsilon \stackrel{~}{\omega }_\epsilon \}`$ un réprésentant dont le déterminant est constant par rapport à $`\stackrel{~}{\omega }_\epsilon `$. Théorème (Yau). Soit $`(X,\beta )`$ une variété kählérienne compacte, de dimension $`n`$, et soit $`dV`$ un élément volume sur $`X`$, tel que $`_X\beta ^n=_X𝑑V`$. Alors il existe $`\rho 𝒞^{\mathrm{}}(X)`$ telle que
(i)$`\beta +i\overline{}\rho >0`$ sur $`X`$.
(ii)$`(\beta +i\overline{}\rho )^n=dV`$. On applique le théorème précédent pour chaque $`X_\epsilon `$, munie de la métrique kählérienne $`\stackrel{~}{\alpha }_\epsilon +\epsilon \stackrel{~}{\omega }_\epsilon `$, et l’élément volume $`dV:=C(\epsilon )\stackrel{~}{\omega }_\epsilon ^n`$. pour une certaine constante adéquate de normalisation $`C(\epsilon )`$, dont le role est de satisfaire la condition cohomologique du théorème de Yau. Ainsi on montre l’existence d’une famille de fonctions $`\rho _\epsilon 𝒞^{\mathrm{}}(X_\epsilon )`$, telle que
(a)$`\stackrel{~}{\alpha }_\epsilon +\epsilon \stackrel{~}{\omega }_\epsilon +i\overline{}\rho _\epsilon >0`$ sur $`X_\epsilon `$.
(b)$`(\stackrel{~}{\alpha }_\epsilon +\epsilon \stackrel{~}{\omega }_\epsilon +i\overline{}\rho _\epsilon )^n=C(\epsilon )\stackrel{~}{\omega }_\epsilon ^n`$.
Nous observons maintenant que, grâce au théorème de Yau, les propriétés numériques de la classe canonique se reflètent dans ses propriétés métriques, i.e. on obtient un minorant pour la constante $`C(\epsilon )`$ en intégrant l’égalité (b) ci-dessus. Ainsi, on montre l’existence d’une constante $`C_1>0`$, telle que $`C(\epsilon )C_1\epsilon ^{n\nu }`$, lorsque $`\epsilon 0`$. En effet, l’égalité (b) implique
$$\begin{array}{cc}\hfill C(\epsilon )Vol(X_\epsilon ,\stackrel{~}{\omega }_\epsilon )=& _{X_\epsilon }\left(\stackrel{~}{\alpha }_\epsilon +\epsilon \stackrel{~}{\omega }_\epsilon \right)^n\hfill \\ \hfill & C_n^\nu \epsilon ^{n\nu }_{X_\epsilon }\stackrel{~}{\alpha }_\epsilon ^\nu \stackrel{~}{\omega }_\epsilon ^{n\nu }C_n^\nu /2^{n\nu }\epsilon ^{n\nu }_{X_\epsilon }\stackrel{~}{\alpha }_\epsilon ^\nu \mu _\epsilon ^{}\omega ^{n\nu }\hfill \\ \hfill & C_0\epsilon ^{n\nu }\hfill \end{array}$$
(dans la suite des inégalités précédentes, on a utilisé la semi-positivité de $`\alpha _\epsilon `$, ainsi que la définition de la dimension numérique dans le cadre pseudo-effectif). L’existence de $`C_1`$ se déduit du fait que le volume de $`(X_\epsilon ,\stackrel{~}{\omega }_\epsilon )`$ est uniformément majoré.
2.5 Considérons le courant image directe $`\mathrm{\Theta }_\epsilon =\mu _{\epsilon ,}(\stackrel{~}{\alpha }_\epsilon +\epsilon \stackrel{~}{\omega }_\epsilon +[E_\epsilon ]+i\overline{}\rho _\epsilon )`$ sur la variété $`X`$. C’est un courant positif fermé et sa classe de cohomologie est $`c_1(K_X)+2\epsilon \{\omega \}`$, comme on le voit immédiatement. De plus, $`\mathrm{\Theta }_\epsilon `$ est non-singulier sur $`X\backslash Y_\epsilon `$ et l’équation de Calabi-Yau (ii) qu’on résout sur $`X_\epsilon `$ montre l’existence d’une constante $`C_2>0`$ telle que
$$\mathrm{\Theta }_\epsilon ^nC_2\epsilon ^{n\nu }\omega ^n$$
$`(4)`$
en tout point de $`X\backslash Y_\epsilon `$. Pour vérifier la relation (4), plaçons-nous en un point $`x_0XY_\epsilon `$. Au voisinage de $`x_0`$, on obtient $`\mathrm{\Theta }_\epsilon =\mu _{\epsilon ,}(\stackrel{~}{\alpha }_\epsilon +\epsilon \stackrel{~}{\omega }_\epsilon +i\overline{}\rho _\epsilon )`$, car le support de $`[E_\epsilon ]`$ est contenu dans l’image inverse de $`Y_\epsilon `$. L’égalité (b) combinée avec la propriété (ii) de la famille de métriques $`(\omega _\epsilon )`$ et le minorant de $`C(\epsilon )`$ montrent que
$$(\stackrel{~}{\alpha }_\epsilon +\epsilon \stackrel{~}{\omega }_\epsilon +i\overline{}\rho _\epsilon )^nC_1\epsilon ^{n\nu }/2^n\mu _\epsilon ^{}\omega ^n.$$
$`(5)`$
Maintenant au voisinage de $`x_0`$, l’application $`\mu _\epsilon `$ est un isomorphisme, et ainsi l’inégalité (5) implique (4), avec une constante $`C_2:=C_1/2^n`$.
2.6 On écrit maintenant $`\mathrm{\Theta }_\epsilon =\alpha +2\epsilon \omega +i\overline{}f_\epsilon `$, pour une certaine fonction $`f_\epsilon L^1(X)𝒞^{\mathrm{}}(X\backslash Y_\epsilon )`$, normalisée de telle sorte que $`{\displaystyle _X}f_\epsilon \omega ^n=0`$ (pour simplifier l’écriture, on a noté $`\alpha :=\mathrm{\Theta }_{det(\omega )}(K_X)`$). Le but de la normalisation de $`f_\epsilon `$ est d’obtenir la relation d’uniformité suivante.
Lemme. Il existe une constante positive $`C_3:=C_3(X,\alpha ,\omega )`$ telle que
$$\mathrm{max}(_X|f_\epsilon |𝑑V_\omega ,\underset{X}{sup}(f_\epsilon ))C_3$$
$`(6)`$
pour tout $`\epsilon >0`$.
Preuve du lemme. Tout d’abord, le lemme est bien connu dans la cadre suivant: soit $`\beta `$ une (1, 1)-forme différentielle de classe $`𝒞^{\mathrm{}}`$ sur $`X`$, et $`f𝒞^{\mathrm{}}(X)`$ une fonction telle que $`\beta +i\overline{}f0`$, et telle que $`_Xf𝑑V_\omega =0`$. Alors
$$\mathrm{max}(_X|f|𝑑V_\omega ,\underset{X}{sup}(f))C(X,\beta ,\omega )$$
$`(7)`$
(voir par exemple ). L’énoncé général est conséquence de ce fait, et du théorème de regularisation suivant, dû à Demailly: Théorème (().) Soit $`X`$ une variété complexe compacte et soit $`T=\alpha +i\overline{}f`$ un courant positif fermé de type $`(1,1)`$. Alors pour chaque entier positif $`k`$, il existe une fonction $`f_k𝒞^{\mathrm{}}(X)`$ telle que:
(1)$`f_kf`$ en norme $`L^1`$ et ponctuellement sur $`X`$. Posant, de plus: $`T_k=\alpha +i\overline{}f_k`$, alors:
(2)$`T_k\{\alpha \}`$, et $`T_k\lambda _k\omega `$, où $`\lambda _k(x)\nu (T,x)`$ lorsque $`k\mathrm{}`$ (autrement dit, en chaque point $`xX`$, la perte de positivité est de la taille du nombre de Lelong $`\nu (T,x)`$ du courant initial $`T`$).
En effet, on applique à chaque $`f_\epsilon `$ le théorème précédent; l’observation importante est la suivante. Etant donné que la classe de cohomologie des courants $`\mathrm{\Theta }_\epsilon `$ est bornée par rapport à $`\epsilon `$, les nombres de Lelong de $`\mathrm{\Theta }_\epsilon `$ le sont également, car ils sont dominés par la masse du courant, et dans le cadre kählérien, cette quantité est cohomologique (ceci reste, en fait, vrai pour les courants de type $`(1,1)`$ sur une variété complexe compacte, grâce à l’existence des métriques de Gauduchon). Donc, la perte de positivité pour les courants régularisés est uniforme par rapport aux paramètres $`\epsilon ,k`$, et ainsi notre lemme est démontré par le résultat rappelé au début de la preuve.
2.7 Revenons maintenant à notre application $`\phi :\mathrm{}^nX`$. Si on note $`\mathrm{\Delta }_1`$ le polydisque unité, alors $`B_\omega (x_0,\delta )\phi (\mathrm{\Delta }_1)`$ pour un certain rayon $`\delta >0`$, et la relation (1) montre l’existence d’un $`x_\epsilon \mathrm{\Delta }_1`$ tel que, pour tout $`\epsilon >0`$:
$$f_\epsilon \phi (x_\epsilon )C_3^{}/\delta ^{2n}$$
$`(7^{})`$
Afin de ne pas trop alourdir les notations, on va supposer que $`x_\epsilon =0`$, l’origine de $`\mathrm{}^n`$ (les arguments presentés par la suite montreront qu’on peut se le permettre, car translater l’origine de $`\mathrm{}^n`$ en $`x_\epsilon `$ revient à changer $`r`$ en $`r+1`$ à la fin de la preuve, et la conclusion désirée n’est pas affectée par ce changement).
Sur $`\mathrm{}^n`$, considérons le jacobien $`J(\phi )H^0(\mathrm{}^n,\phi ^{}K_X^1)`$. La norme de $`J(\phi )`$ par rapport à la métrique $`det(\omega )`$ sur $`K_X^1`$ est donnée par l’égalité $`\phi ^{}\omega ^n=J(\phi )_\omega ^2d\lambda `$, où $`d\lambda `$ est la mesure de Lebesgue de $`\mathrm{}^n`$. Ainsi, on a
$$i\overline{}\mathrm{log}\left(J(\phi )_\omega ^2e^{f_\epsilon \phi }\right)+2\epsilon \phi ^{}\omega \phi ^{}\mathrm{\Theta }_\epsilon $$
car la différence est donnée par la courant d’intégration sur le lieu des zéros du jacobien.
Maintenant, en tout point $`z\mathrm{}^n\backslash \phi ^1(Y_\epsilon )`$, on a l’inégalité
$$\phi ^{}\mathrm{\Theta }_\epsilon (i\overline{}z^2)^{n1}C_2^{1/n}\epsilon ^{1\frac{\nu }{n}}J(\phi )^{\frac{2}{n}}(i\overline{}z^2)^n$$
Pour vérifier cette affirmation, soient $`\lambda _\epsilon ^{(j)}(z)`$ les valeurs propres de $`\phi ^{}\mathrm{\Theta }_\epsilon `$ par rapport à la métrique euclidienne au point $`z\mathrm{}`$. La relation (4) montre que
$$\lambda _\epsilon ^{(j)}(z)C_2\epsilon ^{n\nu }J(\phi ,z)_\omega ^2$$
et par l’inégalité de la moyenne on déduit un minorant pour la somme des valeurs propres, qui est précisement l’inégalité ci-dessus. De plus, on remarque que cette relation reste vraie sur $`\mathrm{}^n`$ tout entier, au sens des distributions. En conclusion on obtient
$$\left(i\overline{}\mathrm{log}J(\phi )_\omega ^2e^{f_\epsilon \phi }+2\epsilon \phi ^{}\omega \right)(i\overline{}z^2)^{n1}\epsilon ^{1\frac{\nu }{n}}J(\phi )_\omega ^{\frac{2}{n}}d\lambda $$
$`(8)`$
2.8 Pour la suite, on voudrait utiliser l’inégalité (8) afin d’appliquer les arguments de courbure négative (cf. Ahlfors, Kodaira, Griffiths) mais malheureusement, on ne peut pas le faire directement, à cause du terme $`\phi ^{}\omega (i\overline{}z^2)^{n1}`$ (la trace de la métrique image inverse par rapport à la métrique euclidienne). Afin d’inclure ce terme dans la derivée logarithmique ci-dessus, on résout le problème de Dirichlet suivant.
$$\{\begin{array}{cc}& \mathrm{\Delta }\mathrm{\Psi }_r=(\phi ^{}\omega )(i\overline{}z^2)^{n1}/d\lambda danszr\hfill \\ & \mathrm{\Psi }_r=0surz=r.\hfill \end{array}$$
La solution $`\mathrm{\Psi }_r`$ est déterminée de manière unique par la formule de Green; on laisse au soin du lecteur de vérifier que la valeur au bord qu’on a considérée est optimale pour ce qui va suivre.
A l’aide de la fonction $`\mathrm{\Psi }_r`$ l’inégalité (8) s’écrit
$$\mathrm{\Delta }\left(\mathrm{log}J(\phi )_\omega ^{\frac{2}{n}}e^{\frac{\epsilon \mathrm{\Psi }_r+f_\epsilon \phi }{n}}\right)\frac{C_2^{1/n}}{n}\epsilon ^{1\frac{\nu }{n}}J(\phi )_\omega ^{\frac{2}{n}}$$
$`(9)`$
Grâce à l’inégalité (6), la dernière relation implique
$$\mathrm{\Delta }\left(\mathrm{log}J(\phi )_\omega ^{\frac{2}{n}}e^{\frac{\epsilon \mathrm{\Psi }_r+f_\epsilon \phi }{n}}\right)C_4\epsilon ^{1\frac{\nu }{n}}J(\phi )_\omega ^{\frac{2}{n}}e^{\frac{f_\epsilon \pi }{n}}$$
$`(10)`$
(l’expression de la constante $`C_4`$ se déduit facilement de $`C_2`$ et $`C_3`$).
Par ailleurs, la fonction $`\mathrm{\Psi }_r`$ solution du problème de Dirichlet est sous-harmonique, donc par le principe du maximum, $`\underset{zr}{\mathrm{max}}\mathrm{\Psi }_r(z)=0`$, et finalement on peut écrire
$$\mathrm{\Delta }\left(\mathrm{log}J(\phi )_\omega ^{\frac{2}{n}}e^{\frac{\epsilon \mathrm{\Psi }_r+f_\epsilon \phi }{n}}\right)C\epsilon ^{1\frac{\nu }{n}}J(\phi )_\omega ^{\frac{2}{n}}e^{\frac{\epsilon \mathrm{\Psi }_r+f_\epsilon \phi }{n}}$$
$`(11)`$
Les arguments de courbure négative suivants sont classiques (voir par exemple , ). Sur $`zr`$, considérons la fonction $`\tau _\epsilon `$ définie par :
$$z\left(1\frac{z^2}{r^2}\right)^2J(\phi ,z)_\omega ^{\frac{2}{n}}e^{\frac{\epsilon \mathrm{\Psi }_r(z)+f_\epsilon \phi (z)}{n}}:=\tau _\epsilon (z)$$
C’est une fonction positive, et sur la sphère $`\left(z=r\right)`$ elle vaut zéro; par conséquent, son point de maximum est atteint en un $`z_{max}`$ de $`z<r`$. Alors on a
$$\mathrm{\Delta }\mathrm{log}\tau _\epsilon (z_{max})0,$$
Par ailleurs, un calcul sans difficulté montre que:
$$i\overline{}\mathrm{log}\left(1\frac{z^2}{r^2}\right)\left(i\overline{}z^2\right)^{n1}\frac{d\lambda }{r^2(1\frac{z^2}{r^2})^2}$$
$`(12)`$
Donc au point de maximum $`z=z_{max}`$ on aura:
$$C_4\epsilon ^{1\frac{\nu }{n}}J(\phi ,z)_\omega ^{\frac{2}{n}}e^{\frac{\epsilon \mathrm{\Psi }_r(z)+f_\epsilon \phi (z)}{n}}\frac{1}{r^2(1\frac{z^2}{r^2})^2},$$
comme conséquence de (11) et (12). Alors on obtient:
$$C_4\epsilon ^{1\frac{\nu }{n}}J(\phi ,z)_\omega ^{\frac{2}{n}}e^{\frac{\epsilon \mathrm{\Psi }_r(z)+f_\epsilon \phi (z)}{n}}(1\frac{z^2}{r^2})^2\frac{1}{r^2}$$
au point $`z_{max}`$, et donc $`\tau _\epsilon (z){\displaystyle \frac{1}{C_4\epsilon ^{1\frac{\nu }{n}}r^2}}`$ pour $`z\mathrm{𝔹}(r)`$. En particulier, si $`z=0`$, l’inégalité précédente devient:
$$J(\phi ,0)_\omega ^{\frac{2}{n}}\epsilon ^{1\frac{\nu }{n}}e^{\frac{\epsilon \mathrm{\Psi }_r(0)+f_\epsilon \phi (0)}{n}}\frac{1}{r^2}.$$
Afin de déterminer la quantité $`\mathrm{\Psi }_r(0)`$, on applique la formule de Green
$$_{z=r}\mathrm{\Psi }_r(\xi )𝑑\sigma \mathrm{\Psi }_r(0)=_0^r\frac{dt}{t^{2n1}}_{B(t)}\mathrm{\Delta }\mathrm{\Psi }_r𝑑\lambda $$
et comme $`\mathrm{\Psi }_r`$ est la solution du problème de Dirichlet, on voit que $`\mathrm{\Psi }_r(0)=T(\phi ,r)`$, l’indicatrice de croissance de la fonction $`\phi `$. En conclusion, on a
$$J(\phi ,0)_\omega ^{\frac{2}{n}}\frac{C_5}{r^2}\frac{1}{\epsilon ^{1\frac{\nu }{n}}}e^{\epsilon T(\phi ,r)},$$
compte-tenu de l’inégalité (7’).
On peut remarquer que dans l’expression précédente, la constante $`C_5`$ dépend uniquement des quantités suivantes: géométrie de $`(X,\omega )`$, et rayon de la boule euclidienne contenue dans l’image $`\phi (\mathrm{\Delta }_1)`$.
Maintenant on observe que les paramètres $`\epsilon `$ et $`r`$ sont indépendants. Par conséquent, on peut prendre $`\epsilon :=T^1(\phi ,r)`$ et on obtient
$$J(\phi ,0)_\omega ^{\frac{2}{n}}\frac{C}{r^2}(T(\phi ,r))^{1\frac{\nu }{n}}$$
La preuve du théorème 1 est ainsi complète.
Remarque. Comme conséquence de la méthode de démonstration précédente, on obtient l’enoncé suivant: Théorème $`2^{}`$ Soit $`(X^n,\omega )`$ une variété kählérienne compacte de dimension $`n`$, dont le fibré canonique $`K_X`$ est est pseudo-effectif, avec $`\nu :=\nu (K_X)`$.
Si $`\phi :B^n(r)X`$ est définie sur une boule de rayon $`r>1`$ de $`\mathrm{}`$<sup>n</sup>, et si $`\phi (B(1))B_\omega (x_0,\delta _0)`$ pour un certain $`x_0X`$, alors
$$T(\phi ,r^{})^{1\nu /n}C(X,\omega ,\delta _0)(r^{})^2J(\phi ,0)_\omega ^{\frac{2}{n}}r^{}r.$$
Pour faire le lien avec la première partie de cet article, il serait intéressant d’analyser les conséquences de la positivité numérique de $`K_X`$ sur la croissance du degré moyen de l’application $`\phi `$ (i.e., d’obtenir l’analogue du théorème 2 pour le degré moyen à la place de la fonction caractéristique de $`\phi `$).
3. Applications quasi-conformes.
Nous rappelons qu’étant donné un point $`x_0X`$, on peut considérer la quantité :
$$r(x_0)=sup\{r>0/f:\mathrm{𝔹}^n(r)Xtelque\left\{\begin{array}{cc}& f(0)=x_0\hfill \\ & J(f,0)_\omega =1\hfill \end{array}\right\}$$
Alors le théorème de Kobayashi-Ochiai peut être reformulé comme suit. Théorème (KO) Soit $`X`$ une variété projective, dont le fibré canonique est big. Alors $`r(x_0)<\mathrm{},x_0X`$. Comme cela nous a été sugéré par Y-T Siu, on peut considérer une quantité analogue où seules les applications quasi-conformes seront impliquées. On définit
$$r_{q,C}(x_0)=sup\{r>0/f:\mathrm{𝔹}^n(r)Xtelque\left\{\begin{array}{cc}& f(0)=x_0\hfill \\ & J(f,0)_\omega =1\hfill \\ & df(z)_\omega <CJ(f,z)_\omega ^{\frac{1}{n}}\hfill \end{array}\right\}$$
Dans ce contexte, la méthode de démonstration du théorème 2 donne le résultat suivant. Corollaire 3.1 Soit $`(X,\omega )`$ une variété kählérienne compacte avec $`K_X`$ nef. Si $`r_{q,C}(x_0)=+\mathrm{}`$ pour un certain $`x_0X`$, alors $`\nu (K_X)=0`$.
Preuve. L’hypothèse implique l’existence d’une famille de boules $`\phi _r:\mathrm{\Delta }_rX`$ telle que
$$\{\begin{array}{cc}& (1)\phi _r(0)=x_0\hfill \\ & (2)J(\phi _r,0)=1\hfill \\ & (3)d\phi _r(z)<CJ(\phi _r,z)^{\frac{1}{n}}\hfill \end{array}$$
où la dernière condition est vérifiée en tout point $`z`$ du polydisque $`\mathrm{\Delta }_r`$.
Le fibré canonique de $`X`$ étant nef, le théorème de Yau (voir la preuve du théorème 2) montre l’existence d’une suite de métriques $`h_\epsilon :=he^{f_\epsilon }`$, telles que $`f_\epsilon 𝒞^{\mathrm{}}(X)`$, et telles que
$$(\mathrm{\Theta }_h(K_X)+i\overline{}f_\epsilon +\epsilon \omega )^n=C_\epsilon \omega ^n$$
$`(13)`$
où la constante de normalisation $`C_\epsilon `$ est de l’ordre de $`𝒪(\epsilon ^{n\nu })`$, et où $`\nu =\nu (K_X)`$ désigne la dimension numérique du fibré canonique de $`X`$.
Pour chaque $`r>0`$, l’inégalité de la moyenne combinée avec (13) implique
$$(i\overline{}\mathrm{log}J(\phi _r)^2e^{f_\epsilon \phi _r}+\epsilon \phi _r^{}\omega )(i\overline{}z^2)^{n1}\epsilon ^{1\nu /n}J(\phi _r)^{2/n}d\lambda $$
$`(14)`$
Supposons à présent que $`\nu (K_X)1`$; on aura alors clairement $`\epsilon ^{1\nu /n}\epsilon `$, lorsque $`\epsilon 0`$. L’inégalité (14) et l’hypothèse sur $`\phi _r`$ donnent:
$$\mathrm{\Delta }\mathrm{log}J(\phi _r)^{2/n}e^{f_\epsilon \phi _r/n}\epsilon ^{1\nu /n}J(\phi _r)^{2/n}e^{f_\epsilon \phi _r/n}.$$
$`(15)`$
(car le terme $`(\phi _r^{}\omega )(i\overline{}z^2)^{n1}`$ est dominé –à une constante universelle près–par $`J(\phi _r)^{\frac{2}{n}}d\lambda `$, grâce à la relation de quasi-conformalité (3)).
Fixons par la suite une valeur de $`\epsilon 1`$ telle que $`(15)`$ soit satisfaite. Le reste de la preuve suit les arguments de courbure négative dejà exposés en 2.9, et on obtient ainsi
$$J(\phi _r,0)_\omega ^{2/n}e^{f_\epsilon \phi _r(0)/n}C/r^2.$$
$`(16)`$
Maintenant les conditions (1) et (2) montrent que la relation (16) est équivalente à l’inégalité: $`e^{f_\epsilon (x_0)}C/r^2`$, et on obtient une contradiction lorsque $`r\mathrm{}`$. La preuve du corollaire est achevée.
Dans le cours de ce dernier paragraphe, nous allons étendre le résultat précédent au cas où $`\phi `$ est seulement supposée quasi-conforme en moyenne.
Soit $`\phi :\mathrm{}^nX`$ une application holomorphe non-dégénérée. On suppose l’existence d’une constante $`\delta >0`$, telle que
$$_0^r\frac{dt}{t^{2n1}}_{B(t)}J(\phi )_\omega ^{2/n}𝑑\lambda CT_\omega (\phi ,r)$$
$`(17)`$
pour tout $`r0`$. Par exemple, si l’application $`\phi `$ est quasi-conforme, la condition $`(17)`$ est automatiquement satisfaite; en général, appelons une telle application quasi-conforme en moyenne. Nous allons maintenant présenter une preuve de l’enoncé suivant.
Théorème 3.2. Soit $`X`$ une variété kählérienne compacte de dimension $`n`$, et $`\phi :\mathrm{}^nX`$ une application holomorphe quasi-conforme en moyenne. Si le fibré canonique de $`X`$ est nef, alors $`\nu (K_X)=0`$.
Preuve. Dans ce qui va suivre, on utilise de façon essentielle des arguments de la théorie de Nevanlinna, et également les métriques sur le fibré canonique qu’on a construites au cours de la preuve du théorème 2. Le point de départ est l’inégalité:
$$i\overline{}\mathrm{log}||J(\phi )||^2e^{f_\epsilon \phi }+\epsilon \phi ^{}\omega )(i\overline{}||z||^2)^{n1}\epsilon ^{1\nu /n}||J(\phi )||^{2/n}d\lambda $$
$`(18)`$
en tout point de $`\mathrm{}^n`$. On intègre cette relation à la manière de Nevanlinna; ainsi, pour tout $`r>0`$ on aura:
$$\begin{array}{cc}\hfill _0^r\frac{dt}{t^{2n1}}_{B(t)}\mathrm{\Delta }\mathrm{log}J(\phi )_\omega ^2e^{f_\epsilon \phi }𝑑\lambda & \end{array}$$
$$\epsilon ^{1\nu /n}_0^r\frac{dt}{t^{2n1}}_{B(t)}J(\phi )_\omega ^{2/n}𝑑\lambda \epsilon T_\omega (\phi ,r).$$
Si $`\nu (K_X)>0`$, alors quitte à prendre $`\epsilon 1`$, (qui sera fixé par la suite), notre hypothèse de quasi-conformalité en moyenne et la relation ci-dessus entrainent:
$$_0^r\frac{dt}{t^{2n1}}_{B(t)}\mathrm{\Delta }\mathrm{log}J(\phi )_\omega ^2e^{f_\epsilon \phi }𝑑\lambda C_\epsilon _0^r\frac{dt}{t^{2n1}}_{B(t)}J(\phi )_\omega ^{2/n}𝑑\lambda $$
$`(19)`$
Maintenant, on rappelle la formule de Jensen suivante:
$$_0^r\frac{dt}{t^{2n1}}_{B(t)}\mathrm{\Delta }\rho 𝑑\lambda =_{|z|=r}\rho 𝑑\sigma \rho (0)$$
valable pour toute fonction $`\rho `$ assez régulière pour que les quantités ci-dessus aient un sens.
Grâce à cette identité, le terme se trouvant de gauche de l’inégalité (19) est égal à $`{\displaystyle _{|z|=r}}\mathrm{log}J(\phi )^2𝑑\sigma +𝒪(1)`$. La concavité de la fonction logarithme et les considérations précédentes impliquent l’inégalité suivante:
$$\mathrm{log}_{|z|=r}J(\phi )_\omega ^{2/n}𝑑\sigma +𝒪(1)C_0^r\frac{dt}{t^{2n1}}_{B(t)}J(\phi )_\omega ^{2/n}𝑑\lambda $$
$`(20)`$
La preuve sera achevée si l’on montre que l’inégalité (20) est impossible à satisfaire, pour des valeurs de $`r`$ assez grandes.
Pour chaque $`r>0`$, soit $`𝒥(r):={\displaystyle _0^r}{\displaystyle \frac{dt}{t^{2n1}}}{\displaystyle _{B(t)}}J(\phi )_\omega ^{2/n}𝑑\lambda `$. La fonction ainsi définie est croissante, et en dérivant succesivement on obtient:
$$r^{2n1}_{|z|=r}J(\phi )^{2/n}𝑑\sigma =\left(r^{2n1}𝒥^{}\right)^{}$$
pour toute valeur de $`r`$. On fait appel maintenant au lemme suivant (du type E. Borel), classique dans la théorie de Nevanlinna (voir ).
Lemme. Soit $`F:[0,\mathrm{})\mathrm{}_+`$ croissante, dérivable. Alors pour tout $`\delta >0`$, il existe un ensemble $`E_\delta \mathrm{}_+`$ tel que $`{\displaystyle _{E_\delta }}d\mathrm{log}t<\mathrm{}`$ et tel que l’on ait: l’inégalité $`rF^{}(r)F^{1+\delta }(r)`$, si $`r\mathrm{}_+E_\delta `$.
Dans notre situation, on applique d’abord l’enoncé précédent à la fonction $`F(r)=r^{2n1}𝒥^{}(r)`$, avec $`\delta =1/2n1`$. Le lemme précédent montre que
$$r\left(r^{2n1}𝒥^{}(r)\right)^{}\left(r^{2n1}𝒥^{}(r)\right)^{1+1/2n1}=r^{2n}\left(𝒥^{}(r)\right)^{1+1/2n1}.$$
Une nouvelle application du lemme, cette fois à la fonction $`F(r)=𝒥(r)`$, montre que $`𝒥^{}(r)𝒥^2`$, en dehors d’un ensemble de mesure logarithmique finie. En conclusion, on aura
$$\left(r^{2n1}𝒥^{}(r)\right)^{}r^{2n1}𝒥^4(r)$$
$`(21)`$
pour tout $`r\mathrm{}_+E`$. On ré-ecrit maintenant l’inégalité (15) comme suit:
$$4\mathrm{log}_0^r\frac{dt}{t^{2n1}}_{B(t)}J(\phi )_\omega ^{2/n}𝑑\lambda +𝒪(1)C_0^r\frac{dt}{t^{2n1}}_{B(t)}J(\phi )_\omega ^{2/n}𝑑\lambda $$
si $`r\mathrm{}_+E`$. On choisit maintenant une suite $`(r_k)\mathrm{}_+E`$ telle que $`r_k\mathrm{}`$. La dernière relation donne alors une contradiction.
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# Untitled Document
Comments on the structure and dynamics
of magnetic fields in stellar convection zones
Habilitationsschrift
zur Erlangung der Lehrbefugnis im Fachgebiet
Astronomie und Astrophysik
vorgelegt dem Fachbereich Physik in der
Mathematisch-Naturwissenschaftlichen Fakultät
der Georg-August-Universität zu Göttingen
von
Manfred Schüßler
aus Weinheim
Göttingen 1990
Preliminary remark
This paper is a “Habilitationsschrift”, a second thesis required until recently by universities in Germany and in a few other countries to obtain the right to lecture. It was accepted by the University of Göttingen in 1990 after review by a number of german and international experts. Although the introduction and the references represent the state of research in 1990, most of the remaining content is still relevant and has never been published elsewhere. The most important part is the derivation of a linear stability formalism for thin magnetic flux tubes following an arbitrary path in a gravitationally stratified medium with a stationary velocity. It was found later (Ferriz-Mas & Schüssler, Geophys. Astrophys. Fluid Dyn. vol. 72, 209; 1995) that, for consistency, the inertial term in the equation of motion for the external medium should be included in Eq. (3.24), which leads to an additional term in the stability equations in the case of a spatially varying external velocity. This term is missing in the present text, but can be easily introduced into the formalism.
Summary
Some aspects of magnetic fields in stellar convection zones are investigated in this contribution. Observational and theoretical results are discussed which support the conjecture that the magnetic field structure in a convection zone is intermittent with most of the magnetic flux being concentrated in small filaments or ‘flux tubes’ surrounded by field-free plasma. These kind of structures can be mathematically described with aid of the ‘approximation of slender flux tubes’ whose general form for flux tubes which follow an arbitrary path in space is rederived and discussed.
The approximation is applied to equilibrium structures of flux tubes determined by hydrostatic equilibrium along the magnetic field lines and by a balance of buoyancy, curvature and drag forces (exerted by external velocity fields like convection, rotation and large-scale flows) perpendicular to the field. Some general properties of static equilibria (without drag forces) are derived and it is shown that such structures are incompatible with the observed properties of solar magnetic fields. We discuss methods to determine equilibrium flux tubes in practice and give analytical examples of stationary tubes in a horizontal velocity field.
In the central part of the contribution we present a linear stability analysis of general flux tube equilibria including arbitrary external velocity fields. The tube may follow a curved path in space, gradients of external velocity and gravitational acceleration are included. The general equations for Lagrangian displacements are derived and for the application to stellar convection zones we give a suitably non-dimensionalized, approximate form for large values of the plasma parameter $`\beta `$ which represents the ratio of gas pressure to magnetic pressure. For static equilibria a symmetric form of the equations is obtained which allows the application of variational methods and a simplified numerical treatment. It also serves as a consistency check for the (somewhat lengthy) algebra. The formalism is applied to analytically tractable cases, namely horizontal flux tubes and symmetric loops with horizontal tangent at the point of extremum (maximum or minimum). Numerical examples on the basis of the properties of the solar convection zone indicate that in a superadiabatically stratified environment and in the absence of an external flow all these structures are monotonically unstable. Stabilizing external velocity fields can be constructed, but they do not seem to be of much practical importance for the case of a stellar convection zone. Additionally, we have found that overstable modes can be excited by an external velocity field in which case the drag force conspires with the curvature force with the result of oscillations with growing amplitude. Overstable modes cannot be stabilized by stratification; they appear also in subadiabatic regions like a layer of overshooting convection.
Finally, we attempt to summarize our present state of understanding of magnetic fields in stellar convection zones and, in particular, the convection zone of the Sun. We favor the picture of a strongly fragmented, intermittent field structure. While in the deep parts the individual magnetic filaments are passive with respect to large-scale velocity fields, surface fields exhibit a peculiar nature owing to thermal effects and the dominance of buoyancy. As to the dynamo problem, we find that observational and theoretical evidence is in favor of a ‘boundary layer dynamo’ operating in an overshoot region below the superadiabatic parts of a convection zone.
Acknowledgements
Many people have contributed to this work, be it by discussion and criticism, by moral and technical support, or by their own work which laid the foundations on which I could build by adding a small brick to the walls of an unfinished building. Let me thank in particular
V. Anton for the good time we spent calculating flux tube models,
W. Deinzer for his encouragement and support,
H. Düker for help with technical problems,
A. Ferriz-Mas for his scrutinity with mathematical details and comments on Ch. 3,
M. Knölker for sharing his knowledge on eigenvalue problems,
D. Schmitt for advice on stability problems and the energy principle,
M. Stix for discussions about ways to determine the stability of loops, and
H. Spruit for clarifying answers concerning the approximation of slender flux tubes.
This work has been carried out and written down partly at the Kiepenheuer-Institut für Sonnenphysik, Freiburg, and partly at the Universitäts-Sternwarte, Göttingen. I like to thank my colleagues at both institutions for their support and for a climate of cooperation and kindness. Moreover, I am grateful to U. Grossmann-Doerth, R. Hammer and H. Schleicher for having taken the burden of additional work in a critical phase of installation and management of the computer network at the Kiepenheuer-Institut while I was on leave in Göttingen or absorbed in writing.
Finally, I want to thank E.N. Parker for his life-long work on solar and stellar magnetic fields which revealed so many basic mechanisms, opened so many doors and possibilities for investigation and understanding – and for continuously stirring up inveterate ways of thinking by pointing out inconsistencies and throwing new ideas into the arena.
Contents
1. Introduction. 1
2. Formation of structures. 5
2.1 Rayleigh-Taylor instability of a magnetic layer . 5
2.2 Flux expulsion . 8
2.3 Instabilities and fragmentation of single flux tubes . 10
3. The approximation of slender flux tubes. 14
4. Flux tubes in equilibrium. 21
4.1 Static equilibrium . 21
4.2 Stationary equilibrium . 26
5. Stability of flux tubes. 31
5.1 Previous work . 31
5.2 Equilibrium . 33
5.3 Perturbation equations . 34
5.4 Non-dimensionalization and the case $`\beta 1`$. 40
5.5 Symmetric form for static equilibrium . 44
5.6 Horizontal tubes with vertical external flow . 46
5.7 Symmetric loops with vertical external flow . 54
5.7.1 Local analysis . 54
5.7.2 Constant vertical displacement . 58
5.7.3 Heuristic approach for perturbations with large wavelength . 61
5.8 Summary of the stability properties . 67
6. Dynamics of flux tubes in a convection zone. 69
6.1 Size distribution . 69
6.2 The relation of the basic forces . 72
6.3 The peculiar state of the surface fields . 74
6.4 Consequences for the dynamo problem . 76
7. Outlook. 79
References. 81
1. Introduction
The activity of the Sun and other stars with outer convection zones and the origin of their hot chromospheres, coronae, and winds is intimately related to the existence of magnetic fields in their atmospheres. Most probably, the source of the magnetic flux which is observed to pervade photospheres of late-type stars is the underlying convection zone. In the case of the Sun, magnetic flux is directly observed to emerge from the convection zone. On the other hand, the 22-year period of the solar magnetic cycle leads to a very small (skin) depth to which the alternating magnetic fields can penetrate within the electrically well-conducting, quiescent radiative region below the solar convection zone.
A complete theory of the structure, dynamics and evolution of magnetic fields in stellar convection zones does not exist – and this work does not attempt to provide one. This lack of a consistent description of a basic astrophysical situation is due to a) the impossibility of direct measurements and b) our unsufficient comprehension of turbulent flows.
In the case of the Sun, although the physics of the photosphere is much more complicated, thanks to a wealth of observational data our understanding of photospheric magnetic fields is much more advanced than the state of theory for the fields within the convection zone. Only indirect observational evidence about their state is available through photospheric observations – and due to the peculiar nature of this layer the results are not necessarily representative for the deeper parts (cf. Sec. 6.3). The situation can be illustrated by the following simile: Imagine a person who is unfamiliar with the purpose and operation of clocks. Imagine further that this person is confronted with a mechanical clock and the task to analyze the internal mechanism without opening it, just from the visual appearance of the dial, the motion of the hands and a spectral analysis of the ticking. Good luck ! A person who tries to understand the structure and dynamics of the convection zone and its magnetic fields is in a similar situation: Only a shallow surface layer can be observed directly (dial and hands) while global oscillations (ticking) supply some indirect information from the deeper layers.
The large spatial extension of a stellar convection zone and the small viscosity of a stellar plasma lead to enormous values of the Reynolds number $`Re=UL/\nu `$ ($`U`$, $`L`$: velocity and spatial scale of the dominant convective flow, $`\nu `$: kinematic viscosity) which describes the ratio of the magnitudes of inertial force and viscous force. For the granular velocity field observed in the solar photosphere ($`U1`$ km$``$s<sup>-1</sup>, $`L10^3`$ km, $`\nu 10^3`$ m$`{}_{}{}^{2}`$s<sup>-1</sup>) we find $`Re10^{12}`$. Consequently, the nonlinear inertial forces dominate and a turbulent cascade of kinetic energy to larger spatial wavenumbers ensues until scales of less than 1 cm are reached at which disspation by molecular viscosity effectively takes place. Stellar convection zones thus span a huge range of scales which reach from their global dimensions and time scales (of the order of $`10^5`$ km and $`10^6`$ s, respectively) to the disspation range of about 1 cm with a related dynamical time scale of less than 1 s. The problem of describing turbulent convection in a stellar convection zone is rendered even more difficult by the influences of rotation and of stratification due to gravity.
The magnetic Reynolds number $`Re=UL/\eta `$ where $`\eta `$ denotes the magnetic diffusivity is also very large. For the solar convection zone it increases from about $`10^3`$ in the upper layers to about $`10^9`$ near the bottom (cf. Stix, 1976). Consequently, field line advection and stretching is much more important than Ohmic dissipation on the dominant scale of convective flows. Similar to the fluid motions, the magnetic field spans large ranges of spatial and temporal scales which extend from the global convective to the dissipative scales.
Given our insufficient understanding of turbulent convection it may seem futile to complicate things even more by introducing magnetic fields or, more precisely, by taking account of the large electrical conductivity of the stellar plasma. Besides the fact that the very presence of magnetic fields in the solar and stellar atmospheres forces us to do so, the additional complication brought about by including magnetic fields must not necessarily be prohibitive since there is a number of indications that magnetic fields do not significantly influence the global structure and dynamics of stellar convection zones. One piece of evidence is the success of the theory of stellar structure and evolution which has been developed without taking account of magnetic fields while the other line of arguments is provided by the comparatively small magnitude of observed variations of the Sun during the activity cycle.
The reversals of the polar fields, the polarity rules for active regions, and the strong variation of the frequency of occurrence of large active regions indicate a major change of the magnetic structure in the convection zone during the activity cycle. On the other hand, the change of the solar convection zone is much less significant: Variations of the surface rotation rate of both plasma and magnetic structures (at a given heliographic latitude) are smaller than a few percent (Howard, 1984; Schröter, 1985) while a velocity structure associated with the activity belts (misnamed ‘torsional oscillation’ by its discoverers, cf. LaBonte and Howard, 1982) has an amplitude of less than 1% of the differential rotation in latitude. Convective flow patterns do not show a significant change during the cycle either, apart from slight variations of the size distribution of granules (Müller and Roudier, 1984) and, possibly, of their temperature structure (Livingston and Holweger, 1982). It is improbable that much larger variations in the deep layers of the convection are effectively ‘screened’ by the surface layers since perturbations of temperature and velocity are transmitted from the bottom of the convection zone to the surface without being significantly attenuated (Stix, 1981b). Solar cycle variations of the solar radius are smaller than $`200`$ km (Wittmann et al., 1981; Parkinson, 1983). The short-term luminosity variation is of the order of $`10^3`$ and corresponds directly to the fraction of the surface covered by sunspots while a variation of about the same magnitude on the time scale of the cycle is indicated (Willson, 1984). However, larger heat flux variations on the time scale of the solar cycle in the deep layers of a convection zone may be efficiently screened due to its large thermal conductivity and heat capacity (Spruit, 1977a, 1982; Stix, 1981b).
All these results support the thesis that magnetic fields do not significantly modify the convection zone on a global scale. However, one must not conclude that magnetic forces can as well be neglected locally, i.e. at any given location. The contrary seems to be the case: A global equipartition of magnetic and kinetic energy which would lead to major changes of the convection zone during the magnetic cycle is avoided by concentration of the magnetic flux into filaments of strong field with large regions of non-magnetic, undisturbed convection between (Parker, 1984a). Such a ‘phase separation’ in a convecting and a magnetic phase is observed for the case of the solar photosphere: Most of the observable magnetic flux (outside sunspots) is in the form of concentrated structures of high flux density which are arranged in a network defined by the downflow regions of granulation and supergranulation, the dominant convective structures (Stenflo, 1989; Solanki, 1990).
Under the conditions prevailing in a convection zone the flux expulsion process (for a discussion and further references see Sec. 2.2) which is responsible for this separation in magnetic and convecting regions leads to local equipartition of the magnetic and kinetic energy densities which gives a flux density of about $`10^4`$ Gauss in the deep parts of the solar convection zone. Since the amount of magnetic flux which emerges during one half cycle (11 years) is about $`10^{24}`$ mx (Howard, 1974) it fills only about 1% of convection zone volume given such a flux density. Consequently, the total magnetic energy of $`E_{mag}`$ $`310^{35}`$ erg amounts to about 1% of the total kinetic energy ($`E_{con}`$) of the convective flows. The total energy of differential rotation ($`E_{dr}`$) is of the same order of magnitude as $`E_{con}`$ since the velocity differences due to differential rotation in depth and in latitude (Duvall et al., 1984, 1987) are of the same order of magnitude as the convective velocities in the deeper parts of the convection zone. Consequently, we find the following scaling of total energies within the convection zone:
$$E_{con}E_{dr}\mathrm{\hspace{0.33em}10}^2E_{mag}.$$
These relations are consistent with the observed percent-level variations of the properties of the solar convection zone during the activity cycle. In the light of these considerations it seems adequate to take the convective flows as given and undisturbed by the magnetic fields for all scales which are large compared to the typical size of the magnetic flux concentrations. Small-scale flows locally are strongly affected by the presence of the field and probably convective heat exchange between a flux concentration and its environment is largely suppressed.
While a complete theory is lacking, a variety of more or less satisfactory physical descriptions and models of certain aspects of the complicated thermodynamical and (magneto)hydrodynamical system represented by a stellar convection zone can be found in the literature. We do not attempt to give a complete overview but we may roughly classify them in two complimentary groups, namely the mean field approach and the model problem approach. The contributions belonging to the first group attempt to describe the large-scale structure and dynamics of a stellar convection zone with the aid of parametrized model equations for appropriately averaged quantities which vary on large scales. Such equations typically contain parameters or functions which represent conjectures about small-scale processes. Prominent examples are Prandtl’s mixing length formalism which has been used successfully in the theory of stellar structure and evolution and the mean-field approach for magnetic fields (e.g. Parker, 1979a; Krause and Rädler, 1980) which led to the presently most developed theory of the solar cycle, the theory of turbulent dynamo action.
Attempts towards a numerical simulation of the hydrodynamical and magnetic structure of the convection zone (Gilman and Miller, 1981; Gilman, 1983; Glatzmaier 1984, 1985a,b; Brandenburg et al., 1990) also fall into this category: Due to limitations of memory capacity and processor speed of presently available computers only a small part of the range of spatial and temporal scales can be covered by the simulation and the influence of the small scales is parametrized by introducing ‘turbulent’ values for viscosity as well as thermal and magnetic diffusivity. The effective (hydrodynamical and magnetic) Reynolds numbers for such simulations are therefore many orders of magnitude smaller than those of the real system. The results are partly in contradiction to observational data: Neither the predicted uniformity of angular velocity on cylindrical surfaces nor the large amplitude of large-scale convective flows has been borne out by measurements (Duvall et al., 1987; Brown et al. 1989; Dziembowski et al., 1989; LaBonte et al., 1981). Furthermore, the characteristic features of the solar cycle could not be reproduced by the simulations. Apparently processes operating on small scales which have not been resolved play an important rôle for the hydrodynamic and the magnetic structure of the convection zone.
In the contributions belonging to the second group, model problem approaches, individual processes are studied in (artificial) isolation and their relevance for the global behavior of the system is evaluated. This may lead to the introduction of new terms in model equations and to a more sensible parametrization in numerical simulations. Furthermore, one might attempt to combine a sample of reasonable well understood processes like a jigsaw puzzle in order to obtain a description of the whole system. The work of E.N. Parker (cf. Parker, 1979a) is a prominent example for the model problem approach. The theory of magnetic flux tubes (e.g. Spruit and Roberts, 1983), the work on magnetoconvection carried out by N.O. Weiss and his colleagues in Cambridge (cf. Proctor and Weiss, 1982), and the simulations of turbulence with magnetic fields performed by the Nice group around U. Frisch and A. Pouquet (e.g. Meneguzzi et al., 1981; Grappin et al., 1982; Pouquet, 1985; Meneguzzi and Pouquet, 1989) also fall in the group of model problems.
The work presented here belongs to the same class of contributions. Its motivation results from the debate on magnetic flux storage in a convection zone and the location of the dynamo process which is responsible for the solar activity cycle. Parker (1975a) argued that magnetic buoyancy leads to rapid flux loss and thus prohibits the storage of magnetic flux within the convection zone for time intervals comparable to the cycle period. This argument was strengthened by Spruit and van Ballegooijen (1982) who showed that toroidal flux tubes are unstable in a superadiabatically stratified region. Following earlier proposals (e.g. Spiegel and Weiss, 1980; Galloway and Weiss, 1981) they suggested that a slightly subadiabatic region of overshooting convection near the bottom of the solar convection zone represents a favorable place for the storage of magnetic flux and the operation of a dynamo mechanism (see also van Ballegooijen 1982a,b; Schüssler, 1983, 1984a; Durney, 1989).
However, are the arguments given so far sufficient to definitely exclude that the major part of the magnetic flux emerging in the solar activity cycle is stored within the convection zone proper ? For instance, Parker (1987a-d; 1988a-c) has argued that ‘thermal shadows’ due to local suppression of convection can keep large portions of magnetic flux in the deep parts of the convection zone. Here we consider another possibility and investigate whether flux tubes could possibly find a stable equilibrium in the convection zone if they follow a curved path in space and/or if the influence of external flows is taken into account. Of particular interest in this respect are loop structures and sequences of loops (‘sea serpents’). To this end we reconsider the general form of the approximation of slender flux tubes (Spruit, 1981a,b), determine general properties of static and stationary flux tube equilibria, and generalize the approach of Spruit and van Ballegooijen (1982, see also van Ballegooijen, 1983; van Ballegooijen and Choudhuri, 1988) to derive a stability analysis formalism for a flux tube which follows an arbitrary path in space. This analysis is embedded in a general discussion of the structure and dynamics of magnetic fields in a stellar convection zone which is necessarily tentative and far from rigorous.
We start by presenting some arguments in favor of the conjecture that magnetic fields in stellar convection zones consist of small, concentrated structures separated by field-free plasma. As a basis for this conjecture, a variety of physical processes which form and maintain such structures is discussed in Ch. 2. The small size of the magnetic flux concentrations resulting from fragmentation and expulsion processes permits their description using the approximation of slender flux tubes. In Ch. 3 we rederive the general form of this approximation for a flux tube which follows an arbitrary path in space. We use a somewhat different approach than that taken by Spruit (1981a,b) in order to elucidate some aspects of the approximation. The formalism is applied in Ch. 4 to obtain some general properties of flux tubes in static equilibrium given by the balance of buoyancy and curvature force and of flux tubes in stationary equilibrium for which the drag force exerted by an external flow field is additionally taken into account. In Ch. 5, the central part of this contribution, a formalism is derived which allows the stability analysis of a general static or stationary flux tube equilibrium with arbitrary path in space. Some basic stability properties are derived by application of this formalism to a number of special cases which can be treated analytically, in particular horizontal tubes and symmetric loops. In Ch. 6 the outcome of these calculations, the results of other investigations and further considerations are tentatively combined in a (hopefully) coherent picture which summarizes the author’s present view of magnetic fields in stellar convection zones. This includes also a discussion of the consequences for the dynamo problem. Finally, an outlook on possibilities for extension and continuation of this work is given in Ch. 7.
2. Formation of structures
The work presented here is based on the hypothesis that, similar to the observed magnetic fields in the photosphere of the Sun, most of the magnetic flux within a stellar convection zone is concentrated into filaments or flux tubes <sup>*</sup><sup>*</sup>Strictly spoken, the term ‘flux tube’ refers to a cylindrically shaped bundle of magnetic field lines. However, in what follows we shall often use this term more loosely to generally denote a magnetic filament or flux concentration of arbitrary shape. embedded in nearly field-free plasma. Although at present this conjecture can neither be proven theoretically in a rigorous way nor undubitably demonstrated by observations it is supported by a number of observational indications and theoretical results.
Observationally, Hale’s polarity rules for solar active regions are known to apply with very few exceptions (e.g. Howard, 1989). This can only be so if the magnetic fields in the convection zone are strong enough to avoid a significant deformation by convective velocity fields, i.e. if the magnetic energy density is equal to or larger than the kinetic energy density of convection. Otherwise the magnetic fields would be passively carried around at random and the erupting active region would not show a preferred orientation. Consequently, the magnetic field strength at least must be of the order of the equipartition field strength
$$B_e=v_c\left(4\pi \rho \right)^{1/2}$$
$`\left(2.1\right)`$
($`v_c`$: convective velocity, $`\rho `$: density). Throughout the whole convection zone $`B_e`$ is much larger than the average field strength of about 100 Gauss which can be estimated from the total magnetic flux emerging during one activity cycle. Consequently, if the magnetic flux in the convection zone has at least equipartition field strength it is strongly intermittent and fills only a small fraction of the volume. The observed properties of emerging active regions and the formation of sunspots by accumulation of fragments also indicate that the magnetic flux already is in a concentrated form before it appears at the solar surface (Zwaan, 1978; McIntosh, 1981; Garcia de la Rosa, 1987).
Theoretical arguments for a filamented nature of the magnetic fields in a stellar convection zone are given in the subsequent sections.
2.1 Rayleigh-Taylor instability of a magnetic layer
Beginning with Parker (1975a) a number of arguments has been put forward which support the assertion that the major part of the magnetic flux which is responsible for the solar activity cycle cannot be kept in the superadiabatic parts of the convection zone for times comparable to the cycle period. The flux rather has to be stored below in a region of overshooting convection where a magnetic layer forms which occasionally ejects magnetic flux into the convection zone proper (Spiegel and Weiss, 1980). Since the thickness of such an overshoot layer probably is less than the local pressure scale height of about $`510^4`$ km (Shaviv and Salpeter, 1973; Schmitt et al., 1984; van Ballegooijen, 1982b; Pidatella and Stix, 1986) and thus a lot of toroidal magnetic flux (about $`10^{24}`$ mx) has to be accommodated in a rather small volume the field there is thought to be densely packed (non-filamented) and, because the flow velocities are small compared to the sound speed, in magnetostatic equilibrium.
If the (horizontal or toroidal) magnetic field decreases rapidly enough with height or even drops to zero discontinuously at some level, this equilibrium which is basically a balance between gravity and the gradient of the total (magnetic + gas) pressure becomes unstable with respect to an interchange of more and less magnetized fluid which lowers the potential energy of the system. This is the magnetic Rayleigh-Taylor instability <sup>*</sup><sup>*</sup>In the case of a discontinuous transition between a non-magnetic plasma and a vacuum magnetic field the instability is known as Kruskal-Schwarzschild instability (cf. Cap, 1976, Ch. 11). of a layer of magnetic field directed perpendicular to the local direction of gravity (Gilman, 1970; Cadez, 1974; Acheson and Gibbons, 1978; Acheson, 1979; Parker, 1979b; Schmitt and Rosner, 1983; Hughes, 1985a,b, 1987; Schmitt, 1985; Hughes and Cattaneo, 1987). For this kind of instability the stratification need not necessarily show a density inversion as in the case of the hydrodynamical Rayleigh-Taylor instability. The nonlinear evolution of the instability in the case of a discontinuity of the magnetic field is nicely illustrated by the numerical simulation of Cattaneo and Hughes (1988) who showed that secondary Kelvin-Helmholtz instabilities lead to the formation of intense vortices whose dynamical interaction dominates the dynamics after the first phases of the instability.
Because of the stabilizing effect of magnetic curvature forces linear stability analysis shows that the fastest growing perturbations are those with large wavelength along and small wavelength perpendicular to the equilibrium magnetic field. Consequently, the formation of thin structures at the upper boundary of a sheet of horizontal magnetic field (where the field strength may rapidly decrease with height) is favored. Quantitative results for the case of a stellar convection zone are difficult to obtain since most of the work done so far is restricted to linear analysis under idealized assumptions. The typical fragment sizes of a destabilized flux sheet depend on diffusive effects (thermal, viscous and magnetic) and on the detailed height dependence of field strength and entropy. Another important effect to consider is rotation which may drastically reduce the growth rates or even entirely suppress the instability (cf. Acheson, 1978; 1979; Roberts and Stewartson, 1977).
We shall not discuss here the complex variety of instability mechanisms which grows out of the interaction of differential rotation, magnetic field, stratification and diffusive effects (see e.g. Schmitt and Rosner, 1983). In order to give a crude estimate of typical temporal and spatial scales we consider a case which is thought to represent the typical properties of a magnetic layer in the overshoot region below the solar convection zone: Adiabatic or weakly subadiabatic temperature stratification, sound speed ($`500`$ km$``$s<sup>-1</sup>) large compared to rotational velocity (1.4 km$``$s<sup>-1</sup>) which, in turn, is large compared to the Alfvén velocity ($`60`$ m$``$s<sup>-1</sup> for an equipartition field of about $`10^4`$ Gauss). Under these circumstances the magnetic Rayleigh-Taylor instability evolves near the top of the sheet, i.e. in a region where the field strength, $`B`$, decreases with height faster than density, $`\rho `$, in the form of growing magnetostrophic waves which propagate along the direction of the equilibrium magnetic field (Acheson, 1979; Schmitt, 1985). For a magnetic layer of thickness $`D`$ the fastest growing wave mode is characterized by the wavenumbers $`n`$ (perpendicular to both gravity and magnetic field) and $`k`$ (parallel to the field) which are given by (Acheson, 1979, Ch. 3)
$$n^2D^2=\frac{\pi }{2}C^{1/2}\left[\frac{H_p}{\gamma }\frac{d}{dz}\mathrm{ln}\left(\frac{B}{\rho }\right)\right]$$
$`\left(2.2\right)`$
$$k^2H_p^2=\frac{1}{2}\left[\frac{H_p}{\gamma }\frac{d}{dz}\mathrm{ln}\left(\frac{B}{\rho }\right)\right]$$
$`\left(2.3\right)`$
($`H_p`$: pressure scale height, $`\gamma `$: ratio of specific heats, $`z`$: height coordinate, antiparallel to the direction of gravity). The quantity $`C`$ is given by
$$C=\frac{v_A^2}{2\mathrm{\Omega }\eta }\frac{D^2}{H_p^2}$$
$`\left(2.4\right)`$
($`v_A`$: Alfvén speed, $`\mathrm{\Omega }`$: angular velocity, $`\eta `$: magnetic diffusivity). Using values of $`B=10^4`$ Gauss, $`\mathrm{\Omega }=`$$`2.710^6`$ s<sup>-1</sup>, the surface equatorial rotation rate, $`\eta =10^4`$ cm<sup>2</sup>s<sup>-1</sup>, $`H_p=610^4`$ km and $`D=10^4`$ km (cf. Schmitt et al., 1984) we find $`C210^7`$. Taking the thickness of the layer, $`D`$, as typical length scale for the decrease of $`B/\rho `$ with height, we find
$$\frac{d}{dz}\mathrm{ln}\left(\frac{B}{\rho }\right)D^1$$
$`\left(2.5\right)`$
and from Eqs. (2.2) and (2.3)
$$\begin{array}{cc}\hfill n^1& 610^3D60\text{km}\hfill \\ \hfill k^1& .75H_p4.510^4\text{km}.\hfill \end{array}$$
Thus the fastest growing wave has a longitudinal wavelength of about a scale height and a much smaller transversal scale. With increasing amplitude the wave penetrates into the convection zone proper and its filaments become subject to the convective flows. Furthermore, they are liable to other instabilities and fragmentation processes (see Sec. 2.3). We conclude that the magnetic Rayleigh-Taylor instability of a magnetic layer in an overshoot region leads to magnetic fragments in the convection zone which have a transversal scale of the order of $`100`$ km or less.
The growth time, $`\tau `$, of the instability is given by
$$\begin{array}{ccc}\hfill \tau & \frac{4\mathrm{\Omega }H_p^2}{v_A^2}\left[\frac{H_p}{\gamma }\frac{d}{dz}\mathrm{ln}\left(\frac{B}{\rho }\right)\right]^1\hfill & \\ & & \\ & \frac{4\gamma \mathrm{\Omega }H_pD}{v_A^2}35\text{days}.\hfill & \left(2.6\right)\hfill \end{array}$$
This rather large time scale reflects the stabilizing influence of rotation. The instability cannot be suppressed by the stable, subadiabatic stratification of the overshoot layer since the radiative thermal diffusivity $`\kappa `$ $`210^7`$ cm<sup>2</sup>s<sup>-1</sup> equalizes temperature differences between a structure with $`l=60`$ km diameter and its environment in a timescale $`l^2/\kappa 41`$ days which is comparable to the growth time of the instability given by Eq. (2.6). The flux loss of a magnetic layer in the overshoot region due to the magnetic Rayleigh-Taylor instability is possibly limited by the ‘turbulent diamagnetism’ of the convection zone which is briefly discussed in Sec. 2.2.
The conclusions drawn above are not significantly changed if an appropriate “turbulent” value for the magnetic diffusivity is taken instead of the molecular value $`\eta =10^4`$ cm<sup>2</sup>s<sup>-1</sup>. This is true even if the suppression of motions by the strong magnetic field in the layer is ignored. We take an appropriate “microscale” $`\delta =100`$ km for the motions to be described by the turbulent diffusivity $`\eta _t`$ which is given by
$$\eta _t\mathrm{\hspace{0.33em}0.1}u\left(\delta \right)\delta $$
$`\left(2.7\right)`$
where $`u\left(\delta \right)`$ is the turbulent velocity on the spatial scale $`\delta `$. Assuming a Kolmogorov spectrum with an external scale $`L=10^{10}`$ cm and $`u\left(L\right)=10^4`$ cm$``$s<sup>-1</sup> as typical for the convective flows in the deep convection zone we find
$$u\left(\delta \right)=u\left(L\right)\left(\frac{\delta }{L}\right)^{1/3}\mathrm{\hspace{0.33em}10}^3\text{cm}\text{s}^1$$
$`\left(2.8\right)`$
and using Eq. (2.7) we have $`\eta _t10^9`$ cm<sup>2</sup>s<sup>-1</sup>. Hence, the quantity $`C`$ is decreased by a factor $`10^5`$ (cf. Eq. 2.4) and we see from Eq. (2.2) that the transversal length scale $`n^1`$ is increased by a factor $`10^{5/4}`$ to about 1000 km which confirms a posteriori our choice of the microscale in Eq. (2.7). The fragment size is still much smaller than the local scale height and the typical length scales of the convective motions. While the growth time of the instability given by Eq. (2.6) is not changed, the thermal diffusion timescale is now about 30 years and radiative heating cannot remove the stabilizing effect of a subadiabatic temperature stratification.
In a much more detailed study Schmitt and Rosner (1983) come to essentially the same conclusions. The magnetic Rayleigh-Taylor instability is of “double-diffusive” nature: For molecular magnetic diffusivity we have $`\eta /\kappa 1`$, i.e. the stabilizing effect of the subadiabatic stratification is connected with a larger diffusivity ($`\kappa `$) than the destabilizing magnetic field gradient ($`\eta `$) and instability ensues. On the other hand, if $`\eta =\eta _t`$ we have $`\eta >\kappa `$ and a sufficiently subadiabatic stratification like the one used by Schmitt and Rosner (1983) removes the magnetic instability while an adiabatic stratification as assumed by Acheson (1979) still leads to instability. Consequently, the stability properties of a magnetic layer depend sensitively on the entropy gradient within the overshoot region. Model results for this quantity have been presented, among others, by Shaviv and Salpeter (1973), Schmitt et al. (1984) and Pidatella and Stix (1986). All these authors find that the subadiabaticity is rather small ($`_{ad}10^6\mathrm{}10^7`$) in the overshoot layer.
Schmitt and Rosner (1983) propose the following scenario: During a first phase magnetic flux is accumulated and amplified by dynamo processes but the field strength stays below the equipartition value and the turbulent motions are not strongly affected by the magnetic field. Consequently, turbulent diffusivities are appropriate and, given a sufficiently large subadiabaticity, the configuration is stable. As the field strength increases (e.g. by differential rotation) the turbulent motions are more and more affected by the field and the magnetic diffusivity approaches its molecular value. As we have seen above, this quenches the stabilizing effect of the stratification, the magnetic layer sheet becomes unstable and is ejected into the convection zone in the form of small filaments.
The important point for this work which mainly deals with the magnetic structure within the convection zone proper is that in any case very small structures with diameters between 100 and 1000 km are formed when the instability sets in.
2.2 Flux expulsion
Independent of having entered from below or being generated in place magnetic flux within a convection zone interacts with the convective flows, a situation which the theory of magnetoconvection attempts to describe (see Proctor and Weiss, 1982, for a review). An important result of this interaction is flux expulsion, a process which has first been demonstrated in the kinematical case (passive magnetic field) by Parker (1963). He showed that in an electrically well-conducting plasma with a stationary velocity field a magnetic field is excluded from the regions of closed streamlines. Starting with the work of Weiss (1966) a number of numerical studies have been performed (e.g. Galloway et al., 1978; Weiss, 1981a,b). They showed that in a convecting medium at high magnetic Reynolds number $`R_m=UL/\eta `$ ($`U`$, $`L`$: Velocity and size of the dominating convective cell, $`\eta `$: magnetic diffusivity) permeated by a magnetic field, the magnetic flux is concentrated into filaments between the convection cells. This effect is related to the phenomenon of ‘intermittency’ for magnetic fields in turbulent flow (e.g. Kraichnan, 1976; Orszag and Tang, 1979; Meneguzzi et al., 1981).
In a compressible, stratified fluid the magnetic flux concentrations formed by flux expulsion are found in the convective downflow regions (Nordlund, 1983; 1986; Hurlburt et al., 1984; Hurlburt and Toomre, 1988). Observations demonstrate that the flux expulsion process operates in the solar (sub-)photosphere: On both the granular (Title et al., 1987) and the supergranular scale (the well-known network fields) magnetic flux is predominantly located and concentrated in the downflow regions.
The nonlinear back-reaction of the magnetic field on the convective flows via the Lorentz force limits the flux density which can be achieved by flux expulsion to a value which is roughly given by the equipartition of magnetic and kinetic energy density. This limit may be modified by compressibility, diffusive and thermal effects (e.g. Galloway et al., 1978; Schüssler, 1990). Furthermore, motion is excluded from strong flux concentrations and a kind of “phase separation” between field-free, convecting fluid and magnetic, almost stagnant regions evolves. Such a situation seems to be favored energetically (Parker, 1984): The interference of the magnetic field with the convective energy transport is minimized and the total energy is smaller than for a state with a diffuse field and the same convective energy flux. The exclusion of motion for fields stronger than equipartition suppresses the convective heat exchange between the magnetic structure and its surroundings. Thermal interaction with the environment is reduced to radiative energy transport.
The properties of nonlinear magnetoconvection are much more complicated than can be discussed here (see Proctor and Weiss, 1982). For our purposes it suffices to state that a magnetic field permeating a convecting fluid at high magnetic Reynolds number inevitably is concentrated into structures of about equipartition field strength. In the convection zone the magnetic Reynolds number for the dominating convective flows is everywhere large: It increases from about $`10^3`$ in the upper layers to about $`10^9`$ near the bottom (cf. Stix, 1976). Consequently, flux expulsion is relevant throughout the whole convection zone and we have to expect a concentrated, filamented magnetic field structure.
The well-known formal analogy between the MHD induction equation and the equation which determines the time evolution of vorticity supports the conjecture that the expulsion effect operates also for vorticity and leads to the formation of intense whirls or vortices. An example in cylindrical geometry has been given by Galloway (1978) while Schüssler (1984a) showed that a full analogy between kinematic expulsion of magnetic field and vorticity holds for two-dimensional flow in cartesian geometry. Numerical simulations of turbulence (e.g. McWilliams, 1984) and laboratory experiments in rotating, turbulent fluids (McEwan, 1973; 1976; Hopfinger et al., 1982) have clearly demonstrated vorticity expulsion for rotationally dominated flows (i.e. flows with small Rossby number $`U/(2\mathrm{\Omega }L)`$ with $`U`$, $`L`$: velocity and size of the dominating eddy, $`\mathrm{\Omega }`$: angular velocity). A similar effect occurs in the simulations of solar granulation carried out by Nordlund (1984b, 1985a) who found that narrow granular ‘fingers’ come into rapid rotation. At least in the surface layers of the Sun magnetic field and vorticity are concentrated in the same locations such that magnetic flux concentrations become surrounded by rapidly rotating, descending whirl flows.
In a general context flux expulsion is related to the idea of turbulent diamagnetism which goes back to Zel’dovich (1957) and Spitzer (1957). This means a transport of magnetic field antiparallel to the gradient of turbulent intensity in inhomogeneous turbulence. The effect tends to expel a magnetic field into the boundary parts of a confined turbulent region (like a convection zone). Adopting a two-scale approach turbulent diamagnetism has been studied by Rädler (1968), Moffatt (1983) and Cattaneo et al. (1988). Again, we may conjecture that the effect operates in a similar way for vorticity. For the illustrative case of spatially periodic velocity field in two dimensions with a perpendicular shear flow Schüssler (1984a) applied a mean-field treatment for the kinematic case and showed that vorticity is transported and expelled in the same way as a magnetic field. He also gave an estimate for the effect of this mechanism on the depth dependence of rotation in the lower half of the solar convection zone where rotation is dominant (Rossby number $`0.3`$). A stationary profile of the angular velocity $`\mathrm{\Omega }`$ could be determined by assuming a balance between the ‘negative viscosity’ effect of vorticity expulsion and normal (turbulent) viscosity.
For the solar and stellar convection zones we may conjecture that flux expulsion has two major effects in a stellar convection zone: It leads to a filamentary state of magnetic fields by generating local concentrations of magnetic flux and vorticity and it pushes both magnetic field and vorticity (angular momentum in a rotating system) to the boundaries. This generates a magnetic shear layer at the bottom of the convection zone, a region which is favorable for the operation of a dynamo mechanism (see Sec. 6.4).
2.3 Instabilities and fragmentation of single flux tubes
Having shown some evidence in favor of a concentrated and filamented state of magnetic fields in a stellar convection zone we may ask for the typical size of a flux concentration. The size distribution of magnetic structures depends not only on the initial sizes of the flux tubes injected into the region but also on fragmentation, accumulation and coagulation processes operating in the convection zone itself (Bogdan, 1985; Bogdan and Lerche, 1985). Large-scale convective flows tend to accumulate magnetic flux by way of the flux expulsion mechanism discussed in the preceding section but a number of processes counteracts this tendency to form large structures in the form of single, coherent flux tubes.
A first mechanism to mention is, again, the magnetic Rayleigh-Taylor instability: If a magnetic structure is oriented mainly horizontally, it may be fragmented by this instability if it is larger than the scale given by Eq. (2.2). However, the rather large growth time (Eq. 2.6) must be compared with the time scales of other processes in order to evaluate the relevance of this instability. Of particular importance in this connection are the hydromagnetic interchange instability and a special form of the Kelvin-Helmholtz instability.
The interchange instability is well known from laboratory plasmas (e.g. Krall and Trivelpiece, 1973, Ch. 5; Cap, 1976, Ch. 11). As a simple example consider the interface between a non-magnetic plasma and a vacuum magnetic field where in equilibrium gas pressure and magnetic pressure are equal. This equilibrium is unstable if, looking from the plasma side, the boundary is concave: the potential energy of the system can be diminished by exchanging magnetic and non-magnetic volume elements because this procedure decreases the magnetic tension. Conversely, if the boundary is convex with respect to the plasma, the equilibrium is stable. These considerations apply in a similar way to a non-vacuum magnetic field. For cylindrical or rotationally symmetric configurations the instability is often referred to as fluting or flute instability since the most rapidly growing perturbations are reminiscent to the flutes of classical columns.
The interchange instability is important for magnetic fields in the convection zone in at least three respects. Firstly, it increases the efficiency of magnetic field line reconnection (Parker, 1979a, Ch. 15); secondly, it can lead to fragmentation of vertical magnetic structures in the uppermost (subphotospheric) layers: Since the gas pressure decreases rapidly with height, a vertical flux tube (e.g. a sunspot) flares out with height. Consequently, the interface becomes concave with respect to the non-magnetic plasma and the configuration is liable to the interchange instability (Parker, 1975b; Piddington, 1975). Meyer et al. (1977) showed that the stratification of the fluid (due to gravity) stabilizes large flux tubes (magnetic flux larger than about $`510^{19}`$ mx): Interchanging magnetic and non-magnetic fluid entails lifting dense material above light material which gives a positive contribution to the potential energy. Similarly, small flux tubes (magnetic flux below about $`510^{17}`$ mx) may be stabilized by surrounding whirl flows (Schüssler, 1984b). In this case, the dynamical stability of the angular momentum distribution at the boundary compensates the destabilizing effect of field line curvature.
The third aspect – which is most important for the discussion here – is the interchange instability of deformed magnetic structures in the deep convection zone. Because of the large electrical conductivity and due to hydrodynamic coupling, magnetic structures follow the fluid motions (convection, differential rotation) until the curvature forces have grown strong enough to resist further deformation. However, for the curvature force to come into play a flux tube must be significantly deformed (cf. Sec. 6.2). As a simple example, consider an initially vertical flux tube subject to a horizontal, localized jet-like flow. The tube is deformed to the shape sketched in Fig. 1 and reaches an equilibrium which, in the absence of gravity, is determined by the balance of the hydrodynamic drag force and the magnetic curvature force (cf. Sec. 4.2).
$$\mathrm{̧\backslash vbox\{\; \{\backslash immediate\}\; \{\; \{\backslash immediate\}\; \backslash vboxto142.26378pt\{\; \backslash hboxto163.60326pt\{\; \backslash hfil\}\; \backslash vfil\}\; \}\}}$$
Fig. 1: Sketch of a flux tube under the influence of a horizontal, jet-like velocity field. An equilibrium is reached if the curvature force balances the hydrodynamic drag force. The interface on the upstream side is liable to the interchange instability.
The interface between the flux tube and the surrounding plasma is unstable on the upstream side which faces the flow and in the absence of stabilizing effects the flux tube splits up into smaller fragments. Linear stability analysis in the case of vanishing magnetic diffusivity gives for the growth time, $`\tau `$, of the instability (cf. Cap, 1976)
$$\tau \left(n\right)=\left(\frac{\rho R}{2\mathrm{\Delta }pn}\right)^{1/2}$$
$`\left(2.9\right)`$
($`\mathrm{\Delta }p`$: gas pressure difference at the interface; $`n`$: perturbation wave number, perpendicular to equilibrium magnetic field; $`R`$: radius of curvature of the interface). Since the gas pressure difference is equal to the magnetic pressure of the flux tube, i.e. $`\mathrm{\Delta }p=B^2/8\pi `$, we find from Eq. (2.9)
$$\tau \left(n\right)=\left(v_A\right)^1\left(\frac{R}{n}\right)^{1/2}$$
$`\left(2.10\right)`$
where $`v_A=B/\left(4\pi \rho \right)^{1/2}`$ is the Alfvén velocity. An upper limit for the growth time can be obtained by specifying a lower limit for $`n`$, i.e. $`n\left(2a\right)^1`$ where $`a`$ is the radius of the flux tube. This leads to
$$\tau \frac{\left(2Ra\right)^{1/2}}{v_A}$$
$`\left(2.11\right)`$
Consequently, the growth time is smaller than the travel time of an Alfvén wave over a distance given by the geometric mean of flux tube diameter and its radius of curvature. For a strongly deformed tube with $`Ra`$ we find $`\tau a/v_A`$. As an example consider a large tube, $`a=10^4`$ km, in the lower convection zone of the Sun with $`B=B_e`$, $`v_A=100`$ m$``$s<sup>-1</sup> which is moderately deformed, $`RH_p610^4`$ km, by convective flows or differential rotation. According to Eq. (2.11) we find a growth time of $`\tau 3`$ days. Hence, even a moderately curved structure fragments within a few days.
Since the growth time decreases for smaller perturbation wavelength (cf. Eq. 2.10), equality with the time scale of magnetic diffusion is reached only for very small fragment size. For this size, $`d=1/n_0`$, inhomogeneities are as rapidly smoothed out by diffusion as they are formed due to interchange instability. $`d`$ can be estimated by equating $`\tau (n_0)`$ with the diffusion time $`\tau _d=d^2\eta ^1`$ which gives
$$d=\left(\frac{2R\eta ^2}{v_A^2}\right)^{1/3}.$$
$`\left(2.12\right)`$
For the large flux tube discussed above we find $`d40`$ km even if we use the turbulent diffusivity derived after Eq. (2.8). The growth time for such a structure is only a few hours. Consequently, the interchange instability constitutes a very efficient fragmentation mechanism which leads to splitting of even moderately deformed flux tubes in a time scale of hours to days. The resulting fragment sizes are of the same order of magnitude as those generated by the magnetic Rayleigh-Taylor instability at the upside of a horizontal flux tube but the growth time of the latter is much larger.
Another important fragmentation mechanism is related to the Kelvin-Helmholtz instability (Tsinganos, 1980). Assume a flux tube embedded in an external velocity field as sketched in Fig. 2a. A small perturbation near the stagnation point leads to the flow geometry sketched in Fig. 2b. The centrifugal force due to the curved streamlines near the stagnation point leads to a pressure gradient which causes a local pressure maximum at the interface. Very similar to the interchange instability, this causes growth of the perturbation and fragmentation of the flux tube. This process is nicely demonstrated by laboratory experiments with rising gas bubbles in liquids (see the photographs reproduced in Tsinganos, 1980) and is also visible in a numerical simulation of buoyant, rising flux tubes (Schüssler, 1979).
$$\mathrm{̧\backslash vbox\{\; \{\backslash immediate\}\; \{\; \{\backslash immediate\}\; \backslash vboxto142.26378pt\{\; \backslash hboxto402.60634pt\{\; \backslash hfil\}\; \backslash vfil\}\; \}\}}$$
Fig. 2: Dynamical fragmentation of a flux tube. a: Flux tube embedded in an external flow field. A small perturbation near the stagnation point $`S`$ leads to the situation sketched in b: The the centrifugal force causes a pressure gradient such that a local pressure maximum evolves at the interface. This leads to growth of the perturbation and fragmentation of the tube.
All these considerations lead us to expect strongly fragmented magnetic structures of the size of a few tens of km in the solar convection zone. Converging flows could possibly accumulate them in loose bundles but they are unable to produce large, monolithic structures since we have seen that the flows themselves cause fragmentation (see also Schüssler, 1984b, and the discussion in Sec. 6.1).
However, could possibly a twisted flux tube escape from being fragmented? With an azimuthal component, $`B_\varphi `$, of the magnetic field the internal pressure gradient is changed and, furthermore, perturbations which do not bend magnetic field lines are no longer possible. A sufficiently large $`B_\varphi `$ may suppress all fragmentation processes discussed so far. For the interchange instability the necessary magnitude can be estimated by equating the (stabilizing) tension force due to $`B_\varphi `$ with the (destabilizing) curvature force exerted by the longitudinal field $`B_z`$
$$\frac{B_\varphi ^2}{4\pi a}\frac{B_z^2}{4\pi R}$$
which gives for the ratio of the field components
$$\frac{B_\varphi }{B_z}\left(\frac{a}{R}\right)^{1/2}.$$
$`\left(2.13\right)`$
For the large flux tube discussed above Eq. (2.13) gives a $`B_\varphi /B_z0.4`$, i.e. the azimuthal field must be of the order of the longitudinal field in order to stabilize a moderately deformed flux tube ($`R=H_p`$) with respect to the interchange instability. With such large azimuthal fields, however, kink instabilities become relevant at perturbation wavelengths along the flux tube which are comparable to the diameter of the structure, i.e. the flux tubes buckles, reconnects and evolves rapidly on a dynamical time scale. Presumably this instability either leads directly to fragmentation or it removes most of the twist and leaves the flux tube unprotected against the other instabilities discussed above. A twisted flux tube is kink unstable for longitudinal wavelengths $`\lambda `$ which satisfy the inequality (Cap, 1976, Ch. 11)
$$\frac{\lambda }{a}>\frac{B_z}{B_\varphi }.$$
$`\left(2.14\right)`$
Using this and Eq. (2.13) we find that a flux tube which is stabilized against interchanging by an azimuthal field is kink unstable for
$$\lambda >\left(aR\right)^{1/2}.$$
$`\left(2.15\right)`$
In our example we have $`\lambda >2.410^4`$ km. The time scale for the kink instability is of the same order of magnitude as that of the interchange instability. On the other hand, slightly deformed small flux tubes can be stabilized by a much smaller amount of twist. For $`a=100`$ km and $`R=H_p=610^4`$ km, a ratio $`B_\varphi /B_z=0.04`$ is sufficient. The helical kink instability would set in on a scale of $`\lambda 2500`$ km which is much larger than the flux tube diameter and could possibly, as conjectured by Parker (1979a, Ch. 9.2), saturate in a stable, cork-screw shaped form of the flux tube. In the uppermost layers of the convection zone, such a modest amount of twisting could conceivably be produced by a surrounding whirl flow which itself exerts a stabilizing influence (Schüssler, 1984b).
3. The approximation of slender flux tubes
The considerations presented in the preceding chapter support the view that the magnetic flux in a convection zone consists of an ensemble of thin (diameter $`<100`$ km $`H_p`$), concentrated (at least equipartition field strength) filaments. Such structures allow a simplification of the MHD equations if all quantities do not vary significantly within each cross section and if the spatial scales of variations along the filament are large compared to its diameter. In particular, this approximation requires that the diameter is small compared to pressure scale height, radius of curvature, and to the longitudinal length scales of all dynamical processes (e.g. longitudinal wavelengths, scale of variation of flows). Under these circumstances the global statics and dynamics of a filament (excluding processes like body waves which involve a significant structure within a cross-section, cf. Ferriz-Mas et al., 1989) can be described using truncated Taylor expansions of the lateral variation of all quantities. A truncation of the resulting set of equations at zeroth order leads to the approximation of slender flux tubes (ASF) which involves only the values of the quantities along a representative curve (e.g. the axis of a flux tube with circular cross section) or averages over the cross section.
For vertical, axisymmetric flux tubes with straight axis the ASF has been systematically derived by Roberts and Webb (1978; see also Defouw, 1976). For this case, the general properties of the expansion procedure are discussed by Ferriz-Mas and Schüssler (1989). Spruit (1981a,b) gives a general form of the ASF for flux tubes with curved axis. For the application to extragalactic jets, Achterberg (1982, 1988) has derived a similar approximation which he calls the ‘firehose limit’. In what follows we shall rederive the ASF in a way somewhat different from Spruit’s approach in order to introduce the formalism and notation to be used in the subsequent chapters and hopefully also to elucidate the nature of the approximation.
\̧vbox{ {\immediate} { {\immediate} \vboxto170.71652pt{ \hboxto360.21184pt{ \hfil} \vfil} }}
Fig. 3: Sketch of a space curve r($`l`$) described in each point $`P`$ by the orthogonal unit vectors $`\widehat{𝐥}`$ (tangent), $`\widehat{𝐧}`$ (normal), $`\widehat{𝐛}`$ (binormal, perpendicular to the plane of the paper), and by the radius of curvature, $`R`$.
Consider a space curve whose path is described by the variation of the radius vector $`𝐫=𝐫\left(l\right)`$ with the arc length, $`l`$, along the curve. <sup>*</sup><sup>*</sup>In general, the path and all other quantities depend explicitly on time. We take this into account by writing all spatial derivatives partial but we do not especially indicate the time dependence unless it actually matters. As indicated in Fig. 3 the curve is described at every point, $`P`$, by the triad of unit vectors $`\widehat{𝐥}`$ (tangent), $`\widehat{𝐧}`$ (normal), and $`\widehat{𝐛}`$ (binormal) which are given by
$$\begin{array}{ccc}& \widehat{𝐥}=\frac{𝐫}{l}\hfill & \left(3.1\right)\hfill \\ & & \\ & \widehat{𝐧}=R\frac{\widehat{𝐥}}{l}\hfill & \left(3.2\right)\hfill \\ & & \\ & \widehat{𝐛}=\widehat{𝐥}\times \widehat{𝐧}\hfill & \left(3.3\right)\hfill \end{array}$$
The radius of curvature is given by
$$R=\left|\frac{\widehat{𝐥}}{l}\right|^1.$$
$`\left(3.4\right)`$
The vectors $`\widehat{𝐥}`$ and $`\widehat{𝐧}`$ define the local plane of the curve (plane of the paper in Fig. 3) which contains the local center of curvature, $`C`$. The change of the local plane of the curve is described by the derivative of the binormal unit vector which defines the radius of torsion, $`R_t`$:
$$\frac{\widehat{𝐛}}{l}=\widehat{𝐥}\times \frac{\widehat{𝐧}}{l}=R_t^1\widehat{𝐧}.$$
$`\left(3.5\right)`$
In the ASF we consider a flux tube as a coherently moving bundle of magnetic field lines whose paths can be represented by one single space curve $`𝐫(l,t)`$. For such a description to be valid, the variation of tangent, $`\widehat{𝐥}`$, radius of curvature, $`R`$, and all physical quantities within the cross section of the tube in the plane perpendicular to $`\widehat{𝐥}`$ has to be sufficiently small. Since it allows a better readable presentation, in what follows we assume a circular cross section of the flux tube and take its axis as representative space curve and as origin of Taylor expansions. However, the following considerations are valid for any shape of the cross section and can easily be generalized as long as the flux tube remains sufficiently thin in any direction perpendicular to its local tangent. Note in particular that we do not (and in general are not allowed to) assume axial symmetry.
The natural coordinates within any cross section of the tube are defined along the normal and binormal directions with the axis as origin. Using $`\xi `$ for the coordinate in the direction of the local normal, $`\widehat{𝐧}`$, we may write for the differences $`\mathrm{\Delta }R`$ and $`|\mathrm{\Delta }\widehat{𝐥}|`$ between axis ($`\xi =0`$) and boundary ($`\xi =a`$) of the flux tube in the normal direction
$$\begin{array}{ccc}\hfill \frac{\left|\mathrm{\Delta }R\right|}{R}& =\left|\frac{R}{\xi }\right|_{\xi =0}\frac{a}{R}\hfill & \left(3.6\right)\hfill \\ & & \\ \hfill \left|\mathrm{\Delta }\widehat{𝐥}\right|& =\left|\frac{\widehat{𝐥}}{\xi }\right|_{\xi =0}a.\hfill & \left(3.7\right)\hfill \end{array}$$
The quantities on the left of Eqs. (3.6) and (3.7), respectively, can be made arbitrarily small if $`a`$ is chosen small enough, in particular if the diameter of the flux tube is everywhere small compared to the local radius of curvature of the axis. Similar relations are valid for the variation in the binormal direction and for the variation of torsion which are satisfied also by a sufficiently thin flux tube.
For the (zeroth order) ASF we consider the values of the variables (and possibly their first spatial derivatives, see below) on the axis of the tube. The combined equations of induction and continuity (Walén equation for zero resistivity) and the equation of motion for the fluid within the flux tube in Lagrangian form read
$$\begin{array}{ccc}& \frac{d}{dt}\left(\frac{𝐁}{\rho }\right)=\left(\frac{𝐁}{\rho }\right)𝐮\hfill & \left(3.8\right)\hfill \\ & & \\ & \rho \frac{d𝐮}{dt}=p+\rho 𝐠+\frac{1}{4\pi }\left(\times 𝐁\right)\times 𝐁+𝐅_D\hfill & \left(3.9\right)\hfill \end{array}$$
where u denotes fluid velocity, $`\rho `$ density, $`p`$ gas pressure, g gravitational acceleration, and B the magnetic flux density vector. $`𝐅_D`$ is the drag force which results from the motion of the flux tube relative to the surrounding, non-magnetic fluid. For a sufficiently thin tube we may write for the flux density on the axis (the representative space curve)
$$𝐁\left(l\right)=B\left(l\right)\widehat{𝐥}$$
$`\left(3.10\right)`$
and Eq. (3.8) can be rewritten
$$\frac{B}{\rho }\frac{d\widehat{𝐥}}{dt}+\widehat{𝐥}\frac{d}{dt}\left(\frac{B}{\rho }\right)=\frac{B}{\rho }\left(\widehat{𝐥}\right)𝐮\frac{B}{\rho }\frac{𝐮}{l}$$
$`\left(3.11\right)`$
By scalar multiplication with $`\widehat{𝐥}`$ and noting that a unit vector is perpendicular to its derivative we find from Eq. (3.11)
$$\frac{d}{dt}\left(\frac{B}{\rho }\right)=\frac{B}{\rho }\widehat{𝐥}\frac{𝐮}{l}=\frac{B}{\rho }\left(\frac{𝐮\widehat{𝐥}}{l}𝐮\frac{\widehat{𝐥}}{l}\right).$$
$`\left(3.12\right)`$
We write $`𝐮\widehat{𝐥}u_l,𝐮\widehat{𝐧}u_n`$ and use Eq. (3.2) to obtain
$$\frac{d}{dt}\left(\frac{\rho }{B}\right)=\frac{\rho }{B}\left(\frac{u_n}{R}\frac{u_l}{l}\right).$$
$`\left(3.13\right)`$
Eq. (3.13) represents the ASF form of the Walén equation. Multiplication of Eq. (3.11) with $`\widehat{𝐧}`$ and $`\widehat{𝐛}`$, respectively, gives the normal and binormal components of the time derivative of the tangent vector which describes the change of the flux tube path in time, namely
$$\begin{array}{ccc}& \widehat{𝐧}\frac{d\widehat{𝐥}}{dt}=\frac{u_n}{l}+\frac{u_l}{R}+\frac{u_b}{R_t}\hfill & \left(3.14\right)\hfill \\ & & \\ & \widehat{𝐛}\frac{d\widehat{𝐥}}{dt}=\frac{u_b}{l}\frac{u_n}{R_t}\hfill & \left(3.15\right)\hfill \end{array}$$
where we have used the general relation $`\widehat{𝐧}=\widehat{𝐥}\times \widehat{𝐛}`$ and defined $`𝐮\widehat{𝐛}u_b`$.
The equation of motion requires somewhat more consideration. Let us first write down the Lorentz force on the axis using Eq. (3.10)
$$\begin{array}{ccc}\hfill \frac{1}{4\pi }\left(\times 𝐁\right)\times 𝐁& =\frac{1}{4\pi }\left[\left(B\right)\times \widehat{𝐥}+B\times \widehat{𝐥}\right]\times B\widehat{𝐥}\hfill & \\ & & \\ & =\frac{1}{4\pi }\left[B\left(B\times \widehat{𝐥}\right)\times \widehat{𝐥}+B^2\left(\times \widehat{𝐥}\right)\times \widehat{𝐥}\right]\hfill & \left(3.16\right)\hfill \end{array}$$
Using $`\left(\times \widehat{𝐥}\right)\times \widehat{𝐥}=R^1\widehat{𝐧}`$ (e.g. Smirnov, 1968, Ch. V) and
$$\left(B\times \widehat{𝐥}\right)\times \widehat{𝐥}=B+\widehat{𝐥}\left(\widehat{𝐥}B\right)\left(B\right)_{}$$
$`\left(3.17\right)`$
where $`\left(B\right)_{}`$ denotes the projection of the gradient on the plane perpendicular to the tangential direction (i.e. the cross section of the tube) we find for the Lorentz force
$$𝐅_L\frac{1}{4\pi }\left(\times 𝐁\right)\times 𝐁=\left[\left(\frac{B^2}{8\pi }\right)\right]_{}+\frac{B^2}{4\pi R}\widehat{𝐧}.$$
$`\left(3.18\right)`$
The projections of $`𝐅_L`$ on the three directions defined by the triad of unit vectors are given by
$$\begin{array}{ccc}& 𝐅_L\widehat{𝐥}=\mathrm{\hspace{0.33em}0}\hfill & \left(3.19\text{a}\right)\hfill \\ & & \\ & 𝐅_L\widehat{𝐧}=\frac{}{n}\left(\frac{B^2}{8\pi }\right)+\frac{B^2}{4\pi R}\hfill & \left(3.19\text{b}\right)\hfill \\ & & \\ & 𝐅_L\widehat{𝐛}=\frac{}{b}\left(\frac{B^2}{8\pi }\right).\hfill & \left(3.19\text{c}\right)\hfill \end{array}$$
Here we have defined $`\widehat{𝐧}/n`$ and $`\widehat{𝐛}/b`$. We recognize the two familiar constituents of the Lorentz force, i.e. the magnetic pressure force in the plane perpendicular to the axis and the curvature force in the normal direction. There is no curvature force in the binormal direction and no magnetic force at all in the tangential direction. Note that the derivatives in the normal and binormal directions generally do not vanish since, in contrast to the ASF for vertical flux tubes, a curved flux tube cannot assumed to be axisymmetric.
Using the expressions derived for the Lorentz force in Eq. (3.19) we write for the components of the equation of motion (Eq. 3.9)
$$\begin{array}{ccc}& \rho \frac{d𝐮}{dt}\widehat{𝐥}=\frac{p}{l}+\rho 𝐠\widehat{𝐥}\hfill & \left(3.20\right)\hfill \\ & & \\ & \rho \frac{d𝐮}{dt}\widehat{𝐧}=\frac{}{n}\left(p+\frac{B^2}{8\pi }\right)+\rho 𝐠\widehat{𝐧}+\frac{B^2}{4\pi R}+𝐅_D\widehat{𝐧}\hfill & \left(3.21\right)\hfill \\ & & \\ & \rho \frac{d𝐮}{dt}\widehat{𝐛}=\frac{}{b}(p+\frac{B^2}{8\pi })+\rho 𝐠\widehat{𝐛}++𝐅_D\widehat{𝐛}.\hfill & \left(3.22\right)\hfill \end{array}$$
The component along the flux tube (Eq. (3.20)) is already in a suitable form for the ASF. The other components contain derivatives in the normal and binormal directions which have to be determined by considering the external fluid. This is achieved by assuming that at the interface between flux tube and its environment instantaneous pressure equilibrium (more precisely, continuity of normal stress) is maintained permanently, viz.
$$p+\frac{B^2}{8\pi }=p_e.$$
$`\left(3.23\right)`$
Since the time required to establish pressure equilibrium is of the order of the travel time of a fast magneto-acoustic wave across the tube given by $`2a/\sqrt{v_A^2+v_S^2}`$ ($`v_A`$: Alfvén speed, $`v_S`$: sound speed) it can be made arbitrarily small compared to any other dynamical time scale of the system if the flux tube is sufficiently thin.
Let us first assume a straight flux tube embedded in a static external fluid of constant pressure $`p_e`$, without gravity. We see from Eq. (3.23) that under these conditions pressure equilibrium entails $`p+B^2/8\pi =\text{const.}`$ in each cross section which is identical with the condition for static equilibrium of the internal fluid since all terms vanish on the r.h.s. of Eqs. (3.21) and (3.22) except the derivatives of total pressure.
In the case of a curved flux tube, non-vanishing gravity and stratification of the external fluid, pressure equilibrium at the interface and static equilibrium of the fluid in the interior of the flux tube generally cannot be reached simultaneously and a lateral acceleration of the fluid in the tube results. Since we deal with the zeroth-order ASF, only the terms involving first derivatives are retained in the Taylor expansion of total pressure within the cross section which is inserted into Eqs. (3.21) and (3.22). On the other hand, these derivatives are already fixed by Eq. (3.23) since a linear profile of total pressure along the normal and binormal directions is determined by the values of $`p_e`$ at the intersections with the boundary of the flux tube. We can therefore formally insert Eq. (3.23) in Eqs. (3.21) and (3.22). Note that this procedure is legitimate only for linear pressure profiles, i.e. if the requirements for the validity of the ASF are met. If we take the external stratification to be hydrostatic (assuming all motions to be far subsonic), i.e.
$$p_e=\rho _e𝐠$$
$`\left(3.24\right)`$
($`\rho _e`$: density of the external fluid), we obtain for the three components of the ASF form of the equation of motion <sup>*</sup><sup>*</sup>The transversal force (per unit length along the tube) can also be obtained by assuming static equilibrium within the flux tube (in the comoving frame) and integrating the resulting total pressure difference over the circumference of the cross section. For a thin tube the result is identical to Eqs. (3.26/27). In contrast to the procedure described in the text, this method is as well applicable for non-thin tubes as long as the assumption of internal hydrostatic equilibrium is valid. It has been used to calculate the buoyancy force on a circular flux tube of arbitrary radius (Schüssler, 1979).
$$\begin{array}{ccc}& \rho \frac{d𝐮}{dt}\widehat{𝐥}=\frac{p}{l}+\rho 𝐠\widehat{𝐥}\hfill & \left(3.20\right)=\left(3.25\right)\hfill \\ & & \\ & \rho \frac{d𝐮}{dt}\widehat{𝐧}=\left(\rho \rho _e\right)𝐠\widehat{𝐧}+\frac{B^2}{4\pi R}+𝐅_D\widehat{𝐧}\hfill & \left(3.26\right)\hfill \\ & & \\ & \rho \frac{d𝐮}{dt}\widehat{𝐛}=\left(\rho \rho _e\right)𝐠\widehat{𝐛}+𝐅_D\widehat{𝐛}\hfill & \left(3.27\right)\hfill \end{array}$$
The first terms on the r.h.s. of Eqs. (3.26/27) represent the components of the buoyancy force which are proportional to the density difference between internal and external fluid. The magnetic field enters explicitly only by way of the curvature force in Eq. (3.26) and in the pressure balance condition (Eq. 3.23).
Let us now discuss the drag term, $`𝐅_D`$, which has been introduced to describe the dynamical effect of a motion of the flux tube relative to the surrounding fluid. Spruit (1981a,b) considered impulsive motions of a flux tube in an initially static environment in which case this effect can be described by a larger effective inertia of the tube with respect to perpendicular motions. This is introduced into the equations by changing $`\rho \rho +\rho _e`$ on the l.h.s. of Eqs. (3.26) and (3.27). In our case we wish to include arbitrary flow fields around the flux tube (convection, differential rotation, dynamical motion of the tube itself). Their effect on the tube has to be described explicitly by considering an aerodynamic drag force (e.g. Parker, 1975a; Schüssler 1977; 1979; Moreno-Insertis, 1983; 1986; Chou and Fisher, 1989). Since the drag force has its origin in pressure differences between the upstream and the downstream sides of the interface between flux tube and its environment, for the ASF there is no component along the tube ($`𝐅_D\widehat{𝐥}=0`$). In the simplest case we may use the expression for the drag force acting on a rigid circular cylinder of radius $`a`$ (cf. Batchelor, 1967), viz.
$$\left(\pi a^2\right)𝐅_D=C_D\rho _eav_{}^2\widehat{𝐤}.$$
$`\left(3.28\right)`$
$`C_D`$ is the (dimensionless) drag coefficient and $`v_{}^2`$ is given by
$$v_{}^2=\left(𝐯\widehat{𝐤}\right)^2=\left[𝐯\widehat{𝐥}\left(𝐯\widehat{𝐥}\right)\right]^2$$
$`\left(3.29\right)`$
where $`𝐯=𝐯_e𝐮`$ is the relative velocity between the flux tube and the surrounding fluid moving with velocity $`𝐯_e`$, and $`\widehat{𝐤}`$ is the unit vector in the direction of the component of $`𝐯`$ perpendicular to the tube. Eq. (3.28) has been derived for laminar flows in which case wind tunnel measurements give $`C_D1`$ for a wide range ($`10^2<Re<10^5`$) of the hydrodynamic Reynolds number $`Re=v_{}a/\nu `$ ($`\nu `$: kinematic viscosity). For turbulent flows the effective Reynolds number of the flow is determined by the turbulent viscosity which can be orders of magnitude larger than the molecular value. Similar to the discussion of the turbulent magnetic diffusivity in the preceding chapter, the turbulent viscosity must be determined taking into account the spatial scale of the flow to be described. For $`v_{}`$ given by large-scale convection on spatial scale $`L`$ and a flux tube of radius $`a`$, the relevant turbulent viscosity is of the order of $`0.1u\left(\delta \right)\delta `$ where $`u\left(\delta \right)`$ is the turbulent velocity on a scale $`\delta `$ which is somewhat smaller than $`a`$. For $`\delta =a/10`$ and using Eq. (2.8) we find for the Reynolds number calculated with turbulent viscosity: $`Re=100\left(L/\delta \right)^{1/3}`$. With $`a=100`$ km and $`L=10^4`$ km we finally estimate $`Re10^3`$. Hence, the effective Reynolds number for turbulent flow is such that we could hope to stay in the range of validity of Eq. (3.28) with $`C_D1`$ (see also the detailed discussion by Moreno-Insertis, 1984).
An enhanced inertia for perpendicular motions as introduced by Spruit (1981a,b) is not used here since it cannot be easily specified for turbulent flows which we expect in a convection zone. This may introduce errors if impulsive perpendicular motions like those connected with transversal tube waves are relevant.
For some applications it is more convenient to write the inertial terms in Eqs. (3.25-3.27) in the form of Lagrangian derivatives of the velocity components. <sup>*</sup><sup>*</sup>Note that there is an error in the expressions given by Chou and Fisher (1989) who treated a plane flux tube ($`R_t\mathrm{},u_b=0`$). The terms involving the derivative of the normal velocity, $`u_n/l`$, are missing in their equations (1) and (2). For example, the longitudinal component of the inertial term in Eq. (3.25) can be rewritten using Eqs. (3.14/15) in the form
$$\begin{array}{ccc}\hfill \rho \frac{d𝐮}{dt}\widehat{𝐥}& =\rho \frac{d𝐮\widehat{𝐥}}{dt}\rho 𝐮\frac{d\widehat{𝐥}}{dt}=\hfill & \\ & & \\ & =\rho \frac{du_l}{dt}\rho \left(u_n\frac{u_n}{l}+\frac{u_nu_l}{R}+\frac{u_nu_b}{R_t}\right)\rho \left(u_b\frac{u_b}{l}\frac{u_bu_n}{R_t}\right).\hfill & \left(3.30\right)\hfill \end{array}$$
Similar expressions can be derived for the other components of the inertial force. The set of equations for the ASF derived so far, i.e. the equations of motion, Eqs. (3.25-3.27), continuity, Eq. (3.13), flux tube shape, Eqs. (3.14/15), and instantaneous pressure balance, Eq. (3.23), are complemented by the condition of magnetic flux conservation along the tube, namely
$$AB=\mathrm{\Phi }_{mag}=\text{const.}$$
$`\left(3.31\right)`$
where $`A`$ is the cross-sectional area. For a flux tube with circular cross section of radius $`a`$ we have $`A=\pi a^2`$. The form of the energy equation and the equation of state is determined by the particular problem to be treated. Since we shall only consider adiabatic changes we do not specify more complicated forms of the energy equation here. In most cases the derivation of the appropriate ASF form is straightforward.
4. Flux tubes in equilibrium
The comparatively long lifetime and slow evolution of large solar active regions after a much shorter phase of flux eruption and dynamical evolution give rise to the conjecture that the magnetic structures in the convection zone reach a static or stationary (i.e. with a surrounding flow) equilibrium which is characterized by a time-independent shape of the tube and hydrostatic equilibrium along its longitudinal direction. Trivial examples of static equilibria (for constant direction of gravity, g) are:
a) a horizontal ($`𝐠\widehat{𝐥}=0`$), straight flux tube with constant pressure and density ($`\rho =\rho _e`$), and
b) a vertical ($`𝐠\widehat{𝐧}=0`$), straight flux tube in hydrostatic equilibrium
($`dp/dz=\rho g;z`$: vertical coordinate). An example for a stationary equilibrium is a straight, horizontal flux tube with a density difference, $`\rho \rho _e`$, such that the resulting buoyancy force compensates the drag force due to a constant vertical velocity, $`𝐯_e`$, in the exterior. The force balance, Eq. (3.26), can be determined using Eq. (3.28) and gives
$$1\frac{\rho }{\rho _e}=\frac{C_Dv_e^2}{\pi ag}\text{sgn(}𝐠𝐯_e)$$
$`\left(4.1\right)`$
with $`v_e\left|𝐯_e\right|`$ and $`g\left|𝐠\right|`$. For a downflow, i.e. $`\text{sgn(}𝐠𝐯_e)=+1`$, the density of the fluid in the tube has to be smaller than that of the surroundings and the resulting upward directed buoyancy force is balanced by the drag force. In the case of an upflow the buoyancy force is directed downwards.
In general, we expect more complicated shapes of equilibrium flux tubes in the solar convection zone. For example, van Ballegooijen (1982a) calculated static and stationary equilibrium solutions for flux tubes forming loops which are “anchored” in a horizontal flux system below the convection zone. Anton (1984) determined equilibrium flux tubes which reside completely within the convection zone for a variety of internal and external temperature stratifications. We shall not perform detailed calculations here but rather derive some general properties of equilibrium flux tubes and give a few illustrative examples.
4.1 Static equilibrium
The equations describing the static equilibrium of a flux tube in a hydrostatically stratified environment are obtained by setting the inertial and the drag terms to zero in Eqs. (3.25-3.27). In the direction of the binormal the equilibrium condition reads
$$\left(\rho \rho _e\right)𝐠\widehat{𝐛}=\mathrm{\hspace{0.33em}0}.$$
$`\left(4.2\right)`$
Unless we have $`𝐠\widehat{𝐛}=0`$ this equation can only be satisfied if the density difference between flux tube and exterior vanishes everywhere along the flux tube. From the equilibrium condition in the normal direction,
$$\left(\rho \rho _e\right)𝐠\widehat{𝐧}+\frac{B^2}{4\pi R}=\mathrm{\hspace{0.33em}0},$$
$`\left(4.3\right)`$
we see that in this case the curvature force vanishes, $`R\mathrm{}`$, i.e. the tube has to be straight. Hydrostatics of the environment, Eq. (3.24), and along the tube, Eq. (3.25), together with $`\rho =\rho _e`$ entail
$$\frac{}{l}\left(pp_e\right)=0$$
$`\left(4.4\right)`$
and from Eq. (3.23) we find
$$\frac{}{l}\frac{B^2}{8\pi }=\mathrm{\hspace{0.33em}0}.$$
$`\left(4.5\right)`$
Consequently, the magnetic field strength is constant. Using the perfect gas law and assuming the molecular weight to be the same inside and outside the flux tube we find from Eq. (3.23) for the ratio of internal temperature, $`T`$, and external temperature, $`T_e`$:
$$\frac{T}{T_e}=\mathrm{\hspace{0.33em}1}\frac{B^2}{8\pi p_e}.$$
$`\left(4.6\right)`$
For a constant direction of gravity, a horizontal flux tube with no variation of $`p_e`$ in its longitudinal direction, and a subadiabatic stratification of the external medium such a temperature reduction might be achieved by an adiabatically expanding, rising flux tube. For an oblique tube with a concomitant variation of $`p_e`$ along its length we can hardly imagine a thermodynamic process which precisely leads to a temperature variation along the tube as prescribed by Eq. (4.6).
In a real star the direction of g varies spatially since gravity is directed towards the center. In this case, a longitudinally uniform flux tube assumes circular (toroidal) shape with finite $`R`$ and therefore we have $`\rho \rho _e`$. Thus in practice straight flux tubes with $`\rho =\rho _e`$ are irrelevant and we can conclude from Eq. (4.2) that $`𝐠\widehat{𝐛}=0`$ is necessary, i.e. static flux tubes lie in planes which contain the vector of gravitational acceleration. For the spherical geometry of a star we may state:
In a spherical star with radial gravity magnetic flux tubes in static equilibrium lie in planes which contain the center of the star. The singular case of straight tubes without density contrast is of no practical importance.
For example, a toroidal flux tube in a plane parallel to but outside the equatorial plane does not fulfill this condition. It cannot find a static equilibrium since the component of the buoyancy force perpendicular to the plane of the tube is not balanced and leads to a poleward drift (Pneuman and Raadu, 1972; Spruit and van Ballegooijen, 1982).
Another important property of static flux tubes is obtained by rewriting the Lorentz force in its familiar form
$$\frac{1}{4\pi }\left(\times 𝐁\right)\times 𝐁=\left(\frac{B^2}{8\pi }\right)+\frac{1}{4\pi }𝐁𝐁.$$
$`\left(4.7\right)`$
We may use this to obtain the static form of the equation of motion, Eq. (3.9), viz.
$$\left(p+\frac{B^2}{8\pi }\right)+\rho 𝐠+\frac{1}{4\pi }𝐁𝐁=\mathrm{\hspace{0.33em}0}$$
$`\left(4.8\right)`$
and with Eqs. (3.23/24) we get
$$\left(\rho \rho _e\right)𝐠+\frac{1}{4\pi }𝐁𝐁=\mathrm{\hspace{0.33em}0}.$$
$`\left(4.9\right)`$
We now consider the general case of a spatially varying direction of gravity and denote by $`\widehat{𝐠}`$ the unit vector in the direction of local gravitational acceleration. We define $`\widehat{𝐡}`$ as unit vector perpendicular to $`\widehat{𝐠}`$ within the local plane of the flux tube (spanned by $`\widehat{𝐧}`$ and $`\widehat{𝐥}`$). Since $`𝐠\widehat{𝐛}=0`$ for a static tube we have
$$\widehat{𝐡}=\widehat{𝐠}\times \widehat{𝐛}.$$
$`\left(4.10\right)`$
$`\widehat{𝐡}`$ defines the local horizontal direction. We multiply Eq. (4.9) by $`\widehat{𝐡}`$, use Eq. (3.10), and find
$$\widehat{𝐡}\left(\widehat{𝐥}\right)B\widehat{𝐥}=\mathrm{\hspace{0.33em}0}.$$
$`\left(4.11\right)`$
This may be written as
$$0=\widehat{𝐡}\frac{B\widehat{𝐥}}{l}=\frac{}{l}B\left(\widehat{𝐡}\widehat{𝐥}\right)B\widehat{𝐥}\frac{\widehat{𝐡}}{l}.$$
$`\left(4.12\right)`$
Using Eqs. (4.10), (3.3) and (3.5) we find
$$\begin{array}{cc}\hfill \frac{\widehat{𝐡}}{l}& =\frac{\widehat{𝐠}}{l}\times \widehat{𝐛}+\widehat{𝐠}\times \frac{\widehat{𝐛}}{l}=\frac{\widehat{𝐠}}{l}\times \left(\widehat{𝐥}\times \widehat{𝐧}\right)+R_t^1\widehat{𝐠}\times \widehat{𝐧}=\hfill \\ & \\ & =\widehat{𝐥}\left(\frac{\widehat{𝐠}}{l}\widehat{𝐧}\right)\widehat{𝐧}\left(\frac{\widehat{𝐠}}{l}\widehat{𝐥}\right)+R_t^1\widehat{𝐠}\times \widehat{𝐧}.\hfill \end{array}$$
Consequently, we have
$$\widehat{𝐥}\frac{\widehat{𝐡}}{l}=\frac{\widehat{𝐠}}{l}\widehat{𝐧}$$
and defining $`B_hB\left(\widehat{𝐡}\widehat{𝐥}\right)=𝐁\widehat{𝐡}`$ we write Eq. (4.12) as
$$\frac{B_h}{l}=B\left(\frac{\widehat{𝐠}}{l}\widehat{𝐧}\right).$$
$`\left(4.13\right)`$
This equation couples the variation of the strength of the field component in the local horizontal direction to the variation of the direction of gravity along the flux tube. In the case of a constant direction of gravitation along the tube we find from Eq. (4.13)
$$\frac{B_h}{l}=\mathrm{\hspace{0.33em}0},$$
$`\left(4.14\right)`$
i.e. the component of the magnetic field perpendicular to gravity is constant in static equilibrium (see also Parker, 1979, Sec. 8.6, and van Ballegooijen, 1982a).
As we have seen above, static flux tubes in a sphere with radially directed gravity lie in planes which contain the center. If we introduce polar coordinates $`(r,\phi )`$ within such a plane, we have $`\widehat{𝐠}=(1,\mathrm{\hspace{0.17em}0})`$, $`\widehat{𝐡}=(0,\mathrm{\hspace{0.17em}1})`$, $`\widehat{𝐥}=(l_r,l_\phi )`$, and
$$\frac{\widehat{𝐠}}{l}=\widehat{𝐥}\widehat{𝐠}=(0,\frac{l_\phi }{r}).$$
Hence, we find from Eq. (4.13) with $`B_h=B\left(\widehat{𝐡}\widehat{𝐥}\right)=Bl_\phi B_\phi `$ :
$$\frac{B_\phi }{l}=B\frac{l_\phi n_\phi }{r}=B_\phi \frac{n_\phi }{r}.$$
$`\left(4.15\right)`$
For $`r\mathrm{}`$ this passes over into Eq. (4.14). We may use Eq. (4.15) to estimate the difference between the spherical and the plane-parallel case. If we denote by $`\delta B_\phi `$ the change of $`B_\phi `$ over a length interval $`\delta L`$ along the tube an upper limit for this quantity given by Eq. (4.15) is
$$\frac{\delta B_\phi }{B_\phi }\frac{\delta L}{r}.$$
$`\left(4.16\right)`$
As an example, for a flux tube extending nearly vertically through the whole depth of the convection zone we have $`\delta L=210^5`$ km and $`r=510^5`$ km. Consequently, the change of $`B_\phi `$ between top and bottom of the convection zone is smaller than $`0.4B_\phi `$. If this flux tube is anchored in a toroidal flux system near the bottom of the convection zone with $`B_\phi B_e10^4`$ Gauss, the toroidal field strength at the top (in the photosphere) cannot be smaller than $`6000`$ Gauss for a tube in static equilibrium. Since photospheric magnetic fields and sunspots are basically vertical with very small net inclination of the magnetic structures as a whole (the mean horizontal field component is less than 100 Gauss) this excludes static flux tubes rooted in a toroidal equipartition field deep in the convection zone as models for sunspots and solar active regions. Since flux expulsion always tends to establish equipartition field strengths which are much larger than 100 Gauss the concept of large-scale static magnetic structures in the convection zone has to be abandoned.
In spite of these pessimistic remarks we shall continue the discussion of static equilibrium in the remainder of this section because this case is well suited to demonstrate the mathematical methods which are used to calculate flux tube equilibria in practice. Furthermore, static equilibrium is only excluded on a large scale for flux tubes extending through the whole convection zone but it might well be locally a reasonable approximation, for instance for nearly vertical flux tubes in the photosphere. For such a small region we neglect the spherical geometry of the star and assume a constant direction of gravity. The subsequent considerations follow the approach first proposed by van Ballegooijen (1982a; see also Parker, 1975c).
Consider cartesian coordinates $`(x,z)`$ in a plane which contains the vector of gravitational acceleration, $`𝐠=(0,g)`$. If the external medium is in plane-parallel hydrostatic equilibrium the path of a static flux tube is given by a curve $`z=z(x)`$ which satisfies Eq. (4.3). This situation is sketched in Fig. 4. Hydrostatic equilibrium along the flux tube is described by Eq. (3.25) which gives <sup>*</sup><sup>*</sup>Since we deal with time-independent quantities we may write non-partial derivatives.
$$\frac{dp}{dl}+\rho 𝐠\widehat{𝐥}=\mathrm{\hspace{0.33em}0}.$$
$`\left(4.17\right)`$
Uniformity of the horizontal field component entails, by virtue of Eq. (4.14)
$$B_x=\text{const.}$$
$`\left(4.18\right)`$
Since we may measure the arc length in both directions along the tube, we remove this ambiguity by requiring $`\widehat{𝐥}\times \widehat{𝐱}>0`$ where $`\widehat{𝐱}`$ is the unit vector in $`x`$-direction. If the shape of the tube is such that it has a vertical tangent somewhere and turns backwards with respect to $`x`$ this part of the tube has to be treated separately in the same way as described below after transforming $`\widehat{𝐥}\widehat{𝐥}`$. We use the angle $`\gamma \left(l\right)`$ between the local tangent of the flux tube and the positive $`x`$-axis defined by
$$\frac{dz}{dx}=\frac{B_z}{B_x}=\mathrm{tan}\gamma ;\frac{d}{dl}=\mathrm{sin}\gamma \frac{d}{dz}$$
$`\left(4.19\right)`$
to write
$$\begin{array}{ccc}\hfill \widehat{𝐥}& =(\mathrm{cos}\gamma ,\mathrm{sin}\gamma )\hfill & \left(4.20\right)\hfill \\ & & \\ \hfill \widehat{𝐧}& =R(\mathrm{sin}\gamma ,\mathrm{cos}\gamma )\frac{d\gamma }{dl}\hfill & \left(4.21\right)\hfill \\ & & \\ \hfill R^1& =\left|\frac{d\gamma }{dl}\right|=\left|\frac{d\mathrm{cos}\gamma }{dz}\right|=\left(\frac{d\mathrm{cos}\gamma }{dz}\right)\text{sgn}\left(\frac{d\gamma }{dl}\right).\hfill & \left(4.22\right)\hfill \end{array}$$
$$\mathrm{̧\backslash vbox\{\; \{\backslash immediate\}\; \{\; \{\backslash immediate\}\; \backslash vboxto199.16928pt\{\; \backslash hboxto239.00311pt\{\; \backslash hfil\}\; \backslash vfil\}\; \}\}}$$
Fig. 4: Flux tube in cartesian geometry with gravity directed downward. $`\gamma (l)`$ is the angle between the flux tube (direction of increasing arc length) and the $`x`$-axis.
Since $`𝐠\widehat{𝐥}=g\mathrm{sin}\gamma `$ we can rewrite Eq. (4.17) using Eq. (4.19)
$$\frac{dp}{dz}=\rho g$$
$`\left(4.23\right)`$
and in the external medium we have
$$\frac{dp_e}{dz}=\rho _eg.$$
$`\left(4.24\right)`$
Consequently, for given external plane-parallel stratification and given internal temperature profile, $`T\left(z\right)`$, pressure and density within the flux tube depend only on $`z`$ and can be determined without prior knowledge of its path. This applies also to the field strength, viz.
$$\frac{d}{dz}\left(\frac{B^2}{8\pi }\right)=\frac{d}{dz}\left(p_ep\right)=\left(\rho \rho _e\right)g.$$
$`\left(4.25\right)`$
With $`𝐠\widehat{𝐧}=g\mathrm{cos}\gamma \text{sgn}\left(d\gamma /dl\right)`$ and using Eq. (4.25) we find for the sum of the curvature and buoyancy forces which determines the force balance in normal direction (cf. Eq. 4.3)
$$0=\frac{d}{dz}\left(\frac{B^2}{8\pi }\right)\mathrm{cos}\gamma \frac{B^2}{4\pi }\frac{d}{dz}\mathrm{cos}\gamma =\frac{B}{4\pi }\frac{d}{dz}\left(B\mathrm{cos}\gamma \right)=\frac{B}{4\pi }\frac{dB_x}{dz}.$$
$`\left(4.26\right)`$
We see that the uniformity of the horizontal component of the magnetic field is sufficient for static equilibrium if the internal stratification along the flux tube is hydrostatic. We can use this property to reduce the calculation of the flux tube path, $`z=z\left(x\right)`$, to a quadrature. Since
$$8\pi \left(p_ep\right)=B^2=B_z^2+B_x^2$$
$`\left(4.27\right)`$
we can write
$$\frac{dz}{dx}=\frac{B_z}{B_x}=\left(\frac{8\pi \left(p_ep\right)}{B_x^2}1\right)^{1/2}.$$
$`\left(4.28\right)`$
Since $`B_x=`$ const. and the pressures are known functions of $`z`$, the function $`x\left(z\right)`$ can be determined by integration:
$$x\left(z\right)x\left(z_0\right)=_{z_0}^z\left(\frac{8\pi \left(p_ep\right)}{B_x^2}1\right)^{1/2}𝑑\stackrel{~}{z}$$
$`\left(4.29\right)`$
The resulting function may be inverted to yield the path, $`z=z\left(x\right)`$. For special cases, analytical solutions of Eq. (4.29) can be obtained (e.g. Parker, 1979a, Sec. 8.6).
4.2 Stationary equilibrium
As we have seen above, flux tubes in static equilibrium do not seem to be particularly relevant for the description of magnetic structures in the convection zone. On the other hand, the slow evolution of active regions after the vigorous dynamical phase of flux eruption indicates some kind of underlying equilibrium structure. Therefore it seems worthwhile to include the effect of large scale external velocity fields (convection, differential rotation) and to consider stationary equilibria of flux tubes. In practice, such an equilibrium can be determined either by direct (numerical) integration of a second order ordinary differential equation or by solving a first order equation and subsequent quadrature. Let us first consider the latter method (van Ballegooijen, 1982a).
We start from the plane-parallel geometry used in the preceding section and introduce a relative velocity, v, between flux tube and external medium. v is assumed to lie in the $`xz`$-plane. The unit vector $`\widehat{𝐤}`$ defined in Eqs. (3.28/29) in this case is given by
$$\widehat{𝐤}=\widehat{𝐧}\text{sgn}\left(𝐯\widehat{𝐧}\right).$$
$`\left(4.30\right)`$
We now determine
$$𝐯\widehat{𝐧}=\left(v_x\mathrm{sin}\gamma +v_z\mathrm{cos}\gamma \right)\text{sgn}\left(\frac{d\gamma }{dl}\right)v_{}\text{sgn}\left(\frac{d\gamma }{dl}\right)$$
and find
$$\text{sgn}\left(𝐯\widehat{𝐧}\right)=\text{sgn}\left(v_{}\right)\text{sgn}\left(\frac{d\gamma }{dl}\right).$$
$`\left(4.31\right)`$
Using Eqs. (3.26) and (3.28) force balance in the normal direction leads to
$$\left(\rho \rho _e\right)𝐠\widehat{𝐧}+\frac{B^2}{4\pi R}+\frac{C_D\rho _ev_{}^2}{\pi a}\widehat{𝐤}\widehat{𝐧}=\mathrm{\hspace{0.33em}0}.$$
$`\left(4.32\right)`$
In analogy to Eq. (4.26) we find using Eqs. (4.30/31)
$$\frac{B}{4\pi }\frac{dB_x}{dz}=\frac{C_D\rho _ev_{}^2}{\pi a}\text{sgn}\left(v_{}\right).$$
$`\left(4.33\right)`$
In principle, this equation can be used to determine the variation of the horizontal field component with height which then can be inserted in Eq. (4.29) to determine the path of the flux tube. However, since $`v_{}`$ depends on $`\gamma `$, the differential equation (4.33) in general is not easily solved, especially if v varies spatially.
As an example, let us consider the case of a purely horizontal velocity field $`𝐯=(v(z),\mathrm{\hspace{0.17em}0})`$ which may represent a depth-dependent differential rotation. We have $`v_{}=v\mathrm{sin}\gamma `$ and assume $`v0`$. Eq. (4.33) now reads
$$\frac{B}{4\pi }\frac{dB_x}{dz}=\frac{C_D\rho _ev^2}{\pi a}\mathrm{sin}^2\gamma .$$
$`\left(4.34\right)`$
With $`\pi a^2B=\mathrm{\Phi }_{mag}=`$ const. and $`B_x/B=\mathrm{cos}\gamma `$ we find
$$\frac{dB_x}{dz}=\mathrm{\hspace{0.17em}4}C_D\rho _ev^2\left(\frac{\pi }{\mathrm{\Phi }_{mag}B}\right)^{1/2}\left(1\frac{B_x^2}{B^2}\right).$$
$`\left(4.35\right)`$
In the special case of a constant magnetic field, $`B\left(z\right)=B_0`$, the variables can be separated and defining $`yB_x/B_0`$ we may write
$$\frac{dy}{1y^2}=\alpha \rho _ev^2dz$$
$`\left(4.36\right)`$
with
$$\alpha =\mathrm{\hspace{0.33em}4}C_D\left(\frac{\pi }{\mathrm{\Phi }_{mag}B_0^3}\right)^{1/2}.$$
$`\left(4.37\right)`$
As we have discussed in the preceding section the case of a constant magnetic field generally is of not much practical interest because it requires $`\rho =\rho _e`$ and, therefore, a very special internal temperature profile. However, if the scale height is much larger than the height range covered by the flux tube equilibrium path the variation of density is small and the assumption of a constant field may be tenable. Such a situation can be expected near the bottom of the solar convection zone. Integration of Eq. (4.36) yields
$$\frac{1}{2}\mathrm{ln}\left(\frac{1}{c_0}\frac{y1}{y+1}\right)=_0^z\alpha \rho _e\left(\stackrel{~}{z}\right)v^2\left(\stackrel{~}{z}\right)𝑑\stackrel{~}{z}f\left(z\right)$$
$`\left(4.38\right)`$
where $`c_0`$ is a constant of integration which is determined by a boundary condition at $`z=0`$. Obviously we must have $`y1`$ in Eq. (4.38). The case $`y=1`$ is a singular solution of Eq. (4.36) and represents a horizontal flux tube for which all forces vanish individually. Excluding this case we can determine $`y\left(z\right)`$ from Eq. (4.38)
$$y\left(z\right)=\frac{1+c_0e^{2f\left(z\right)}}{1c_0e^{2f\left(z\right)}}$$
$`\left(4.39\right)`$
and determine $`c_0`$ by specifying $`y\left(0\right)=y_0<1`$, i.e.
$$c_0=\frac{y_01}{y_0+1}.$$
$`\left(4.40\right)`$
For $`y_0[0,1)`$ the value of $`c_0`$ passes through the interval $`c_0[1,0)`$. In the case $`y_0=0`$ (vertical tube at $`z=0`$) we have $`c_0=1`$ and Eq. (4.39) gives
$$y\left(z\right)=\mathrm{tanh}\left[f\left(z\right)\right].$$
$`\left(4.41\right)`$
Since we have
$$z^{}\left(x\right)=\mathrm{tan}\gamma =\frac{\left(B^2B_x^2\right)^{1/2}}{B_x}=\frac{\left(1y^2\right)^{1/2}}{y}$$
$`\left(4.42\right)`$
integration yields
$$x\left(z\right)=_0^z\frac{y}{\left(1y^2\right)^{1/2}}𝑑\stackrel{~}{z}$$
$`\left(4.43\right)`$
where we have assumed $`x\left(0\right)=0`$. Inserting Eq. (4.39) into Eq. (4.43) we find
$$x\left(z\right)=_0^z\frac{1+c_0e^{2f\left(\stackrel{~}{z}\right)}}{2\sqrt{c_0}e^{f\left(\stackrel{~}{z}\right)}}𝑑\stackrel{~}{z}$$
$`\left(4.44\right)`$
and the special case $`c_0=1`$ gives
$$x\left(z\right)=_0^z\mathrm{sinh}\left[f\left(\stackrel{~}{z}\right)\right]d\stackrel{~}{z}.$$
$`\left(4.45\right)`$
We give two examples for which Eq. (4.44) can be directly integrated. First assume a velocity field with constant kinetic energy density, i.e. $`\left(\rho _ev^2\right)\left(z\right)=`$ const. Consequently, we have from Eq. (4.38) that $`f\left(z\right)=\alpha \rho _ev^2z\widehat{\alpha }z`$ and thus
$$x\left(z\right)=\frac{1}{2\widehat{\alpha }\sqrt{c_0}}\left(c_0e^{\widehat{\alpha }z}e^{\widehat{\alpha }z}+1c_0\right).$$
$`\left(4.46\right)`$
For an initially vertical tube, i.e. $`c_0=1`$, we have
$$x\left(z\right)=\widehat{\alpha }^1\left(1\mathrm{cosh}\left(\widehat{\alpha }z\right)\right)$$
$`\left(4.47\right)`$
and normalizing length by $`1/\widehat{\alpha }`$, viz. $`\widehat{z}\widehat{\alpha }z`$ and $`\widehat{x}\widehat{\alpha }x`$, we find
$$\widehat{x}=\mathrm{\hspace{0.33em}1}\mathrm{cosh}\widehat{z}.$$
$`\left(4.48\right)`$
This solution has earlier been given by Parker (1979d). Another directly solvable example is the case of kinetic energy density being proportional to $`z`$, viz.
$$\left(\rho _ev^2\right)\left(z\right)=\left(\rho _ev^2\right)\left(z_0\right)\left(\frac{z}{z_0}\right)$$
$`\left(4.49\right)`$
with some reference level $`z_0`$. Inserting Eq. (4.49) into Eq. (4.38) we have
$$f\left(z\right)=\frac{\alpha \left(\rho _ev^2\right)\left(z_0\right)}{2z_0}z^2\stackrel{~}{\alpha }z^2$$
$`\left(4.50\right)`$
and thus write Eq. (4.44) as
$$x\left(z\right)=\frac{1}{2\sqrt{\stackrel{~}{\alpha }c_0}}\left(_0^{\sqrt{\stackrel{~}{\alpha }}z}e^{w^2}𝑑w+c_0_0^{\sqrt{\stackrel{~}{\alpha }}z}e^{w^2}𝑑w\right).$$
$`\left(4.51\right)`$
The first integral essentially represents the error function $`𝚽\left(\sqrt{\stackrel{~}{\alpha }}z\right)`$ and the second is related to the Dawson integral $`𝚿`$ (cf. Abramowitz and Stegun, 1965)
$$𝚿\left(u\right)=e^{u^2}_0^ue^{w^2}𝑑w.$$
$`\left(4.52\right)`$
For $`c_0=1`$ we have
$$\widehat{x}=\frac{1}{2}\left(\frac{\sqrt{\pi }}{2}𝚽\left(\widehat{z}\right)e^{\widehat{z}^2}𝚿\left(\widehat{z}\right)\right)$$
$`\left(4.53\right)`$
with $`\widehat{z}=\sqrt{\stackrel{~}{\alpha }}z`$, $`\widehat{x}=\sqrt{\stackrel{~}{\alpha }}x`$. For both examples, Fig. 5 shows the resulting flux tube shape.
̧\vbox{ {\immediate} { {\immediate} \vboxto284.52756pt{ \hboxto398.33853pt{ \hfil} \vfil} }}
Fig. 5: Stationary equilibrium shape of flux tubes under the influence of a horizontal velocity field with kinetic energy density $`\rho v^2`$ constant with height $`z`$ (full line) and proportional to $`z`$ (dashed line). The direction of the flow is from left to right. Note that generally the length scales are not equal for both examples.
We see that the flux tube turns towards the flow because only in this way a balance of forces is possible. This is in contrast to a tree bent by a storm for which the differential tension force due to bending acts in the opposite direction of the normal vector. The path given by Eq. (4.53) and shown by the dashed line has a smaller curvature for small $`\widehat{z}`$ than that for constant kinetic energy density (Eq. 4.48, full line) but as $`\widehat{z}`$ increases it quickly bends over. In the case $`2z_0=1/\widehat{\alpha }`$ the length scales for both cases are equal and the curves can be directly compared. We may use these results to estimate the influence of velocity fields on slender flux tubes in the lower parts of the solar convection zone. Using the values
$`\rho _e=0.2`$ g$``$cm<sup>-3</sup>
$`B_0=10^4`$ Gauss (equipartition)
$`\mathrm{\Phi }_{mag}=10^{18}\mathrm{}10^{20}`$ mx ($`a=610^6\mathrm{}610^7`$ cm)
$`v_0=10^4`$ cm$``$s<sup>-1</sup>(convective flow, differential rotation)
$`z_0=610^9`$ cm (equal to the pressure scale height)
we find
$`\widehat{\alpha }^110^7\mathrm{}10^8`$ cm
$`\stackrel{~}{\alpha }^{1/2}310^8\mathrm{}10^9`$ cm.
We see in Fig. 5 that the flux tube paths become almost horizontal at typical heights $`\widehat{z}1\mathrm{}3`$ which refers to heights of the order of $`\widehat{\alpha }^1`$ and $`\stackrel{~}{\alpha }^{1/2}`$, respectively, for the two cases. This means that in equilibrium an initially vertical flux tube cannot intrude significantly into a layer of horizontal flow unless either the field strength is much larger than the equipartition value, the horizontal flow speed is much smaller than the typical convective velocities, or the radius of the flux tube is of the order of the scale height. The dominant rôle of external flow fields in the dynamics of thin flux tubes is an important effect which must be taken into account when discussing the properties of magnetic structures in stellar convection zones (see Ch. 6).
For other velocity fields $`𝐯(x,z)`$, Eq. (4.33) in most cases has to be solved by numerical forward integration in height starting from a suitable initial point. The appropriate value of $`v_{}`$ for each point is calculated using the earlier determined angle $`\gamma `$ and location $`x(z)`$. The value of $`B_x`$ which results from Eq. (4.33) can then be used to calculate the values of $`\gamma `$ and $`x(z)`$ at the next point using Eqs. (4.19) and (4.28/29). Both steps of this procedure which correspond to two integrations can be combined in the solution of a second order differential equation for the path $`z(x)`$ (cf. Schüssler, 1980a; Parker, 1982c; Anton, 1984). To this end we use the relation between radius of curvature and the derivatives, $`z^{}dz/dx`$, $`z^{\prime \prime }d^2z/dx^2`$, and write for the curvature force
$$\frac{B^2}{4\pi R}\widehat{𝐧}=\frac{B^2}{4\pi }\frac{z^{\prime \prime }}{\left(1+z^2\right)^{3/2}}\widehat{𝐧}.$$
$`\left(4.54\right)`$
Since we have $`\mathrm{cos}\gamma =\left(1+z^2\right)^{1/2}`$ the condition for stationary equilibrium, Eq. (4.32), can be written as
$$\left(\rho \rho _e\right)g+\frac{B^2}{4\pi }\frac{z^{\prime \prime }}{1+z^2}+\frac{C_D\rho _ev_{}^2}{\pi a}\left(1+z^2\right)^{1/2}\text{sgn}\left(v_{}\right)=\mathrm{\hspace{0.33em}0}$$
$`\left(4.55\right)`$
Sometimes it proves useful to rewrite the buoyancy term with aid of the relation
$$\frac{4\pi \left(\rho _e\rho \right)g}{B^2}=\frac{4\pi }{B^2}\frac{d}{dz}\left(\frac{B^2}{8\pi }\right)=\frac{1}{2}\frac{d}{dz}\mathrm{ln}\left(\frac{B^2}{4\pi }\right).$$
$`\left(4.56\right)`$
It depends on the properties of the particular problem which of the two possible ways to calculate the path $`z\left(x\right)`$, i.e. Eqs. (4.29/4.33) or Eq. (4.55), is more appropriate. Other examples for the calculation of stationary flux tube equilibria have been given by Parker (1979d; 1982c,d) for horizontal flows and for idealized cellular velocity fields, by van Ballegooijen (1982a) for a constant horizontal drift of flux tubes in the solar convection zone, and by Anton and Schüssler (unpublished) for giant convective cell patterns (see also Moreno-Insertis, 1984).
5. Stability of flux tubes
Static or stationary equilibrium configurations of flux tubes can only have a practical relevance if they are at least linearly stable, i.e. if the flux tube returns to its equilibrium position and shape after a small displacement. In Ch. 2 we have discussed mechanisms which lead to fragmentation of large magnetic structures and to the formation of flux tubes which are much smaller than the scale height in the deep parts of a convection zone. Such tubes are strongly influenced by motions in their environment like convection, differential rotation, and meridional circulation. Even if a flux tube is stable with respect to fragmentation and reaches a static or stationary equilibrium characterized by a balance of buoyancy, tension, drag and rotationally induced (Coriolis, centrifugal) forces, this equilibrium might be unstable due to
a) superadiabaticity of the environment,
b) loop formation with downflows from the crests to the troughs with concomitant perturbations of magnetic buoyancy (akin to the instability discussed by Parker, 1966),
c) gradients of the drag force exerted by external flows,
d) differential rotation.
These instabilities lead to motion and deformation of the flux tube as a whole and can be treated within the framework of the approximation of slender flux tubes. This allows a considerable simplification of the mathematics and simultaneously excludes the fragmentation instabilities (Rayleigh-Taylor, Kelvin-Helmholtz) discussed in Ch. 2.
5.1 Previous work
Vertical flux tubes in the (sub)photosphere of the Sun are liable to a convective instability (convective collapse) caused by the superadiabatic stratification of the surrounding fluid (Parker, 1978; Webb and Roberts, 1978; Spruit and Zweibel, 1979; Unno and Ando, 1979). It is believed that this process is responsible for the amplification of small-scale solar magnetic fields far beyond the equipartition field strength (for a more detailed discussion see Schüssler, 1990).
Apart from preliminary studies (e.g. Schüssler, 1980a) the first detailed stability analysis of flux tubes within a convection zone has been presented by Spruit and van Ballegooijen (1982) who analyzed the stability of horizontal, non-buoyant ($`\rho =\rho _e`$) flux tubes in cartesian geometry and of toroidal flux tubes in static equilibrium between buoyancy and tension force in spherical geometry. They found instabilities which represent a mixture between the convective and the (Parker type) kink <sup>*</sup><sup>*</sup>Not to be confused with the “classical” kink instability of a plasma pinch caused by the azimuthal magnetic field component: The kink instability discussed here is driven by magnetic buoyancy and superadiabaticity. instabilities (a and b above). The perturbations leading to kink instability must have a finite wavelength in the direction along the tube. It turns out that all flux tubes embedded in a superadiabatic environment are unstable and that kink instability occurs even for slightly subadiabatic stratification. Moreno-Insertis (1984, 1986) performed numerical simulations of the nonlinear evolution of kink-unstable horizontal flux tubes at the bottom of the solar convection zone. He found that while the upper part of the unstable loop rises towards the solar surface, the lower part sinks down and enters the subadiabatic region below the convection zone where it reaches a stable equilibrium. Moreno-Insertis (1984) also investigated the influence of the drag force exerted by a prescribed giant convective velocity cell on the evolution of the instability. Recently, Choudhuri (1989; see also Choudhuri and Gilman, 1987) has included rotation in a numerical study of kink-unstable flux tubes in the solar convection zone. He found in most cases that the Coriolis force dominates the dynamics and leads to a trajectory of the rising loop which is parallel to the axis of rotation. Consequently, the unstable loops break through the surface far away from the equatorial regions where solar activity predominantly is observed. A more radial eruption of unstable loops is only achieved for quite strong fields (about $`10^5`$ Gauss) for which buoyancy becomes the dominating force.
Van Ballegooijen (1983) continued the analytical stability study of toroidal flux tubes by including (differential) rotation of the external medium and a difference between the rotation rate of the gas within and outside of the flux tube. Such a difference (a longitudinal flow along the tube with respect to a coordinate system which corotates with the external gas) may arise due to conservation of angular momentum if the flux tube is carried by an equatorward meridional circulation in the lower part of the convection zone and thus increases its distance from the axis of rotation (cf. van Ballegooijen, 1982b). The Coriolis force caused by slower rotation of the internal gas helps to balance the component of the buoyancy force perpendicular to the axis of rotation. For flux tubes situated in the equatorial plane, rigid rotation and longitudinal flow have a stabilizing effect. However, in the parameter regime relevant for the lower convection zone of the Sun the stability properties are determined by differential rotation: The flux tube is stable if $`\mathrm{\Omega }_e/r>0`$, i.e. if the angular velocity of the external medium increases radially outward; it is unstable with respect to non-axisymmetric disturbances (growing waves along the tube) if $`\mathrm{\Omega }_e`$ decreases. Similar to the buoyancy-driven kink instability it is the downflow along the legs of a loop which triggers the instability due to differential rotation, in this case by introducing differential Coriolis forces in the radial direction.
The more general case of toroidal flux tubes outside the equatorial plane including meridional circulation has been treated by van Ballegooijen and Choudhuri (1988). In equilibrium, the component of the drag force perpendicular to the plane of the tube balances the corresponding component of the buoyancy force and thus removes the ‘poleward slip instability’ (Pneuman and Raadu, 1972; Spruit and van Ballegooijen, 1982). The authors found that an increase of the velocity of meridional circulation in radial direction has a stabilizing effect on the flux tube. However, their analysis is restricted to rigid rotation of the exterior and to axisymmetric perturbations such that the whole class of kink instabilities induced by buoyancy or differential rotation is neglected. The stability properties of toroidal flux tubes under these conditions remain to be investigated.
For a totally different application, the stability and interaction of jets from active galactic nuclei, Achterberg (1982, 1988) has derived a formalism which is similar to these approaches and also to the formalism developed in this work.
In this chapter we present a formalism which can be used to analyze the stability of general static and stationary flux tube equilibria in a plane with constant direction of gravity. It can be applied to any given plane equilibrium path $`𝐫(l)`$ along which all quantities may vary. Such a formalism is needed in order to determine the stability properties of a number of non-trivial flux tube equilibria like the examples given in Sec. 4.2 or the loop structures calculated by van Ballegooijen (1982a) and Anton (1984). We include a velocity field of arbitrary structure but, in order to limit the complication of the already somewhat involved formalism, we have refrained from treating the spherical case and also ignored rotation. However, this restriction is not fundamental and will be dropped in future work.
5.2 Equilibrium
We assume a plane flux tube in static or stationary equilibrium described by the time-independent form of Eqs. (3.25) and (3.26). We take $`𝐠\widehat{𝐛}=0`$ and assume that the external velocity is restricted to the plane of the equilibrium tube such that both terms on the r.h.s. of Eq. (3.27) (binormal direction) vanish. We continue to use the notation introduced in Chs. 3 and 4. Denoting all equilibrium quantities by a suffix ‘0’, the equilibrium path of the flux tube is given by $`𝐫_0\left(l_0\right)`$ and the normal, $`\widehat{𝐧}_0\left(l_0\right)`$, and tangential, $`\widehat{𝐥}_0\left(l_0\right)`$, unit vectors as well as the radius of curvature, $`R_0`$, are defined in the usual way as functions of the equilibrium path length, $`l_0`$:
$$\widehat{𝐥}_0=\frac{𝐫_0}{l_0},\widehat{𝐧}_0=R_0\frac{\widehat{𝐥}_0}{l_0},R_0=\left|\frac{\widehat{𝐥}_0}{l_0}\right|^1.$$
$`\left(5.1\right)`$
The equilibrium state is characterized by hydrostatic equilibrium in the longitudinal direction (cf. Eq. 3.25)
$$0=\frac{p_0}{l_0}+\rho _0𝐠_0\widehat{𝐥}_0$$
$`\left(5.2\right)`$
and by a balance of buoyancy, curvature and drag force in the normal direction (cf. Eq. 3.26)
$$0=\left(\rho _0\rho _{e0}\right)g_0\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)+\frac{B_0^2}{4\pi R_0}+\frac{C_D\rho _{e0}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}{\pi a_0}$$
$`\left(5.3\right)`$
$`\left(𝐠_0=g_0\widehat{𝐠}_0\right)`$. Eqs. (5.2/3) are complemented by the conditions of pressure balance, Eq. (3.23), and flux conservation, Eq. (3.31), and by hydrostatic equilibrium of the external medium, Eq. (3.24). Note that all quantities may depend on $`l_0`$. Since the external velocity, $`𝐯_e`$, has no component perpendicular to the plane of the tube the drag force given by Eqs. (3.28/29) can be written in the simpler form shown in Eq. (5.3). We may express the equilibrium condition in terms of the relative density contrast between exterior and interior of the flux tube by rewriting Eq. (5.3) in the form
$$\beta \left(\frac{\rho _{e0}}{\rho _0}1\right)\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)=\frac{2H_{p0}}{R_0}+\beta r\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)$$
$`\left(5.4\right)`$
where $`\beta =8\pi p_0/B_0^{\mathrm{\hspace{0.17em}2}}`$, $`H_{p0}`$ is the internal pressure scale height defined by $`p_0/H_{p0}=\rho _0g_0`$, and $`r`$ is given by
$$r=\frac{\rho _{e0}}{\rho _0}\frac{C_D\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2}{\pi a_0g_0}.$$
$`\left(5.5\right)`$
The product $`\beta r`$ can be written in the form
$$\beta r=\left(\frac{2C_D}{\pi }\right)\left(\frac{\rho _{e0}}{\rho _0}\right)\left(\frac{\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2}{v_{A0}^2}\right)\left(\frac{H_{p0}}{a_0}\right)$$
$`\left(5.6\right)`$
where $`v_{A0}=B_0/\sqrt{4\pi \rho _0}`$ is the Alfvén velocity. For equipartition fields ($`B_0=10^4\mathrm{}10^5`$ Gauss) near the bottom of the solar convection zone ($`p_0=610^{13}\text{dyn}\text{cm}^2`$) we have very large values of $`\beta 1.510^5\mathrm{}1.510^7`$. We shall use this property in Sec. 5.3 to derive a simplified version of the perturbation equations in the limit $`\beta 1`$. In most cases relevant for the deep convection zone, the density contrast in Eq. (5.4) is very small such that $`\beta \left(\rho _{e0}/\rho _01\right)`$, $`H_{p0}/R_0`$, and $`\beta r`$ are all of order unity or at least have moderate values.
5.3 Perturbation equations
In order to determine the linear stability of the equilibrium we introduce Lagrangean displacements in the normal and tangential directions described by the functions $`\epsilon \left(l_0\right)`$ and $`\eta \left(l_0\right)`$, respectively, and write for the perturbed path <sup>*</sup><sup>*</sup>The perturbations of all quantities which have a non-vanishing equilibrium value are indicated by a suffix ‘1’. Perturbed quantities (equilibrium value $`+`$ perturbation) are written without suffix.
$$𝐫=𝐫_0+𝐫_1𝐫_0\left(l_0\right)+\epsilon \left(l_0\right)\widehat{𝐧}_0+\eta \left(l_0\right)\widehat{𝐥}_0.$$
$`\left(5.7\right)`$
Similar to the analyses of Spruit and van Ballegooijen (1982) and of van Ballegooijen (1983), a displacement in the binormal direction decouples from the rest of the equations and gives rise to (stable) transversal flux tube waves which are of no further interest here. We therefore consider only displacements within the plane of the equilibrium flux tube path.
First we determine the perturbed geometry of the flux tube. The relation between the arc length of the perturbed tube, $`l`$, and the arc length in equilibrium, $`l_0`$, to first order in the perturbations is given by
$$\frac{l}{l_0}=\mathrm{\hspace{0.33em}1}\frac{\epsilon }{R_0}+\frac{\eta }{l_0}.$$
$`\left(5.8\right)`$
The perturbed tangent vector, $`\widehat{𝐥}`$, is
$$\widehat{𝐥}=\frac{𝐫}{l}=\frac{𝐫}{l_0}\frac{l_0}{l}=\widehat{𝐥}_0+\widehat{𝐧}_0\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right).$$
$`\left(5.9\right)`$
Now we can calculate the curvature vector, $`𝐜\widehat{𝐧}/R`$ :
$$𝐜=\frac{\widehat{𝐥}}{l}=\frac{\widehat{𝐥}}{l_0}\frac{ł_0}{l}=\frac{\widehat{𝐧}_0}{R_0}\left(1+\frac{\epsilon }{R_0}+R_0\frac{d^2\epsilon }{dl_0^{\mathrm{\hspace{0.17em}2}}}\frac{\eta }{R_0}\frac{R_0}{l_0}\right)\frac{\widehat{𝐥}_0}{R_0}\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right).$$
$`\left(5.10\right)`$
Taking the absolute value of Eq. (5.10) gives the perturbed radius of curvature, $`R`$, and the perturbed normal vector, viz.
$$R=R_0\left(1+\frac{\epsilon }{R_0}+R_0\frac{d^2\epsilon }{dl_0^{\mathrm{\hspace{0.17em}2}}}\frac{\eta }{R_0}\frac{R_0}{l_0}\right)^1$$
$`\left(5.11\right)`$
$$\widehat{𝐧}=\widehat{𝐧}_0\widehat{𝐥}_0\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right).$$
$`\left(5.12\right)`$
We now proceed by determining the gravitational acceleration, g, in the perturbed state. We assume a power law for the dependence of $`g_0`$ on height, $`z`$, and write
$$𝐠_0=g_m\left(\frac{z_0}{z_m}\right)^s\widehat{𝐳}g_0\widehat{𝐳}$$
$`\left(5.13\right)`$
($`\widehat{𝐳}`$: unit vector in direction of height; $`g_m`$: gravitational acceleration at some reference height, $`z_m`$; $`z_0\left(l_0\right)`$: height of equilibrium flux tube). With the height displacement, $`z_1`$, given by
$$z_1𝐫_1\widehat{𝐳}=𝐫_1\widehat{𝐠}_0=\epsilon \left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\eta \left(\widehat{𝐠}_0\widehat{𝐥}_0\right)$$
$`\left(5.14\right)`$
where $`\widehat{𝐠}_0=\widehat{𝐳}`$ denotes the unit vector in the direction of gravity we find, to first order in $`z_1`$,
$$\begin{array}{ccc}\hfill 𝐠=𝐠_0+𝐠_1& =g_m\left(\frac{z_0+z_1}{z_m}\right)^s\widehat{𝐳}=g_m\left(\frac{z_0}{z_m}\right)^s\left(1+\frac{z_1}{z_0}\right)^s\widehat{𝐳}\hfill & \\ & & \\ & =𝐠_0\left(1s\frac{\epsilon \left(\widehat{𝐠}_0\widehat{𝐧}_0\right)+\eta \left(\widehat{𝐠}_0\widehat{𝐥}_0\right)}{z_0}\right).\hfill & \left(5.15\right)\hfill \end{array}$$
Defining
$$\mathrm{\Delta }\frac{z_1}{z_0}=\frac{\epsilon \left(\widehat{𝐠}_0\widehat{𝐧}_0\right)+\eta \left(\widehat{𝐠}_0\widehat{𝐥}_0\right)}{z_0}$$
$`\left(5.16\right)`$
we write for the component of gravity along $`\widehat{𝐥}`$ with aid of Eqs. (5.9) and (5.15)
$$𝐠\widehat{𝐥}=\left(1s\mathrm{\Delta }\right)𝐠_0\left[\widehat{𝐥}_0+\widehat{𝐧}_0\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right)\right]$$
$`\left(5.17\right)`$
which gives for the perturbation of longitudinal component of g
$$\left(𝐠\widehat{𝐥}\right)_1=g_0\left[\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right)\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)s\mathrm{\Delta }\right].$$
$`\left(5.18\right)`$
For the perturbation of the normal component of g we find using Eqs. (5.12) and (5.15)
$$\left(𝐠\widehat{𝐧}\right)_1=g_0\left[\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)s\mathrm{\Delta }+\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right)\right].$$
$`\left(5.19\right)`$
We now consider the longitudinal component of the equation of motion, Eq. (3.25). To first order in the displacements, the inertial term on its l.h.s. can be written using $`𝐮=d𝐫_1/dt(\dot{\eta },\dot{\epsilon })`$ in the form
$$\rho \frac{d𝐮}{dt}\widehat{𝐥}=\rho \frac{d𝐮\widehat{𝐥}}{dt}\rho 𝐮\frac{d\widehat{𝐥}}{dt}=\rho \left(\frac{d}{dt}\left(\dot{\eta }\widehat{𝐥}_0+\dot{\epsilon }\widehat{𝐧}_0\right)\right)\widehat{𝐥}\rho 𝐮\left(\frac{d\widehat{𝐥}_0}{dt}+\frac{d\widehat{𝐥}_1}{dt}\right)=\rho _0\ddot{\eta }.$$
$`\left(5.20\right)`$
The perturbation of the longitudinal pressure gradient is given by
$$\left(\frac{p}{l}\right)_1=\frac{p_1}{l_0}+\frac{p_0}{l_0}\left(\frac{\epsilon }{R_0}\frac{\eta }{l_0}\right)$$
$`\left(5.21\right)`$
and the perturbation of the gravity force reads
$$\begin{array}{ccc}\hfill \left(\rho 𝐠\widehat{𝐥}\right)_1& =\rho _0\left(𝐠\widehat{𝐥}\right)_1+\rho _1g_0\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\hfill & \\ & & \\ & =\rho _0g_0\left[\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right)\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)s\mathrm{\Delta }\right]+\frac{\rho _1}{\rho _0}\frac{p_0}{l_0}.\hfill & \left(5.22\right)\hfill \end{array}$$
where we have used Eqs. (5.2) and (5.18). Inserting Eqs. (5.20-5.22) into Eq. (3.25) we find for the longitudinal equation of motion, to first order in the perturbations:
$$\begin{array}{ccc}\hfill \ddot{\eta }=& \frac{1}{\rho _0}\frac{p_1}{l_0}\frac{1}{\rho _0}\frac{p_0}{l_0}\left(\frac{\epsilon }{R_0}\frac{\eta }{l_0}\frac{\rho _1}{\rho _0}\right)\hfill & \\ & & \\ & +g_0\left[\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right)\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)s\mathrm{\Delta }\right].\hfill & \left(5.23\right)\hfill \end{array}$$
Similar to Eq. (5.20) the inertial term in Eq. (3.26), the normal component of the equation of motion, reduces to $`\rho _0\ddot{\epsilon }`$. <sup>*</sup><sup>*</sup>An enhanced inertia of the flux tube with respect to transversal motions due to the acceleration of material in the exterior is not considered here. It would only affect growth rates (or oscillation frequencies) but does not change the stability criteria. In the case $`\beta =8\pi p_0/B_0^21`$ which we consider later the enhanced inertia can be crudely taken account of by multiplying the inertial term of the normal component of the equation of motion by a factor 2 (cf. Spruit and van Ballegooijen, 1982). Using Eq. (5.19) the perturbation of the buoyancy force on the r.h.s. of Eq. (3.26) is given by
$$\begin{array}{ccc}\hfill \left[\left(\rho \rho _e\right)𝐠\widehat{𝐧}\right]_1=& \left(\rho _1\rho _{e1}\right)g_0\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\hfill & \\ & & \\ & \left(\rho _0\rho _{e0}\right)g_0\left[\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)s\mathrm{\Delta }+\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right)\right].\hfill & \left(5.24\right)\hfill \end{array}$$
The perturbation of the external density, $`\rho _{e1}`$, in our Lagrangian approach is
$$\rho _{e1}=\rho _{e0}z_1H_{\rho e}^1=\rho _{e0}z_1\left(1\right)H_{pe}^1$$
$`\left(5.25\right)`$
($`H_{\rho e}`$: external density scale height; $`H_{pe}`$: external pressure scale height; $``$: logarithmic temperature gradient in the external medium; all these quantities are taken at $`z=z_0`$). The relation $`H_{pe}=\left(1\right)H_{\rho e}`$ is valid if the molecular weight is constant, an assumption which is well justified in the lower parts of the solar convection zone. With aid of Eq. (5.11) the perturbation of the curvature force can be written
$$\left(\frac{B^2}{4\pi R}\right)_1=\frac{B_0B_1}{2\pi R_0}+\frac{B_0^2}{4\pi R_0}\left(\frac{\epsilon }{R_0}+R_0\frac{d^2\epsilon }{dl_0^{\mathrm{\hspace{0.17em}2}}}\frac{\eta }{R_0}\frac{R_0}{l_0}\right).$$
$`\left(5.26\right)`$
The perturbation of the drag force is slightly more complicated to determine. We start by noting that for linear analysis we may take $`\text{sgn}\left(𝐯_e\widehat{𝐧}\right)=\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)`$ since the displacement can always be made sufficiently small. In the case $`𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0=0`$ the perturbation of the drag term vanishes identically as one can see in Eqs. (5.31/33) below. The perturbed relative velocity between flux tube and environment is given by
$$𝐯_e=𝐯_{e}^{}{}_{0}{}^{}+\left(𝐫_1\right)𝐯_e|_{𝐫_0}𝐮.$$
$`\left(5.27\right)`$
The third term on the r.h.s. represents the motion of the tube due to the displacement while the second term describes the spatial change of $`𝐯_e`$ which is written in cartesian coordinates $`(x,z)`$:
$$𝐰_{e1}\left(𝐫_1\right)𝐯_e|_{𝐫_0}=x_1\frac{𝐯_e}{x}|_{𝐫_0}+z_1\frac{𝐯_e}{z}|_{𝐫_0}$$
$`\left(5.28\right)`$
with
$$\begin{array}{ccc}\hfill x_1& =𝐫_1\widehat{𝐱}=\epsilon \left(\widehat{𝐱}\widehat{𝐧}_0\right)+\eta \left(\widehat{𝐱}\widehat{𝐥}_0\right)\hfill & \\ \hfill z_1& =𝐫_1\widehat{𝐳}=\epsilon \left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\eta \left(\widehat{𝐠}_0\widehat{𝐥}_0\right).\hfill & \left(5.29\right)\hfill \end{array}$$
$`\widehat{𝐱}`$ is the unit vector in direction of the $`x`$ coordinate, i.e. the horizontal direction. Next we determine the perturbation of $`\left(𝐯_e\widehat{𝐧}\right)^2`$. Using Eqs. (5.12) and (5.27) we write
$$\left(𝐯_e\widehat{𝐧}\right)^2=\left[\left(𝐯_{e}^{}{}_{0}{}^{}+𝐰_{e1}𝐮\right)\left(\widehat{𝐧}_0\widehat{𝐥}_0\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right)\right)\right]^2$$
$`\left(5.30\right)`$
which gives
$$v_1\left(𝐯_e\widehat{𝐧}\right)_1^2=\mathrm{\hspace{0.33em}2}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)\left[\left(𝐰_{e1}\widehat{𝐧}_0\right)\dot{\epsilon }\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐥}_0\right)\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right)\right].$$
$`\left(5.31\right)`$
Note that $`v_1`$ may have positive or negative sign. Using Eq. (5.28/29) we obtain
$$\begin{array}{ccc}\hfill \left(𝐰_{e1}\widehat{𝐧}_0\right)& =\left[\epsilon \left(\widehat{𝐱}\widehat{𝐧}_0\right)+\eta \left(\widehat{𝐱}\widehat{𝐥}_0\right)\right]\left(\frac{𝐯_e}{x}|_{𝐫_0}\widehat{𝐧}_0\right)\left[\epsilon \left(\widehat{𝐠}_0\widehat{𝐧}_0\right)+\eta \left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\right]\left(\frac{𝐯_e}{z}|_{𝐫_0}\widehat{𝐧}_0\right)\hfill & \\ & & \\ & \left[\epsilon \left(\widehat{𝐱}\widehat{𝐧}_0\right)+\eta \left(\widehat{𝐱}\widehat{𝐥}_0\right)\right]w_x\left[\epsilon \left(\widehat{𝐠}_0\widehat{𝐧}_0\right)+\eta \left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\right]w_z.\hfill & \left(5.32\right)\hfill \end{array}$$
Now we can write down the perturbation of the drag term:
$$\left(\frac{C_D\rho _e\left(𝐯_e\widehat{𝐧}\right)^2\text{sgn}\left(𝐯_e\widehat{𝐧}\right)}{\pi a}\right)_1=\frac{C_D\rho _{e0}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}{\pi a_0}\left(\frac{v_1}{\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2}+\frac{\rho _{e1}}{\rho _{e0}}\frac{a_1}{a_0}\right).$$
$`\left(5.33\right)`$
The last term on the r.h.s. can be rewritten using the conservation of magnetic flux, $`Ba^2=`$ const., which yields
$$\frac{a_1}{a_0}=\frac{B_1}{2B_0}.$$
$`\left(5.34\right)`$
We combine Eqs. (5.24), (5.26) and (5.33/34) and obtain for the normal component of the perturbed equation of motion by inserting into Eq. (3.26):
$$\begin{array}{ccc}\hfill \ddot{\epsilon }=& \frac{\rho _1\rho _{e1}}{\rho _0}g_0\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\left(1\frac{\rho _{e0}}{\rho _0}\right)g_0\left[\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)s\mathrm{\Delta }+\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\left(\frac{\eta }{R_0}+\frac{\epsilon }{l_0}\right)\right]\hfill & \\ & & \\ & +\frac{B_0^2}{4\pi R_0\rho _0}\left(\frac{2B_1}{B_0}+\frac{\epsilon }{R_0}+R_0\frac{d^2\epsilon }{dl_0^{\mathrm{\hspace{0.17em}2}}}\frac{\eta }{R_0}\frac{R_0}{l_0}\right)\hfill & \\ & & \\ & +\frac{C_D\rho _{e0}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}{\pi a_0\rho _0}\left(\frac{v_1}{\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2}+\frac{\rho _{e1}}{\rho _{e0}}+\frac{B_1}{2B_0}\right).\hfill & \left(5.35\right)\hfill \end{array}$$
Since it is our aim to obtain a set of two coupled equations for the displacements $`\eta `$ and $`\epsilon `$ alone, we have to eliminate the other perturbations ($`B_1,\rho _1,\rho _{e1},p_1,\mathrm{}`$) in favor of $`\eta `$, $`\epsilon `$, and equilibrium quantities. This has already been achieved for $`\rho _{e1}`$ with Eqs. (5.25) and (5.14) and for $`v_1`$ in Eqs. (5.31/32). For the remaining quantities we use Eq. (3.13), the equation of continuity, which can be written to first order in the perturbations (after a time integration) in the form
$$\frac{\rho _1}{\rho _0}\frac{B_1}{B_0}+\frac{\eta }{l_0}\frac{\epsilon }{R_0}=\mathrm{\hspace{0.33em}0},$$
$`\left(5.36\right)`$
the equation of state for adiabatic perturbations,
$$\frac{dp}{dt}=\frac{\gamma p}{\rho }\frac{d\rho }{dt},$$
$`\left(5.37\right)`$
which yields (to first order, after time integration)
$$\frac{p_1}{p_0}=\gamma \frac{\rho _1}{\rho _0},$$
$`\left(5.38\right)`$
and the condition of instantaneous pressure equilibrium, Eq. (3.23), which gives
$$\frac{p_1}{p_0}+\frac{2}{\beta }\frac{B_1}{B_0}=\frac{p_{e1}}{p_0}$$
$`\left(5.39\right)`$
with $`\beta =8\pi p_0/B_0^2`$. Assuming hydrostatic equilibrium in the environment and using the equilibrium pressure balance condition, $`p_0+B_0^2/8\pi =p_{e0}`$, we rewrite Eq. (5.39) in the form
$$\frac{p_1}{p_0}+\frac{2}{\beta }\frac{B_1}{B_0}=\left(1+\frac{1}{\beta }\right)\frac{z_1}{H_{pe}}.$$
$`\left(5.40\right)`$
Eqs. (5.36) and (5.38/39) are used to obtain
$$\begin{array}{ccc}\hfill \frac{B_1}{B_0}& =\left(\frac{\beta \gamma }{\beta \gamma +2}\right)\left(\frac{\eta }{l_0}\frac{\epsilon }{R_0}\right)\left(\frac{\beta +1}{\beta \gamma +2}\right)\frac{z_1}{H_{pe}}\hfill & \left(5.41\right)\hfill \\ & & \\ \hfill \frac{\rho _1}{\rho _0}& =\left(\frac{\beta \gamma }{\beta \gamma +2}1\right)\left(\frac{\eta }{l_0}\frac{\epsilon }{R_0}\right)\left(\frac{\beta +1}{\beta \gamma +2}\right)\frac{z_1}{H_{pe}}\hfill & \left(5.42\right)\hfill \\ & & \\ \hfill \frac{p_1}{p_0}& =\gamma \left(\frac{\beta \gamma }{\beta \gamma +2}1\right)\left(\frac{\eta }{l_0}\frac{\epsilon }{R_0}\right)\gamma \left(\frac{\beta +1}{\beta \gamma +2}\right)\frac{z_1}{H_{pe}}.\hfill & \left(5.43\right)\hfill \end{array}$$
Using the abbreviations
$$\frac{\beta \gamma }{\beta \gamma +2}\alpha _1,\frac{\beta +1}{\beta \gamma +2}\alpha _2$$
$`\left(5.44\right)`$
we insert Eqs. (5.42/43) into Eq. (5.23) and obtain for the equation which determines the time evolution of the longitudinal displacement $`\eta `$ (primes denote derivatives with respect to $`l_0`$):
$$\begin{array}{ccc}\hfill \ddot{\eta }=& \epsilon \{(\alpha _11)\frac{\gamma p_0}{\rho _0}\frac{R_0^{}}{R_0^2}+\frac{1}{\rho _0R_0}\left[\gamma p_0(\alpha _11)\right]^{}+\frac{1}{\rho _0}[\frac{\alpha _2p_0^{}}{H_{pe}}\left(\frac{\gamma p_0\alpha _2}{H_{pe}}\right)^{}](\widehat{𝐠}_0\widehat{𝐧}_0)\hfill & \\ & & \\ & \frac{\alpha _1p_0^{}}{\rho _0R_0}+\frac{\gamma p_0\alpha _2}{\rho _0H_{pe}R_0}(\widehat{𝐠}_0\widehat{𝐥}_0)\frac{g_0s}{z_0}(\widehat{𝐠}_0\widehat{𝐥}_0)(\widehat{𝐠}_0\widehat{𝐧}_0)\}\hfill & \\ & & \\ & +\epsilon ^{}\left\{\left(\alpha _11\right)\frac{\gamma p_0}{\rho _0R_0}+\left(g_0\frac{\gamma p_0\alpha _2}{\rho _0H_{pe}}\right)\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\right\}\hfill & \\ & & \\ & +\eta \{[\frac{\alpha _2p_0^{}}{\rho _0H_{pe}}\frac{1}{\rho _0}\left(\frac{\gamma p_0\alpha _2}{H_{pe}}\right)^{}](\widehat{𝐠}_0\widehat{𝐥}_0)\hfill & \\ & & \\ & +(\frac{g_0}{R_0}\frac{\gamma p_0\alpha _2}{\rho _0H_{pe}R_0})(\widehat{𝐠}_0\widehat{𝐧}_0)\frac{g_0s}{z_0}(\widehat{𝐠}_0\widehat{𝐥}_0)^2\}\hfill & \\ & & \\ & +\eta ^{}\left\{\frac{1}{\rho _0}\left[\gamma p_0\left(\alpha _11\right)\right]^{}+\frac{p_0^{}\alpha _1}{\rho _0}\frac{\gamma p_0\alpha _2}{\rho _0H_{pe}}\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\right\}\hfill & \\ & & \\ & +\eta ^{\prime \prime }\left\{\left(\alpha _11\right)\frac{\gamma p_0}{\rho _0}\right\}.\hfill & \left(5.45\right)\hfill \end{array}$$
In order to derive Eq. (5.45) we have used Eqs. (5.14), (5.16) and the relation
$$\left[\epsilon \left(\widehat{𝐠}_0\widehat{𝐧}_0\right)+\eta \left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\right]^{}=\left(\eta ^{}\frac{\epsilon }{R_0}\right)\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)+\left(\epsilon ^{}+\frac{\eta }{R_0}\right)\left(\widehat{𝐠}_0\widehat{𝐧}_0\right).$$
$`\left(5.46\right)`$
In a similar way we obtain the equation which determines the time evolution of $`\epsilon `$, the displacement in normal direction, by inserting into Eq. (5.35):
$$\begin{array}{ccc}\hfill \ddot{\epsilon }=& \eta \{\frac{2p_0R_0^{}}{\rho _0\beta R_0^2}+\frac{4p_0\alpha _2}{\rho _0\beta R_0H_{pe}}(\widehat{𝐠}_0\widehat{𝐥}_0)+\frac{g_0}{H_{pe}}[\alpha _2+\frac{\rho _{e0}}{\rho _0}(1)](\widehat{𝐠}_0\widehat{𝐥}_0)(\widehat{𝐠}_0\widehat{𝐧}_0)\hfill & \\ & & \\ & \left(1\frac{\rho _{e0}}{\rho _0}\right)\left[\frac{g_0s}{z_0}\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)+\frac{g_0}{R_0}\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\right]\hfill & \\ & & \\ & +\frac{C_D\rho _{e0}\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2}{\pi a_0\rho _0}[\frac{2}{\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}((\widehat{𝐱}\widehat{𝐥}_0)w_x(\widehat{𝐠}_0\widehat{𝐥}_0)w_z)\hfill & \\ & & \\ & \frac{2\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐥}_0\right)}{R_0\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}\frac{1}{H_{pe}}(\mathrm{\hspace{0.17em}1}\frac{\alpha _2}{2})(\widehat{𝐠}_0\widehat{𝐥}_0)]\}\hfill & \\ & & \\ & +\eta ^{}\left\{\frac{4p_0\alpha _1}{\rho _0\beta R_0}+g_0\left(\alpha _11\right)\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)+\frac{C_D\rho _{e0}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)\alpha _1}{2\pi a_0\rho _0}\right\}\hfill & \\ & & \\ & +\epsilon \{\frac{4p_0\alpha _1}{\rho _0\beta R_0^2}+\frac{4p_0\alpha _2}{\rho _0\beta R_0H_{pe}}(\widehat{𝐠}_0\widehat{𝐧}_0)+\frac{2p_0}{\rho _0R_0^2\beta }\frac{g_0s}{z_0}(1\frac{\rho _{e0}}{\rho _0})(\widehat{𝐠}_0\widehat{𝐧}_0)^2\hfill & \\ & & \\ & \left[\frac{g_0}{R_0}\left(\alpha _11\right)\frac{g_0}{H_{pe}}\left(\alpha _2+\frac{\rho _{e0}}{\rho _0}\left(1\right)\right)\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\right]\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\hfill & \\ & & \\ & +\frac{C_D\rho _{e0}\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)^2}{\pi a_0\rho _0}[\frac{2}{\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}((\widehat{𝐱}\widehat{𝐧}_0)w_x(\widehat{𝐠}_0\widehat{𝐧}_0)w_z)\hfill & \\ & & \\ & \frac{1}{H_{pe}}(\mathrm{\hspace{0.17em}1}\frac{\alpha _2}{2})(\widehat{𝐠}_0\widehat{𝐧}_0)\frac{\alpha _1}{2R_0}]\}\hfill & \\ & & \\ & +\epsilon ^{}\left\{g_0\left(1\frac{\rho _{e0}}{\rho _0}\right)\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\frac{2C_D\rho _{e0}\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}{\pi a_0\rho _0}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐥}_0\right)\right\}\hfill & \\ & & \\ & +\epsilon ^{\prime \prime }\left\{\frac{2p_0}{\rho _0\beta }\right\}\hfill & \\ & & \\ & \dot{\epsilon }\left\{\frac{2C_D\rho _{e0}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}{\pi a_0\rho _0}\right\}.\hfill & \left(5.47\right)\hfill \end{array}$$
In the case of a non-constant molecular weight, the term $`\left(1\right)/H_{pe}`$ has to be replaced by $`1/H_{\rho e}`$ (cf. Eq. 5.25). A form of Eqs. (5.45/5.47) which is more convenient for our purposes is obtained below by considering the limit $`\beta 1`$ and a suitable non-dimensionalization.
5.4 Non-dimensionalization and the case $`\beta 1`$
As we have discussed in Sec. 5.2, the value of $`\beta `$ for equipartition flux tubes in the deep layers of the solar convection zone is very large. Having this application in mind, it is convenient to rewrite and simplify Eqs. (5.45/47) in the limit $`\beta 1`$ and to introduce non-dimensional quantities. As units for all quantities we use their values in the interior of the flux tube at the reference level, $`z=z_m`$, which has already been defined in Eq. (5.13). As length scale we choose the internal scale height, $`H_{pm}H_{p0}\left(z_m\right)`$, and the units of magnetic field, pressure, density etc. are defined by $`B_mB_0\left(z_m\right)`$, $`p_mp_0\left(z_m\right)`$, $`\rho _m\rho _0\left(z_m\right)`$ etc., respectively. As time unit, $`\tau `$, we take, similar to Spruit and van Ballegooijen (1982)
$$\tau =\left(\frac{\beta _mH_{pm}}{g_m}\right)^{1/2}$$
$`\left(5.48\right)`$
with $`\beta _m=8\pi p_m/B_m^{\mathrm{\hspace{0.17em}2}}`$. It is easy to see that $`\tau `$ is $`\sqrt{2}`$ times the Alfvén travel time across one scale height. For values of $`\beta _m`$ of the order $`10^5\mathrm{}10^7`$, $`H_{pm}=10^9`$ cm, $`g_m=`$$`610^4`$ cm$`{}_{}{}^{2}`$s<sup>-1</sup> we find $`\tau =10^5\mathrm{}10^6`$s $`1\mathrm{}2`$ days. The velocity unit is defined as $`V=H_{pm}/\tau =v_A/\sqrt{2}`$ where $`v_A=B_m/\sqrt{4\pi \rho _m}`$ is the Alfvén speed at the reference height. Since Eqs. (5.45/47) are linear and homogeneous, we can separate the time dependence by writing $`f\mathrm{exp}\left(i\omega t\right)`$ for all quantities where $`\omega `$ is a (generally complex) frequency. This leads to
$$(\ddot{\eta },\ddot{\epsilon })=\omega ^2(\eta ,\epsilon ).$$
$`\left(5.49\right)`$
We insert Eq. (5.49) into Eqs. (5.45/47) and take the limit $`\beta 1`$ such that only terms of order unity and order $`\beta ^1`$ are retained. It turns out that the terms of order unity cancel. Consequently, non-dimensionalization by multiplication of the resulting equations with $`\tau ^2=\beta _mH_{pm}/g_m`$ leads to terms of order unity and to terms of the form $`\beta \delta `$, $`\beta r`$, and $`\beta \left(1\rho _{e0}/\rho _0\right)`$ which are also of order unity ($`\delta `$ is the superadiabaticity of the external gas and $`r`$ is defined in Eq. 5.5). For example, a flux tube in temperature equilibrium with its surroundings ($`T_0=T_{e0}`$) has $`\beta \left(1\rho _{e0}/\rho _0\right)=1`$ and, using a standard mixing-length model, Spruit and van Ballegooijen (1982) found a value of $`\beta \delta =3.6`$ for an equipartition flux tube in the deep parts of the solar convection zone. Thus for $`\beta \delta 1`$ we have fields strong compared to the equipartition value, for $`\beta \delta 1`$ we have weak fields.
We give some examples of how the transformation of the coefficients in Eqs. (5.45/47) is carried out but avoid a presentation of the whole calculation. The algebra is straightforward although lengthy and tedious, filling dozens of pages. Let us first write down the quantities $`\alpha _1`$ and $`\alpha _2`$ up to first order in $`\beta ^1`$ (cf. Eq. 5.44), viz.
$$\begin{array}{ccc}\hfill \alpha _1& =\mathrm{\hspace{0.33em}1}\frac{2}{\beta \gamma }+O\left(\beta ^2\right)\hfill & \\ & & \\ \hfill \alpha _2& =\frac{1}{\gamma }+\frac{1}{\beta \gamma }\left(1\frac{2}{\gamma }\right)+O\left(\beta ^2\right).\hfill & \left(5.50\right)\hfill \end{array}$$
It is easy to show that for $`\alpha _1`$ and $`\alpha _2`$ the operations of taking the limit $`\beta 1`$ and taking the derivative with respect to $`l_0`$ can be interchanged. This simplifies the calculation of the coefficients which contain derivatives of these quantities. We now consider the first coefficient on the r.h.s. of Eq. (5.45) and use Eq. (5.50) to obtain the limit for large $`\beta `$:
$$\left(\alpha _11\right)\frac{\gamma p_0}{\rho _0}\frac{R_0^{}}{R_0^2}=\frac{B_0^{\mathrm{\hspace{0.17em}2}}}{4\pi \rho _0}\frac{R_0^{}}{R_0^{\mathrm{\hspace{0.17em}2}}}+O\left(\beta ^2\right).$$
$`\left(5.51\right)`$
The non-dimensional form is found by multiplication with $`\beta _mH_{pm}/g_m`$ which yields
$$\frac{B_0^{\mathrm{\hspace{0.17em}2}}}{4\pi \rho _0}\frac{R_0^{}}{R_0^{\mathrm{\hspace{0.17em}2}}}\frac{8\pi p_mH_{pm}}{B_m^{\mathrm{\hspace{0.17em}2}}g_m}=\frac{2\stackrel{~}{B}_0^{\mathrm{\hspace{0.17em}2}}}{\stackrel{~}{\rho }_0}\frac{\stackrel{~}{R}_0^{}}{\stackrel{~}{R}_0^{\mathrm{\hspace{0.17em}2}}}.$$
$`\left(5.52\right)`$
The tilde denotes dimensionless quantities (e.g. $`\stackrel{~}{\rho }_0=\rho _0/\rho _m`$) and we have used the hydrostatic relation $`p_m/g_m=\rho _mH_{pm}`$. A similar procedure is carried out for all other coefficients. Sometimes it is helpful to use the relation between the derivative along the flux tube and the derivative with respect to $`z`$ for quantities which only depend on height such as $`\rho _0`$, $`p_0`$ and $`B_0`$:
$$f^{}=\widehat{𝐥}_0f=\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\frac{df}{dz}.$$
$`\left(5.53\right)`$
An often needed quantity is the $`z`$-derivative of a scale height (external or internal) which can be written for the case of constant mean molecular weight, $`\overline{\mu }`$, in the form
$$\frac{dH_p}{dz}=\frac{d}{dz}\left(\frac{T}{\overline{\mu }g}\right)=\frac{}{\overline{\mu }g}\frac{dT}{dz}+\frac{T}{\overline{\mu }}\frac{d}{dz}\left(\frac{1}{g}\right)=\frac{sH_p}{z}$$
$`\left(5.54\right)`$
where we have used the equation of state for a perfect gas. $``$ is the gas constant and $`T`$ denotes the temperature. Another helpful relation is given by the condition of internal hydrostatic equilibrium, i.e. $`p_0/H_{p0}=\rho _0g_0`$. We give another example for the transformation of a coefficient which is not quite straightforward. Consider the large $`\beta `$ limit for the third term on the r.h.s. of Eq. (5.45):
$$\begin{array}{ccc}\hfill \frac{1}{\rho _0}& \left[\frac{\alpha _2p_0^{}}{H_{pe}}\left(\frac{\gamma p_0\alpha _2}{H_{pe}}\right)^{}\right]\frac{1}{\rho _0}\left\{\left[\frac{1}{\gamma }+\frac{1}{\beta \gamma }\left(1\frac{2}{\gamma }\right)\right]\frac{p_0^{}}{H_{pe}}\left[\frac{\gamma p_0}{H_{pe}}\left(\frac{1}{\gamma }+\frac{1}{\beta \gamma }\left(1\frac{2}{\gamma }\right)\right)\right]^{}\right\}\hfill & \\ & & \\ & =\frac{1}{\rho _0}\left\{\frac{p_0^{}}{\gamma H_{pe}}+\frac{p_0^{}}{\beta \gamma H_{pe}}\left(1\frac{2}{\gamma }\right)\left(\frac{p_0}{H_{pe}}\right)^{}\left[\frac{p_0}{\beta H_{pe}}\left(1\frac{2}{\gamma }\right)\right]^{}\right\}.\hfill & \left(5.55\right).\hfill \end{array}$$
We denote the superadiabaticity of the external gas by $`\delta =_{ad}`$, take $`_{ad}=\left(\gamma 1\right)/\gamma `$, such that $`1/\gamma =1+\delta `$ and write using Eqs. (5.53/54)
$$\frac{p_0^{}}{\gamma H_{pe}}\left(\frac{p_0}{H_{pe}}\right)^{}=\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\left[\frac{\delta }{H_{pe}}\frac{dp_0}{dz}+\frac{p_0}{H_{pe}^2}\left(\frac{H_{pe}}{H_{p0}}1\right)\frac{p_0s}{H_{pe}z_0}\right].$$
$`\left(5.56\right)`$
Inserting Eq. (5.56) and using $`dp_0/dz=p_0/H_{p0}=\rho _0g_0`$ we find for the r.h.s. of Eq. (5.55):
$$\begin{array}{ccc}\hfill \mathrm{}& =(\widehat{𝐠}_0\widehat{𝐥}_0)\frac{1}{\rho _0}\{\frac{\delta p_0}{H_{pe}H_{p0}}+\frac{p_0}{H_{pe}^2}(\frac{H_{pe}}{H_{p0}}1)\frac{p_0s}{H_{pe}z_0}\hfill & \\ & & \\ & (1\frac{2}{\gamma })\frac{p_0}{\beta \gamma H_{pe}H_{p0}}(1\frac{2}{\gamma })[\frac{1}{H_{pe}}\frac{d}{dz}\left(\frac{p_0}{\beta }\right)+\frac{p_0}{\beta }\frac{d}{dz}\left(\frac{1}{H_{pe}}\right)]\}\hfill & \\ & & \\ & =(\widehat{𝐠}_0\widehat{𝐥}_0)\{\frac{g_0}{H_{pe}}[\delta (1\frac{H_{p0}}{H_{pe}})+\frac{sH_{p0}}{z_0}]\hfill & \\ & & \\ & +(1\frac{2}{\gamma })\frac{1}{\rho _0H_{pe}}[\frac{d}{dz}\left(\frac{B_0^{\mathrm{\hspace{0.17em}2}}}{8\pi }\right)+\frac{B_0^{\mathrm{\hspace{0.17em}2}}}{8\pi H_{pe}}(\frac{H_{pe}}{\gamma H_{p0}}++\frac{sH_{pe}}{z_0})]\}.\hfill & \left(5.57\right)\hfill \end{array}$$
By multiplication with $`\beta _mH_{pm}/g_m`$ which can be rewritten as
$$\frac{\beta _mH_{pm}}{g_m}=\frac{\beta _mH_{pm}}{g_m\beta }\beta =\frac{H_{pm}^{\mathrm{\hspace{0.17em}2}}\rho _mB_0^{\mathrm{\hspace{0.17em}2}}}{g_0H_{p0}\rho _0B_m^{\mathrm{\hspace{0.17em}2}}}\beta $$
$`\left(5.58\right)`$
and
$$\frac{\beta _mH_{pm}}{g_m}=\frac{8\pi H_{pm}^{\mathrm{\hspace{0.17em}2}}\rho _m}{B_m^{\mathrm{\hspace{0.17em}2}}},$$
$`\left(5.59\right)`$
Eq. (5.57) is transferred to non-dimensional form:
$$\begin{array}{ccc}\hfill \mathrm{}& (\widehat{𝐠}_0\widehat{𝐥}_0)\frac{\stackrel{~}{B}_0^{\mathrm{\hspace{0.17em}2}}}{\stackrel{~}{\rho }_0\stackrel{~}{H}_{pe}}\{\frac{1}{\stackrel{~}{H}_{p0}}[\beta \delta \beta (1\frac{\stackrel{~}{H}_{po}}{\stackrel{~}{H}_{pe}})+\frac{\beta s\stackrel{~}{H}_{p0}}{\stackrel{~}{z}_0}]\hfill & \\ & & \\ & +(1\frac{2}{\gamma })[\frac{1}{\stackrel{~}{B}_0^{\mathrm{\hspace{0.17em}2}}}\frac{d\stackrel{~}{B}_0^{\mathrm{\hspace{0.17em}2}}}{d\stackrel{~}{z}}+\frac{1}{\stackrel{~}{H}_{pe}}(\frac{\stackrel{~}{H}_{pe}}{\gamma \stackrel{~}{H}_{p0}}++\frac{s\stackrel{~}{H}_{pe}}{\stackrel{~}{z}_0})]\}.\hfill & \left(5.60\right)\hfill \end{array}$$
In what follows we shall omit the tildes and tacitly assume that all quantities are non-dimensional unless the contrary is explicitly stated. The procedures described above have been applied to each term in Eqs. (5.45/47). For equipartition flux tubes in the deep solar convection zone we have $`H_{pe}/H_{p0}1=O\left(\beta ^1\right),\rho _{e0}/\rho _01=O\left(\beta ^1\right)`$ and we can take $`H_{pe}/H_{p0}1,\rho _{e0}/\rho _01`$ unless such a term is multiplied by $`\beta `$. Under these conditions we have
$$\beta \left(1\frac{H_{pe}}{H_{p0}}\right)+\mathrm{\hspace{0.17em}1}=\beta \left(\frac{\rho _{e0}}{\rho _0}1\right)+O\left(\beta ^2\right).$$
$`\left(5.61\right)`$
Another useful relation is
$$\frac{1}{B_0^{\mathrm{\hspace{0.17em}2}}}\frac{dB_0^{\mathrm{\hspace{0.17em}2}}}{dz}=\frac{8\pi }{B_0^{\mathrm{\hspace{0.17em}2}}}\frac{d}{dz}\left(p_{e0}p_0\right)=\frac{\beta }{H_{p0}}\left(1\frac{\rho _{e0}}{\rho _0}\right).$$
$`\left(5.62\right)`$
After some tedious algebra, Eqs. (5.45/5.47) in their final form are given by
$$\begin{array}{ccc}\hfill \omega ^2\frac{\rho _0}{B_0^{\mathrm{\hspace{0.17em}2}}}\eta & =\epsilon \{(\widehat{𝐠}_0\widehat{𝐥}_0)[\frac{1}{R_0H_{p0}}\beta (1\frac{\rho _{e0}}{\rho _0})\frac{2}{R_0^{\mathrm{\hspace{0.17em}2}}}\frac{dR_0}{dz}]\hfill & \\ & & \\ & +(\widehat{𝐠}_0\widehat{𝐥}_0)(\widehat{𝐠}_0\widehat{𝐧}_0)(\frac{s}{z_0H_{p0}}[\beta (1\frac{\rho _{e0}}{\rho _0})+\frac{2}{\gamma }]+\frac{1}{H_{p0}^{\mathrm{\hspace{0.17em}2}}}[\beta \delta \frac{1}{\gamma }\frac{\beta }{\gamma }(1\frac{\rho _{e0}}{\rho _0})])\}\hfill & \\ & & \\ & +\epsilon ^{}\left\{\frac{2}{R_0}+\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\frac{1}{H_{p0}}\left[\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)+\frac{2}{\gamma }\right]\right\}\hfill & \\ & & \\ & +\eta \{(\widehat{𝐠}_0\widehat{𝐥}_0)^2(\frac{s}{z_0H_{p0}}[\beta (1\frac{\rho _{e0}}{\rho _0})+\frac{2}{\gamma }]+\frac{1}{H_{p0}^{\mathrm{\hspace{0.17em}2}}}[\beta \delta \frac{1}{\gamma }\frac{\beta }{\gamma }(1\frac{\rho _{e0}}{\rho _0})])\hfill & \\ & & \\ & +(\widehat{𝐠}_0\widehat{𝐧}_0)\frac{1}{R_0H_{p0}}[\beta (1\frac{\rho _{e0}}{\rho _0})+\frac{2}{\gamma }]\}\hfill & \\ & & \\ & +\eta ^{}\left\{\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\frac{1}{H_{p0}}\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)\right\}+\mathrm{\hspace{0.33em}2}\eta ^{\prime \prime }\hfill & \left(5.63\right)\hfill \end{array}$$
for the longitudinal displacement and
$$\begin{array}{ccc}\hfill \omega ^2\frac{\rho _0}{B_0^{\mathrm{\hspace{0.17em}2}}}\epsilon & =\eta \{(\widehat{𝐠}_0\widehat{𝐥}_0)[\frac{2}{R_0^{\mathrm{\hspace{0.17em}2}}}\frac{dR_0}{dz}+\frac{2}{\gamma R_0H_{p0}}\frac{1}{R_0H_{p0}}\beta (1\frac{\rho _{e0}}{\rho _0})]\hfill & \\ & & \\ & +\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\frac{1}{H_{p0}^{\mathrm{\hspace{0.17em}2}}}\left[\beta \delta \frac{sH_{p0}}{z_0}+\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)\right]\hfill & \\ & & \\ & +\frac{\beta rm}{H_{p0}}(\frac{2}{𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0}[(\widehat{𝐱}\widehat{𝐥}_0)w_x(\widehat{𝐠}_0\widehat{𝐥}_0)w_z\frac{𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐥}_0}{R_0}]+\frac{1}{2\gamma H_{p0}}(\widehat{𝐠}_0\widehat{𝐥}_0))\}\hfill & \\ & & \\ & +\eta ^{}\left\{\frac{4}{R_0}\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)\frac{2}{\gamma H_{p0}}+\frac{\beta rm}{2H_{p0}}\right\}\hfill & \\ & & \\ & +\epsilon \{\frac{2}{R_0^{\mathrm{\hspace{0.17em}2}}}+(\widehat{𝐠}_0\widehat{𝐧}_0)\frac{4}{\gamma R_0H_{p0}}+(\widehat{𝐠}_0\widehat{𝐧}_0)^2\frac{1}{H_{p0}^{\mathrm{\hspace{0.17em}2}}}[\beta \delta +\frac{1}{\gamma }(1\frac{2}{\gamma })\frac{sH_{p0}}{z_0}\beta (1\frac{\rho _{e0}}{\rho _0})]\hfill & \\ & & \\ & +\frac{\beta rm}{H_{p0}}(\frac{2}{𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0}[(\widehat{𝐱}\widehat{𝐧}_0)w_x(\widehat{𝐠}_0\widehat{𝐧}_0)w_z]+\frac{1}{2\gamma H_{p0}}(\widehat{𝐠}_0\widehat{𝐧}_0)\frac{1}{2R_0})\}\hfill & \\ & & \\ & +\epsilon ^{}\left\{\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)\frac{1}{H_{p0}}\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)\frac{2\beta rm\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐥}_0\right)}{H_{p0}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}\right\}\hfill & \\ & & \\ & +\mathrm{\hspace{0.33em}2}\epsilon ^{\prime \prime }i\omega \epsilon \frac{2\beta rm}{H_{p0}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)}\hfill & \left(5.64\right)\hfill \end{array}$$
for the displacement in normal direction. In Eq. (5.64) we have abbreviated $`\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)m`$ and used the equilibrium condition, Eq. (5.4), in some places. In the special case of a horizontal flux tube in static equilibrium and uniform gravity ($`R_0\mathrm{}`$, $`\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)=0`$, $`\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)=1`$, $`s=0`$, $`𝐯_e=0`$, $`\rho _0=\rho _{e0}`$, $`B_0^{\mathrm{\hspace{0.17em}2}}=\rho _0=1`$) Eqs. (5.63/64) take the form
$$\begin{array}{ccc}\hfill \omega ^2\eta & =\mathrm{\hspace{0.33em}2}\eta ^{\prime \prime }\epsilon ^{}\frac{2}{\gamma }\hfill & \\ & & \\ \hfill 2\omega ^2\epsilon & =\mathrm{\hspace{0.33em}2}\epsilon ^{\prime \prime }+\eta ^{}\frac{2}{\gamma }+\epsilon \left(\beta \delta +\frac{1}{\gamma }\frac{2}{\gamma ^2}\right).\hfill & \left(5.65\right)\hfill \end{array}$$
In order to compare with the result of Spruit and van Ballegooijen (1982) we have multiplied the inertial term in the equation for $`\epsilon `$ by a factor 2 intending to describe the enhanced inertia of the flux tube with respect to transversal motions in the same way as done by these authors. Since the coefficients are constant we are permitted to write $`(\eta ,\epsilon )\mathrm{exp}\left(ikl_0\right)`$ with wavenumber $`k`$ in the longitudinal (here: horizontal) direction. Inserting this into Eq. (5.65) we obtain the dispersion relation
$$\omega ^4+\omega ^2\left[3k^2+\frac{\beta \delta }{2}+\frac{1}{2\gamma }\left(1\frac{2}{\gamma }\right)\right]+\mathrm{\hspace{0.33em}2}k^2\left(k^2\frac{\beta \delta }{2}\frac{1}{2\gamma }\right)=\mathrm{\hspace{0.33em}0}$$
$`\left(5.66\right)`$
which is identical to the relation found by Spruit and van Ballegooijen (1982, their Eq. 39).
5.5 Symmetric form for static equilibrium
If we ignore the drag force there is no dissipation in the system and the force operator is self-adjoint (Spruit, 1981b). Hence, in the case $`𝐯_{e}^{}{}_{0}{}^{}=0`$ (static equilibrium) the perturbation equations (5.63/64) must lead to a self-adjoint eigenvalue problem with real eigenvalues. <sup>*</sup><sup>*</sup>Since the drag force is quadratic in the relative velocity between flux tube and environment, its perturbation vanishes if $`𝐯_{e}^{}{}_{0}{}^{}=0`$. Symbolically, we can write Eqs. (5.64/65) in the form
$$\omega ^2\left(\begin{array}{cc}& \eta \hfill \\ & \epsilon \hfill \end{array}\right)=𝐅\left(\begin{array}{cc}& \eta \hfill \\ & \epsilon \hfill \end{array}\right)$$
$`\left(5.67\right)`$
where the operator F is defined by
$$𝐅\left(\begin{array}{cc}& \eta \hfill \\ & \epsilon \hfill \end{array}\right)=\left(\begin{array}{cc}& A\epsilon +B\epsilon ^{}+C\eta +D\eta ^{}+E\eta ^{\prime \prime }\hfill \\ & F\eta +G\eta ^{}+H\epsilon +D\epsilon ^{}+E\epsilon ^{\prime \prime }\hfill \end{array}\right).$$
$`\left(5.68\right)`$
The functions $`A\left(l_0\right),\mathrm{},H\left(l_0\right)`$ can be obtained from the r.h.s. of Eqs. (5.64/65) by multiplication with $`B_0^{\mathrm{\hspace{0.17em}2}}/\rho _0`$ and setting $`𝐯_{e}^{}{}_{0}{}^{}=0`$. Using the equilibrium condition given by Eq. (5.4) is readily shown that in this (static) case we have $`B=G`$. However, the operator F is not yet written in a form in which its self-adjointness becomes apparent. Such a form is necessary for application of variational methods like the energy principle (Bernstein et al., 1958). It is also useful for numerical calculations since it leads to symmetric matrices which are much easier to deal with. A symmetric form can be obtained by transforming the eigenvector as
$$\left(\begin{array}{cc}& \eta \hfill \\ & \epsilon \hfill \end{array}\right)=\left(\begin{array}{cc}& ax\hfill \\ & by\hfill \end{array}\right)$$
$`\left(5.69\right)`$
where $`x`$ and $`y`$ are the new variables and the functions $`a\left(l_0\right)`$ and $`b\left(l_0\right)`$ are determined such that the eigenvalue problem given by Eqs. (5.67/68) attains the form
$$\omega ^2\left(\begin{array}{cc}& x\hfill \\ & y\hfill \end{array}\right)=\left(\begin{array}{cc}& Ky+\frac{1}{2}\left[\left(By\right)^{}+By^{}\right]+Mx+\frac{1}{2}\left[\left(Ex\right)^{\prime \prime }+Ex^{\prime \prime }\right]\hfill \\ & Kx\frac{1}{2}\left[\left(Bx\right)^{}+Bx^{}\right]+Ny+\frac{1}{2}\left[\left(Ey\right)^{\prime \prime }+Ey^{\prime \prime }\right]\hfill \end{array}\right)$$
$`\left(5.70\right)`$
which possesses the same eigenvalues as Eqs. (5.67/68). Inserting Eq. (5.69) into Eq. (5.67) leads to the following pair of equations:
$$\begin{array}{ccc}\hfill \omega ^2x=y& \frac{b}{a}\left(A+B\frac{b^{}}{b}\frac{1}{2}B^{}\right)+\frac{b}{2a}\left[\left(By\right)^{}+By^{}\right]+x\left(C+D\frac{a^{}}{a}+E\frac{a^{\prime \prime }}{a}\frac{1}{2}E^{\prime \prime }\right)\hfill & \\ & & \\ & +x^{}\left(D+2E\frac{a^{}}{a}E^{}\right)+\frac{1}{2}\left[\left(Ex\right)^{\prime \prime }+Ex^{\prime \prime }\right]\hfill & \left(5.71\right)\hfill \end{array}$$
$$\begin{array}{ccc}\hfill \omega ^2y=x& \frac{a}{b}\left(FB\frac{a^{}}{a}+\frac{1}{2}B^{}\right)\frac{a}{2b}\left[\left(Bx\right)^{}+Bx^{}\right]+y\left(H+D\frac{b^{}}{b}+E\frac{b^{\prime \prime }}{b}\frac{1}{2}E^{\prime \prime }\right)\hfill & \\ & & \\ & +y^{}\left(D+2E\frac{b^{}}{b}E^{}\right)+\frac{1}{2}\left[\left(Ey\right)^{\prime \prime }+Ey^{\prime \prime }\right].\hfill & \left(5.72\right)\hfill \end{array}$$
By comparison with the desired form given by Eq. (5.70) we find that the functions $`a\left(l_0\right)`$ and $`b\left(l_0\right)`$ are determined by the requirement that the terms which multiply $`x^{}`$ in Eq. (5.71) and $`y^{}`$ in Eq. (5.72), respectively, must vanish:
$$\begin{array}{ccc}\hfill D+\mathrm{\hspace{0.33em}2}E\frac{a^{}}{a}E^{}& =\mathrm{\hspace{0.33em}0}\hfill & \\ & & \\ \hfill D+\mathrm{\hspace{0.33em}2}E\frac{b^{}}{b}E^{}& =\mathrm{\hspace{0.33em}0}.\hfill & \left(5.73\right)\hfill \end{array}$$
Consequently, we have $`a^{}/a=b^{}/b`$ and since $`a`$ and $`b`$ may be multiplied by any constant number we can take $`ab`$ without loss of generality. Eq. (5.73) may be integrated to give
$$\mathrm{ln}a=\frac{1}{2}\mathrm{ln}E\frac{D}{2E}𝑑l_0.$$
$`\left(5.74\right)`$
The coefficient functions $`M`$ and $`N`$ in Eq. (5.70) are determined using Eq. (5.73) which leads to the condition
$$\begin{array}{ccc}\hfill M& =C+\frac{\left(E^{}\right)^2D^2}{4E}+E\left(\frac{E^{}D}{2E}\right)^{}\frac{E^{\prime \prime }}{2}\hfill & \\ & & \\ \hfill N& =H+\frac{\left(E^{}\right)^2D^2}{4E}+E\left(\frac{E^{}D}{2E}\right)^{}\frac{E^{\prime \prime }}{2}.\hfill & \left(5.75\right)\hfill \end{array}$$
The requirement that the function which multiplies $`y`$ in the first term on the r.h.s. of Eq. (5.71) is identical to that which multiplies $`x`$ in the first term on the r.h.s. of Eq. (5.72) in order to conform with Eq. (5.70) leads to
$$A+B\frac{b^{}}{b}\frac{B^{}}{2}=K=FB\frac{a^{}}{a}+\frac{B^{}}{2}.$$
$`\left(5.76\right)`$
This implies the compatibility relation
$$FA=\frac{B\left(E^{}D\right)}{E}B^{}$$
$`\left(5.77\right)`$
which must be fulfilled if the system is self-adjoint. Hence, we may use Eq. (5.77) in order to check the correctness of the algebraic manipulations and the consistency of the approximation which led to Eqs. (5.63/64). It turns out, after some lengthy algebra, that Eq. (5.77) indeed is valid for our problem, i.e. the self-adjointness of the problem is reflected in the symmetric structure of the resulting system of equations. Defining
$$\xi \left(\begin{array}{cc}& x\hfill \\ & y\hfill \end{array}\right)$$
it is easy to show that Eq. (5.70) indeed is selfadjoint, i.e. the equality
$$\stackrel{~}{\xi }𝐆\left(\xi \right)𝑑V=\xi 𝐆\left(\stackrel{~}{\xi }\right)𝑑V$$
$`\left(5.78\right)`$
holds. G denotes the operator on the r.h.s. of Eq. (5.70), the tilde indicates the complex conjugate, and the integration is taken over an appropriate volume. Since integrations by parts are necessary to demonstrate the validity of Eq. (5.78) the volume of integration must be chosen large enough such that the displacements can be assumed to vanish at the boundaries.
5.6 Horizontal tubes with vertical external flow
The preceding sections give a basis for the determination of the stability properties of rather general flux tube equilibrium structures. For realistic convection zone models and flow patterns an equilibrium can only be determined by numerical means (e.g. van Ballegooijen, 1982a; Anton, 1984; Anton and Schüssler, unpublished) and thus the perturbation equations, Eqs. (5.63/64), have to be transformed into a numerically tractable matrix eigenvalue problem. Such an undertaking is outside the scope of the present investigation and has to be deferred to future work. In what follows we shall consider a couple of analytically tractable cases instead to which we nevertheless attribute some general relevance.
As a first simple application of the formalism we consider the case $`R_0\mathrm{}`$, $`(\widehat{𝐠}_0\widehat{𝐥}_0)=0`$, $`(\widehat{𝐠}_0\widehat{𝐧}_0)=1`$, $`s=0`$, $`𝐯_e=v(z)\widehat{𝐳}`$, i.e. a horizontal, straight flux tube whose equilibrium is determined by a balance between the buoyancy and drag forces (cf. Eq. 5.4/5):
$$\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)=\beta rm.$$
$`\left(5.79\right)`$
Here we have $`m=\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)=\text{sgn}\left(v\left(z_0\right)\right)`$, i.e. $`m=1`$ describes an upflow, $`m=1`$ a downflow. Defining
$$\frac{C_Dv_0^2}{\pi a_0g_0}q$$
$`\left(5.80\right)`$
with $`v_0=\left|v\left(z_0\right)\right|`$ we obtain the equilibrium density contrast as function of the positive number $`q`$ as
$$1\frac{\rho _{e0}}{\rho _0}=\frac{q}{m+q}.$$
$`\left(5.81\right)`$
It is clear from Eq. (5.81) that $`q`$ is limited to $`q<1`$ in the case of a downflow ($`m=1`$) since the tube becomes completely evacuated for $`q1`$ and the buoyancy force cannot be increased beyond that limit. On the other hand, for an upflow ($`m=1`$) the range of values for $`q`$ is not restricted because the internal density can be made large enough to balance any upward directed drag force.
If we insert Eq. (5.79) into the perturbation equations, Eqs. (5.63/64), and have regard to the properties of the equilibrium as described in the beginning of this section we obtain the following pair of equations:
$$\begin{array}{ccc}\hfill \omega ^2\eta & =\epsilon ^{}\left(\beta rm+\frac{2}{\gamma }\right)+\mathrm{\hspace{0.33em}2}\eta ^{\prime \prime }\hfill & \left(5.82\right)\hfill \\ & & \\ \hfill \omega ^2\epsilon & =\eta ^{}\left(\frac{\beta rm}{2}+\frac{2}{\gamma }\right)+\epsilon \left[\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)+\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)\right]\hfill & \\ & & \\ & +\mathrm{\hspace{0.33em}2}\epsilon ^{\prime \prime }i\omega \epsilon \frac{2\beta r}{v_0}\hfill & \left(5.83\right)\hfill \end{array}$$
Here we have written $`\left|v\left(z_0\right)\right|v_0`$ and the notation $`dv_0/dz`$ has been used to abbreviate $`d\left|v\left(z\right)\right|/dz`$ taken at $`z=z_0`$. Let us first consider displacements which do not depend on the horizontal coordinate $`l_0x`$. All derivatives vanish in this case and we find $`\eta 0`$ from Eq. (5.82). The equation for $`\epsilon `$ is transformed into the dispersion relation
$$\omega ^2i\omega \frac{2\beta r}{v_0}+\left[\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)+\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)\right]=\mathrm{\hspace{0.33em}0}.$$
$`\left(5.84\right)`$
Writing
$$\frac{\beta r}{v_0}S0$$
$`\left(5.85\right)`$
and
$$\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)+\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)T$$
$`\left(5.86\right)`$
we find from Eq. (5.84)
$$\omega _\pm =iS\pm \left(S^2T\right)^{1/2}$$
$`\left(5.87\right)`$
and by multiplication with the imaginary unit we obtain
$$i\omega _\pm =S\left(S^2+T\right)^{1/2}.$$
$`\left(5.88\right)`$
We now have to consider two cases. Firstly, if $`S^2+T<0`$, the square root in Eq. (5.88) is imaginary and the growth rate (real part of $`i\omega _\pm `$) becomes equal to $`S`$. Consequently, we have a damped oscillation in this case and the equilibrium is stable. In the second case, $`S^2+T>0`$, both the square root and $`i\omega _\pm `$ are real. Depending on the sign of the latter, the perturbation will grow or decay monotonically. It is easy to see that the stability in this case depends solely on the sign of $`T`$: If $`T<0`$ we have $`i\omega _\pm <0`$ and the displacement decays, if $`T>0`$ we have $`i\omega _{}>0`$ and a monotonic growth of the perturbation ensues. Combining the results for both cases we find
$$T\{\begin{array}{cc}>0:\text{ monotonic growth}\hfill & \\ & \\ <0:\{\begin{array}{cc}S^2+T<0:\hfill & \text{damped oscillation}\hfill \\ & \\ S^2+T>0:\hfill & \text{monotonic decay}\hfill \end{array}\hfill & \end{array}$$
Hence, the stability of the flux tube depends only on the sign of $`T`$ while the imaginary term in Eq. (5.83) influences the way in which stable perturbations decay, monotonically or oscillatory. This behavior is plausible from the fact that this term, being proportional to $`\dot{\epsilon }`$, describes the drag force caused by the perturbation itself and therefore only has a damping effect. Similar to a damped oscillator the case of ‘creeping motion’ is achieved if the damping rate (described by $`S`$) exceeds a critical value. The equilibrium is unstable if $`T>0`$, i.e. if the superadiabaticity satisfies the inequality
$$\beta \delta >\frac{1}{\gamma }\left(\frac{2}{\gamma }1\right)+\beta rm\left(\frac{1}{2\gamma }\frac{2}{v_0}\frac{dv_0}{dz}\right).$$
$`\left(5.89\right)`$
Consequently, positive terms on the r.h.s. of Eq. (5.89) have a stabilizing influence. For constant velocity, i.e. $`dv_0/dz=0`$, we find that an upflow ($`m=1`$) has a stabilizing effect while a downflow ($`m=1`$) is destabilizing. This behavior is caused by the changes in flux tube radius and external density. An upward displacement leads to an expansion of the tube and a decrease of external density. Consequently, the drag force decreases (cf. Eq. 5.33). Similarly, for a downward displacement we find an increase of the drag force. If we now have a flux tube which is in equilibrium with a downflow, for both upward and downward displacements the perturbation of the (downward) drag force tends to increase the displacement and thus favors instability while the reverse is true for an upflow.
For $`dv_0/dz0`$, the situation is more complicated. Whether the change of velocity with height acts stabilizing or destabilizing depends on its sign and on the direction of the flow. For example, an upflow ($`m=1`$) whose velocity increases with height ($`dv_0/dz>0`$) leads to a perturbation of the drag force which tends to increase an initial displacement and thus favors instability. The influence of the velocity term on the stability properties in the case of purely vertical displacements is summarized Tab. 1:
| | $`m`$ | $`dv_0/dz`$ | $`1/2\gamma `$ | $`\left(1/2\gamma \right)\left(2/v_0\right)\left(dv_0/dz\right)`$ |
| --- | --- | --- | --- | --- |
| | 1 | $`>0`$ destabilizing | stabilizing | ? |
| | 1 | $`<0`$ stabilizing | stabilizing | stabilizing |
| | -1 | $`>0`$ stabilizing | destabilizing | ? |
| | -1 | $`<0`$ destabilizing | destabilizing | destabilizing |
Tab. 1: Influence of the velocity-related terms in the stability criterion (Eq. 5.89) for horizontal magnetic flux tubes and purely vertical displacement (infinite longitudinal wavelength).
The combined effect of both terms in Eq. (5.89) which are affected by the external velocity is indicated in the last column. A question mark indicates that one term is stabilizing and the other destabilizing such that it depends on their relative sizes which one dominates. In dimensionalized form, the term involving the velocity derivative can be written as
$$\frac{2}{\stackrel{~}{v}_0}\frac{d\stackrel{~}{v}_0}{d\stackrel{~}{z}}\frac{2H_{p0}}{v_0}\frac{dv_0}{dz}\frac{2H_{p0}}{H_v}$$
$`\left(5.90\right)`$
where $`H_v`$ is the scale height of the external velocity. If we assume a flow with constant mass flux density, $`\rho v`$, and assume further that the temperature varies much more slowly with height than the density, we have $`p_0v_0=`$ const. and thus $`H_{p0}H_v`$. The relation between the two terms in the second bracket on the r.h.s. of Eq. (5.89) is then given by
$$\frac{\left(1/2\gamma \right)}{\left(2/v_0\right)\left(dv_0/dz\right)}\frac{1}{4\gamma }=\frac{3}{20}$$
$`\left(5.91\right)`$
for $`\gamma =5/3`$. We see that for a flow with constant mass flux density the stability properties are mainly determined by the velocity gradient. Since density decreases with height, the flow velocity increases and we see from the table above that an upflow has a destabilizing influence, while a downflow stabilizes.
If we now consider all terms in Eq. (5.89) we find that a flow can stabilize/destabilize a convectively unstable/stable flux tube. For a flux tube in temperature equilibrium ($`T_0=T_{e0}`$, $`\beta rm=1`$) and a downflow with constant mass flux density instability requires
$$\beta \delta >\frac{2}{\gamma ^2}\frac{3}{2\gamma }+2=\mathrm{\hspace{0.33em}1.82}$$
$`\left(5.92\right)`$
Since equipartition flux tubes in the deep convection zone of the Sun have $`\beta \delta 3.6`$ (Spruit and van Ballegooijen, 1982) they cannot be stabilized by a flow of this kind. On the other hand, Eq. (5.92) predicts stability for flux tubes located in a downflow within an overshoot region which have $`\beta \delta 0`$.
The case of purely vertical displacements which do not depend on the horizontal coordinate treated so far excludes an important destabilizing mechanism, the flow from the crests to the troughs of a wavelike disturbance of the flux tube which is the mechanism for the so-called Parker instability (Parker, 1966). If we allow for a dependence of the displacements $`\eta `$ and $`\epsilon `$ on the horizontal coordinate, $`x`$, we may write
$$\eta ,\epsilon e^{ikx}$$
$`\left(5.93\right)`$
with (real) wavenumber $`k`$ since the coefficients in Eqs. (5.82/83) are constant. Using the symbols $`\eta `$ and $`\epsilon `$ again for the (constant) amplitudes of the perturbations we obtain by inserting Eq. (5.93) into Eqs. (5.82/83):
$$\begin{array}{ccc}\hfill \omega ^2\eta & =ik\epsilon \left(\beta rm+\frac{2}{\gamma }\right)\mathrm{\hspace{0.33em}2}k^2\eta \hfill & \left(5.94\right)\hfill \\ & & \\ \hfill \omega ^2\epsilon & =ik\eta \left(\frac{\beta rm}{2}+\frac{2}{\gamma }\right)+\epsilon \left[\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)+\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)\right]\hfill & \\ & & \\ & \mathrm{\hspace{0.33em}2}k^2\epsilon i\omega \epsilon \frac{2\beta r}{v_0}.\hfill & \left(5.95\right)\hfill \end{array}$$
This linear, homogeneous system of equations has non-trivial solutions for eigenvalues $`\omega `$ which satisfy the dispersion relation
$$\begin{array}{ccc}\hfill \left(\omega ^22k^2\right)& \left[\omega ^2i\omega \frac{2\beta r}{v_0}2k^2+\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)+\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)\right]\hfill & \\ & & \\ & k^2\left(\frac{2}{\gamma }+\frac{\beta rm}{2}\right)\left(\frac{2}{\gamma }+\beta rm\right)=\mathrm{\hspace{0.33em}0}.\hfill & \left(5.96\right)\hfill \end{array}$$
Eq. (5.96) can be written as a fourth order polynomial in $`i\omega \widehat{\omega }`$ with real coefficients:
$$\widehat{\omega }^4+A\widehat{\omega }^3+B\widehat{\omega }^2+C\widehat{\omega }+D=\mathrm{\hspace{0.33em}0}$$
$`\left(5.97\right)`$
with
$$\begin{array}{cc}\hfill A& =\frac{2\beta r}{v_0}\hfill \\ \hfill B& =\mathrm{\hspace{0.33em}4}k^2\beta \delta \frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)\hfill \\ \hfill C& =\frac{4k^2\beta r}{v_0}\hfill \\ \hfill D& =\mathrm{\hspace{0.33em}2}k^2\left[2k^2\beta \delta \frac{1}{\gamma }\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}+\frac{1}{\gamma }\right)\frac{\left(\beta r\right)^2}{4}\right]\hfill \end{array}$$
Since the coefficients are real, the roots of this polynomial can be real numbers and pairs of complex conjugates. A positive, real root means monotonic instability, i.e. exponential growth of the perturbation, while a complex root with positive real part represents overstability, i.e. oscillations or waves with exponentially growing amplitude. In order to establish a sufficient condition for monotonic instability we may utilize Descartes’ sign rule:
The number of positive, real roots of a fourth order polynomial with real coefficients is smaller than or equal to the number of sign changes in the sequence of its coefficients. The difference is an even number.
Since in our case we have $`A0`$ and $`C0`$ the number of sign changes is determined by the signs of $`B`$ and $`D`$ as indicated in Tab. 2:
| | Case | $`B`$ | $`D`$ | Sign changes | Positive, real roots |
| --- | --- | --- | --- | --- | --- |
| | 1 | $`>0`$ | $`<0`$ | 1 | 1 |
| | 2 | $`>0`$ | $`>0`$ | 0 | 0 |
| | 3 | $`<0`$ | $`<0`$ | 3 | 1 or 3 |
| | 4 | $`<0`$ | $`>0`$ | 2 | 0 or 2 |
Tab. 2: Estimate of the number of positive, real roots (i.e. monotonically unstable modes) of Eq. (5.97) using the number of sign changes in the sequence of its coefficients (Descartes’ rule). Since $`A`$ and $`C`$ are positive, this number depends only on the signs of $`B`$ and $`D`$.
In the cases 1 and 3 we have at least one positive root. Consequently, $`D<0`$ is a sufficient condition for monotonic instability. In case 2 there is no monotonic instability (but possibly overstability) while for case 4 no definite statement can be made at this stage. The condition $`D<0`$ can be transformed into a condition for the wavenumber, $`k`$:
$$k^2<k_0^2\frac{1}{2}\left[\beta \delta +\frac{1}{\gamma }+\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}+\frac{1}{\gamma }\right)+\frac{\left(\beta r\right)^2}{4}\right]$$
$`\left(5.98\right)`$
Since any value of $`k`$ is allowed in our cartesian model (in contrast to a spherical system where the wavelength cannot be larger than the circumference), the flux tube is monotonically unstable if only $`k_0^2>0`$, i.e.
$$\beta \delta >\frac{1}{\gamma }+\beta rm\left(\frac{1}{\gamma }\frac{2}{v_0}\frac{dv_0}{dz}\right)\frac{\left(\beta r\right)^2}{4}.$$
$`\left(5.99\right)`$
We immediately see that a flow with constant velocity ($`dv_0/dz=0`$) cannot stabilize a tube within a superadiabatic environment: For an upflow ($`m=+1`$) both remaining velocity terms are negative while it is easy to see that for a downflow ($`m=1`$) the r.h.s. of Eq. (5.99) is always negative since $`\gamma >1`$. If the flow speed depends on depth, stabilization is possible. A downflow whose speed decreases with depth such that the mass flux density stays constant \[$`m=1`$, $`\left(2/v_0\right)\left(dv_0/dz\right)=2`$\] leads to positive values of the r.h.s. of Eq. (5.99) provided that $`0.23<\beta r<10.16`$, i.e. within a certain range of values for the parameter $`\beta r`$.
Due to the last term in Eq. (5.99) a flux tube is always monotonically unstable for sufficiently large values of $`\beta r`$, irrespective of the flow direction. This is caused by the change of the flux tube radius in the course of a wave-like displacement, a mechanism which can be understood with aid of the heuristic approach followed in Sec. 5.7.3.
For the example discussed in conjunction with the criterion Eq. (5.92) (temperature equilibrium, $`T_0=T_{e0}`$, and downflow with constant mass flux density, $`\beta rm=1`$) we find from Eq. (5.99) the criterion $`\beta \delta >1.75`$ for monotonic instability, i.e. a slightly smaller critical value than that obtained for the case of purely vertical displacements. Generally we find that the condition given in Eq. (5.99) leads to instability for smaller values of the $`\beta \delta `$, i.e. monotonic instability is easier to excite for wave-like perturbations (even if $`k0`$) than for those with $`k0`$. This is similar to the case without external flow treated by Spruit and van Ballegooijen (1982) whose results may be recovered by setting $`\beta r=0`$. However, in our case flux tubes embedded in a subadiabatic region ($`\beta \delta 0`$) are stable with respect to wave-like perturbations while in the case without external flow instability results if $`\beta \delta >1/\gamma `$.
The general properties of the roots of Eq. (5.97) unfortunately are not as easily obtained as the sufficient condition for positive, real roots. However, since the coefficients $`A`$ and $`C`$ result from the imaginary term in Eq. (5.83) which mainly has a damping effect and does not affect the stability conditions for the case $`k=0`$, we may conjecture that this is the case for $`k0`$ as well. We therefore consider a reduced dispersion relation by setting $`A=C=0`$ in Eq. (5.97). The conclusions drawn from the discussion of this equation have been verified by (numerical) determination of the roots of the full equation (5.97). <sup>*</sup><sup>*</sup>The general algorithm for the determination of roots of quartic equations in closed form is not used here because, given the complicated structure of the coefficients, it leads to lengthy algebraic expressions which give no direct insight in the properties of the stability problem. For $`A=C=0`$ Eq. (5.97) is transformed into a biquadratic equation which is readily solved, viz.
$$\widehat{\omega }^2=\frac{B}{2}\pm \left(\frac{B^2}{4}D\right)^{1/2}.$$
$`\left(5.100\right)`$
We see immediately that we recover the sufficient condition for monotonic instability, $`D<0`$, since in this case always one solution $`\widehat{\omega }^2>0`$ exists. A second possibility for monotonic instability is the case $`B<0`$ and $`0DB^2/4`$ while the case $`B>0`$, $`0DB^2/4`$ leads to stable solutions ($`\widehat{\omega }^2<0`$). Complex roots which signify oscillatory instability (overstability) appear for $`B^2/4D<0`$, irrespective of the sign of $`B`$. Since we have
$$\frac{B^2}{4}D=k^2\left[\frac{4}{\gamma ^2}+\frac{\left(\beta r\right)^2}{2}+\frac{3}{\gamma }\beta rm\right]+\frac{1}{4}\left[\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)+\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)\right]^2$$
$`\left(5.101\right)`$
complex roots exist only if the first term on the r.h.s. is negative. This requires $`m=1`$ (downflow) and
$$\frac{2}{\gamma }<\beta r<\frac{4}{\gamma }.$$
$`\left(5.102\right)`$
For $`\gamma =5/3`$ this range is given by $`1.2<\beta r<2.4`$. If both conditions are satisfied we see from Eq. (5.101) that the wavenumber can always be made large enough to give $`B^2/4D<0`$. Hence, the condition for overstability is $`k^2>k_1^2`$ with
$$k_1^2=\frac{1}{4}\left[\frac{4}{\gamma ^2}\frac{\left(\beta r\right)^2}{2}\frac{3}{\gamma }\beta rm\right]^1\left[\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)+\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)\right]^2.$$
$`\left(5.103\right)`$
Note that in order to observe the limits set by the approximation of slender flux tubes the radius of the flux tube has to be much smaller than the perturbation wavelength. This can be a severe restriction if overstability appears only for very small wavelength (very large wavenumber). If we have $`k^2>k_1^2`$ the four solutions of Eq. (5.100) are given by
$$\widehat{\omega }=\pm \frac{1}{2}\left(2D^{1/2}B\right)^{1/2}\pm \frac{1}{2}i\left(2D^{1/2}+B\right)^{1/2}.$$
$`\left(5.104\right)`$
We see that if the conditions for complex roots are satisfied there are always unstable solutions, i.e. roots with a positive real part. The overstability is caused by the drag force which is added to the restoring force due to magnetic tension such that an oscillation with growing amplitude results. The conditions for overstability depend only on the sign and the amplitude of the velocity (cf. Eq. 5.102). In particular, they are independent of the velocity gradient, a fact which is due to the linear approach.
Numerical solutions of the exact dispersion relation, Eq. (5.97), basically confirm the criteria derived on the basis of the discussion of the reduced form, Eq. (5.100). In most cases, finite values of $`A`$ and $`C`$ only affect the growth or decay rates or change oscillatory into monotonic decay. A notable exception is the result that finite values of $`A`$ and $`C`$ lead to positive growth rates of oscillatory modes whenever the conditions $`m=1`$ and Eq. (5.102) are fulfilled, even for wave numbers which give $`B^2/4D>0`$. The second possibility for monotonic instability mentioned above, $`B<0`$ and $`0DB^2/4`$, does not appear in practice since it requires that the conditions $`m=1`$ and Eq. (5.102) are fulfilled which immediately lead to oscillatory instability. Hence, we can summarize the stability properties as follows:
A horizontal flux tube whose equilibrium is determined by a balance between buoyancy and drag force exerted by an external, vertical flow is monotonically unstable if the condition $`D<0`$ (cf. Eq. 5.99) is satisfied. In particular, flux tube equilibria with a sufficiently large value of $`\beta r`$ are always unstable, irrespective of the flow direction. A flux tube may exhibit oscillatory instability (overstability) if the external flow is a downflow and the value of $`\beta r`$ is within a certain range (cf. Eq. 5.102).
Oscillatory instability can be excited whenever the conditions $`m=1`$ and $`2/\gamma <\beta r<4/\gamma `$ are fulfilled. This excludes isothermal flux tubes which require $`\beta rm=1`$. Overstability cannot be stabilized by the stratification, however subadiabatic it may be. Consequently, it is potentially relevant for all flux tubes, irrespective of their location in convection zones, overshoot layers or regions in radiative equilibrium. Since $`\beta r`$ depends on the field strength and radius of the tube as well as on the flow velocity, for given values of two of these parameters Eq. (5.102) defines a range of the third parameter which leads to overstability. The tube radius, however, must always be small enough to stay consistent with the utilization of the approximation of slender flux tubes.
While the convective and Parker-type instabilities are most easily excited for large wavelength along the flux tube, overstable modes have their largest growth rates for small wavelengths. Consequently, we expect that overstability is not only relevant for horizontal flux tubes but also for more general forms of equilibrium. If we restrict ourselves to perturbation wavelengths which are much smaller than the wavelength of the equilibrium path of a flux tube a kind of local analysis is possible. Overstable modes can be revealed by this kind of treatment while the convective and Parker-type instabilities are suppressed since they only appear for large enough wavelengths. An example of such a local analysis is given in Sec. 5.7.1 below.
Let us finally consider the influence of the coefficients $`A`$ and $`C`$ on the growth rates of unstable modes. Tab. 3 below illustrates the dependence of $`\widehat{\omega }`$ on the value of $`v_0`$ which determines the size of $`A`$ and $`C`$ for given values of $`\beta r`$ (cf. Eq. 5.97). We consider two examples: Monotonic instability of equipartition flux tubes ($`\beta \delta =3.6`$) in temperature equilibrium with the environment ($`\beta rm=1`$) and overstability of equipartition tubes with $`\beta rm=1.8`$. <sup> </sup><sup> </sup>* Strictly spoken, $`v_0\mathrm{}`$ is inconsistent with the assumption of a finite value for $`\beta r`$ unless the flux tube diameter becomes infinite, too. We nevertheless consider this case (which gives $`A=C=0`$) for the purpose of comparison with the results for finite values of $`v_0`$.
$`v_0`$ $`\beta rm=1.`$, $`k=0.`$ (monotonic) $`\beta rm=1.8`$, $`k=1.`$ (overstable) Re($`\widehat{\omega }`$) Im($`\widehat{\omega }`$) Re($`\widehat{\omega }`$) Im($`\widehat{\omega }`$) $`\mathrm{}^{}`$ 1.334 0. 0.137 1.345 10. 1.238 0. 0.077 1.348 1. 0.667 0. 0.012 1.413 0.1 0.089 0. 0.001 1.414
Tab. 3: Dependence of growth rate, Re($`\widehat{\omega }`$), and oscillation frequency, Im($`\widehat{\omega }`$), on the the vertical velocity $`v_0`$ which determines the size of the coefficients of the damping terms $`A`$ and $`C`$. Two cases are considered: Monotonic instability of equipartition flux tubes ($`\beta \delta =3.6`$) in temperature equilibrium with the environment ($`\beta rm=1`$) and overstability of equipartition tubes with $`\beta rm=1.8`$. While Im($`\widehat{\omega }`$) is barely affected, the growth rate decreases drastically for small values of $`v_0`$ and fixed $`\beta r`$.
As explained in Sec. 5.4, the velocity $`v_0`$ is measured in units of $`v_A/\sqrt{2}`$ and the time unit is of the order of one day for the deep layers of the solar convection zone. Smaller values of $`v_0`$ (for fixed values of $`\beta r`$, see discussion below) entail larger values of $`A`$ and $`C`$ which lead to smaller growth rates of both monotonic and oscillatory instabilities. The overstable modes which typically have much smaller growth rates than the monotonic modes are more strongly affected by this damping mechanism. The growth time $`2\pi /\text{Re}\left(\widehat{\omega }\right)`$ of the overstable mode in our example increases from about 80 days for $`v_0=10.`$ to 17 years for $`v_0=0.1`$ while for the monotonic mode the numbers are 7 days and 70 days, respectively. On the other hand, the oscillation frequencies are hardly affected.
We may estimate the value of $`v_0`$ associated with a given value of $`\beta r`$ with aid of Eq. (5.6) which may be simplified by assuming $`\beta 1`$ and $`2C_D/\pi 1`$. Written in non-dimensional quantities we obtain
$$\beta r\frac{v_0^2}{a_0}.$$
$`\left(5.105\right)`$
Here $`a_0`$ is the flux tube radius in units of the external pressure scale height. For the interesting range $`\beta r=O\left(1\right)`$ we have
$$v_0\sqrt{a}_0$$
$`\left(5.106\right)`$
We see that for fixed $`\beta r`$ the quantity $`v_0`$ effectively is a measure of the flux tube radius. Tab. 3 shows that small flux tubes are more strongly influenced by the drag forces and suffer from a stronger damping than those with larger radius. For $`\beta r=1`$ and $`H_{p0}`$ = $`510^4`$ km we find the range $`v_00.03\mathrm{}0.3`$ for flux tube radii between 50 km and 5000 km. Consequently, we have to expect a strong effect of the damping terms on the growth rates, especially in the case of overstable modes.
5.7 Symmetric loops with vertical external flow
As an example of a flux tube equilibrium with a curved path we investigate a planar loop which is symmetric with respect to a vertical line through its maximum or minimum. Since we shall consider the stability only in the vicinity of the point of extremum, the results are relevant also for flux tubes which wind like a serpentine line, i.e. consist of a sequence of (symmetric) loops with alternating maxima and minima. Such configurations are of practical interest since one could imagine that the kink instability of a horizontal (or toroidal) flux tube in a convection zone leads to ‘sea serpent’ structures, i.e. a series of erupted active regions connected by loops with minima in the convection zone. Another interesting question is whether ‘dived sea serpents’, a series of minimum and maximum loops fully within the convection zone represent a stable alternative to the unstable horizontal tubes.
As in the preceding section we assume a purely vertical external velocity field which does not depend on $`x`$, the coordinate in the horizontal direction, and continue to use the notation of the velocity terms introduced there. The equilibrium is determined by a balance of buoyancy, curvature and drag force which is described by Eq. (5.4). For the point of extremum with $`(\widehat{𝐠}_0\widehat{𝐧}_0)=1`$ in the case of a minimum and $`(\widehat{𝐠}_0\widehat{𝐧}_0)=+1`$ in the case of a maximum we have (in non-dimensional form)
$$\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)=\frac{2}{R_0}+\beta rm.$$
$`\left(5.107\right)`$
The upper sign on the l.h.s. applies for a maximum, the lower sign for a minimum. In both cases we have $`\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)=0`$, i.e. the tangent vector has a horizontal direction at the point of extremum.
5.7.1 Local analysis
Overstable modes of a symmetric loop can be treated analytically by way of a local stability analysis assuming perturbations with small wavelength (large wavenumber) in the direction along the flux tube. We consider the extremum point (maximum or minimum) of a symmetric loop and investigate its stability with respect to growing oscillations of short wavelength in the neighborhood of this point.
We assume wave-like perturbations, i.e. $`\epsilon ,\eta \mathrm{exp}(ikx)`$ whose wavelength $`\lambda =2\pi /k`$ is small enough such that we can take all quantities describing the equilibrium flux tube ($`R_0,H_{p0},\rho _0,(\widehat{𝐠}_0\widehat{𝐥}_0),(\widehat{𝐠}_0\widehat{𝐧}_0)\mathrm{})`$ to be constant within a wavelength. On the other hand, the approximation of slender flux tubes demands $`\lambda a_0`$ where $`a_0`$ is the radius of the tube. Consequently, we require $`a_0\lambda R_0`$ and also $`\lambda (8R_0H_{p0})^{1/2}`$ which results from the requirement that the height increment of the path of the equilibrium flux tube within one wavelength must be small compared to the scale height. Both conditions lead to
$$\frac{2\pi }{a_0}k\frac{2\pi }{\mathrm{min}(R_0,\mathrm{\hspace{0.17em}8}R_0H_{p0})^{1/2}}.$$
$`\left(5.108\right)`$
Since overstability often requires that the wavenumber exceeds some critical value one has to keep in mind that the following results are only applicable within the limits of the approximation of slender flux tubes if the tube radius satisfies the left part of the above relation. For wavenumbers which satisfy Eq. (5.108) we may use Eq. (5.107) to rewrite the general perturbation equations, Eqs. (5.63/64), taking the extremum as reference point, $`z_m`$, for the non-dimensionalization of all quantities and obtain
$$\begin{array}{ccc}\hfill \omega ^2\eta & =ik\epsilon \left(\frac{4}{R_0}\beta rm\pm \frac{2}{\gamma }\right)+\eta \frac{1}{R_0}\left(\frac{2}{R_0}\beta rm\pm \frac{2}{\gamma }\right)\mathrm{\hspace{0.33em}2}k^2\eta \hfill & \left(5.109\right)\hfill \\ & & \\ \hfill \omega ^2\epsilon & =ik\eta \left(\frac{4}{R_0}+\frac{\beta rm}{2}+\frac{2}{\gamma }\right)\mathrm{\hspace{0.33em}2}k^2\epsilon i\omega \epsilon \frac{2\beta r}{v_0}\hfill & \\ & & \\ & +\epsilon \left[\frac{2}{R_0^{\mathrm{\hspace{0.17em}2}}}\pm \frac{4}{\gamma R_0}+\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)+\beta rm\left(\frac{2}{v_0}\frac{dv_0}{dz}\pm \frac{1}{2\gamma }\frac{1}{2R_0}\right)\right].\hfill & \left(5.110\right)\hfill \end{array}$$
Note that for the lower signs ($`\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)=1`$) and in the limit $`R_0\mathrm{}`$ these equations pass over to Eqs. (5.94/95) for a horizontal flux tube. We follow the same procedure as in the preceding section and determine from Eqs. (5.109/110) a dispersion relation for $`i\omega \widehat{\omega }`$, again of a fourth order polynomial with real coefficients:
$$\widehat{\omega }^4+A\widehat{\omega }^3+B\widehat{\omega }^2+C\widehat{\omega }+D=\mathrm{\hspace{0.33em}0}$$
$`\left(5.111\right)`$
$$\begin{array}{cc}& A=\frac{2\beta r}{v_0}\hfill \\ & \\ & B=\mathrm{\hspace{0.33em}4}k^2+\frac{2}{R_0}\left(\frac{2}{R_0}\frac{3}{\gamma }\right)\beta \delta \frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)+\beta rm\left(\pm \frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }+\frac{3}{2R_0}\right)\hfill \\ & \\ & C=\frac{2\beta r}{v_0}\left[2k^2+\frac{1}{R_0}\left(\frac{2}{R_0}+\beta rm\frac{2}{\gamma }\right)\right]\hfill \\ & \\ & D=\mathrm{\hspace{0.33em}2}k^2\left[2k^2\beta \delta \frac{1}{\gamma }+\frac{2}{R_0}\left(\pm \frac{1}{\gamma }\frac{2}{R_0}\right)+\beta rm\left(\pm \frac{2}{v_0}\frac{dv_0}{dz}\pm \frac{1}{\gamma }\frac{3}{2R_0}\right)\frac{\left(\beta r\right)^2}{4}\right]\hfill \\ & \\ & +\frac{1}{R_0}\{\frac{4}{R_0^{\mathrm{\hspace{0.17em}2}}}(\frac{1}{R_0}\frac{3}{\gamma })+\frac{2}{\gamma R_0}(\frac{6}{\gamma }1)\pm \frac{2}{\gamma ^2}(1\frac{2}{\gamma })+\beta \delta (\frac{2}{R_0}\beta rm\pm \frac{2}{\gamma })\hfill \\ & \\ & +\beta rm[\frac{2}{v_0}\frac{dv_0}{dz}(\frac{2}{\gamma }\frac{2}{R_0})+\frac{1}{\gamma }(\frac{3}{\gamma }\frac{6}{R_0}1)+\frac{3}{R_0^{\mathrm{\hspace{0.17em}2}}}]+\left(\beta r\right)^2(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }+\frac{1}{2R_0})\}\hfill \end{array}$$
In analogy to the treatment in Sec. 5.6 we consider a reduced dispersion relation obtained by setting $`A=C=0`$ in Eq. (5.111) which leads to
$$\widehat{\omega }^2=\frac{B}{2}\pm \left(\frac{B^2}{4}D\right)^{1/2}.$$
$`\left(5.112\right)`$
We have obtained a number of numerical solutions of the full dispersion relation, Eq. (5.111), which confirm the criteria derived below on the basis of Eq. (5.112). Finite values of $`A`$ and $`C`$ only affect the growth or decay rates. Monotonic instability for large values of $`k`$ requires extreme values for $`\beta \delta `$ or $`\beta r`$ which are unrealistic for a convection zone. Oscillatory instability, on the other hand, which sets in if $`B^2/4D<0`$ can be well described by local analysis. This expression is given by
$$\begin{array}{ccc}\hfill \frac{B^2}{4}D& =k^2\left[\frac{4}{\gamma ^2}+\frac{\left(\beta r\right)^2}{2}+3\beta rm\left(\frac{2}{R_0}\frac{1}{\gamma }\right)+\frac{16}{R_0}\left(\frac{1}{R_0}\frac{1}{\gamma }\right)\right]\hfill & \\ & & \\ & +\frac{1}{4}\left[\frac{2}{R_0}\left(\frac{2}{R_0}\frac{3}{\gamma }\right)+\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)\beta rm\left(\pm \frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }+\frac{3}{2R_0}\right)\right]^2\hfill & \\ & & \\ & \frac{1}{R_0}\left\{\frac{4}{R_0^{\mathrm{\hspace{0.17em}2}}}\left(\frac{1}{R_0}\frac{3}{\gamma }\right)+\mathrm{}+\left(\beta r\right)^2\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }+\frac{1}{2R_0}\right)\right\}\hfill & \left(5.113\right)\hfill \end{array}$$
where the term in braces is the same as the last term in the above definition of $`D`$. Since local analysis demands large values of $`k`$ we conclude that the sign of $`B^2/4D`$ is determined by the sign of the term multiplied by $`k^2`$ in Eq. (5.113). Note that the superadiabaticity does not enter into this term, i.e. the excitation condition for oscillatory instability is independent of the stability of the external stratification which may only affect the wavelength range of overstable modes. This applies also to horizontal flux tubes (cf. Eqs. 5.101/103). A simple calculation shows that the term under consideration is negative if
$$\left(\beta r\right)_{\mathrm{min}}<\beta r<\left(\beta r\right)_{\mathrm{max}}$$
$`\left(5.114\right)`$
where
$$\left(\beta r\right)_{\mathrm{min}}=\mathrm{min}[\left(\beta r\right)_1,\left(\beta r\right)_2],\left(\beta r\right)_{\mathrm{max}}=\mathrm{max}[\left(\beta r\right)_1,\left(\beta r\right)_2]$$
$`\left(5.115\right)`$
with
$$\begin{array}{ccc}\hfill \left(\beta r\right)_1& =\left(3m1\right)\left(\pm \frac{1}{\gamma }\frac{2}{R_0}\right)\hfill & \\ & & \\ \hfill \left(\beta r\right)_2& =\left(3m+1\right)\left(\pm \frac{1}{\gamma }\frac{2}{R_0}\right).\hfill & \left(5.116\right)\hfill \end{array}$$
Since $`\beta r`$ is a positive quantity and $`\left(\beta r\right)_1`$ and $`\left(\beta r\right)_2`$ always have the same sign, overstability is only possible if both are positive, too. We now have to distinguish 4 cases, i.e. maximum/minimum and upflow/downflow, respectively. Note that since $`m=\text{sgn}\left(𝐯_{e}^{}{}_{0}{}^{}\widehat{𝐧}_0\right)`$ the value $`m=1`$ signifies an upflow for a minimum and a downflow for a maximum while $`m=1`$ corresponds to a downflow for a minimum and an upflow for a maximum. Tab. 4 summarizes the stability properties for the four possible cases. The fourth column gives necessary conditions on $`R_0`$ for overstability while the last column shows the range of values of $`\beta r`$ which lead to overstability. Cases 3 and 4 with $`R_0\mathrm{}`$ pass over to the case of a horizontal flux tube treated in the preceding section (cf. Eq. 5.102).
| | Case | Extremum | $`m`$ | Condition | Oscillatory instability for |
| --- | --- | --- | --- | --- | --- |
| | 1 | maximum | $`1`$ (up) | $`R_0<2\gamma `$ | $`2\left(2/R_01/\gamma \right)<\beta r<4\left(2/R_01/\gamma \right)`$ |
| | 2 | maximum | $`+1`$ (down) | $`R_0>2\gamma `$ | $`2\left(1/\gamma 2/R_0\right)<\beta r<4\left(1/\gamma 2/R_0\right)`$ |
| | 3 | minimum | $`1`$ (down) | | $`2\left(1/\gamma +2/R_0\right)<\beta r<4\left(1/\gamma +2/R_0\right)`$ |
| | 4 | minimum | $`+1`$ (up) | | no overstability |
Tab. 4: Range of values for $`\beta r`$ which lead to overstability of symmetric loops with fixed radius of curvature, $`R_0`$. The four possible cases (maximum/minimum, upflow/downflow) are given. For the first pair of cases necessary conditions for overstability exist which are indicated in the fourth column.
We must keep in mind, however, that the quantities $`R_0`$ and $`\beta r`$ generally cannot be chosen independently since they are related through the equilibrium condition given by Eq. (5.107). If the relation between internal and external temperature and the value of $`\beta rm`$ are given, the radius of curvature is fixed. Hence, Eq. (5.116) and Tab. 4 are only applicable if the internal temperature and $`\beta r`$ are chosen in such a way that $`R_0`$ stays constant. For more realistic cases we must prescribe the relation between external and internal temperature and determine $`R_0`$ according to Eq. (5.107) for any given value of $`\beta r`$.
In what follows we discuss the case of an isothermal flux tube, i.e. internal equal to external equilibrium temperature. In this case the equilibrium (written in dimensionless form) is given by
$$\frac{2}{R_0}+\beta rm=\pm 1.$$
$`\left(5.117\right)`$
We notice that no equilibrium is possible for a minimum with upflow. This is due to the fact that an isothermal flux tube has an upward directed buoyancy force and the curvature force in case of a minimum has the same direction. Consequently, force balance can only be achieved by a downflow. In other cases there are restrictions on $`\beta r`$: For a minimum with downflow we must have $`\beta r1`$ while a maximum with downflow requires $`\beta r1`$ and $`R_01`$. In the case of a maximum with upflow any value of $`\beta r`$ is permitted. We insert Eq. (5.117) into Eq. (5.113) and determine the range of oscillatory instability by the same procedure which led to Eqs. (5.114-116). We find
$$\left(\beta r\right)_1=\frac{\left(\pm 5m1\right)}{3}\left(1\frac{1}{\gamma }\right),\left(\beta r\right)_2=\frac{\left(\pm 5m+1\right)}{3}\left(1\frac{1}{\gamma }\right).$$
$`\left(5.118\right)`$
Since we certainly have $`\gamma 1`$ the sign of $`\left(\beta r\right)_{1,2}`$ for a given equilibrium is determined only by $`m`$. The stability criteria for the isothermal case as they follow from Eqs. (5.114/118) are summarized in Tab. 5:
| | Case | Extremum | $`m`$ | Oscillatory instability for |
| --- | --- | --- | --- | --- |
| | 1 | maximum | $`1`$ (up) | no overstability |
| | 2 | maximum | $`+1`$ (down) | $`\left(4/3\right)\left(11/\gamma \right)<\beta r<2\left(11/\gamma \right)`$ |
| | 3 | minimum | $`1`$ (down) | $`\left(4/3\right)\left(11/\gamma \right)<\beta r<2\left(11/\gamma \right)`$ |
| | 4 | minimum | $`+1`$ (up) | no equilibrium |
Tab. 5: Range of values for $`\beta r`$ which lead to overstability of isothermal, symmetric loops. Only configurations with an external downflow are liable to overstability.
For $`\gamma =5/3`$ the range of values of $`\beta r`$ which lead to overstability in cases 2 and 3 is given by
$$0.53<\beta r<\mathrm{\hspace{0.33em}0.8}$$
$`\left(5.119\right)`$
Since the equilibrium condition given by Eq. (5.117) requires $`\beta r1`$ for case 3 (minimum with downflow) there is no possibility for overstability in this case if the flux tube is isothermal. Consequently, the minimum region of a symmetric loop formed by a flux tube which is in temperature equilibrium with the external gas is always stable with respect to growing oscillations and a maximum with a downflow is the only configuration which can lead to overstability of an isothermal loop. For $`R_0\mathrm{}`$ in all cases we have $`\beta r1`$ and we thus recover the result of the preceding section, namely that an isothermal, horizontal flux tube does not show overstability.
Another interesting special case is the neutrally buoyant flux tube, $`\rho _0=\rho _{e0}`$, whose equilibrium is determined by a balance between curvature force and drag force. Eq. (5.107) then leads to (non-dimensional form)
$$\frac{2}{R_0}=\beta rm$$
$`\left(5.120\right)`$
and it is clear that this kind of equilibrium requires $`m=1`$, i.e. either a minimum with a downflow or a maximum with an upflow. Inserting Eq. (5.120) into Eq. (5.113) with $`m=1`$ we now obtain
$$\begin{array}{ccc}\hfill \left(\beta r\right)_1& =\frac{\pm 51}{3\gamma }\hfill & \\ & & \\ \hfill \left(\beta r\right)_2& =\frac{\pm 5+1}{3\gamma }.\hfill & \end{array}$$
Consequently, the minimum region of a neutrally buoyant flux loop is always stable with respect to growing oscillations while a maximum may show overstability. For a maximum loop and $`\gamma =5/3`$ we find $`\left(\beta r\right)_1=0.8`$ and $`\left(\beta r\right)_2=1.2`$. Consequently, overstability of a neutrally buoyant loop can occur if it has a maximum with a radius of curvature of about 2 scale heights.
5.7.2 Constant vertical displacement
If we assume that the perturbation wavelength along the equilibrium flux tube is infinite and consider purely vertical displacements (in $`z`$-direction) the perturbation equations can be simplified considerably and are analytically solvable in the case of symmetric loops. In this way we can investigate monotonic instability. As in the preceding section we consider an extremum point (local minimum or maximum) of a flux tube whose path has the form of a symmetric loop. We assume that the whole structure is displaced vertically by a constant amount, $`z_1`$. Since the tube is not stretched by this operation we have $`l=l_0`$ and conclude from Eq. (5.8) that
$$\eta ^{}=\frac{\epsilon }{R_0}.$$
$`\left(5.121\right)`$
At the point of extremum we have for reasons of symmetry
$$\eta =\eta ^{\prime \prime }=\epsilon ^{}=\mathrm{\hspace{0.33em}0}$$
$`\left(5.122\right)`$
and since the radius of curvature is unaffected by this kind of displacement $`\left(R=R_0\right)`$ we find from Eq. (5.11)
$$\epsilon ^{\prime \prime }=\frac{\epsilon }{R_0^{\mathrm{\hspace{0.17em}2}}}.$$
$`\left(5.123\right)`$
Furthermore, we have $`\left(\widehat{𝐠}_0\widehat{𝐥}_0\right)=0`$ and $`\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)=\pm 1`$ at an extremum point. Using this together with Eqs. (5.121-123) we find that both sides of Eq. (5.63) vanish and Eq. (5.64) has non-trivial solutions provided that
$$\omega ^2i\omega \frac{2\beta r}{v_0}+\left[\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)\pm \frac{2}{\gamma R_0}+\beta rj\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)\right]=\mathrm{\hspace{0.33em}0}$$
$`\left(5.124\right)`$
where (as before) the upper ($`+`$) sign applies for a maximum and the lower ($``$) sign for a minimum. In this equation we have introduced $`j\text{sgn}\left(𝐯_{e0}\widehat{𝐳}\right)=m`$ such that for both cases (minimum and maximum) $`j=+1`$ indicates an upflow and $`j=1`$ a downflow. The quantities have been non-dimensionalized with respect to their values at the point of extremum. As in the preceding sections we have assumed a constant gravitational acceleration ($`s=0`$) in order to simplify the discussion. The influence of a variation of g on the stability properties is marginal, however. Note that in the case of a horizontal tube $`\left(R_0\mathrm{}\right)`$ Eq. (5.124) transforms into Eq. (5.84), the dispersion relation for vertical displacements of a horizontal tube. <sup>*</sup><sup>*</sup>In the derivation of Eq. (5.84) we have assumed $`\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)=1`$ such that we have $`j=m`$ in this case. The same line of arguments as in that case shows that a necessary and sufficient criterion for instability is that the expression in square brackets in Eq. (5.124) is positive, viz.
$$\beta \delta +\frac{1}{\gamma }\left(1\frac{2}{\gamma }\right)\pm \frac{2}{\gamma R_0}+\beta rj\left(\frac{2}{v_0}\frac{dv_0}{dz}\frac{1}{2\gamma }\right)>\mathrm{\hspace{0.33em}0}.$$
$`\left(5.125\right)`$
The effect of the terms which depend on the external velocity is identical to the case of a horizontal tube, i.e. the curvature of the tube does not change the effect of the flow on the stability of the point of extremum. This is not obvious from the original equations (5.63/64) since we find terms there which depend on both velocity and curvature. In the special case of purely vertical displacements these terms cancel, however.
For the heuristic argumentation presented in the following subsection it is important to note here that the stability criterion given by Eq. (5.125) can be also obtained by considering the perturbations of the buoyancy, curvature and drag forces brought about by a vertical displacement of the point of extremum. The sign of the resulting perturbation of the total force then determines the stability properties of the equilibrium. In the present case we assume that a purely vertical displacement does not lead to a flow of matter along the tube $`(\eta 0)`$ such that the ratio $`B/\rho `$ is constant and the perturbations (indicated by an index 1) are related to the equilibrium quantities as $`B_1/B_0=\rho _1/\rho _0`$. Together with the adiabaticity of the perturbations (Eq. 5.38) and the condition of pressure equilibrium (Eq. 5.39) this leads to the following relation for the magnetic field perturbation as function of the displacement $`z_1`$:
$$\frac{B_1}{B_0}=\left(\frac{\beta +1}{\beta \gamma +2}\right)\frac{z_1}{H_{pe}}.$$
$`\left(5.126\right)`$
Consequently, for $`\beta 1`$ (ignoring terms of order $`\beta ^1`$) we find for the perturbation of the curvature force (with $`R_1=0`$):
$$F_{C1}=\frac{B_0B_1}{2\pi R_0}=\pm \frac{B_0^{\mathrm{\hspace{0.17em}2}}z_1}{2\pi \gamma R_0H_{p0}}.$$
$`\left(5.127\right)`$
In a similar way we determine the perturbations of the buoyancy and drag forces as functions of $`z_1`$ and the equilibrium quantities. We add all force perturbations together, take the limit $`\beta 1`$, and non-dimensionalize with respect to the values of the quantities at the point of extremum. Taking then $`z_1\mathrm{exp}\left(i\omega t\right)`$ we obtain exactly the stability criterion given by Eqs. (5.124/125). This result lends some support to the treatment presented in the subsequent section where a similar heuristic approach is used for displacements with large but finite wavelength.
The perturbation of the curvature force given by Eq. (5.127) is helpful for understanding the effect of curvature on the stability of the loop which is expressed by the term $`\pm 2/\gamma R_0`$ in Eq. (5.124). Since the magnetic field decreases for an upward displacement and increases for a downward displacement (cf. Eq. 5.126), the absolute value of the curvature force always decreases for an upward displacement and increases for a downward displacement. Now consider a minimum: The curvature force is directed upward (positive) and is reduced by an upward displacement and increased by a downward displacement; the result is a restoring force which acts against the displacement – this stabilizing effect of curvature shows up in the negative sign of the term $`2/\gamma R_0`$ in the case of a minimum. An analogous consideration shows that in the case of a maximum the effect of curvature is destabilizing. Tab. 6 below summarizes the effect of the different terms in Eq. (5.125) on the stability properties of the loop. Regarding the stratification of the convection zone we assume that the flow speed decreases with depth.
| | Extremum | Flow | $`\pm 2/\left(\gamma R_0\right)`$ | $`\left(2/v_0\right)\left(dv_0/dz\right)>0`$ | $`1/\left(2\gamma \right)<0`$ |
| --- | --- | --- | --- | --- | --- |
| | minimum | up | stabilizing | destabilizing | stabilizing |
| | maximum | up | destabilizing | destabilizing | stabilizing |
| | minimum | down | stabilizing | stabilizing | destabilizing |
| | maximum | down | destabilizing | stabilizing | destabilizing |
Tab. 6: Influence of individual terms in the stability criterion (Eq. 5.125) for various configurations of a symmetric loop in the case of purely vertical displacements. The term given in the third column describes the direct effect of curvature while the terms in the fourth and fifth column represent the effect of the external velocity.
According to Eq. (5.91) we expect that the velocity gradient term dominates over the last term in Eq. (5.125) such that in particular minimum loops with a downflow are possible examples of stable configurations.
In order to estimate the quantitative effect of the curvature term we consider the case of a flux tube in thermal equilibrium with the environment, i.e. $`T_0=T_{e0}`$ and $`\beta (\rho _{e0}/\rho _01)=1`$, such that the equilibrium condition, Eq. (5.107), can be written in non-dimensional form as
$$1+\beta rj=\pm \frac{2}{R_0}.$$
$`\left(5.128\right)`$
We see that for a minimum (lower sign) we must have a downflow $`\left(j=1\right)`$ of sufficient strength since in temperature equilibrium the internal density is always smaller than the external density leading to an upward directed buoyancy force. Inserting Eq. (5.128) into the instability criterion Eq. (5.125) and assuming a constant flow velocity ($`dv_0/dz=0`$) we find instability for
$$\beta \delta +\frac{2}{\gamma }\left(1\frac{1}{\gamma }\right)+\frac{\beta rj}{2\gamma }>\mathrm{\hspace{0.33em}0}.$$
$`\left(5.129\right)`$
The fact that in this case a downflow always stabilizes while an upflow destabilizes can be understood by considering Eq. (5.128): With the exception of the small range $`1\beta rj0`$ a downflow leads to a minimum and an upflow to a a maximum. In both cases the effect of the flow is to decrease $`R_0`$, the radius of curvature. Thus the stabilizing effects of curvature for a minimum (downflow) and its destabilizing effects for a maximum (upflow) are amplified by the flow.
In the absence of a flow we only can have a maximum and the criterion for instability reads
$$\beta \delta >\frac{2}{\gamma }\left(1\frac{1}{\gamma }\right)=0.48$$
$`\left(5.130\right)`$
for $`\gamma =5/3`$. If we compare this with the criterion for a horizontal tube without flow, i.e. $`\beta \delta >+0.12`$ (cf. Eq. 5.89) we see that the effect of curvature is quite significant: While for the horizontal tube a positive superadiabaticity was necessary in the case of purely vertical perturbations, the static, isothermal loop is unstable even in a moderately subadiabatic environment. On the other hand, an external downflow of sufficient amplitude leads to stable local minimum configurations for any value of the superadiabaticity. Take for example the value $`\beta \delta =3.6`$ for equipartition flux tubes in the deep solar convection zone. With $`j=1`$ we find from Eq. (5.129) that such an equilibrium is stable if $`\beta r>13.6`$. From Eq. (5.105) we see that this condition is realized for flux tube radii smaller than about $`H_{p0}/13.64000`$ km if the external velocity is of the order of the typical convective velocities obtained from mixing length theory.
5.7.3 Heuristic approach for perturbations with large wavelength
As far as monotonic instability is concerned, it has already been shown by Spruit and van Ballegooijen (1982) that for horizontal flux tubes without external flow the most unstable perturbations are those which lead to wave-like displacements of very large, but finite, wavelength. In this case a Parker-type/convective instability with a flow along the tube can be excited which is largely undisturbed by curvature forces. We have generalized this result to horizontal tubes with external flow and we suspect that this kind of perturbation is also decisive for the monotonic instability of symmetric loops for which we until now have only considered very small and infinite wavelength (subsections 5.7.1 and 5.7.2, respectively). Unfortunately, finite wavelength perturbations for non-horizontal flux tube equilibria in most cases do not permit analytical treatment since the coefficients in the perturbation equations, Eqs. (5.63/64), become variable. One then has to resort to a numerical solution of the eigenvalue problem for any given equilibrium tube. However, for symmetric loops and displacements with very large, but finite, wavelength which turned out to be most unstable kind of perturbation for horizontal tubes we can extend the heuristic approach sketched in the preceding subsection. Since in contrast to the treatment there we cannot check against the exact result, no definite proof can be given that the approach described below is correct. However, we will show that in the limit of horizontal tubes ($`R_0\mathrm{}`$) the result becomes identical to the exact criterion for monotonic instability and $`k0`$ (but finite, cf. Eqs. 5.98/99). Together with the success of the method for purely vertical displacements demonstrated in the preceding subsection this gives some confidence in its validity.
We consider a symmetric loop in a vertical flow and determine the perturbations of the various forces (buoyancy, curvature, drag) at the point of extremum which are brought about by a vertical displacement. In contrast to the previous treatment we now assume a perturbation with finite wavelength such that the flux tube is not displaced as a whole – parts of the equilibrium tube are lifted while other parts are displaced downward. Since in our cartesian model the wavelength of the perturbation can be made arbitrarily large, it can in particular be chosen large enough such that the perturbations of the tube geometry (arc length, radius of curvature etc.) can be neglected in comparison with the relative perturbations of the other quantities. For example, for a displacement $`z_1`$ with wavenumber $`k`$, the perturbation of the radius of curvature at the extremum point of a symmetric loop is given by
$$\left|R_1\right|=k^2R_0^{\mathrm{\hspace{0.17em}2}}\left|z_1\right|$$
$`\left(5.131\right)`$
where we have used Eqs. (5.11/14) and the symmetry properties. We see that for any given displacement we can make $`R_1`$ as small as we want by decreasing $`k`$ sufficiently, i.e. by increasing the wavelength of the perturbation.
The important difference to the treatment in the preceding subsection lies in the determination of the perturbations of internal density and pressure. Since the wavelength of the displacement is now finite (albeit very large) parts of the equilibrium tube are lifted while other parts are displaced downward such that an internal flow along the tube sets in which tends to establish hydrostatic equilibrium along the magnetic field lines according to the principle of communicating tubes. We therefore determine the perturbed internal density and pressure by assuming that the flow along the tube has already restored hydrostatic equilibrium at the point of extremum, viz.
$$p_1=\frac{p_0}{H_{p0}}z_1$$
$`\left(5.132\right)`$
which entails for adiabatic perturbations (cf. Eq. 5.38)
$$\rho _1=\frac{\rho _0}{\gamma H_{p0}}z_1.$$
$`\left(5.133\right)`$
By assuming that hydrostatic equilibrium is reestablished we probably loose information about growth rates and we also cannot reproduce the overstable modes but we conjecture that the stability criteria for monotonic instabilities are correctly described by this approach. We shall prove this below for the special case of horizontal tubes for which the exact solution is available.
We continue by determining $`F_{B1}`$, the perturbation of the buoyancy force, which is given by
$$F_{B1}=\left(\rho _{e1}\rho _1\right)g_0$$
$`\left(5.134\right)`$
where we have assumed that the gravitational acceleration is constant ($`s=0`$). With $`_{ad}=\left(\gamma 1\right)/\gamma `$, $`\delta =_{ad}`$ and using Eqs. (5.25) and (5.133) we find, after some algebra
$$F_{B1}=\frac{B_0^{\mathrm{\hspace{0.17em}2}}}{8\pi H_{p0}^{\mathrm{\hspace{0.17em}2}}}\left\{\left(1\right)\left[\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)+\frac{\rho _{e0}}{\rho _0}\frac{H_{p0}}{H_{pe}}\beta \left(\frac{H_{pe}}{H_{p0}}1\right)\right]+\beta \delta \right\}z_1.$$
$`\left(5.135\right)`$
We now perform the same procedure as in Sec. 5.4 and take the limit $`\beta 1`$ such that all terms of order $`\beta ^1`$ are neglected unless they are multiplied by $`\beta `$. In this limit we have
$$\beta \left(\frac{H_{pe}}{H_{p0}}1\right)=1+\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)+O\left(\beta ^1\right)$$
$`\left(5.136\right)`$
as well as $`1=1/\gamma +O\left(\beta ^1\right)`$, $`\rho _{e0}/\rho _0=1+O\left(\beta ^1\right)`$, $`H_{p0}/H_{pe}=1+O\left(\beta ^1\right)`$, and Eq. (5.135) is transformed into
$$F_{B1}=\frac{B_0^{\mathrm{\hspace{0.17em}2}}}{8\pi H_{p0}^{\mathrm{\hspace{0.17em}2}}}\left[\beta \delta +\frac{1}{\gamma }+\frac{2}{\gamma }\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)\right]z_1.$$
$`\left(5.137\right)`$
The perturbation of the magnetic field is determined by the condition of pressure balance, Eq. (5.39), which yields together with Eqs. (5.25) and (5.132/233)
$$\frac{B_1}{B_0}=\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)\frac{z_1}{2H_{p0}}.$$
$`\left(5.138\right)`$
Inserting into the first (general) part of Eq. (5.127) we obtain the perturbation of the curvature force
$$F_{C1}=\frac{B_0^{\mathrm{\hspace{0.17em}2}}}{4\pi H_{p0}R_0}\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)z_1.$$
$`\left(5.139\right)`$
Note that we have assumed a sufficiently large wavelength of the displacement such that the perturbation of the radius of curvature can be neglected. Finally we determine $`F_{D1}`$, the perturbation of the drag force. As in the preceding sections we assume a purely vertical external velocity field and continue to use the notation introduced in Eqs. (5.83) and (5.124). Since our present approach cannot adequately describe impulsive motion of the flux tube we omit the contribution to $`F_{D1}`$ due to the motion of the tube itself which gives rise to the last term in Eq. (5.64). We have discussed at some length in Sec. 5.6 that this term does only affect the growth rates but not the stability criteria. A derivation along the lines of Eqs. (5.27-5.33) applied to the special case discussed here yields
$$F_{D1}=\frac{C_D\rho _{e0}v_0^{\mathrm{\hspace{0.17em}2}}j}{\pi a_0}\left(\frac{v_1}{v_0^{\mathrm{\hspace{0.17em}2}}}+\frac{\rho _{e1}}{\rho _{e0}}+\frac{B_1}{2B_0}\right).$$
$`\left(5.140\right)`$
Using Eqs. (5.25), (5.90) and (5.138) and taking the limit $`\beta 1`$ we obtain, after some algebra
$$F_{D1}=\frac{B_0^{\mathrm{\hspace{0.17em}2}}}{8\pi H_{p0}^{\mathrm{\hspace{0.17em}2}}}\beta rj\left[\frac{2H_{p0}}{H_{v0}}\frac{1}{\gamma }+\frac{1}{4}\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)\right]z_1$$
$`\left(5.141\right)`$
where $`H_{v0}v_0\left(dv_0/dz\right)^1`$ denotes the scale height of the external velocity. The total force perturbation is obtained by adding Eqs. (5.137), (5.139), (5.141), and using the equilibrium condition, Eq. (5.4), which gives
$$\begin{array}{ccc}\hfill F_{B1}+F_{C1}+F_{D1}=& \frac{B_0^{\mathrm{\hspace{0.17em}2}}}{8\pi H_{p0}^{\mathrm{\hspace{0.17em}2}}}[\beta \delta +\frac{1}{\gamma }+\frac{4H_{p0}}{R_0}(\frac{H_{p0}}{R_0}\frac{1}{\gamma })+\frac{\left(\beta r\right)^2}{4}\hfill & \\ & & \\ & +\beta rj(\frac{2H_{p0}}{H_{v0}}+\frac{1}{\gamma }\pm \frac{5}{2}\frac{H_{p0}}{R_0})]z_1.\hfill & \left(5.142\right)\hfill \end{array}$$
The equilibrium is unstable if the expression within square brackets is positive which means that the perturbation of total force has the same sign as the displacement and thus tends to increase the latter. The exact result for a horizontal tube obtained in Sec. 5.6 is recovered by taking $`R_0\mathrm{}`$ which yields as condition for monotonic instability
$$\beta \delta +\frac{1}{\gamma }+\frac{\left(\beta r\right)^2}{4}+\beta rj\left(\frac{2H_{p0}}{H_{v0}}+\frac{1}{\gamma }\right)>\mathrm{\hspace{0.33em}0}.$$
$`\left(5.143\right)`$
Regarding Eq. (5.90) we find that this is identical to Eq. (5.99), <sup>*</sup><sup>*</sup>Since we have taken $`\left(\widehat{𝐠}_0\widehat{𝐧}_0\right)=1`$ in Sec. 5.6 we have $`j=m`$. the exact criterion in the limit of very large wavelength. For $`\beta r=0`$ the result of Spruit and van Ballegooijen (1982) is recovered. For finite radius of curvature the general criterion for monotonic instability is
$$\beta \delta +\frac{1}{\gamma }+\frac{4H_{p0}}{R_0}\left(\frac{H_{p0}}{R_0}\frac{1}{\gamma }\right)+\frac{\left(\beta r\right)^2}{4}+\beta rj\left(\frac{2H_{p0}}{H_{v0}}+\frac{1}{\gamma }\pm \frac{5}{2}\frac{H_{p0}}{R_0}\right)>\mathrm{\hspace{0.33em}0}.$$
$`\left(5.144\right)`$
This could easily be transformed to a non-dimensional form by formally taking $`H_{p0}1`$ but we do not change the notation here. Let us now discuss the influence of a finite curvature of the equilibrium flux tube on its stability which is expressed in the terms containing the ratio $`H_{p0}/R_0`$. The third term on the l.h.s. of Eq. (5.144) is always positive and therefore destabilizing for a minimum (lower sign) while for a maximum (upper sign) it is stabilizing if $`H_{p0}/R_0<1/\gamma `$. However, since $`\gamma 1`$ it is easy to show that the sum of the second and the third term on the l.h.s. of Eq. (5.144) is always positive, i.e.
$$\frac{1}{\gamma }+\frac{4H_{p0}}{R_0}\left(\frac{H_{p0}}{R_0}\frac{1}{\gamma }\right)>\mathrm{\hspace{0.33em}0}$$
$`\left(5.145\right)`$
such that without external velocity curvature cannot stabilize a flux tube in a superadiabatic environment ($`\beta \delta >0`$). For a minimum the effect of curvature always is destabilizing, in contrast to the case of a constant vertical displacement (cf. Eq. 5.125). This is caused by the perturbation of the curvature force. If we take $`\beta r=0`$ the equilibrium condition (Eq. 5.4) for a minimum reads
$$\beta \left(1\frac{\rho _{e0}}{\rho _0}\right)=\frac{2H_{p0}}{R_0}$$
$`\left(5.146\right)`$
and we find from Eq. (5.139)
$$F_{C1}=+\frac{B_0^{\mathrm{\hspace{0.17em}2}}z_1}{2\pi R_0^{\mathrm{\hspace{0.17em}2}}}$$
$`\left(5.147\right)`$
such that the perturbation of the curvature always tends to increase the displacement.
We summarize the influence of the velocity terms in Eq. (5.144) on the stability properties of a loop in Tab. 7 below, assuming that the flow speed decreases with depth ($`H_{v0}>0`$).
Extremum Flow $`2H_{p0}/H_{v0}+1/\gamma >0`$ $`\pm \left(5/2\right)\left(H_{p0}/H_{v0}\right)`$ minimum up destabilizing stabilizing maximum up destabilizing destabilizing minimum down stabilizing destabilizing maximum down stabilizing stabilizing
Tab. 7: Influence of velocity-related terms in the criterion for monotonic instability (Eq. 5.144) for various configurations of a symmetric loop in the case of displacements with large wavelength. The fourth term in the criterion, $`(\beta r)^2/4`$, is always positive and destabilizing.
The sum of the second, third and fourth term in Eq. (5.144) is always positive (destabilizing) such that in a superadiabatic environment a maximum loop with an upflow cannot be stabilized by a flow with $`H_{v0}>0`$. The other three configurations can be stabilized provided that certain conditions are fulfilled.
As example let us consider the isothermal case, $`T_0=T_{e0}`$, for which the equilibrium is determined by Eq. (5.128). If we insert this condition into the criterion given by Eq. (5.144) we find
$$\beta \delta +\mathrm{\hspace{0.17em}1}\frac{1}{\gamma }+\frac{5}{2}\left(\beta r\right)^2+\beta rj\left(\frac{2H_{p0}}{H_{v0}}\frac{1}{\gamma }+\frac{13}{4}\right)>\mathrm{\hspace{0.33em}0}.$$
$`\left(5.148\right)`$
In the static case ($`\beta r=0`$) which leads to a loop with a maximum (cf. Eq. 5.128) we have
$$\beta \delta >\frac{1}{\gamma }\mathrm{\hspace{0.33em}1}=0.4$$
$`\left(5.149\right)`$
for $`\gamma =5/3`$ (this value will also be used in the following numerical examples). Consequently, a static isothermal loop is unstable in both superadiabatic and slightly subadiabatic regions. If we include the velocity terms and assume a constant velocity ($`H_{v0}\mathrm{}`$) we find, similar to the case of displacements with infinite wavelength (Eq. 5.129), that an upflow ($`j=+1`$) always destabilizes while a downflow ($`j=1`$) may exert a stabilizing influence provided that
$$\frac{5}{2}\left(\beta r\right)^2\beta r\left(\frac{13}{4}\frac{1}{\gamma }\right)<\mathrm{\hspace{0.33em}0}$$
$`\left(5.150\right)`$
which means
$$0<\beta r<\mathrm{\hspace{0.33em}1.06}$$
$`\left(5.151\right)`$
For $`\beta r>1`$ the loop form changes from a maximum to a minimum. The smallest value that the expression on the l.h.s. of Eq. (5.150) can reach is $`0.7`$ (for $`\beta r=0.53`$) which leads to the criterion $`\beta \delta >0.4+0.7=0.3`$. Although some stabilization has been achieved, a constant downflow cannot stabilize an isothermal equipartition loop in the deep convection zone of the Sun where we have $`\beta \delta =3.6`$.
In which way does a velocity gradient affect the stability of an isothermal loop? From Eq. (5.148) we see that a velocity which increases with height ($`H_{v0}>0`$) has a stabilizing effect in the case of a downflow and a destabilizing effect for an upflow (and vice versa). In the case of a flow with constant mass flux density we have seen in Sec. 5.6 that $`H_{v0}H_{p0}`$ such that we find the criterion
$$\beta \delta +\mathrm{\hspace{0.17em}0.4}+\frac{5}{2}\left(\beta r\right)^2+\mathrm{\hspace{0.17em}4.65}\left(\beta rj\right)>\mathrm{\hspace{0.33em}0}.$$
$`\left(5.152\right)`$
Consequently, a downflow has a stabilizing effect for $`0<\beta r<1.86`$. The velocity terms attain their most negative value for $`\beta r=0.93`$ which gives as condition for instability
$$\beta \delta >\mathrm{\hspace{0.33em}1.76}$$
$`\left(5.153\right)`$
Thus an equipartition tube with $`\beta \delta =3.6`$ is still unstable. In principle we may expect larger velocity gradients for convective downflows near the bottom of the convection zone or within an overshoot region where the flows are strongly decelerated due to the strong subadiabaticity of the radiative region below. Take for example a large equipartition flux tube ($`a_0=10^4`$ km) near the bottom of the solar convection zone ($`H_{p0}=510^4`$ km, $`v_0=v_{A0}=100`$ m/s which gives $`\beta r=2.5`$, $`\beta \delta =3.6`$). We find from Eq. (5.148) that a minimum loop formed by such a tube is stable in a downflow provided that $`H_{v0}/H_{p0}<0.4`$ or $`H_{v0}<210^4`$ km, a value which does not appear unrealistic. On the other hand, a smaller tube with $`a_0=10^3`$ km and $`\beta r=25`$ already requires a value of $`H_{v0}<160`$ km in order to be stabilized which is clearly unrealistic. We conclude that minimum loops formed by relatively large, isothermal flux tubes can possibly be stabilized by a strongly decelerating downflow near the bottom of the solar convection zone.
The considerations above were for the isothermal case, i.e. a flux tube which is in temperature equilibrium with its environment. Although there is a natural tendency towards this state due to radiative energy exchange the relevant time scale becomes very large in the deep layers of a convection zone. If a loop has evolved out of an initially horizontal tube and hydrostatic equilibrium along the field lines has been established adiabatically, a temperature difference with respect to a superadiabatic environment is the consequence. A loop with a minimum would be somewhat cooler and a loop with a local maximum somewhat hotter than its surroundings. In order to assess the influence of such a temperature difference on the stability of a loop we define the parameter
$$\alpha \beta \left(\frac{\rho _{e0}}{\rho _0}1\right)$$
$`\left(5.154\right)`$
which expresses the relation between external and internal temperature: If we assume that the mean molecular weight is the same inside and outside the tube we find from Eq. (5.136):
$$\beta \left(\frac{\rho _{e0}}{\rho _0}1\right)=\mathrm{\hspace{0.33em}1}\beta \left(\frac{T_{e0}}{T_0}1\right)$$
$`\left(5.155\right)`$
where $`T_{e0}`$ and $`T_0`$ denote external and internal temperature, respectively. Consequently, temperature equilibrium entails $`\alpha =1`$ while for a cooler interior we find $`\alpha <1`$ and for a hotter interior we have $`\alpha >1`$. If the temperature difference $`\mathrm{\Delta }TT_0T_{e0}`$ is small compared to the external temperature we can approximate Eq. (5.155) as
$$\alpha \mathrm{\hspace{0.33em}1}+\beta \frac{\mathrm{\Delta }T}{T_{e0}}$$
$`\left(5.156\right)`$
and for equipartition flux tubes near the bottom of the solar convection zone ($`T_{e0}=210^6`$ K, $`\beta =10^6`$) we have
$$\alpha \mathrm{\hspace{0.33em}1}+\frac{\mathrm{\Delta }T}{2}$$
$`\left(5.157\right)`$
where $`\mathrm{\Delta }T`$ is assumed to be given in degrees Kelvin. Using Eq. (5.154) the equilibrium condition (Eq. 5.4) at the point of extremum is written as
$$\alpha +\beta rj=\pm \frac{2H_{p0}}{R_0}.$$
$`\left(5.158\right)`$
If we insert Eq. (5.158) into Eq. (5.144) we obtain the following condition for instability:
$$\beta \delta +\alpha ^2+\frac{1}{\gamma }\left(12\alpha \right)+\frac{5}{2}\left(\beta r\right)^2+\beta rj\left(\frac{2H_{p0}}{H_{v0}}\frac{1}{\gamma }+\frac{13}{4}\alpha \right)>\mathrm{\hspace{0.33em}0}.$$
$`\left(5.159\right)`$
For given values of $`\beta \delta `$ and $`\beta rj`$ we can determine a range of values of $`\alpha `$ for which the l.h.s. of Eq. (5.159) is negative describing a stable loop. For the case $`H_{v0}=H_{p0}`$ and $`\gamma =5/3`$ this is achieved by solving the quadratic inequality
$$\alpha ^2+\alpha P+Q<\mathrm{\hspace{0.33em}0}$$
$`\left(5.160\right)`$
where
$$\begin{array}{cc}& P=\frac{13}{4}\left(\beta rj\right)\frac{6}{5}\hfill \\ & Q=\beta \delta +\frac{5}{2}\left(\beta r\right)^2+\frac{7}{5}\left(\beta rj\right)+\frac{3}{5}\hfill \end{array}$$
This inequality has a range of real solutions provided that $`WP^2/4Q0`$ which leads to the condition
$$\frac{9}{64}\left(\beta r\right)^2\frac{67}{20}\left(\beta rj\right)\frac{6}{25}\beta \delta \mathrm{\hspace{0.33em}0}.$$
$`\left(5.161\right)`$
For an equipartition flux tube in an external downflow ($`j=1`$) within the deep layers of the convection zone ($`\beta \delta =3.6`$) we find that Eq. (5.161) is satisfied for $`\beta r>1.1`$. For example, if we take $`\beta rj=2`$, Eq. (5.160) is fulfilled for $`2.<\alpha <5.7`$, i.e. if the tube is slightly hotter than its environment. The equilibrium condition (Eq. 5.158) shows that this case refers to a loop with a maximum. Minima can be stabilized by an upflow ($`j=+1`$) in which case Eq. (5.161) is satisfied for $`\beta r>24.9`$. For example, if we take $`\beta rj=25.`$ we find that Eq. (5.160) is fulfilled for $`40.57<\alpha <39.48`$, i.e. a cool loop which forms a minimum. Generally we can conclude from Eqs. (5.160/161) and (5.158) that in the case $`H_{v0}=H_{p0}`$ and $`\beta \delta =3.6`$ only cool minima in an upflow and hot maxima in a downflow represent a stable configuration provided that the temperature difference and the flow velocity correspond to the relationship expressed by the above inequalities.
From the bottom to the middle parts of the solar convection zone the integrated temperature difference between adiabatic and actual stratification amounts to only $`\mathrm{\Delta }T1`$K (Parker, 1987) such that Eq. (5.157) gives rather small values for the parameter $`\alpha `$. Consequently, the estimates discussed above suggest that unless more efficient cooling or heating takes place (e.g. if a loop sinks down from the top layers of the convection zone) thermal effects in conjunction with drag forces cannot effectively stabilize flux tubes with or without loops in a stellar convection zone.
5.8 Summary of the stability properties
The formalism derived in Secs. 5.3/4 provides a tool which can be used to investigate the stability of a wide class of flux tube equilibrium structures. In most cases, however, a numerical treatment of the resulting eigenvalue problem is necessary. Such an undertaking is intended for the future but outside the scope of the work presented here. However, we have been able to determine analytically the stability properties in a number of cases which are not without general relevance. We have considered in particular
\- horizontal flux tubes with purely vertical and wave-like displacements (Sec. 5.6),
\- symmetric loops with perturbations of small wavelength (Sec. 5.7.1),
\- symmetric loops with purely vertical displacements (infinite wavelength) (Sec. 5.7.2), and
\- symmetric loops with displacements of large wavelength (Sec. 5.7.3).
For horizontal flux tubes we have generalized the results of Spruit and van Ballegooijen (1982) to the case of a vertical external flow. The monotonic instability found by these authors can be stabilized if the flow speed has a large gradient, for example by a downflow whose velocity strongly decreases with depth. If perturbations with finite wavelength along the flux tube are considered the value of the quantity $`\beta r`$ has to be within a specific range for this stabilization to become effective. Irrespective of the direction or gradient of the flow, sufficiently large values of $`\beta r`$ (due to small field strength, small radius, or large velocity) provoke monotonic instability caused by the radius change of the tube during its displacement. An estimate based on the properties of the deep layers of the solar convection zone shows that the effect of a downflow with constant mass flux density is insufficient to prevent isothermal equipartition flux tubes from monotonic instability due to the superadiabatic stratification.
The monotonic mode has been investigated also for the case of a symmetric loop structure with a horizontal tangent vector at the point of extremum (maximum or minimum). We have derived the exact solution for constant vertical displacements (infinite longitudinal wavelength) and used it as a guideline and test for a heuristic approach which allowed to treat also the case of large (but finite) wavelength. This represents the most unstable perturbation for horizontal tubes and it turned out that this is true also for symmetric loop structures. In the absence of an external flow, all loops (maximum or minimum) are monotonically unstable in a superadiabatic or slightly subadiabatic environment. A flow may exert a stabilizing influence: For an isothermal tube, a downflow with values of $`\beta r`$ within a certain range stabilizes. In particular, a strongly decelerating downflow leads to a stable minimum loop even if $`\beta \delta =3.6`$. If the tube is non-isothermal, cool minima in an upflow or hot maxima in a downflow may be stabilized in a superadiabatic region provided that the temperature difference and, again, the value of $`\beta r`$ are within specific intervals. As for horizontal tubes, sufficiently large values of $`\beta r`$ always lead to instability.
While monotonic instability preferentially evolves for displacements with large (but finite) wavenumber in the longitudinal direction, another mode of instability preferentially appears for large wavenumbers, i.e. overstable transversal oscillations. In the case of overstability, the drag force conspires with the magnetic tension force such that oscillations with growing amplitude result. For a horizontal tube this requires a downflow and is restricted to a certain interval of values for $`\beta r`$. The fact that overstability occurs preferentially for large wavenumbers suggests a local analysis for non-horizontal equilibrium tubes. We have carried out such an analysis for symmetric loops and found that from the four possible combinations of loop geometry (minimum/maximum) and flow (up/down) overstability is excluded for the minimum loop with an upflow. For the other cases overstability may appear within specific intervals for $`\beta r`$. In the case of an isothermal flux tube, however, overstability is restricted to maximum loops with downflow while a neutrally buoyant loop may only become overstable if it represents a maximum with an upflow. Since the approximation of slender flux tubes demands that the perturbation wavelength is large compared to the flux tube radius while overstability requires small wavelengths, the applicability of the present results is restricted to tubes of sufficiently small radius. Apart from the effects discussed so far, the introduction of external flows and drag forces leads to a decrease of the growth rates of unstable perturbations which is most significant for oscillatory instability.
The excitation of the overstable mode depends only on the direction of the flow and the value of $`\beta r`$; in particular, it cannot be stabilized by the stratification and therefore may appear also in convectively stable layers like overshoot zones or regions in radiative equilibrium. Excited locally, for instance in a loop formed by a downflow or an upflow, such oscillations could propagate as transversal tube waves. Under the influence of rotationally induced Coriolis forces these waves may even exhibit helicity and contribute to the field-regeneration mechanism which is necessary for the operation of a dynamo. In order to investigate this conjecture we have to extend the present formalism by moving to spherical geometry and including a (differential) rotation.
Convective flows in the deep parts of the solar convection zone have a typical time scale of the order of a month. Even a stable, stationary flux tube configuration within the convection zone cannot be expected to exist for a significantly longer time. This is in general accordance with the lifetime of large active regions. We have found that stable flux tube equilibria within a superadiabatic region require fine-tuned relations between the various parameters which determine the equilibrium such that they probably are not of great practical importance. A realistic convection zone, of course, is much more complicated than can be expressed by the simple analytical examples treated here. For example, the superadiabaticity probably shows significant spatial variations, be it between upflow and downflow regions or related to differential rotation (Durney, 1989). Hence, the estimates given in the preceding sections must not be taken too serious; however, they strongly indicate that, in the long run, the unstable stratification of the convection zone itself cannot be overcome.
6. Dynamics of flux tubes in a convection zone
In the preceding chapters we have investigated in some detail certain aspects of the structure and dynamics of concentrated fields. We may have obtained some pieces of a yet unfinished jigsaw puzzle in this way but we certainly are not in a position to present a full theory. In the present chapter we shall nevertheless try to sketch a tentative picture of the magnetic field dynamics in stellar convection zones, based upon our own results and on the work of other researchers. This picture is largely based on heuristic arguments, sometimes supported by more solid results. It is not thought as a comprehensive model but more as an orientation and stimulus for further work.
6.1 Size distribution
We have already discussed in Ch. 2 that the magnetic Rayleigh-Taylor instability and other fragmentation processes tend to produce magnetic structures with sizes of less than 100 km within the convection zone. Fragmentation proceeds until the fragments are so small that they merge by magnetic diffusion as fast as they are formed, i.e. until the diffusive time scale becomes equal to the growth time of the instability considered. For the case of the interchange instability this minimum fragment size, $`d_i`$, is given by Eq. (2.12) which we repeat here:
$$d_i=\left(\frac{2R\eta ^2}{v_A^{\mathrm{\hspace{0.17em}2}}}\right)^{1/3}$$
$`\left(2.12\right)=\left(6.1\right)`$
($`R`$: radius of curvature, $`\eta `$: magnetic diffusivity, $`v_A`$: Alfvén velocity). Taking equipartition fields, i.e. Alfvén velocity equal to $`v_c`$, the typical velocity of convective flows, and $`R`$ equal to $`L`$, the length scale of the convective flows, we can determine $`d_i`$ for different depths in the solar convection zone using the model of Spruit (1977b). For this rough estimate we use the depth as typical size of the dominant convective cell. The result is given in the following table:
| | Depth (cm) | $`\eta `$ (cm$`{}_{}{}^{2}`$s<sup>-1</sup>) | $`v_c`$ (cm$``$s<sup>-1</sup>) | $`d_i`$ (cm) | $`d_r`$ (cm) | $`\tau _r`$ (s) | $`\tau _s`$ (s) |
| --- | --- | --- | --- | --- | --- | --- | --- |
| | $`1.010^8`$ | $`1.510^6`$ | $`1.310^5`$ | $`3.010^3`$ | $`3.410^4`$ | $`7.710^2`$ | $`3.010^1`$ |
| | $`1.010^9`$ | $`1.010^5`$ | $`3.210^4`$ | $`2.710^3`$ | $`5.610^4`$ | $`3.110^4`$ | $`9.110^2`$ |
| | $`1.010^{10}`$ | $`7.810^3`$ | $`9.110^3`$ | $`2.510^3`$ | $`9.310^4`$ | $`1.110^6`$ | 2.1 |
| | $`1.810^{10}`$ | $`7.510^3`$ | $`4.910^3`$ | $`4.410^3`$ | $`1.710^5`$ | $`3.710^6`$ | 2.3 |
Tab. 8: Properties of magnetic filaments as a function of depth in the solar convection zone. $`\eta `$: molecular magnetic diffusivity; $`v_c`$: convective velocity; $`d_i`$: diffusive scale for the interchange instability; $`d_r`$: resistive boundary layer thickness; $`\tau _r=d_r^2/\eta `$: resistive diffusion time; $`\tau _s`$: radiative diffusion time on spatial scale $`d_i`$.
The table also gives $`d_r`$, the thickness of the resistive boundary layer:
$$d_r=\left(\frac{L\eta }{v_c}\right)^{1/2}.$$
$`\left(6.2\right)`$
This quantity is determined by the balance of magnetic diffusion and advection by convective flows and represents the scale of structures formed by kinematical magnetic flux expulsion. The relevant time scale for this process is the eddy turnover time, $`L/v_c`$, which is equal to the resistive diffusion time on the spatial scale $`d_r`$, viz.
$$\tau _r=\frac{d_r^2}{\eta }.$$
$`\left(6.3\right)`$
In the dynamical case Eq. (6.2) still gives a lower limit: Structures smaller than $`d_r`$ diffuse too rapidly to be held together by the convective cell and therefore cannot be sustained individually. On the other hand, we see from Table 6.1 that always $`d_i<d_r`$ which means that the growth time for the interchange instability for a structure of size $`d_r`$ is shorter than its diffusion time, i.e. the time scale of flux expulsion. Consequently, even if flux expulsion produces structures with sizes larger than $`d_r`$, these are actually bundles of smaller filaments with a typical size $`d_i`$. It is the very influence of the collecting flow itself which leads to fragmentation: It provokes interchanging by deforming the magnetic structures as well as by exerting a destabilizing pressure gradient at the interface between magnetic and non-magnetic gas (see the discussion in Schüssler, 1984b).
The last column in Table 6.1 gives the radiative diffusion time, the time scale which determines the heat exchange between a fragment of size $`d_i`$ and its environment,
$$\tau _s=\frac{d_i^2}{\eta _s},$$
$`\left(6.4\right)`$
with the radiative diffusivity $`\eta _s`$ given by
$$\eta _s=\frac{16\sigma T^3}{3\rho ^2c_V\kappa _R}$$
$`\left(6.5\right)`$
($`\sigma `$: Stefan’s radiation constant; $`T`$: temperature; $`c_V`$: specific heat capacity; $`\kappa _r`$: Rosseland mean opacity). For both resistive and thermal diffusion we have assumed that turbulence on length scales smaller than the filament size which could give rise to turbulent diffusivities is suppressed by the magnetic field. If the field is in equipartition with flows on the dominant convective scale it will always be stronger than the respective equipartition field on any other scale of the convective/turbulent velocity field.
We find from Table 6.1 that $`\tau _s`$ is very small compared to any relevant dynamical time scale throughout the whole convection zone. This statement remains valid if the larger spatial scale $`d_r`$ is used in Eq. (6.4). Furthermore, since we always have $`d_i<d_r`$, the diffusion time on the scale of the filaments, $`d_i`$, is shorter than the dynamical time scale $`L/v_c=\tau _r`$. Consequently, for magnetic structures which have been fragmented down to the resistive diffusion limit, $`d_i<100`$ m, as well as for structures of the size of the resistive boundary layer, $`d_r1`$ km, which are formed by flux expulsion, two statements can be made which become important in the following sections:
\- The magnetic diffusion time is small enough to allow mass exchange between magnetic structures and their surroundings within a dynamical time scale (e.g. while being stretched by differential rotation), and
\- the thermal diffusion time is sufficiently small to ensure temperature equilibrium between interior and exterior of a filament during its dynamical evolution.
As we have seen above, convective flows accumulate magnetic flux and form larger structures which appear in the form of bundles of small filaments, not monolithic flux tubes. For spatial scales which are much larger than the filament size a mean field treatment can be employed (cf. Parker, 1982b) and appropriate turbulent diffusivities have to be used. Knobloch (1981; see also Knobloch and Rosner, 1981, and references therein) has taken this approach and used the scaling laws given by Galloway et al. (1978) for nonlinear Boussinesq magnetoconvection to estimate a size spectrum of magnetic structures in a turbulent fluid for which he assumed a Kolmogorov spectrum with a (viscous) cut-off. Even without taking into account fragmentation processes he found that most of the structures have sizes at or below the length scale defined by the cut-off. He comes to the conclusion that larger structures can only be formed in the form of flux tube bundles whose size distribution is difficult to obtain.
This may not be the whole story though: A mechanism not mentioned so far is coalescence of two parallel twisted flux tubes with the same sense of twist (Parker, 1982e, 1983a,b; Bogdan, 1984; Choudhuri, 1988). If the tubes collide a neutral sheet forms at their interface, the azimuthal field reconnects and builds a binding sheath of magnetic flux about both tubes. An azimuthal field component at the same time exerts a stabilizing influence with respect to the interchange and Kelvin-Helmholtz instabilities. Consequently, besides advection by convective flows and fragmentation, coalescence may be another important factor which determines the size spectrum of flux tubes in the convection zone. First attempts to include this effect in a consistent treatment of the size distribution have been made by Bogdan and Lerche (1985) and by Bogdan (1985) who found that, in principle, sunspot-size structures can be produced within the convection zone. However, the models used so far are very idealized and it is not clear whether the results hold under more realistic conditions, i.e. non-straight flux tubes, inclusion of fragmentation processes, convective and Parker-type instability, advection by large scale flows, magnetic buoyancy. All these effects tend to either fragment large structures or to quickly remove them from the deeper layers of the convection zone. It seems doubtful that under these conditions sunspot-size structures can be formed but this claim can be substantiated or disproved only by more detailed models which include also the more subtle effects of flux tubes and flux tube arrays in a convectively unstable medium like ‘convective counterflow’ (Parker, 1985a,b) and ‘convective propulsion’ (Parker, 1979e).
How can we reconcile our picture of the magnetic field in the deep convection zone as an ensemble of very thin flux tubes with observations of the surface fields where a whole spectrum of structures from large spots to small magnetic elements is present? The key to an answer seems to be the fact that except for the first phase of flux emergence in an active region only fragmentation of large structures into smaller structures is observed, never the opposite process: Old active regions do not ‘rejuvenate’ and again form pores or spots unless new magnetic flux erupts. The lognormal distribution of sunspot umbral areas found by Bogdan et al. (1988) also is consistent with their origin in a fragmentation sequence.
We can understand the observation of large sunspots in view of a convection zone which continually shreds and fragments magnetic structures if the magnetic flux does not originate there but is injected in large portions from below (where it resides in a non-fragmented form) and rises rapidly to the surface. Indeed, flux emergence and the appearance of large structures at the surface always take place within the very first days of the life of an active region. The simulations of Moreno-Insertis (1984, 1986) show that a kink-unstable large flux tube traverses the convection zone and breaks out at the surface within a few days. Even within this short time fragmentation apparently has occurred since a whole spectrum of structures appears at the photosphere and large sunspots always form out of a number of pieces (McIntosh, 1981) which seem to know where to go (the “rising magnetic tree” of Zwaan, 1978): Initially the rise and emergence of only weakly fragmented parts of the erupting flux tube is faster than the ongoing fragmentation processes.
After flux emergence has come to an end, fragmentation proceeds until all magnetic flux at the surface is in the form of small network elements with a size of about 100 km <sup>*</sup><sup>*</sup>This size is much larger than the resistive scale because surface fields can temporarily achieve stable configurations and resist further fragmentation. This is further discussed in Sec. 6.3.. In the deep convection zone, all flux tends to become fragmented down to the diffusion limit. The formation of filament bundles by convective flows and flux tube coalescence may lead to somewhat larger structures but their actual size distribution is difficult to determine. In any case, these structures are very fragile: They are always subject to the various instabilities and fragmentation processes and filament bundles are closely coupled to the changing pattern of convective flows. We do not believe that large active regions and sunspots are formed in this way: It seems impossible to store even moderately large flux tubes within the convection zone for a sufficiently long time against their inherent buoyancy and instability. Large active regions probably are direct evidence for the genuine dynamo process operating on a more ordered field in a less unstable environment than the convection zone proper.
6.2 The relation of the basic forces
We have argued in Ch. 2 and Sec. 6.1 that the magnetic field in a stellar convection zone consists of an ensemble of small filaments whose dynamics can be described with aid of the approximation of slender flux tubes discussed in Ch. 3. We found that the evolution of a flux tube is determined by three basic forces, the buoyancy force $`F_B`$, the curvature force $`F_C`$ and the drag force $`F_D`$, all of which are directed perpendicular to the flux tube axis while gravity and gas pressure determine the momentum balance in the direction along the field lines.
Let us discuss the relative importance of these forces under the physical conditions in a convection zone. If we presume temperature equilibrium between the flux tube and its environment (a reasonable assumption for small filaments as we have seen in the preceding section) we may use Eqs. (3.26) and (3.28) to obtain the order of magnitude of the various forces per unit length of the flux tube:
$$\begin{array}{ccc}\hfill F_B& =\frac{B^2a^2}{8H_{pe}}\hfill & \left(6.6\right)\hfill \\ & & \\ \hfill F_C& =\frac{B^2a^2}{4R}\hfill & \left(6.7\right)\hfill \\ & & \\ \hfill F_D& =C_D\rho _{e0}av_{}^2\hfill & \left(6.8\right)\hfill \end{array}$$
($`a`$: flux tube radius, $`H_{pe}`$: external pressure scale height, $`\rho _{e0}`$: external density, $`v_{}`$: velocity component perpendicular to the flux tube due to large scale convection, differential rotation or the motion of the flux tube itself, $`R`$: radius of curvature). Ignoring factors of order unity we may write for the ratio of buoyancy and curvature force to drag force, respectively:
$$\begin{array}{ccc}\hfill \frac{F_B}{F_D}& =\left(\frac{a}{H_{pe}}\right)\left(\frac{B}{B_e}\right)^2\left(\frac{v_c}{v_{}}\right)^2\hfill & \left(6.9\right)\hfill \\ & & \\ \hfill \frac{F_C}{F_D}& =\left(\frac{a}{R}\right)\left(\frac{B}{B_e}\right)^2\left(\frac{v_c}{v_{}}\right)^2.\hfill & \left(6.10\right)\hfill \end{array}$$
Here $`B_e=v_c\sqrt{4\pi \rho _{e0}}`$ is the equipartition field strength with respect to convective flows of typical velocity $`v_c`$. We see from these ratios that for equipartition fields and $`v_{}v_c`$ the drag force dominates the dynamics of a thin flux tube: $`a/H_{pe}1`$ and $`a/R1`$ in fact are conditions for the applicability of the approximation of slender flux tubes. For strongly fragmented fields in the deep convection zone they are well met: $`a<1`$ km, $`H_{pe}10^4`$ km; the curvature force comes into play only for extreme distortions ($`Ra`$) of the filament. Consequently, a magnetic field structured in this way follows passively any large-scale flow with relative velocity of the order of the convective velocity, be it convection itself, differential rotation or, possibly, meridional circulation. Only structures with $`aH_{pe}`$ or larger (flux tubes of sunspot size) or strong fields with $`B>B_e`$ (partially evacuated flux tubes at the surface) can avoid being severely distorted by large scale flows.
On the other hand, the growing distortion of the tube leads to an increase of its length and, if the mass content stays constant, to an increase of the field strength. In the case $`\beta 1`$ we find that the change of the magnetic field strength is proportional to the length increment, i.e.
$$\frac{\mathrm{\Delta }B}{B}=\frac{\mathrm{\Delta }l}{l}.$$
$`\left(6.11\right)`$
Doubling the flux tube length doubles the field strength and the curvature force increases by a factor 4. Consequently, a flux tube with constant mass content quickly becomes dynamically active and resists to further distortion and stretching. The curvature force increases rapidly and a balance between curvature and drag force establishes itself. However, the assumption of a constant mass content does not hold for two important cases:
\- The distortion of flux tubes is strongest in the upper layers of the convection zone where the flow velocities are large. Here the density is small and if parts of the tube are still located in the deep layers they represent a large mass reservoir: Matter can flow along the tube to fill the volume created by stretching the tube, and
\- for small filaments formed by the various fragmentation processes and instabilities we have seen in the preceding section that the diffusive time scale is smaller than the dynamical time scale which governs the distortion process. Consequently, matter can diffuse into the flux tube rapidly enough to fill the volume created by its distortion. This is particularly relevant for the deep layers where no further mass reservoir exists for flux tubes contained within the convection zone.
Of course, both mechanisms can be relevant for the same flux tube. Their effect is to keep the flux tube in a state of passiveness with respect to further stretching. We conclude that due to instability and fragmentation the magnetic field structures that reside in the convection zone for an extended period of time (i.e. longer than about a month, the dynamical time scale in the deep layers) are passive with respect to external flows (convection, differential rotation, meridional circulation), i.e. they are strongly coupled to and distorted/fragmented by these flows and probably never reach an equilibrium.
An exception from this rule are large flux tubes whose radius is of the order of the local scale height or larger. Here all three basic forces may become comparable and the results concerning equilibrium structure and stability obtained in Chs. 4 and 5 can be applied. We have found that in most cases such a flux tube is unstable and parts of it rapidly erupt towards the surface layers while other parts sink down below superadiabatic layers of the convection zone. Although the applicability of the approximation of slender flux tubes is somewhat doubtful for such large tubes we believe that the results are qualitatively correct also in this case. Another exception from the rule of passive magnetic fields are the flux tubes in the surface layers. They will be discussed in the following section.
6.3 The peculiar state of the surface fields
The observable surface layers of a star like the Sun are quite different from the state of the deep convection zone. They represent the transition region between the strongly superadiabatic top of the convection zone and the stably stratified photosphere where most of the energy carried by the overshooting convective motions escapes into space by means of radiation. The gas is strongly stratified since the relatively low temperature leads to a small scale height while the temperature fluctuations become a significant fraction of the temperature itself.
Under these conditions magnetic fields can be concentrated to field strengths far beyond equipartition with convective motions. Magnetic fields are swept to the granular downflow regions and concentrated to about equipartition by the horizontal flows of granular convection. These flows at the same time are responsible for carrying heat to the downflow regions which compensates for the radiative losses. Since the growing magnetic field retards the motions it throttles the energy supply and the magnetic region cools down. This leads to an increase of the magnetic field since the gas pressure in the magnetic region becomes smaller. Furthermore, hydrostatic equilibrium on the basis of a reduced temperature tends to reestablish itself via a downflow which gives rise to the superadiabatic effect (Parker, 1978): An adiabatic downflow in a magnetic flux tube which is thermally isolated from its surroundings leads to a cooling of the interior with respect to the superadiabatically stratified surroundings and a partial evacuation of the the upper layers ensues. Pressure equilibrium with the surrounding gas is maintained by a contraction of the flux tube which increases the magnetic pressure. In this way, the magnetic field can be locally intensified to values which are only limited by the confining pressure of the external gas.
It has been shown by a number of authors (Webb and Roberts, 1978; Spruit and Zweibel, 1979; Unno and Ando, 1979) that the superadiabatic effect in the case of a vertical flux tube which is in magnetostatic and temperature equilibrium with its environment drives a convective instability in the form of a monotonically increasing up- or downflow. Consequently, the initial downflow due to the radiative cooling is enhanced by this effect leading to an even stronger amplification of the magnetic field, a process which is often referred to as convective collapse. A more detailed discussion of this process has been given elsewhere (Schüssler, 1990).
Cooling due to suppression of convective motions and the superadiabatic effect have the consequence that most of the observed magnetic flux in the solar photosphere has a field strength in excess of the equipartition value of a few hundred Gauss. It is organized in a hierarchy of structures which have a magnetic pressure comparable to the gas pressure in their apparently non-magnetic environment. This hierarchy extends from large sunspots (diameter $`>`$ $`510^4`$ km, field strength up to 4000 Gauss) to small magnetic elements (diameter $`<`$ 200 km, field strength about 2000 Gauss at optical depth unity). In contrast to the circumstances in the deep convection zone which have been discussed in the preceding section, under these conditions buoyancy becomes the dominating force. Taking 500 Gauss for the equipartition field strength and using Eq. (6.9) we find that $`F_B/F_D`$, the ratio of the buoyancy to the drag force due to convective motions, is of the order of 10 for magnetic elements ($`a=100`$ km, $`B=2000`$ Gauss) and amounts to about 3000 for large sunspots ($`a=10^4`$ km, $`B=4000`$ Gauss). By the dominance of buoyancy the magnetic structures are forced to become straight and vertical in the surface layers: A very small inclination suffices to compensate any drag force exerted by the external velocity fields (Schüssler, 1986, 1987).
Besides the ability to resist from being distorted and deformed by convective motions another peculiarity of the surface fields is the existence of configurations which are (at least temporarily) stable with respect to the interchange instability and other fragmentation processes. Buoyancy and the large field strength prevent the flux concentrations from being disrupted by local convection. Large structures like sunspots (magnetic flux larger than about $`10^{19}`$ mx) are stabilized against the interchange instability by buoyancy (Meyer et al., 1977) while small magnetic elements (magnetic flux less than a few times $`10^{17}`$ mx) may be stabilized by whirl flows which surround them (Schüssler, 1984b). Such whirls form naturally in the narrow downflow regions of convection in a strongly stratified medium (e.g. Nordlund, 1985a). These downflows are enhanced around small magnetic flux concentrations by thermal effects (Deinzer et al., 1984).
The existence of two regimes of stable magnetic flux tubes connected by a range of magnetic configurations which are subject to the interchange instability certainly should have an influence on the observable size distribution of magnetic structures. On the one hand large sunspots may live for an extended period of time during which they only show a slow decay while, on the other hand, small features exist with sizes up to the limit where the stabilization due to a surrounding whirl flow becomes inefficient. The verticality of the magnetic structures due to the dominating buoyancy force facilitates their organization in a network pattern defined by the downflow regions which are the loci of convergence of the horizontal convective flows. The probability of encounters and coalescence is much larger within such a network than for a more random spatial distribution. For all these reasons we think that observed size distributions of surface fields (e.g. Spruit and Zwaan, 1981) are by no means representative for the conditions in the deep layers of the convection zone where most of the effects which determine the particular properties of surface fields become irrelevant.
The special properties of the surface fields are restricted to the photosphere and the rather narrow layer of strong superadiabaticity in the upper convection zone of about $`1000`$ km depth. As discussed in more detail elsewhere (Schüssler, 1987) the non-observation of a systematic dependence of the dynamics of observed flux concentrations on their size together with a consideration of the dominating forces strongly support the cluster model of sunspots (Parker, 1979c; Spruit, 1981b). This model assumes that a spot is a conglomerate of a large number of small magnetic filaments with a diameter of 100 to 1000 km which are closely packed due to buoyancy in the uppermost layers of the convection zone to form the visible spot umbra. Below some merging level (which is situated not much deeper than 1000 km) the spot splits into its fibril components which are passive with respect to the convective motions which continually stretch and distort them; this causes the slow decay of the spot and the fragmentation of its magnetic flux into the network fields.
We have seen that the observable surface fields probably are in many respects different from the conjectured properties of the magnetic fields in the deep layers of the convection zone. The strong superadiabaticity of the uppermost convective layers and the thermal effects associated with the radiative release of the transported energy into free space entail a number of peculiar effects which cannot be expected to operate in the main body of the convection zone. Only at the surface can we expect the field strength to significantly exceed the equipartition limit and only there do we find mechanisms which can stabilize magnetic flux tubes from being quickly disrupted by instabilities. Consequently, it is improbable that a flux concentration observed at the surface maintains its identity as a single flux tube throughout the whole or even a significant part of the convection zone. Visible sunspots and all other observed structures probably fragment into small filaments below about 1000 km depth. These filaments with sizes down to the diffusion limit are passive with respect to the convective motions, they are stretched and deformed, accumulated and dispersed by the action of the dominating drag forces. Due to their inherent stability, the surface fields can tolerate the distortion of their ‘roots’ for a certain amount of time until they are forced to react accordingly, be it by slowly dispersing and eventually breaking apart as in the case of sunspots, be it by a continuous rearrangement of the small-scale fields in the supergranular network.
6.4 Consequences for the dynamo problem
The presently favored tool to describe the origin of solar activity is the theory of dynamo action in a turbulent medium which started with the pioneering work of Parker (1955). Beginning in the 1960s the next landmark was placed by the Potsdam group (Krause, Rädler, Rüdiger, Steenbeck and others; see Krause and Rädler, 1980, for an overview) who used a statistical approach (mean field theory) in the kinematical limit (passive magnetic field, no back-reaction on the turbulent flows). They found the so-called $`\alpha `$-effect which can lead to dynamo action in turbulent fluids which lack mirror symmetry, e.g. due to the influence of rotation. Many large-scale properties of the solar cycle could be successfully reproduced by ‘$`\alpha \omega `$-dynamos’ which are based on the combined induction effects of turbulence ($`\alpha `$) and differential rotation ($`\omega `$). Reviews of these models have been given, among others, by Stix (1976, 1981a, 1982), Parker (1979a) and Yoshimura (1981).
It was again Parker (1975a) who raised doubt concerning the theory of turbulent dynamo action within the convection zone. He argued that, if most of the toroidal magnetic flux is in the form of large flux tubes as indicated by the existence of sunspots and active regions, it cannot be stored in the convection zone for times comparable with the solar cycle period: Buoyancy removes magnetic flux from the convection zone much too fast for the induction mechanisms to operate efficiently, in particular for the differential rotation to generate a sufficient amount of toroidal magnetic flux within a cycle period. The instabilities of flux tubes due to superadiabaticity, flows along the field lines, rotation and external flows investigated by Spruit and van Ballegooijen (1982), van Ballegooijen (1983), Moreno-Insertis (1984, 1986), van Ballegooijen and Choudhuri (1988), and in Ch. 5 of this work aggravate the problem even more.
The assumption of very small flux tubes or of a large turbulent viscosity may reduce the velocity of buoyant rise drastically (Unno and Ribes, 1976; Schüssler, 1977; 1979; Kuznetsov and Syrovatskii, 1979) but, on the other hand, leads to a strong coupling between the flux tubes and the convective flows which would destroy their azimuthal orientation and quickly raise them to the solar surface within the convective time scale of about one month. Again, the flux cannot be contained within the convection zone for a sufficiently long time.
Recently it has been argued by Parker (1987a,d) that the magnetic flux observed to emerge in large complexes of activity of relatively small latitude extension would fill as an equipartion field a considerable part of the underlying convection zone and interfer significantly with the convective energy transport. He proposed that the suppression of vertical convective heat transport by a band of horizontal (azimuthal) field in a convection zone leads to an overlying ‘thermal shadow’, i.e. a cool region of enhanced density, whose weight could balance the magnetic buoyancy and thus keep the field down within the convection zone. He studied in some detail the convective flows set up by a thermal shadow (Parker, 1987b) and the Rayleigh-Taylor instability caused by the pile-up of heat beneath the magnetic layer (Parker, 1987c). He envisages that a thermal relaxation oscillation evolves which leads to the intermittent eruption of magnetic flux in the form of rising, hot plumes (Parker, 1988a). The dynamical instability of a flux sheet of finite lateral extent with respect to sideways displacements leads to a lateral velocity of the order of a few m/s which he connects with the observed latitudinal motions of the solar activity belts (Parker, 1988b).
Parker’s conjecture is based on a number of illustrative calculations of idealized problems in order to permit an analytical treatment. While the thermal shadows may well have important dynamical effects it is not shown that they can keep large amounts of magnetic flux in the superadiabatic parts of the convection zone for time intervals comparable to the period of the solar cycle. Actually, in view of the instabilities discussed by Parker himself (Rayleigh-Taylor instability of top and bottom, lateral dynamical instability) and the quick disruption of a magnetic layer even in a stably stratified region (Cattaneo and Hughes, 1988), this seems to be rather improbable.
However we turn the problem, we run into difficulties if we assume the dynamo to operate within the convection zone proper. If the magnetic field is organized in large, dynamically active structures, these are unstable, buoyant and rapidly lost from the deep layers of the convection zone. If the field is diffuse or consists of very small structures (the latter view is favored in this work), it is passive with respect to the motion of the fluid (convection, rotation, …) and will be carried to the surface within the convective time scale. Moreover, such a passive field is not in accordance with basic features of solar activity: In a large active region, magnetic flux erupts coherently in large quantities and within a few days; large sunspots form which comprise a significant fraction of the total flux erupting during the cycle; Hale’s polarity rules are obeyed nearly strictly, not in a statistical sense. On the other hand, the relative velocities due to differential rotation and the convective velocities are of the same order of magnitude as is well known from surface observations (e.g. Schröter, 1985) and also shown by the rotational splitting of $`p`$-modes for the deeper layers of the convection zone (e.g. Duvall et al., 1987; Brown et al., 1989; Dziembowski et al., 1989). How can passive fields be predominantly toroidal and obey the polarity rules under such conditions ?
We cannot avoid the conclusion that the original source region at least of the large, sunspot-forming active regions cannot be the convection zone proper. It cannot be the radiative core of the Sun either since the time scale of 22 years for the magnetic cycle does not allow a penetration into the radiative interior because of its large electrical conductivity which leads to a skin effect. Consequently, a number of authors (Spiegel and Weiss, 1980; Galloway and Weiss, 1981; van Ballegooijen, 1982a,b; Schmitt and Rosner, 1983; Schüssler, 1983; DeLuca and Gilman, 1986; Durney, 1989; Durney et al., 1990) have proposed a boundary layer of overshooting convection below the convection zone proper as a favorable site of the solar dynamo. There, in a region of ‘mild’ convection and turbulence, we may suppose that differential rotation dominates all other velocity fields and generates predominantly toroidal magnetic fields. The failure of large-scale simulations to reproduce the observed characteristics of solar activity (in particular, the direction of the latitude drift of the activity belts) by dynamically consistent 3D-models of the convection zone (Gilman and Miller, 1981; Gilman, 1983; Glatzmaier, 1985a) led their proponents to the same conclusion (Glatzmaier, 1985b).
The subadiabatic stratification of an overshoot region alleviates the stability problems and a number of mechanisms is available which may hold a growing toroidal magnetic flux there for a time comparable with the cycle period (cf. Schüssler, 1983). We may note in particular the ‘turbulent diamagnetism’ briefly discussed in Sec. 2.2, which is akin to the flux expulsion process and transports magnetic flux antiparallel to a gradient of turbulent intensity. Krivodubskii (1984, 1987) has given quantitative estimates of the diamagnetic effect using models of the solar convection zone. Since, at least in the kinematical limit, this mechanism operates equally well for vorticity (Schüssler, 1984a) it leads to the formation and maintenance of a magnetic shear layer at the bottom of the convection zone (including the overshoot region), just the situation we envisage as a favorable setup for the operation of the solar dynamo.
Besides differential rotation we need a regeneration mechanism for the poloidal field component in order to close the dynamo cycle. Such a mechanism might be provided by the usual $`\alpha `$-effect due to cyclonic convection or by an analogous electromotive force generated by waves propagating along the toroidal field in a rotating system. Examples of such waves are slow magnetostrophic waves driven by the magnetic Rayleigh-Taylor instability (Schmitt, 1984; 1985) or transversal flux tube waves excited by differential rotation (van Ballegooijen, 1983) or overstability in an external flow as described in Ch. 5. A boundary layer dynamo model may thus show many of the properties of the ‘classical’ $`\alpha \omega `$-dynamo models which have been so successful in describing basic features of the solar cycle. This supports the conjecture that the mathematical description of the field regeneration mechanism (i.e. a mean current parallel to the mean field) in the mean induction equation is basically correct. Note, however, that a kinematical approach is not justified for a boundary layer dynamo as envisaged above since the magnetic and kinetic energy densities are comparable for equipartition fields. A consistent theory for an $`\alpha `$-effect in an overshoot region with a strong magnetic field, shear flow and weak turbulence remains to be developed.
Eventually, magnetic flux is released into the convection zone proper by instability or buoyancy. The length and time scales of the formation of large active regions are well in agreement with the characteristics of the eruption of a kink-unstable large flux tube originally situated near the bottom of the convection zone (Moreno-Insertis, 1984; 1986). Moreover, we could think about a modification of Parker’s thermal shadow scenario: If the amount of magnetic flux is small enough that it ‘fits’ into the subadiabatic overshoot layer there is almost no interference (except from opacity effects, cf. Parker, 1984b) with the energy flux which is mainly carried by radiation – neither a shadow nor a significant pile-up of heat evolve. If the magnetic flux layer intrudes significantly into the superadiabatic convection zone proper a thermal shadow and a pile-up of heat in the magnetic overshoot region are the consequence. Both effects provoke Rayleigh-Taylor instability and the whole magnetic layer is disrupted and erupts towards the surface: a large active region is born. In this picture the thermal shadow does not primarily keep the flux submerged but, on the contrary, is the agent of the eruption of magnetic flux. It is tempting to speculate whether the larger efficiency of differential rotation in generating azimuthal magnetic field near the equator might provide the excess azimuthal flux that drives the magnetic layer unstable and produces large active regions while the azimuthal flux in the polar regions always fits into the subadiabatic region.
As shown by Choudhuri and Gilman (1987) and Choudhuri (1989) the influence of rotation on the dynamics of rising loops might be quite significant. In order to avoid the flux erupting in the polar regions either the flux density in the overshoot region must be quite high, i.e. of the order of $`10^5`$ Gauss – which alleviates the storage space problems pointed out by Parker (1987a) but increases the buoyancy problems – or the transport is dominated by the drag of predominantly radial convective flows whose motion the magnetic fields passively follow. Obviously, the rôle of rotation in the transfer of magnetic fields in a convection zone deserves further consideration.
Altogether, the results obtained here support the picture of a boundary layer dynamo sketched above. We have not been able to find a plausible way by which magnetic structures within the convection zone proper could escape from being extremely distorted, dismembered down to the diffusion limit, and eventually becoming completely passive with respect to the convective flows. Large active regions with their specific properties cannot arise from such a kind of field. On the other hand, if a less fragmented magnetic structure at some instant exists in the convection zone (e.g. having been injected from below) in most cases it will immediately become unstable by one of the mechanisms discussed in Ch. 5, part of it will erupt at the surface while another part sinks down into the subadiabatic boundary layer. After flux emergence has come to an end, the magnetic structures are more and more fragmented and shredded until they are merge into the extremely filamented and distorted genuine convection zone fields.
It has been already remarked by Parker (1982b, see also Weiss, 1981c) that an ensemble of passive flux tubes can be described by a kinematic mean field theory in much the same way as in the theory developed by the Potsdam group (see Krause and Rädler, 1980). We may speculate (e.g. Durney, 1989) that part of the field within the convection zone is maintained by a turbulent dynamo in the classical sense which contributes to the background fields at the surface while a boundary layer dynamo provides the source of the big active regions and the large-scale features of solar activity. Anyway, a major revision of the conventional picture of a convection zone dynamo seems to be in order. How the dynamic boundary layer dynamo operates is not understood yet and the detailed study of the dynamics of concentrated fields in a convection zone and its adjacent lower overshoot layer has just started.
7. Outlook
As it turns out so often, most problems remain to be solved. In view of the unsatisfactory state of our understanding of stellar convection and turbulence, a closed and complete theory of magnetic fields in convection zones is not in sight. The theory of photospheric magnetic fields is much more advanced thanks to the availability of detailed measurements and the close connection between theoretical and observational work. Some information about the internal magnetic field and large-scale velocity structure will be obtained in the future by helioseismology but, perhaps with the exception of differential rotation, we are not too optimistic about the prospects of obtaining much more stringent observational boundary conditions. Consequently, besides the ongoing efforts of numerical simulation and the detailed analysis of model problems this field of research will remain open for conjecture, speculation and the presentation of more or less ingenious scenarios.
Numerical simulations will partially substitute unavailable observations. As the development of ever faster computers and sophisticated numerical methods to adequately use them proceeds, simulations will grow more realistic as three-dimensional, compressible MHD calculations with high spatial resolution become available and will provide an immensely useful tool for understanding the magnetic field dynamics in the convection zone. However, in contrast to some fashionable folklore existing and forthcoming numerical simulations do not make other approaches obsolete. The dynamics of motions and magnetic fields in the solar convection zone extends over huge ranges of temporal and spatial scales which in both cases comprise more than ten decades. Since only a small part of these can be covered by any simulation in the foreseeable future, artificial boundaries have to be introduced, certain scales are ignored and others are included only in a parametrized form. Such parametrizations can only be made in a sensible way if they are based on a sound understanding of processes which determine the properties of flows and fields on the scales which they attempt to describe.
We have given some arguments in favor of the conjecture that magnetic fields in stellar convection zones are strongly fragmented and can be treated on the basis of the approximation of slender flux tubes. In the deep parts of a convection zone the scale height is very large such that the approximation is justified even for structures containing the magnetic flux of a whole active region. The investigation of flux tube dynamics in a realistic convection zone is a promising path for future research. Equilibrium structures, stability and dynamical evolution of flux tubes in prescribed velocity fields can be determined on the basis of the methods described in this work, guided by forthcoming 3D simulations of the large-scale convective flows and observational results on differential rotation. As far as the linear stability analysis is concerned, this requires a transformation of the formalism presented in Ch. 5 to spherical geometry and the inclusion of (differential) rotation. For realistic models of convection zones, the equations describing flux tube equilibria and the perturbation equations will have to be treated numerically. Work in this direction is in progress, being done in cooperation with Antonio Ferriz-Mas.
A possibility to gain information about the size distribution of magnetic structures in a convection zone is the application of methods taken from statistical mechanics. The evolution of the properties of an ensemble of flux tubes may be derived from a collisional Boltzmann equation which includes the effect of large-scale flows, diffusion, fragmentation, coagulation and other processes. Moreover, this approach opens a possibility to put the vague notion of a ‘flux tube dynamo’ on a firm theoretical basis. A cooperation with Tom Bogdan (Boulder), Antonio Ferriz-Mas and Michael Knölker on such a project is arranged.
Finally, the boundary layer or overshoot layer dynamo will remain a challenge. Kinematical theory probably cannot be applied since kinetic and magnetic energy density are of the same order of magnitude. Future work will focus on two approaches, i.e. the quantitative determination of a field regeneration process and the angular momentum distribution in an overshoot layer, and the investigation of nonlinear dynamo models with a given regeneration process which take into account the particular geometry and thermodynamics of the region as well as expulsion processes and magnetic instabilities. As for most of the discussed problem areas here again comprehensive 3D simulations and idealized/simplified approaches will play complimentary parts: The simulations help to identify the relevant processes and allow us a glimpse at phenomena which are observationally unreachable. They can guide us in picking the relevant pieces of physics to study in detail without falling into the trap of oversimplified or prejudiced concepts. An understanding of the physics governing these processes, of their general properties and the validity of their description in the simulation can only come from a detailed study in the spirit of theoretical physics.
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# Kinematic Properties and Stellar Populations of Faint Early-Type Galaxies. I. Velocity Dispersion Measurements of Central Coma Galaxies
## 1 Introduction
Dwarf elliptical galaxies (dEs) are low luminosity spheroidal systems with $`M_B>18`$ mag (Sandage & Binggeli 1984) that have a low surface brightnesses, $`\mu _{e,V}>22`$ mag arcsec<sup>-2</sup> (Ferguson & Binggeli 1994). Over the past two decades there has been considerable interest in studying dEs, despite the difficulty associated with observing these objects. In particular, there has been a number of extensive photometric studies that have concentrated on comparing dEs to the bright elliptical galaxies with $`M_B<20.50`$ mag (Graham & Guzmán 2003). Two main arguments support the hypothesis that they are structurally distinct classes.
Firstly, many authors in the past argued that de Vaucouleur’s $`R^{1/4}`$ law (de Vaucouleurs 1948) best fits the light profiles of bright Es. Meanwhile, Faber & Lyn (1983) and Binggeli, Sandage & Tarenghi (1984) showed that an exponential profile describes dwarf ellipticals. Secondly, as Kormendy (1985) notes, dEs and bright Es fall almost perpendicular to each other in the effective surface brightness–luminosity plot, $`\mu _eL`$ and are clearly different in the luminosity–central surface brightness, $`L\mu _0`$, plot. Furthermore, dEs and bright Es follow a different $`\mu _eR_e`$ relation (e.g. Wirth & Gallagher 1984; Capaccioli, Caon & C’Onforio 1992).
More evidence to strengthen the proposed dichotomy exists. Dwarf ellipticals seem to lie off the Fundamental Plane, the relation between the surface brightness at the effective radius, $`\mu _e`$, effective half-light radius, $`R_e`$, and the velocity dispersion, $`\sigma `$ (Bender, Burstein, & Faber 1992; de Carvalho, & Djorgovski 1992; Peterson, & Caldwell 1993). This disparity has been interpreted as a difference in the formation mechanism for dwarf and bright Es.
Nonetheless, there are many studies arguing for continuity in the dwarf-bright family. Just to mention a few: Caldwell 1983; Caldwell 1987; Caldwell & Bothun 1987; Ferguson & Sandage 1988; Hudson et al. 1997; Jerjen & Binggeli 1997; Jerjen, Binggeli, & Freeman 2000; and Karachentsev, et al. 1995. Some of these studies show that dEs exhibit the same central surface brightness and absolute magnitude relation as bright Es. Caldwell (1983) also pointed out that a continuous trend exists between colour and luminosity, while Caldwell & Bothun (1987) show the same continuity for the luminosity–metallicity relation.
Graham & Guzmán (2003; hereafter GG03)(see also Guzmán et al. 2003) offer a possible resolution of the differing views. They point out that the dichotomy in the luminosity–effective surface brightness relation, $`M_B\mu _e`$, and the luminosity–effective radius relation, $`M_BR_e`$, is a direct consequence of the linear relations between the luminosity, the central surface brightness, $`\mu _0`$, and the light profile shape, $`n`$. Furthermore, they argue that dEs and intermediate luminosity ellipticals follow a continuous sequence up to $`M_B20.5`$ mag. At this point bright ellipticals start showing evidence of evacuated cores, possibly coalescing black holes, causing their central surface brightness to decrease with increasing luminosity (GG03; Graham 2004 and references therein). The most massive Es may thus be the exception and not the rule to the empirical correlations defined by early–type galaxies that include the low luminosity galaxies and range over $`8`$ magnitudes (see GG03).
Up until the past couple of years most information on spectroscopic properties of dEs came from their line strength indices (Held & Mould 1994; Gorgas et al. 1997; Mobasher et al. 2002; Moore et al. 2002), and a handful of velocity dispersion measurements. Although more difficult to obtain due to low surface brightness of these objects, the number of papers in the literature which include velocity dispersion measurements of dEs in different clusters has increased (Bender & Nieto 1990; Brodie & Huchra 1991; Held et al. 1992; Bender, Burstein, & Faber 1992; Peterson & Caldwell 1993; Bernardi et al. 1998; Melhert et al. 2002; Hudson et al. 2001; De Rijcke et al. 2001; Simien & Prugniel 2002; Pedraz et al. 2002; Moore, Lucey, Kuntschner, & Colless 2002 (hereafter MLKC02); Geha, Guhathakurta, & van der Marel 2002, 2003; Guzmán et al. 2003; Bernardi et al. 2003; Smith et al. 2004, van Zee, Skillman, & Haynes 2004; De Rijcke et al. 2004). The NOAO Fundamental Plane survey (NFPS) of Smith et al. 2004 is the largest compilation including velocity dispersion measurements in low-redshift galaxy clusters up to date. Our Coma sample, although significantly smaller than that of the NFPS is complementary to this study. It includes a statistically representative sample of faint early-type galaxies in different environments within a cluster. Such sample is essential for testing current ideas on the formation and evolution of dEs.
Via recent studies, dEs have also been linked to the Butcher-Oemler effect (Butcher, & Oemler 1978; Butcher, & Oemler 1984). Observations of distant clusters reveal the existence of numerous star-forming, low-mass ’blue disk’ galaxies in clusters at $`z0.4`$. These galaxies, as shown by HST, are distorted small spirals which have disappeared from the present day clusters. The fate of these galaxies remains one of the most important unanswered questions in modern cosmology. The ’galaxy harassment’ model Moore, Lake & Katz (1996; 1998) explains how the dwarf spiral galaxies in clusters may evolve into today’s population of cluster dEs due to encounters with brighter galaxies and the cluster tidal field. The galaxy harassment model predicts differences in the properties of dEs located in the inner, high density, and outer, low density, regions of the clusters.
This is the first paper of a series in which we will characterise the kinematics and stellar populations of dEs and other low luminosity early–type galaxies as a function of the environment. Given the difficulty in distinguishing between dEs, dS0s, and dwarf spirals in Coma from ground-based images, we refer to all these objects as ‘early–type’ galaxies. For convenience, we define ‘faint’ early–type galaxies as those with $`M_B>20.50`$ mag, and the ‘bright’ ellipticals with $`M_B<20.50`$ mag. Here we describe spectroscopic observations and velocity dispersion ($`\sigma `$) measurements of $`69`$ faint early–type galaxies in the central $`30\mathrm{}\times 30\mathrm{}`$ region of the Coma cluster. We investigate whether these galaxies follow the luminosity–velocity dispersion ($`L\sigma `$) relation derived for bright ellipticals, and discuss the constraints on their formation epochs provided by the colour–$`\sigma `$ ($`C\sigma `$) relation. In future papers, we will test the implications of the ’galaxy harassment’ scenario by comparing the internal kinematics and stellar populations of the early–type galaxies in the core and the outskirts of the Coma cluster. Section 2 describes the photometric and spectroscopic observations of our sample galaxies, including the sample selection. In Section 3, we describe the data reduction technique and the velocity dispersion measurements. We investigate the $`L\sigma `$ relation in Section 4, and the $`C\sigma `$ relation in Section 5. A summary of our results is provided in Section 6.
## 2 Observations
### 2.1 Sample Selection
The selection of faint early–type galaxies in the Coma cluster was done using the photometry in $`U,B`$ and $`R`$ bands. We obtained the images with the WIYN/MiniMo and INT/Wide Field Camera. To select the faint early–type galaxy candidates in the central $`30\mathrm{}\times 30\mathrm{}`$ region of the Coma cluster we used the B–R vs. B colour-magnitude plane. We applied an absolute luminosity cutoff at $`M_B>17.3`$ mag (Ferguson, & Binggeli 1994; corresponding to apparent $`B>17.5`$ mag at the distance of the Coma cluster<sup>1</sup><sup>1</sup>1Throughout the paper we use $`H_0`$ = 70 km s<sup>-1</sup> Mpc<sup>-1</sup>, and a distance modulus for Coma cluster of 35.078, ($`d=99`$ Mpc).). To minimize contamination at the faint end by the background field disk galaxies at $`z<0.2`$, we applied another cutoff using the (U–B) vs. (B–R) colour-colour diagram at $`0.2<(UB)<0.6`$ mag, and $`1.3<(BR)<1.5`$ mag.
### 2.2 Spectroscopic Observations
We observed the Coma cluster faint early–type galaxies during 1998 May 23–26, and 1999 May 14–19 on the 3.5 m WIYN telescope at Kitt Peak National Observatory with the multi-fibre spectrograph HYDRA. We used the 600 l mm<sup>-1</sup> grating in the $`2^{\mathrm{nd}}`$ order, and the blue fibre cable, which we chose for its transmission at the desired wavelengths. The selected grating allowed us to observe in the wavelength range of $`\mathrm{\Delta }\lambda =41205600`$ Å, which is optimal for discerning some of the most prominent absorption features of the faint early–type galaxies including molecular G-band, H$`\gamma `$, H$`\beta `$, Mg<sub>2</sub>, and Fe$`\lambda 5350`$. With this setup we achieved the dispersion of $`0.705`$ Å px<sup>-1</sup>, while our instrumental resolution of FWHM $`=1.91`$ Å allowed us to detect velocity dispersions down to 35 km s<sup>-1</sup>. This is assuming that we can measure velocity dispersions up to 30 percent better than the instrumental resolution for galaxies with SNR $`>15`$.
The HYDRA multi-fibre spectrograph has $`100`$ fibres each with 3$`\mathrm{}`$ diameter. Thus, it is well suited for the detection of the faint early–type galaxies, which typically have a half light radius of $`2\mathrm{}`$ at the distance of Coma (GG03). We observed $`45`$ galaxies and $`45`$ adjacent sky spectra for each HYDRA setup. The galaxy sample was divided into 3 different groups depending on the exposure times to achieve SNR $`>15`$. We observed the brightest galaxies $`b_j<17.5`$ mag for a total integration time of 4 hours, objects with $`17.50<b_j<18.5`$ mag for 8 hours, and the faintest objects $`b_j>18.5`$ mag for 16 hours. In addition, we obtained spectra of template stars representative of the prevailing stellar-population of dEs, primarily G and K–type stars. The sample also included $`30`$ bright elliptical galaxies observed in previous studies in order to assess any systematic effects. Sample spectra of 4 galaxies with different luminosities and SNR are given in Figure 1. All the candidates had enough signal for radial velocity measurements from which we conclude that $`100`$ percent have the range of recession velocities of 4,000–10,000 km s<sup>-1</sup>, consistent with membership in the Coma cluster (Colless & Dunn 1996).
## 3 Data Reduction and Measurements
### 3.1 Basic Data Reduction
We used the basic and multi-fibre spectral reduction tasks within IRAF<sup>2</sup><sup>2</sup>2Image Reduction and Analysis Facility. Distributed by the National Optical Astronomy Observatories, which is operated by AURA (Association of Universities for Research in Astronomy, Inc) under cooperative agreement with the National Science Foundation. to reduce our observations. The data was first trimmed accordingly and corrected for the overscan region. We removed the cosmic rays with FIGARO within the STARLINK software program. For the remaining reduction we used the ‘dohydra’ task within IRAF, which provided: aperture identification; tracing of the apertures; flat field correction for the pixel to pixel sensitivity and the different throughput from fibre to fibre; removal of internal reflections within the spectrograph (scattered light); the wavelength calibration; and the sky subtraction. The largest rms for the wavelength calibration was 0.02 Å.
### 3.2 Velocity Dispersion and Radial Velocity Measurements
Velocity dispersions, $`\sigma `$, were measured from galaxy spectra with the software REDUCEME (Cardiel 1999). The program implements the Fourier quotient method (described by González-González 1993), originally introduced by Sargent et al. (1977), to measure the velocity dispersions. The Fourier quotient method assumes that observed galaxy spectra can be described as a convolution between the spectral characteristics of the stellar population, the broadening function, and the effective response function of the instrument. Using an initial guess of the velocity dispersion, the program calculates a broadening function described by the dispersion relation and models a galaxy spectrum as a convolution of the broadening function and an optimal stellar synthesis spectrum. This template for the galaxy is produced by combination of the different star templates. Via $`\chi ^2`$ minimization between the galaxy and the model spectra, the best value of the velocity dispersion for the broadening function is determined.
We tested the stability of the software for a range of values of the involved parameters. The only difference in $`\sigma `$ measurements occurred when we changed the wavelength range used to calculate the velocity dispersions and when we altered the tapering fraction. The measured velocity dispersion varied by $`10`$ percent with the choice of the spectral region, and up to $`20`$ percent for the faintest galaxies when varying the tapering fractions. We conducted a set of tests altering either the tapering fraction for a specific wavelength range, or the wavelength range for a specific tapering fraction. We found that the most stable solutions occurred for a tapering fraction of 0.25 and when we trimmed the spectra at the edges. We chose to trim the spectra by preserving the largest rest frame wavelength range possible: $`41505400`$ Å.
### 3.3 Uncertainty Measurements
We estimated the uncertainties in velocity dispersion measurements by the ‘boot–strapping’ technique implemented by J. Gorgas (private communication) within the REDUCEME software. Assuming that the noise in the galaxy spectrum is dominated by Poisson noise, the program uses Monte-Carlo realizations to produce simulated galaxy spectra with similar SNR to the original galaxy. For each galaxy we ran 50 simulations, after confirming that the results were the same for 100 simulations. The program then measures the velocity dispersion of each simulated galaxy and via $`\chi ^2`$ minimization estimates the error in the velocity dispersion measurements, which also accounts for the template mismatch.
We have checked internally and externally for systematical offsets in the velocity dispersion measurements in our sample of the core faint early–type galaxies. The external check consisted of comparing our measurements with those comparable in SNR and resolution in the literature, (MLKC02 and NFPS). Figure 2 shows the comparison of $`\sigma `$ measurements and includes the uncertainty in the $`\sigma /\sigma _{Lit}`$. This plot shows no systematic offsets between our and the two literature samples. Furthermore, the median of the comparison, 1.01, and the rms scatter of the points, 0.10, indicate that our $`\sigma `$ measurements are in good agreement with those of both NFPS (open triangles) and MLKC02 (solid triangles). We also note that some galaxies seem to be inconsistent with the average value when taking their uncertainties into account. This would imply that the uncertaities in the $`\sigma `$ measurements are possibly underestimated. However, it is not possible to determine whether it is the literature, ours, or both uncertaities that are underestimated.
The internal check was done in the same way as the external, except that in this case $`\sigma `$ was measured using only half of their exposures in our own data. We found no inconsistencies in the $`\sigma `$ measurements nor their error measurements for our sample of faint early–type galaxies, but the error measurements for our bright ellipticals were, on average, underestimated by 30 percent. We have therefore increased the uncertainties of the bright Es by 30 percent.
In addition, our analysis excludes galaxies with SNR $`<15`$, since below this value the fit of the model spectra to the galaxy was uncertain. We also investigate the effects noise has on the velocity dispersion measurements. The following method closely resembles the analysis of J$`ø`$rgensen, Franx & Kj$`\ae `$rgaard (1995; hereafter JFK95). We selected a template star that best fits a typical galaxy with a high SNR, and broadened it by convolving it by Gaussians with $`\sigma `$ ranging from 35 to 100 km s<sup>-1</sup>. Different amounts of noise were added to each spectra so the SNR would yield values ranging 10–50. The unaltered spectra of the star was used as a template, while we measured the velocity dispersion of the broadened and lower SNR spectra. The final $`\sigma `$ was derived after 1000 simulations using the bootstraping method. Figure 3 shows the percentage difference between the ’galaxy’ spectra (template spectra broadened to 50 km s<sup>-1</sup>) and the observed $`\sigma `$, i.e. the same broadened galaxy with different amount of noise. All our simulated spectra have a slightly overestimated measurement of the velocity dispersion. The effect is larger for galaxies with a smaller velocity dispersion. For example, a galaxy with $`\sigma =35`$ km s<sup>-1</sup> and SNR of 15 has $`\sigma `$ overestimated by $`6`$ percent. JFK95 find that a galaxy with $`\sigma =65`$ km s<sup>-1</sup> is overestimated by $`4`$ percent, and a galaxy with $`\sigma =100`$ km s<sup>-1</sup> by $`1`$ percent which is in good agreement with our measurements. We chose not to correct for this systematic effect since it is significantly smaller than the uncertainties in the $`\sigma `$ measurements for majority of the galaxies in our sample. Implications of this effect on the $`L\sigma `$ relation are discussed in section 5.
## 4 Results
The $`\sigma `$ and radial velocity measurements for 87 early–type galaxies we observed in the Coma cluster are presented in Table 1. For the following analysis we only used 72 objects with SNR $`>15`$. We classify galaxies with $`M_R<22.17`$ mag as bright Es. Our sample includes 69 galaxies with $`M_R>22.17`$ mag classified here as faint early–type galaxies and 3 bright Es.We include other samples from the literature with $`\sigma `$ measurements and plot the galaxies in diagnostic diagrams to characterize their kinematic properties.
In Figure 4 (left) we show the $`L\sigma `$ relation including galaxies with $`\sigma `$ measurements from MLKC02, Hudson et al. (2001), EFAR (Colless et al. 2001), and our sample. The total number of galaxies is 167, where 24 galaxies are classified as bright Es. In case of multiple $`\sigma `$ measurements from the literature, we used a minimum variance weighted average of velocity dispersions and calculated the error in the weighted average. When those galaxies were in common with our sample, we used the same process but adopted the symbol used for this paper (triangles) in the figure. Table 1 lists the $`\sigma `$ values measured in our sample.
We have derived the photometry in $`M_R`$ from Gutiérrez et al. (2004) for the MLKC02 and our sample, while the other catalogues included their own photometry. We averaged the actual luminosities for galaxies with multiple photometry. The uncertainty in the average of the magnitudes was 0.015 mag, except for 3 galaxies, which we excluded from the sample since their uncertainty was large and most likely due to a catalogue mismatch. Out of our sample of 72 galaxies with $`\sigma `$ measurements, 17 were not in the Gutierrez list. In this case we have derived $`R`$ (Johnson filter) by using the $`b_j`$ magnitudes listed in the GMP catalogue. We used a least squares fit to objects that had both $`b_j`$ and $`R`$ and obtained the transformation $`R=(1.050\pm 0.030)b_j2.422\pm 0.529`$ between these two magnitudes.
The $`L\sigma `$ relation in Figure 4 (left) exhibits a curvature or a change of slope. To allow for a comparison with earlier studies we perform least-squares fits to the bright Es and faint early–type galaxies separately. We partition the bright Es from the other early–type galaxies with the dotted line at $`M_R>22.17`$ mag. The dashed line represents a most recent $`L\sigma ^n`$ from the literature (Forbes, & Ponman 1999), where $`\mathrm{log}\sigma =0.102M_B+0.243`$, corresponding to $`n=3.92`$. We obtained the ordinary least squares fit (OLS, as described by Feigelson, & Babu 1992) for all the galaxies excluding bright Es. The dash-dot line represents the fit which minimizes the residuals in $`M_R`$, the dash-dot-dot minimization in $`\mathrm{log}\sigma `$, and the solid line is the bisector line. The details of the fits are summarized in Table 2. Only the galaxy with a very large observational error from MLKC02 (see Figure 4) is excluded from the linear regression.
Another diagnostic diagram used for the characterization of the low luminosity early–type galaxies and a comparison of their properties to those of the bright Es is the $`C\sigma `$ relation. Figure 5 shows the Johnson B–R mag (from Trenham, unpublished data) vs. $`\mathrm{log}\sigma `$ for galaxies in our sample, including some galaxies with $`\sigma `$ measurements from the literature. In case of the galaxies in common with the literature, we performed a weighted average on the velocity dispersions and their errors. We plot only galaxies that have $`BR`$ measurements, except in the case of the Hudson et al. (2001) sample where we have derived the colours from the GMP catalogue. The transformation between the GMP $`br`$ colours and $`BR`$ of our sample was: $`BR=(1.128\pm 0.098)(br)_j0.535\pm 0.178`$. The derived $`C\sigma `$ relations minimizing residuals in either $`\mathrm{log}\sigma `$, B–R, or using the bisector for all the galaxies in the figure are found in Table 3. The right-hand panel includes galaxies from clusters observed by Faber et al. (1989) where we used the colour transformation $`(BR)=(BV)+0.71`$ from Fukugita, Shimasaku & Ichikawa (1995).
## 5 Discussion
### 5.1 L-$`\sigma `$ Relation for Faint Early–type Galaxies
Faber & Jackson (1976) showed that luminosity of bright Es correlates well with velocity dispersion ($`\sigma `$) for these galaxies. The $`L\sigma `$, or Faber-Jackson, relation can be expressed as $`L\sigma ^n`$, where $`n`$ was originally $`4`$ (Faber & Jackson 1976; Sargent et al. 1977; Schechter & Gunn 1979; Schechter 1980; Tonry & Davis 1981; Terlevich et al. 1981). Tonry (1981) was the first to note a slight change of slope in the $`L\sigma `$ relation suggesting that $`n4`$ for more luminous objects while $`n3`$ for fainter galaxies. This result was confirmed by Davies et al. (1983) who found $`n=4.2\pm 0.9`$ for galaxies brighter than $`M_B=20`$ ($`M_R<21.67`$) and $`n=2.4\pm 0.9`$ for those fainter than this magnitude, and Held et al. (1992) who found $`n=2.5`$ for dEs.
Unfortunately, the data samples of Tonry (1981), Davies et al. (1983) and Held et al. (1992) only included a dozen of the faint early–type galaxies. To further investigate the $`L\sigma `$ relation for a wide range of luminosities we present a sample of 143 galaxies with $`22<M_R<17.5`$ mag.
The $`L\sigma `$ relation (Figure 4) derived for this large sample exhibits a change of slope; the slope of faint early–type galaxies is shallower than that of bright ellipticals. Following the results of previous studies and including these in our data set we find that the value of $`n4`$ fits the bright E end of the diagram. In contrast to bright Es, we obtain $`L\sigma ^{2.01\pm 0.36}`$ for faint early–type galaxies (adopting the bisector fit). This relation spans a range of 4.5 magnitudes fainter than $`M_R=22.17`$ mag, the lower limit of bright Es, and is the largest sample of the faint early–type galaxies in a single cluster thus far. Our result derived for 143 galaxies is consistent with $`L\sigma ^{2.4\pm 0.9}`$ derived by Davies et al. (1983) for their 14 faint ellipticals and that of De Rijcke et al. (2004), and it is inconsistent with the standard Faber-Jackson relation. This raises intriguing questions concerning the physical processes responsible for the change of slope in the $`L\sigma `$ relation.
We note that our galaxies exhibit a small systematic offset depending on their SNR, as discussed in section 3.2. However, if corrected for this effect of $`6`$ percent at the lowest $`\sigma `$ and lowest SNR galaxies, the slope of the faint early-type galaxies would be even shallower than derived in this paper.
A feasible explanation of the different slope between bright Es and faint early–type galaxies may be due to the presence of other types of galaxies at lower luminosities. We used the photometry of Gutiérrez et al. (2004) to classify galaxies from the present combined sample. In cases where the bulge-to-total (B/T) luminosity ratios from Gutiérrez et al.’s data were unreliable, we used the NASA/IPAC Extragalactic Database (NED). We defined galaxies with $`B/T=1.0`$ as bulge-dominated, $`0.5<B/T<1.0`$ as bulge+single exponential component galaxies, and $`B/T<0.5`$ as single exponential component. This notation is to avoid confusion of calling galaxies with exponential light profiles strictly as disk galaxies, since dEs tend to have exponential profiles (see GG03). The distribution of these galaxies in the $`L\sigma `$ plot (Figure 4, right) indicates that there is a slight difference in the slopes for each type. Although a different $`L\sigma `$ relation can be derived for each galaxy type, the individual relations are still within $`3\sigma `$ of each other (Table 2). It is not surprising that a significant number of single exponential component galaxies appear to be present in this sample since this light profile best describes the dwarf elliptical galaxies. However, classifying these low-luminosity galaxies is difficult without resolved photometry and we also note discrepancies in classification depending on the literature source. Nonetheless, the slopes of all three galaxy types, single exponential component, bulge+exponential component and bulge-dominated galaxies, are still in disagreement with the FJ relation, but consistent with $`L\sigma ^{2.01\pm 0.36}`$ as derived earlier for the low-luminosity early–type galaxies.
Tonry (1981) attributed the change of slope in the $`L\sigma `$ relation to less luminous systems having significant rotation. Subsequently, Davies et al. (1983) investigated rotational properties of about a dozen faint Es with $`20.5<M_B<18`$ mag ($`22.17<M_R<19.67`$ mag), and showed that the faint Es rotate more rapidly than most of the bright Es. More recent studies of the low-luminosity early–type galaxies, i.e. Simien, & Prugniel (2001), Geha, Guhathakurta, & van der Marel (2003), and van Zee, Skillman, & Haynes (2004) indicate that some of these objects have a rotational component while others show little or no rotation.
We investigate whether rotational effects in the faint early–type galaxies can cause the change of slope in the $`L\sigma `$ relation using the following approach. We assume that there is a universal relation between the luminosity and the kinetic energy per unit mass (or KE) such as: $`LKE^2`$, which is common to both faint early–type galaxies and bright ellipticals. Since the kinetic energy per unit mass for a spheroid is (from Busarello, Longo & Feoli 1992):
$$\mathrm{KE}=\frac{1}{2}v^2=\frac{1}{2}V_{rot}^2+\frac{3}{2}\sigma ^2$$
the assumed universal relation between luminosity and kinetic energy per unit mass translates into a correlation between luminosity, $`\sigma `$, and the anisotropy parameter ($`V_{rot}/\sigma `$):
$$L=a\sigma ^4(\frac{V_{rot}^2}{\sigma ^2}+3)^2$$
For bright ellipticals, the existence of the Faber-Jackson relation ($`L=b\sigma ^4`$) implies that $`V_{rot}=0`$, as it is indeed the case since bright ellipticals are anisotropy-supported stellar systems and show no or little rotation. For faint early–type galaxies, however, the existence of a relation such as $`L=c\sigma ^2`$, implies that there is a systematic increase in the anisotropy parameter (i.e., in the amount of rotation) as the velocity dispersion (or luminosity) decreases:
$$\frac{V_{rot}^2}{\sigma ^2}=\sqrt{\frac{c}{b}}\frac{1}{\sigma }3$$
This expression yields the amount of rotation a faint early–type galaxy would need to have in order to follow the same relation between luminosity and $`v^2`$ followed by bright Es. Under these assumptions, the observed change of slope in the $`L\sigma `$ diagram would simply be the result of not including the rotational component in the kinematic energy of faint early–type systems.
Using the expression above we can calculate the expected $`V_{rot}/\sigma `$ for galaxies at different luminosities and velocity dispersions, and compare them to the observed values. We show this comparison through a set of graphs (see Figure 6). The left panel shows $`V_{rot}/\sigma `$ vs. $`M_R`$, and the right panel $`V_{rot}/\sigma `$ vs. $`\sigma `$. The solid line in both panels represents the predicted value for $`V_{rot}/\sigma `$, while the points are the most recent data from the literature (Davies et al. 1983; Simien & Prugniel 2002; Pedraz et al. 2002; Geha, Guhathakurta & van der Marel 2003; van Zee, Skillman, & Haynes 2004).
According to the predicted relation between the $`V_{rot}/\sigma `$ and $`M_R`$, $`V_{rot}/\sigma `$ would have to increase steadily toward the faint end (left panel of Figure 6). For example, a faint early–type galaxy with $`M_R=18`$ mag and $`\sigma =39`$ km s<sup>-1</sup> would have $`V_{rot}/\sigma =3.2`$, or $`V_{rot}=123`$ km s<sup>-1</sup>. This value seems unreasonably large when compared to observed $`V_{rot}/\sigma `$ values from Geha, Guhathakurta & van der Marel (2003), which range from as low as 0.01 to $`0.5`$. The discrepancy is worse for the faintest galaxies. We conclude that the predicted $`V_{rot}/\sigma `$ is inconsistent with observations of early–type galaxies fainter than $`M_R=20.5`$ mag. Therefore, it is implausible that rotation is solely responsible for the difference in the $`L\sigma `$ slope between faint early–type and bright elliptical galaxies.
An independent confirmation of $`L\sigma ^2`$ has recently been provided by De Rijcke et al. (2004). They investigate how well different galaxy formation scenarios reproduce this slope difference between the bright ellipticals and bulges of spirals and dEs. The semi-analytical models which include quiescent star formation, post-merger star-bursts and gas-loss triggered by supernova winds seem to describe this effect well.
Next, we address the scatter around the $`L\sigma `$ relation for the faint early–type galaxies. The rms scatter of the bisector line is 0.52 mag, where 0.22 mag can be attributed to the observational errors, assuming $`\delta M_R=0.1`$ mag. The scatter in the $`L\sigma `$ relation of faint early–type galaxies is, therefore, intrinsic. A possible reason for this scatter could be that the faint early–type galaxies in question are not yet relaxed. This could indicate that in their recent history these galaxies have encountered interactions with other galaxies in the cluster. It is also possible that the age of the galaxy could have an effect on the scatter, as investigated by Forbes, & Ponman (1999). We do not have direct age measurements for these objects yet. However, in the following section we plot the colour$`\sigma `$ relation for our galaxies. Under assumptions described in the following section we are able to investigate the effects age and metallicity have on this relation.
In conclusion, faint early–type galaxies follow a well-defined $`L\sigma ^{2.01\pm 0.36}`$ relation. This relation is distinct from the traditional Faber-Jackson relation defined for bright E galaxies and might indicate that bright Es should no longer be viewed as canonical early–type galaxies. We also conclude that rotation in the faint early–type galaxies is not responsible for the change in the slope relative to that derived for bright Es.
### 5.2 Colour-$`\sigma `$ Relation
The colour-magnitude relation (CMR) is a well established relation for early–type galaxies. It is characterized by the more luminous galaxies displaying redder colours. This relation was investigated by many authors in the past who determined that the slope of the relation arises because the more massive galaxies are redder and more metal rich than the less massive galaxies (e.g., Terlevich, Caldwell, & Bower 2001, and references therein). Caldwell (1983) and Prugniel et al. (1993) found that faint early–type galaxies roughly follow the CMR for bright Es. A similar, distance-independent, relation is the correlation between colour and $`\sigma `$, $`C\sigma `$.
In Figure 5, we show the $`C\sigma `$ relation for faint early–type galaxies and bright ellipticals. All of the galaxies seem to follow the same relation, although we note the lack of bright Es in the plot containing only Coma galaxies (left panel). We confirmed the uniformity of the $`C\sigma `$ relation by checking our result with the U–V colours of Terlevich, Caldwell, & Bower (2001). The right panel includes galaxies from different clusters (Faber et al. 1989) and indicates that both faint and bright early–type galaxies follow the same $`C\sigma `$ relation.
Since luminosity and $`\sigma `$ are related, the $`C\sigma `$ relation is equivalent to the CMR, which in turn suggests a more fundamental relation between galaxy metallicity and mass. Although colours depend both on metallicity and age changes in the stellar populations, the evidence so far supports that metallicity changes are responsible for the slope of the CMR, while age differences contribute to the scatter observed around that relation. In the following analysis we assume that the same applies to the $`C\sigma `$ relation. Note however, that it is likely that both age and metallicity affect the slope and the scatter. Bernardi et al. (2005) show that galaxies with large velocity dispersion tend to be older. They also show that at a specific $`\sigma `$, galaxies have a wide range in both age and metallicity in a way that the older galaxies have smaller metallicities and the younger galaxies larger metallicities. Assuming that the intrinsic scatter in the $`C\sigma `$ relation is predominantly due to age, it is possible to constrain the age variations at a given formation epoch for faint early–type galaxies.
Bower, Lucey and Ellis (1992) (hereafter BLE92) showed that it is possible to determine minimum ages for galaxies with different formation scenarios by implementing evolutionary stellar population synthesis models and the intrinsic scatter in the CMR of these galaxies. We closely follow their method but use the intrinsic scatter derived from the $`C\sigma `$ relation for our galaxies instead. Since the uncertainty in the observed parameters is 0.024 mag, the intrinsic scatter is 0.067.
We used the evolutionary stellar population synthesis models of Bruzual, & Charlot (2003) to simulate a galaxy with an exponentially declining star-burst, $`\tau =1`$ Gyr, and a Chabrier IMF. In Figure 7, top panel, we show how the $`BR`$ colour of a galaxy with metalicitiy $`Z=1`$ or $`0.4`$ $`Z_{}`$ evolves with time. Hence, we find how the rate of change of colour varies with the age of the galaxy (middle panel). The colour change can be related to the intrinsic scatter in the colour and the range of epochs for major star-formation events by the relation:
$$\sigma _{(BR)}=\frac{d(BR)}{dt}\times \sigma _{SFE}$$
(1)
where $`\sigma _{(BR)}`$ is the scatter in $`BR`$ colour and $`\sigma _{SFE}`$ is the range in the star formation epoch (BLE92). This is in accordance with the assumption that the scatter in $`C\sigma `$ is only due to age variation, or in this case, to the scatter in the star formation epoch.
After finding the d(B-R)/dt at different ages of the galaxy, and using the intrinsic scatter of 0.067 in the $`C\sigma `$ relation, we derive values for the maximum range in the star formation epoch. In Figure 7 (bottom panel) we show the maximum range in the star formation epochs as a function of the formation time, constrained by the intrinsic scatter in $`C\sigma `$ for our sample. If the average age of our galaxy sample is 10 Gyr old, for example, the maximum range in its star formation epoch will be $`3`$ Gyr. However, in order to determine the upper limit on variations in the ages of galaxies we must take into account that the scatter in the SF epoch will also depend on how the galaxies were formed. Different formation scenarios can be parametrised by a parameter $`\beta `$, which describes the degree of coordination of galaxies during their formation. BLE92 define $`\beta `$ as ‘the ratio of the spread in formation time to the characteristic time-scale at formation,’ where the galaxy formation times range as $`\beta (t_Ht_F)`$ ($`t_H`$ is the age of the universe, and $`t_F`$ is the time of formation of the galaxy). For example, $`\beta =1`$ for uncoordinated galaxy formation, and $`\beta =0.1`$ for strong coordination. Assuming solar metallicity, the minimum average age for the faint early type galaxies in our sample with $`\beta =0.1`$ is $`6`$ Gyr (Figure 7, middle panel) with a scatter in the star formation epoch of $`\pm 1`$ Gyr (where $`t_H=13.7`$ Gyr, $`\mathrm{\Omega }_M=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$). As a reference, a galaxy that is 6 Gyr old must have formed at redshift $`z0.7`$. Note that we can only put a lower limit on the ages of galaxies and an upper limit on the scatter in their SF epoch, provided that we know the level of coordination of galaxies during their formation.
Although we initially assumed that the scatter of the $`C\sigma `$ relation is due to an age spread around the formation epoch in a single burst, secondary bursts of star formation will also contribute to the scatter. In fact, observations of galaxies in clusters at redshifts $`z0.5`$ point to a possibility of secondary star-bursts (Butcher & Oemler 1978; Butcher & Oemler 1984) in what may become today’s population of faint early–type galaxies in clusters. By modeling the secondary star-bursts, we can place upper limits on the star-burst strengths.
Assuming that all galaxies formed at $`t>10`$ Gyrs and had a secondary starburst $`5`$ Gyrs ago, we again use the models of Bruzual, & Charlot (2003) with two exponentially declining bursts to make this constraint. For this purpose we also follow the discussion of BLE92. Using the scatter in the $`BR`$ colour for various uniformly distributed burst strengths, solar metallicity and a Chabrier IMF, we find the typical rms burst strength $`r_{typ}=\mathrm{\Delta }(BR)/0.17`$ (the ratio between the stellar mass of the secondary burst to that of the initial burst). Our observed scatter of $`(BR)_{rms}=0.07`$ places an upper limit on the secondary burst of 40 percent by stellar mass of the first burst. Unfortunately, this constraint is not very robust since it strongly depends on the assumptions of the age of the starburst and the modeling components.
In conclusion, we have shown that there is a well defined relation between colour and $`\sigma `$ for faint early–type systems. Assuming that metallicity changes are responsible for the slope of this correlation while age variations are the main contributor to the scatter, it is possible to constrain the age range of major star formation events for a given formation epoch. However, it is difficult to decouple the effects of age and metallicity using colours. In future papers we will study the detailed stellar population properties of faint early–type galaxies, both age and metallicity, using line strength indices and stellar population synthesis models. Furthermore, we will test if the galaxies in the centre of the cluster are more metal-rich than those in the outskirts, since this is predicted by the galaxy harassment model (Moore, Lake, & Katz 1998).
## 6 Summary
We present velocity dispersion measurements for 69 faint early–type galaxies in the centre of the Coma cluster with $`22.17<M_R<17.5`$ mag. We derive the $`L\sigma `$ relation for faint early–type galaxies as $`L\sigma ^{2.01\pm 0.36}`$, which differs from the Faber-Jackson relation, $`L\sigma ^4`$, defined for bright ellipticals. Rotation in these objects is investigated as a possible cause for the difference in the slope. Although rotation may contribute to the scatter in this relation, it is not the main cause for the different slope derived for these galaxies.
We also investigate whether the slope change is due to the presence of different classes of early–type galaxies in our sample. Although our sample includes bulge-dominated, bulge+single exponential component and a few single exponential component galaxies, all three types essentially follow the same relation.
We find that faint early–type galaxies follow a well-defined $`C\sigma `$ relation. By assuming that this relation is mostly driven by an increased metallicity with increasing galaxy mass, while the scatter reflects age differences, we investigated how we can constrain either the ages, the range of star formation epoch, or the strength of secondary bursts for the faint early–type galaxies for various galaxy formation scenarios.
## Acknowledgments
We thank our referee, Michael Hudson, for his suggestions which have significantly improved a previous version of this paper. We also thank Alister Graham for many insightful discussions, Javier Gorgas, Nicolás Cardiel and Patricia Sánchez-Blázquez for help and instructions on velocity dispersion measurements; Marla Geha for an independent check of velocity dispersions measurements; Neil Trentham for providing the B-R colours for our galaxies; Eric McKenzie for discussions and help with stellar synthesis models; Nicolas Gruel for help and use of his catalogue cross-correlation program; and Nelson Caldwell for providing the U-V coluor catalog. R. G. acknowledges funding from the archive HST proposal HF01092.01-97A. R.G. also thanks the Yale TAC for generous allocation of time on the WIYN telescope.
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# Solvable relativistic quantum dots with vibrational spectra
## 1 Klein-Gordon equation
### 1.1 Pseudo-Hermitian Feshbach-Villars Hamiltonian
As long as the most common relativistic Klein-Gordon (KG) operators are partial differential operators of the second order with respect to time, the time evolution of the wave functions $`\mathrm{\Psi }^{(KG)}(x,t)`$ must be studied together with their first time derivatives $`i_t\mathrm{\Psi }^{(KG)}(x,t)`$. After the routine Fourier transformation we arrive at the Feshbach-Villars (FV, ) non-Hermitian eigenvalue problem
$$\widehat{H}^{(FV)}|\psi =E|\psi ,\widehat{H}^{(FV)}=\left(\begin{array}{cc}0& \widehat{h}^{(KG)}\\ 1& 0\end{array}\right)$$
(1)
where the two wave-function components may be marked as $`|D\}`$ (=“down component”) and $`|U\}`$ (=“up component”). For the description of the bound states in one dimension the two-by-two partitioning in (1) allows us to extract $`|U\}=E|D\}`$ and to replace our Klein-Gordon equation by its reduced form
$$\widehat{h}^{(KG)}|D_n\}=\epsilon _n|D_n\},n=1,2,\mathrm{}$$
(2)
with squared energy $`E^2`$ abbreviated as $`\epsilon `$ and with the “large” Hilbert space $``$ of kets $`|\psi `$ reduced to the “smaller” Hilbert space $`_{(c)}`$ of the curly-ket “down” components $`|D_n\}`$ .
### 1.2 Biorthogonal bases
The “right” eigenkets $`|D_n\}`$ will not carry all information about $`\widehat{h}^{(KG)}`$ whenever $`[\widehat{h}^{(KG)}][\widehat{h}^{(KG)}]^{}`$. Then, the parallel Schrödinger-type problem generates different eigenkets marked by the double curly ket symbol. The latter sequence may be re-read as the left eigenvectors of our original operator $`\widehat{h}^{(KG)}`$, related to the same (by assumption, real) eigenvalues $`\epsilon _n\kappa _n^2`$,
$$\{\{L_n|\widehat{h}^{(KG)}=\kappa _n^2\{\{L_n|,n=1,2,\mathrm{}.$$
(3)
It is well known that the set of the bras $`\{\{L_n|`$ and kets $`|D_n\}`$ is bi-orthogonal ,
$$\{\{L_m|D_n\}=\{\begin{array}{c}0\hfill \\ \varrho _n0\hfill \end{array}\mathrm{for}\{\begin{array}{c}mn,\hfill \\ m=n,\hfill \end{array}$$
and that it forms, usually, a basis in the infinite-dimensional Hilbert space $`_{(c)}`$. Then, we may decompose the unit operator and/or derive the bi-orthogonal spectral representation of the Hamiltonian in $`_{(c)}`$,
$$I_{(c)}=\underset{n=1}{\overset{\mathrm{}}{}}|D_n\}\frac{1}{\varrho _n}\{\{L_n|,\widehat{h}^{(KG)}=\underset{n=1}{\overset{\mathrm{}}{}}|D_n\}\frac{\kappa _n^2}{\varrho _n}\{\{L_n|.$$
(4)
The overlaps $`\varrho _n`$ need not be all of the same sign.
## 2 Relativistic observables
### 2.1 $`\mathrm{\Theta }`$quasi-Hermiticity
In the space $`=_{(c)}_{(c)}`$ of the eigenstates of $`H^{(FV)}`$ we have to consider the pair of conjugate equations
$$\widehat{H}^{(FV)}|n^{(\pm )}=\pm \kappa _n|n^{(\pm )},n^{(\pm )}|\widehat{H}^{(FV)}=\pm \kappa _nn^{(\pm )}|.$$
(5)
Both the left and right eigenstates have the two-component structure,
$$|m^{(\pm )}=\left(\begin{array}{c}|L_m\}\}\\ \pm \kappa _m|L_m\}\}\end{array}\right),|n^{(\pm )}=(\begin{array}{c}\pm \kappa _n|D_n\}\\ |D_n\}\end{array})$$
and form the bi-orthogonal set in the “bigger” space $``$,
$$m^{(\nu )}|n^{(\nu )}=\delta _{mn}\delta _{\nu \nu }\mu _m^{(\nu )},\mu _m^{(\pm )}=\pm 2\kappa _m\varrho _m,\nu ,\nu =\pm 1.$$
It is expected to be complete and useful,
$$I=\underset{\tau =\pm 1}{}\underset{n=1}{\overset{\mathrm{}}{}}|n^{(\tau )}\frac{1}{\mu _n^{(\tau )}}n^{(\tau )}|,$$
(6)
$$H^{(FV)}=\underset{\tau =\pm 1}{}\underset{n=1}{\overset{\mathrm{}}{}}|n^{(\tau )}\frac{\tau \kappa _n}{\mu _n^{(\tau )}}n^{(\tau )}|=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(|n^{(+)}n^{(+)}|)+(|n^{()}n^{()}|)}{2\varrho _n}.$$
Let us now assume that at a given $`\widehat{H}^{(FV)}`$, equation
$$\left[\widehat{H}^{(FV)}\right]^{}=\eta \widehat{H}^{(FV)}\eta ^1$$
(7)
possesses a positive and Hermitian solution $`\eta _+=\mathrm{\Theta }>0`$. Such an operator may play the role of a metric and induces the following specific scalar product in $``$,
$$\left(|\psi _1|\psi _2\right)=\psi _1|\mathrm{\Theta }|\psi _2=\psi _1|\psi _2_{(physical)},|\psi _1,|\psi _2.$$
(8)
This product generates the norm, $`\psi =\sqrt{\psi |\psi _{(physical)}}`$. In terms of the later product and metric we may call all the operators $`A`$ with the property $`A^{}=\mathrm{\Theta }A\mathrm{\Theta }^1`$ quasi-Hermitian and treat them as observables (see for a deeper outline of some more sophisticated mathematical details). Indeed, we have
$$\left(|\psi _1|A\psi _2\right)\left(|A\psi _1|\psi _2\right)$$
(9)
so that the probabilistic expectation values $`\psi |A|\psi _{(physical)}`$ are mathematically unambigously defined.
### 2.2 Explicit constructions of the metric $`\mathrm{\Theta }`$
Let us assume non-Hermiticity of the type $`\widehat{h}^{(KG)}[\widehat{h}^{(KG)}]^{}=𝒫\widehat{h}^{(KG)}𝒫`$ in the smaller space $`_{(c)}`$ (here, $`𝒫`$ is operator of parity). Then, a consistent physical meaning may still be assigned to all the relativistic bound states, provided only that in the bigger space $``$ we find a suitable physical metric $`\mathrm{\Theta }`$. For this purpose we may employ the ansatz
$$\mathrm{\Theta }=\underset{\tau ,\tau =\pm 1}{}\underset{m,n=1}{\overset{\mathrm{}}{}}|n^{(\tau )}M_{nm}^{(\tau \tau )}m^{(\tau )}|,$$
the backward insertion of which in (7) gives the condition
$$\tau \kappa _nM_{nm}^{(\tau \tau )}=M_{nm}^{(\tau \tau )}\tau \kappa _m$$
with the set of solutions $`M_{nm}^{(\tau \tau )}=\omega _n^{(\tau )}\delta _{nm}\delta _{(\tau \tau )}`$ numbered by the free parameters $`\stackrel{}{\omega }^{(\pm )}`$. The Hermiticity and positivity constraints restrict the freedom of the choice of both the optional sequences $`\stackrel{}{\omega }^{(\pm )}`$ to the real and positive values, $`\omega _n^{(\pm )}>0`$. Vice versa, any choice of the latter two sequences defines an eligible operator of the metric
$$\mathrm{\Theta }=\mathrm{\Theta }_{\stackrel{}{\omega }^{(\pm )}}=\underset{\tau =\pm 1}{}\underset{n=1}{\overset{\mathrm{}}{}}|n^{(\tau )}\omega ^{(\tau )}_nn^{(\tau )}|.$$
(10)
Its inverse
$$\mathrm{\Theta }^1=\underset{\tau =\pm 1}{}\underset{n=1}{\overset{\mathrm{}}{}}|n^{(\tau )}\frac{1}{\omega _n^{(\tau )}|\mu _n^{(\tau )}|^2}n^{(\tau )}|$$
(11)
is similar. In terms of the metric $`\mathrm{\Theta }`$, the formal bound-state wave functions re-acquire the standard probabilistic interpretation.
## 3 Models with complex point interactions
In a way inspired by the success of several non-relativistic studies of $`𝒫𝒯`$symmetric models with point interactions and by the encouraging experience we made in our paper we shall combine the infinitely deep square-well real part of the potential \[$`V(x)=\mathrm{}`$ for all $`x(1,1)`$\] with the following purely imaginary delta-function formula for its remaining part,
$$V(x)=\underset{\mathrm{}=1}{\overset{}{}}\left[i\xi _{\mathrm{}}\delta \left(xa_{\mathrm{}}\right)i\xi _{\mathrm{}}\delta \left(x+a_{\mathrm{}}\right)\right],x(1,1),$$
(12)
at real couplings $`\xi _{\mathrm{}}`$ and ordered points $`0<a_1<a_2<\mathrm{}<a_1<a_{}<1`$.
### 3.1 Wave functions
The key advantage of our $`V(x)`$ in (12) is that the $`𝒫𝒯`$symmetrically normalized coordinate representants $`\psi (x)=\psi ^{}(x)`$ of $`|D\}`$ in eq. (2) remains piecewise trigonometric. At each real and positive bound-state energy $`\epsilon =\kappa ^2`$ we shall have
$$\psi (x)=\{\begin{array}{cc}\psi _L^{()}(x)=(\alpha _{}i\beta _{})\mathrm{sin}\kappa (1+x),\hfill & x(1,a_{}),\hfill \\ \multicolumn{2}{c}{\psi _L^{(\mathrm{})}(x)=(\alpha _{\mathrm{}}i\beta _{\mathrm{}})\mathrm{sin}\kappa (a_{\mathrm{}+1}+x)+(\gamma _{\mathrm{}}i\delta _{\mathrm{}})\mathrm{cos}\kappa (a_{\mathrm{}+1}+x),}\\ & x(a_{\mathrm{}+1},a_{\mathrm{}}),\hfill \\ \psi _C^{(0)}(x)=\mu \mathrm{cos}\kappa x+i\nu \mathrm{sin}\kappa x,\hfill & x(a_1,a_1),\hfill \\ \multicolumn{2}{c}{\psi _R^{(\mathrm{})}(x)=(\alpha _{\mathrm{}}+i\beta _{\mathrm{}})\mathrm{sin}\kappa (a_{\mathrm{}+1}x)+(\gamma _{\mathrm{}}+i\delta _{\mathrm{}})\mathrm{cos}\kappa (a_{\mathrm{}+1}x),}\\ & x(a_{\mathrm{}},a_{\mathrm{}+1}),\hfill \\ \psi _R^{()}(x)=(\alpha _{}+i\beta _{})\mathrm{sin}\kappa (1x),\hfill & x(a_{},1),1\mathrm{}<.\hfill \end{array}$$
(13)
Its differentiation as well as continuity conditions
$$\begin{array}{c}\psi _L^{(\mathrm{}1)}(a_{\mathrm{}})=\psi _L^{(\mathrm{})}(a_{\mathrm{}}),\mathrm{}=,1,\mathrm{},2,\\ \psi _C^{(0)}(a_1)=\psi _L^{(1)}(a_1),\psi _R^{(1)}(a_1)=\psi _C^{(0)}(a_1),\\ \psi _R^{(\mathrm{}+1)}(a_{\mathrm{}+1})=\psi _R^{(\mathrm{})}(a_{\mathrm{}+1}),\mathrm{}=1,2,\mathrm{},1,\end{array}$$
(14)
enter the definition of the action of the delta functions,
$$\begin{array}{c}\left[\psi _L^{(\mathrm{}1)}(a_{\mathrm{}})\right]^{}\left[\psi _L^{(\mathrm{})}(a_{\mathrm{}})\right]^{}=i\xi _{\mathrm{}}\psi _L^{(\mathrm{})}(a_{\mathrm{}}),\mathrm{}=,1,\mathrm{},2,\\ \left[\psi _C^{(0)}(a_1)\right]^{}\left[\psi _L^{(1)}(a_1)\right]^{}=i\xi _1\psi _C^{(0)}(a_1),\\ \left[\psi _R^{(1)}(a_1)\right]^{}\left[\psi _C^{(0)}(a_1)\right]^{}=i\xi _1\psi _C^{(0)}(a_1),\\ \left[\psi _R^{(\mathrm{}+1)}(a_{\mathrm{}+1})\right]^{}\left[\psi _R^{(\mathrm{})}(a_{\mathrm{}+1})\right]^{}=i\xi _{\mathrm{}+1}\psi _R^{(\mathrm{})}(a_{\mathrm{}+1}),\mathrm{}=1,2,\mathrm{},1,\end{array}$$
(15)
After the insertion of the ansatz (13), the set of formulae (14) and (15) may be read as a homogeneous linear algebraic system of $`4`$ equations for the $`4`$ unknown wave-function coefficients $`\alpha _{},\beta _{},\mathrm{},\nu `$. The secular determinant $`𝒟(\kappa )`$ of this system must vanish so that the not too complicated transcendental equation
$$𝒟(\kappa )=0$$
(16)
determines finally the set of the bound-state roots $`\kappa =\kappa _n`$ at $`n=1,2,\mathrm{}`$.
### 3.2 Energies at the simplest choice of $`=1`$
At $`=1`$, potential (12) degenerates to the most elementary double-well model with the single coupling $`\xi _1=\xi `$ and one displacement $`a_1=a`$ . Out of the related eight real constraints (14) and (15) only four are independent and define the four real coefficients $`\alpha _1=\alpha `$, $`\beta _1=\beta `$ and $`\mu `$ and $`\nu `$ as an eigenvector of a four-by-four matrix with the secular determinant
$$𝒟(\kappa )=\frac{1}{2}\left\{\mathrm{sin}\mathrm{\hspace{0.17em}2}\kappa +\frac{\xi ^2}{\kappa ^2}\mathrm{sin}\mathrm{\hspace{0.17em}2}\kappa a\mathrm{sin}^2[\kappa (1a)]\right\}.$$
(17)
Numerically, the first term would give us the well-known square-well spectrum at $`\xi =0`$, the completeness of which is controlled by the Sturm-Liouville oscillation theory . As long as all the roots $`\kappa _n=\kappa _n(\xi )`$ are smooth and real functions of $`\xi `$ at the smallest couplings, $`\kappa _n(\xi )n\pi /2+𝒪(\xi ^2/n)`$, our explicit construction confirms the general mathematical prediction that the influence of the non-Hermiticity will be most pronounced at the lowest part of the spectrum.
### 3.3 The next choice of $`=2`$
In the quadruple-well potential (12) with $`=2`$ we may shorten $`a_1=a,a_2=b`$ and drop the two redundant subscripts in $`\gamma _1=\gamma ,\delta _1=\delta `$. In the eight-dimensional matrix of the system the elimination of four unknowns is either trivial \[$`\gamma =\alpha _2\mathrm{sin}\kappa (1b)`$, $`\delta =\beta _2\mathrm{sin}\kappa (1b)`$\] or easy \[$`\alpha _1=\alpha _1(\alpha _2,\beta _2)`$, $`\beta _1=\beta _1(\alpha _2,\beta _2)`$\]. We end up with a four-by-four matrix problem and with the secular determinant
$$𝒟(\kappa )=𝒟_{(0)}(\kappa )+𝒟_{(\xi _1)}(\kappa )+𝒟_{(\xi _2)}(\kappa )+𝒟_{(\xi _1\xi _2)}(\kappa ),$$
(18)
$$𝒟_{(0)}(\kappa )=\frac{1}{2}\mathrm{sin}\mathrm{\hspace{0.17em}2}\kappa ,𝒟_{(\xi _j)}(\kappa )=\frac{\xi _j^2}{2\kappa ^2}\mathrm{sin}\mathrm{\hspace{0.17em}2}\kappa a_j\mathrm{sin}^2[\kappa (1a_j)],j=1,2,$$
$$𝒟_{(\xi _1\xi _2)}(\kappa )=\left\{\frac{\xi _1\xi _2}{\kappa ^2}\mathrm{sin}\mathrm{\hspace{0.17em}2}\kappa a+\frac{\xi _1^2\xi _2^2}{\kappa ^4}\mathrm{sin}^2[\kappa (ba)]\right\}\mathrm{sin}^2[\kappa (1b)].$$
This secular determinant correctly degenerates to the previous $`=1`$ formula in both the independent limits of $`\xi _10`$ and $`\xi _20`$.
### 3.4 Simplifications at the rational $`a_j`$
Let us return to the secular eq. (17) with $`=1`$ and choose $`a=1/2`$ . This leads to a factorization of $`𝒟(\kappa )`$ and to the pair of the eigenvalue conditions
$$\mathrm{cos}\kappa _{2m1}=\frac{\xi ^2}{\xi ^24\kappa _{2m1}^2},\mathrm{sin}\kappa _{2m}=0,m=1,2,\mathrm{}$$
(19)
with the second series of equations being exactly solvable, $`\kappa _{2m}=m\pi `$.
At the next choice of $`a=1/3`$ we factorize eq. (17) in the similar manner and get the series of the $`\xi `$dependent roots specified by the implicit definitions
$$\mathrm{cos}\frac{4}{3}\kappa _p=\frac{\xi ^2+2\kappa _p^2}{\xi ^24\kappa _p^2},p=1,2,4,5,7,8,10,\mathrm{}$$
(20)
complemented by the closed formula for all the skipped roots of the factor $`\mathrm{sin}2\kappa /3`$ which remain $`\xi `$independent and read $`\kappa _{3m}=3m\pi /2`$ with $`m=1,2,\mathrm{}`$. The regularity of such a pattern of the $`\xi `$independent roots is easily prolonged to the decreasing sequence of $`a`$ with $`\kappa _{4m}=2m\pi `$ at $`a=1/4`$ and all $`m=1,2,\mathrm{}`$, etc.
The less elementary composite choice of $`a=2/3`$ may be observed to give the same factor as at $`a=1/3`$ and, hence, the same $`\xi `$independent series of the roots $`\kappa _{3m}=3m\pi /2`$ with $`m=1,2,\mathrm{}`$. The implicit formula for the remaining roots is a slightly more complicated quadratic equation in the trigonometric unknown $`X=\mathrm{cos}2\kappa /3`$,
$$\left(4\kappa ^2\xi ^2\right)X^2+\xi ^2X\kappa ^2=0.$$
(21)
Its trigonometric part $`X`$ may be eliminated in the form resembling eq. (20).
One of the important consequences of the existence of the elementary formulae for the rational $`a`$ is that they allow us to perform an elementary analysis of the qualitative features of the $`n`$th root $`\kappa _n`$ during the growth of the strength $`\xi `$ of the non-Hermiticity. During such an analysis one discovers that these levels are either “robust” (marked by a superscript, $`\kappa _n^{(R)}`$, and remaining real for all $`\xi `$) or “fragile” (such a $`\kappa _n^{(F)}`$ will merge with another $`\kappa _m^{(F)}`$ at a “critical” $`\xi _{n,m}^{(C)}`$ while the pair will complexify beyind this “exceptional” point). For illustration let us display this pattern in the three simplest spectra,
$$\kappa _1^{(F)},\kappa _2^{(R)},\kappa _3^{(F)},\kappa _4^{(R)},\kappa _5^{(F)},\kappa _6^{(R)},\mathrm{},a=1/2$$
$$\kappa _1^{(F)},\kappa _2^{(F)},\kappa _3^{(R)},\kappa _4^{(F)},\kappa _5^{(F)},\kappa _6^{(R)},\mathrm{},a=1/3$$
$$\kappa _1^{(F)},\kappa _2^{(F)},\kappa _3^{(R)},\kappa _4^{(R)},\kappa _5^{(R)},\kappa _6^{(F)},\kappa _7^{(F)},\kappa _8^{(R)},\kappa _9^{(R)},\kappa _{10}^{(R)},\kappa _{11}^{(F)},\mathrm{},a=1/4.$$
Acknowledgements. Work partially supported by AS CR (GA grant Nr. A1048302 and IRP AV0Z10480505).
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# Darwin-Lagrangian Analysis for the Interaction of a Point Charge and a Magnet: Considerations Related to the Controversy Regarding the Aharonov-Bohm and Aharonov-Casher Phase Shifts
## I Introduction
The interaction of a point charge and a magnet is a complicated and controversial problem of electromagnetism. The problem is ignored by the classical physics textbooks and is discussed in the research literature in connection with the Shockley-James paradox,S-J and in connection with the Aharonov-BohmAB and Aharonov-CasherAC phase shifts for particles. The problem in understanding arises because the interaction involves relativistic terms of order $`1/c^2`$ (where $`c`$ is the speed of light in vacuum) which are not nearly so familiar as nonrelativistic mechanics. Writing regarding the interaction of a point charge and a magnet in 1968, Coleman and Van Vleck remarked in an oft-cited article,CVV ”Unfortunately, the equations which we have obtained are singularly resistant to a simple physical interpretation in terms of particles exchanging forces; …” However, despite the complications and in line with the controversy, the problem is an important one which reflects back on our understanding of classical electromagnetism and on the connections between classical and quantum physics.
## II The Problem and the Controversy
There are no electric or magnetic fields outside a long neutral solenoid or toroid when the currents are maintained constant. Therefore when a charged particle passes a long solenoid or a toroid, there are no electric or magnetic fields at the position of the passing charge due to the unperturbed charge and current densities of the magnet. On the other hand, there are clearly electric and magnetic fields due to the passing charge at the position of the magnet. The electric fields of the passing charge will cause accelerations of the charges which carry the currents which create the flux of the magnet. Also, the magnetic fields of the passing charge will cause a net Lorentz force on the magnet. Thus far the description would be approved by all physicists. However, the response of the multiparticle magnet seems so complicated that no one has calculated the magnet’s response in detail.
Since it does not seem possible at present to carry out a complete multiparticle calculation starting from accepted theory, we are left with suggestive partial calculations and hence with competing points of view depending upon which aspects of the partial calculations are favored. At present, there are two competing interpretations for the behavior of a magnet and a passing point charge.
#### II.0.1 The No-Velocity-Change Point of View
The supporters of the quantum topological viewAB AC APV Vaidman P of the Aharonov-Bohm phase shift claim that there are no velocity changes for the interacting charged particle or the magnet. Indeed, the supporters of this view say that there are no significant changes in the charge or current densities in the magnet. Therefore the passing charge never experiences a Lorentz force and never changes velocity. Furthermore, although the magnet does indeed experience a net Lorentz force due to the magnetic field of the passing charge, nevertheless the electric field of the passing charge penetrates into the magnet giving a ”hidden momentum in magnets” whose change ”cancels” the net magnetic Lorentz force on the magnet so that the center of energy of the magnet is never disturbed. In this point of view, the electromagnetic fields of the passing charge may cause confusion behind the scenes inside the magnet, but there is no change in the magnet’s center of energy and there is no feedback signal sent to the passing charge which is causing the confusion in the magnet.nature
#### II.0.2 The Classical-Lag Point of View
The classical-lag point of viewLieb B2 B3 B4 B5 B6 B7 B8 takes a totally different perspective on the changes in the charge and current densities induced in the magnet. In this view, the induced densities lead to a Lorentz force back on the passing charge which is equal in magnitude and opposite in direction to the net magnetic Lorentz force which the magnetic field of the passing charge places on the magnet. The electric charges on the surface of the magnet screen the electric field of the passing charge out from the interior of the magnet, and therefore there is no significant change in the momentum of the electromagnetic fields. On the other hand, the magnetic field of the passing charge penetrates into the magnet, and it is the magnetic energy change associated with the overlapping magnetic fields which gives the magnitude of the energy change of the passing charge due to the back force. This view fits with what we know of the penetration of electric and magnetic velocity fields into ohmic conductors. In this scenario, we have explicit ideas concerning conservation of energy, linear momentum, and constant motion of the center of energy. We also have the validity of Newton’s third law for the net Lorentz forces between the magnet and the passing charge.
Both points of view predict the Aharonov-Bohm and Aharonov-Casher phase shifts. The no-velocity-change point of view claims that, in the light of their interpretation, the phase shifts represent completely new quantum topological effects occurring in the absence of classical forces, and there are no classical analogues. The classical-lag point of view claims that the phase shifts present classical velocity shifts analogous to those occurring when only one beam of light passes through a piece of glass before two coherent beams interfere. The conflict between the two points of view has existed for thirty years without ever being put to experimental test to determine whether or not there are velocity changes for the electrons passing through a toroid or past a long solenoid. The no-velocity-change point of view has been widely accepted because most physicists do not think of the possibility of induced charge and current densities in magnets; they do consider induced charge densities only in electrostatic situations. Furthermore, the proponents of the no-velocity-change point of view have declared that the lag point of view is impossible because i) the electromagnetic fields of the passing charge would not penetrate into a conductor surrounding a toroid or solenoid, and ii) the back electric field at the passing charge could not be of order $`1/c^2`$ and proportional to the magnetic flux of the magnet. The objection i) has been shown to be groundless.B5 Magnetic velocity fields do indeed penetrate into good conductors in exactly the required form which is completely different from the exponential skin-depth form taken by electromagnetic wave fields.B1999
The objection ii) is addressed in the present article. In 1968 Coleman and Van VleckCVV discussed the interaction of a stationary point charge and a magnet using the Darwin Lagrangian. We will be following their approach in the following analysis. We will discuss the interaction of a passing point charge and a magnetic moment where the magnetic moment is modeled as a classical hydrogen atom and where the electromagnetic interactions are carried to order $`1/c^2`$ by using the Darwin Lagrangian. This is a well-defined classical electromagnetic system which is relativistic through order $`1/c^2`$. In order to separate out the electrostatic effects (which are independent of the magnetic moment) from magnetic effects dependent upon the magnetic moment, we will sometimes average over atoms and anti-atoms with the same magnetic moment. We will describe the motion and check all the conservation laws. We will find that in this case the induced currents are important and that there are electric Lorentz forces back on the passing charge which indeed are of order $`1/c^2`$ and are proportional to the magnetic moment. There is also a displacement of the center of energy of the magnetic moment. This behavior contradicts the suggestions of the proponents of the no-velocity-change point of view.APV Vaidman Next we will discuss the passage to the limit of a multiparticle magnet. Finally, in this multiparticle limit, we discuss the conservation-law aspects which are mentioned above.
## III The Darwin Lagrangian and Electromagnetic Fields
The Darwin Lagrangian for particles of charge $`e_a`$, mass $`m_a`$, displacement $`𝐫_a`$, and velocity $`𝐯_a`$ is given byCVV Jackson
$`L`$ $`={\displaystyle \underset{a}{}}\left({\displaystyle \frac{1}{2}}m_a𝐯_a^2+{\displaystyle \frac{1}{8c^2}}m_a𝐯_a^4\right){\displaystyle \frac{1}{2}}{\displaystyle \underset{a}{}}{\displaystyle \underset{ba}{}}{\displaystyle \frac{e_ae_b}{r_{ab}}}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{a}{}}{\displaystyle \underset{ba}{}}{\displaystyle \frac{e_ae_b}{2c^2r_{ab}}}\left[𝐯_a𝐯_b+{\displaystyle \frac{(𝐯_a𝐫_{ab})(𝐯_b𝐫_{ab})}{r_{ab}^2}}\right]`$ (1)
where $`𝐫_{ab}=𝐫_a𝐫_b`$ and $`r_{ab}=|𝐫_a𝐫_b|`$. Lagrange’s equations of motion give a canonical momentum
$$𝐩_a^{canonical}=\frac{L}{𝐯_a}=m_a𝐯_a\left(1+\frac{𝐯_a^2}{2c^2}\right)+\underset{ba}{}\frac{e_ae_b}{2c^2r_{ab}}\left[𝐯_b+\frac{𝐫_{ab}(𝐫_{ab}𝐯_b)}{r_{ab}^2}\right]$$
(2)
and a time derivative
$`{\displaystyle \frac{d}{dt}}𝐩_a^{canonical}`$ $`={\displaystyle \frac{L}{𝐫_a}}={\displaystyle \underset{ba}{}}{\displaystyle \frac{e_ae_b𝐫_{ab}}{2c^2r_{ab}^3}}{\displaystyle \underset{ba}{}}{\displaystyle \frac{e_ae_b𝐫_{ab}}{2c^2r_{ab}^3}}\left[𝐯_a𝐯_b+{\displaystyle \frac{3(𝐯_a𝐫_{ab})(𝐯_b𝐫_{ab})}{r_{ab}^2}}\right]`$
$`+{\displaystyle \underset{ba}{}}{\displaystyle \frac{e_ae_b}{2c^2r_{ab}^3}}\left[𝐯_a(𝐯_b𝐫_{ab})+𝐯_b(𝐯_a𝐫_{ab})\right]`$ (3)
The Darwin Lagrangian accurately reflects the classical electromagnetic interaction of charged particles through order $`1/c^2`$. To lowest order in $`1/c^2,`$ the interaction among the charges is given by the Coulomb force and the nonrelativistic form of Newton’s second law $`𝐅=m𝐚`$. This $`0`$-order behavior can then be inserted back into the equations of motion to allow calculation of the higher-order corrections.
It is sometimes revealing to rewrite the Lagrangian equations of motion in terms of the mechanical momentum
$$𝐩_a=m_a𝐯_a[1+𝐯_a^2/(2c^2)]$$
(4)
Then Newton’s second law
$$d𝐩_a/dt=\frac{d}{dt}\{m_a𝐯_a[1+𝐯_a^2/(2c^2)]\}=e_a𝐄(𝐫_a,t)+e_a(𝐯_a/c)\times 𝐁(𝐫_a,t)$$
(5)
is obtained by carrying out the time derivative in the Darwin equations of motion (3) and recognizing the electric field asP-A
$$𝐄(𝐫_a,t)=\underset{ba}{}\left\{\frac{e_b𝐫_{ab}}{r_{ab}^3}\left[1+\frac{1}{2}\frac{𝐯_b^2}{c^2}\frac{3}{2}\frac{(𝐯_b𝐫_{ab})^2}{c^2r_{ab}^2}\right]\frac{e_b}{2c^2r_{ab}}\left[𝐚_b+\frac{𝐫_{ab}(𝐫_{ab}𝐚_b)}{r_{ab}^2}\right]\right\}$$
(6)
and the magnetic field as
$$𝐁(𝐫_a,t)=\underset{ba}{}\frac{e_b}{c}\frac{𝐯_b\times 𝐫_{ab}}{r_{ab}^3}$$
(7)
where $`𝐚_b`$ is the acceleration of particle $`b`$. In Eq. (6), the terms of order $`1/c^2`$ provide the familiar effects of Faraday induction. We can also write the electromagnetic fields in terms of electromagnetic potentials as
$$𝐄(𝐫_a,t)=_a\mathrm{\Phi }(𝐫_a,t)\frac{1}{c}\frac{}{t}𝐀(𝐫_a,t)\text{ and }𝐁(𝐫_a,t)=_a\times 𝐀(𝐫_a,t)$$
(8)
whereJ-2
$$\mathrm{\Phi }(𝐫_a,t)=\underset{ba}{}\frac{e_b}{r_{ab}}\text{ and }𝐀(𝐫_a,t)=\underset{ba}{}\frac{e_b}{2cr_{ab}}\left[𝐯_b+\frac{𝐫_{ab}(𝐫_{ab}𝐯_b)}{r_{ab}^2}\right]$$
(9)
We recognize from Eq. (2) and Eq. (9) that
$$𝐩_a^{canonical}=m_a𝐯_a[1+𝐯_a^2/(2c^2)]+(e_a/c)𝐀(𝐫_a,t)$$
(10)
where $`𝐀(𝐫_a,t)`$ is the vector potential due to all the other charges evaluated at the position $`𝐫_a`$ of the charge $`e_a.`$
## IV Two-Particle Model for a Magnetic Moment
Our model for a magnetic moment will consist of two charge particles of different mass in Coulomb orbit around each other (a classical hydrogen atom). There is no electromagnetic radiation in the Darwin Lagrangian, and thus the orbiting charges do not lose energy in this $`1/c^2`$ approximation. Furthermore, for our model, we will average over the phases of orbital motion and also average over the configurations where the both the charges and the velocities of the charges are reversed in sign. In this fashion one maintains the magnetic moment behavior while averaging out the irrelevant electrostatic aspects.
In this article, the motion of the magnetic moment charges is considered extensively. Therefore, for simplicity of notation (and in contrast to the notation of Coleman and Van Vleck), the magnetic moment consists of a particle of charge $`e`$, small mass $`m,`$ displacement $`𝐫`$, velocity $`𝐯`$, and acceleration $`𝐚`$ in orbit around a massive particle of charge $`e,`$ mass $`M`$ (with $`M>>m`$), displacement $`𝐑m𝐫/M0`$, velocity $`𝐕m𝐯/M`$, and acceleration $`d𝐕/dt.`$ Since the mass $`M`$ is large compared to $`m`$, the displacement $`𝐑`$, velocity $`𝐕`$, and acceleration $`d𝐕/dt`$ are all small compared to $`𝐫`$, $`𝐯`$, and $`𝐚`$ respectively. The distant point charge with which the magnetic moment interacts has charge $`q`$, mass $`m_q`$, displacement $`𝐫_q`$, velocity $`𝐯_q,`$ and acceleration $`d𝐯_q/dt.`$ Then from equations (4)-(7), our equations of motion for the charge $`e`$ in orbit, the massive particle $`e`$, and the distant charge $`q`$ are respectively
$`{\displaystyle \frac{d}{dt}}\left[m𝐯\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯^2}{c^2}}\right)\right]`$ $`=e𝐄_e(𝐫,t)+e𝐄_q(𝐫,t)+e{\displaystyle \frac{𝐯}{c}}\times 𝐁_q(𝐫,t)`$
$`={\displaystyle \frac{e^2𝐫}{r^3}}+{\displaystyle \frac{eq𝐫_{eq}}{r_{eq}^3}}\left[1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯_q^2}{c^2}}{\displaystyle \frac{3}{2}}{\displaystyle \frac{(𝐯_q𝐫_{eq})^2}{c^2r_{eq}^2}}\right]+e{\displaystyle \frac{𝐯}{c}}\times \left({\displaystyle \frac{q}{c}}{\displaystyle \frac{𝐯_q\times 𝐫_{eq}}{r_{eq}^3}}\right)`$ (11)
$`{\displaystyle \frac{d}{dt}}\left(M𝐕\right)`$ $`=e𝐄_e(0,t)e𝐄_q(0,t)`$
$`={\displaystyle \frac{e^2𝐫}{r^3}}\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯^2}{c^2}}{\displaystyle \frac{3}{2}}{\displaystyle \frac{(𝐯𝐫)^2}{c^2r^2}}\right)+{\displaystyle \frac{e^2}{2c^2r}}\left(𝐚+{\displaystyle \frac{(𝐚𝐫)𝐫}{r^2}}\right)`$
$`{\displaystyle \frac{eq𝐫_q}{r_q^3}}\left[1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯_q^2}{c^2}}{\displaystyle \frac{3}{2}}{\displaystyle \frac{(𝐯_q𝐫_q)^2}{c^2r_q^2}}\right]`$ (12)
and
$`{\displaystyle \frac{d}{dt}}\left[m_q𝐯_q\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯_q^2}{c^2}}\right)\right]`$ $`=q𝐄_e(𝐫_q,t)+q𝐄_e(𝐫_q,t)+q{\displaystyle \frac{𝐯_q}{c}}\times 𝐁_e(𝐫_q,t)`$
$`=q{\displaystyle \frac{e𝐫_q}{r_q^3}}+q{\displaystyle \frac{e𝐫_{qe}}{r_{qe}^3}}\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯^2}{c^2}}{\displaystyle \frac{3}{2}}{\displaystyle \frac{(𝐯𝐫_{qe})^2}{c^2r_{qe}^2}}\right)`$
$`q{\displaystyle \frac{e}{2c^2r_{qe}}}\left(𝐚+{\displaystyle \frac{(𝐚𝐫_{qe})𝐫_{qe}}{r_{qe}^2}}\right)+q{\displaystyle \frac{𝐯_q}{c}}\times \left({\displaystyle \frac{e}{c}}{\displaystyle \frac{𝐯\times 𝐫_{qe}}{r_{qe}^3}}\right)`$ (13)
where $`𝐫_{qe}=𝐫_q𝐫=𝐫_{eq}`$, and we have assumed that $`𝐕^2/c^2<<1.`$
### IV.1 Nonrelativistic Interaction
In order to understand the interaction represented by these equations of motion (11)-(13), we consider first the nonrelativistic approximation 0-order in $`1/c^2`$ where the equations become
$$m𝐚=\frac{e^2𝐫}{r^3}+e𝐄_q^{(0)}(𝐫,t)$$
(14)
$$M\frac{d𝐕}{dt}=\frac{e^2𝐫}{r^3}e𝐄_q^{(0)}(0,t)$$
(15)
and
$$m_q\frac{d𝐯_q}{dt}=q\frac{e𝐫_q}{r_q^3}+q\frac{e𝐫_{qe}}{r_{qe}^3}$$
(16)
Here the small electrostatic field of the charge $`q`$ is essentially uniform across the magnetic moment
$$𝐄_q^{(0)}(𝐫,t)=\frac{q(𝐫𝐫_q)}{|𝐫𝐫_q|^3}\frac{q𝐫_q}{r_q^3}=𝐄_q^{(0)}(0,t)$$
(17)
since the charge $`q`$ is distant from the magnetic moment at the origin of coordinates, $`r/r_q<<1.`$ The electrostatic field at the charge $`q`$ appearing on the right-hand side in Eq. (16) is an electric dipole field and is even smaller (for $`q`$ and $`e`$ of the same magnitude) because the magnetic moment is electrically neutral.
In this nonrelativistic approximation, the interaction of the distant point charge $`q`$ with this magnetic moment depends crucially upon the orientation of the magnetic moment. i)If the magnetic moment $`\stackrel{}{\mu }`$ at the origin is aligned parallel to the displacement $`𝐫_q`$ to the point charge, $`\stackrel{}{\mu }||𝐫_q`$, we find the stable electrostatic polarizability aspect. ii)If the magnetic moment $`\stackrel{}{\mu }`$ is aligned perpendicular to the displacement $`𝐫_q`$ to the point charge, $`\stackrel{}{\mu }𝐫_q`$, then we find Solem’sSolem unstable ”strange polarizability” aspect. It is the second, unfamiliar aspect which is crucial for understanding the electric forces which are proportional to the magnetic moment.
#### IV.1.1 Stable Electrostatic Polarizability
If the distant charge $`q`$ lies along the axis perpendicular to the orbital motion and through its center, $`\stackrel{}{\mu }||𝐫_q`$, then the electric field $`𝐄_q^{(0)}`$ will cause a displacement $`l`$ of the orbital plane relative to the massive particle $`M.`$ The equilibrium situation for the orbital motion with angular frequency $`\omega `$ corresponds to Newton’s equations of motion in the radial and axial directions giving
$$m\omega ^2r=e^2r(r^2+l^2)^{3/2}\text{ and }eE_q^{(0)}=e^2l(r^2+l^2)^{3/2}$$
(18)
Eliminating $`r`$ between the equations, we find $`e^2/(m\omega ^2)E_q^{(0)}=el=𝔭,`$ where $`𝔭`$ is the average electric dipole moment of the two-particle magnetic moment. Thus the magnetic moment in this orientation has an electrostatic polarizability
$$\alpha =e^2/(m\omega ^2)\text{ where }\stackrel{}{𝔭}=\alpha 𝐄_q^{(0)}$$
(19)
a form for $`\alpha `$ which is familiar for a dipole harmonic oscillator.J3 We notice that the polarizability is even in the charge $`e`$ and in the frequency $`\omega `$ and has no relation to the sign of the magnetic moment
$$\stackrel{}{\mu }=e𝐋/((2mc)=e\stackrel{}{\omega }r^2/(2c)$$
(20)
#### IV.1.2 Solem’s Unstable ”Strange” Polarization
If the magnetic moment is oriented perpendicular to the displacement to the distant charge $`q,`$ $`\stackrel{}{\mu }𝐫_q`$, then we find behavior which is mentioned only rarely in the physics literature.Solem It does not appear in Coleman and Van Vleck’s article,CVV but it is crucial to understanding the classical electromagnetic interactions associated with the Aharonov-Bohm and Aharonov-Casher phase shifts. In this case when the angular momentum $`𝐋`$ of the orbit for the magnetic moment is perpendicular to the electric field $`𝐄_q^{(0)}`$ of the distant charge $`q`$,$`\stackrel{}{\mu }𝐫_q,`$ the initial circular orbit is transformed into an elliptical orbit of ever-changing ellipticity with its semi-major axis oriented perpendicular to both the angular momentum $`𝐋`$ and the electric field $`𝐄_q^{(0)}`$.Solem In order to analyze this motion, it is useful to introduce the Laplace-Runge-Lenz vector $`𝐊`$ for the Coulomb orbit of the charge $`e.`$Gold We assume that the much larger mass $`M`$ is at the origin, $`𝐑0,`$ so that the charge $`e`$ moves with a displacementSolem
$$𝐫=\frac{3}{2}\frac{𝐊}{(2mH_0)^{1/2}}+\frac{1}{4H_0}\frac{d}{dt}[m(𝐫\times 𝐯)\times 𝐫+m𝐯r^2]$$
(21)
where $`𝐊`$ is the Laplace-Runge-Lenz vectorGold
$$𝐊=\frac{1}{(2mH_0)^{1/2}}\left([𝐫\times (m𝐯)]\times (m𝐯)+me^2\frac{𝐫}{r}\right)$$
(22)
and $`H_0`$ is the particle energy
$$H_0=mv^2/2e^2/r$$
(23)
The equation (21) can be checked by carrying out the time derivative and then inserting the equation of motion $`𝐚=e^2𝐫/(mr^3)`$ for every appearance of the acceleration $`𝐚=d𝐯/dt=d^2𝐫/dt^2.`$ The Laplace-Runge-Lenz vector is constant in time for a Coulomb orbit, and the second term of (21) involving a time derivative shows how the displacement $`𝐫`$ varies in time. On time-averaging, the time derivative vanishes leaving
$$<𝐫>=\frac{3}{2}\frac{𝐊}{(2mH_0)^{1/2}}$$
(24)
The average electric dipole moment $`\stackrel{}{𝔭}`$ is given by
$$\stackrel{}{𝔭}=e<𝐫>=\frac{3}{2}\frac{e𝐊}{(2mH_0)^{1/2}}$$
(25)
We assume that initially the magnetic moment has a circular orbit for the charge $`e`$, and therefore initially the electric dipole moment vanishes, $`\stackrel{}{𝔭}=e<𝐫>=0`$ and $`𝐊=0.`$ However, in the presence of the electric field $`𝐄_q^{(0)}`$ of the distant charge $`q`$, the equation of motion for $`e`$ is given in Eq. (14). We assume that the field $`𝐄_q^{(0)}`$ is small so that the orbit remains Coulombic but now with a slowly changing Laplace-Runge-Lenz vector. The time rate of change of $`𝐊`$ can be obtained by differentiating both sides of equation (22) and the use of the equation of motion (14),
$`{\displaystyle \frac{d𝐊}{dt}}`$ $`={\displaystyle \frac{1}{(2mH_0)^{1/2}}}m\left\{\left[𝐫\times \left({\displaystyle \frac{e^2𝐫}{r^3}}+e𝐄_q\right)\right]\times 𝐯+(𝐫\times 𝐯)\times \left({\displaystyle \frac{e^2𝐫}{r^3}}+e𝐄_q\right)+e^2\left[{\displaystyle \frac{𝐯}{r}}{\displaystyle \frac{𝐫(𝐫𝐯)}{r^3}}\right]\right\}`$
$`=(2mH_0)^{1/2}me[2𝐫(𝐯𝐄_q^{(0)})+𝐄_q^{(0)}(𝐫𝐯)+𝐯(𝐫𝐄_q^{(0)})]`$ (26)
We note again that the Laplace-Runge-Lenz vector would be constant in time were it not for the external electric field $`𝐄_q^{(0)}.`$ Since we assume that the distant charge $`q`$ is causing a small perturbation, we may average the particle displacement $`𝐫`$ and velocity $`𝐯`$ over an orbit of the unperturbed motion. Now if $`f(𝐫,𝐯)`$ is any function of the displacement and velocity of the unperturbed orbit, then it is a periodic function in time with period given by the orbital period $`T`$. Therefore, the time average of the time derivative vanishes
$$\frac{d}{dt}f(𝐫,𝐯)=\frac{1}{T}\underset{0}{\overset{t=T}{}}𝑑t\frac{d}{dt}f(𝐫,𝐯)=0$$
In particular for $`f(𝐫,𝐯)=x_ix_j`$ where $`x_i`$and $`x_j`$ are the $`i`$th and $`j`$th components of $`𝐫`$, then we have
$$\frac{d}{dt}(x_ix_j)=x_iv_j+x_jv_i=0$$
(27)
so that
$`𝐫𝐯`$ $`=0`$
$`𝐫\left(𝐯𝐄_q^{(0)}\right)`$ $`=𝐯\left(𝐫𝐄_q^{(0)}\right)=\left(𝐫\times 𝐯\right)\times 𝐄_q^{(0)}/2`$ (28)
This result allows us to average over the unperturbed motion to obtain
$`(2mH_0)^{1/2}d𝐊/dt`$ $`=me[2𝐫(𝐯𝐄_q^{(0)})+𝐄_q^{(0)}(𝐫𝐯)+𝐯(𝐫𝐄_q^{(0)})]`$
$`=(3/2)me[𝐫\times 𝐯\times 𝐄_q^{(0)}]=(3/2)e𝐋\times 𝐄_q^{(0)}=3m(c\stackrel{}{\mu })\times 𝐄_q^{(0)}`$ (29)
where $`𝐋`$ is the angular momentum of the orbit and $`\stackrel{}{\mu }=e𝐋/(2mc)`$. Thus from Eqs. (25) and (29), the electric dipole moment is changing as
$$\frac{d\stackrel{}{𝔭}}{dt}=\frac{9}{4}\frac{e^2}{(2mH_0)}𝐋\times 𝐄_q^{(0)}$$
(30)
This is a very strange polarization indeed. The initially unpolarized orbit does indeed develop an electrical polarization with time, but the predominant electric dipole moment depends upon the orbital angular momentum and is in a direction perpendicular to the applied electric field $`𝐄_q^{(0)}`$. Since the angular momentum $`𝐋`$ is related to the magnetic moment as $`\stackrel{}{\mu }=e𝐋/(2mc)`$, we have the developing polarization related to the magnetic moment. However, if we average over both the orbital positions and over both signs $`\pm e`$ of charge while maintaining the direction of the magnetic moment $`\stackrel{}{\mu }=e\stackrel{}{\omega }r^2/(2c),`$ then we see that the time rate of change of the Laplace-Runge-Lenz vector $`𝐊`$ does not average to zero while the average rate of change of polarization $`d\stackrel{}{𝔭}/dt`$ actually vanishes, since the direction of angular momentum in Eq. (20) reverses as the sign of the charge $`e`$ is reversed. We also notice that the rate of change of the Laplace-Runge-Lenz vector and of the electric dipole moment for an individual orbit depends upon the value of the field $`𝐄_q^{(0)}`$ alone and is independent of any rate of change of the electric field $`𝐄_q^{(0)}`$. This is completely different from the electrical polarization $`\stackrel{}{𝔭}`$ found from the electrostatic polarizability in Eq. (19) where there is no change in the polarization unless the field $`𝐄_q^{(0)}`$ changes in time.
There are additional observations which should be made regarding the behavior of the magnetic moment under the action of the electric field $`𝐄_q^{(0)}`$ of the distant point charge $`q`$. The sum of the particle kinetic energy plus electrostatic potential energy is conserved. Indeed, while the average displacement $`<𝐫>`$ of the charge $`e`$ is initially zero and increases in time, the length of the semimajor axis of the orbit does not change and is oriented in a direction perpendicular to the electric field $`𝐄_q^{(0)}`$; the work done by the electric field $`𝐄_q^{(0)}`$ on the orbiting charge $`e`$ vanishes when averaged over the Coulomb orbit. The average position of the heavier mass $`M`$ with charge $`e`$ also shifts slightly so as to maintain the position of the center of (rest) mass of the magnetic moment system at the origin; since the average electrostatic force on the magnetic moment (due to the uniform electric field $`𝐄_q^{(0)}`$ of the point charge $`q`$) vanishes, the position of the center of (rest) mass does not change. As the orbiting system develops an electric dipole moment $`\stackrel{}{𝔭}`$, there are balancing electrostatic forces and torques on the magnetic moment due to the point charge and on the point charge due to the magnetic moment. However, when we average over magnetic moments carrying opposite charges $`\pm e`$ but the same magnetic moment $`\stackrel{}{\mu }=e\stackrel{}{\omega }r^2/(2c)`$, all of the dipole-associated electrostatic forces and torques vanish in the average.
### IV.2 Electromagnetic Forces on the Distant Point Charge
#### IV.2.1 Force Associated with the Stable Electrostatic Polarization
Having obtained the behavior of the magnetic moment in the 0-order nonrelativistic system, we now wish to consider the electromagnetic forces $`𝐅_{onq}=q𝐄_\mu +q(𝐯_q/c)\times 𝐁_\mu `$ acting on the distant point charge $`q`$ due to the magnetic moment $`\stackrel{}{\mu }`$. The forces are different depending upon the orientation of the magnetic moment. When the magnetic moment $`\stackrel{}{\mu }`$ is parallel to the displacement $`𝐫_q`$ to the distant charge $`q,`$ $`\stackrel{}{\mu }||𝐫_q,`$ then we saw in Eq. (19) that the magnetic moment has an induced electric dipole moment $`\stackrel{}{𝔭}=\alpha 𝐄_q.`$ Accordingly, the electrically polarized magnetic moment creates an electrostatic dipole field $`𝐄_𝔭(𝐫_q,t)`$ which causes an electrostatic force $`𝐅_{on\text{ }q}`$ on $`q`$ given by
$$𝐅_{on\text{ }q}=q𝐄_𝔭(𝐫_q,t)=q\{2\stackrel{}{𝔭}\}r_q^3=q\{2[e^2/(m\omega ^2)]𝐄_q(0,t)\}r_q^3=𝐫_qq^2e^2/(m\omega ^2r_q^7)$$
(31)
The electrostatic force back at the charge $`q`$ is independent of the sign of the charge $`q,`$ or of the sign of the charge $`e,`$ or of the direction of rotation $`\omega `$. When averaged over the orbital motion and over both signs of charge $`\pm e`$ for the magnetic moment, the only force on $`q`$ is this electrostatic dipole force. There is no additional force of order $`1/c^2`$. As an aside, we note that for this orientation of the magnetic moment, $`\stackrel{}{\mu }||𝐫_q`$, the magnetic vector potential $`𝐀_\mu `$ vanishes along the axis through the magnetic moment parallel to the magnetic moment direction.
#### IV.2.2 Force Associated with Solem’s Unstable ”Strange” Polarization
The situation is completely different when the magnetic moment is oriented perpendicular to the displacement $`𝐫_q,`$ $`\stackrel{}{\mu }𝐫_q`$. In this case we saw that after carrying out the averaging for the magnetic moment, there were no electric monopole or dipole contributions to a force back on the point charge $`q`$. Since these 0-order back forces vanish, the back forces in order $`1/c^2`$ caused by the 0-order changes of the magnetic moment are of considerable interest. The alteration in the shape of the Coulomb orbit leads to unbalanced accelerations $`𝐚`$ which lead to new contributions to the electric field according to Eq. (6). The vector potential in the Coulomb gauge of a point charge $`e`$ is given in Eq. (9), and we see that the last term in Eq. (6) corresponds to the electric field contribution from $`𝐀_e/t=𝐀_\mu /t`$. Now the magnetic moment model corresponds to a magnetic moment given initially by $`\stackrel{}{\mu }=e\stackrel{}{\omega }r^2/(2c)`$ in Eq. (20), and therefore to a vector potential
$$𝐀_\mu (𝐫,t)=\frac{\stackrel{}{\mu }\times 𝐫}{cr^3}=\frac{e}{2mc^2}\frac{𝐋\times 𝐫}{r^3}$$
(32)
Thus for our magnetic moment model, the average electric field $`𝐄_\mu `$ back at the charged particle $`q`$ will be related to the change in the angular momentum $`𝐋`$ of the orbit. Now the change in angular momentum $`𝐋`$ of the orbit of the charge $`e`$ is due solely to the presence of the external charge $`q`$ which gives $`d𝐋/dt=𝐫\times e𝐄_q^{(0)}`$, and, when averaged over one period of the motion, becomes from Eq. (24)
$$\frac{d𝐋}{dt}=𝐫\times e𝐄_q^{(0)}=e𝐄_q^{(0)}\times 𝐫=e𝐄_q^{(0)}\times \frac{3}{2}\frac{𝐊}{(2mH_0)^{1/2}}$$
(33)
Thus the electric field back at the charge $`q`$ is given by
$`𝐄_\mu (𝐫_q,t)`$ $`={\displaystyle \frac{1}{c}}{\displaystyle \frac{}{t}}𝐀_\mu (𝐫_q,t)={\displaystyle \frac{e}{2mc}}{\displaystyle \frac{𝐫_q}{r_q^3}}\times {\displaystyle \frac{d𝐋}{dt}}`$
$`={\displaystyle \frac{e}{2mc^2}}{\displaystyle \frac{𝐫_q}{r_q^3}}\times \left(e𝐄_q^{(0)}\times {\displaystyle \frac{3}{2}}{\displaystyle \frac{𝐊}{(2mH_0)^{1/2}}}\right)`$ (34)
Now our magnetic moment model is initially in a circular orbit with $`𝐊=0,`$ and $`𝐊`$ changed as in Eq. (29) only because of the presence of the electric field $`𝐄_q^{(0)}`$ due to the distant charge $`q.`$ Thus the force $`𝐅_{on\text{ }q}`$ on $`q`$ due to the electric field $`𝐄_\mu `$ of the magnetic moment is
$$𝐅_{on\text{ }q}=q𝐄_\mu (𝐫_q,t)=q\frac{e}{2mc^2}\frac{𝐫_q}{r_q^3}\times \left(\frac{3}{2}\frac{e𝐄_q^{(0)}(0,t^{})}{(2mH_0)^{1/2}}\times \underset{0}{\overset{t}{}}𝑑t^{}\{3m[c\stackrel{}{\mu }(t^{})]\times 𝐄_q^{(0)}(0,t^{})\}\right)$$
(35)
where $`𝐄_q^{(0)}`$ is the electrostatic field in Eq. (17) of the distant charge $`q`$ acting on the magnetic moment. We notice that this force back on the charge $`q`$ due to the magnetic moment $`\stackrel{}{\mu }`$ is proportional to $`q^3e^2\mu `$; it changes sign with the external charge $`q`$, changes sign with the magnetic moment $`\stackrel{}{\mu },`$ but does not depend upon the sign of the charge $`e`$. Furthermore it changes sign with the reversal of the position $`𝐫_q`$ of the charge $`q`$. Finally, it does not depend upon any velocity of the charge $`q`$. It arises from the 0-order acceleration of the orbiting magnetic moment charge due to the electrostatic field $`𝐄_q^{(0)}`$ of the distant charge $`q.`$ These properties are in total contrast with those found for electrostatic forces such as in Eq. (31).
## V Conservation Laws
In our model, the (zero-order) electrostatic field of the passing charge causes a change in the magnetic moment which then produces an (order $`1/c^2`$) electric field back at the position of the passing charge. Since this back electric field is unanticipated by treatments (such as in the no-velocity-change point of view) which do not allow for changes in the charge and current densities of magnetic moments, it seems appropriate to discuss all the conservation laws associated with electromagnetic theory and to see how they are upheld by the present model.
### V.1 Linear Momentum in the Electromagnetic Field
The Darwin Lagrangian conserves linear momentumMomentum . For our magnetic moment and passing charge, the total linear momentum is
$`𝐏`$ $`=𝐏_\mu +𝐏_{em\mu q}+𝐩_q`$
$`=\left[M𝐕+m𝐯\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯^2}{c^2}}\right){\displaystyle \frac{e^2}{2c^2r}}\left(𝐯+{\displaystyle \frac{[𝐯𝐫]𝐫}{r^2}}\right)\right]`$
$`+\left[{\displaystyle \frac{qe}{2c^2|𝐫_q𝐫|}}\left(𝐯+{\displaystyle \frac{[𝐯(𝐫_q𝐫)](𝐫_q𝐫)}{|𝐫_q𝐫|^2}}\right)\right]+\left[m_q𝐯_q\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯_q^2}{c^2}}\right)\right]`$ (36)
Here we have grouped the total momentum into three terms which can be assigned to the magnetic moment, the electromagnetic fields between the magnetic moment and the charge $`q`$, and the mechanical momentum of the passing charge $`q`$. When averaged over the orbital motion of the magnetic moment, the system carries an average linear momentum in the electromagnetic field given by
$`𝐏_{em\mu q}`$ $`={\displaystyle \frac{1}{4\pi c}}{\displaystyle d^3r𝐄_q\times 𝐁_\mu }`$
$`={\displaystyle \frac{qe}{2c^2|𝐫_q𝐫|}}\left(𝐯+{\displaystyle \frac{[𝐯(𝐫_q𝐫)](𝐫_q𝐫)}{|𝐫_q𝐫|^2}}\right)`$
$`={\displaystyle \frac{q}{c}}{\displaystyle \frac{\stackrel{}{\mu }\times 𝐫_q}{r_q^3}}={\displaystyle \frac{q}{c}}𝐀_\mu (𝐫_q,t)`$ (37)
where, from Eq. (9), $`𝐀_\mu (𝐫_q,t)`$ is the vector potential in the Coulomb gauge due to the magnetic moment and evaluated at the position of the point charge $`q`$. Any contribution from the other electromagnetic field combination $`𝐄_\mu \times 𝐁_q`$ is very small since the magnetic moment $`\stackrel{}{\mu }`$ has no net charge.
Now the time derivative of the electromagnetic field momentum $`𝐏_{em\mu q}`$ in Eq. (37) involves changes connected with the particle position $`𝐫_q`$ and with the magnetic moment $`\stackrel{}{\mu }`$. We can write
$`{\displaystyle \frac{d}{dt}}𝐏_{em\mu q}`$ $`={\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{q}{c}}𝐀_\mu (𝐫_q,t)\right)`$
$`=(𝐯_q_q)\left({\displaystyle \frac{q}{c}}𝐀_\mu (𝐫_q,t)\right)+{\displaystyle \frac{}{t}}\left({\displaystyle \frac{q}{c}}𝐀_\mu (𝐫_q,t)\right)`$
$`=_q\left({\displaystyle \frac{q}{c}}𝐯_q𝐀_\mu (𝐫_q,t)\right){\displaystyle \frac{q}{c}}𝐯_q\times \left[_q\times 𝐀_\mu (𝐫_q,t)\right]+{\displaystyle \frac{}{t}}\left({\displaystyle \frac{q}{c}}𝐀_\mu (𝐫_q,t)\right)`$
$`=_q\left({\displaystyle \frac{q}{c}}𝐯_q{\displaystyle \frac{\stackrel{}{\mu }\times 𝐫_q}{r_q^3}}\right){\displaystyle \frac{q}{c}}𝐯_q\times \left[_q\times \left({\displaystyle \frac{\stackrel{}{\mu }\times 𝐫_q}{r_q^3}}\right)\right]+{\displaystyle \frac{q}{c}}\left({\displaystyle \frac{d\stackrel{}{\mu }}{dt}}\times {\displaystyle \frac{𝐫_q}{r_q^3}}\right)`$
$`=𝐅_{on\mu }^{Lorentz}𝐅_{onq}^{Lorentz}`$ (38)
where
$$𝐅_{on\mu }^{Lorentz}=_𝐫[\stackrel{}{\mu }𝐁_q(𝐫,t)]_{r=0}=_q(\frac{q}{c}𝐯_q𝐀_\mu (𝐫_q,t))$$
(39)
and
$$𝐅_{onq}^{Lorentz}=q𝐄_\mu (𝐫_q,t)+q\frac{𝐯_q}{c}\times 𝐁_\mu (𝐫_q,t)$$
(40)
with
$$𝐄_\mu (𝐫_q,t)=\frac{}{t}𝐀_\mu (𝐫_q,t)=\frac{d\stackrel{}{\mu }}{dt}\times \frac{𝐫_q}{r_q^3}$$
(41)
and
$$𝐁_\mu (𝐫_q,t)=_q\times 𝐀_\mu (𝐫_q,t)=_q\times \left(\frac{q}{c}\frac{\stackrel{}{\mu }\times 𝐫_q}{r_q^3}\right)$$
(42)
Thus the average electromagnetic linear momentum $`𝐏_{em\mu q}`$ in Eq. (37) changes with respect to time for two reasons: the change in $`\stackrel{}{\mu }`$ (due to the change in the orbital shape of the magnetic moment) and the change in the separation $`𝐫_q`$. As the shape changes for the orbit of the charge $`e`$ in the magnetic moment, the magnetic moment $`\stackrel{}{\mu }`$ changes creating an electric field at the position of the passing particle $`q`$. Thus due to this changing-$`\mu `$ effect, the linear momentum $`𝐏_{em\mu q}`$ in the electromagnetic field decreases at the same rate that the linear momentum of the point charge $`q`$ increases due to the force from the electric field of the changing magnetic moment. The change in the electromagnetic linear momentum $`𝐏_{em\mu q}`$ due to the changing position $`𝐫_q`$ is associated with the magnetic Lorentz forces on the magnetic moment and on the passing charge.
Next we average the total system momentum in Eq. (36) over the orbital motion and differentiate with respect to time to find
$`{\displaystyle \frac{d𝐏}{dt}}`$ $`=0=\left[{\displaystyle \frac{d𝐏_\mu }{dt}}+_q\left({\displaystyle \frac{q}{c}}{\displaystyle \frac{\stackrel{}{\mu }\times 𝐫_q}{r_q^3}}\right)\right]`$
$`+\left[{\displaystyle \frac{q}{c}}\left({\displaystyle \frac{d\stackrel{}{\mu }}{dt}}\right)\times {\displaystyle \frac{𝐫_q}{r_q^3}}{\displaystyle \frac{q}{c}}𝐯_q\times \left[_q\times \left({\displaystyle \frac{q}{c}}{\displaystyle \frac{\stackrel{}{\mu }\times 𝐫_q}{r_q^3}}\right)\right]+{\displaystyle \frac{d𝐩_q}{dt}}\right]`$
$`=\left[{\displaystyle \frac{d𝐏_\mu }{dt}}𝐅_{on\mu }^{Lorentz}\right]+\left[{\displaystyle \frac{d𝐩_q}{dt}}𝐅_{onq}^{Lorentz}\right]`$ (43)
The equations of motion tell us that each of the quantities in square brackets vanishes. Note that the sum of the average Lorentz forces $`𝐅_{on\mu }^{Lorentz}+𝐅_{onq}^{Lorentz}`$ does not vanish, but rather (according to Eq. (38)) is equal to the negative rate of change of the electromagnetic field linear momentum $`𝐏_{em\mu q}`$. Thus in the conservation law for linear momentum, the changing electromagnetic field momentum $`𝐏_{em\mu q}`$ is partially balanced by the changing momentum of the magnetic moment and partially balanced by the changing momentum of the passing particle.
### V.2 Energy Conservation
The Darwin lagrangian conserves energy.energy For our magnetic moment and passing charge, the total energy through order $`1/c^2`$ is
$`U`$ $`=U_\mu +U_{em\mu q}+U_q`$
$`=[Mc^2+mc^2(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{v^2}{c^2}}+{\displaystyle \frac{3}{8}}{\displaystyle \frac{v^4}{c^4}}){\displaystyle \frac{e^2}{r}}]+[{\displaystyle \frac{eq}{r_q}}+{\displaystyle \frac{eq}{r_{eq}}}`$
$`+{\displaystyle \frac{eq}{2c^2r_{eq}}}(𝐯𝐯_q+{\displaystyle \frac{(𝐯𝐫_{eq})(𝐯_q𝐫_{eq})}{r_{eq}^2}})]+[m_qc^2(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{v_q^2}{c^2}}+{\displaystyle \frac{3}{8}}{\displaystyle \frac{v_q^4}{c^4}})]`$ (44)
When averaged over the orbital motion of the magnetic moment, the electrostatic energy $`eq/r_q+eq/r_{eq}`$ involves only quadrupole energies, which vanish when averaged $`\pm e,`$ $`\pm \omega `$ with $`\stackrel{}{\mu }`$ held constant. The system carries an average magnetic energy in the electromagnetic field given by
$`U_{em\mu q}`$ $`={\displaystyle \frac{1}{8\pi }}{\displaystyle d^3r𝐁_q\times 𝐁_\mu }`$
$`={\displaystyle \frac{eq}{2c^2r_{eq}}}\left(𝐯𝐯_q+{\displaystyle \frac{(𝐯𝐫_{eq})(𝐯_q𝐫_{eq})}{r_{eq}^2}}\right)`$
$`=\stackrel{}{\mu }𝐁_q(0,t)=\stackrel{}{\mu }\left({\displaystyle \frac{q}{c}}{\displaystyle \frac{𝐯_q\times (𝐫_q)}{r_q^3}}\right)`$
$`={\displaystyle \frac{q}{c}}𝐯_q{\displaystyle \frac{\stackrel{}{\mu }\times 𝐫_q}{r_q^3}}={\displaystyle \frac{q}{c}}𝐯_q𝐀_\mu (𝐫_q,t)`$ (45)
The time derivative of the magnetic field energy $`U_{em\mu q}`$ can be written using the calculations in Eq. (38) for $`d𝐀_\mu /dt`$
$`{\displaystyle \frac{d}{dt}}U_{em\mu q}`$ $`={\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{q}{c}}𝐯_q𝐀_\mu (𝐫_q,t)\right)=𝐯_q{\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{q}{c}}𝐀_\mu (𝐫_q,t)\right)`$
$`=𝐯_q\left(𝐅_{on\mu }^{Lorentz}𝐅_{onq}^{Lorentz}\right)`$ (46)
since $`U_{em\mu q}`$ is already of order $`1/c^2`$ and any change in $`𝐯_q`$ due to the changing magnetic moment $`\stackrel{}{\mu }`$ is also of order $`1/c^2.`$
Since the total energy in Eq. (44) is constant in time, it follows from averaging over the orbital motion and differentiating with respect to time that
$`{\displaystyle \frac{dU}{dt}}`$ $`=0={\displaystyle \frac{dU_\mu }{dt}}+{\displaystyle \frac{d}{dt}}U_{em\mu q}+{\displaystyle \frac{dU_q}{dt}}`$
$`=\left[{\displaystyle \frac{dU_\mu }{dt}}+{\displaystyle \frac{q}{c}}𝐯_q\left(\stackrel{}{\mu }\times {\displaystyle \frac{d}{dt}}{\displaystyle \frac{𝐫_q}{r_q^3}}\right)\right]+\left[{\displaystyle \frac{q}{c}}𝐯_q\left({\displaystyle \frac{d\stackrel{}{\mu }}{dt}}\times {\displaystyle \frac{𝐫_q}{r_q^3}}\right)+{\displaystyle \frac{dU_q}{dt}}\right]`$
$`=\left[{\displaystyle \frac{dU_\mu }{dt}}𝐯_q𝐅_{on\mu }^{Lorentz}\right]+\left[{\displaystyle \frac{dU_q}{dt}}𝐯_q𝐅_{onq}^{Lorentz}\right]`$ (47)
Here we have used the calculations in Eqs. (45) and (46); we also note that the dot product of $`𝐯_q`$ with the term involving $`𝐯_q\times [_q\times 𝐀_\mu ]`$ in Eq. (38) vanishes. The average energy in the magnetic field $`U_{em\mu q}`$ changes because of the changing magnetic moment $`\stackrel{}{\mu }`$ and also due to the changing position $`𝐫_q`$ of the passing charge $`q`$. Just as above in Eq. (41), the changing magnetic moment is associated with an electric field $`𝐄_\mu (𝐫_q,t)`$ back at the passing charge which changes the kinetic energy of the passing charge.
$`{\displaystyle \frac{dU_q}{dt}}`$ $`=𝐯_q𝐅_{onq}^{Lorentz}=q𝐄_\mu (𝐫_q,t)𝐯_q`$
$`=𝐯_q{\displaystyle \frac{q}{c}}{\displaystyle \frac{}{t}}𝐀_\mu (𝐫_q,t)=𝐯_q\left[{\displaystyle \frac{q}{c}}\left({\displaystyle \frac{d}{dt}}\stackrel{}{\mu }\right)\times {\displaystyle \frac{𝐫_q}{r_q^3}}\right]`$ (48)
The change in the magnetic field energy associated with the changing position $`𝐫_q`$ of the passing charge is compensated by the change in the kinetic energy (in order $`1/c^2`$) of the orbiting charge of the magnetic moment. This energy change can be written in various forms
$`{\displaystyle \frac{dU_\mu }{dt}}`$ $`=𝐯_q𝐅_{on\mu }^{Lorentz}={\displaystyle \frac{q}{c}}𝐯_q\stackrel{}{\mu }\times {\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{𝐫_q}{r_q^3}}\right)=\stackrel{}{\mu }{\displaystyle \frac{}{t}}𝐁_q(0,t)`$
$`=\left(𝐯_q_q\right)\left({\displaystyle \frac{q}{c}}𝐯_q𝐀_\mu (𝐫_q,t)\right)`$
$`=\left(𝐯_q_q\right){\displaystyle \frac{eq}{2c^2r_{eq}}}\left(𝐯𝐯_q+{\displaystyle \frac{(𝐯𝐫_{eq})(𝐯_q𝐫_{eq})}{r_{eq}^2}}\right)`$
$`=e𝐯{\displaystyle \frac{q𝐫_{eq}}{c^2r_{eq}^3}}\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{v_q^2}{c^2}}+{\displaystyle \frac{3}{2}}{\displaystyle \frac{(𝐯_q𝐫_{eq})^2}{c^2r_{eq}^2}}\right)=e𝐯[𝐄_q(𝐫,t)𝐄_q^{(0)}(𝐫,t)]`$ (49)
and corresponds to energy delivered to a moving charge by the emf of the changing magnetic field of the passing charge. We notice that it is the relativistic $`v_q^2/c^2`$ terms in the electric field $`𝐄_q`$ which deliver the power to the orbiting charge. Thus in the energy conservation law, the changing magnetic field energy $`U_{em\mu q}`$ is associated with the changing energy of both the magnetic moment and the passing charge.
### V.3 Center-of-Energy Motion for the Magnetic Moment
In this section we will discuss the motion of the center of energy of the magnetic moment from two points of view. First we connect its motion to the motion of the passing charge using the conservation law for the constant motion of the system center of energy. Second we use the particle equations of motion to obtain what has been called ”the equation of motion of the magnet,” referring to the center of energy motion of the magnetic moment
The Darwin Lagrangian gives constant velocity to the system center of energy to order $`1/c^2`$, the same order to which the Darwin Lagrangian is invariant under Lorentz transformations.CVV The center of energy $`\stackrel{}{𝐗}`$ to order $`1/c^2`$ involves only rest mass energy and electrostatic energy
$`{\displaystyle \frac{U}{c^2}}\stackrel{}{𝐗}`$ $`={\displaystyle \frac{U_\mu }{c^2}}\stackrel{}{𝐗}_\mu +\left[{\displaystyle \frac{U_{emeq}}{c^2}}\left({\displaystyle \frac{𝐫+𝐫_q}{2}}\right)+{\displaystyle \frac{U_{emeq}}{c^2}}\left({\displaystyle \frac{𝐫_q}{2}}\right)\right]+{\displaystyle \frac{U_q}{c^2}}𝐫_q`$
$`={\displaystyle \frac{1}{c^2}}\left(U_\mu ^{[q=0]}e𝐫𝐄_q(0)\right)\left(\stackrel{}{𝐗}_\mu ^{[q=0]}+\delta \stackrel{}{𝐗}_\mu \right)`$
$`+\left[{\displaystyle \frac{eq}{r_{eq}}}\left({\displaystyle \frac{𝐫+𝐫_q}{2}}\right)+{\displaystyle \frac{eq}{r_q}}\left({\displaystyle \frac{𝐫_q}{2}}\right)\right]+\left[m_q\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯_q^2}{c^2}}\right)\right]𝐫_q`$ (50)
where $`U_\mu ^{[q=0]}`$ and $`\stackrel{}{𝐗}_\mu ^{[q=0]}`$ correspond to the energy and center of energy of the magnetic moment when the passing charge is not present. When averaged over the orbital motion of the magnetic moment, the electromagnetic field contribution in Eq. (50) yields a quadrupole contribution, corresponding to the neutrality of the magnetic moment,
$`\left[{\displaystyle \frac{eq}{r_{eq}}}\left({\displaystyle \frac{𝐫+𝐫_q}{2}}\right)+{\displaystyle \frac{eq}{r_q}}\left({\displaystyle \frac{𝐫_q}{2}}\right)\right]`$ $`={\displaystyle \frac{eq}{2}}{\displaystyle \frac{𝐫}{r_q}}{\displaystyle \frac{(𝐫𝐫_q)𝐫_q}{r_q^2}}+O\left({\displaystyle \frac{r^2}{r_q^2}}\right)`$
$`={\displaystyle \frac{eq}{2}}O\left({\displaystyle \frac{r^2}{r_q^2}}\right)`$ (51)
since $`𝐫=0`$, and this contribution vanishes entirely if we average over $`\pm e`$ and $`\pm \omega `$ so as to keep only the magnetic moment contribution. Furthermore, the magnetic momentum contribution in Eq. (50) can be averaged over the orbital motion to give (through first order in the interaction perturbation)
$`{\displaystyle \frac{U_\mu }{c^2}}\stackrel{}{𝐗}_\mu `$ $`={\displaystyle \frac{1}{c^2}}\left(U_\mu ^{[q=0]}e𝐫𝐄_q(0)\right)\left(\stackrel{}{𝐗}_\mu ^{[q=0]}+\delta \stackrel{}{𝐗}_\mu \right)`$
$`={\displaystyle \frac{1}{c^2}}U_\mu ^{[q=0]}\left(\stackrel{}{𝐗}_\mu ^{[q=0]}+\delta \stackrel{}{𝐗}_\mu \right)={\displaystyle \frac{U_\mu }{c^2}}\stackrel{}{𝐗}_\mu `$
since $`𝐫=0`$. It follows that Eq. (50) becomes
$$\frac{U}{c^2}\stackrel{}{𝐗}=\frac{U_\mu }{c^2}\stackrel{}{𝐗}_\mu +\frac{U_q}{c^2}𝐫_q$$
(52)
Now differentiating twice with respect to time and noting that the energies $`U`$, $`U_\mu ,`$ and $`U_q`$ are all constant in time through 0-order in $`1/c^2,`$ while $`d^2\stackrel{}{𝐗}/dt^2=0`$, we find
$$0=\frac{U_\mu }{c^2}\frac{d^2\stackrel{}{𝐗}_\mu }{dt^2}+\frac{U_q}{c^2}\frac{d^2𝐫_q}{dt^2}$$
(53)
Thus the motions of the centers of energy of the magnetic moment and the passing charge are coupled together. Our equations (52) and (53) here correspond to Eqs. (14) and (15) in Coleman and Van Vleck’s discussion of the interaction of a point charge and a magnet.
For the magnetic moment alone, the center of energy $`\stackrel{}{𝐗}_\mu `$ is defined as
$$\frac{U_\mu }{c^2}\stackrel{}{𝐗}_\mu =m\left(1+\frac{1}{2}\frac{𝐯^2}{c^2}\right)𝐫+M𝐑\frac{e^2}{c^2r}\left(\frac{𝐫}{2}\right)$$
(54)
where the energy $`U_\mu `$ of the magnetic moment through 0-order in $`1/c^2`$ is
$$U_\mu =mc^2\left(1+\frac{1}{2}\frac{𝐯^2}{c^2}\right)+Mc^2\frac{e^2}{r}$$
(55)
and where we have taken the displacement $`𝐑`$ of the large mass $`M`$ as small compared to $`𝐫`$. In the nonrelativistic (0-order $`1/c`$) limit, the center of energy $`\stackrel{}{𝐗}_\mu ^{(0)}`$ corresponds to the center of (rest) mass
$$(m+M)\stackrel{}{𝐗}_\mu ^{(0)}=m𝐫+M𝐑$$
(56)
which in our example has been chosen so $`\stackrel{}{𝐗}_\mu ^{(0)}=0.`$ Furthermore, the center of (rest) mass remains at rest since differentiating Eq. (56) with respect to time leads to the nonrelativistic statement regarding the momentum of the magnetic moment
$$(m+M)\frac{d}{dt}\stackrel{}{𝐗}_\mu ^{(0)}=m𝐯+M𝐕=0$$
(57)
The 0-order (nonrelativistic) linear momentum of the magnetic moment indeed vanishes since the internal Coulomb forces within the magnetic moment satisfy Newton’s third law and the nonrelativistic Coulomb forces on the two oppositely charged particles of the magnetic moment due to the distant point charge $`q`$ are equal and opposite in the approximation of Eq. (17).
If we differentiate Eq. (54) for the center of energy of the magnetic moment, we obtain,
$$\frac{U_\mu }{c^2}\frac{d\stackrel{}{𝐗}_\mu }{dt}+\frac{1}{c^2}\frac{dU_\mu }{dt}\stackrel{}{𝐗}_\mu =m\left(1+\frac{1}{2}\frac{𝐯^2}{c^2}\right)𝐯+M𝐕\frac{e^2}{2c^2r}\left(𝐯\frac{(𝐫𝐯)𝐫}{r^2}\right)+\frac{m𝐫(𝐯𝐚)}{c^2}$$
(58)
The acceleration $`𝐚`$ of the orbiting charge is given in Eq. (14) and the time derivative of the energy is related to the work done by the electric field of the passing charge $`dU_\mu /dt=e𝐯𝐄_q`$. Then averaging over the orbital motion, equation (58) becomes
$`{\displaystyle \frac{U_\mu }{c^2}}{\displaystyle \frac{d\stackrel{}{𝐗}_\mu }{dt}}+{\displaystyle \frac{1}{c^2}}{\displaystyle \frac{dU_\mu }{dt}}\stackrel{}{𝐗}_\mu `$ $`={\displaystyle \frac{U_\mu }{c^2}}{\displaystyle \frac{d\stackrel{}{𝐗}_\mu }{dt}}+{\displaystyle \frac{1}{c^2}}e𝐯𝐄_q^{(0)}\stackrel{}{𝐗}_\mu ={\displaystyle \frac{U_\mu }{c^2}}{\displaystyle \frac{d\stackrel{}{𝐗}_\mu }{dt}}`$
$`=<m(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯^2}{c^2}})𝐯+M𝐕{\displaystyle \frac{e^2}{2c^2r}}(𝐯{\displaystyle \frac{(𝐫𝐯)𝐫}{r^2}})`$
$`+{\displaystyle \frac{𝐫}{c^2}}\left[𝐯\left({\displaystyle \frac{e^2𝐫}{r^3}}+e𝐄_q^{(0)}(𝐫,t)\right)\right]`$ $`>`$ (59)
where we have noted $`e𝐯𝐄_q^{(0)}=0`$. Now combining the terms involving $`e^2`$, and rewriting the average of $`e𝐫\left[𝐯𝐄_q^{(0)}\right]/c^2`$ as in Eq. (28), we have
$`{\displaystyle \frac{U_\mu }{c^2}}{\displaystyle \frac{d\stackrel{}{𝐗}_\mu }{dt}}`$ $`=m\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯^2}{c^2}}\right)𝐯+M𝐕{\displaystyle \frac{e^2}{2c^2r}}\left(𝐯+{\displaystyle \frac{(𝐫𝐯)𝐫}{r^2}}\right){\displaystyle \frac{1}{c}}\stackrel{}{\mu }\times 𝐄_q^{(0)}(0,t)`$
$`=𝐏_\mu {\displaystyle \frac{1}{c}}\stackrel{}{\mu }\times 𝐄_q^{(0)}(0,t)`$ (60)
This result (60) corresponds to Eq. (26) of the work by Coleman and Van Vleck.CVV Next differentiating Eq. (60) with respect to time so as to obtain a second derivative of $`\stackrel{}{𝐗}_\mu `$
$$\frac{U_\mu }{c^2}\frac{d^2\stackrel{}{𝐗}_\mu }{dt^2}=\frac{d}{dt}𝐏_\mu \frac{d}{dt}\left(\frac{1}{c}\stackrel{}{\mu }\times 𝐄_q^{(0)}(0,t)\right)$$
(61)
This equation is sometimes called ”the equation of motion for a magnetic moment.”APV Vaidman
### V.4 The Argument over Hidden Momentum in Magnets
Because the interaction of a magnet and a passing point charge is so poorly understood, there can arise certain notions which are used as ”explanations” but are not explored in detail. ”Hidden momentum in magnets” is such a notion. We will illustrate the situation using our calculations for the interaction of a point charge and a hydrogen-atom magnetic moment which we have calculated above.
Because the proponents of the no-velocity-change point of view are so sure that there is no force back on a charged particle passing a magnet, they also feel sure that there must be no change in the center of energy of the magnet. Thus if the center of energy of the magnet did change position, then according to our Eq. (53) (and according to Coleman and Van Vleck’s Eq. (15)), the passing charge must accelerate. Moreover, there is clearly a possibility of acceleration for the magnet’s center of energy since there is an obvious magnetic Lorentz force on the magnet given by $`𝐅_{on\mu }=(\stackrel{}{\mu }𝐁_q)`$. Now fundamental classical theorems connect the force and changes in system momentum so that we must have $`𝐅_{on\mu }=(\stackrel{}{\mu }𝐁_q)=d𝐏_\mu /dt.`$ But our equation (61) gives an escape from motion for the center of energy of the magnet because there is a second term in the expression for the acceleration of the center of energy. Thus the proponents of the no-velocity-change point of view decide that the quantity $`(1/c)\stackrel{}{\mu }\times 𝐄_q`$ represents a ”hidden momentum in magnets” whose change ”cancels” the classical applied force. Indeed, a mechanical momentum of the required form is mentioned in a footnote in Coleman and Van Vleck’s workCVV and now appears in an electromagnetism text book.Griffiths However, no one who speaks of ”hidden momentum in magnets” has ever given any relativistic calculation which shows how this momentum carries out this cancellation without continuing changes in the charge and current densities of the magnet. ”Hidden momentum in magnets” (as used by the proponents of the no-velocity-change point of view) seems to be an idea which exists simply to prevent the motion of the center of energy of a magnet. As we see above in our explicit model of a hydrogen-atom magnetic moment and a point charge, there is indeed a force back on the passing charge and there is indeed motion of the center of energy of the magnet. Both of these results are contrary to the claims of the proponents of the no-velocity-change point of view.
## VI Transition to a Multiparticle Magnet
Experimental observation of the interaction of a magnet and a point charge (such as in the Aharonov-Bohm phase shift) involves not two-particle magnetic moments but rather multiparticle magnets. We are interested in understanding the experimental situation based upon the insight gained from the fundamental interaction involving a two-particle magnetic dipole moment.
Within classical electromagnetism, the transition to a multiparticle system is most familiar for the electrostatics of polarizable particles. In our calculation above, we found that our magnetic moment oriented in the direction of the displacement $`𝐫_q,`$ $`\stackrel{}{\mu }||𝐫_q,`$ acted like a polarizable particle producing a back force of magnitude $`F_{on\text{ }q}=q^2e^2/(m\omega ^2r_q^6)`$ back on the point charge $`q.`$ When the polarizability is larger (for example, $`m`$ is smaller for fixed $`\omega `$), then the force back on the distant charge is larger. Also, when we have many polarizable particles present, the force back on the distant particle does not disappear but rather increases to a well-defined limit. Thus if we consider a dielectric wall formed by polarizable particles, then the mutual interaction among the polarizable particles changes the functional dependence of the force over toward $`F_{on\text{ }q}=q^2/(2r_q)^2,`$ which holds for a conducting wall where the force is independent of the polarizability in the limit of large polarizability. This occurs because polarizable particles which are next to each other in the wall form electric dipole moments which tend to cancel the external electric field $`𝐄_q`$ at the position of the other electric dipoles in the wall.
In an analogous fashion, we expect multiparticle interactions within a magnet to alter the back force on a passing charge found in Eq. (35). We note that the force back at the passing charge $`q`$ due to our model magnetic moment can be varied by changing the mass of the orbiting charge while keeping the magnetic moments fixed. When the magnetic moment involves a small mass $`m`$ (and thus is easily influenced by the external electric field $`𝐄_q)`$, the force back at the passing charge is larger, just as is true for a polarizable particle in the electrostatic situation. The most symmetrical multiparticle arrangement of magnetic moments involves $`N`$ magnetic moments arranged around a circle as a toroid with the distant charged particle $`q`$ located along the axis of the toroid. The 0-order (nonrelativistic) electrostatic force on each of the orbiting charges of the toroid due to the charge $`q`$ is $`e𝐄_q^{(0)}`$ just as before, while the back force on the charge $`q`$ is now $`N`$ times as large. Again, in analogy with the electrostatic situation, we expect that due to multiparticle interactions within the toroid the back force on a passing charge will not disappear but rather will increase to a limit. Now there will be nonrelativistic electrostatic forces between the charges of the $`N`$ magnetic moments. Also, each of the orbiting charges $`e`$ produces acceleration fields of order $`1/c^{2\text{ }}`$which act on all of the other orbiting charges of the magnetic moment. Since the $`1/c^2`$-acceleration fields act on each of the other orbiting charges of the toroid, the back force on each orbiting charge $`e`$ increases as the number $`N`$ of two-particle magnetic moments increases. These acceleration fields always cause forces such as to oppose any change in the currents of the toroid. This corresponds to a self-inductance effect which increases as $`N^2`$ when there are $`N`$ current-carrying loops.
It is important to notice that the present situation does not correspond to the elementary mutual-inductance problem of electromagnetism texts. In mutual inductance effects, the self-induced emf is such as to oppose any change in magnetic flux introduced externally and the magnitude of the back emf is independent of the current which is flowing in the toroid winding. In our case here, the initial accelerations tending to change the magnetic flux through the toroid do not arise from any induced emf through the toroid. Indeed in the limit $`𝐯_q=0`$ there is no emf at all in the toroid. Furthermore, the back force on the charge $`q`$ does not behave as in Lenz’s law. Rather, the tendency to change the currents of the toroid arises from Solem’s strange polarization associated with the electrostatic field of the external charge $`q`$ treated as a uniform electric field across each magnetic moment; the change in the magnetic moment is proportional to the magnetic moment and changes sign with the sign of $`q`$ as seen in Eq. (35).
We expect that in the multiparticle limit, the electrostatic interactions within the toroid will tend to screen the field of the passing charge $`q`$ out of the toroid and the back force on the passing charge will be limited by the magnetic energy of interaction. Indeed, calculations for ohmic conductors suggest that the electric fields of a passing charge are screened out of the body of the conductor by surface charges while the magnetic fields of the passing charge penetrate into the body of the conductor.B1999 We note that if the point charge is held at rest outside a conductor, then the electric fields of the point charge are screened out of the body of the conductor by surface charges. If the charged particle is moving, we do not expect this electric-field screening to suddenly disappear. On the other hand, it has been shown that magnetic fields due to moving charges penetrate into an ohmic conductor giving a time-integral of the magnetic field which is independent of the conductivity of the materials.B1999 As was suggested earlier, this is precisely the result which is needed to account for the Aharonov-Bohm phase shift as a classical lag associated with energy-related classical forces.B5
### VI.1 Energy, Momentum, and Forces in the Multiparticle Limit
Let us now consider the momentum, energy, and forces when a charged particle $`q`$ moves with velocity $`𝐯_q`$ down the axis of a magnet in the form of a toroid which is initially at rest. The screening of the electric field of the passing charge out of the body of the magnet implies that the electric field vanishes inside the toroid and therefore there is no significant contribution to momentum from the electromagnetic field of the form $`𝐏_{em\mu q}`$ discussed above, and no significant energy flow across the magnet. It follows from Eq. (36), that now the total system momentum consists of only two contributions, one each from the magnet and the passing charge
$$𝐏=𝐏_\mu +m_q𝐯_q\left(1+\frac{1}{2}\frac{𝐯_q^2}{c^2}\right)\text{ multiparticle limit}$$
(62)
The Lorentz forces on the magnet and on the passing charge then satisfy Newton’s third law
$`0`$ $`={\displaystyle \frac{d𝐏}{dt}}={\displaystyle \frac{d𝐏_\mu }{dt}}+{\displaystyle \frac{d}{dt}}\text{ }\left[m_q𝐯_q\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{𝐯_q^2}{c^2}}\right)\right]`$
$`=𝐅_{on\mu }^{Lorentz}+𝐅_{onq}^{Lorentz}\text{ multiparticle limit}`$ (63)
Furthermore, since the electric field is screened out of the body of the magnet, the center of energy motion of the magnet in Eq. (61) becomes the familiar Newton’s second law connecting the center of mass motion with the net Lorentz force
$$\frac{U_\mu }{c^2}\frac{d^2\stackrel{}{𝐗}_\mu }{dt^2}=\frac{d}{dt}𝐏_\mu =𝐅_{on\mu }^{Lorentz}$$
(64)
The net Lorentz force on the magnet is exactly the original standard classical magnetic Lorentz forceJ5 on the magnet due to the magnetic fields of the passing charge,
$`𝐅_{on\mu }^{Lorentz}`$ $`=\left[_𝐫\left\{\stackrel{}{\mu }𝐁_q(𝐫,t)\right\}\right]_{𝐫=0}=_q\left\{\stackrel{}{\mu }\left(q{\displaystyle \frac{𝐯_q}{c}}\times {\displaystyle \frac{(𝐫_q)}{r_q^3}}\right)\right\}`$
$`=_q\left\{{\displaystyle \frac{q}{c}}𝐯_q[\stackrel{}{\mu }\times {\displaystyle \frac{𝐫_q}{r_q^3}}]\right\}={\displaystyle \frac{q}{c}}\left(𝐯_q_q\right)𝐀_\mu (𝐫_q)`$ (65)
where we have written the magnetic field of the charged particle evaluated at the origin as $`𝐁_q=q𝐯\times (𝐫_q)c^1r_q^3`$, have used standard vector identities, have recognized the magnetic vector potential $`𝐀_\mu (𝐫_q)=\stackrel{}{\mu }\times 𝐫_q/r_q^3`$ of the magnet at the position of the charged particle $`q`$, and have dropped the magnetic Lorentz force $`(q/c)𝐯_q\times 𝐁_\mu `$ which vanishes for a point charge $`q`$ on the axis of a toroid. Newton’s third law in Eq. (63) for the forces between the toroid and the passing charge requires that
$$𝐅_{onq}^{Lorentz}=\frac{q}{c}\left(𝐯_q_q\right)𝐀_\mu (𝐫_q)$$
(66)
While the electric velocity field of a passing charge is screened out of a good conductor, the magnetic field penetrates into a good conductor with a time integral which is independent of the conductivity of the ohmic material of the conductor.B1999 Thus the magnetic field energy $`U_{em}`$ associated with the overlap of the toroid magnetic field and the point charge magnetic fieldB1 is
$$U_{em\mu q}=\frac{1}{8\pi }d^3r\mathrm{\hspace{0.17em}2}𝐁_q𝐁_\mu =q\frac{𝐯_q}{c}𝐀_\mu (𝐫_q)$$
(67)
just what was given for $`U_{em\mu q}`$ in Eq. (45). Let us assume that this magnetic field energy is equal to the change in kinetic energy of the passing charge due to the electric fields from the changing charge and current densities of the magnet. Since the change in magnetic field energy is of order $`1/c^2`$, we need to consider only the nonrelativistic approximation to the passing particle kinetic energy. Then we find
$`{\displaystyle \frac{1}{2}}m_q𝐯_q^2{\displaystyle \frac{1}{2}}m_q𝐯_{q0}^2`$ $`=U_{em\mu q}`$
$`m_q𝐯_{q0}\mathrm{\Delta }𝐯_q`$ $`={\displaystyle \frac{q}{c}}𝐯_{q0}𝐀_\mu (𝐫_q)`$ (68)
where $`𝐯_{q0}`$ is the velocity of the charged particle $`q`$ when far from the magnet where $`𝐀_\mu (𝐫_q)`$ vanishes, and $`\mathrm{\Delta }𝐯_q`$ is the change in the velocity of the passing charge. Thus we find
$$m_q\mathrm{\Delta }𝐯_q=(q/c)𝐀_\mu (𝐫_q)$$
(69)
and the force on the passing charge is therefore
$$𝐅_{onq}=m_qd𝐯_q/dt=m_qd(\mathrm{\Delta }𝐯_q)/dt=(q/c)\left(𝐯_q_q\right)𝐀_\mu (𝐫_q)$$
(70)
exactly as found in Eq. (66) from Newton’s third law. Thus there is a certain consistency between our momentum and energy considerations. However, it should be noted that the kinetic energy change for the passing charge is assumed to be of the same sign as the change in energy of the magnetic field. Energy conservation thus requires that the charges carrying the currents of the toroid must absorb twice the kinetic energy change of the passing charge. If the currents of the toroid act in a fashion analogous to a battery in magnetic systems involving mechanical work, then such an energy balance is consistent with what is found for familiar magnetic systems.flat We note that the energy absorbed by the center of mass motion of the magnet is of order $`1/c^4`$ and hence is negligible, since the recoil velocity of the center of energy of the toroidal magnet (which was initially at rest) is of the order of $`1/c^2`$ from Eq. (65).
One should note the difference in perspectives between the analysis given here in the classical-lag point of view and that suggested by proponents of the no-velocity-change point of view (those who support the quantum topological interpretation of the Aharonov-Bohm phase shift). It was pointed out by Coleman and Van Vleck,CVV and repeated above in Eq. (53), that the accelerations of the centers of energy for the toroid and the passing charge must be related as in Newton’s third law. We have assumed that the electric field of the passing charge is screened out of the magnet, have obtained the force on the passing charge $`q`$ by assuming that it is the third law partner of the usual magnetic Lorentz force on the toroidal magnet, and then have shown that this force is directly related to the energy change in the magnetic fields which penetrate into the magnet. The no-velocity-change point of view claims that there is no force back on the passing charge, that the magnetic moment of the magnet does not change, and that the changing electromagnetic field momentum is associated with ”hidden momentum in magnets” whose change ”cancels” the magnetic Lorentz force on the magnet. This requires that the electric field of the passing charge should penetrate into the magnet so as to give the ”hidden momentum,” a penetration which seems contrary to the screening of electric fields by conductors. Furthermore, this point of view tells us nothing about magnetic energy changes between the passing charge and a toroid.
## VII Discussion
Although the Aharonov-Bohm phase shift is well known and is now standard in all the recent quantum mechanics texts, most physicists seem unaware of the long-standing controversy regarding the interpretation of the phase shift. In 1959, Aharonov and BohmAB solved the Schroedinger equation and predicted their phase shift. The phase shift has been observed experimentally.Chamb Aharonov and Bohm attracted attention to their phase shift by claiming that their predicted phase shift occurred in the absence of classical electromagnetic forces and velocity changes and represented a new quantum topological effect with no analogue in classical theory. There is no experimental evidence for this claim. Indeed, the interpretation has aroused controversy. Most of the initial controversy regarding the Aharonov-Bohm phase shift centered on a distraction, whether or not the shift was caused by stray magnetic fields outside the solenoid or toroid. This aspect of the controversy has been removed by the toroidal experiments of TonomuraT which allow very little stray magnetic flux.
The suggestion that the Aharonov-Bohm phase shift might be based upon a classical lag effect involving classical electromagnetic forces and velocity changes (the suggestion repeated here) depends upon our understanding of classical electromagnetism. The conventional attitude regarding the Aharonov-Bohm phase shift is best stated by Aharonov, Pearle, and Vaidman:APV ”In the Aharonov-Bohm effect it is obvious that the electron is not subject to any electromagnetic force, because the magnetic field lies wholly within the filament and so is zero at the electron’s location.” This naive statement omits the crucial possibility of induced charge or current densities in the magnet leading to forces back on the passing charge. Indeed, induced currents do lead to forces back on passing charges; the phase shifts may well arise from classical lag effects.
In the 1970s, it was suggested that the possible influence of the electromagnetic fields of the passing charge could be removed by surrounding the solenoid or toroid by a conductor which would screen out the electromagnetic fields.pre Experiment showed that the phase shift persisted even when the solenoid was surrounded by a conductor.T However, it was realized that although electric fields are indeed well screened by a conductor, magnetic velocity fields penetrate into an ohmic conductor (and also into superconductors at high frequencies) in a form which is completely different from the skin-depth behavior of wave fields, and indeed there is an invariant time integral which has precisely the correct form to account for the Aharonov-Bohm phase shift as an energy-related lag effect based on classical forces.B1999 The experiments to date do no not remove the possibility of a classical electromagnetic basis for the Aharonov-Bohm phase shift.B5 In addition, it was pointed out that electrostatic forces can give interference pattern shifts which take exactly the same form as the Aharonov-Bohm phase shift.B2 Matteucci and Pozzi confirmed this experimentally in 1985.MP
In 1984, Aharonov and CasherAC suggested a second phase shift, this time for a magnetic moment passing a line charge, which they claimed was the duel of the Aharonov-Bohm phase shift and again occurred in the absence of classical forces and velocity changes. However, it was pointed out that conventional classical electromagnetic theory clearly predicted a force on a passing magnetic moment treated as a current loop, and Newton’s second law suggested a lag effect.B4 To counter this observation, Aharonov, Pearle and VaidmanAPV introduced a new analysis for the interaction of a magnetic moment and a point charge, and claimed that the magnetic moment, although indeed experiencing a net Lorentz force, nevertheless moved as though it experienced no forces whatsoever, because of changes in ”hidden momentum in magnets” cancelling the applied Lorentz force.
For the Aharonov-Bohm phase shift, the Aharonov-Casher phase shift, and the Shockley-James paradox, the heart of the controversy and paradox involves the interaction between a point charge and a magnetic moment through order $`1/c^2.`$ Although the literature of the Aharonov-Bohm phase shift is full of statements about the interaction which claim to exclude any possibility of an explanation based upon classical electromagnetic forcesP , the claims often depend upon nonrelativistic modelsPTT or point to familiar effects, such as aspects of mutual inductance, which indeed will not give the desired behavior,mut but overlook the 0-order forces on the charges of the magnet because the magnet is neutral. Coleman and Van Vleck have treated the interaction consistently relativistically using the Darwin Lagrangian. In the present work, we have followed the Darwin Lagrangian analysis. We have modeled the magnetic moment as a classical hydrogen atom interacting with the passing charge through the Darwin Lagrangian, and have noted particularly the nonrelativistic behavior of the magnetic moment pointed out by Solem. The model is unambiguous in its prediction of classical electromagnetic forces, energies, and changes of the center of energy. It is the 0-order accelerations which cause electric fields in order $`1/c^2`$ which act strongly on the passing charge.
The transition to a multiparticle limit still allows ambiguities. However, the assumption that in this limit the electric fields are screened out of the magnet while the magnetic fields penetrate into the magnet both fits with what is known for ohmic conductors and also allows for a consistent treatment of the conservation laws of relativistic theory. The discussion given here represents a refutation of the suggestions of Aharonov, Pearle, and Vaidman regarding the role of ”hidden momentum in magnets” and confirms the semiclassical calculations of both the Aharonov-Bohm and Aharonov-Casher phase shifts based upon classical lag effects.B3 B4 What is needed now are experiments to test whether or not the Aharonov-Bohm and Aharonov-Casher phase shifts occur in the presence or absence of velocity changes for the passing particles.BCC
### VII.1 Acknowledgement
I wish to thank Professor Joel Gersten for a number of helpful discussions.
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# Laser-cluster interaction: x-ray production by short laser pulses
## Abstract
We investigate the heating of the quasi-free electrons in large rare gas clusters ($`N`$ exceeding $`10^5`$ atoms) by short laser pulses at moderate intensities ($`I10^{15}\mathrm{Wcm}^2`$). We identify elastic large-angle backscattering of electrons at ionic cores in the presence of a laser field as an efficient heating mechanism resembling the Fermi shuttle. Its efficiency as well as the effect of collective electron motion, electron-impact ionization and cluster charging, are studied employing a mean-field classical transport simulation. Results for the absolute x-ray yields are in surprisingly good quantitative agreement with recent experimental results.
The interaction of short and ultra-short intense laser pulses with clusters has become an important area of laser-matter research bridging the gap between gas-phase and solid-state processes krain . The observation of characteristic x-ray emission from laser-irradiated clusters paris1 ; parishci ; paristhese ; parisprl suggested its potential as an x-ray source through a highly non-linear conversion of IR radiation that combines advantages of both solid and gaseous targets. Like solids, clusters provide large x-ray yields, yet unlike solids they are debris-free, just like gaseous targets. Characteristic x-ray emission also provides important time-differential information on the laser-induced electronic dynamics on a femtosecond scale. The charge state as well as the vacancy distribution of the cluster ions at the instant of emission can be extracted from high-resolution x-ray spectra. Moreover, as the vacancy production in deeply bound shells (e.g. K-shell in argon or L-shell in xenon) proceeds via impact ionization by energetic electrons, characteristic x-rays provide an “in situ” thermometer of the temperature of the heated electron gas in the cluster. Recent experiments parishci ; paristhese ; parisprl found an unexpectedly low laser intensity threshold for the production of x-ray radiation. When irradiating large argon clusters with $`N>10^5`$ atoms with laser pulses with a short pulse duration of $`\tau =60\mathrm{fs}`$ at full width half maximum (FWHM), the intensity threshold for the production of characteristic K-x-rays lies at $`I_{th}=2.210^{15}\mathrm{Wcm}^2`$. By comparison, the ponderomotive energy, $`U_P=F^2/(4\omega ^2)`$, of a free electron in a laser field of this intensity is as low as $`U_P130\mathrm{eV}`$, more than an order of magnitude below the binding energy $`E_K3.2\mathrm{keV}`$ of the K-shell electrons. This observation raises puzzling questions as to the efficient heating mechanism for electrons in large clusters at such moderate intensities of very short pulses with $`40`$ optical cycles.
Several theoretical models for intense laser-cluster interaction have been proposed krain ; ditmire ; last1 ; rose ; saalmann ; brabec ; smirnov , none of which appears to provide a satisfactory explanation for such rapid acceleration of electrons. A theoretical description of intense laser-cluster interaction represents a considerable challenge in view of the many-body nature of this process. Molecular dynamic simulations last1 ; rose ; saalmann are limited to about 1000 atoms, and results obtained for small clusters are difficult to scale to larger sizes. The recently proposed microscopic particle in cell (MPIC) method brabec reaches clusters of $`10^4`$ atoms. Larger clusters with $`N>10^5`$ particles appear still not in reach, and quantitative predictions for x-ray emission and inner-shell processes have not yet been attempted. We propose in this letter an efficient heating mechanism of electrons in large clusters that is operational within a few optical cycles and at moderate laser intensities. It is based on the observation that *elastic* large-angle scattering of electrons at cluster atoms (ions) in the simultaneous presence of a laser field provides an efficient route to electron acceleration. Elastic backscattering at the core potential of the ions can flip the velocity vector of an electron, allowing it with non-negligible probability to remain synchronized with the alternating laser field vector (Fig. 1). Consequently, the electron will rapidly accumulate rather than lose momentum during subsequent half cycles. Within a few optical cycles, electrons can thus be accelerated to high kinetic energies well beyond the quiver energy $`E_p=2U_p`$. This heating mechanism resembles the Fermi shuttle acceleration fermi ; burg91 and is also related to the lucky-electron model proposed for IR photoemission from metallic surfaces lucky . The acceleration of a particle by successive backscatterings from a moving and a stationary target in the original proposal by Fermi for acceleration of cosmic particles is here modified such that the cluster ions represent stationary targets, while the alternating force field of the laser plays the role of the moving target. A realistic estimate for the efficiency of this heating mechanism hinges on a proper description of the differential elastic scattering cross sections for electrons, $`\mathrm{d}\sigma _e/\mathrm{d}\theta `$, into backward angles $`\theta 90^{}`$, which are determined by the non-Coulombic short-ranged potentials of the ionic cores. $`\mathrm{d}\sigma _e/\mathrm{d}\theta `$ was calculated for electron scattering at argon ions for all relevant charge states between $`q=1`$ and $`q=16`$ over a wide range of energies using parametrized Hartree-Fock potentials szydlik ; salvat . $`\mathrm{d}\sigma _e/\mathrm{d}\theta `$ is typically dominated by few low-order partial waves giving rise to generalized Ramsauer-Townsend minima mott and diffraction oscillations burg95 (Fig. 2). We assumed for simplicity the electronic ground state occupation for each $`q`$. Extensions to core-excited configurations would be straight forward. For the interstitial region a muffin-tin potential is employed in order to account for solid-state effects salvat . The potential shape in this region has, however, no significant influence on the cross sections at backward angles. The latter exceeds the pure Coulomb case by several orders of magnitude for all charge states and over a wide range of electron energies ($`keV`$). The frequent usage in simulations of unrealistic (softened) Coulomb potentials last1 ; saalmann ; brabec , which grossly underestimates backscattering, is quite likely one reason why this route of electron acceleration has not yet been accounted for. Moreover, this process becomes much more important for large clusters, as the mean-free path for elastic scattering becomes comparable to the cluster size. The important role of realistic core potentials has recently also been identified in the quantum analogue of this process, inverse Bremsstrahlung ($`U_P\mathrm{}\omega `$), for clusters in a vacuum ultraviolet (VUV) laser field santra .
A full ab initio simulation for large clusters $`(N10^5`$ particles) appears still impractical. In the following we present a simplified theoretical description of the electronic ensemble that allows to tackle its short-time dynamics ($`\tau 60\mathrm{fs}`$). It employs a generalization of classical transport theory (CTT burg90 ) for open systems, in which the electronic dynamics is represented by a classical phase-space distribution $`f(𝒓,\dot{𝒓},t)`$ whose evolution is determined by test-particle discretization, i.e. by solving the corresponding Langevin equation for representative trajectories. In the present case, the ensemble consists of $`N_e`$ quasi-free electrons, liberated inside the cluster after ionization of the cluster atoms, represented by $`N_{test}=\alpha N_e`$ particles, where the scaling parameter $`\alpha `$ is limited by computational feasibility. Each test particle is subject to a Langevin equation (atomic units are used unless otherwise stated),
$`\ddot{𝒓}_i`$ $`=`$ $`𝑭_L(t)\mathbf{}V(𝒓_i,t)`$ (1)
$`𝑭_{mean}(𝒓_i,t)+𝑭_{stoc}(𝒓_i,\dot{𝒓}_i,t)`$
with $`i=1,\mathrm{}N_{test}(t)`$. Eq. 1 describes a dynamical system open to both particle number variation, $`N_{test}(t)`$, due to successive ionization events, and energy exchange with the many-particle reservoir (atoms, ions, and electrons) as well as with the laser field taken to be linearly polarized with a temporal envelope
$$𝑭_L(t)=F_0\widehat{𝒛}\mathrm{sin}(\omega t)\mathrm{sin}^2\left(\frac{\pi t}{2\tau }\right).$$
(2)
Eq. (1) provides a computational starting point for treating many-body collisional correlation effects through stochastic forces $`𝑭_{stoc}`$, which can be determined either from independent ab-initio quantum calculations or experimental data burg90 . Forces resulting from static conservative potentials $`V`$ inside the cluster last1 , can be included as well. Because of the coherent motion and high ionization density, the present extension of the CTT (Eq. 1) goes beyond the independent-particle description by including dynamic electron-electron and electron-ion interactions on a time-dependent mean field level. Accordingly, $`𝑭_{mean}`$ depends on the entire ensemble of test-particle coordinates $`\left\{𝒓_i(t)\right\}`$. We therefore propagate the elements of the ensemble by self-consistently coupling a given trajectory to the mean field of other trajectories running in parallel by continuously updating the forces.
Effects of fluctuations on the electronic dynamics can be taken into account through stochastic forces which are determined from Poissonian random processes. For example, electron-ion scattering, electron-impact ionization, and core-hole excitation are determined by the probability per unit pathlength for scattering
$$\lambda _{scatt}^1=\frac{\mathrm{d}P_{scatt}}{\mathrm{d}x}=\sigma _{scatt}(q,E)\rho (t),$$
(3)
controlled by the energy ($`E`$) and charge state ($`q`$) dependent integral cross-section for this process, $`\sigma _{scatt}(q,E)`$, and the instantaneous ionic target density $`\rho (t)`$ of a given charge state. Each stochastic scattering process results in “jumps” (classical trajectory jumps and jumps in occupation) at discrete times. A jump in momentum, $`\mathrm{\Delta }\dot{𝒓}`$, signifies elastic scattering determined by the differential cross section, $`\mathrm{d}\sigma _e(q,E,\theta )/\mathrm{d}\theta `$, a simultaneous jump in test-particle number, $`\mathrm{\Delta }N_{test}`$, represents ionizing collisions, and a simultaneous jump in the number of inner-shell vacancies, $`\mathrm{\Delta }N_K`$, results from core-exciting collisions. The key point is that the necessary input data, $`\sigma _{scatt}(q,E)`$, can be determined and tabulated independently from the simulation at any desired level of sophistication. In the present simulation, the electron-impact ionization cross sections are determined from a modified Lotz formula paristhese ; lotz $`\sigma _i(q,E)=A_q^{}\mathrm{ln}(E/W_q^{})/(EW_q^{})`$ (for $`EW_q^{}`$), where the empirical parameters $`A_q^{}`$ and $`W_q^{}`$ were obtained by a fit to experimental ion-atom collision data paristhese ; zhang . The cluster-specific effects of suppression of the work function can be incorporated by modifying the effective work function $`W_q^{}`$ in the presence of nearby ionic cores.
For the mean field we perform a multipole expansion keeping only the monopole and dipole terms. The monopole term is given by $`𝑭_{mean}^{(0)}(\stackrel{}{r},t)=Q(r,t)𝒓/r^3`$ where $`Q(r,t)`$ is the instantaneous charge of the sphere of radius $`r`$ resulting from the displacement of the ensemble of test particles relative to the ionic background. Analogously, the dipole field inside the cluster ($`r<R(t)`$) is $`𝑭_{mean}^{(1)}(t)=𝒑(t)/R(t)^3`$, while outside it is that of a central dipole. The dipole moment, $`𝒑`$ is determined by $`𝒑(t)\frac{1}{\alpha }_i𝒓_i`$, where the sum extends over the subset of test particles with $`r_i<R(t)`$.
As the ionic and electronic dynamics proceed on different time scales, the onset of cluster expansion can be taken into account through the parametric variation of the radius $`R(t)`$ of the uniform spherical charge background representing the ions of mass $`M`$ in their time-dependent mean charge state $`q(t)`$:
$$M\frac{\mathrm{d}^2R(t)}{\mathrm{d}t^2}=\frac{q(t)Q(R,t)}{R^2(t)}.$$
(4)
We solve (Eq. 1) for a cluster with $`N=2.8\times 10^5`$ argon atoms with initial atomic number density $`\rho (t=0)=2.66\times 10^{22}\mathrm{cm}^3`$ and initial radius $`R(0)=258\mathrm{a}.\mathrm{u}.`$ irradiated by a laser pulse of length (FWHM) $`\tau =60fs`$, wavelength $`\lambda =\mathrm{\hspace{0.17em}800}\mathrm{nm}`$ and peak intensities $`I10^{15}\mathrm{Wcm}^2`$. The ponderomotive energy of the electrons is of the order of $`U_P70\mathrm{eV}2.5\mathrm{a}.\mathrm{u}.`$ As the laser field reaches for the first time the threshold field for over-barrier ionization $`F_L^1(t_1)0.08`$, the first $`N_{test}(t_1)`$ test particles with zero velocity randomly distributed over the cluster provide initial conditions for the propagation of Eq. (1). In the present case $`N_{test}=0.1N_e`$ (or $`\alpha =0.1`$). Contributions from tunneling ionization can be included but are in the present case negligible. The test-particle number subsequently increases by further ionization events (Fig. 3(a)). Impact ionization by electrons meanwhile accelerated is highly effective, making further field ionization events unlikely. The mean charge state rapidly increases to $`7`$, which most likely still underestimates the ionization efficiency, as non-radiative core-hole relaxation and enhanced ionization by suppression of the work function by ion proximity are not yet included. The Coulomb expansion of the cluster sets in slowly due to the large inertia of the ions. Even after the laser pulse is switched off ($`t2\tau `$), the cluster has expanded by less than a factor 2. Our simulation shows that the charge resulting from electrons leaving the cluster (Fig. 3(b)) is concentrated on the surface of the cluster, the ions in the inside of the cluster being well shielded by the quasi-free electrons, in agreement with the MPIC simulation brabec . After the first ionization burst the electronic plasma frequency is given by $`\omega _p^2=N_{test}(t_1)/(\alpha R(t_1)^3)=\rho (t_1)4\pi /30.016\mathrm{a}.\mathrm{u}.5\omega ^2`$. As electron-impact ionization produces more quasi-free electrons, $`\omega _p`$ grows rapidly before diminishing again as the cluster expansion sets in (Fig. 3(c)). During the evolution, the effective field inside the cluster consisting of both the laser and the polarization field is approximately given by
$`𝑭_{eff}(t)`$ $``$ $`Re\{{\displaystyle _{\omega \mathrm{\Delta }\omega }^{\omega +\mathrm{\Delta }\omega }}\stackrel{\mathbf{~}}{𝑭}_L(\omega ^{})`$
$`\times (1{\displaystyle \frac{\omega _p^2}{\omega _p^2\omega ^2i\omega ^{}\gamma }})e^{i\omega ^{}t}\mathrm{d}\omega ^{}\},`$
where $`\mathrm{\Delta }\omega `$ is the Fourier width caused by the temporal profile of the pulse (Eq. 2) and $`\gamma `$ stands for the damping due to scattering events. $`F_{eff}`$ is significantly reduced compared to the bare laser field (Fig. 3(d)) due to the combined effect of collective electron motion and electron-impact ionization. A significant resonant enhancement ditmire ; saalmann is absent. Ref. brabec suggested laser dephasing heating (LDH), in which the phase shift between dipole moment and laser field caused by the macroscopic electric field inside the cluster, would lead to a net electron energy absorption $`𝑱𝑭_{eff}𝑑t`$, where $`𝑱`$ stands for the electron current. However, evaluating this integral using the simulated current and field shows that heating by LDH is not sufficient in the present case.
The efficiency of heating by elastic electron-ion scattering is directly reflected in the simulated absolute x-ray yields (Fig. 4). The latter are determined by the number of K-shell vacancies created, corrected for the mean fluorescence yield $`\eta `$ taken to be $`\eta 0.12`$ bhalla for argon with partially filled L-shell but empty M-shell. It should be emphasized that the simulation contains no freely adjustable parameter. To compare the simulation results to the experiments, an ensemble average over the spatial intensity profile of the laser beam, which is Gaussian to a good degree of approximation paristhese , is performed. To quantify the significance of the Fermi-shuttle acceleration by repeated elastic backscattering, we performed an otherwise identical simulation with elastic electron-ion scattering switched off. In this case, for $`I1.5\times 10^{15}\mathrm{Wcm}^2`$ a small fraction of quasi-free electrons gains sufficient energy to produce K-shell vacancies. Their mean kinetic energy can be estimated from the potential energy of the charged-up cluster with charge $`Q(t)`$ (i.e. the monopole term of the mean effective field). However, including elastic electron-ion scattering drastically increases the x-ray yield by a factor 3 to 6. We then find surprisingly close agreement with the experimental results (Fig. 4).
In summary, we have analyzed the heating of the quasi-free electrons in large rare-gas clusters ($`N10^5`$ atoms) at moderate laser intensities ($`I=10^{15}10^{16}\mathrm{Wcm}^2`$). We have identified a novel, highly efficient electron heating mechanism operative at short times within a few optical cycles in terms of elastic large-angle scattering resembling the Fermi shuttle. Other processes such as heating in a plasma-resonant field are found to be less effective. In particular, the polarization of the cluster leads to a reduction rather than an enhancement of the effective field. While the surprisingly good quantitative agreement on an absolute scale with experimental data may be, in part, fortuitous, the importance of this route to fast electron acceleration appears unambiguously established.
###### Acknowledgements.
Work is supported by FWF SFB-16 (Austria).
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# Coordinate Systems and Bounded Isomorphisms for Triangular Algebras
## 1. Introduction
This paper presents the results of our study of bounded isomorphisms of coordinatized (nonselfadjoint) operator algebras. Isometric isomorphisms have been extensively studied (see, for example, ) and are quite natural, as they include restrictions of $``$-isomorphisms. Isometric isomorphisms of $`C^{}`$-algebras preserve adjoints. Bounded isomorphisms, in contrast, need not preserve adjoints or map $``$-subalgebras to $``$-subalgebras. Nonetheless, we obtain structural results, most notably, that bounded isomorphisms are completely bounded and that they factor into diagonal-fixing, spatial, and isometric parts, analogous to Arveson-Josephson’s factorization of bounded isomorphisms of analytic crossed products.
Coordinates have been used in the categories of $`W^{}`$-algebras and $`C^{}`$-algebras with $``$-homomorphisms for decades, going back at least to the work of Feldman and Moore, for von Neumann algebras and Renault’s construction of $`C^{}`$-algebras for a wide range of topological groupoids . Of particular interest to us is Kumjian’s construction, in , of a certain $`𝕋`$-groupoid which he called a twist, which he showed is a classifying invariant for a diagonal pair, a separable $`C^{}`$-algebra with a distinguished MASA satisfying various properties. Renault describes twists in terms of the dual groupoid and this perspective is often helpful, see for example the work of Thomsen in . Twists have been used by various authors, most notably Muhly, Qiu, and Solel, and Muhly and Solel, to study varied categories of subalgebras and submodules of groupoid $`C^{}`$-algebras with isometric morphisms .
To apply coordinate methods to bounded isomorphisms, we found it necessary to revisit these coordinate constructions, eliminating as much as possible the use of adjoints and clarifying the role of the extension property. We define coordinates for bimodules over an abelian $`C^{}`$-algebra which are intrinsic to the bimodule structure and not a priori closely tied to the $``$-structure. These definitions allow us to simplify and extend some of the structural results in the literature. In particular, we obtain a number of results for algebras containing the abelian $`C^{}`$-algebra: e.g. isomorphism of the coordinates is equivalent to diagonal-preserving isomorphism of the algebras. This analysis of coordinate systems will be useful, we expect, in applying coordinate constructions to more general settings.
Bounded isomorphisms play a role in the study of norm-closed operator algebras which is parallel to similarity transforms in the study of weakly-closed subalgebras of $`()`$. During the mid-1980’s and early 1990’s, there was considerable interest in the structural analysis of such algebras via their similarity theory. This was particularly successful with the class of nest algebras (see, for example, ) and, to a lesser degree, the CSL algebras. Interestingly, by using certain faithful representations of $`C^{}`$-diagonals, we can employ similarity theory for atomic CSL algebras to obtain structural results for bounded isomorphisms between triangular subalgebras of $`C^{}`$-diagonals.
We turn now to a more detailed outline of the paper. Throughout the paper, we consider bimodules over an abelian unital $`C^{}`$-algebra $`𝒟`$. Our view is that the set of coordinates for such a bimodule $``$ is the collection $`^1()`$ of norm one elements of the Banach space dual $`^\mathrm{\#}`$ which are eigenvectors for the bimodule action of $`𝒟`$ on $`^\mathrm{\#}`$. We use eigenfunctional for such elements, and we topologize $`^1()`$ using the weak-$``$ topology. With this structure, we call $`^1()`$ the coordinate system for $``$. In Section 2, we establish some very general, but useful, properties. For example, Theorem 2.6 shows that when $`_1`$ is a submodule of $``$, then an element of $`^1(_1)`$ can be extended (but not necessarily uniquely) to an element of $`^1()`$.
A key step in minimizing the use of adjoints is replacing normalizers with intertwiners, that is, elements $`m`$ so that $`m𝒟=𝒟m`$. Section 3 starts by showing that intertwiners and normalizers are closely related, at least when $`𝒟`$ is a MASA in a unital $`C^{}`$-algebra $`𝒞`$ containing $``$ (Propositions 3.2 and 3.3). The extension property is used to construct intertwiners and to slightly strengthen a key technical result of Kumjian. In this generality, $`^1()`$ need not separate points of $``$.
In Section 4 we work in the context of $`𝒟`$-bimodules $`𝒞`$, where $`𝒞`$ is a unital $`C^{}`$-algebra and $`𝒟𝒞`$ is a regular MASA with the extension property. Such a pair $`(𝒞,𝒟)`$ we term a regular $`C^{}`$-inclusion. Here the coordinates are better behaved: eigenfunctionals on submodules of $``$ extend uniquely to $``$. Nevertheless, the coordinates for regular $`C^{}`$-inclusions again are not sufficiently rich to separate points. However, the failure to separate points is intimately related to a certain ideal of $`𝔑𝒞`$, and Theorem 4.9 shows that the quotient of $`𝒞`$ by $`𝔑`$ is a $`C^{}`$-diagonal, which is a mild generalization of Kumjian’s notion of diagonal pair due to Renault . Essentially, a $`C^{}`$-diagonal is a regular $`C^{}`$-inclusion where the conditional expectation $`E:𝒞𝒟`$ arising from the extension property is faithful. There are an abundance of $`C^{}`$-diagonals: crossed products of abelian $`C^{}`$ algebras by freely acting amenable groups are $`C^{}`$-diagonals, and Theorem 4.24 shows that inductive limits of $`C^{}`$-diagonals are again $`C^{}`$-diagonals when the connecting maps satisfy a regularity condition. In particular, AF-algebras and circle algebras can be viewed as $`C^{}`$-diagonals. While our primary interest in this paper is the use of coordinate methods to study nonselfadjoint algebras, Theorem 4.9 and Theorem 4.24 are examples of results in the theory of $`C^{}`$-algebras obtained using our perspective.
When $``$ is a $`𝒟`$-bimodule contained in a $`C^{}`$-diagonal, the elements of $`^1()`$ do separate points, and when the bimodule is an algebra, also have a continuous product. In Section 4, we use the extension property to show that the coordinate system $`^1(𝒞)`$ for a $`C^{}`$-diagonal agrees with Kumjian’s twist. Our methods provide some simplifications and generalizations of Kumjian’s results. One of the interesting features of our approach is that it allows us to show that the coordinate systems for bimodules $`𝒞`$ are intrinsic to the bimodule alone, and not dependent on the choice of the embedding into the particular $`C^{}`$-diagonal. This, and the agreement of our construction with Kumjian’s is achieved in Theorems 4.14, Corollary 4.16, and Proposition 4.18. An interesting application of our coordinate methods is Theorem 4.22, which shows that if the pair $`(𝒞,𝒟)`$ is a $`C^{}`$-diagonal, and $`𝒜`$ is a norm closed algebra with $`𝒟𝒜𝒞`$, then the $`C^{}`$-envelope of $`𝒜`$ coincides with the $`C^{}`$-subalgebra of $`𝒞`$ generated by $`𝒜`$.
We study the representation theory of $`C^{}`$-diagonals in Section 5, obtaining a faithful atomic representation compatible with the $`C^{}`$-diagonal structure. Our methods are reminiscent of Gardner’s work on isomorphisms of $`C^{}`$-algebras in . The significance to us of these representations is that they carry subalgebras containing the diagonal to algebras weakly dense in a CSL-algebra, see Theorem 5.9. This theorem enables us to prove, with fewer hypotheses than previously needed, that the conditional expectation is multiplicative when restricted to a triangular subalgebra, Theorem 5.10.
Section 6 considers diagonal-preserving bounded isomorphisms, those that map the diagonal of one algebra onto the diagonal of the other. Theorem 6.3 shows that in this case, there is an isomorphism between coordinate systems arising naturally from the algebra isomorphism. Consequently, we are able to prove several results, such as Theorem 6.7, which shows an automorphism of a triangular algebra which fixes the diagonal pointwise arises from a cocycle, and Theorems 6.11 and 6.14, which show that coordinates are invariant under diagonal-preserving bounded isomorphisms, extending previous results for isometric isomorphisms.
We then turn to bounded isomorphisms of triangular algebras which do not preserve the diagonal. A main result, Theorem 7.7 shows that a given bounded isomorphism of triangular subalgebras induces an algebraic isomorphism $`\gamma `$ of their coordinate systems, but our methods are not strong enough to show that this isomorphism is continuous everywhere. Nevertheless, this does show that the algebraic structure of coordinate systems for triangular algebras is invariant under bounded isomorphism, a fact we believe would be difficult to show using previously existing methods.
While we are not able to prove continuity of the map $`\gamma `$ on coordinate systems in general, we can prove it in various special cases. A bounded isomorphism of triangular algebras induces a canonical $``$-isomorphism of the diagonals. If this $``$-isomorphism extends to a $``$-isomorphism of the $`C^{}`$-envelopes, then Corollary 7.8 shows that the product of $`\gamma `$ with an appropriate cocycle yields a continuous isomorphism. Also, Theorem 7.11 shows that bounded isomorphism of triangular algebras implies an isomorphism of their coordinate systems when the triangular algebras are generated by their algebra-preserving normalizers. This is a new class of algebras which contains a variety of known classes, including those generated by order-preserving normalizers or those generated by monotone $`G`$-sets.
Another of our main results, Theorem 8.2, shows that if boundedly isomorphic triangular subalgebras are represented in the faithful representation constructed in Theorem 5.9, then the isomorphism extends to an isomorphism of the weak closures of the triangular algebras. A crucial ingredient in proving this is Theorem 7.7. Theorem 8.2 allows us to use known results about isomorphisms of CSL algebras to prove another main result, Theorem 8.8, which asserts that every bounded isomorphism of triangular subalgebras of $`C^{}`$-diagonals is completely bounded. Consequently, we are able to prove that every isometric isomorphism of triangular algebras extends to an isometric isomorphism of the corresponding $`C^{}`$-envelopes. Another application of Theorem 8.2 is Theorem 8.7, which gives the factorization into diagonal-fixing, spatial, and isometric parts mentioned earlier.
Together with Power, we asserted \[12, Theorem 4.1\] that two limit algebras are isomorphic if and only if a certain type of coordinates for the algebras, namely their spectra, are isomorphic. Unfortunately, there is a serious gap in the proof, and another of our motivations for the work in this paper was an attempt to provide a correct proof. While we have not yet done this, our results provide evidence that \[12, Theorem 4.1\] is true. The main result of Section 9, Theorem 9.9, shows that an (algebraic) isomorphism of a limit algebra $`𝒜_1`$ onto another limit algebra $`𝒜_2`$ implies the existence of a $``$-isomorphism $`\tau `$ of their $`C^{}`$-envelopes. If we could choose $`\tau `$ so that $`\tau (𝒜_1)=𝒜_2`$, then \[12, Theorem 4.1\] would follow easily, but we do not know this. However, any isomorphism of $`𝒜_1`$ onto $`𝒜_2`$ induces a $``$-isomorphism $`\alpha `$ of $`𝒜_1𝒜_1^{}`$ onto $`𝒜_2𝒜_2^{}`$. If $`\tau `$ can be chosen so that $`\tau `$ extends $`\alpha `$, Corollary 7.8 shows that $`\tau `$ carries $`𝒜_1`$ onto $`𝒜_2`$, and moreover, there is an isomorphism of the corresponding coordinate systems. It is somewhat encouraging that Theorem 9.9 shows that there is no $`K`$-theoretic obstruction to the existence of a $`\tau `$ which extends $`\alpha `$. This is as close as we have been able to come in our efforts to provide a correct proof of \[12, Theorem 4.1\].
## 2. Intertwiners and Eigenfunctionals
In this section we provide a very general discussion of coordinates. Although our focus in this paper is on $`C^{}`$-diagonals and regular $`C^{}`$-inclusions, defined in Section 4, we start in a more general framework, with a view to extending coordinate methods beyond our focus here. Indeed, there are several useful general results, most notably Theorem 2.6, which shows that eigenfunctionals can be extended from one bimodule to another bimodule containing the first. From this result, we characterize when an eigenfunctional with a given range and source exists, using a suitable seminorm.
Notational Convention. Given a Banach space $`X`$, we denote its Banach space dual by $`X^\mathrm{\#}`$, to minimize confusion with adjoints.
Throughout this section, $`𝒟`$ will be a unital, abelian $`C^{}`$-algebra, and $``$ will be a Banach space which is also a bounded $`𝒟`$-bimodule, that is, there exists a constant $`K>0`$ such that for every $`d,f𝒟`$ and $`m`$,
$$dmfKdmf.$$
As usual, $`^\mathrm{\#}`$ becomes a Banach $`𝒟`$-module with the action,
$$m,f\varphi d=dmf,\varphi d,f𝒟,m,\text{ and }\varphi ^\mathrm{\#}.$$
###### Definition 2.1.
A nonzero element $`m`$ is a $`𝒟`$-intertwiner, or more simply, an intertwiner if
$$m𝒟=𝒟m.$$
If $`m`$ is an intertwiner such that for every $`d𝒟`$, $`dmm`$, we call $`m`$ a minimal intertwiner.
A minimal intertwiner of $`^\mathrm{\#}`$ will be called a eigenfunctional; when necessary for clarity, we use $`𝒟`$-eigenfunctional. That is, an eigenfunctional is a nonzero linear functional $`\varphi :`$ so that, for all $`d𝒟`$, $`x\varphi (dx)`$, $`x\varphi (xd)`$ are multiples of $`\varphi `$.
Denote the set of all $`𝒟`$-eigenfunctionals by $`_𝒟()`$ (or $`()`$, if the context is clear). We equip $`()`$ with the relative weak topology (i.e. the relative $`\sigma (^\mathrm{\#},)`$-topology).
Denote the set of all norm-one $`𝒟`$-eigenfunctionals by $`_𝒟^1()`$ or $`^1()`$.
Given an eigenfunctional $`\varphi _𝒟()`$, the associativity of the maps $`d𝒟d\varphi `$ and $`d𝒟\varphi d`$ yields the existence of unique multiplicative linear functionals $`s(\varphi )`$ and $`r(\varphi )`$ on $`𝒟`$ satisfying $`s(\varphi )(d)\varphi =d\varphi `$ and $`r(\varphi )(d)\varphi =\varphi d`$, that is,
(1)
$$\varphi (xd)=\varphi (x)\left[s(\varphi )(d)\right],\varphi (dx)=\left[r(\varphi )(d)\right]\varphi (x).$$
###### Definition 2.2.
We call $`s(\varphi )`$ and $`r(\varphi )`$ the source and range of $`\varphi `$ respectively.
There is a natural action of the nonzero complex numbers on $`()`$, sending $`(\lambda ,\varphi )`$ to the functional $`m\lambda \varphi (m)`$; clearly $`s(\lambda \varphi )=s(\varphi )`$ and $`r(\lambda \varphi )=r(\varphi )`$.
We next record a few basic properties of eigenfunctionals.
###### Proposition 2.3.
With the weak-topology, $`()\{0\}`$ is closed. Further, $`r:()\widehat{𝒟}`$ and $`s:()\widehat{𝒟}`$ are continuous.
###### Proof.
Suppose $`(\varphi _\lambda )_{\lambda \mathrm{\Lambda }}`$ is a net in $`()\{0\}`$ and $`\varphi _\lambda \stackrel{\text{w}}{}\varphi ^\mathrm{\#}`$.
If $`\varphi =0`$, there is nothing to do, so we assume that $`\varphi 0`$. Choose $`m`$ such that $`\varphi (m)0`$. Then $`\varphi _\lambda (m)\varphi (m)`$, so there exists $`\lambda _0\mathrm{\Lambda }`$ such that $`\varphi _\lambda (m)0`$ for every $`\lambda \lambda _0`$. For any $`d𝒟`$, and $`\lambda \lambda _0`$, we have
$$s(\varphi _\lambda )(d)=\frac{\varphi _\lambda (md)}{\varphi _\lambda (m)}\frac{\varphi (md)}{\varphi (m)}$$
Thus, $`s(\varphi _\lambda )`$ converges weak-$``$ to the functional $`\rho `$ in $`\widehat{𝒟}`$ given by $`\rho (d)=\varphi (md)/\varphi (m)`$. Similarly, $`r(\varphi _\lambda )`$ converges weak-$``$ to $`\tau \widehat{𝒟}`$ given by $`\tau (d)=\varphi (dm)/\varphi (m).`$ It now follows that $`\varphi `$ is a eigenfunctional with $`r(\varphi )=\tau `$ and $`s(\varphi )=\rho `$. Thus, $`()\{0\}`$ is closed.
In particular, if $`\varphi _\lambda \varphi ()`$, then $`s(\varphi _\lambda )s(\varphi )`$ and similarly for the ranges, so $`s`$ and $`r`$ are continuous. ∎
To be useful, there should be many eigenfunctionals. This need not occur for arbitrary bimodules, as Example 2 shows. However, for bimodules of $`C^{}`$-diagonals, which are the bimodules of principal interest in the present paper, Proposition 4.21 below will show that eigenfunctionals exist in abundance.
We need the following seminorm to extend eigenfunctionals and to characterize when eigenfunctionals exist.
###### Definition 2.4.
For $`\sigma ,\rho \widehat{D}`$, define $`B_{\sigma ,\rho }:`$ by
$$B_{\sigma ,\rho }(m):=inf\{dmf:d,f𝒟,\sigma (d)=\rho (f)=1\}.$$
These infima do not increase if we restrict to elements $`d`$ or $`f`$ of norm one. Indeed, for any elements $`d`$ and $`f`$ as above, since $`|\sigma (d)|d`$ and $`|\rho (f)|f`$, we can replace them with $`d/d`$ and $`f/f`$ and this will only decrease the norm of $`dmf`$. Thus,
$$B_{\sigma ,\rho }(m)=inf\{dmf:d,f𝒟,\sigma (d)=\rho (f)=1=d=f\}.$$
In particular, we have $`B_{\sigma ,\rho }(m)m`$.
A variant of this seminorm was used by Steve Power in to distinguish families of limit algebras associated to singular MASAs.
###### Proposition 2.5.
For $`\sigma ,\rho \widehat{𝒟}`$, we have the following:
1. $`B_{\sigma ,\rho }`$ is a seminorm.
2. For $`m`$ and $`\varphi ()`$, $`|\varphi (m)|\varphi B_{r(\varphi ),s(\varphi )}(m).`$
3. For $`m`$, $`d,f𝒟`$, $`B_{\sigma ,\rho }(dmf)=|\sigma (d)|B_{\sigma ,\rho }(m)|\rho (f)|`$.
4. If $`f^\mathrm{\#}\{0\}`$ satisfies $`|f(m)|B_{\sigma ,\rho }(m)`$ for all $`m`$, then $`f()`$ with $`s(f)=\rho `$, $`r(f)=\sigma `$, and $`f1`$.
###### Proof.
For (1), it is immediate that $`B_{\sigma ,\rho }(\lambda m)=|\lambda |B_{\sigma ,\rho }(m)`$, for $`\lambda `$ and $`m`$.
To show subadditivity, let $`a,b`$ and choose $`\epsilon >0`$. Pick norm one elements $`d_1`$, $`d_2`$, $`f_1`$, $`f_2`$ of $`𝒟`$, satisfying $`\sigma (d_1)=\sigma (d_2)=1=\rho (f_1)=\rho (f_2)`$, $`d_1af_1<B_{\sigma ,\rho }(a)+\epsilon `$ and $`d_2bf_2<B_{\sigma ,\rho }(b)+\epsilon `$. Then
$`d_1d_2(a+b)f_1f_2`$ $`d_2d_1af_1f_2+d_1d_2bf_2f_1`$
$`d_1af_1+d_2bf_2`$
$`<B_{\sigma ,\rho }(a)+B_{\sigma ,\rho }(b)+2\epsilon ,`$
whence $`B_{\sigma ,\rho }(a+b)B_{\sigma ,\rho }(a)+B_{\sigma ,\rho }(b)`$.
For (2), suppose $`d,f𝒟`$ with $`s(\varphi )(d)=r(\varphi )(f)=1`$. Then $`|\varphi (m)|=|\varphi (dmf)|\varphi dmf`$. Taking the infimum over all such $`d`$ and $`f`$ gives the inequality.
For (3), we show that $`B_{\sigma ,\rho }(dm)=|\sigma (d)|B_{\sigma ,\rho }(m)`$; the proof for $`\rho `$ and $`f`$ is similar.
For any $`a,b𝒟`$ with $`a=b=1=\sigma (a)=\rho (b)`$ we have,
$`admb`$ $`\sigma (d)amb+(ad\sigma (d)a)mb`$
$`|\sigma (d)|amb+ad\sigma (d)am.`$
Hence $`B_{\sigma ,\rho }(dm)|\sigma (d)|amb+ad\sigma (d)am.`$ The definition of $`B_{\sigma ,\rho }`$ and the fact that $`inf\{ad\sigma (d)a:a𝒟,a=1=\sigma (d)\}=0`$, imply that given $`\epsilon >0`$, we may find norm one elements $`a_1,b_1`$ and $`a_2`$ of $`𝒟`$ such that $`\sigma (a_1)=\sigma (a_2)=\rho (b_1)=1`$ and
$$a_1mb_1<B_{\sigma ,\rho }(m)+\epsilon \text{and}a_2d\sigma (d)a_2m<\epsilon .$$
Then
$`B_{\sigma ,\rho }(dm)`$ $`|\sigma (d)|a_1a_2mb_1+a_1a_2d\sigma (d)a_1a_2m`$
$`|\sigma (d)|a_1mb_1+a_2d\sigma (d)a_2m`$
$`|\sigma (d)|(B_{\sigma ,\rho }(m)+\epsilon )+\epsilon ,`$
whence $`B_{\sigma ,\rho }(dm)|\sigma (d)|B_{\sigma ,\rho }(m).`$
To obtain $`|\sigma (d)|B_{\sigma ,\rho }(m)B_{\sigma ,\rho }(dm)`$, observe that for $`a,b𝒟`$ and $`a=b=1=\sigma (a)=\rho (b)`$,
$$|\sigma (d)|B_{\sigma ,\rho }(m)\sigma (d)amb\sigma (d)aadm+admb,$$
and argue as above.
Finally, suppose $`f`$ is a nonzero linear functional on $``$ satisfying $`|f(m)|B_{\sigma ,\rho }(m)`$ for every $`m`$. As $`B_{\sigma ,\rho }(x)x`$, we see that $`f`$ is bounded and $`f1`$.
Suppose $`d𝒟`$, and let $`k=d\rho (d)I`$. Clearly $`\rho (k)=0`$, and so, by for $`x`$, $`|f(xk)|B_{\rho ,\sigma }(xk)=0`$. Therefore $`f(xd)=f(x)\rho (d)+f(xk)=f(x)\rho (d)`$. Similarly, $`f(dx)=\sigma (d)f(x)`$. Thus, $`f`$ is an eigenfunctional with range $`\sigma `$ and source $`\rho `$. ∎
###### Theorem 2.6.
Suppose $``$ is a norm-closed $`𝒟`$-bimodule and $`𝒩`$ is a norm-closed sub-bimodule. Given $`\varphi (𝒩)`$, there is $`\psi ()`$ with $`\psi |_𝒩=\varphi `$ and $`\varphi =\psi `$.
Necessarily, $`\psi `$ has the same range and source as $`\varphi `$.
###### Proof.
Let $`\rho =s(\varphi )`$ and $`\sigma =r(\varphi )`$. From Proposition 2.5 (2), $`|\varphi (n)|\varphi B_{\sigma ,\rho }(n)`$ for all $`n𝒩`$. By the Hahn-Banach Theorem, there exists an extension of $`\varphi `$ to a linear functional $`\psi `$ on $``$ satisfying $`|\psi (x)|\varphi B_{\sigma ,\rho }(x)`$ for all $`x`$. Now apply the last part of Proposition 2.5. ∎
We would like to be able to say that the extension in Theorem 2.6 is unique, but this need not be true. For example, if $`_1=𝒟`$ and $`_2=𝒞`$, then we are considering extensions of pure states, which need not be unique (see , for example). However, in the context of regular $`C^{}`$-inclusions the extension is unique, as we show in Section 4.
We now characterize the existence of eigenfunctionals in terms of the $`B_{\sigma ,\rho }`$ seminorms.
###### Theorem 2.7.
Suppose $`\sigma ,\rho \widehat{𝒟}`$. There is $`\varphi ()`$ with $`r(\varphi )=\sigma `$ and $`s(\varphi )=\rho `$ if and only if $`B_{\sigma ,\rho }|_{}0.`$
###### Proof.
If $`\varphi ()`$ with $`s(\varphi )=\rho `$ and $`r(\varphi )=\sigma `$, then Proposition 2.5 (2) implies $`B_{\sigma ,\rho }|_{}0`$.
Conversely, suppose $`B_{\sigma ,\rho }(m)0`$. Define a linear functional $`f`$ on $`m`$ by $`f(\lambda m)=\lambda B_{\sigma ,\rho }(m)`$. Now use the Hahn-Banach Theorem to extend $`f`$ to a linear functional on all of $``$ which satisfies $`|f(m)|B_{\sigma ,\rho }(m)`$, and apply Proposition 2.5. ∎
Example 2.8. Here is an example where $`()`$ is empty. In $`(L^2[0,1])`$, let $`𝒟`$ be the operators of multiplication by elements of $`C[0,1]`$, and let $``$ be the compact operators. If $`\sigma [0,1]`$ and $`\mathrm{\Lambda }=\{d𝒟:\widehat{d}(\sigma )=1\text{ and }0dI\}`$, then $`\mathrm{\Lambda }`$ becomes a directed set with the direction $`de`$ if and only if $`de0`$. Viewing $`\mathrm{\Lambda }`$ as a net indexed by itself, then $`\mathrm{\Lambda }`$ is a bounded net converging strongly to zero. Hence given any compact operator $`K`$, the net $`\{dK\}_{d\mathrm{\Lambda }}`$ converges to zero in norm. It follows that $`B_{\sigma ,\rho }(K)=0`$ for all $`\sigma ,\rho \widehat{𝒟}`$, so by Theorem 2.7, the set of eigenfunctionals is $`\{0\}`$.
Not surprisingly, eigenfunctionals behave appropriately under bimodule maps.
For $`i=1,2`$, let $`_i`$ be $`𝒟`$-bimodules and let $`\theta :_1_2`$ be a bounded $`𝒟`$-bimodule map. Recall the Banach adjoint map $`\theta ^\mathrm{\#}:_2^\mathrm{\#}_1^\mathrm{\#}`$, given by $`\theta ^\mathrm{\#}(\varphi )=\varphi \theta `$. If $`\varphi (_2)`$, then $`\theta ^\mathrm{\#}\varphi (_1)`$, and we have $`s(\varphi \theta )=s(\varphi )`$, and $`r(\varphi \theta )=r(\varphi )`$.
We include the following simple result for reference purposes; the proof is left to the reader.
###### Proposition 2.9.
For $`i=1,2`$, let $`_i`$ be $`𝒟`$-bimodules and suppose $`\theta :_1_2`$ is a bounded linear map which is also a $`𝒟`$-bimodule map. Then $`\theta ^\mathrm{\#}|_{(_2)}`$ is a continuous map of $`(_2)`$ into $`(_1)\{0\}`$.
If $`\theta `$ is bijective, then $`\theta ^\mathrm{\#}|_{(_2)}`$ is a homeomorphism of $`(_2)`$ onto $`(_1)`$. If $`\theta `$ is isometric, $`\theta ^\mathrm{\#}|_{^1(_2)}`$ is a homeomorphism of $`^1(_2)`$ onto $`^1(_1)`$.
## 3. Normalizers and the Extension Property
We begin this section with a discussion of the relationship between normalizers and intertwiners. Together, Propositions 3.2 and 3.3 show that all intertwiners of a maximal abelian $`C^{}`$-algebra are normalizers and every normalizer can be approximated as closely as desired by intertwiners. This shows that for our purposes, there is no disadvantage in using intertwiners instead of normalizers; moreover, the fact that intertwiners behave well under bounded isomorphism is a considerable advantage.
In the second part of the section, we consider a $`C^{}`$-algebra $`𝒞`$ containing a MASA $`𝒟`$ which has the extension property (see Definition 3.5) and a $`𝒟`$-bimodule $`𝒞`$. We use the extension property and a Theorem from to construct intertwiners in $``$ from intertwiners in $`𝒞`$, and strengthen a key technical result of Kumjian, Proposition 3.10. We see these results as a step towards extending coordinates from $`C^{}`$-diagonals to more general settings.
Context for Section 3. Throughout this section, $`𝒞`$ will be a unital $`C^{}`$-algebra and $`𝒟𝒞`$ will be an abelian $`C^{}`$ subalgebra containing the unit of $`𝒞`$. Bimodules considered in this section will be closed subspaces $``$ of $`𝒞`$ which are $`𝒟`$-bimodules under multiplication in $`𝒞`$.
###### Definition 3.1.
An element $`v`$ is a normalizer of $`𝒟`$ if $`v𝒟v^{}v^{}𝒟v𝒟`$. The set of all such elements is denoted $`𝒩_𝒟()`$ or, if $`𝒟`$ is clear, $`𝒩()`$. Recall that $``$ is said to be regular if the closed span of $`𝒩_𝒟()`$ equals $``$.
Typically, normalizers play a major role in constructing coordinates for operator algebras. Since normalizers depend on the involution, it can be difficult to determine if isomorphisms that are not $``$-extendible or isometric preserve normalizers. Intertwiners (Definition 2.1) are not defined in terms of the involution and so it can be easier to decide if they are preserved by such isomorphisms. We begin with a comparison of normalizers and intertwiners.
It is easy to find examples of intertwiners of abelian $`C^{}`$-algebras which are not normalizers: for a simple example, observe that every element of $`M_2()`$ is an intertwiner for $`I_2`$. However, the next theorem shows that when the abelian algebra is a MASA, intertwiners are normalizers.
###### Proposition 3.2.
If $`v𝒞`$ is an intertwiner for $`𝒟`$, then $`v^{}v,vv^{}𝒟^{}𝒞`$. If $`𝒟`$ is maximal abelian in $`𝒞`$, then $`v`$ is a normalizer of $`𝒟`$.
###### Proof.
Let $`v`$ be an intertwiner. Let $`J_s(v):=\{d𝒟:vd=0\}`$ and $`J_r(v):=\{d𝒟:dv=0\}`$. Then $`J_s`$ and $`J_r`$ are norm-closed ideals in $`𝒟`$. Define a mapping $`\alpha _v`$ between $`𝒟/J_s`$ and $`𝒟/J_r`$ by $`\alpha _v(d+J_s)=d^{}+J_r`$, where $`d^{}𝒟`$ is chosen so that $`vd=d^{}v`$. It is easy to check that $`\alpha _v`$ is a well-defined $``$-isomorphism of $`𝒟/J_s`$ onto $`𝒟/J_r`$.
Let $`d=d^{}𝒟`$. Then $`\alpha _v(d+J_s)=d^{}+J_r`$ where $`d^{}`$ is chosen so that $`d^{}=d^{}𝒟`$. Thus, we have the equality of sets,
$$\begin{array}{cc}\hfill \{v^{}vd\}=v^{}v(d+J_s)& =v^{}(d^{}+J_r)v=[(d^{}+J_r)v]^{}v\hfill \\ & =[v(d+J_s)]^{}v=(d+J_s)v^{}v=\{dv^{}v\}.\hfill \end{array}$$
Hence $`v^{}v`$ commutes with the selfadjoint elements of $`𝒟`$ and so commutes with $`𝒟`$. Since $`v^{}`$ is also an intertwiner, we conclude similarly that $`vv^{}𝒟^{}`$.
If $`𝒟`$ is maximal abelian, then $`v𝒟v^{}=𝒟vv^{}𝒟`$ and $`v^{}𝒟v=v^{}v𝒟𝒟`$, and $`v`$ is a normalizer as desired. ∎
For $`v𝒩(𝒞)`$, let $`S(v):=\{\varphi \widehat{𝒟}:\varphi (v^{}v)>0\}`$; note this is an open set in $`\widehat{𝒟}`$. As observed by Kumjian (see \[21, Proposition 6\]), there is a homeomorphism $`\beta _v:S(v)S(v^{})`$ given by
$$\beta _v(\varphi )(d)=\frac{\varphi (v^{}dv)}{\varphi (v^{}v)}.$$
It is easy to show that $`\beta _v^1=\beta _v^{}`$.
###### Proposition 3.3.
For $`v𝒩(𝒞)`$, if $`\beta _v^{}`$ extends to a homeomorphism of $`\overline{S(v^{})}`$ onto $`\overline{S(v)}`$, then $`v`$ is an intertwiner. Moreover, if $`:=\{w𝒞:w𝒟=𝒟w\}`$ is the set of intertwiners, then $`𝒩(𝒞)`$ is contained in the norm-closure of $``$, and when $`𝒟`$ is a MASA in $`𝒞`$, $`𝒩(𝒞)=\overline{}`$.
###### Proof.
It is clear that the set of normalizers is norm-closed.
Regard $`𝒞`$ as sitting inside its double dual $`𝒞^{\mathrm{\#}\mathrm{\#}}`$. Let $`v=u|v|=|v^{}|u`$ be the polar decomposition for $`v`$. Since $`u`$ is the strong-$``$ limit of $`u_n:=v(1/n+|v|)^1`$, we find that $`u`$ also normalizes $`𝒟^{\mathrm{\#}\mathrm{\#}}`$.
Therefore, given any $`d𝒟`$, $`vdv^{}=u|v|du^{}|v^{}|=udu^{}vv^{}`$. Hence for $`\varphi S(v^{})`$, we have
(2)
$$\beta _v^{}(\varphi )(d)=\varphi (udu^{}).$$
Suppose now that $`\beta _v^{}`$ extends to a homeomorphism of $`\overline{S(v^{})}`$ onto $`\overline{S(v)}`$. By Tietze’s Extension Theorem, we may then find an element $`d_1𝒟`$ such that for every $`\varphi S(v^{})`$, $`\varphi (d_1)=\beta _v^{}(\varphi )(d).`$ Thus, for every $`\varphi \widehat{𝒟}`$, $`\varphi (vdv^{})=\varphi (udu^{})\varphi (vv^{})=\varphi (d_1vv^{}),`$ so that
$$(udu^{}d_1)vv^{}(ud^{}u^{}d_1^{})=0.$$
This shows that $`udu^{}v=d_1v`$ and so
$$vd=u|v|d=uu^{}ud|v|=udu^{}u|v|=udu^{}v=d_1v.$$
Hence $`v𝒟𝒟v`$. Since the adjoint of a normalizer is again a normalizer and $`\beta _v=(\beta _v^{})^1`$, we may repeat this argument to obtain $`v^{}𝒟𝒟v^{}`$. Taking adjoints yields $`𝒟vv𝒟`$. Hence $`v`$ is an intertwiner.
Given a general normalizer $`v`$, let $`\epsilon >0`$ and let $`K=\{\varphi \widehat{𝒟}:\varphi (vv^{})\epsilon ^2\}`$. Then $`K`$ is a compact subset of $`S(v^{})`$. Choose $`d_0𝒟`$ so that $`0d_0I`$ and $`\widehat{d_0}`$ is compactly supported in $`S(v^{})`$ and $`\widehat{d_0}|_K=1`$. Since $`\beta _{v^{}d_0}=\beta _v^{}|_{S(v^{}d_0)}`$ and $`\beta _v^{}`$ is a homeomorphism, it extends to a homeomorphism of $`\overline{S(v^{}d_0)}`$ onto $`\overline{S(d_0v)}`$. Thus $`d_0v`$ is an intertwiner. Further, $`d_0vv=(d_01)vv^{}(d_01)^{1/2}<\epsilon `$ so $`d_0v`$ approximates $`v`$ to within $`\epsilon `$. Thus, $`𝒩(𝒞)\overline{}`$. When $`𝒟`$ is a MASA in $`𝒞`$, Proposition 3.2 shows every intertwiner of $`𝒟`$ is a normalizer, so $`𝒩(𝒞)=\overline{}`$. ∎
Remarks 3.4. (1) Taken together, these two propositions show that for a MASA $`𝒟`$ in a $`C^{}`$-algebra $`𝒞`$, a partial isometry $`v`$ is a normalizer if and only if it is an intertwiner. Related results for partial isometries are known \[30, Lemma 3.2\].
(2) The $``$-isomorphism $`\alpha _v:𝒟/J_s𝒟/J_r`$ appearing in the proof of Proposition 3.2 induces a homeomorphism from the zero set $`Z_r:=\{\rho \widehat{𝒟}:\rho |_{J_r}=0\}`$ onto the zero set $`Z_s:=\{\rho \widehat{𝒟}:\rho |_{J_s}=0\}`$. It is not hard to show that when $`v`$ is an intertwiner and $`𝒟`$ is a MASA, $`Z_s=\overline{S(v)}`$, $`Z_r=\overline{S(v^{})}`$ and that $`\alpha _v`$ is the extension of $`\beta _v^{}=\beta _v^1`$ to $`Z_r`$. Thus, it is possible to describe $`\beta _v`$ without explicit reference to the $``$-structure.
We turn to constructing intertwiners in a module using the following property.
###### Definition 3.5.
Let $`𝒞`$ be a unital $`C^{}`$-algebra. A $`C^{}`$-subalgebra $`𝒟𝒞`$ is said to have the extension property if every pure state of $`𝒟`$ has a unique extension to a state on $`𝒞`$ and no pure state of $`𝒞`$ annihilates $`𝒟`$.
If $`𝒟𝒞`$ is abelian and $`𝒟`$ has the extension property relative to $`𝒞`$, then the Stone-Weierstrass Theorem implies that $`𝒟`$ is a MASA \[20, p. 385\]. We shall make essential use of the following result characterizing the extension property for abelian algebras.
###### Theorem 3.6 (\[1, Corollary 2.7\]).
Let $`𝒞`$ be a unital $`C^{}`$-algebra and let $`𝒟`$ be an abelian $`C^{}`$-subalgebra of $`𝒞`$ which contains the unit of $`𝒞`$. Then $`𝒟`$ has the extension property if and only if
$$\overline{\mathrm{co}}\{gxg^1:g𝒟\text{ and }g\text{ is unitary}\}𝒟\mathrm{}.$$
Furthermore, when this occurs, $`𝒟`$ is a MASA and there exists a conditional expectation $`E:𝒞𝒟`$ such that
$$\overline{\mathrm{co}}\{gxg^1:g𝒟\text{ and }g\text{ is unitary}\}𝒟=\{E(x)\};$$
Remark. Archibold, Bunce and Gregson in also show that the condition
$$𝒞=𝒟+\overline{\mathrm{span}}\{cddc:c𝒞,d𝒟\}$$
also characterizes the extension property for an abelian subalgebra $`𝒟𝒞`$. This characterization was important in Kumjian’s work on $`C^{}`$-diagonals.
###### Definition 3.7.
For $`v𝒩(𝒞)`$, define $`E_v:𝒞𝒩(𝒞)`$ by
$$E_v(x)=vE(v^{}x).$$
Our first application of Theorem 3.6 is essentially contained in the proof of \[24, Proposition 4.4\], but Theorem 3.6 provides a different (and simpler) proof. The crucial implication of Proposition 3.8 is that bimodules contain many normalizers.
###### Proposition 3.8.
Suppose $`𝒟`$ is a abelian $`C^{}`$-subalgebra of the unital $`C^{}`$-algebra $`𝒞`$ which has the extension property, and $`𝒞`$ is a $`𝒟`$-bimodule. If $`v𝒞`$ is a normalizer (resp. intertwiner) and $`x`$, then $`E_v(x)`$ and is a normalizer (resp. intertwiner).
###### Proof.
Fix $`x`$ and let $`G`$ be the unitary group of $`𝒟`$. Since $`vN_𝒟(𝒞)`$, for every $`gG`$, $`vgv^{}𝒟`$, so that $`(vgv^{})xg^1`$. Thus the norm-closed convex hull,
$$H:=\overline{\text{co}}\{(vgv^{})xg^1:gG\}.$$
By Theorem 3.6, $`E(v^{}x)`$ belongs to $`K:=\overline{\text{co}}\{gv^{}xg^1:gG\}`$. Since $`vKH`$, we conclude that $`vE(v^{}x)=E_v(x)`$. ∎
Although we have no application for it, the following result, which is a corollary of Proposition 3.8, provides another means of constructing normalizers in a bimodule.
###### Proposition 3.9.
If $`v𝒩(𝒞)`$, $`m`$ and $`|\varphi (v)||\varphi (m)|`$ for all $`\varphi ^1(𝒞)`$, then $`v`$.
###### Proof.
For $`\rho \widehat{𝒟}`$ with $`\rho (v^{}v)0`$, $`\rho (v^{}v)`$ equals
$$\rho (v^{}v)^{1/2}[v,\rho ](v)|\rho (v^{}v)^{1/2}[v,\rho ](m)|=|\rho (v^{}m)|=|\rho (E(v^{}m))|,$$
and so, for $`\rho \widehat{𝒟}`$, $`\rho (v^{}v)|\rho (E(v^{}m))|`$. Thus, there exists $`d𝒟`$ with $`v^{}v=E(v^{}m)d`$. Hence for $`n`$, $`(v^{}v)^{1/n}`$ belongs to the closed ideal of $`𝒟`$ generated by $`E(v^{}m)`$. By Proposition 3.8, $`vE(v^{}m)`$ and we conclude $`v(v^{}v)^{1/n}`$. But $`v=\underset{n\mathrm{}}{lim}v(v^{}v)^{1/n}`$, whence $`v`$. ∎
Our second application of Theorem 3.6 is to provide an alternate proof of, and slightly strengthen, a result of Kumjian.
###### Proposition 3.10 (\[21, Lemma 9, p. 972\]).
Let $`𝒟`$ be an abelian $`C^{}`$-subalgebra of the unital $`C^{}`$-algebra $`𝒞`$ with the extension property. For $`v𝒩(𝒞)`$, $`v^{}E(v)`$ and $`vE(v^{})`$ both belong to $`𝒟`$. If $`\rho S(v)`$, then the following are equivalent:
1. $`\rho (v)0`$,
2. $`\beta _v(\rho )=\rho `$,
3. $`\rho (E(v^{}))0`$.
###### Proof.
Taking $`x=I`$ and $`=𝒟`$ in Proposition 3.8 we see that $`v^{}E(v)`$ and $`vE(v^{})`$ both belong to $`𝒟`$.
Suppose $`\rho (v)0`$. An easy calculation shows that when $`d𝒟`$ and $`\rho (d)0`$, then $`\beta _{vd}=\beta _v`$. Let $`d=v^{}E(v)`$. By hypothesis, we find $`\rho (d)=|\rho (v)|^20`$, and another calculation shows that $`\beta _{vd}=\beta _v`$, so (1) implies (2).
Assume (2) holds. Letting $`G`$ again be the unitary group of $`𝒟`$, we have, for all $`gG`$,
$$\rho (v^{}gvg^1)=\rho (v^{}gv)\rho (g^1)=\beta _v(\rho )(g)\rho (g^1)\rho (v^{}v)=\rho (v^{}v).$$
Thus, $`\rho (v^{}\overline{\text{co}}\{gvg^1gG\})\{\rho (v^{}v)\}`$, and therefore by Theorem 3.6 we obtain $`\rho (v^{})\rho (v)=\rho (v^{}E(v))0`$. Thus $`\rho (E(v^{}))0`$ and by our earlier remarks, $`vE(v^{})𝒟`$.
Finally, it is evident that (3) implies (1). ∎
## 4. Regular $`C^{}`$-inclusions and $`C^{}`$-diagonals
We now turn to the key context for our subsequent work, that of $`C^{}`$-diagonals and regular $`C^{}`$-inclusions. After recalling Kumjian’s twist, we show that the elements of the twist are eigenfunctionals on the $`C^{}`$-algebra and conversely (Proposition 4.8 and Theorem 4.10). This allows us to show that every regular $`C^{}`$-inclusion has a quotient which is a $`C^{}`$-diagonal with the same coordinate system, Theorem 4.9, as well as strengthening various results from Section 2 in this context. The crucial result for subsequent sections is Theorem 4.14, which shows that for a bimodule $``$, the intrinsically defined eigenfunctionals on $``$ and the restriction of the twist are the same.
###### Definition 4.1.
The pair $`(𝒞,𝒟)`$ will be called a regular $`C^{}`$-inclusion if $`𝒟`$ is a maximal abelian $`C^{}`$-subalgebra of the unital $`C^{}`$-algebra $`𝒞`$ such that
1. $`𝒟`$ has the extension property in $`𝒞`$;
2. $`𝒞`$ is regular (as a $`𝒟`$-bimodule).
Always, $`E`$ denotes the (unique) conditional expectation of $`𝒞`$ onto $`𝒟`$. We call $`(𝒞,𝒟)`$ a $`C^{}`$-diagonal if, in addition,
1. $`E`$ is faithful.
Remarks 4.2.
1. Renault gives, without proof, an example of an algebra satisfying only the first two conditions but not the third, i.e., a regular $`C^{}`$-inclusion that is not a $`C^{}`$-diagonal.
2. By Propositions 3.2 and 3.3, regularity of a bimodule $``$ is equivalent to norm-density of the $`𝒟`$-intertwiners.
3. As observed by Renault , this definition of $`C^{}`$-diagonal is equivalent to Kumjian’s original definition, namely,
1. there is a faithful conditional expectation $`E:𝒞𝒟`$
2. the closed span of the free normalizers in $`𝒞`$ is $`\mathrm{ker}E`$.
A normalizer $`v`$ of $`𝒟`$ is free if $`v^2=0`$. Kumjian also required that $`𝒞`$ is separable and $`\widehat{𝒟}`$ is second countable, but shows this is not necessary. It is often easier to verify Kumjian’s axioms when working with particular examples.
Unless explicitly stated otherwise, for the remainder of this section, we work in the following setting.
Context 4.3. Let $`(𝒞,𝒟)`$ be a regular $`C^{}`$-inclusion and $`𝒞`$ be a norm-closed $`𝒟`$-bimodule, where the module action is multiplication. That is, for $`d,e𝒟`$ and $`m`$, $`dme:=dme`$.
For $`\rho \widehat{𝒟}`$, the unique extension on $`𝒞`$ is $`\rho E`$, which we will again denote by $`\rho `$. Thus, we regard $`\rho `$ as either a multiplicative linear functional on $`𝒟`$ or a pure state on $`𝒞`$ which satisfies $`\rho E=\rho `$. In particular, for $`d𝒟`$ and $`x𝒞`$, $`\rho (dx)=\rho (xd)=\rho (d)\rho (x)`$.
We now summarize some results and definitions from and then relate eigenfunctionals to the elements of Kumjian’s twist. Note that these results from hold for regular $`C^{}`$-inclusions.
###### Definition 4.4 (Kumjian).
For $`v𝒩(𝒞)`$ and $`\rho \widehat{𝒟}`$ with $`\rho (v^{}v)>0`$, define a linear functional on $`𝒞`$, $`[v,\rho ]`$, by
$$[v,\rho ](x)=\frac{\rho (v^{}x)}{\rho (v^{}v)^{1/2}}.$$
Kumjian denotes by $`\mathrm{\Gamma }`$ the collection of all such linear functionals. We shall see in Proposition 4.8 below shows that $`\mathrm{\Gamma }=^1(𝒞)`$.
We follow Kumjian \[21, p. 982\] in pointing out that Proposition 3.10 implies
###### Corollary 4.5 (Kumjian).
The following are equivalent:
1. $`[v,\rho ]=[w,\rho ]`$,
2. $`\rho (v^{}w)>0`$,
3. there are $`d,e𝒟`$ with $`\rho (d),\rho (e)>0`$ so that $`vd=we`$.
Remark 4.6. Using Theorem 3.6 and the techniques of the proof of Proposition 3.8, one can show that when $`\rho (v^{}w)>0`$, we may take $`d=w^{}wE(v^{}w)`$ and $`e=w^{}vE(v^{}w)`$ in the third part of Corollary 4.5.
Kumjian shows that $`\mathrm{\Gamma }`$, with a suitable operation and the relative weak-topology, is a groupoid and admits a natural $`𝕋`$-action, given by $`\lambda [v,\rho ]=[\overline{\lambda }v,\rho ]`$. The range map for $`\mathrm{\Gamma }`$ is $`[v,\rho ]\beta _v(\rho )`$ and the source map for $`\mathrm{\Gamma }`$ is $`[v,\rho ]\rho .`$ The map $`[v,\rho ]\mathrm{\Gamma }(\beta _v(\rho ),\rho )`$ sends $`\mathrm{\Gamma }`$ to a Hausdorff equivalence relation (principal groupoid) on $`\widehat{𝒟}`$, denoted $`\mathrm{\Gamma }\backslash 𝕋`$, and $`\mathrm{\Gamma }`$ is a locally trivial principal $`𝕋`$-bundle over $`\mathrm{\Gamma }\backslash 𝕋`$.
While the range and source maps of this groupoid have, a priori, no connection to the range and source maps for eigenfunctionals, they turn out to be the same.
We recall the multiplication on $`\mathrm{\Gamma }`$ and use it to define multiplication of eigenfunctionals for regular $`C^{}`$-inclusions. The partially defined multiplication on $`\mathrm{\Gamma }`$ has $`[v_1,\rho _1],[v_2,\rho _2]`$ composable if $`\rho _1=\beta _{v_2}(\rho _2)`$, in which case,
$$[v_1,\varphi _1][v_2,\rho _2]=[v_1v_2,\rho _2].$$
###### Definition 4.7 (Kumjian).
A twist is a proper $`𝕋`$-groupoid $`\mathrm{\Gamma }`$ so that $`\mathrm{\Gamma }/𝕋`$ is an r-discrete principal groupoid.
Kumjian constructs, for each twist $`\mathrm{\Gamma }`$, a $`C^{}`$-diagonal, $`(A(\mathrm{\Gamma }),B(\mathrm{\Gamma }))`$. The main result of is that for $`(𝒞,𝒟)`$ a $`C^{}`$-diagonal, there is a unique (up to isomorphism) twist $`\mathrm{\Gamma }`$ and an isomorphism $`\mathrm{\Phi }:A(\mathrm{\Gamma })𝒞`$ such that $`\mathrm{\Phi }(B(\mathrm{\Gamma }))=𝒟`$. Thus, every twist arises as a $`\mathrm{\Gamma }`$ as above, for some $`C^{}`$-diagonal $`(𝒞,𝒟)`$.
From our point of view, justified by the following proposition, the twist $`\mathrm{\Gamma }`$ associated to $`(𝒞,𝒟)`$ is $`^1(𝒞)`$ equipped with this groupoid operation, $`𝕋`$-action, and topology.
###### Proposition 4.8.
For all $`\rho \widehat{𝒟}`$ and $`v𝒩(𝒞)`$ with $`\rho (v^{}v)>0`$, $`[v,\rho ]^1(𝒞)`$. Moreover, the range and source maps agree, that is, viewed as an element of $`^1(𝒞)`$, we have $`\rho =s([v,\rho ])`$ and $`\beta _v(\rho )=r([v,\rho ])`$.
Conversely, if $`\varphi (𝒞)`$ there exists a normalizer $`v𝒞`$ such that $`\varphi (v)0`$, and $`v`$ may be taken to be an intertwiner if desired. For any normalizer (or intertwiner) $`v`$ with $`\varphi (v)0`$, we have $`s(\varphi )(v^{}v)0`$ and
$$\varphi (v)[v,s(\varphi )]=[v,s(\varphi )](v)\varphi .$$
In particular, if $`\varphi ^1(𝒞)`$ then $`\lambda :=\frac{\varphi (v)}{[v,s(\varphi )](v)}𝕋`$ and $`\varphi =\lambda [v,s(\varphi )].`$
###### Proof.
Easily, we have $`[v,\rho ](xd)=[v,\rho ](x)\rho (d)`$ and, as $`v^{}v(v^{}dx)=(v^{}dv)v^{}x`$, we can apply $`\rho `$ to this equation and divide by $`\rho (v^{}v)^{3/2}`$ to obtain
$$[v,\rho ](dx)=\frac{\rho (v^{}dv)}{\rho (v^{}v)}[v,\rho ](x)=\beta _v(d)[v,\rho ](x),$$
so $`[v,\rho ]`$ is a $`𝒟`$-eigenfunctional with range $`\beta _v(\rho )`$ and source $`\rho `$. Letting $`w=\rho (v^{}v)^{1/2}v`$, $`[v,\rho ](w)=1`$, so $`[v,\rho ]^1(𝒞)`$.
For the converse, let $`\varphi ^1(𝒞)`$ and set $`\rho =s(\varphi )`$. Since the span of $`N_𝒟(𝒞)`$ is norm dense in $`𝒞`$, there exists $`vN_𝒟(𝒞)`$ with $`\varphi (v)0`$, which by Proposition 3.3 may be taken to be an intertwiner if desired. Fix such a $`v`$.
For any $`n`$, $`\varphi (v(v^{}v)^{1/n})=\varphi (v)\rho (v^{}v)^{1/n}`$. Since $`v(v^{}v)^{1/n}`$ converges to $`v`$, $`\rho (v^{}v)>0`$. Also, for any $`d𝒟`$ we have,
$`r(\varphi )(d)\varphi (v)=\varphi (dv)`$ $`={\displaystyle \frac{\varphi (dv(v^{}v))}{\rho (v^{}v)}}={\displaystyle \frac{\varphi (v(v^{}dv))}{\rho (v^{}v)}}`$
$`={\displaystyle \frac{\varphi (v)\rho (v^{}dv)}{\rho (v^{}v)}}=\varphi (v)\beta _v(\rho )(d),`$
so $`r(\varphi )=\beta _v(\rho )`$.
For any unitary element $`g𝒟`$ we have,
$`\varphi (vgv^{}xg^1)`$ $`=r(\varphi )(vgv^{})\varphi (x)\rho (g^1)`$
$`={\displaystyle \frac{\rho (v^{}vgv^{}v)}{\rho (v^{}v)}}\varphi (x)\rho (g)^1=\rho (v^{}v)\varphi (x).`$
Theorem 3.6 implies that $`vE(v^{}x)`$ belongs to the closed convex hull of $`\{vgv^{}xg^1:g𝒰(𝒟)\}`$, and thus $`\varphi (vE(v^{}x))=\rho (v^{}v)\varphi (x)`$. Hence for any $`x𝒞`$,
$`\varphi (v)[v,\rho ](x)`$ $`={\displaystyle \frac{\varphi (v)\rho (v^{}x)}{\rho (v^{}v)^{1/2}}}={\displaystyle \frac{\varphi (vE(v^{}x))}{\rho (v^{}v)^{1/2}}}`$
$`=\rho (v^{}v)^{1/2}\varphi (x)=[v,\rho ](v)\varphi (x).`$
This equality, together with the fact that $`[v,\rho ]=1,`$ shows that $`\lambda 𝕋`$ when $`\varphi `$ has unit norm. ∎
Proposition 4.8 leads to a description of $`()`$ in terms of $`[v,\rho ]`$’s, for a norm-closed $`𝒟`$-bimodule $``$ (Theorem 4.14). Before giving this description, we use Proposition 4.8 as a tool in the following result, which gives the precise relationship between the concepts of regular $`C^{}`$-inclusion and $`C^{}`$-diagonal.
###### Theorem 4.9.
Let $`(𝒞,𝒟)`$ be a regular $`C^{}`$-inclusion, and let $`𝔑:=\{x𝒞:E(x^{}x)=0\}`$ be the left kernel of $`E`$. Then $`𝔑`$ is a closed (two-sided) ideal of $`𝒞`$ and
$$𝔑=\{x𝒞:\varphi (x)=0\text{ for all }\varphi (𝒞)\}.$$
Let $`\pi :𝒞𝒞/𝔑`$ be the quotient map. Then $`(\pi (𝒞),\pi (𝒟))`$ is a $`C^{}`$-diagonal, $`\pi |_𝒟`$ is a $``$-isomorphism of $`𝒟`$ onto $`\pi (𝒟)`$, and the restriction of the adjoint map $`\pi ^\mathrm{\#}|_{(\pi (C))}`$ is an isometric isomorphism of $`(\pi (𝒞))`$ onto $`(𝒞)`$.
In particular, the first part of the theorem shows that, in a $`C^{}`$-diagonal, $`^1(𝒞)`$ separates points.
###### Proof.
Since $`𝔑=\{x𝒞:E(x^{}x)=0\}`$ is clearly a closed left ideal, to show that it is also a right ideal, it suffices to show that, for $`x𝔑`$ and a normalizer $`v𝒞`$, $`xv𝔑`$. To do this, we shall prove that for every $`\rho \widehat{𝒟}`$, $`\rho (v^{}x^{}xv)=0`$. When $`\rho (v^{}v)=0`$, this holds since $`\rho (v^{}x^{}xv)x^2\rho (v^{}v)`$. When $`\rho (v^{}v)0`$, let $`\psi \widehat{𝒟}`$ be given by
$$\psi (z)=\frac{\rho (v^{}zv)}{\rho (v^{}v)}$$
and observe that, as $`x𝔑`$, $`\psi (x^{}x)=0`$ and hence $`\rho (v^{}x^{}xv)=0`$ in this case as well. Thus $`𝔑`$ is a closed two-sided ideal.
We now show that $`𝔑=\{x𝒞:\varphi (x)=0\text{ for all }\varphi (𝒞)\}.`$ Suppose that $`x𝔑`$. Given $`\varphi (𝒞)`$, by Proposition 4.8 we may assume $`\varphi =[v,\rho ]`$ where $`\rho \widehat{𝒟}`$ and $`v`$ is a normalizer with $`\rho (v^{}v)>0`$. By the Cauchy-Schwarz inequality we have $`|\rho (v^{}x)|^2\rho (v^{}v)\rho (x^{}x)`$. But $`\rho (x^{}x)=\rho (E(x^{}x))=0`$, so $`[v,\rho ](x)=0`$.
Conversely, suppose $`x0`$ and $`\varphi (x)=0`$ for every $`\varphi ^1(𝒞)`$. We shall show that $`\rho (x^{}x)=0`$ for every $`\rho \widehat{𝒟}`$. So fix $`\rho \widehat{𝒟}`$. Notice that for every normalizer $`v𝒞`$, we have $`\rho (v^{}x)=0`$: this follows from the Cauchy-Schwartz inequality when $`\rho (v^{}v)=0`$ and from the hypothesis and Proposition 4.8 when $`\rho (v^{}v)0`$. Let $`\epsilon >0`$. Since the span of the normalizers is dense in $`𝒞`$, we may find normalizers $`v_1,\mathrm{},v_n`$ so that $`x_{i=1}^nv_i<\epsilon /x`$. Thus,
$$|\rho (x^{}x)|=\left|\rho \left(x^{}x\underset{j=1}{\overset{n}{}}v_j^{}x\right)\right|x^{}\underset{j=1}{\overset{n}{}}v_j^{}x<\epsilon ,$$
and we conclude that $`\rho (x^{}x)=0`$. Since this holds for every $`\rho \widehat{D}`$, we have $`E(x^{}x)=0`$.
Clearly, $`𝔑𝒟=0`$, so $`\pi |_𝒟`$ is an isomorphism of $`𝒟`$ onto $`\pi (𝒟)`$. To see that $`(\pi (𝒞),\pi (𝒟))`$ is a $`C^{}`$-diagonal, observe first that $`E`$ gives rise to a faithful conditional expectation on $`𝒞/𝔑`$, by the definition of $`𝔑`$. Given a pure state $`\rho `$ of $`\pi (𝒟)`$, let $`\tau _1`$ and $`\tau _2`$ be pure states of $`\pi (𝒞)`$ which extend $`\rho `$. Now $`\tau _i\pi `$ are extensions of the pure state $`\rho \pi `$ of $`𝒟`$ and hence coincide because $`𝒟`$ has the extension property in $`𝒞`$. Therefore, as $`\pi `$ is onto, $`\tau _1=\tau _2`$, and $`\pi (𝒟)`$ has the extension property in $`\pi (𝒞)`$.
To show that $`\pi (𝒞)`$ is regular (relative to $`\pi (𝒟)`$), let $`x+𝔑\pi (𝒞)`$, and $`\epsilon >0`$. Since $`𝒞`$ is regular, we may find normalizers $`v_i𝒞`$ such that $`y:=_{i=1}^nv_i`$ satisfies $`xy<\epsilon `$. Then $`\pi (v_i)`$ is a normalizer, and since $`\pi `$ is contractive, $`(x+𝔑)(y+𝔑)<\epsilon `$. Hence $`\pi (𝒞)`$ is regular. It is also clear that $`\pi ^\mathrm{\#}|_{(\pi (𝒞))}`$ is a homeomorphism of $`(\pi (𝒞))`$ onto $`(𝒞)`$ which preserves the partially defined product structure. Finally, since the adjoint of a quotient map is always isometric, the proof is complete. ∎
We turn now to additional consequences of Proposition 4.8. Recall that for any $`f𝒞^\mathrm{\#}`$, $`f^{}`$ is the bounded linear functional given by $`f^{}(x)=\overline{f(x^{})}`$. It is easy to see that if $`\varphi `$ is an eigenfunctional on $`𝒞`$, then so is $`\varphi ^{}`$ and also that $`s(\varphi ^{})=r(\varphi )`$ and $`r(\varphi ^{})=s(\varphi )`$. Thus, $`ff^{}`$ provides an involution on $`^1(𝒞)`$. The inverse of $`[v,\rho ]`$ is $`[v^{},\beta _v(\rho )]`$. Thus, we can summarize Proposition 4.8 and the discussion preceding it.
###### Theorem 4.10.
For a regular $`C^{}`$-inclusion $`(𝒞,𝒟)`$, $`^1(𝒞)=\mathrm{\Gamma }`$ and the range and source maps, the involution, and topology are all the same.
###### Corollary 4.11.
If $`\varphi ,\psi ()`$ satisfy $`r(\varphi )=r(\psi )`$ and $`s(\varphi )=s(\psi )`$, then there exists $`\lambda `$ such that $`\lambda 0`$ and $`\varphi =\lambda \psi `$.
###### Proof.
Without loss of generality, assume $`\varphi =\psi =1`$. Let $`\rho =s(\varphi )=s(\psi )`$. Theorem 2.6 shows $`\varphi `$ and $`\psi `$ extend to norm-one eigenfunctionals on $`𝒞`$, which we denote by the same symbols. By Proposition 4.8, we have $`\varphi =[v,\rho ]`$ and $`\psi =[w,\rho ]`$ for some $`v,w𝒩(𝒞)`$. It follows from the hypothesis that $`\beta _v(\rho )=\beta _w(\rho )`$. Thus $`\beta _{v^{}w}(\rho )=\beta _v(\beta _w(\rho ))=\rho `$, so we can find $`\lambda 𝕋`$ with $`\rho (\lambda v^{}w)>0`$. Hence by Corollary 4.5, $`[\overline{\lambda }v,\rho ]=[w,\rho ]`$ and therefore $`\lambda \varphi =\psi `$. ∎
Remark 4.12. Using the expression for $`\varphi =[v,s(\varphi )]`$ from Proposition 4.8 and a short calculation with $`r(\varphi )(x)=\rho (v^{}xv)/\rho (v^{}v)`$, where $`\rho =s(\varphi )`$, we obtain for all $`x𝒞`$,
$$\varphi (x)=\frac{s(\varphi )(v^{}x)}{[s(\varphi )(v^{}v)]^{1/2}}=\frac{r(\varphi )(xv^{})}{[r(\varphi )(vv^{})]^{1/2}}.$$
Also, we can sharpen the inequality of Proposition 2.5 (2) to an equality.
###### Corollary 4.13.
If $`\varphi ()`$, then $`|\varphi (m)|=\varphi B_{r(\varphi ),s(\varphi )}(m)`$ for all $`m`$.
###### Proof.
By Proposition 2.5, we have $`|\varphi (m)|\varphi B_{r(\varphi ),s(\varphi )}(m)`$ for every $`m`$.
To obtain the reverse inequality, fix $`m`$ such that $`B_{r(\varphi ),s(\varphi )}(m)0`$. Given $`\epsilon (0,1)`$, we may find elements $`a,b𝒟`$ such that $`\sigma (a)=\rho (b)=1`$ and $`B_{r(\varphi ),s(\varphi )}(m)>(1\epsilon )amb`$. Let $`m_0:=amb`$ and define a linear functional $`f`$ on $`m_0`$ by
$$f(tm_0)=t\varphi B_{r(\varphi ),s(\varphi )}(m_0)=t\varphi B_{r(\varphi ),s(\varphi )}(m).$$
By the Hahn-Banach Theorem and Proposition 2.5, $`f`$ extends to an eigenfunctional $`F`$ on $``$ such that for every $`x`$,
$$|F(x)|\varphi B_{r(\varphi ),s(\varphi )}(x).$$
Thus $`F\varphi `$, $`s(F)=s(\varphi )`$ and $`r(F)=r(\varphi )`$. By Corollary 4.11, there exists a nonzero scalar $`\lambda `$ with $`|\lambda |1`$ and $`F=\lambda \varphi `$. We obtain,
$`|\varphi (m)|`$ $`|F(m)|=|F(m_0)|=|f(m_0)|=\varphi {\displaystyle \frac{B_{r(\varphi ),s(\varphi )}(m)}{amb}}amb`$
$`>(1\epsilon )\varphi amb`$
$`(1\epsilon )\varphi B_{r(\varphi ),s(\varphi )}(m).`$
Letting $`\epsilon 0`$, we obtain the result. ∎
We now extend Proposition 4.8 to bimodules and show that the set of eigenfunctionals on a bimodule can be written in the form $`[v,\rho ]`$ with $`v`$.
###### Theorem 4.14.
If $`\varphi ^1()`$, then there is a intertwiner $`v`$ and $`\rho \widehat{𝒟}`$ so that $`\varphi =[v,\rho ]|_{}`$.
Conversely, if $`v𝒩()`$, and $`\rho \widehat{𝒟}`$ satisfies $`\rho (v^{}v)>0`$. then $`[v,\rho ]|_{}^1()`$.
###### Proof.
We prove the second assertion first. Given $`v𝒩()`$ and $`\rho \widehat{𝒟}`$ with $`\rho (v^{}v)0`$, Proposition 4.8 shows that $`\varphi :=[v,\rho ]|_{}`$ is an eigenfunctional on $``$. Clearly, $`\varphi [v,\rho ]=1`$. To show that $`\varphi =1`$, fix $`m`$ with $`\varphi (m)0`$ and set $`w=E_v(m)/\rho (v^{}v)`$. Proposition 3.8 shows that $`w𝒩()`$. Calculation then yields
(3)
$$\rho (w^{}w)=|[v,\rho ](m)|^20,[v,\rho ](w)=[v,\rho ](m).$$
Let $`\epsilon >0`$. Pick a norm-one positive element $`d𝒟`$ with $`\rho (d)=1`$ so that $`d^2w^{}w(1+\epsilon )\rho (w^{}w)`$. If we let $`s=wd/\rho (w^{}w)^{1/2}`$, then $`s(1+\epsilon )^{1/2}`$ and the equations (3) give
$$\left|\varphi \left(\frac{wd}{\rho (w^{}w)^{1/2}}\right)\right|=\left|[v,\rho ]\left(\frac{wd}{\rho (w^{}w)^{1/2}}\right)\right|=\left|\frac{[v,\rho ](w)}{\rho (w^{}w)^{1/2}}\right|=1.$$
To prove the first assertion, let $`\varphi ^1()`$ and fix $`m`$ with $`\varphi (m)0.`$ By Theorem 2.6, $`\varphi `$ extends to a norm-one eigenfunctional, also called $`\varphi `$, on $`𝒞`$. By Proposition 4.8, there is an intertwiner $`v`$ and $`\rho \widehat{𝒟}`$ with $`\rho (v^{}v)>0`$ so that $`\varphi =[v,\rho ]`$. By Proposition 3.8,
$$w:=\frac{vE(v^{}m)}{\varphi (m)\rho (v^{}v)^{1/2}}.$$
and $`w`$ is an intertwiner. Thus,
$$\varphi (w)=\varphi (v)\frac{\rho (v^{}m)}{\varphi (m)\rho (v^{}v)^{1/2}}=\varphi (v)\frac{[v,\rho ](m)}{\varphi (m)}=\varphi (v),$$
so that $`[w,\rho ]=[v,\rho ]=\varphi `$. ∎
Remark 4.15. Let $`𝒜𝒞`$ be a norm-closed algebra which is also a $`𝒟`$-bimodule. We can use the multiplication defined for elements of $`\mathrm{\Gamma }`$ to give a (also partially defined) multiplication on $`^1(𝒜)`$ and hence on $`(𝒜)`$. Indeed, call $`\varphi _1,\varphi _2^1(𝒜)`$ composable if $`s(\varphi _1)=r(\varphi _2)`$. By Theorem 4.14, we can write $`\varphi _i=[v_i,\rho _i]`$, where $`v_1`$ and $`v_2`$ are in $`𝒩(𝒜)`$. Define $`\varphi _1\varphi _2`$ to be the product of the $`[v_i,\rho _i]`$ restricted to $`𝒜`$, namely $`[v_1v_2,\rho _2]|_𝒜`$.
We can also improve several of the results of the previous section. First, we immediately have a unique extension in Theorem 2.6.
###### Corollary 4.16.
Suppose $`_1,_2`$ are norm-closed $`𝒟`$-bimodules with $`_1_2`$. There is a unique isometric map $`\iota :(_1)(_2)`$ so that, for every $`\varphi (_1)`$, $`\iota (\varphi )|__1=\varphi `$. The image $`\iota ((_1))`$ is an open subset of $`(_2)`$.
If in addition, $`_2`$ is regular, then $`\iota `$ is $`\sigma (_1^\mathrm{\#},_1)`$$`\sigma (_2^\mathrm{\#},_2)`$ continuous on bounded subsets of $`(_1)`$.
###### Proof.
Fix $`\varphi (_1)`$. Theorem 4.14 shows $`\varphi `$ extends uniquely to an eigenfunctional $`\iota (\varphi )`$ on $`_2`$ of the same norm. Writing $`\varphi =[v,\rho ]`$ for some normalizer $`v_1`$ and $`\rho \widehat{𝒟}`$, then $`\{\psi (_2):\psi (v)0\}`$ is a $`\sigma (_2^\mathrm{\#},_2)`$-open set containing $`\iota (\varphi )`$, so $`\iota ((_1))`$ is an open set in $`(_2)`$.
It remains to prove that $`\iota `$ is continuous on bounded subsets of $`(_1)`$ when $`_2`$ is regular. Suppose $`\varphi _\lambda `$ is a bounded net in $`(_1)`$ converging $`\sigma (_1^\mathrm{\#},_1)`$ to $`\varphi (_1)`$. Let $`\rho =s(\varphi )`$ and $`\rho _\lambda =s(\varphi _\lambda )`$. Then $`\rho _\lambda `$ converges in the $`\sigma (𝒞^\mathrm{\#},𝒞)`$-topology to $`\rho `$.
By Theorem 4.14, there exists a normalizer $`v_1`$ such that $`\varphi =\varphi [v,\rho ]`$. For large enough $`\lambda `$, $`\varphi _\lambda (v)0`$, so there exist scalars $`t_\lambda `$ with $`|t_\lambda |=\varphi _\lambda `$ and $`\varphi _\lambda =t_\lambda [v,\rho _\lambda ]`$. Since $`\varphi _\lambda (v)\varphi (v)`$ we have $`t_\lambda \varphi `$.
For any normalizer $`w_2`$, we have
$$\iota (\varphi _\lambda )(w)=t_\lambda [v,\rho _\lambda ](w)=t_\lambda \frac{\rho _\lambda (v^{}w)}{\rho _\lambda (v^{}v)^{1/2}}\varphi [v,\rho ](w)=\iota (\varphi )(w).$$
As $`_2`$ is the span of the normalizers it contains and $`\varphi _\lambda `$ is a bounded net, we conclude that for any $`x_2`$, $`\iota (\varphi _\lambda )(x)\iota (\varphi )(x)`$. ∎
Remark 4.17. Given a bimodule $`𝒞`$, the set $`\mathrm{\Gamma }_{}:=\{\varphi \mathrm{\Gamma }:\varphi |_{}0\}`$ is an open subset of $`\mathrm{\Gamma }`$ which plays a crucial role in the study of $``$ (see, for example, ). Theorem 4.14, together with Corollary 4.16, shows that the restriction map $`[v,\rho ]\mathrm{\Gamma }_{}[v,\rho ]|_{}`$ is a homeomorphism of $`\mathrm{\Gamma }_{}`$ onto $`^1()`$. Thus, $`\mathrm{\Gamma }_{}`$ can be defined directly in terms of the bimodule structure of $``$, without explicit reference to $`𝒞`$.
Since the norm is only lower semi-continuous for weak-$``$ convergence, it is not possible to show that $`^1()`$ is locally compact for general modules. However, for regular $`C^{}`$-inclusions, we can show this.
###### Proposition 4.18.
With the relative weak-topology, $`^1()\{0\}`$ is compact. Thus, $`^1()`$ is a locally compact Hausdorff space.
###### Proof.
Suppose that $`\varphi _\lambda `$ is a net in $`^1()\{0\}`$ which converges to $`\varphi (^\mathrm{\#})_1.`$ If $`\varphi =0`$, there is nothing to do, so we assume that $`\varphi 0`$ and show that $`\varphi =1`$.
Fix a normalizer $`v`$ with $`\varphi (v)>0`$. From Theorem 4.14, $`\varphi =\varphi [v,\rho ]`$. Choose a positive element $`d𝒟`$ so that $`0dv^{}vdI`$ and $`\widehat{dv^{}vd}=1`$ in a neighborhood of $`\rho `$. For large enough $`\lambda `$, $`\varphi _\lambda (vd)0`$, so there exists $`t_\lambda 𝕋`$ such that $`\varphi _\lambda =t_\lambda [vd,s(\varphi _\lambda )]`$. Thus, $`1=|\varphi _\lambda (vd)|`$. As $`\varphi _\lambda `$ converges to $`\varphi `$, we obtain $`|\varphi (vd)|=1=vd`$, so $`\varphi =1`$. ∎
As usual, we may regard an element $`m`$ as a function on $`^1()`$ via $`\widehat{m}(\varphi )=\varphi (m),`$ and $`^1()`$ can be regarded as a set of coordinates for $``$. Thus we make the following definition.
###### Definition 4.19.
For a norm-closed $`𝒟`$-bimodule $``$, we call the set $`_𝒟^1()`$, equipped with the relative weak-$``$ topology, the $`𝕋`$-action, and the range and source mappings, a coordinate system for $``$.
When $`𝒜`$ is both a norm-closed algebra and a $`𝒟`$-bimodule, the coordinate system $`^1(𝒜)`$ also has the additional structure of a continuous partially defined product as described in Remark 4.14. In this case we will sometimes refer to the coordinate system as a semitwist.
###### Definition 4.20.
If $`𝒞`$ is a $`𝒟`$-bimodule, let $`R():=\{|\varphi |:\varphi ^1()\}`$. Then $`R()`$ may be identified with the quotient $`^1()\backslash 𝕋`$ of $`^1()`$ by the natural action of $`𝕋`$. Obviously, $`\varphi |\varphi |`$ is the quotient map, and the topology on $`R()`$ is the quotient topology. Corollary 4.16 shows that we may regard $`R()`$ as a subset of $`R(𝒞)`$, and, as $`𝒞`$ is regular, Corollary 4.16 also implies that if $`v`$ is an intertwiner, then $`G_v:=\{\varphi R():|\varphi (v)|>0\}`$ is an open set and $`\{G_v:v\text{ is an intertwiner}\}`$ is a base for the topology of $`R()`$. Thus $`R()`$ is a locally compact Hausdorff space.
We shall sometimes find it useful to view $`R()`$ as a topological relation on $`\widehat{𝒟}`$. The map $`|\varphi |(r(\varphi ),s(\varphi ))`$ is a bijection between $`R()`$ and $`\{(r(\varphi ),s(\varphi ))\widehat{𝒟}\times \widehat{𝒟}:\varphi ^1()\}`$, and we will sometimes identify these two sets under this bijection. With this identification, $`(\sigma ,\rho )R()`$ if and only if there is an intertwiner $`v`$ with $`\rho (v^{}v)0`$ and $`\sigma =\beta _v(\rho )`$; moreover, the set $`G_v`$ is the graph of $`\beta _v`$ and the collection of such sets gives a base for the topology. We call $`R()`$ the spectral relation of $``$. Also, $`R()`$ is reflexive if $`𝒟`$, is symmetric if $`=^{}`$ and is transitive if $``$ is a subalgebra.
A topological equivalence relation is a principal topological groupoid. If $`v,w`$ normalize $`𝒟`$, then $`G_{vw}=\{|\varphi \psi |:s(\varphi )=r(\varphi ),\varphi G_v,\psi G_w\}`$ and $`G_v^{}=\{|\varphi ^{}|:\varphi G_v\}`$. It follows that the topology on $`R(𝒞)`$ is compatible with the groupoid operations, so $`R(𝒞)`$ is a topological equivalence relation. We will sometimes write $`\sigma _𝒞\rho `$, or simply $`\sigma \rho `$, when $`(\sigma ,\rho )R(𝒞)`$.
We now show that the regularity of $`(𝒞,𝒟)`$ and the faithfulness of $`E`$ imply that the span of eigenfunctionals is weak-$``$ dense in $`^\mathrm{\#}`$, for any norm closed $`𝒟`$-bimodule $``$.
###### Proposition 4.21.
Suppose $`(𝒞,𝒟)`$ is a $`C^{}`$-diagonal, and let $`𝒞`$ be a norm-closed $`𝒟`$-bimodule. Then $`\mathrm{span}^1()`$ is $`\sigma (^\mathrm{\#},)`$-dense in $`^\mathrm{\#}`$.
###### Proof.
Let $`W`$ be the $`\sigma (^\mathrm{\#},)`$-closure of $`\mathrm{span}^1().`$ If $`W^\mathrm{\#}`$, then there exists a nonzero $`\sigma (^\mathrm{\#},)`$-continuous linear functional $`\psi `$ on $`^\mathrm{\#}`$ which annihilates $`W`$. Since $`\psi `$ is $`\sigma (^\mathrm{\#},)`$-continuous, there exists $`m`$ such that $`\psi (f)=f(m)`$ for all $`f^\mathrm{\#}`$. But then for every $`\varphi ^1()`$, we have $`\psi (\varphi )=\varphi (m)=0`$. Since $`^1(𝒞)`$ separates points (Theorem 4.9), there exists an eigenfunctional $`\overline{\varphi }^1(𝒞)`$ so that $`\overline{\varphi }(m)0`$. The restriction $`\varphi :=\overline{\varphi }|_{}`$ is a eigenfunctional on $``$, so $`0=\psi (\varphi )=\varphi (m)=\overline{\varphi }(m)0`$, a contradiction. Therefore, $`W=^\mathrm{\#}`$. ∎
We conclude this section with two applications of the results in this section. For our first application, we give a description of the $`C^{}`$-envelope of an algebra satisfying $`𝒟𝒜𝒞`$.
###### Theorem 4.22.
Let $`(𝒞,𝒟)`$ be a $`C^{}`$-diagonal and suppose $`𝒜`$ is a norm closed algebra satisfying $`𝒟𝒜𝒞`$. If $``$ is the $`C^{}`$-subalgebra of $`𝒞`$ generated by $`𝒜`$, then $``$ is the $`C^{}`$-envelope of $`𝒜`$.
If in addition, $`=𝒞`$, then $`R(𝒞)`$ is the topological equivalence relation generated by $`R(𝒜)`$.
###### Proof.
Let $`_e`$ be the $`C^{}`$-envelope of $`𝒜`$ and let $`j:𝒜_e`$ be the canonical embedding (see, for example, \[4, Section 4.3\]). Then there exists a unique $``$-epimorphism $`\pi :_e`$ such that $`\pi (a)=j(a)`$ for every $`a𝒜`$. In particular, $`\pi `$ is faithful on $`𝒟`$. Assume, to get a contradiction, that $`\pi `$ is not injective. Then $`\mathrm{ker}\pi `$ is a $`𝒟`$-bimodule and let $`x`$ be a non-zero element of $`\mathrm{ker}\pi `$. By Theorem 4.9, eigenfunctionals separate points, so there exists an element $`\varphi ^1(𝒞)`$ with $`\varphi (x)0`$. Writing $`\varphi =[v,\rho ]`$, Proposition 3.8 shows that $`u:=E_v(x)`$ is a nonzero normalizer belonging to $`\mathrm{ker}\pi `$. But then $`u^{}u`$ is a nonzero element of $`𝒟\mathrm{ker}\pi `$, a contradiction. Thus, $`\pi `$ is faithful on $``$, and hence $``$ is the $`C^{}`$-envelope of $`𝒜`$.
Suppose now that $`=𝒞`$. By Corollary 4.16, there is an inclusion $`^1(𝒜)^1(𝒞)`$. So $`R(𝒜)R(𝒞)`$, and hence $`R(𝒞)`$ contains the equivalence relation generated by $`R(𝒜)`$.
For the other direction, assume that $`(\sigma ,\rho )R(𝒞)`$. Then there is $`\varphi ^1(𝒞)`$ with source $`\rho `$ and range $`\sigma `$. Let $`𝒲`$ be the set of all finite products of intertwiners belonging to $`𝒜`$ or to $`𝒜^{}`$. Then $`𝒲𝒩_𝒟(𝒞)`$ and the set of finite sums from $`𝒲`$ is a $``$-algebra which, because $`𝒜`$ generates $`𝒞`$, is dense in $`𝒞`$. Hence there is some $`w𝒲`$ so that $`\varphi (w)0`$. By Proposition 4.8, $`\varphi =[w,\rho ]`$.
Suppose that $`w`$ factors as $`v_{2n}^{}v_{2n1}\mathrm{}v_2^{}v_1`$, where each $`v_i`$ is an intertwiner in $`𝒜`$. Let $`\rho _1=\rho `$ and for $`i=2,\mathrm{},2n`$, let $`\rho _i`$ be image of $`\rho `$ under conjugation by the rightmost $`i1`$ factors in the factorization of $`w`$. It follows that $`\varphi `$ is the product
$$[v_{2n},\rho _{2n}]^{}[v_{2n1},\rho _{2n1}]\mathrm{}[v_2,\rho _2]^{}[v_1,\rho _1]$$
and each $`[v_i,\rho _i]`$ is in $`^1(𝒜)`$. Thus, the equivalence relation generated by $`R(𝒜)`$ contains $`(\sigma ,\rho )`$. Similar arguments apply for the other possible factorizations of $`w`$, so the equivalence relation generated by $`R(𝒜)`$ contains $`R(𝒞)`$.
It remains to show that the usual topology on $`R(𝒞)`$ equals that generated by $`R(𝒜)`$, i.e., the smallest topology containing the topology of $`R(𝒜)`$ which makes $`R(𝒞)`$ into a topological equivalence relation.
As we noted in Definition 4.20, $`R(𝒞)`$ is already a topological equivalence relation, and so its topology contains the topology generated by $`R(𝒜)`$. Since the norm-closed span of $`𝒲`$ is $`𝒞`$, it follows that $`\{G_w:w𝒲\}`$ is a base for the topology of $`R(𝒞)`$, where, as before, $`G_w=\{\varphi R(𝒞):|\varphi (w)|>0\}`$. For a topological equivalence relation, the inverse map is a homeomorphism. Further, given two precompact open G-sets (i.e., a subset of $`R(𝒞)`$ on which the two natural projection maps into $`\widehat{𝒟}`$ are injective), one can show that their product is again a precompact open G-set (for example, adapt the proof of Proposition I.2.8 in ). Since, for $`v`$ an intertwiner in $`𝒜`$, each $`G_v`$ is a precompact open G-set in $`R(𝒜)`$, and each $`wW`$ is a finite product of such $`v`$’s and their inverses, it follows that each $`G_w`$, $`wW`$, is open in the topology generated by $`R(𝒜)`$. Thus, the topology generated by $`R(𝒜)`$ contains $`R(𝒞)`$. ∎
Our second application is an application of Theorem 4.9. We show show that inductive limits of $`C^{}`$-diagonals are again $`C^{}`$-diagonals, when the connecting maps satisfy a certain condition, which we now define. The difficulty in showing that these inductive limits are again $`C^{}`$-diagonals is in showing that the expectation is faithful, and this is where Theorem 4.9 provides a key tool.
###### Definition 4.23.
Given regular $`C^{}`$-inclusions $`(𝒞_i,𝒟_i)`$, $`i=1,2`$, and a $``$-homomorphism $`\pi :𝒞_1𝒞_2`$, we say $`\pi `$ is regular if $`\pi (𝒩(𝒞_1))𝒩(𝒞_2)`$.
Of course, if $`\pi `$ is regular, then $`\pi (𝒟_1)𝒟_2`$. Indeed, for $`D𝒟_1`$ with $`D0`$, $`D^{1/2}𝒩(𝒞_2)`$ and so $`\pi (D)=\pi (D^{1/2})1\pi (D^{1/2})𝒟_2`$.
###### Theorem 4.24.
Let $`(𝒞_\lambda ,𝒟_\lambda )`$, $`\lambda \mathrm{\Lambda }`$, be a directed net of regular $`C^{}`$-inclusions with regular $``$-monomorphisms $`\pi _{\lambda ,\mu }:𝒞_\mu C_\lambda `$. Then $`(\underset{}{\mathrm{lim}}(𝒞_\lambda ,\pi _{\lambda ,\mu }),\underset{}{\mathrm{lim}}(𝒟_\lambda ,\pi _{\lambda ,\mu }))`$ is a regular $`C^{}`$-inclusion. Moreover, if each $`(𝒞_\lambda ,𝒟_\lambda )`$ is a $`C^{}`$-diagonal, then so is $`(\underset{}{\mathrm{lim}}(𝒞_\lambda ,\pi _{\lambda ,\mu }),\underset{}{\mathrm{lim}}(𝒟_\lambda ,\pi _{\lambda ,\mu }))`$.
###### Proof.
The first part of the proof is routine. We regard the $`𝒞_\lambda `$ as a $``$-subalgebras of $`𝒞:=\underset{}{\mathrm{lim}}(𝒞_\lambda ,\pi _{\lambda ,\mu })`$ and identify $`\pi _{\lambda ,\mu }`$ with the inclusion map from $`𝒞_\lambda `$ to $`𝒞_\mu `$. Then $`𝒟=\underset{}{\mathrm{lim}}(𝒟_\lambda ,\pi _{\lambda ,\mu })`$ is a subalgebra of $`𝒞`$.
Given a normalizer $`v𝒞_\lambda `$, by the regularity of the inclusion maps, $`v`$ normalizes $`𝒟_\mu `$ for all $`\mu \lambda `$. Thus, $`v`$ normalizes $`𝒟`$ and so $`𝒩_{𝒟_\lambda }(𝒞_\lambda )𝒩_𝒟(𝒞)`$. Since $`𝒞`$ is the closed union of the $`𝒞_\lambda `$, and each $`𝒞_\lambda `$ is the span of $`𝒩_{𝒟_\lambda }(𝒞_\lambda )`$, $`𝒞`$ is regular in $`𝒟`$.
Given $`\rho \widehat{𝒟}`$, suppose $`\varphi `$ and $`\psi `$ are extensions of $`\rho `$ to states of $`𝒞`$. Then, for each $`\lambda \mathrm{\Lambda }`$, $`D_\lambda 𝒟`$ and so $`\varphi |_{𝒞_\lambda }`$ and $`\psi |_{𝒞_\lambda }`$ are extensions of the pure state $`\rho |_{𝒟_\lambda }\widehat{𝒟}_\lambda `$ and so agree on $`𝒞_\lambda `$. Since $`𝒞`$ is the closed union of the $`𝒞_\lambda `$, $`\varphi =\psi `$. Thus, $`(𝒞,𝒟)`$ is a regular $`C^{}`$-inclusion.
Let $`E:𝒞𝒟`$ be the expectation. By Theorem 4.9, $`𝔑:=\{x𝒞:E(x^{}x)=0\}`$ is an ideal of $`𝒞`$, and, if $`q:𝒞𝒞/𝔑`$ is the quotient map, then $`(q(𝒞),q(𝒟))`$ is a $`C^{}`$-diagonal. If $`x𝒞_\lambda `$ and $`E(x^{}x)=0`$, then $`\rho (x^{}x)=0`$ for all $`\rho \widehat{𝒟}`$. Since every $`\sigma \widehat{𝒟}_\lambda `$ has at least one extension to an element of $`\widehat{D}`$, we have $`E_\lambda (x^{}x)=0`$, where $`E_\lambda `$ is the expectation for $`(𝒞_\lambda ,𝒟_\lambda )`$. Thus, if $`(𝒞_\lambda ,𝒟_\lambda )`$ is a $`C^{}`$-diagonal, then we have $`x=0`$, that is, $`𝔑𝒞_\lambda =(0)`$, so $`q`$ is faithful on $`𝒞_\lambda `$.
Therefore, when each $`(𝒞_\lambda ,𝒟_\lambda )`$ is a $`C^{}`$-diagonal, $`q(𝒞)`$ contains isomorphic copies of each $`𝒞_\lambda `$, and when $`\lambda \mu `$, $`q(𝒞_\lambda )q(𝒞_\mu )`$. By the minimality of the inductive limit, $`q`$ is an isomorphism of $`𝒞`$ onto $`q(𝒞)`$, i.e. $`𝔑=0`$. Thus, $`(𝒞,𝒟)`$ is a $`C^{}`$-diagonal. ∎
## 5. Compatible Representations of $`C^{}`$-diagonals
Our goal in this section is to produce a faithful representation $`\pi `$ of a $`C^{}`$-diagonal $`(𝒞,𝒟)`$. Because we require the faithfulness of the expectation, we work with $`C^{}`$-diagonals instead of regular $`C^{}`$-inclusions.
Standing Assumptions for Section 5. We assume $`(𝒞,𝒟)`$ is a $`C^{}`$-diagonal. For $`(𝒞,𝒟)`$ a $`C^{}`$-diagonal, we write $`𝒜(𝒞,𝒟)`$ if $`𝒜𝒞`$ is a norm-closed subalgebra with $`𝒟𝒜`$.
For $`\rho \widehat{𝒟}`$, we use $`(_\rho ,\pi _\rho )`$ for the GNS representation of $`𝒞`$ associated to the (unique) extension of $`\rho `$.
Eigenfunctionals can be viewed as normal linear functionals on $`𝒞^{\mathrm{\#}\mathrm{\#}}`$ and we start by using the polar decomposition for such functionals to obtain a ‘minimal’ partial isometry for each eigenfunctional. Although these results are implicit in the development of dual groupoids (see \[40, p. 435\]), we give a (mostly) self-contained treatment.
Fix a norm-one eigenfunctional $`\varphi `$ on $`𝒞`$. By the polar decomposition for linear functionals (see \[42, Theorem III.4.2, Definition III.4.3\]), there is a partial isometry $`u^{}𝒞^{\mathrm{\#}\mathrm{\#}}`$ and a positive linear functional $`|\varphi |𝒞^\mathrm{\#}`$ so that $`\varphi =u^{}|\varphi |=|\varphi ^{}|u^{}.`$ Applying the characterization given in \[42, Proposition III.4.6\], we find that
(4)
$$r(\varphi )=|\varphi |\text{and}s(\varphi )=|\varphi ^{}|.$$
Moreover, $`uu^{}`$ and $`u^{}u`$ are the smallest projections in $`𝒞^{\mathrm{\#}\mathrm{\#}}`$ which satisfy,
$`u^{}us(\varphi )`$ $`=s(\varphi )u^{}u=s(\varphi )`$
and
$`uu^{}r(\varphi )`$ $`=r(\varphi )uu^{}=r(\varphi ).`$
###### Definition 5.1.
For $`\varphi ^1(𝒞)`$, we call the partial isometry $`u`$ above the partial isometry associated to $`\varphi `$ and denote it by $`v_\varphi `$. If $`\varphi \widehat{𝒟}`$, then $`u`$ is a projection and we denote it by $`p_\varphi `$.
Remark 5.2. The above equations show that $`v_\varphi ^{}v_\varphi =p_{s(\varphi )}`$ and $`v_\varphi v_\varphi ^{}=p_{r(\varphi )}`$. Moreover, given $`\varphi ^1(𝒞),`$ $`v_\varphi `$, Proposition 5.3 below implies that may be characterized as the unique minimal partial isometry $`w𝒞^{\mathrm{\#}\mathrm{\#}}`$ such that $`\varphi (w)>0`$.
Our first goal is to show that the initial and final projections of this partial isometry are minimal projections in $`𝒞^{\mathrm{\#}\mathrm{\#}}`$ and compressing by them gives $`\varphi `$, in the following sense.
###### Proposition 5.3.
For $`\rho \widehat{𝒟}`$, $`p_\rho =p_{\rho E}`$ is a minimal projection in $`𝒞^{\mathrm{\#}\mathrm{\#}}`$. For all $`\varphi ^1(𝒞)`$ and $`x𝒞^{\mathrm{\#}\mathrm{\#}}`$,
$$p_{r(\varphi )}xp_{s(\varphi )}=\varphi (x)v_\varphi .$$
###### Proof.
First, we show that $`p_\rho `$ is a minimal projection in $`𝒟^{\mathrm{\#}\mathrm{\#}}`$. We know that $`p_\rho `$ is the smallest projection in $`𝒟^{\mathrm{\#}\mathrm{\#}}`$ such that $`p_\rho \rho =\rho p_\rho =\rho `$. Suppose, to get a contradiction, that $`p_1,p_2`$ are nonzero projections in $`𝒟^{\mathrm{\#}\mathrm{\#}}`$ with $`0p_1,p_2p`$ and $`p_\rho =p_1+p_2`$.
If $`\rho (p_1)=0`$, then for $`d0`$, $`\rho (p_1d)=\rho (p_1dp_1)\rho (p_1)d=0`$ and so, for all $`d𝒟`$, $`\rho (p_1d)=0`$. But then $`p_2\rho =p\rho =\rho =\rho p_2`$, which yields $`p=p_2`$, contrary to hypothesis. Hence $`\rho (p_1)0`$ and, similarly, $`\rho (p_2)0`$.
This implies that $`\rho `$ can be written as a nontrivial convex combination of states on $`𝒟`$, for
$$\rho =\rho (p_1)\frac{p_1\rho }{\rho (p_1)}+\rho (p_2)\frac{p_2\rho }{\rho (p_2)}.$$
But this is a contradiction, since elements of of $`\widehat{𝒟}`$ are pure states.
Since $`p_\rho `$ is minimal in $`𝒟^{\mathrm{\#}\mathrm{\#}}`$, for $`d𝒟`$, $`p_\rho dp_\rho =\rho (d)p_\rho `$.
Now suppose that $`q𝒞^{\mathrm{\#}\mathrm{\#}}`$ is a projection with $`0<qp_\rho `$. Since $`q0`$, there exists a state $`g𝒞^\mathrm{\#}`$ such that $`g(q)>0`$. Define
$$f:=\frac{qgq}{g(q)}.$$
Then $`f`$ is a state on $`𝒞`$ with $`f(q)=1`$. As $`p_\rho q=q`$, we have for $`d𝒟`$,
$$f(d)=\frac{g(q(p_\rho dp_\rho )q)}{g(q)}=\rho (d).$$
Since pure states on $`𝒟`$ extend uniquely to pure states on $`𝒞`$, we conclude that $`f=\rho E`$.
If $`p_\rho `$ is not minimal, write $`p_\rho =q_1+q_2`$ where $`q_i𝒞^{\mathrm{\#}\mathrm{\#}}`$ are projections with $`0<q_1,q_2p_\rho `$. Apply the argument of the previous paragraph to find states $`h_1`$ and $`h_2`$ on $`𝒞`$ such that $`h_i(q_i)=1`$ and $`q_ih_iq_i=h_i`$. Since $`q_1q_2=0`$, $`h_1(q_2)=h_2(q_1)=0`$. But the previous paragraph shows that $`h_1=\rho E=h_2`$, contradicting the extension property. So $`p_\rho `$ is minimal in $`𝒞^{\mathrm{\#}\mathrm{\#}}`$.
The uniqueness of polar decompositions implies that $`p_\rho =p_{\rho E}`$.
To prove the second statement, first note that $`p_{r(\varphi )}𝒞^{\mathrm{\#}\mathrm{\#}}p_{s(\varphi )}`$ has dimension one, since $`p_{r(\varphi )},p_{s(\varphi )}`$ are minimal projections in $`𝒞^{\mathrm{\#}\mathrm{\#}}`$ and $`p_{r(\varphi )}v_\varphi p_{s(\varphi )}=v_\varphi 0`$. Hence there is a linear functional $`g`$ on $`𝒞`$ such that for every $`x𝒞^{\mathrm{\#}\mathrm{\#}}`$, $`p_{r(\varphi )}xp_{s(\varphi )}=g(x)v_\varphi `$. Then $`g|_𝒞`$ is an eigenfunctional with the same source and range as $`\varphi `$. Since $`g(v_\varphi )=\varphi (v_\varphi )`$, $`g=\varphi `$. ∎
Recall (\[42, Lemma III.2.2\]) that any $``$-representation $`\pi `$ of a $`C^{}`$-algebra $`𝒞`$ has a unique extension to a $``$-representation $`\stackrel{~}{\pi }:𝒞^{\mathrm{\#}\mathrm{\#}}\pi (𝒞)^{\prime \prime }`$, continuous from the $`\sigma (𝒞^{\mathrm{\#}\mathrm{\#}},𝒞^\mathrm{\#})`$-topology to the $`\sigma `$-weak topology on $`\pi (𝒞)^{\prime \prime }`$, i.e., the $`\sigma ((),()_{})`$-topology. Let $`=\mathrm{ker}\stackrel{~}{\pi }𝒞^{\mathrm{\#}\mathrm{\#}}`$. By the continuity of $`\stackrel{~}{\pi }`$, $``$ is $`\sigma (𝒞^{\mathrm{\#}\mathrm{\#}},𝒞^\mathrm{\#})`$ closed, and hence (see \[41, Proposition 1.10.5\]) there exists a unique central projection $`P𝒞^{\mathrm{\#}\mathrm{\#}}`$ such that $`=𝒞^{\mathrm{\#}\mathrm{\#}}(IP).`$ Further, $`\stackrel{~}{\pi }|_{𝒞^{\mathrm{\#}\mathrm{\#}}P}`$ is one-to-one (see \[41, Definition 1.21.14\]) and is onto $`()`$. The projection $`P`$ is called the support projection for $`\pi `$.
Recall that $`\sigma ,\rho \widehat{𝒟}`$ have $`(\rho ,\sigma )R(𝒞)`$ if and only if there is $`\varphi ^1(𝒞)`$ with $`r(\varphi )=\rho `$ and $`s(\varphi )=\sigma `$. For brevity, we write $`\sigma \rho `$ in this case.
###### Proposition 5.4.
If $`\rho \widehat{𝒟}`$, then $`\pi _\rho (𝒟)^{\prime \prime }`$ is an atomic MASA in $`(_\rho )`$ and the support projection for $`\pi _\rho `$ is $`_{\sigma \rho }p_\sigma `$. Moreover, the map from $`\pi _\rho (𝒞)`$ to $`\pi _\rho (𝒟)`$ given by $`\pi _\rho (c)\pi _\rho (E(c))`$ is well-defined and extends to a faithful normal expectation $`\stackrel{~}{E}:\pi _\rho (𝒞)^{\prime \prime }\pi _\rho (𝒟)^{\prime \prime }`$.
###### Proof.
Let $`\pi =\pi _\rho `$ and $`=_\rho `$. Since the extension of $`\rho `$ to $`𝒞`$ is pure, the representation $`\pi `$ is irreducible, so $`\pi (𝒞)^{\prime \prime }=()`$.
Recall from \[19, Corollary 1\] that if $`:=\{x𝒞:\rho (x^{}x)=0\}`$, then $`𝒞/`$ is complete relative to the norm induced by the inner product $`x+,y+=\rho (y^{}x)`$, and thus $`=𝒞/`$. Our first task is to obtain a convenient orthonormal basis for $``$ and, towards this end, we require the following observation.
If $`v,w𝒩(𝒞)`$ with $`\rho (v^{}v)=1`$ and $`\rho (v^{}w)0`$, then
(5)
$$w[v,\rho ](w)v.$$
To see this, let $`w_1=\rho (w^{}v)w`$. Then $`\rho (v^{}w_1)>0`$, so that by Corollary 4.5, $`[v,\rho ](w_1)=[w_1,\rho ](w_1)`$, and thus $`\rho (v^{}w_1)=\rho (w_1^{}w_1)^{1/2}`$. It follows that $`|\rho (v^{}w)|=\rho (w^{}w)^{1/2}.`$ A calculation then shows that $`w[v,\rho ](w)v=w\rho (v^{}w)v`$, as required.
Choose a set $`𝒵𝒞`$ of normalizers such that for each $`z𝒵`$, $`\rho (z^{}z)=1`$ and the map $`zr([z,\rho ])`$ is a bijection of $`𝒵`$ onto $`𝒪:=\{\sigma \widehat{𝒟}:\sigma \rho \}`$.
We claim that if $`X=_{j=1}^nw_j𝒞`$ with each $`w_j𝒩_𝒟(𝒞)`$, then
(6)
$$X+=\underset{z𝒵}{}[z,\rho ](X)(z+).$$
To see that the summation is well-defined, first observe that if, for some $`z𝒵`$, $`\rho (z^{}w_j)0`$, then $`r([w_j,\rho ])=r([z,\rho ])𝒪`$. Thus, for each $`j`$, $`\{z𝒵:\rho (z^{}w_j)0\}`$ is a singleton, and so $`\{z𝒵:\rho (z^{}X)0\}`$ is finite. To prove the equality, interchange the order of summation and use (5) as follows:
$`{\displaystyle \underset{z𝒵}{}}[z,\rho ](X)(z+)`$ $`={\displaystyle \underset{z𝒵}{}}{\displaystyle \underset{j=1}{\overset{n}{}}}[z,\rho ](w_j)(z+)`$
$`={\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{z𝒵}{}}[z,\rho ](w_j)(z+)`$
$`={\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{\begin{array}{c}z𝒵\\ \rho (z^{}w_j)0\end{array}}{}}[z,\rho ](w_j)(z+)`$
$`={\displaystyle \underset{j=1}{\overset{n}{}}}w_j+`$ by (5)
$`=X+.`$
Next, we show that $`\{z+:z𝒵\}`$ is an orthonormal basis for $``$. By Corollary 4.5, it is an orthonormal set. Given any $`Y𝒞`$ and $`\epsilon >0`$, we may find a finite sum of normalizers $`X`$ so that $`XY_𝒞<\epsilon `$. Then $`XY+__\rho ^2=\rho ((XY)^{}(XY))XY_𝒞^2<\epsilon ^2`$. As $`X+`$ is in $`\mathrm{span}𝒵`$ and $`\epsilon `$ is arbitrary, $`Y+\overline{\mathrm{span}𝒵}`$, so that $`\{z+:z𝒵\}`$ is an orthonormal basis.
For $`d𝒟`$ and $`z𝒵`$, using (5)
(7)
$$\pi (d)(z+)=dz+=[z,\rho ](dz)z+=\rho (z^{}dz)(z+).$$
Thus, $`\pi (d)`$ is diagonal with respect to the basis $`\{z+\}_{z𝒵}`$. Fixing $`z𝒵`$, let $`\mathrm{\Lambda }=\{d𝒟:d0,\text{ and }r([z,\rho ])(d)=1\}`$. Then $`\mathrm{\Lambda }`$ is a directed set under the ordering $`d_1d_2`$ if and only if $`d_1d_2`$. It is easy to see that the net $`\{\pi (d)\}_{d\mathrm{\Lambda }}`$ decreases to the rank-one projection onto $`z+`$. Thus $`\pi (𝒟)^{\prime \prime }`$ is an atomic MASA.
Fix $`z𝒵`$ and let $`P_z`$ be the orthogonal projection of $``$ onto $`z+`$. Then a calculation shows that for $`x𝒞`$,
(8)
$$P_z\pi (x)P_z=\rho (z^{}xz)P_z.$$
Since $`r([z,\rho ])`$ is the vector state corresponding to $`z+`$ and $`\stackrel{~}{\pi }`$ is normal, (8) holds when $`\pi `$ is replaced by $`\stackrel{~}{\pi }`$ and $`x𝒞^{\mathrm{\#}\mathrm{\#}}`$. As $`p_{r([z,\rho ])}`$ is a minimal projection in $`𝒞^{\mathrm{\#}\mathrm{\#}}`$, $`\stackrel{~}{\pi }(p_{r([z,\rho ])})=P_z`$. If $`Q=_{\sigma \rho }p_\sigma `$, then $`\stackrel{~}{\pi }(Q)=I`$. Letting $`P`$ be the support projection of $`\pi `$, this implies that $`PQ`$. If $`\sigma \rho `$, then $`p_\sigma `$ is a a minimal projection satisfying $`\stackrel{~}{\pi }(p_\sigma )0`$. Thus, $`p_\sigma P`$ and so $`QP`$. This shows that $`Q`$ is the support projection for $`\pi `$ and that $`Q𝒟^{\mathrm{\#}\mathrm{\#}}`$.
Finally, (7) implies that for $`x𝒞`$,
$`\pi (E(x))(z+)`$ $`=\rho (z^{}E(x)z)(z+M)=r([z,\rho ])(E(x))(z+)`$
$`=r([z,\rho ])(x)(z+)=\rho (z^{}xv)(z+).`$
Hence, for $`x𝒞`$, we obtain
$$\pi (E(x))=\underset{z𝒵}{}P_z\pi (x)P_z.$$
If $`\stackrel{~}{E}:()()`$ is defined by $`\stackrel{~}{E}(T)=_{z𝒵}P_zTP_z`$, this shows that $`\stackrel{~}{E}`$ is a faithful normal conditional expectation of $`()=\pi (𝒞)^{\prime \prime }`$ onto $`\pi (𝒟)^{\prime \prime }`$ satisfying $`\stackrel{~}{E}(\pi (x))=\pi (E(x))`$ for $`x𝒞`$. ∎
We next record a simple consequence of the proof of Proposition 5.4.
###### Corollary 5.5.
Suppose $`\rho \widehat{𝒟}`$ and $`\varphi ^1(𝒞)`$ satisfies $`\rho s(\varphi )`$. Then there exist unit orthogonal unit vectors $`\xi ,\eta _\rho `$ such that for every $`x𝒞`$, $`\varphi (x)=\pi _\rho (x)\xi ,\eta `$.
###### Proof.
With the same notation as in the proof of Proposition 5.4, there exists $`z,w𝒵`$ such that $`r([z,\rho ])=s(\varphi )`$ and $`r([w,\rho ])=r(\varphi )`$. Let $`\xi _0=z+`$ and $`\eta =w+`$. For $`x𝒞`$, let $`\psi (x)=\pi (x)\xi _0,\eta `$. Observe that $`\psi (wz^{})=\rho (w^{}wz^{}z)=1`$, so $`\psi `$ is nonzero. Also, for $`x𝒞`$ and $`d𝒟`$, we have
$$\psi (xd)=\rho (w^{}xdz)=\frac{\rho (w^{}xdzz^{}z)}{\rho (z^{}z)}=\rho (w^{}xz)\frac{\rho (z^{}dz)}{\rho (z^{}z)}=\psi (x)s(\varphi )(d).$$
Similarly, $`\psi (dx)=r(\varphi )(d)\psi (x)`$. Thus $`\psi `$ is an eigenfunctional with the same range and source as $`\varphi `$. Hence, there exists $`t𝕋`$ so that $`\varphi =t\psi `$. Take $`\xi =t\xi _0`$. ∎
###### Definition 5.6.
Given a $`C^{}`$-diagonal $`(𝒞,𝒟)`$, a representation $`\pi `$ of $`𝒞`$ is $`𝒟`$-compatible, (or simply compatible) if $`\pi (𝒟)^{\prime \prime }`$ is a MASA in $`\pi (𝒞)^{\prime \prime }`$ and there exists a faithful conditional expectation $`\stackrel{~}{E}:\pi (𝒞)^{\prime \prime }\pi (𝒟)^{\prime \prime }`$ such that for every $`x𝒞`$, $`\stackrel{~}{E}(\pi (x))=\pi (E(x))`$.
Proposition 5.4 shows that the GNS representation of $`𝒞`$ associated to an element of $`\widehat{𝒟}`$ is compatible.
Remark 5.7. Suppose that $`X`$ is an r-discrete locally compact principal groupoid with unit space $`X^0`$ and let $`\sigma `$ be a 2-cocycle. Then Drinen shows that $`(C_r(X,\sigma ),C_0(X^0))`$ is a $`C^{}`$-diagonal. We expect that when $`\lambda ^u`$ is a Haar system on $`X`$, and $`\mu `$ is a measure on $`X^0`$, the induced representation $`\mathrm{Ind}(\mu )`$ of $`C_r(X,\sigma )`$ (see for example, \[26, page 44\]) is a compatible representation of $`(C_r(X,\sigma ),C_0(X^0))`$, and it would not be surprising if every compatible representation for $`(C_r(X,\sigma ),C_0(X^0))`$ arises in this way. We do not pursue this issue here, however.
###### Lemma 5.8.
Let $`\rho ,\sigma \widehat{𝒟}`$. Then $`\rho \sigma `$ if and only if the GNS representations $`\pi _\sigma `$ and $`\pi _\rho `$ are unitarily equivalent.
###### Proof.
If $`U(_\rho ,_\sigma )`$ is a unitary such that $`U\pi _\rho U^{}=\pi _\sigma `$, the fact that $`𝒞/_\rho =_\rho `$ allows us to find $`X𝒞`$ such that $`U^{}(I+_\sigma )=X+_\rho `$. Then for all $`x𝒞`$, we have $`\sigma (x)=\pi _\sigma (x)(I+_\sigma ),(I+_\sigma )=\rho (X^{}xX)`$. As $`\sigma `$ and $`\rho `$ are normal states on $`𝒞^{\mathrm{\#}\mathrm{\#}}`$, this equality is valid for $`x𝒞^{\mathrm{\#}\mathrm{\#}}`$ as well; in particular, $`1=\rho (X^{}p_\sigma X)=\sigma (p_\sigma )`$. For $`x𝒞`$, define $`\varphi (x)=\rho (X^{}p_\sigma x)`$. Then $`\varphi 0`$ because $`\varphi (X)=1`$ and $`\varphi `$ is an eigenfunctional with $`s(\varphi )=\rho `$ and $`r(\varphi )=\sigma `$. Hence $`\rho \sigma `$. Conversely, if $`\rho \sigma `$, then find a normalizer $`v`$ with $`\rho (v^{}v)=1`$ and $`\sigma (x)=\rho (v^{}xv)`$ for every $`x𝒞`$. Then $`\sigma (x)=\pi _\rho (x)v+_\rho ,v+_\rho `$. By Kadison’s Transitivity Theorem, there exists a unitary $`V𝒞`$ such that $`V^{}(1+N_\rho )=v+_\rho .`$ Then $`\pi _\sigma =V\pi _\rho V^{}.`$
###### Theorem 5.9.
If $`𝒳\widehat{𝒟}`$ contains exactly one element from each equivalence class in $`R(𝒞)`$, then $`\pi =_{\rho 𝒳}\pi _\rho `$ on $`=_{\rho 𝒳}_\rho `$ is a faithful compatible representation of $`𝒞`$ on $`()`$ and $`\pi (𝒟)^{\prime \prime }`$ is an multiplicity-free atomic MASA in $`()`$.
Moreover, if $`𝒜(𝒞,𝒟)`$ and $`P`$ is the support projection for $`\rho `$, then $`\mathrm{ker}\stackrel{~}{\pi }|_{𝒜^{\mathrm{\#}\mathrm{\#}}}=P^{}𝒜^{\mathrm{\#}\mathrm{\#}}`$ and $`\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})`$ is a CSL algebra.
###### Proof.
For $`\rho 𝒳`$, let $`P_\rho 𝒞^{\mathrm{\#}\mathrm{\#}}`$ be the support projection of $`\pi _\rho `$. Then for $`\sigma \widehat{𝒟}`$,
$$P_\rho p_\sigma =\{\begin{array}{cc}p_\sigma \hfill & \text{if }\sigma \rho \hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
Thus $`\stackrel{~}{\pi }(p_\sigma )`$ is a minimal projection for every $`\sigma \widehat{𝒟}`$. Since
$$I_{}=\underset{\rho 𝒳}{}\stackrel{~}{\pi }(P_\rho )=\underset{\rho 𝒳}{}\underset{\sigma \rho }{}\stackrel{~}{\pi }(p_\sigma ),$$
which is a sum of minimal projections, $`\pi (𝒟)^{\prime \prime }`$ is an multiplicity-free atomic MASA in $`()`$. If, for each $`\rho 𝒳`$, $`E_\rho ^{}:(_\rho )\pi _\rho (𝒟)^{\prime \prime }`$ is the expectation on $`(_\rho )`$ induced by $`E`$, the map $`E^{}=_{\rho 𝒳}E_\rho ^{}`$ is faithful and satisfies $`E^{}\pi =\pi E`$ so $`\pi `$ is compatible. For any $`x𝒞`$ such that $`\pi (x^{}x)=0`$,
$$0=E^{}(\pi (x^{}x))=\pi (E(x^{}x)),$$
so $`E(x^{}x)=0`$ since $`\pi `$ is faithful on $`𝒟`$. Thus, as $`E`$ is faithful, we conclude that $`x^{}x=0`$, hence $`\pi `$ is faithful on $`𝒞`$.
Suppose $`𝒜(𝒞,𝒟)`$. As $`P=_{\rho 𝒳}P_\rho `$, where $`P_\rho `$ is the support projection for $`\pi _\rho `$, Proposition 5.4 implies that $`P𝒟^{\mathrm{\#}\mathrm{\#}}𝒜^{\mathrm{\#}\mathrm{\#}}`$. Thus $`\mathrm{ker}\stackrel{~}{\pi }|_{𝒜^{\mathrm{\#}\mathrm{\#}}}=P^{}𝒜^{\mathrm{\#}\mathrm{\#}}`$.
We claim that $`\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})`$ is weak-$``$ closed in $`()`$. By the Krein-Smullian Theorem, it suffices to prove the (norm-closed) unit ball of $`\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})`$ is weak-$``$ closed. Suppose $`(y_\lambda )`$ is a net in $`\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})`$ converging weak-$``$ to $`y()`$ with with $`y_\lambda 1`$ for all $`\lambda `$. For each $`\lambda `$, choose $`x_\lambda 𝒜^{\mathrm{\#}\mathrm{\#}}`$ with $`\stackrel{~}{\pi }(x_\lambda )=y_\lambda `$. Since $`\stackrel{~}{\pi }`$ is isometric on $`P𝒞^{\mathrm{\#}\mathrm{\#}}`$, the net $`(Px_\lambda )`$ in $`𝒜^{\mathrm{\#}\mathrm{\#}}`$ satisfies $`Px_\lambda 1`$ for every $`\lambda `$. Hence a subnet $`Px_{\lambda _\mu }`$ converges weak-$``$ to some $`x𝒜^{\mathrm{\#}\mathrm{\#}}`$. Then $`\stackrel{~}{\pi }(x)=y`$, as required.
Since $`\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})`$ contains the atomic MASA $`\pi (𝒟)^{\prime \prime }`$, $`:=\mathrm{Lat}(\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}}))`$ is an atomic CSL. We claim that $`\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})=\mathrm{Alg}`$. As $`\mathrm{Alg}`$ is the largest algebra whose lattice of invariant subspaces is $``$, $`\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})\mathrm{Alg}`$.
To obtain the reverse inclusion, first observe that for minimal projections $`p,q\pi (𝒟)^{\prime \prime }`$, $`q\mathrm{Alg}p=q\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})p`$. Indeed, the subspaces $`\overline{\mathrm{Alg}p}`$ and $`\overline{\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})p}`$ are each the smallest element of $``$ containing the range of $`p`$, so the two subspaces coincide. Thus $`q\overline{\mathrm{Alg}p}=q\overline{\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})p}`$, which yields the observation.
Let $`𝔸`$ be the set of minimal projections of $`\pi (𝒟)^{\prime \prime }`$. Given $`Y\mathrm{Alg}`$, we may write $`Y`$ as the weak-$``$ convergent sum, $`Y=_{q,p𝔸}qYp`$. As each $`qYp\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})`$, and $`\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})`$ is weak-$``$ closed, $`Y\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})`$, as desired. ∎
Muhly, Qiu and Solel \[24, Theorem 4.7\] prove that if $`𝒞`$ is nuclear and $`𝒜(𝒞,𝒟)`$ with $`𝒜`$ triangular, then the expectation $`E|_𝒜`$ is a homomorphism. The connection with CSL algebras provided by Theorem 5.9 allows us to remove the hypothesis of nuclearity.
###### Theorem 5.10.
If $`𝒜`$ is a triangular subalgebra of the $`C^{}`$-diagonal $`(𝒞,𝒟)`$, then $`E|_𝒜`$ is a homomorphism of $`𝒜`$ onto $`𝒟`$.
###### Proof.
Let $`\pi `$ be the faithful compatible representation of $`𝒞`$ provided by Theorem 5.9, and again write $`\mathrm{Alg}`$ for $`\stackrel{~}{\pi }(𝒜^{\mathrm{\#}\mathrm{\#}})`$. Theorem 5.9 also shows that $`\stackrel{~}{\pi }(𝒟^{\mathrm{\#}\mathrm{\#}})`$ is a MASA in $`()`$. We claim that $``$ is multiplicity free. To show this, we prove that $`\mathrm{Alg}(\mathrm{Alg})^{}=\stackrel{~}{\pi }(𝒟^{\mathrm{\#}\mathrm{\#}})`$. Clearly, $`\stackrel{~}{\pi }(𝒟^{\mathrm{\#}\mathrm{\#}})\mathrm{Alg}(\mathrm{Alg})^{}`$. To show the reverse implication, suppose $`q_1,q_2\stackrel{~}{\pi }(𝒟^{\mathrm{\#}\mathrm{\#}})`$ are distinct nonzero minimal projections and $`q_2(\mathrm{Alg})q_1(0).`$ It suffices to show that
(9)
$$q_1(\mathrm{Alg})q_2=(0).$$
We may find $`\rho ,\sigma \widehat{𝒟}`$ so that $`q_1=\stackrel{~}{\pi }(p_\rho )`$ and $`q_2=\stackrel{~}{\pi }(p_\sigma )`$, where $`p_\rho `$ and $`p_\sigma `$ are as in Definition 5.1. Since $`\pi (𝒜)`$ is weak-$``$ dense in $`\mathrm{Alg}`$, we see that $`p_\sigma 𝒜p_\rho (0)`$. Therefore, $`B_{\sigma ,\rho }`$ is nonzero on $`𝒜`$, so that there exists an eigenfunctional $`\varphi ^1(𝒜)`$ with $`s(\varphi )=\rho `$ and $`r(\varphi )=\sigma `$.
We claim that $`\varphi ^{}`$ vanishes on $`𝒜`$. Indeed, as $`q_1q_2`$, $`\rho \sigma `$, so that we may find a normalizer $`v𝒜`$ so that $`(v^{}v)(vv^{})=0`$ and $`\varphi =[v,\rho ]`$. Suppose to obtain a contradiction, that $`\varphi ^{}=[v^{},\sigma ]`$ does not vanish on $`𝒜`$, and let $`y𝒜`$ satisfy $`\varphi ^{}(y)0`$. Proposition 3.8 shows that $`w:=v^{}E(vy)`$ is a nonzero element of $`𝒜`$. But since $`v𝒜`$ and $`E(vy)𝒟`$, we also have $`w^{}𝒜`$. However,
$`ww^{}`$ $`=v^{}E(vy)E(vy)^{}vE(vy)^2v^{}v,`$
$`w^{}w`$ $`=E(vy)^{}vv^{}E(vy)E(vy)^2vv^{}.`$
Thus $`w`$ is a nonnormal element of $`𝒜𝒜^{}`$, violating triangularity of $`𝒜`$.
Therefore, $`\varphi ^{}`$ vanishes on $`𝒜`$, so that $`p_\rho 𝒜p_\sigma =(0)`$. Applying $`\stackrel{~}{\pi }`$, we obtain $`q_1(\mathrm{Alg})q_2=(0)`$. Therefore, $`\mathrm{Alg}`$ is multiplicity free as desired.
Each minimal projection $`e`$ in $`\stackrel{~}{\pi }(𝒟^{\mathrm{\#}\mathrm{\#}})`$ is the difference of elements of $``$, and hence the compression $`xexe`$ is a homomorphism on $`\mathrm{Alg}`$. The extension of $`\stackrel{~}{E}`$ of $`E`$ to all of $`()`$ is faithful and normal, and is the sum of such compressions. Thus, $`\stackrel{~}{E}`$ is a homomorphism on $`\mathrm{Alg}`$ and, by restriction, on $`𝒜`$. ∎
We describe the maximal ideals of $`𝒜`$ and identify $`\mathrm{ker}E|_𝒜`$ in algebraic terms.
###### Proposition 5.11.
Let $`(𝒞,𝒟)`$ be a $`C^{}`$-diagonal and $`𝒜(𝒞,𝒟)`$ be triangular. The map $`JJ𝒟`$ is a bijection between the proper maximal ideals of $`𝒜`$ and the proper maximal ideals of $`𝒟`$. Further,
$$\mathrm{ker}E|_𝒜=\{J𝒜:J\text{ is a maximal ideal of }𝒜\}.$$
###### Proof.
Let $`J`$ be a proper maximal ideal of $`𝒜`$. For any $`x𝒞`$, $`E(x)`$ belongs to the norm-closed convex hull of $`\{gxg^1:g𝒟,g\text{ is unitary}\}`$ (Theorem 3.6). Since $`J`$ is a closed ideal, we see that $`E(J)J𝒟E(J)`$, so $`E(J)=J𝒟`$. Since $`J`$ is proper, $`J𝒟𝒟`$. Hence there exists $`\rho \widehat{𝒟}`$ such that $`\mathrm{ker}\rho J𝒟`$. The unique extension of $`\rho `$ to $`𝒜`$ is $`\rho E|_𝒜`$, so we have $`J\mathrm{ker}(\rho E|_𝒜)`$. Since $`J`$ is maximal, $`J=\mathrm{ker}\rho E|_𝒜`$. Therefore, $`J𝒟=\mathrm{ker}\rho `$, which is a proper maximal ideal of $`𝒟`$.
To show that the map $`JJ𝒟`$ is a bijection, we need only consider proper maximal ideals. The previous paragraph shows that if $`J_1`$ is another proper maximal ideal of $`𝒜`$ and $`J_1𝒟=JD`$, then $`J=J_1`$. Also, if $`K𝒟`$ is a proper maximal ideal, then $`K=\mathrm{ker}\rho `$ for some $`\rho \widehat{𝒟}`$, and then $`J:=\mathrm{ker}\rho E_𝒜`$ is a proper maximal ideal of $`𝒜`$ with $`J𝒟=K`$. Thus the map is onto.
For $`x𝒜`$, we have $`x\mathrm{ker}E|_𝒜`$ if and only if $`x\mathrm{ker}\rho E|_𝒜`$ for every $`\rho \widehat{𝒟}`$. The map above shows this is equivalent to $`x\{J𝒜:J\text{ is a maximal ideal of }𝒜\}`$. ∎
## 6. Invariance under Diagonal-Preserving Isomorphisms
As an application of the results obtained so far, we will show that coordinate systems are preserved under isomorphisms of algebras which preserve the diagonal. These results extend those for isometric isomorphisms and we compare our results with them.
###### Definition 6.1.
For $`i=1,2`$, suppose $`(𝒞_i,𝒟_i)`$ are regular $`C^{}`$-inclusions and that $`𝒜_i`$ are subalgebras with $`𝒜_i(𝒞_i,𝒟_i)`$, i.e., $`𝒜_i𝒞_i`$ is a norm-closed subalgebra with $`𝒟_i𝒜_i`$. We say that a (bounded) isomorphism $`\theta :𝒜_1𝒜_2`$ is diagonal preserving if $`\theta (𝒟_1)=𝒟_2`$.
Remarks 6.2.
1. If $`\theta `$ is diagonal preserving, then $`\theta |_{𝒟_1}`$ is a $``$-isomorphism of $`𝒟_1`$ onto $`𝒟_2`$.
2. While we are interested in coordinates for modules, when studying invariance properties, we note that it suffices to consider subalgebras of regular $`C^{}`$-inclusions. This is because of the well-known “$`2\times 2`$ matrix trick.” If $`(𝒞,𝒟)`$ is a regular $`C^{}`$-inclusion, so is $`(M_2(𝒞),𝒟𝒟)`$. For a $`𝒟`$-bimodule $`𝒞`$, let $`𝔗()`$ be the subalgebra of $`M_2(𝒞)`$,
$$𝔗():=\{\left[\begin{array}{cc}d_1& m\\ 0& d_2\end{array}\right]:d_1,d_2𝒟,m\}$$
contained in $`(M_2(𝒞),𝒟𝒟)`$. An isomorphism of $`𝒟_i`$-bimodules $`_i`$ ($`i=1,2`$), that is, a bounded map $`\theta :_1_2`$ together with an isomorphism $`\alpha :𝒟_1𝒟_2`$ satisfying $`\theta (dme)=\alpha (d)\theta (m)\alpha (e)`$ ($`d,e𝒟_1,m`$) can be equivalently described as a diagonal preserving isomorphism of $`𝔗(_1)`$ onto $`𝔗(_2)`$.
We have noted that $`\theta ^\mathrm{\#}`$ is a bicontinuous isomorphism from $`(𝒜_2)`$ onto $`(𝒜_1)`$ (Proposition 2.9). The next result shows that normalizing $`\theta ^\mathrm{\#}`$ pointwise gives an isomorphism of the norm-one eigenfunctionals.
###### Theorem 6.3.
For $`i=1,2`$, let $`(𝒞_i,𝒟_i)`$ be regular $`C^{}`$-inclusions, let $`𝒜_i(𝒞_i,𝒟_i)`$, and suppose $`\theta :𝒜_1𝒜_2`$ is a bounded diagonal-preserving isomorphism. There exists a bicontinuous isomorphism of coordinate systems $`\gamma :^1(𝒜_1)^1(𝒜_2)`$ given by
$$\gamma (\varphi )=\frac{\varphi \theta ^1}{\varphi \theta ^1}.$$
Moreover, if $`\varphi ^1(𝒜_1)`$ is written as $`[v,\rho ]`$, then $`(\rho \theta ^1)(\theta (v)^{}\theta (v))`$ is nonzero, and
$$\gamma (\varphi )=[\theta (v),\rho \theta ^1].$$
When necessary for clarity, we use $`\gamma _\theta `$ to denote the dependence of $`\gamma `$ on $`\theta `$.
Remark 6.4. As a special case, if $`\theta :𝒜_1𝒜_2`$ is a contractive isomorphism, then it is diagonal preserving, as its restriction to $`𝒟_1`$ is then a $``$-homomorphism. Thus, Theorem 6.3 extends previous work for isometric isomorphisms of TAF algebras and of subalgebras of (nuclear) groupoid $`C^{}`$-algebras; see \[34, Theorem 3\] and \[25, Theorem 2.1\].
###### Proof of Theorem 6.3.
Given $`\varphi ^1(𝒜_1)`$, write $`\varphi =[v,\rho ]`$ where $`v𝒜_1`$ is an intertwiner and $`\rho (v^{}v)0`$. That $`r(\gamma (\varphi ))=r(\varphi )\theta ^1`$ and $`s(\gamma (\varphi ))=s(\varphi )\theta ^1`$ follows from the definition of $`\gamma `$.
Clearly $`\theta (v)`$ is a $`𝒟_2`$-intertwiner, so $`\theta (v)^{}\theta (v)𝒟_2`$. For $`\rho \widehat{𝒟}_1`$ with $`\rho (v^{}v)0`$,
$`(\rho \theta ^1)(\theta (v)^{}\theta (v))`$ $`=inf\{\theta (d)^{}\theta (v)^{}\theta (v)\theta (d):d𝒟_1,\rho (d)=1\}`$
$`=inf\{\theta (vd)^2:d𝒟_1,\rho (d)=1\}`$
$`\theta ^1^2inf\{vd^2:d𝒟_1,\rho (d)=1\}`$
$`=\theta ^1^2\rho (v^{}v)0.`$
Simple calculations show that $`[\theta (v),\rho \theta ^1]`$ and $`\varphi \theta ^1/\varphi \theta ^1`$ are elements of $`_{𝒟_2}^1(𝒜_2)`$ with the same range. Since both are positive on $`\theta (v)`$,
$$\gamma (\varphi )=[\theta (v),\rho \theta ^1].$$
This formula implies that $`\gamma (\varphi _1\varphi _2)=\gamma (\varphi _1)\gamma (\varphi _2)`$ whenever $`\varphi _1\varphi _2`$ is defined, so that $`\gamma `$ is an algebraic isomorphism of coordinate systems.
To show continuity, let $`(\varphi _\lambda )`$ be a net in $`^1(𝒜_1)`$ converging weak-$``$ to $`\varphi ^1(𝒜_1)`$. Write $`\varphi =[v,\rho ]`$ for some intertwiner $`v𝒜_1`$ and $`\rho \widehat{𝒟}_1`$. Since $`\rho (v^{}v)0`$, there exists $`d𝒟_1`$ and a neighborhood of $`\rho `$, $`G\widehat{𝒟}_1`$, with $`d0`$ and $`\sigma (d^{}v^{}vd)=1`$ for all $`\sigma G`$. Since $`\rho (d)0`$, $`\varphi =[vd,\rho ]`$. Replacing $`v`$ with $`vd`$, we may assume that $`\sigma (v^{}v)=1`$ for every $`\sigma `$ in $`G\widehat{𝒟}_1`$, a neighborhood of $`\rho `$.
Since $`\rho _\lambda :=s(\varphi _\lambda )`$ converges weak-$``$ to $`s(\varphi )=\rho `$, by deleting the first part of the net, we may assume that $`\rho _\lambda G`$ and $`\varphi _\lambda (v)0`$ for all $`\lambda `$. By Proposition 4.8, there exist scalars $`t_\lambda 𝕋`$ such that $`\varphi _\lambda =[t_\lambda v,\rho _\lambda ].`$ Since
$$\overline{t}_\lambda =[t_\lambda v,\rho _\lambda ](v)=\varphi _\lambda (v)\varphi (v)=1,$$
and $`[\theta (v),\rho _\lambda \theta ^1]`$ converges weak-$``$ to $`[\theta (v),\rho \theta ^1]`$, we conclude
$$\gamma (\varphi _\lambda )=[t_\lambda \theta (v),\rho _\lambda \theta ]=\overline{t}_\lambda [\theta (v),\rho _\lambda \theta ^1][\theta (v),\rho \theta ^1]=\gamma (\varphi ).$$
Thus, $`\gamma `$ is continuous. Similarly, $`\gamma ^1`$ is continuous. ∎
We give several applications of Theorem 6.3 to isomorphisms of subalgebras. The following corollary is immediate.
###### Corollary 6.5.
For $`i=1,2,3`$, let $`(𝒞_i,𝒟_i)`$ be regular $`C^{}`$-inclusions, and let $`𝒜_i(𝒞_i,𝒟_i)`$ be norm-closed algebras. For $`j=1,2`$, let $`\theta _j:𝒟_j𝒟_{j+1}`$ be bounded diagonal preserving isomorphisms. Then $`\gamma _{\theta _2\theta _1}=\gamma _{\theta _2}\gamma _{\theta _1}.`$
###### Definition 6.6.
Suppose $`(𝒞,𝒟)`$ is a regular $`C^{}`$-inclusion and $`𝒜(𝒞,𝒟)`$. By a (1-)cocycle on $`^1(𝒜)`$, we mean a map $`c:^1(𝒜)_{}`$, the group of nonzero complex numbers under multiplication, satisfying, for all composable elements $`\varphi ,\psi ^1(𝒜)`$,
$$c(\varphi \psi )=c(\varphi )c(\psi ).$$
###### Corollary 6.7.
Suppose $`(𝒞,𝒟)`$ is a regular $`C^{}`$-inclusion, $`𝒜(𝒞,𝒟)`$ is a norm-closed algebra, and $`\theta :𝒜𝒜`$ is a bounded automorphism fixing $`𝒟`$ elementwise. Assume further that $`\gamma _\theta `$ is the identity map on $`^1(𝒜)`$. Then, for all $`\varphi ^1(𝒜)`$,
(10)
$$\varphi (\theta (x))=\varphi (x)\varphi \theta .$$
If $`c:^1(𝒜)`$ is defined by $`c(\varphi )=\varphi \theta `$, then $`c`$ is a positive cocycle on $`^1(𝒜)`$.
If, in addition, $`\theta `$ is isometric, then $`\theta =\mathrm{id}_𝒜`$.
Remark 6.8. In essence, (10) shows that $`\theta `$ is given by multiplication by a cocycle on $`^1(𝒜)`$. Thus, we will call $`\theta `$ a cocycle automorphism.
###### Proof.
Formula (10) follows immediately from Theorem 6.3 applied to $`\theta ^1`$. If $`\theta `$ is isometric, then (10) shows that $`\varphi (x)=\varphi (\theta (x))`$ for all $`\varphi ^1(𝒜)`$. By Theorem 4.9, $`^1(𝒜)`$ separates points and so $`x=\theta (x)`$ for all $`x𝒜`$.
It remains to show that $`c`$ is a cocycle. Observe that if $`v𝒞`$ is a $`𝒟`$-intertwiner, then $`v^{}\theta (v)𝒟`$. Indeed, for a self-adjoint $`d𝒟`$, using Remark 3(2), we may find $`d^{}𝒟`$ with $`d^{}`$ also self-adjoint and so that $`vd=d^{}v`$. Then $`v^{}\theta (v)d=v^{}d^{}\theta (v)=dv^{}\theta (d)`$, so that $`v^{}\theta (v)`$ commutes with the self-adjoint elements of $`𝒟`$ and hence belongs to $`𝒟`$.
Finally, suppose for $`i=1,2`$, that $`\varphi _i=[v_i,\rho _i]^1(𝒜)`$ such that the product $`\varphi _1\varphi _2`$ is defined. Then $`\varphi _1\varphi _2=[v_1v_2,\rho _2]`$, and $`\rho _1=r([v_2,\rho _2])`$, i.e., $`\rho _1(x)=\rho _2(v_2^{}xv_2)/\rho _2(v_2^{}v_2)`$. Using (10) and these facts, we have
$`c(\varphi _1\varphi _2)`$ $`={\displaystyle \frac{[v_1v_2,\rho _2](\theta (v_1v_2))}{[v_1v_2,\rho _2](v_1v_2)}}`$
$`={\displaystyle \frac{\rho _2(v_2^{}v_1^{}\theta (v_1)\theta (v_2))}{\rho _2(v_2^{}v_1^{}v_1v_2)}}`$
$`={\displaystyle \frac{\rho _2(v_2^{}v_2)\rho _2(v_2^{}v_1^{}\theta (v_1)\theta (v_2))}{\rho _2(v_2^{}v_2)\rho _2(v_2^{}v_1^{}v_1v_2)}}`$
$`={\displaystyle \frac{\rho _1(v_1^{}\theta (v_1))\rho _2(v_2^{}\theta (v_2))}{\rho _2(v_2^{}v_1^{}v_1v_2)}}`$
$`={\displaystyle \frac{\rho _1(v_1^{}\theta (v_1))\rho _2(v_2^{}\theta (v_2))}{\rho _1(v_1^{}v_1)\rho _2(v_2^{}v_2)}}=c(\varphi _1)c(\varphi _2)`$
For our next application, we need two technical lemmas.
###### Lemma 6.9.
For $`i=1,2`$, let $`(𝒞_i,𝒟_i)`$ be regular $`C^{}`$-inclusions and let $`_i(𝒞_i,𝒟_i)`$ be selfadjoint. If $`\gamma :^1(_1)^1(_2)`$ is an isomorphism of coordinate systems, then, for all $`\varphi ^1(_1)`$, $`\gamma (\varphi ^{})=\gamma (\varphi )^{}`$.
###### Proof.
Since each $`_i`$ is selfadjoint, $`\tau ^{}^1(_i)`$ if $`\tau ^1(_i)`$, and so $`\gamma (\varphi ^{})`$ and $`\gamma (\varphi )^{}`$ are defined. As $`\gamma (\varphi ^{})`$ and $`\gamma (\varphi )^{}`$ have the same range and domain, by Corollary 4.11, there is $`\lambda 𝕋`$ with $`\gamma (\varphi ^{})=\lambda \gamma (\varphi )^{}`$. As $`\varphi \varphi ^{}=r(\varphi )`$, we have
$$\gamma (r(\varphi ))=\gamma (\varphi \varphi ^{})=\gamma (\varphi )\gamma (\varphi ^{})=\gamma (\varphi )(\lambda \gamma (\varphi )^{})=\lambda r(\gamma (\varphi )).$$
As $`\gamma (r(\varphi ))=r(\gamma (\varphi )`$, the result follows. ∎
###### Lemma 6.10.
For $`i=1,2`$, let $`(𝒞_i,𝒟_i)`$ be regular $`C^{}`$-inclusions, let $`𝒜_i(𝒞_i,𝒟_i)`$ be regular, and let $`\pi :𝒞_1C_2`$ be a $``$-isomorphism with $`\pi (𝒟_1)=𝒟_2`$. If $`\gamma _\pi `$ maps $`^1(𝒜_1)`$ into $`^1(𝒜_2)`$, then $`\pi (𝒜_1)𝒜_2`$.
###### Proof.
Since each $`𝒜_i`$ is regular, it is enough to show that for an intertwiner $`v𝒜_1`$, we have $`\pi (v)𝒜_2`$. Clearly, $`\pi (v)`$ is an intertwiner in $`𝒞_2`$, and it is easy to see that the closed $`𝒟_2`$-bimodule generated by $`\pi (v)`$ is isometrically isomorphic to $`\overline{|\pi (v)|𝒟_2}`$ via the map $`\pi (v)d|\pi (v)|d`$. Given $`\rho \widehat{𝒟}_2`$ with $`\rho (\pi (v^{}v))0`$, we have $`\gamma ([v,\rho \pi ^1])=[\pi (v),\rho ]^1(𝒜_2)`$. Applying Proposition 3.8, we conclude that there exists $`d𝒟_2`$ with $`d0`$, $`\rho (d)=1`$, and $`\pi (v)d𝒜_2`$. Given $`\epsilon >0`$, let $`f_\epsilon 𝒟`$ be chosen so that $`0f_\epsilon I`$, $`\widehat{f}_\epsilon `$ is compactly supported in $`G:=\{\rho \widehat{𝒟}_2:\rho (|\pi (v)|)>0\}`$ and $`\pi (v)\pi (v)f_\epsilon <\epsilon `$. Compactness of $`F_\epsilon :=\overline{\mathrm{supp}\mathrm{f}_\epsilon }`$ ensures that there exist $`d_1,\mathrm{},d_n𝒟_2`$ such that $`\pi (v)d_i𝒜_2`$ and $`\rho (_{i=1}^nd_i)>0`$ for every $`\rho F_\epsilon `$. Hence there exists an element $`g𝒟_2`$ such that for every $`\rho F_\epsilon `$, $`\rho (g_{i=1}^nd_i)=1`$. Then $`\pi (v)f_\epsilon =_{i=1}^n\pi (v)d_igf_\epsilon 𝒜_2`$. Letting $`\epsilon 0`$, we conclude that $`\pi (v)𝒜_2`$ as well. ∎
Recall that a subalgebra $`𝒜𝒞`$ is said to be Dirichlet if $`𝒜+𝒜^{}`$ is norm dense in $`𝒞`$. This implies that $`^1(𝒞)=^1(𝒜)^1(𝒜)^{}`$. To see this, let $`\varphi ^1(𝒞)`$. By density, $`\varphi `$ does not vanish on one of $`𝒜`$ or $`𝒜^{}`$ and hence is in either $`^1(𝒜)`$ or $`^1(𝒜)^{}`$.
###### Theorem 6.11.
For $`i=1,2`$, let $`(𝒞_i,𝒟_i)`$ be regular $`C^{}`$-inclusions and let $`𝒜_i`$ be Dirichlet subalgebras with $`𝒟_i𝒜_i𝒞_i`$. Consider the following statements.
1. $`𝒜_1`$ and $`𝒜_2`$ are isometrically isomorphic.
2. There exists a bounded isomorphism $`\theta :𝒜_1𝒜_2`$ such that $`\theta (𝒟_1)=𝒟_2`$.
3. $`^1(𝒜_1)`$ and $`^1(𝒜_2)`$ are isomorphic coordinate systems.
4. $`^1(𝒞_1)`$ and $`^1(𝒞_2)`$ are isomorphic twists and the isomorphism maps $`^1(𝒜_1)`$ onto $`^1(𝒜_2)`$.
Then statement $`(n)`$ implies statement $`(n+1)`$, $`n=1,2,3`$.
If, in addition, $`(𝒞_i,𝒟_i)`$ are $`C^{}`$-diagonals and $`𝒜_i`$ are regular, then the four statements are equivalent.
###### Proof.
That (1) implies (2) is obvious. To show that (2) implies (3), let $`\alpha =\theta |_{𝒟_1}`$, a $``$-isomorphism, and apply Theorem 6.3.
To show (3) implies (4), suppose $`\gamma :^1(𝒜_1)^1(𝒜_2)`$ is a continuous isomorphism. We extend $`\gamma `$ to a map, call it $`\delta `$, on all of $`^1(𝒞_1)`$ by mapping $`\varphi ^1(𝒜_1)^{}=^1(𝒜_1^{})`$ to $`\gamma (\varphi ^{})^{}`$. By Lemma 6.9, $`\delta `$ is well-defined on $`^1(𝒜_1𝒜_1^{})`$, and hence well-defined on all of $`^1(𝒞_1)`$.
Since the adjoint map is continuous, $`\delta `$ is a continuous homomorphism restricted to $`^1(𝒜_1)`$ or to $`^1(𝒜_1^{})`$. By Corollary 4.16, these are open sets in $`^1(𝒞_1)`$, and so $`\delta `$ is continuous on their union, $`^1(𝒞_1)`$. Applying the same argument to $`\delta ^1,`$ we see that $`\delta ^1`$ is continuous as well.
The restriction of $`\delta `$ to $`^1(𝒜_1^{})`$ or to $`^1(𝒜_1)`$ is a homomorphism. To show that $`\delta `$ is a homomorphism, fix $`\varphi ,\psi ^1(𝒜_1)`$. We claim that if $`\varphi ^{}\psi `$ is defined, then $`\delta (\varphi ^{}\psi )=\delta (\varphi ^{})\delta (\psi )`$ and if $`\varphi \psi ^{}`$ is defined, then $`\delta (\varphi \psi ^{})=\delta (\varphi )\delta (\psi )^{}`$. We only show the first equality; the proof of the second is similar.
As $`\delta `$ maps $`\widehat{𝒟}_1`$ to $`\widehat{𝒟}_2`$, we have, for $`\varphi ^1(𝒜_1)`$,
(11)
$$\delta (\varphi \varphi ^{})=\gamma (r(\varphi ))=r(\gamma (\varphi ))=\gamma (\varphi )\gamma (\varphi ^{}).$$
Let $`\eta =\varphi ^{}\psi `$. If $`\eta ^1(𝒜_1)`$, we have $`\varphi \eta =\varphi \varphi ^{}\psi =\psi `$, and hence $`\gamma (\varphi )\gamma (\eta )=\gamma (\psi )`$. Multiplying each side of this equality by $`\gamma (\varphi ^{})=\delta (\varphi )^{}`$ and using (11) yields $`\delta (\varphi ^{}\psi )=\delta (\varphi )^{}\delta (\psi )`$. If $`\eta ^1(𝒜_1^{})`$, then apply the previous argument to $`\eta ^{}=\psi ^{}\varphi `$ and take adjoints to obtain the equation. Thus $`\delta `$ is a homomorphism and part (4) holds.
Finally, if $`(𝒞_i,𝒟_i)`$ are $`C^{}`$-diagonals and (4) holds, then Kumjian’s theorem implies that there is an isomorphism $`\mathrm{\Phi }:𝒞_1𝒞_2`$ with $`\mathrm{\Phi }(𝒟_1)=𝒟_2`$. Since $`\mathrm{\Phi }`$ is induced by an isomorphism of twists mapping $`^1(𝒜_1)`$ onto $`^1(𝒜_2)`$, Lemma 6.10 shows that $`\mathrm{\Phi }(𝒜_1)=𝒜_2`$. Thus, we obtain (1) with $`\mathrm{\Phi }|_{𝒜_1}`$ the isometric isomorphism. ∎
Remark 6.12. Theorem 6.11 is related to a result of Muhly, Qiu and Solel, \[25, Theorem 2.1\]. Their result, expressed in terms of eigenfunctionals, takes the following form. When the $`𝒞_i`$ are nuclear and $`(𝒞_i,𝒟_i)`$ are $`C^{}`$-diagonals, the following are equivalent, for subalgebras $`𝒜_i(𝒞_i,𝒟_i)`$ so that $`𝒜_i`$ generates $`𝒞_i`$ as $`C^{}`$-algebras (with no Dirichlet hypothesis):
1. there is an isometric isomorphism $`\theta :𝒜_1𝒜_2`$ (necessarily, $`\theta (𝒟_1)=𝒟_2`$);
2. there is a coordinate system isomorphism, $`\gamma :^1(𝒜_1)^1(𝒜_2)`$ which extends to a coordinate system isomorphism $`\gamma ^{}:^1(𝒞_1)^1(𝒞_2)`$;
3. there is a $``$-isomorphism $`\tau :𝒞_1𝒞_2`$ such that $`\tau |_{𝒜_1}`$ is an isomorphism of $`𝒜_1`$ onto $`𝒜_2`$.
Theorem 6.11 extends the Muhly-Qiu-Solel result to not-necessarily-isometric diagonal preserving isomorphisms, assuming the Dirichlet condition instead of the hypothesis in (b) that $`\gamma `$ extends to an isomorphism of $`^1(𝒞_1)`$ onto $`^1(𝒞_2)`$. Example 6.11 shows that in the absence of the Dirichlet condition, isomorphisms of coordinate systems need not extend to isomorphisms of the enveloping twists. Thus, the hypothesis that $`\gamma `$ extends in (b) is essential.
Also, since Theorem 6.11 did not use the Spectral Theorem for Bimodules \[24, Theorem 4.1\], we do not need the $`𝒞_i`$ to be nuclear. In Theorem 8.9, we prove the full Muhly-Qiu-Solel result, without requiring nuclearity.
Example 6.13. Without the Dirichlet hypothesis, an isomorphism of coordinate systems need not induce an isometric isomorphism of the algebras. Let $`(𝒞,𝒟)=(M_4(),𝒟_4)`$, where $`𝒟_4`$ is the algebra of diagonal matrices.
Let $`𝒜_1=𝒜_2`$ be the algebra
$$𝒜:=\{\left(\begin{array}{cccc}t_{11}& 0& t_{13}& t_{14}\\ 0& t_{22}& t_{23}& t_{24}\\ 0& 0& t_{33}& 0\\ 0& 0& 0& t_{44}\end{array}\right):t_{ij}\}.$$
The automorphism
$$\left(\begin{array}{cccc}t_{11}& 0& t_{13}& t_{14}\\ 0& t_{22}& t_{23}& t_{24}\\ 0& 0& t_{33}& 0\\ 0& 0& 0& t_{44}\end{array}\right)\left(\begin{array}{cccc}t_{11}& 0& t_{13}& t_{14}\\ 0& t_{22}& t_{23}& t_{24}\\ 0& 0& t_{33}& 0\\ 0& 0& 0& t_{44}\end{array}\right)$$
is not isometric, and induces an automorphism of $`^1(𝒜)`$ which does not extend to an automorphism of $`^1(𝒞)`$.
The Dirichlet condition can be removed if one assumes a continuous section from $`R(𝒞_i)`$ into $`^1(𝒞_i)`$. Since TAF algebras always admit such sections, the following result generalizes \[37, Theorem 7.5\].
###### Theorem 6.14.
Suppose for $`i=1,2`$, $`(𝒞_i,𝒟_i)`$ are $`C^{}`$-diagonals and $`𝒜_i(𝒞_i,𝒟_i)`$ are norm closed subalgebras such that $`𝒜_i`$ generates $`𝒞_i`$ as a $`C^{}`$-algebra. Consider the following statements:
1. $`𝒜_1`$ and $`𝒜_2`$ are isometrically isomorphic;
2. there exists a bounded isomorphism $`\theta :𝒜_1𝒜_2`$ such that $`\theta (𝒟_1)=𝒟_2`$;
3. $`R(𝒜_1)`$ and $`R(𝒜_2)`$ are isomorphic topological binary relations,
4. $`R(𝒞_1)`$ and $`R(𝒞_2)`$ are isomorphic topological equivalence
relations, and the isomorphism maps $`R(𝒜_1)`$ onto $`R(𝒜_2)`$.
Then for $`i=1,2,3`$, statement $`(i)`$ implies statement $`(i+1)`$.
If, in addition, $`𝒜_i`$ are regular and there exist continuous sections $`h_i:R(𝒞_i)^1(𝒞_i)`$, then the statements are equivalent.
###### Proof.
That (1)$``$(2) is obvious and (2)$``$(3) follows as in the proof of Theorem 6.11.
Suppose (3) holds. By Proposition 4.22, the isomorphism of $`R(𝒜_1)`$ onto $`R(𝒜_2)`$ extends uniquely to an isomorphism of $`R(𝒞_1)`$ onto $`R(𝒞_2)`$, so (4) holds.
To complete the proof, we show that, when the $`𝒜_i`$ are regular and there exist continuous sections $`h_i:R(𝒞_i)^1(𝒞_i)`$, then (4) implies (1). The existence of the sections and (4) gives a coordinate system isomorphism $`\gamma `$ of $`^1(𝒞_1)`$ and $`^1(𝒞_2)`$ such that $`\gamma |_{^1(𝒜_1)}`$ is an isomorphism of $`^1(𝒜_1)`$ onto $`^1(𝒜_2)`$. By Kumjian’s theorem, there is a (regular) $``$-isomorphism $`\pi `$ of $`(𝒞_1,𝒟_1)`$ onto $`(𝒞_2,𝒟_2)`$.
Finally, Lemma 6.10 implies $`\pi (𝒜_1)=𝒜_2`$. ∎
Remarks 6.15. A modification of Example 6.11 shows that in general, there may exist a section for $`R(𝒜)`$ which cannot be extended to a section of $`R(𝒞)`$. Thus, one cannot replace the hypothesis of a section for $`R(𝒞_i)`$ with a hypothesis of a section for $`R(𝒜_i)`$ in Theorem 6.14.
To drop the Dirichlet condition in Theorem 6.11 without assuming the existence of a continuous section, we would need to replace $`\gamma `$ with a new isomorphism from $`^1(𝒜_1)`$ to $`^1(𝒜_2)`$ that could extend to the twist of $`𝒞_1`$, which would be larger than $`^1(𝒜_1)^1(𝒜_1)^{}`$. In particular, while isometric bimodule maps on finite relations are always $``$-extendible (a key ingredient in ), this is not true for general bimodule maps \[7, Theorem 1.2, Proposition 1.4\]. Indeed, based on , there will homological obstructions to be considered.
## 7. Invariance under General Isomorphisms
In this section and the next, we come to the core of the paper, the study of bounded isomorphisms of triangular algebras which need not map the diagonal to the diagonal. The principal result of this section is Theorem 7.7, which shows that such isomorphisms induce algebraic isomorphisms of the corresponding coordinate systems. We would particularly like to know if this algebraic isomorphism is always continuous. We can prove that it is continuous in certain cases, Theorem 7.11, extending results of Donsig-Hudson-Katsoulis.
Standing Assumption for Section 7. For $`i=1,2`$, let $`(𝒞_i,𝒟_i)`$ be $`C^{}`$-diagonals, and let $`𝒜_i(𝒞_i,𝒟_i)`$ be (norm-closed) triangular subalgebras. Let $`E_i:𝒞_i𝒟_i`$ be the unique conditional expectations. Suppose $`\theta :𝒜_1𝒜_2`$ is a bounded isomorphism.
Our first task is to show that $`\theta `$ induces an algebraic isomorphism $`\gamma :^1(𝒜_1)^1(𝒜_2)`$. We have been unable to show that $`\gamma `$ is continuous in general. Since we do not assume $`\theta (𝒟_1)=𝒟_2`$, it is not possible to use Theorem 6.3, so we proceed along different lines.
###### Definition 7.1.
Define $`\alpha :𝒟_1𝒟_2`$ by $`\alpha (d)=E_2(\theta (d))`$.
###### Proposition 7.2.
The map $`\alpha `$ is a $``$-isomorphism of $`𝒟_1`$ onto $`𝒟_2`$.
###### Proof.
Theorem 5.10 shows $`E_2|_{𝒜_2}`$ is a homomorphism; hence $`\alpha `$ is a homomorphism. Since any algebraic isomorphism of commutative $`C^{}`$-algebras is a $``$-isomorphism, it suffices to show that $`\alpha `$ is bijective.
Since $`E_1`$ is idempotent, $`𝒟_1\mathrm{ker}E_1=\{0\}`$, so $`𝒜_1=𝒟_1+\mathrm{ker}E_1|_{𝒜_1}`$ is a direct sum decomposition. By Proposition 5.11, $`\theta (\mathrm{ker}E_1|_{𝒜_1})=\mathrm{ker}E_2|_{𝒜_2}`$, so that we have two direct sum decompositions of $`𝒜_2`$:
$$𝒜_2=𝒟_2+\mathrm{ker}E_2|_{𝒜_2}=\theta (𝒟_1)+\mathrm{ker}E_2|_{𝒜_2}.$$
Therefore, $`\mathrm{ker}E_2|_{\theta (𝒟_1)}=\mathrm{ker}\alpha `$ is trivial. If $`d𝒟_2`$ we may write $`d=x+y`$ where $`x\theta (𝒟_1)`$ and $`y\mathrm{ker}E_2`$. Then $`E_2(x)=d`$, so $`\alpha `$ is onto. ∎
Given Banach spaces $`X`$ and $`Y`$, and a bounded linear map $`R:XY`$, the double transpose map $`R^{\mathrm{\#}\mathrm{\#}}:X^{\mathrm{\#}\mathrm{\#}}Y^{\mathrm{\#}\mathrm{\#}}`$ is a norm-continuous extension of $`R`$ which is also $`\sigma (X^{\mathrm{\#}\mathrm{\#}},X^\mathrm{\#})`$$`\sigma (Y^{\mathrm{\#}\mathrm{\#}},Y^\mathrm{\#})`$ continuous. In light of our standing assumptions, $`\theta ^{\mathrm{\#}\mathrm{\#}}:𝒜_1^{\mathrm{\#}\mathrm{\#}}𝒜_2^{\mathrm{\#}\mathrm{\#}}`$ and $`\alpha ^{\mathrm{\#}\mathrm{\#}}:𝒟_1^{\mathrm{\#}\mathrm{\#}}𝒟_2^{\mathrm{\#}\mathrm{\#}}`$ are also isomorphisms. Similarly, the $`E_i^{\mathrm{\#}\mathrm{\#}}`$ are homomorphisms of $`𝒜_i^{\mathrm{\#}\mathrm{\#}}`$ onto $`𝒟_i^{\mathrm{\#}\mathrm{\#}}`$. For notational ease, we will sometimes identify the double transpose map with the original map. Thus we often write $`\theta `$ or $`\alpha `$ instead of $`\theta ^{\mathrm{\#}\mathrm{\#}}`$ or $`\alpha ^{\mathrm{\#}\mathrm{\#}}`$.
Remark 7.3. Notice that $`𝒜_2`$ may be regarded as a $`𝒟_1`$-bimodule in two ways: for $`d,e𝒟_1`$ and $`x𝒜_2`$, define $`d_\alpha x_\alpha e:=\alpha (d)x\alpha (e)`$ and $`d_\theta x_\theta e:=\theta (d)x\theta (e)`$. When these two modules are (boundedly) isomorphic, methods similar to those used in the proof of Theorem 6.3 show there exists an isomorphism of the coordinate systems $`^1(𝒜_1)`$ and $`^1(𝒜_2)`$. Unfortunately, we do not know in general whether the $`\alpha `$ and $`\theta `$ module actions of $`𝒟_1`$ on $`𝒜_2`$ are isomorphic. However, there is enough structure present to show that these modules are “virtually” isomorphic.
###### Definition 7.4.
Let $`G=𝒰(𝒟_1)`$ be the group of unitary elements of $`𝒟_1`$, regarded as a discrete abelian group and fix, once and for all, an invariant mean $`\mathrm{\Lambda }`$ on $`G`$.
Define elements $`S,T𝒜_2^{\mathrm{\#}\mathrm{\#}}`$ by requiring that for every $`f𝒜_2^\mathrm{\#}`$,
$$f(S)=\underset{gG}{\mathrm{\Lambda }}f(\theta (g)\alpha (g^1))\text{and}f(T)=\underset{gG}{\mathrm{\Lambda }}f(\alpha (g)\theta (g^1)).$$
Notice that $`S\overline{\text{co}}^\sigma \{\theta (g)\alpha (g^1):gG\}`$ and $`T\overline{\text{co}}^\sigma \{\alpha (g)\theta (g^1):gG\}`$, where $`\overline{\text{co}}^\sigma Z`$ is the $`\sigma (𝒜_2^{\mathrm{\#}\mathrm{\#}},𝒜_2^\mathrm{\#})`$-closed convex hull of the set $`Z`$. (We implicitly embed $`𝒜_2`$ into $`𝒜_2^{\mathrm{\#}\mathrm{\#}}`$ using the canonical inclusion.)
We next collect some properties of $`S`$ and $`T`$. Of particular interest to us is the fact that they intertwine $`\alpha (𝒟_1)`$ and $`\theta (𝒟_1)`$.
###### Proposition 7.5.
For $`S`$ and $`T`$ as above, we have
1. For every $`d𝒟_1`$,
$$T\theta (d)=\alpha (d)T\text{and}S\alpha (d)=\theta (d)S.$$
2. $`E_2^{\mathrm{\#}\mathrm{\#}}(S)=I=E_2^{\mathrm{\#}\mathrm{\#}}(T)`$ and $`E_1^{\mathrm{\#}\mathrm{\#}}(\theta ^1(S))=I=E_1^{\mathrm{\#}\mathrm{\#}}(\theta ^1(T))`$.
3. Given $`\rho \widehat{𝒟}_1`$, let $`p=p_\rho `$ (see Definition 5.1). Then
$`\alpha (p)T`$ $`=\alpha (p)T\theta (p)=T\theta (p)=\alpha (p)\theta (p)\text{and}`$
$`\theta (p)S`$ $`=\theta (p)S\alpha (p)=S\alpha (p)=\theta (p)\alpha (p).`$
4. For all $`x𝒜_2`$ and for all $`\varphi ^1(𝒜_1)`$, $`\varphi (\theta ^1(STxST))=\varphi (\theta ^1(x)).`$
###### Proof.
If $`hG=𝒰(𝒟_1)`$, then $`T\theta (h)=\alpha (h)T`$ follows from the invariance of $`\mathrm{\Lambda }`$. Indeed, for every $`f𝒜_2^\mathrm{\#}`$ we have,
$`f(\alpha (h)T)`$ $`=\mathrm{\Lambda }_g(f\alpha (h))(\alpha (g)\theta (g)^1)`$
$`=\mathrm{\Lambda }_gf(\alpha (hg)\theta (hg)^1\theta (h))`$
$`=\mathrm{\Lambda }_gf(\alpha (g)\theta (g^1)\theta (h))`$
$`=\mathrm{\Lambda }_g(\theta (h)f)(\alpha (g)\theta (g^1))`$
$`=f(T\theta (h)).`$
Since $`𝒜_2^\mathrm{\#}`$ separates points of $`𝒜_2^{\mathrm{\#}\mathrm{\#}}`$, we see that $`T\theta (h)=\alpha (h)T`$ for every $`hG`$. But the span of $`G`$ is norm dense in $`𝒟_1`$, which yields $`T\theta (d)=\alpha (d)T`$ for $`d𝒟_1`$. The proof that $`\theta (d)S=S\alpha (d)`$ is similar.
To prove $`E_1^{\mathrm{\#}\mathrm{\#}}(\theta ^1(S))=I`$, first observe that weak-$``$ continuity of $`E_1^{\mathrm{\#}\mathrm{\#}}`$ (and $`\theta _{}^{1}{}_{}{}^{\mathrm{\#}\mathrm{\#}}`$) implies
$$E_1^{\mathrm{\#}\mathrm{\#}}(\theta ^1(S))E_1^{\mathrm{\#}\mathrm{\#}}\overline{\text{co}}^\sigma \{g\theta ^1(\alpha (g^1))\}=\overline{\text{co}}^\sigma \{E_1(g\theta ^1(\alpha (g^1)))\}.$$
Modifying the proof of Proposition 7.2 yields $`\alpha ^1=E_1\theta ^1`$ and $`E_1(g\theta ^1(\alpha (g^1)))=I`$ for each $`gG`$, and the equality follows. The remaining equalities in part (2) have similar proofs.
For part (3), the first two equalities follow from statement (1), as $`p`$ is a $`\sigma (𝒜_1^{\mathrm{\#}\mathrm{\#}},𝒜_1^\mathrm{\#})`$-limit of elements of $`𝒟_1`$. For the third equality, first observe that for $`gG`$, Theorem 5.3 gives $`\alpha (p)\alpha (g)=\rho (g)\alpha (p)`$; similarly, $`\theta (g)^1\theta (p)=\rho (g^1)\theta (p).`$ Hence,
$$\alpha (p)\alpha (g)\theta (g)^1\theta (p)=\alpha (p)\theta (p).$$
Since $`T\overline{\text{co}}^\sigma \{\alpha (g)\theta (g)^1:gG\}`$, we find that
$$\alpha (p)T\theta (p)\overline{\text{co}}^\sigma \{\alpha (p)\alpha (g)\theta (g)^1\theta (p):gG\}.$$
This set is a singleton, so $`\alpha (p)T\theta (p)=\alpha (p)\theta (p)`$. The proofs of the equalities involving $`S`$ are similar.
For part (4), fix $`\varphi ^1(𝒜_1)`$, and let $`q`$ and $`p`$ be the minimal projections in $`𝒟_1^{\mathrm{\#}\mathrm{\#}}`$ corresponding to $`r(\varphi )`$ and $`s(\varphi )`$, respectively. Part (1) implies that $`p`$ and $`q`$ commute with $`\theta ^1(ST)`$ and by part (2), we have $`r(\varphi )(\theta ^1(ST))=r(\varphi )(E_1(\theta ^1(ST)))=r(\varphi )(I)=1`$. Hence by Proposition 5.3, $`q\theta ^1(ST)=q`$. Likewise $`p\theta ^1(ST)=p`$. As, again by Proposition 5.3, $`\varphi (a)=\varphi (qap)`$, we have
$`\varphi (\theta ^1(STxST))`$ $`=\varphi (q\theta ^1(ST)\theta ^1(x)\theta ^1(ST))p)`$
$`=\varphi (q\theta ^1(x)p)=\varphi (\theta ^1(x)),`$
as desired. ∎
We now obtain a bijective mapping between the eigenfunctionals of $`𝒜_1`$ and those of $`𝒜_2`$.
###### Proposition 7.6.
For $`\varphi (𝒜_1)`$, let
$$f=T(\varphi \theta ^1)S.$$
Then $`f`$ is an eigenfunctional for $`𝒜_2`$ with $`r(f)=r(\varphi )\alpha ^1`$ and $`s(f)=s(\varphi )\alpha ^1`$.
Moreover, $`\varphi \theta ^1=SfT`$.
###### Proof.
For clarity, let $`\psi =\varphi \theta ^1`$. For all $`d,e𝒟_1`$ and $`x𝒜_2`$,
$`f(\alpha (d)x\alpha (e))`$ $`=\psi (S\alpha (d)x\alpha (e)T)`$
$`=\psi (\theta (d)SxT\theta (e))`$
$`=r(\varphi )(d)\psi (SxT)s(\varphi )(e)`$
$`=(r(\varphi )\alpha ^1)(\alpha (d))f(x)(s(\varphi )\alpha ^1)(\alpha (e)),`$
showing $`f`$ is an eigenfunctional with the claimed range and source.
The last equality follows from part (4) of Proposition 7.5. ∎
We now show the existence of an algebraic isomorphism between the coordinate systems which is the non-diagonal preserving analog of Theorem 6.3.
###### Theorem 7.7.
The map $`\gamma :^1(𝒜_1)^1(𝒜_2)`$ given by
$$\gamma (\varphi )=\frac{T(\varphi \theta ^1)S}{T(\varphi \theta ^1)S}$$
is an algebraic isomorphism of coordinate systems such that for every $`\rho \widehat{𝒟}_1`$, $`\gamma (\rho )=\rho \alpha ^1`$.
###### Proof.
The fact that $`\gamma `$ is a bijection between $`^1(𝒜_1)`$ and $`^1(𝒜_2)`$ such that for every $`\varphi ^1(𝒜_1)`$, $`r(\gamma (\varphi ))=r(\varphi )\alpha ^1`$ and $`s(\gamma (\varphi ))=s(\varphi )\alpha ^1`$ follows immediately from Proposition 7.6, so we need only show $`\gamma `$ is multiplicative on composable elements.
Given $`\varphi ^1(𝒜_1)`$ with minimal partial isometry $`v_\varphi `$ (see Definition 5.1), we first identify the minimal partial isometry for $`\gamma (\varphi )`$. In fact, we claim that
(12)
$$v_{\gamma (\varphi )}=\frac{T\theta (v_\varphi )S}{T\theta (v_\varphi )S}.$$
To see this, first observe that part (4) of Proposition 7.5 also holds for $`x𝒜_2^{\mathrm{\#}\mathrm{\#}}.`$ Thus,
$$(T(\varphi \theta ^1)S)(T\theta (v_\varphi )S)=\varphi (\theta ^1(ST\theta (v_\varphi )ST))=\varphi (v_\varphi )=1.$$
Therefore,
$$\gamma (\varphi )(T\theta (v_\varphi )S)>0.$$
Moreover, if $`q=p_{r(\varphi )}`$ and $`p=p_{s(\varphi )}`$, then by part (3) of Proposition 7.5,
$$T\theta (v_\varphi )S=T\theta (q)\theta (v_\varphi )\theta (p)S=\alpha (p)\theta (v_\varphi )\alpha (q).$$
Hence $`{\displaystyle \frac{T\theta (v_\varphi )S}{T\theta (v_\varphi )S}}`$ is a minimal partial isometry in $`𝒜_2^{\mathrm{\#}\mathrm{\#}}`$ on which $`\gamma (\varphi )`$ takes a positive value. Thus, equation (12) holds by Remark 5.
Now suppose $`\varphi _1,\varphi _2^1(𝒜_1)`$ are such that $`\varphi _1\varphi _2`$ is defined. Notice that the minimal partial isometry for the product $`\varphi _1\varphi _2`$ is the product of the minimal partial isometries for $`\varphi _1`$ and $`\varphi _2`$, that is, $`v_{\varphi _1\varphi _2}=v_{\varphi _1}v_{\varphi _2}`$. To show that $`\gamma (\varphi _1\varphi _2)=\gamma (\varphi _1)\gamma (\varphi _2)`$, it suffices to show that
(13)
$$v_{\gamma (\varphi _1\varphi _2)}=v_{\gamma (\varphi _1)}v_{\gamma (\varphi _2)}.$$
To do this, we first show that for all $`\rho \widehat{𝒟}_1`$, we have
(14)
$$\theta (p_\rho )\alpha (p_\rho )\theta (p_\rho )=\theta (p_\rho ).$$
Indeed, by Proposition 5.3,
$`p_\rho \theta ^1(\alpha (p_\rho ))p_\rho `$ $`=\rho (\theta ^1(\alpha (p_\rho )))p_\rho `$
$`=\rho ((E_1\theta ^1)(\alpha (p_\rho ))p_\rho `$
$`=\rho (\alpha ^1(\alpha (p_\rho )))p_\rho =\rho (p_\rho )p_\rho =p_\rho .`$
Applying $`\theta `$ to the ends of this equality yields (14).
Noting that $`p_{s(\varphi _1)}=p_{r(\varphi _2)}`$, we have
$`(T\theta (v_{\varphi _1})S)(T\theta (v_{\varphi _2})S)`$
$`=[\alpha (p_{r(\varphi _1)})\theta (v_{\varphi _1})\alpha (p_{s(\varphi _1)})][\alpha (p_{r(\varphi _2)})\theta (v_{\varphi _2})\alpha (p_{s(\varphi _2)})]`$
$`=[\alpha (p_{r(\varphi _1)})\theta (v_{\varphi _1})\theta (p_{s(\varphi _1)})\alpha (p_{s(\varphi _1)})][\alpha (p_{r(\varphi _2)})\theta (p_{r(\varphi _2)})\theta (v_{\varphi _2})\alpha (p_{s(\varphi _2)})]`$
$`=\alpha (p_{r(\varphi _1)})\theta (v_{\varphi _1})[\theta (p_{s(\varphi _1)})\alpha (p_{s(\varphi _1)})\alpha (p_{r(\varphi _2)})\theta (p_{r(\varphi _2)})]\theta (v_{\varphi _2})\alpha (p_{s(\varphi _2)})`$
$`=\alpha (p_{r(\varphi _1)})\theta (v_{\varphi _1})\theta (v_{\varphi _2})\alpha (p_{s(\varphi _2)})`$
$`=\alpha (p_{r(\varphi _1)})\theta (v_{\varphi _1\varphi _2})\alpha (p_{s(\varphi _2)})`$
$`=T\theta (v_{\varphi _1\varphi _2})S.`$
This relation, together with equation (12), shows that equation (13) holds, and the proof is complete. ∎
The continuity of the map $`\gamma `$ appearing in Theorem 7.7 is a particularly vexing issue; in general, we do not know whether it is continuous. Theorem 7.7 does imply the restriction of $`\gamma `$ to $`\widehat{𝒟}`$ is continuous, a fact we will use in Example 7.
In the following corollary, we show that in some circumstances, $`\gamma `$ is “nearly continuous”, in the sense that it is possible to alter $`\gamma `$ by multiplying by an appropriate $`𝕋`$-valued cocycle to obtain a continuous isomorphism of coordinate systems. The key hypothesis, that $`\alpha `$ from Definition 7.1 extends to a $``$-isomorphism of $`C^{}`$-envelopes, is in part motivated by Theorem 8.9, which shows that an isometric isomorphism $`\theta `$ between triangular algebras is $``$-extendible to their $`C^{}`$-envelopes, and in particular, $`\alpha =\theta |_{𝒟_1}`$ is $``$-extendible to the envelopes.
###### Corollary 7.8.
If there is a $``$-isomorphism $`\pi :𝒞_1𝒞_2`$ so that $`\pi |_{𝒟_1}=\alpha `$, where $`\alpha =E_2\theta |_{𝒟_1}`$, then the map $`\delta :^1(𝒜_1)^1(𝒜_2)`$, defined by $`\varphi \varphi \pi ^1`$, is an isomorphism of the coordinate systems. Moreover, there exists a cocycle $`c:^1(𝒜_1)𝕋`$ such that for every $`\varphi ^1(𝒜_1)`$,
$$\delta (\varphi )=c(\varphi )\gamma (\varphi ).$$
###### Proof.
Clearly, $`\varphi \varphi \pi ^1`$ is a bicontinuous isomorphism of $`^1(𝒞_1)`$ onto $`^1(𝒞_2)`$. We must show that this map sends $`^1(𝒜_1)`$ into $`^1(𝒜_2)`$. Fix $`\varphi ^1(𝒜_1)`$ and let $`\gamma `$ be the map from Theorem 7.7. By Proposition 7.6,
$$r(\varphi \pi ^1)=r(\varphi )\pi ^1=r(\varphi )\alpha ^1=r(\gamma (\varphi )),$$
and similarly, $`s(\varphi \pi ^1)=s(\gamma (\varphi ))`$. By Corollary 4.11, there is $`c(\varphi )𝕋`$ so that $`\varphi \pi ^1=c(\varphi )\gamma (\varphi )`$, and so $`\varphi \pi ^1^1(𝒜_2)`$. Since both $`\gamma `$ and $`\delta `$ are multiplicative on composable elements of $`^1(𝒜_1)`$, so is $`c`$, whence $`c`$ is a cocycle. ∎
We now introduce a new class of algebras for which $`\gamma `$ is continuous for bounded isomorphisms between triangular algebras in the class. This class includes those algebras $`𝒜(𝒞,𝒟)`$ where $`𝒞`$ admits a cover by monotone $`G`$-sets with respect to $`𝒜`$ (see \[26, p. 57\]). In the context of limit algebras, this class includes limit algebras generated by their order-preserving normalizers (see ).
As the definition of the class does not require our standing assumptions for the section, we relax them momentarily.
###### Definition 7.9.
Let $`(𝒞,𝒟)`$ be a $`C^{}`$-diagonal and $`𝒜(𝒞,𝒟)`$ be a subalgebra (not necessarily triangular). We say a normalizer $`v𝒜`$ is algebra-preserving if either $`v𝒜v^{}𝒜`$ or $`v^{}𝒜v𝒜`$.
This is related to the notion of order-preserving normalizers, which can be described as those $`v𝒩_𝒟(𝒜)`$ satisfying both $`v𝒜v^{}𝒜`$ and $`v^{}𝒜v𝒜`$.
As a trivial example of an algebra-preserving normalizer that is not order preserving, let $`𝒞`$ be $`M_4()`$, $`𝒟`$ the diagonal matrices, and $`𝒜`$ the span of all upper-triangular matrix units except $`e_{1,2}`$. Then $`v=e_{1,3}+e_{2,4}`$ normalizes $`𝒟`$ but is not order preserving, since $`ve_{3,4}v^{}=e_{1,2}𝒜`$. However, $`v`$ is algebra preserving, since $`v^{}𝒜v=v^{}𝒟v𝒜`$.
###### Lemma 7.10.
Let $`(𝒞,𝒟)`$ be a $`C^{}`$-diagonal and let $`𝒜(𝒞,𝒟)`$ be triangular. If $`\varphi ^1(𝒜)`$ and $`v𝒜`$ is an algebra-preserving normalizer with $`\varphi (v)0`$, then for $`x,y𝒜^{\mathrm{\#}\mathrm{\#}}\mathrm{ker}E^{\mathrm{\#}\mathrm{\#}}`$,
$$\varphi (xv)=\varphi (vy)=\varphi (xvy)=0.$$
###### Proof.
By Remark 4 and the fact that $`\mathrm{ker}E^{\mathrm{\#}\mathrm{\#}}𝒜^{\mathrm{\#}\mathrm{\#}}`$ is a $`𝒟`$-bimodule,
$$\varphi (xv)=\frac{r(\varphi )(xvv^{})}{[r(\varphi )(vv^{})]^{1/2}}=0,\varphi (vy)=\frac{s(\varphi )(v^{}vy)}{[s(\varphi )(v^{}v)]^{1/2}}=0,$$
as $`xvv^{},v^{}vy\mathrm{ker}E^{\mathrm{\#}\mathrm{\#}}`$.
For the last equality, assume first that $`v𝒜v^{}𝒜`$. Then $`v𝒜^{\mathrm{\#}\mathrm{\#}}v^{}𝒜^{\mathrm{\#}\mathrm{\#}}`$ and, as $`\mathrm{ker}E^{\mathrm{\#}\mathrm{\#}}𝒜^{\mathrm{\#}\mathrm{\#}}`$ is an ideal in $`𝒜^{\mathrm{\#}\mathrm{\#}}`$ (Theorem 5.10),
$$\varphi (xvy)=\frac{r(\varphi )(x(vyv^{}))}{[r(\varphi )(vv^{})]^{1/2}}=\frac{r(\varphi )(E^{\mathrm{\#}\mathrm{\#}}(x(vyv^{})))}{[r(\varphi )(vv^{})]^{1/2}}=0.$$
If $`v^{}𝒜v𝒜`$, then, similarly, $`\varphi (xvy)=s(\varphi )((v^{}xv)y)/[s(\varphi )(v^{}v)]^{1/2}`$ shows that $`\varphi (xvy)=0`$. ∎
We now reimpose the Standing Assumptions for Section 7; they remain in force through the remainder of the section.
###### Theorem 7.11.
If $`\varphi ^1(𝒜_1)`$ and $`v𝒜_1`$ is an algebra-preserving normalizer such that $`\varphi (v)0`$, then $`\gamma ^1`$ is continuous at $`\gamma (\varphi )`$.
In particular, if $`𝒜_1`$ and $`𝒜_2`$ are the closed span of their algebra-preserving normalizers, then $`\gamma `$ is a homeomorphism.
###### Proof.
Letting $`\rho =s(\varphi )`$, $`\varphi =\lambda [v,\rho ]`$ for some $`\lambda 𝕋`$, by Theorem 4.8. Without loss of generality, we may replace $`v`$ by $`\lambda v`$.
Fix $`\psi ^1(𝒜_1)`$. By Proposition 7.5, $`\theta ^1(S)=I+X`$, $`\theta ^1(T)=I+Y`$ where $`X,Y\mathrm{ker}E_1^{\mathrm{\#}\mathrm{\#}}|_{𝒜_1^{\mathrm{\#}\mathrm{\#}}}`$. Thus,
$$\psi (\theta ^1(S\theta (v)T))=\psi (v+Xv+vY+XvY)=\psi (v),$$
by Lemma 7.10. Putting $`n(\psi )=T(\psi \theta ^1)S^1`$, we can conclude that, for all $`\psi ^1(𝒜_1)`$,
(15)
$$\gamma (\psi )(\theta (v))=n(\psi )\psi (v).$$
Let $`(\varphi _\lambda )`$ be a net in $`^1(𝒜_1)`$ such that $`\gamma (\varphi _\lambda )\gamma (\varphi )`$ and let $`\rho _\lambda =s(\varphi _\lambda )`$. Then $`s(\gamma (\varphi _\lambda ))s(\gamma (\varphi ))`$, and so, by Proposition 7.6, $`\rho _\lambda \rho `$. Since $`\gamma (\varphi _\lambda )(\theta (v))\gamma (\varphi )(\theta (v))=n(\varphi )\varphi (v)>0`$, by (15), we may assume that $`\varphi _\lambda (v)0`$ for all $`\lambda `$. Thus, there exist $`t_\lambda 𝕋`$ such that $`\varphi _\lambda =t_\lambda [v,\rho _\lambda ]`$. Using (15) and the convergence of $`\gamma (\varphi _\lambda )`$,
$`n(\varphi _\lambda )t_\lambda [v,\rho _\lambda ](v)`$ $`=\gamma (\varphi _\lambda )(\theta (v))`$
$`\gamma (\varphi )(\theta (v))=n(\varphi )\varphi (v)=n(\varphi )[v,\rho ](v).`$
Since $`\rho _\lambda \rho `$, $`[v,\rho _\lambda ](v)[v,\rho ](v)0`$, and hence $`t_\lambda n(\varphi _\lambda )n(\varphi )`$. Taking absolute values shows that $`n(\varphi _\lambda )n(\varphi )`$, and hence $`t_\lambda 1`$. Therefore,
$$\varphi _\lambda =t_\lambda [v,\rho _\lambda ][v,\rho ]=\varphi ,$$
as desired. ∎
###### Corollary 7.12.
Suppose, in addition, that each $`𝒜_i`$ is Dirichlet and is the norm closure of the span of its algebra-preserving normalizers. Then $`𝒜_1`$ and $`𝒜_2`$ are boundedly isomorphic if and only if they are isometrically isomorphic.
###### Proof.
If $`𝒜_1`$ and $`𝒜_2`$ are boundedly isomorphic, Theorem 7.11 shows $`^1(𝒜_1)`$ and $`^1(𝒜_2)`$ are isomorphic coordinate systems. As the $`𝒜_i`$ are Dirichlet and regular by hypothesis, the result follows from an application of Theorem 6.11. ∎
This corollary extends \[9, Theorem 2.3\], which proves the corresponding result for strongly maximal TAF algebras generated by their order-preserving normalizers. The cited theorem does somewhat more, as it shows that algebraic isomorphism implies isometric isomorphism.
In light of Theorem 7.7, it is natural to ask what implications can be drawn from the existence of the algebraic isomorphism of coordinate systems. It has been known for more than a decade that algebraic isomorphism of coordinate systems does not imply isometric isomorphism \[18, Remark on page 120\]. It is easily shown that the algebras in their example fail to be boundedly isomorphic and the algebraic isomorphism of coordinates they exhibit is continuous on the diagonals.
We give a somewhat different example, where the algebras are anti-isomorphic, have no minimal projections, and, most importantly, continuity on the diagonal can be exploited to show that the algebras are not boundedly isomorphic.
Example 7.13. Let $`M_k`$, $`D_k`$, $`T_k`$ be the algebra of $`2^k\times 2^k`$ matrices and the subalgebras in $`M_k`$ of diagonal and upper-triangular matrices, respectively. Where necessary, we will equip an $`M_k`$ with a matrix unit system $`\{e_{i,j}\}`$. Let $`A_k`$ (resp., $`B_k`$) be the permutation unitary in $`M_k`$ that interchanges the first (resp., last) two entries of a vector in $`^{2^k}`$. We consider three embeddings from $`M_k`$ to $`M_{k+1}`$. First, we have $`\pi _k`$ that sends a matrix $`[a_{ij}]`$ to the block matrix $`[a_{ij}I_2]`$ where $`I_2`$ is the $`2\times 2`$ identity matrix. Let $`\alpha _k=\mathrm{Ad}A_{k+1}\pi _k`$ and $`\beta _k=\mathrm{Ad}B_{k+1}\pi _k`$. Then $`_a=\underset{}{\mathrm{lim}}(M_k,\alpha _k)`$ and $`_b=\underset{}{\mathrm{lim}}(M_k,\beta _k)`$ are both the $`2^{\mathrm{}}`$ UHF $`C^{}`$-algebra. Since $`\alpha _k|_{D_k}=\beta _k|_{D_k}=\pi _k|_{D_k}`$, we have $`\underset{}{\mathrm{lim}}(D_k,\alpha _k|_{D_k})=\underset{}{\mathrm{lim}}(D_k,\alpha _k|_{D_k})`$, which we denote $`D`$.
The operator algebras $`T_a=\underset{}{\mathrm{lim}}(T_k,\alpha _k|_{T_k})`$ and $`T_b=\underset{}{\mathrm{lim}}(T_k,\beta _k|_{T_k})`$ are anti-isomorphic. Indeed, if $`\varphi _k:T_kT_k`$ sends $`e_{i,j}`$ to $`e_{2^k+1j,2^k+1i}`$, then $`\beta _{k+1}\varphi _k=\varphi _{k+1}\alpha _k`$, so the limit of the $`\varphi _k`$ defines an anti-isomorphism between $`T_a`$ and $`T_b`$.
Suppose that $`T_a`$ and $`T_b`$ were boundedly isomorphic. By Proposition 7.2, we would have a $``$-isomorphism $`\alpha :DD`$ and by Theorem 7.7, there would be an algebraic isomorphism $`\gamma :^1(T_a)^1(T_b)`$. This induces $`\delta :R(T_a)R(T_b)`$, an (algebraic) isomorphism of spectral relations, namely $`\delta (|\varphi |)=|\gamma (\varphi )|`$. Although $`\delta `$ need not be continuous, we know that on the diagonals of $`R(T_a)`$ and $`R(T_b)`$, $`\delta `$ can be identified with $`\widehat{\alpha }:\widehat{D}\widehat{D}`$, the map induced by $`\alpha `$, and so is continuous. Moreover, by the range and source condition in Theorem 7.7, $`(\rho ,\sigma )R(T_a)`$ if and only if $`(\widehat{\alpha }(\rho ),\widehat{\alpha }(\sigma ))R(T_b)`$.
Let $`f`$ (resp., $`l`$) be the element of $`\widehat{D}`$ that equals $`1`$ on the $`e_{1,1}`$ matrix unit (resp., $`e_{2^k,2^k}`$ matrix unit) in each $`D_k`$. Now $`(f,l)`$ is in both $`R(T_a)`$ and $`R(T_b)`$. Let $`𝒪_a=\{\rho \widehat{𝒟}:(f,\rho )R(T_a)`$ and define $`𝒪_b`$ similarly. The existence of $`\delta `$ implies that $`\widehat{\alpha }`$ maps $`𝒪_a`$ onto $`𝒪_b`$. We claim that this is impossible for a continuous $`\alpha `$. The essence of the following argument is that every element of $`𝒪_a\backslash \{f\}`$ has a neighborhood $`N`$ where it is maximal in $`N𝒪_a`$, while no element of $`𝒪_b\backslash \{l\}`$ has such a neighborhood.
The basic neighborhoods for $`f`$ are given by elements of $`\widehat{D}`$ that are nonzero on a $`e_{1,1}`$ matrix unit in some $`D_k`$. If we consider some $`\rho 𝒪_a\backslash \{f\}`$, then there is some $`k`$ and some $`j\{2,\mathrm{},2^k\}`$ so that $`\rho =e_{j,1}fe_{1,j}`$. where $`e_{1,j}`$ is a matrix unit in $`M_k`$. In particular, $`N:=\{\sigma \widehat{D}:\sigma (e_{j,j})0\}`$ is a basic neighborhood of $`\rho `$. Every element of $`N𝒪_a`$ is smaller than $`\rho `$ in the ordering induced by $`R(T_a)`$, since $`N`$ is given by conjugating a basic neighborhood of $`f`$ by $`e_{1,j}`$ and conjugation by $`e_{1,j}`$ reverses the diagonal ordering for pairs $`(f,\psi )`$, $`\psi 𝒪_a\backslash \{f\}`$. This last fact follows from considering the image of $`e_{1,j}`$ in $`M_l`$, $`l>k`$.
On the other hand, if $`\sigma 𝒪_b\backslash \{l\}`$, then every neighborhood of $`\sigma `$ contains elements $`\tau 𝒪_b`$ with $`(\sigma ,\tau )R(T_b)`$. This follows from repeating the argument of the previous paragraph, observing that conjugation by $`e_{1,j}`$ preserves the $`R(T_b)`$-ordering for pairs $`(f,\psi )`$.
Pick some $`\rho 𝒪_a\backslash \{f,l\}`$ and a neighborhood, $`N`$, of $`\rho `$ with all elements of $`N𝒪_a`$ less than $`\rho `$ in the diagonal order. Now $`\widehat{\alpha }`$ maps $`\rho `$ to an element of $`𝒪_b\backslash \{f,l\}`$, call it $`\sigma `$. By the previous paragraph, every neighborhood of $`\sigma `$, including $`\widehat{\alpha }(N)`$, contains points of $`𝒪_b`$ greater than $`\sigma `$ in the $`R(T_b)`$-ordering. This contradicts $`\widehat{\alpha }`$ mapping $`R(T_a)`$ onto $`R(T_b)`$ and so $`T_a`$ and $`T_b`$ are not boundedly isomorphic. In fact, by the automatic continuity result of , there is not even an algebraic isomorphism between $`T_a`$ and $`T_b`$.
We show there is an algebraic isomorphism from $`R(T_a)`$ to $`R(T_b)`$, and hence, using the continuous section, between $`^1(T_a)`$ and $`^1(T_b)`$. It suffices to construct a map $`h:\widehat{D}\widehat{D}`$ so that $`(\sigma ,\tau )R(T_a)`$ if and only if $`(h(\sigma ),h(\tau ))R(T_b)`$.
For $`u\{a,b\},`$ define
$$X_u:=\{(\sigma ,\tau )R(T_u):\sigma =f\}\{(\sigma ,\tau )R(T_u):\tau =l\}.$$
Before defining $`h`$, we observe that $`R(T_a)\backslash X_a=R(T_b)\backslash X_b`$. To see this, first let $`\stackrel{~}{T}_k:=(e_1e_1^{})^{}T_k(e_{2^k}e_{2^k}^{})^{}`$. If $`(\sigma ,\tau )R(T_a)\backslash X_a`$, there is some $`p`$ and some matrix unit $`e\stackrel{~}{T}_p`$ such that $`\sigma =e\tau e^{}`$. For $`kp`$, $`\alpha _k`$ and $`\beta _k`$ agree on the image of $`e`$ in $`\stackrel{~}{T}_k.`$ Thus, the image of $`e`$ in $`T_a`$ and the image of $`e`$ in $`T_b`$ induce the same partial homeomorphism of $`\widehat{D}`$. In particular, $`(\sigma ,\tau )R(T_b)\backslash X_b`$. The reverse inclusion is similar.
Thus, each of $`𝒪_a\backslash \{f,l\}`$ and $`𝒪_b\backslash \{f,l\}`$ is ordered the same way by both $`R(T_a)`$ and $`R(T_b)`$. In fact, we can show that each set is ordered like $``$—for example, given $`(\sigma ,\tau )𝒪_a\backslash \{f,l\}`$ with $`\sigma \tau `$, find, as above, an off-diagonal matrix unit $`e`$ and use the images of $`ee^{}`$ and $`e^{}e`$ in a later matrix algebra to show there is $`\eta \widehat{D}`$ with $`(\sigma ,\eta ),(\eta ,\tau )R(T_a)`$ and $`\eta `$ different from $`\sigma `$ and $`\tau `$.
Define $`h`$ to be the identity map everywhere except $`𝒪_a\backslash \{f,l\}`$ and $`𝒪_b\backslash \{f,l\}`$. Since these two sets have the same order type, there is a bijection $`g:𝒪_a\backslash \{f,l\}𝒪_b\backslash \{f,l\}`$ so that $`(\sigma ,\tau )𝒪_a\backslash \{f,l\}`$ if and only if $`(g(\sigma ),g(\tau ))𝒪_b\backslash \{f,l\}`$. As $`𝒪_a\{f,l\}`$ and $`𝒪_b\{f,l\}`$ are disjoint sets, we define $`h`$ on $`𝒪_a\backslash \{f,l\}`$ to be $`g`$ and on $`𝒪_b\backslash \{f,l\}`$ to be $`g^1`$. Using the observation above and the definitions of $`𝒪_a`$ and $`𝒪_b`$, it is straightforward to show that $`h`$ has the required property.
## 8. Structure of General Isomorphisms
We now turn to the structure of isomorphisms of triangular algebras. First, we build on the representation results from Section 5, obtaining Theorem 8.2, which extends such an isomorphism to an isomorphism of CSL algebras. After several results about isomorphisms of CSL algebras, we obtain an analogue of a factorization result of Arveson-Josephson, Theorem 8.7. A main result of the paper is Theorem 8.8, which shows that such an isomorphism is completely bounded. Finally, we extend a result of Muhly-Qiu-Solel, Theorem 8.9, showing that an isometric isomorphism extends to $``$-isomorphism of the $`C^{}`$-diagonals.
Our standing assumptions are the same as those of the previous section.
We start with a technical lemma.
###### Lemma 8.1.
Suppose that $`C^{}(𝒜_i)=𝒞_i`$. Given $`\rho _2\widehat{𝒟}_2`$, let $`\rho _1=\rho _2\alpha \widehat{𝒟}_1`$, and let $`(\pi _i,_i)`$ be the $`(`$compatible$`)`$ GNS representations of $`(𝒞_i,𝒟_i)`$ on $`_i`$ corresponding to $`\rho _i`$.
If, for $`i=1,2`$, $`P_i`$ is the support projection for $`\pi _i`$, then $`\theta (P_1)=P_2`$.
###### Proof.
By Proposition 5.4, $`P_i𝒟_i^{\mathrm{\#}\mathrm{\#}}`$ and $`P_i=_{q𝒪_i}q`$, where
$$𝒪_i:=\{q𝒟_i^{\mathrm{\#}\mathrm{\#}}:q=p_\sigma \text{ for some }\sigma \widehat{𝒟}_i\text{ such that }(\rho _i,\sigma )R(𝒞_i)\}.$$
Since $`𝒜_i`$ generates $`𝒞_i`$, Theorem 4.22 shows that the equivalence relation generated by $`R(𝒜_i)`$ is $`R(𝒞_i)`$. By Theorem 7.7, we have $`(\sigma ,\rho _2)R(𝒜_2)`$ if and only if $`(\sigma \alpha ,\rho _1)R(𝒜_1)`$. Hence $`(\sigma ,\rho _2)R(𝒞_2)`$ if and only if $`(\sigma \alpha ,\rho _1)R(𝒞_1)`$. Theorem 7.7 also shows that for $`\sigma _1\widehat{𝒟}_1`$, $`\alpha (p_{\sigma _1})=p_{\sigma _1\alpha ^1}.`$ Therefore, $`\alpha (𝒪_1)=𝒪_2`$ and we obtain,
(16)
$$P_2=\underset{p𝒪_1}{}\alpha (p).$$
Fix $`p𝒪_1`$. Since $`p`$ is a minimal projection in $`𝒞_1^{\mathrm{\#}\mathrm{\#}}`$, it is a minimal idempotent in $`𝒜_1^{\mathrm{\#}\mathrm{\#}},`$ so $`\theta (p)`$ is a minimal idempotent in $`𝒜_2^{\mathrm{\#}\mathrm{\#}}`$. As $`P_2`$ is a central projection in $`𝒞_2^{\mathrm{\#}\mathrm{\#}}`$, $`\theta (p)P_2`$ is an idempotent in $`𝒜_2^{\mathrm{\#}\mathrm{\#}}`$ and so $`\theta (p)P_2`$ is either $`\theta (p)`$ or $`0`$. As done earlier, we again use $`\stackrel{~}{\pi }_i`$ for the unique extension of $`\pi _i`$ to $`𝒞_i^{\mathrm{\#}\mathrm{\#}}`$. Since
$$\stackrel{~}{\pi }_2(E_2(\theta (p)P_2))=\stackrel{~}{\pi }_2(\alpha (p))0,$$
we must have $`\theta (p)P_2=\theta (p)`$.
The $`\sigma (𝒜_1^{\mathrm{\#}\mathrm{\#}},𝒜_1^\mathrm{\#})`$-$`\sigma (𝒜_2^{\mathrm{\#}\mathrm{\#}},𝒜_2^\mathrm{\#})`$ continuity of $`\theta `$ yields
(17)
$$\theta (P_1)=\theta \left(\underset{p𝒪_1}{}p\right)=\underset{p𝒪_1}{}\theta (p)=\underset{p𝒪_1}{}P_2\theta (p)=\theta (P_1)P_2.$$
Similar considerations show that for every $`p𝒪_1`$, $`\theta ^1(\alpha (p))P_1=\theta ^1(\alpha (p))`$ and
(18)
$$\theta ^1(P_2)=\underset{p𝒪_1}{}P_1\theta ^1(\alpha (p))=P_1\theta ^1(P_2).$$
Applying $`\theta `$ to (18) and using (17) yields $`\theta (P_1)=P_2`$. ∎
The support projection of a direct sum of inequivalent representations of a $`C^{}`$-algebra is the sum of the support projections of the individual representations. Thus, Lemma 8.1 and Theorem 5.9 combine to produce the following result, which connects our context to the theory of CSL algebras.
When the $`C^{}`$-envelope of $`𝒜_i`$ is $`𝒞_i`$, Theorem 4.22 shows that $`R(𝒞_1)`$ and $`R(𝒞_2)`$ are isomorphic as topological equivalence relations. Thus, the assumption on $`𝒳_2`$ below implies that $`𝒳_2`$ also has exactly one element from each $`R(𝒞_2)`$-equivalence class.
###### Theorem 8.2.
Suppose that $`C^{}(𝒜_i)=𝒞_i`$. Let $`𝒳_2\widehat{𝒟}_2`$ contain exactly one element from each $`R(𝒞_2)`$ equivalence class and let $`𝒳_1=\{\rho \alpha :\rho 𝒳_2\}`$. Let $`\pi _i=_{\rho 𝒳_i}\pi _\rho `$ be the faithful, compatible representations of $`(𝒞_i,𝒟_i)`$ as constructed in Theorem 5.9.
If $`\theta ^{}:\pi _1(𝒜_1)\pi _2(𝒜_2)`$ is the map given by $`\theta ^{}(\pi _1(a))=\pi _2(\theta (a))`$, then $`\theta ^{}`$ extends uniquely to a $`(`$bounded $`)`$ isomorphism $`\overline{\theta }:\stackrel{~}{\pi }_1(𝒜_1^{\mathrm{\#}\mathrm{\#}})\stackrel{~}{\pi }_2(𝒜_2^{\mathrm{\#}\mathrm{\#}})`$.
###### Proof.
Let $`P_i`$ be the support projections of $`\pi _i`$. By Proposition 5.4 and Lemma 8.1, $`P_i𝒟_i^{\mathrm{\#}\mathrm{\#}}`$ and $`\theta (P_1)=P_2`$. By Theorem 5.9, $`\mathrm{ker}\stackrel{~}{\pi }_i|_{𝒜_i^{\mathrm{\#}\mathrm{\#}}}=P_i^{}𝒜_i^{\mathrm{\#}\mathrm{\#}}`$ and so $`\stackrel{~}{\pi }_i`$ is faithful on $`P_i𝒜_i^{\mathrm{\#}\mathrm{\#}}`$. As $`\stackrel{~}{\pi }_i|_{P_i𝒜_i^{\mathrm{\#}\mathrm{\#}}}`$ has image $`\stackrel{~}{\pi }_i(𝒜_i^{\mathrm{\#}\mathrm{\#}})`$, the map $`\overline{\theta }:\stackrel{~}{\pi }_1(𝒜_1^{\mathrm{\#}\mathrm{\#}})\stackrel{~}{\pi }_2(𝒜_2^{\mathrm{\#}\mathrm{\#}})`$ given by $`\stackrel{~}{\pi }_1(a)\stackrel{~}{\pi }_2(\theta (a))`$ is well defined.
Uniqueness follows from the weak-$``$ density of $`𝒜_i`$ in $`𝒜_i^{\mathrm{\#}\mathrm{\#}}`$. ∎
For a representation, $`\pi `$, of $`𝒞_2`$, we suspect that $`\stackrel{~}{\pi }(S)`$ and $`\stackrel{~}{\pi }(T)`$ are inverses of each other whenever $`\pi `$ is a compatible representation. The next two propositions offer some evidence for this. Indeed, Theorem 8.4 proves it when $`\pi _2`$ is the faithful compatible atomic representations of Theorem 8.2.
###### Proposition 8.3.
If $`\pi `$ is a compatible representation of $`𝒞_2`$ on $``$, then $`\stackrel{~}{\pi }(TS)=I`$.
###### Proof.
From Proposition 7.5, we have $`\stackrel{~}{E_2}(\stackrel{~}{\pi }(S))=\stackrel{~}{E_2}(\stackrel{~}{\pi }(T))=I`$. Proposition 7.5 also shows that $`\stackrel{~}{\pi }(TS)\pi (\alpha (𝒟_1))^{}`$, so $`\stackrel{~}{\pi }(TS)\pi (𝒟_2)^{\prime \prime }`$ since $`\pi (𝒟_2)^{\prime \prime }`$ is a MASA in $`\pi (𝒞)^{\prime \prime }`$. Since $`\stackrel{~}{E}_2`$ is a homomorphism on $`\pi (𝒜_2)`$ it is also a homomorphism on $`\stackrel{~}{\pi }(𝒜_2^{\mathrm{\#}\mathrm{\#}})`$. Hence,
$$\stackrel{~}{\pi }(TS)=\stackrel{~}{E}_2(\stackrel{~}{\pi }(TS))=\stackrel{~}{E}_2(\stackrel{~}{\pi }(T))\stackrel{~}{E}_2(\stackrel{~}{\pi }(S))=I.$$
###### Theorem 8.4.
For $`\pi _2`$ as in Theorem 8.2, $`\stackrel{~}{\pi }_2(S)^1=\stackrel{~}{\pi }_2(T)`$.
###### Proof.
We use the same notation as in Lemma 8.1, Theorem 8.2 and their proofs. Fix $`\rho _2\widehat{𝒟}_2`$. We claim that $`\stackrel{~}{\pi }_{\rho _2}(S)`$ is invertible and $`\stackrel{~}{\pi }_{\rho _2}(S)^1=\stackrel{~}{\pi }_{\rho _2}(T)`$.
Applying $`\stackrel{~}{\pi }_{\rho _2}`$ to $`_{p𝒪_1}\theta (p)=P_2`$ (obtained from Lemma 8.1 and the first equality in (17)) yields the important equality,
(19)
$$\underset{p𝒪_1}{}\stackrel{~}{\pi }_{\rho _2}(\theta (p))=I_{_{\rho _2}}.$$
By Proposition 7.5, $`S\alpha (p)=\theta (p)\alpha (p)`$ for all $`p`$, and using (19) gives
$$\stackrel{~}{\pi }_{\rho _2}(S)=\underset{p𝒪_1}{}\stackrel{~}{\pi }_{\rho _2}(S\alpha (p))=\underset{p𝒪_1}{}\stackrel{~}{\pi }_{\rho _2}(\theta (p)\alpha (p))$$
A similar calculation with $`T`$ gives
$$\stackrel{~}{\pi }_{\rho _2}(T)=\underset{p𝒪_1}{}\stackrel{~}{\pi }_{\rho _2}(\alpha (p)\theta (p)).$$
Finally, we then have
$$\stackrel{~}{\pi }_{\rho _2}(ST)=\underset{p𝒪_1}{}\stackrel{~}{\pi }_{\rho _2}(\theta (p))=I_{_{\rho _2}}.$$
Proposition 8.3 established $`\stackrel{~}{\pi }_{\rho _2}(TS)=I_{_{\rho _2}}`$, and hence our claim holds.
As $`\pi _2=_{\rho _2𝒳_2}\pi _{\rho _2}`$, the result follows. ∎
We need two structural results for CSL algebras. The factorization result, Lemma 8.6, is well known, and we only sketch its proof.
###### Theorem 8.5 (\[17, Theorem 2.1\]).
<sup>1</sup><sup>1</sup>1Gilfeather and Moore attribute this result to Ringrose in the nest algebra case and to Hopenwasser for CSL algebras. However, Gilfeather and Moore show that $`\beta `$ is a bounded automorphism.
Suppose $`_1`$ and $`_2`$ are CSLs on Hilbert spaces $`_1`$ and $`_2`$ and that $`\pi :\mathrm{Alg}_1\mathrm{Alg}_2`$ is an algebra isomorphism. Then, given a MASA $`𝔐(_1)`$ which is also contained in $`\mathrm{Alg}_1`$, there exist an invertible operator $`X(_1,_2)`$ and an automorphism $`\beta :\mathrm{Alg}_1\mathrm{Alg}_1`$ such that, for every $`T\mathrm{Alg}_1`$,
$$\pi (T)=X\beta (T)X^1\text{and}\beta |_𝔐=\text{Id}_𝔐.$$
###### Lemma 8.6.
Suppose $`_1`$ and $`_2`$ are atomic CSLs on Hilbert spaces $`_1`$ and $`_2`$ and that $`X\mathrm{Alg}_1X^1=\mathrm{Alg}_2`$ for an invertible operator $`X(_1,_2)`$. There exists a unitary operator $`U(_1,_2)`$ and an invertible operator $`A\mathrm{Alg}_1`$ such that $`A^1\mathrm{Alg}_1`$ and $`X=UA`$.
###### Proof.
Regard $`_i`$ as a commuting family of projections in $`(_i)`$. Let $`𝔸_i`$ be the set of minimal projections in $`_i^{\prime \prime }`$. By hypothesis, $`I=_{a𝔸_i}a`$. Define $`\mu :_1_2`$ by $`\mu (P)=[XP]`$, where $`[XP]`$ denotes the projection onto the range of $`XP`$. Then $`\mu `$ is a complete lattice isomorphism. As each minimal projection in $`𝔸_1`$ has the form $`PQ^{}`$ for some $`P,Q_1`$, we see that $`\mu `$ induces a map, $`\mu ^{}:𝔸_1𝔸_2`$ given by $`\mu ^{}(PQ^{})=\mu (P)\mu (Q)^{}`$. This map is well defined and bijective. Also, for each $`P_1`$,
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$$P=\underset{\{a𝔸_1:aP\}}{}a,\mu (P)=\underset{\{a𝔸_1:aP\}}{}\mu ^{}(a).$$
Since $`_i`$ atomic and the isomorphism between $`\mathrm{Alg}_1`$ and $`\mathrm{Alg}_2`$ is given by an invertible element $`X`$, $`dima_1=dim\mu ^{}(a)_2`$ for every $`a𝔸_1`$. For $`a𝔸_1`$, let $`u_a:_1_2`$ be a partial isometry with $`u_au_a^{}=\mu ^{}(a)`$ and $`u_a^{}u_a=a`$. Put
$$U=\underset{a𝔸_1}{}u_a.$$
Then $`U`$ is unitary and it follows from (20) that $`A:=U^{}X`$ satisfies $`A,A^1\mathrm{Alg}_1`$. ∎
The next two results are structural results for bounded isomorphisms of triangular algebras. The first is an analog of a result of Arveson and Josephson \[2, Theorem 4.10\] appropriate to our setting. Briefly, Arveson and Josephson study a variant of the crossed product algebra associated to a homeomorphism of a locally compact Hausdorff space. If the homeomorphism has no periodic points, then results in show easily that the resulting algebra is a triangular subalgebra of a $`C^{}`$-diagonal (see also \[28, Section 4\]). Arveson and Josephson show that a bounded isomorphism of these algebras factors into three maps, the first an isometric map arising from a homeomorphism of the underlying spaces, the second an isometric map arising from a diagonal unitary, and the third a weakly inner automorphism, i.e., one implemented by an invertible in the ultraweak closure of a suitable representation.
The main difference in the form of the Arveson-Josephson factorization and the factorization in Theorem 8.7 below is that we do not know if the approximately inner part of our factorization carries $`𝒜_1`$ to itself, so we need to introduce an algebra $`𝒜_3`$. We also remark that isomorphisms which fix the diagonal pointwise are essentially cocycle automorphisms (see Definition 6.6).
###### Theorem 8.7.
Assume that $`C^{}(𝒜_i)=𝒞_i`$. Let $`\pi :𝒞_1()`$ be the faithful compatible representation of $`𝒜_1`$ constructed in Theorem 5.9, and let $`\mathrm{Alg}`$ be the weak-$``$ closure $`(`$in $`()`$$`)`$ of $`\pi (𝒜_1)`$. Then $`\theta `$ factors as
$$\theta =\tau \mathrm{Ad}A\beta \pi |_{𝒜_1},$$
where $`\beta \mathrm{Aut}(\mathrm{Alg})`$ with $`\beta (x)=x`$ for $`x\pi (𝒟_1)^{\prime \prime }`$, $`A\mathrm{Alg}`$ with $`A`$ invertible and $`A^1\mathrm{Alg}`$, and, finally, if $`𝒜_3:=(\mathrm{Ad}A\beta )(\pi (𝒜_1)),`$ then $`𝒜_3\mathrm{Alg}`$ and $`\tau :𝒜_3𝒜_2`$ is an isometric isomorphism.
###### Proof.
Apply Theorem 8.2 and Theorem 1 to obtain an invertible operator $`X(_1,_2)`$ which implements a similarity between $`\stackrel{~}{\pi }_1(𝒜_1^{\mathrm{\#}\mathrm{\#}})`$ and $`\stackrel{~}{\pi }_2(𝒜_2^{\mathrm{\#}\mathrm{\#}}).`$ Factor $`X`$ as $`UA`$ where $`U`$ is unitary and $`A,A^1\mathrm{Alg}`$, as in Lemma 8.6. Take $`\tau =\mathrm{Ad}U|_{𝒜_3}`$. The result follows from Theorem 1. ∎
We come now to a main result.
###### Theorem 8.8.
Suppose that $`C^{}(𝒜_i)=𝒞_i`$. If $`\theta :𝒜_1𝒜_2`$ is a bounded isomorphism, then $`\theta `$ is completely bounded and $`\theta _{cb}=\theta `$.
###### Proof.
Using Theorem 8.2 (and its notation), we have $`\overline{\theta }:\stackrel{~}{\pi }_1(𝒜_1^{\mathrm{\#}\mathrm{\#}})\stackrel{~}{\pi }_2(𝒜_2^{\mathrm{\#}\mathrm{\#}})`$. By Theorem 1, $`\overline{\theta }`$ factors as $`\overline{\theta }=\mathrm{Ad}X\beta `$, where $`X:_{\pi _1}_{\pi _2}`$ is a bounded invertible operator and $`\beta `$ is an automorphism of $`\stackrel{~}{\pi }_1(A_1^{\mathrm{\#}\mathrm{\#}})`$ fixing $`\stackrel{~}{\pi }_1(𝒟_1^{\mathrm{\#}\mathrm{\#}})`$ pointwise. By Lemma 8.6, $`X=UA`$ where $`A`$ and $`A^1`$ both belong to $`\stackrel{~}{\pi }_1(𝒜_1^{\mathrm{\#}\mathrm{\#}})`$ and $`U`$ is a unitary operator. Then $`\mathrm{Ad}A\beta `$ is an automorphism of $`\stackrel{~}{\pi }_1(𝒜_1^{\mathrm{\#}\mathrm{\#}})`$ whose norm is $`\theta `$.
By \[7, Corollary 2.5 and Theorem 2.6\], $`\mathrm{Ad}A\beta _{cb}=\mathrm{Ad}A\beta =\theta `$. Thus, $`\overline{\theta }_{cb}=\theta `$. Therefore, for $`\theta ^{}`$ as in Theorem 8.2, $`\theta ^{}`$ is completely bounded. Since $`\theta ^{}_{cb}\overline{\theta }_{cb}`$, we have $`\theta ^{}_{cb}=\theta `$. Noting that each $`\pi _i`$ is a complete isometry of $`𝒜_i`$ onto its respective range completes the proof. ∎
Finally, we use the universal property of $`C^{}`$-envelopes to generalize a result of Muhly, Qiu, and Solel, \[25, Theorem 1.1\]. Their result includes a corresponding statement for anti-isomorphisms, which can be deduced from the statement below by considering appropriate opposite algebras. Our generalization does not require nuclearity of the $`𝒞_i`$ or the second countability of the $`\widehat{𝒟}_i`$, as we do not use the spectral theorem for bimodules, \[24, Theorem 4.1\].
This result also generalizes Corollary 6.11 for isometric $`\theta `$ from triangular subdiagonal algebras to general triangular subalgebras.
###### Theorem 8.9.
For $`i=1,2`$, let $`𝒜_i`$ be a triangular subalgebra of the $`C^{}`$-diagonal $`(𝒞_i,𝒟_i)`$ and assume that $`𝒜_i`$ generates $`𝒞_i`$ as a $`C^{}`$-algebra. If $`\theta :𝒜_1𝒜_2`$ is an isometric isomorphism, then there is a unique $``$-isomorphism $`\pi :𝒞_1𝒞_2`$ with $`\pi |_{𝒜_1}=\theta `$.
###### Proof.
By Proposition 4.22, we know that $`𝒞_i`$ is the $`C^{}`$-envelope of $`𝒜_i`$. Since $`\theta `$ is completely isometric by Theorem 8.8, the universal property for $`C^{}`$-envelopes shows that there exist unique $``$-epimorphisms $`\pi _{12}:𝒞_2𝒞_1`$ and $`\pi _{21}:𝒞_1𝒞_2`$ such that
$$\pi _{12}i_2=i_1\theta ^1\text{and}\pi _{21}i_1=i_2\theta .$$
where, for $`k=1,2`$, $`i_k`$ is the inclusion mapping of $`𝒜_k`$ into $`𝒞_k`$. Thus, $`\pi _{12}\pi _{21}i_1=\text{Id}|_{i_1(𝒜_1)}`$, and hence $`\pi _{12}\pi _{21}=\text{Id}|_{𝒞_1}`$. Thus $`\pi _{21}`$ is injective, and is the required $``$-isomorphism of $`𝒞_1`$ onto $`𝒞_2`$. ∎
Remark 8.10. In Theorems 8.28.78.8 and 8.9, we require that $`𝒞_i`$ is the $`C^{}`$-envelope of $`𝒜_i`$, which is somewhat unsatisfying, as we would prefer conditions in terms of $`𝒜_i`$ alone. The hypothesis that $`C^{}(𝒜_i)=𝒞_i`$ could be removed if we knew that every $`C^{}`$-algebra $`(𝒞,𝒟)`$ is regular, for then $`(C^{}(𝒜_i),𝒟_i)`$ would again be a $`C^{}`$-diagonal.
## 9. Bounded Isomorphism to $``$-Extendible Isomorphism
Roughly speaking, Theorem 8.9 states that an isometric isomorphism of triangular algebras is $``$-extendible. Clearly a bounded, nonisometric isomorphism between triangular algebras cannot be extended to a $``$-isomorphism of the $`C^{}`$-envelopes, but it still may be the case that the $`C^{}`$-envelopes of the triangular algebras are $``$-isomorphic.
Question 9.1. Suppose $`𝒜_i(𝒞_i,𝒟_i)`$ are triangular algebras such that $`C^{}(𝒜_i)=𝒞_i`$. If $`𝒜_1`$ and $`𝒜_2`$ are (boundedly) isomorphic, are $`𝒞_1`$ and $`𝒞_2`$ $``$-isomorphic?
In view of Theorems 6.14 and 7.7, one might expect an affirmative answer when there exists a continuous section of $`R(𝒞_i)`$ into $`^1(𝒞_i)`$, since $`R(𝒜_i)`$ generates $`R(𝒞_i)`$ by Theorem 4.22. However, Theorem 7.7 only implies an algebraic isomorphism of $`R(𝒜_1)`$ onto $`R(𝒜_2)`$; to apply Theorem 6.14, we need to know that the isomorphism of $`R(𝒜_1)`$ onto $`R(𝒜_2)`$ is continuous. Establishing the continuity of the map $`\gamma `$ from Theorem 7.7 would immediately do this.
In this section, we provide an affirmative answer to Question 9 for the class of triangular limit algebras. Perhaps surprisingly, our proof of this result uses $`K`$-theory. We do not know whether the isomorphism obtained satisfies the hypotheses of Corollary 7.8, so we cannot use that corollary to establish the existence of a continuous mapping between the coordinate systems or spectral relations of the triangular limit algebras.
We start with a theorem about Murray-von Neumann equivalence. The proof uses the ideas developed in the previous section. Recall that for any Banach algebra $``$, two idempotents $`e,f`$ are algebraically equivalent if there exist $`x,y`$ such that $`xy=e`$ and $`yx=f`$.
###### Theorem 9.2.
For $`i=1,2`$, suppose $`(𝒞_i,𝒟_i)`$ are $`C^{}`$-diagonals, $`𝒜_i𝒞_i`$ are triangular subalgebras, and $`\theta :𝒜_1𝒜_2`$ is a bounded isomorphism. If $`u𝒜_1`$ is a partial isometry intertwiner, then $`\theta (uu^{})`$ and $`\theta (u^{}u)`$ are algebraically equivalent in $`𝒞_2`$.
To prove the theorem, we need the following well-known result. We give a proof to be self-contained.
###### Lemma 9.3.
Let $`𝒞()`$ be a concrete unital $`C^{}`$-algebra. If $`e`$ is an idempotent in $`𝒞`$, then the projection $`P_e`$ onto the range of $`e`$ is given by $`P_e=(ee^{}+(1e)^{}(1e))^1ee^{}`$ and so is in $`𝒞`$. Moreover, if $`z=IeP_e^{}`$, then $`z`$ is an invertible element of $`𝒞`$ and $`zP_ez^1=e`$.
Suppose $`e`$ and $`f`$ are idempotents in $`𝒞`$ and there exists an element $`x𝒞`$ so that $`xe=x=fx`$ and, as a map from $`e`$ to $`f`$, $`x`$ is invertible. If $`x=v|x|`$ is the polar decomposition for $`x`$, then $`v𝒞`$.
###### Proof.
As $`ee^{}`$ and $`(1e)^{}(1e)`$ commute and their product is 0, they are self-adjoint elements which are supported on orthogonal subspaces of $``$. Since $`ee^{}`$ is bounded below on $`e^{}`$, it is invertible on $`e^{}`$. Similarly, $`(1e)^{}(1e)`$ is invertible on $`(1e)`$. Since the kernel of $`e^{}`$ is the range of $`(1e)`$, $`ee^{}+(1e)^{}(1e)`$ is invertible. By decomposing $``$ into the sum of the ranges of $`e`$ and $`1e^{}`$ and looking at the block matrix decomposition of $`ee^{}+(1e)^{}(1e)`$, we see that $`(ee^{}+(1e)^{}(1e))^1ee^{}`$ is the projection onto the range of $`ee^{}`$, which is the range of $`e`$.
Since $`e`$ and $`P_e`$ have the same range, $`eP_e=P_e`$ and $`P_ee=e`$. A calculation now shows that $`z^1=I+eP_e^{}`$ and $`zP_e=P_e=ez.`$
Turning to $`x`$, let $`P𝒞`$ be the range projection of $`e`$. Then $`|x|`$ is invertible from $`e`$ to $`e`$, so $`|x|+(1P)`$ is an invertible operator on $``$. As $`v^{}v`$ is the projection onto $`e=|x|`$, $`v(1P)=0`$. Thus, $`v=x(|x|+(1P))^1𝒞`$. ∎
###### Proof of Theorem 9.2.
By Theorem 8.2, there exists an isomorphism $`\overline{\theta }`$ between the atomic CSL algebras $`\stackrel{~}{\pi }_1(𝒜_1^{\mathrm{\#}\mathrm{\#}})`$ and $`\stackrel{~}{\pi }_2(𝒜_2^{\mathrm{\#}\mathrm{\#}})`$. Invoking Theorem 1, we can factor $`\overline{\theta }`$ as $`\mathrm{Ad}X\beta `$ where $`\beta `$ is an automorphism of $`\stackrel{~}{\pi }_1(𝒜_1^{\mathrm{\#}\mathrm{\#}})`$ that fixes $`\pi (𝒟_1)^{\prime \prime }`$ pointwise and $`X`$ is invertible.
For ease of notation, identify $`𝒞_i`$ with its image $`\pi _i(𝒞_i)`$; in particular, we write $`u`$ instead of $`\pi _1(u)`$, etc..
Since $`\beta `$ fixes $`𝒟_1^{\prime \prime }`$ and $`u`$ is a partial isometry intertwiner, for every $`d𝒟_1`$, we have $`\beta (u)d=udu^{}\beta (u)`$; and hence $`u^{}\beta (u)𝒟_1^{}=𝒟_1^{\prime \prime }`$. Let $`r=u^{}\beta ^1(u)u^{}.`$ We claim that
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$$r\beta (u)=u^{}u\text{and}\beta (u)r=uu^{}.$$
Indeed,
$$r\beta (u)=u^{}\beta ^1(u)u^{}\beta (u)=u^{}\beta ^1(uu^{}\beta (u))=u^{}\beta ^1(\beta (u))=u^{}u.$$
The other equality is similar. We have $`Xu^{}uX^1=X\beta (u^{}u)X^1=\theta (u^{}u)`$, and similarly, $`Xuu^{}X^1=\theta (uu^{})`$. Thus (21) yields,
$$(XrX^1)\theta (u)=\theta (u^{}u)\text{and}\theta (u)(XrX^1)=\theta (uu^{}),$$
so that $`\theta (u)`$ is invertible as an operator from the range of $`\theta (u^{}u)`$ onto the range of $`\theta (uu^{})`$.
Invoking the second part of Proposition 9.3, $`\theta (u)=v|\theta (u)|`$ with $`v𝒞_2`$. Thus, the range projections of $`\theta (u^{}u)`$ and $`\theta (uu^{})`$ are algebraically equivalent, and hence $`\theta (u^{}u)`$ and $`\theta (uu^{})`$ are algebraically equivalent in $`𝒞_2`$. ∎
Remark 9.4. Uniqueness of inverses shows that actually $`XrX^1𝒞_2`$.
Definition 9.5. For $`n`$, let $`(𝒞_n,𝒟_n)`$ be a $`C^{}`$-diagonal, where $`𝒞_n`$ is a unital finite dimensional $`C^{}`$-algebra, and suppose each $`\alpha _n:𝒞_n𝒞_{n+1}`$ is a regular $``$-monomorphism. Theorem 4.24 shows that the inductive limit, $`(\underset{}{\mathrm{lim}}𝒞_n,\underset{}{\mathrm{lim}}𝒟_n)`$ is a $`C^{}`$-diagonal, which we call an AF-$`C^{}`$-diagonal. The MASA $`\underset{}{\mathrm{lim}}𝒟`$ is often called a canonical MASA. A norm-closed subalgebra $`𝒜(\underset{}{\mathrm{lim}}𝒞_n,\underset{}{\mathrm{lim}}𝒟_n)`$ is called a limit algebra.
We reprise some of the results on limit algebras we require. For a more detailed exposition, see or the introduction to .
Let $`(𝒞,𝒟)=(\underset{}{\mathrm{lim}}𝒞_n,\underset{}{\mathrm{lim}}𝒟_n)`$ be an AF $`C^{}`$-diagonal, and let $`𝒜(𝒞,𝒟)`$ be a limit algebra. For $`vN_𝒟(𝒞)`$, there is some $`i`$ so that we can write $`v=dw`$ where $`d𝒟`$ and $`wN_{𝒟_i}(𝒞_i)`$. We should also point out that, if $`v_{k=1}^{\mathrm{}}N_{𝒟_k}(𝒞_k)`$, the sets $`S(v)`$ and $`S(v^{})`$ are closed and open, so by Propositions 3.3 and 3.2, $`v`$ is also an intertwiner. Also, $`𝒜_k:=𝒞_k𝒜`$ is a finite-dimensional CSL algebra in $`𝒞_k`$, and $`𝒜`$ is the closed union of the $`𝒜_k`$.
The $`C^{}`$-subalgebra $``$ of $`𝒞`$ generated by $`𝒜`$ is again an AF-algebra containing $`𝒟`$, and $`(,𝒟)`$ is again an AF-$`C^{}`$-diagonal. By Proposition 4.22, $``$ is the $`C^{}`$-envelope of $`𝒜`$. Thus by replacing $`𝒞`$ with $``$ if necessary, we may, and shall, always assume that $`𝒜`$ generates $`𝒞`$ as a $`C^{}`$-algebra.
The spectrum, or fundamental relation, of $`𝒜`$, was first defined in , as pairs $`(\sigma ,\rho )\widehat{𝒟}\times \widehat{𝒟}`$ for which there is a partial isometry normalizer $`v𝒜`$ with $`\sigma =v\rho v^{}`$. In our notation, this is $`R(𝒜)`$. The spectrum can also be described by picking systems of matrix units for each $`𝒞_n`$ so that matrix units in $`𝒞_n`$ are sums of matrix units in $`𝒞_{n+1}`$ and then considering those elements of $`𝒜^\mathrm{\#}`$ that are either $`0`$ or $`1`$ on all matrix units. These elements of $`𝒜^\mathrm{\#}`$ are eigenfunctionals and this description provides a continuous section from $`(𝒜)`$ to $`^1(𝒜)`$.
We require a technical result on normalizing idempotents in a triangular subalgebra of a finite dimensional $`C^{}`$-diagonal. The method is similar to that of Proposition 7.5, and is in fact what led to the constructions of $`S`$ and $`T`$.
###### Lemma 9.6.
Suppose that $`(𝒞,𝒟)`$ is a $`C^{}`$-diagonal with $`𝒞`$ finite dimensional, and $`𝒜(𝒞,𝒟)`$ is triangular. Let $`𝔅𝒜`$ be a (necessarily finite) Boolean algebra of commuting idempotents. Then there exists an invertible element $`A𝒜`$ such that $`A𝔅A^1𝒟`$ is a Boolean algebra of idempotents.
###### Proof.
Let $`G=\{I2e:e𝔅\}`$; then $`G`$ is a finite group whose elements are all square roots of the identity. Clearly $`G`$ is in bijective correspondence with $`𝔅`$. Define
$$S=\frac{1}{|G|}\underset{gG}{}E(g)g^1.$$
A calculation shows that for any $`gG`$, $`E(g)S=Sg`$, and we have $`E(S)=I`$. Thus, $`S=I+Y`$ where $`Y𝒜`$ is nilpotent, and we conclude that $`S`$ is invertible. Then for every $`e𝔅`$, $`SeS^1=E(e)`$, and we are done. ∎
###### Corollary 9.7.
Suppose $`(𝒞,𝒟)`$ is an AF-$`C^{}`$-diagonal and $`𝒜(𝒞,𝒟)`$ is a triangular subalgebra. If $`e𝒜`$ is an idempotent, then there exist $`A𝒜`$ such that $`AeA^1=E(e).`$
###### Proof.
Write $`(𝒞,𝒟)=\underset{}{\mathrm{lim}}(𝒞_n,𝒟_n)`$ where $`(𝒞_n,𝒟_n)`$ are finite dimensional $`C^{}`$-diagonals, and let $`𝒜_n=𝒞_n𝒜`$, so that $`𝒜=\underset{}{\mathrm{lim}}𝒜_n`$. By \[3, Proposition 4.5.1\], there exists $`n`$, an idempotent $`f𝒜_n`$ and an invertible element $`X𝒜`$ so that $`XeX^1=f`$. Lemma 9.6 (applied to $`\{0,e,Ie,I\}`$) shows that there exists $`S𝒜_n`$ so that $`SfS^1=E(f)`$. Thus $`(SX)e(SX)^1=E(f)`$. Since $`E|_𝒜`$ is a homomorphism, applying $`E`$ to the previous equality yields $`E(e)=E(f)`$, and the proof is complete. ∎
Remark 9.8. Let $`\iota :𝒟𝒜`$ be the inclusion map of $`𝒟`$ into the triangular limit algebra $`𝒜`$. As in , Corollary 9.7 implies $`\iota _{}:K_0(𝒟)K_0(𝒜)`$ is an isomorphism of scaled dimension groups and $`\iota _{}^1=E_{}.`$
We now show that an isomorphism of triangular limit algebras implies the existence of a $``$-isomorphism of the $`C^{}`$-envelopes.
###### Theorem 9.9.
Suppose $`\theta :𝒜_1𝒜_2`$ is an algebra isomorphism of the triangular limit algebras $`𝒜_i`$. For $`i=1,2`$, let $`𝒞_i`$ be the $`C^{}`$-envelope of $`𝒜_i`$ and let $`h_i:𝒟_i𝒜_i`$ and $`k_i:𝒜_i𝒞_i`$ be the inclusion maps. Then there exists a $``$-isomorphism $`\tau :𝒞_1𝒞_2`$ such that the following diagram of scaled dimension groups commutes.
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$$\begin{array}{ccccc}K_0(𝒟_1)& \stackrel{h_1}{}& K_0(𝒜_1)& \stackrel{k_1}{}& K_0(𝒞_1)\\ \alpha _{}& & \theta _{}& & \tau _{}& & \\ K_0(𝒟_2)& \stackrel{h_2}{}& K_0(𝒜_2)& \stackrel{k_2}{}& K_0(𝒞_2)\end{array}$$
###### Proof.
Recall from that algebraic isomorphisms of limit algebras are necessarily bounded.
That $`\theta _{}h_1=h_2\alpha _{}`$ follows from the fact that $`\alpha =E_2\theta |_{𝒟_1}`$ and Remark 9. For $`j=1,2`$, let $`\iota _j=k_jh_j`$ be the inclusion mapping of $`𝒟_j`$ into $`𝒞_j`$. To complete the proof, we shall show the existence of $`\tau _{}:K_0(𝒞_1)K_0(𝒞_2)`$ so that $`\tau _{}\iota _1=\iota _2\alpha _{}`$.
Write $`𝒞_j=\underset{}{\mathrm{lim}}_k(𝒞_{jk},𝒟_{jk})`$ where $`(𝒞_{jk},𝒟_{jk})`$ are finite dimensional $`C^{}`$-inclusions. Without loss of generality, we may assume that $`𝒜_{jk}:=𝒜_j𝒞_{jk}`$ satisfies $`C^{}(𝒜_{jk})=𝒞_{jk}`$. Any projection $`p𝒞_1`$ is algebraically equivalent to a projection in $`𝒞_{1k}`$ for some $`k`$, so $`p`$ is algebraically equivalent to a projection $`p^{}𝒟_{1k}`$. It follows that the induced mapping of scaled dimension groups, $`(\iota _j)_{}:K_0(𝒟_j)K_0(𝒞_j)`$ is onto.
We claim that if $`p`$ and $`q`$ are projections in $`𝒟_1`$ which are algebraically equivalent in $`𝒞_1,`$ then $`\alpha (p)`$ and $`\alpha (q)`$ are algebraically equivalent in $`𝒞_2`$, and we modify ideas of \[36, Lemma 2.2\] for this. We may assume $`p,q𝒟_{1k}`$ for some $`k`$, and are algebraically equivalent in $`𝒞_{1k}`$. In fact, we shall show that they are equivalent via an element of $`𝒩_{𝒟_{1k}}(𝒞_{1k}).`$
Since $`p`$ and $`q`$ are algebraically equivalent, they have the same center-valued trace, and hence there exists a positive integer $`r`$ and minimal projections $`p_i`$, $`q_i`$ belonging to $`𝒟_{1k}`$ so that $`p=p_1+\mathrm{}+p_r`$, $`q=q_1+\mathrm{}+q_r`$. By relabeling if necessary, we may assume that for each $`i`$ with $`1ir`$, $`p_i`$ and $`q_i`$ are algebraically equivalent. Let $`w_i𝒞_{1k}`$ be a partial isometry so that $`q_iw_ip_i=w_i`$, $`w_i^{}w_i=p_i`$ and $`w_iw_i^{}=q_i`$. Since $`p_i`$ and $`q_i`$ are minimal projections in $`𝒞_{1k}`$, $`w_i`$ are minimal partial isometries in $`𝒞_{1k}`$. Moreover, $`w=_{i=1}^rw_i𝒩_{𝒟_{1k}}(𝒞_{1k})`$ satisfies $`w^{}w=p`$ and $`ww^{}=q`$.
Since $`C^{}(𝒜_{1k})=𝒞_{1k}`$, we can write $`w_i=v_{i_1}v_{i_2}\mathrm{}v_{i_l}`$ as a finite product of partial isometries, with each $`v_{i_j}`$ a normalizer of $`𝒟_{1k}`$ and belonging to either $`𝒜_{1k}`$ or $`𝒜_{1k}^{}`$. Theorem 9.2 shows that $`\alpha (v_{i_j}^{}v_{i_j})`$ and $`\alpha (v_{i_j}v_{i_j}^{})`$ are algebraically equivalent in $`𝒞_2.`$ Thus, $`\alpha (p_i)`$ and $`\alpha (q_i)`$ are equivalent in $`𝒞_2`$ also, whence $`\alpha (p)`$ and $`\alpha (q)`$ are algebraically equivalent in $`𝒞_2`$ as desired.
Thus, if $`p𝒞_1`$ is a projection, we may define $`\tau _{}([p])=(\iota _2\alpha _{})([p^{}])`$, where $`p^{}𝒟_1`$ is any projection in $`𝒟_1`$ with $`\iota _1([p^{}])=[p].`$ The previous paragraph shows $`\tau _{}`$ is well-defined, and it determines an isomorphism of the scaled dimension groups $`K_0(𝒞_1)`$ and $`K_0(𝒞_2)`$ satisfying (22).
An application of Elliott’s Theorem now completes the proof. ∎
We would very much like to know whether it is possible to choose $`\tau `$ in the conclusion of Theorem 9.9 so that $`\tau |_{𝒟_1}=\alpha `$. When this is the case, Corollary 7.8 implies the existence of a continuous isomorphism of coordinate systems, and hence spectra. The next example shows that more than the $`K`$-theoretic data provided by the conclusion of Theorem 9.9 is required to prove the existence of such a $``$-isomorphism.
Example 9.10. Suppose, for $`j=1,2`$, that $`(𝒞_j,𝒟_j)`$, are AF $`C^{}`$-diagonals and $`i_j:𝒟_j𝒞_j`$ are the natural inclusions. Given an isomorphism $`\alpha :𝒟_1𝒟_2`$ with an isomorphism of scaled dimension groups, $`h:K_0(𝒞_1)K_0(𝒞_2)`$ with $`i_2\alpha _{}=hi_1`$, there need not exist a $``$-isomorphism $`\tau :𝒞_1𝒞_2`$ with $`\tau |_{𝒟_1}=\alpha `$.
To see this, we use two well-known direct systems for triangular AF algebras, the refinement system and refinement with twist system. Define $`\rho _k:M_{2^k}M_{2^{k+1}}`$ by sending a matrix $`A=[a_{ij}]`$ to $`\rho _k(A)=[a_{ij}I_2]`$, i.e., replacing each entry of $`A`$ with the corresponding multiple of a $`2\times 2`$ identity matrix. Define $`\varphi _k:M_{2^k}M_{2^{k+1}}`$ to be $`\mathrm{Ad}U_k\rho _k`$, where $`U_k`$ is the $`2^{k+1}\times 2^{k+1}`$ permutation unitary which is the direct sum of a $`2^{k+1}2`$ identity matrix and $`\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right]`$.
Let $`𝒞_1=\underset{}{\mathrm{lim}}(M_{2^k},\rho _k)`$ and $`𝒞_2=\underset{}{\mathrm{lim}}(M_{2^k},\varphi _k)`$. Let $`𝒟_1`$ and $`𝒟_2`$ be the direct limits of the diagonal matrices in each direct system. Since $`\rho _k`$ and $`\varphi _k`$ agree on $`D_{2^k}`$, the direct limit of the identity maps $`\mathrm{id}:D_{2^k}D_{2^k}`$ defines an isomorphism, $`\alpha `$, from $`𝒟_1`$ to $`𝒟_2`$. Now, $`𝒞_1`$ and $`𝒞_2`$ are isomorphic, as they are UHF $`C^{}`$-algebras with the same ‘supernatural’ number, $`2^{\mathrm{}}`$. Further, $`K_0(𝒞_j)`$ can be identified with $`G=\{k/2^n:k,n\}`$, with the usual order and scale $`G[0,1]`$. With this identification, $`i_j`$ is the usual trace from $`K_0(𝒟_j)`$ into $`G`$. Since $`\alpha `$ is the identity on $`K_0(𝒟_j)`$, we have $`i_2\alpha _{}=i_1`$. Thus, we can take $`h`$ to be the identity map on $`G`$.
It remains to show that there is no $``$-isomorphism $`\tau :𝒞_1𝒞_2`$ with $`\tau |_{𝒟_1}=\alpha `$. We argue by contradiction, so assume such a $`\tau `$ exists.
We may build an intertwining diagram as follows. For brevity, let $`C_i`$ be $`M_{2^i}`$, $`\rho _{i,j}`$ denote $`\rho _i\rho _{i+1}\mathrm{}\rho _{j1}`$, and define $`\varphi _{i,j}`$ similarly. By \[5, Theorem 2.7\], there are sequences $`(m_i)`$ and $`(n_i)`$ and $``$-monomorphisms $`(\psi _i)`$ and $`(\eta _i)`$ so that the following diagram commutes:
$$\begin{array}{cccccccccc}C_1& \stackrel{\rho _{1,m_1}}{}& C_{m_1}& \stackrel{\rho _{m_1,m_2}}{}& C_{m_2}& \stackrel{\rho _{m_2,m_3}}{}& C_{m_3}& & \mathrm{}& 𝒞_1\\ & \psi _1& \eta _1& \psi _2& \eta _2& \psi _3& \eta _3& \psi _4& & \tau \\ C_1& \stackrel{\varphi _{1,n_1}}{}& C_{n_1}& \stackrel{\varphi _{n_1,n_2}}{}& C_{n_2}& \stackrel{\varphi _{n_2,n_3}}{}& C_{n_3}& & \mathrm{}& 𝒞_2\end{array}$$
Since $`\tau `$ maps $`𝒟_1`$ onto $`𝒟_2`$, we can use this diagram to show that each $`\psi _k`$ and $`\eta _k`$ are restrictions of $`\alpha `$ and $`\alpha ^1`$, respectively, and so are the identity map at the level of matrix algebras.
To obtain the contradiction, first fix $`C_k=M_{2^k}`$ and observe that if $`e`$ is the $`(1,1)`$ matrix unit and $`f`$ the $`(2^k,2^k)`$ matrix unit in $`C_k`$, then for any $`l<k`$, $`e\rho _{l,k}(C_l)f=0`$ while $`e\varphi _{l,k}(C_l)f0`$. Now consider the two maps $`\lambda =\varphi _{n_1,n_3}`$ and $`\mu =\psi _3\rho _{m_1,m_2}\eta _1`$ from $`C_{n_1}`$ into $`C_{n_3}`$. Letting $`e`$ and $`f`$ be the $`(1,1)`$ and $`(2^{n_3},2^{n_3})`$ matrix units in $`C_{n_3}`$, the observation implies that $`e\lambda (C_{n_1})f0`$. To see that $`e\mu (C_{n_1})f=0`$, let $`e^{}`$ and $`f^{}`$ be the $`(1,1)`$ and $`(2^{m_2},2^{m_2})`$ matrix units in $`C_{m_2}`$ and observe that $`e^{}\rho _{m_1,m_2}(\eta _1(C_{n_1}))f^{}=0`$ by the observation. Applying $`\varphi _3`$ and noting that $`e`$,$`f`$ are subprojections of $`\varphi _3(e^{})`$,$`\varphi _3(f^{})`$ respectively completes the argument.
Thus, no such diagram exists, and hence no such $`\tau `$ exists.
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# A directed walk model of a long chain polymer in a slit with attractive walls
## 1 Introduction
When a long linear polymer molecule in dilute solution is confined between two parallel plates the polymer loses configurational entropy and exerts a repulsive force on the confining plates. If the monomers are attracted to one of the two confining plates, the force is still repulsive, but less is magnitude because the polymer can adsorb on one plate and therefore does not extend as far into solution. The effect of the second confining plate is then less and the loss of configurational entropy is less. If the monomers are attracted to both plates the net force can be attractive at large distances and repulsive at smaller distances. At large distances the energy term is dominant while at smaller distances the entropy loss is dominant. These phenomena are related to the stabilization of colloidal dispersions by adsorbed polymers (steric stabilization) and the destabilization when the polymer can adsorb on surfaces of different colloidal particles (sensitized flocculation).
One would hope to be able to say something about this problem for a relatively realistic model of a polymer in a good solvent such as a self-avoiding walk and the self-avoiding walk case has been considered by several authors (see for instance Wall *et al* 1977 and 1978, Hammersley and Whittington 1985). When the self-avoiding walk is simply confined between two parallel lines or planes, and does not otherwise interact with the confining lines or planes, a number of results are available. Let $`c_n(w)`$ be the number of self-avoiding walks on the simple cubic lattice $`^3`$, starting at the origin and confined to have the $`z`$-coordinate of each vertex in the slab $`0zw`$. Then it is known that the limit
$$\underset{n\mathrm{}}{lim}n^1\mathrm{log}c_n(w)\kappa (w)$$
(1.1)
exists. If $`c_n`$ is the number of $`n`$-edge self-avoiding walks with no geometric constraint then the connective constant of the lattice is given by
$$\underset{n\mathrm{}}{lim}n^1\mathrm{log}c_n\kappa .$$
(1.2)
It is known that
1. $`\kappa (w)`$ is monotone increasing in $`w`$,
2. $`lim_w\mathrm{}\kappa (w)=\kappa `$, and
3. $`\kappa (w)`$ is a concave function of $`w`$.
In two dimensions, the exact values of $`\kappa (w)`$ are known for small $`w`$ (Wall *et al* 1977). When the walk interacts with the confining lines or planes almost nothing is known rigorously but the problem has been studied by exact enumeration methods (Middlemiss *et al* 1976).
To make progress with the situation in which the walk interacts with the confining lines or planes one has to turn to simpler models. In a classic paper DiMarzio and Rubin (1971) considered a random walk model on a regular lattice where the random walk is confined between two parallel lines or planes, and interacts with one or both of the confining surfaces. They see both attractive and repulsive regimes in their model, though they confine themselves to the case where either the polymer interacts with only one surface or equally with both surfaces.
In this paper we study a directed version of the self-avoiding walk model. We consider directed self-avoiding walks on the square lattice, confined between two lines ($`y=0`$ and $`y=w`$). We consider the cases where the walk starts and ends in $`y=0`$ (confined Dyck paths, or *loops*), starts in $`y=0`$ and ends in $`y=w`$ (*bridges*) and starts in $`y=0`$ but has no condition on the other endpoint beyond the geometrical constraint (*tails*). Our model is related to that of DiMarzio and Rubin (1971) but we consider the more general situation where the interaction with the two confining lines or planes can be different. The techniques which we use are quite different from those of DiMarzio and Rubin.
We derive results for the behaviour of these three classes of directed walks as a function of the interaction of the vertices with the two confining lines, and the width $`w`$. In particular, we find the generating functions exactly for each of the three models. We find the free energy and hence the phase diagram in the limit of large slit width for polymers much larger than the slit width. We demonstrate that this phase diagram is *not* the one obtained by considering the limit of large slit width for any finite polymer. We calculate the asymptotics for the free energy in wide slits and so deduce the forces between the walls for large slit widths. This allows us to give a “force diagram” showing the regions in which the attractive, repulsive, short-ranged and long-ranged forces act.
## 2 The model
The directed walk model that we consider is closely related to Dyck paths (see Stanton and White 1986 or Deutsch 1999). These classical objects are closely related to Ballot paths which were studied as a voting problem (Andre 1887 and Bertrand 1887). Dyck paths are directed walks on $`^2`$ starting at $`(0,0)`$ and ending on the line $`y=0`$, which have no vertices with negative $`y`$-coordinates, and which have steps in the $`(1,1)`$ and $`(1,1)`$ directions. We impose the additional geometrical constraint that the paths lie in the slit of width $`w`$ defined by the lines $`y=0`$ and $`y=w`$. We refer to Dyck paths that satisfy this slit constraint as *loops* (see figure 1). We also consider two similar sets of paths; *bridges* (see figure 2) and *tails* (see figure 3). A bridge is a directed path lying in the slit which has its zeroth vertex at $`(0,0)`$ and last vertex in $`y=w`$. A tail is a directed path in the slit whose zeroth vertex is at $`(0,0)`$ and whose last vertex may lie at any $`0yw`$.
Let $`_w,_w`$ and $`𝒯_w`$ be the sets of loops, bridges and tails, respectively, in the slit of width $`w`$. We define the generating functions of these paths as follows:
$`L_w(z,a,b)`$ $`=`$ $`{\displaystyle \underset{p_w}{}}z^{n(p)}a^{u(p)}b^{v(p)};`$ (2.1)
$`B_w(z,a,b)`$ $`=`$ $`{\displaystyle \underset{p_w}{}}z^{n(p)}a^{u(p)}b^{v(p)};`$ (2.2)
$`T_w(z,a,b)`$ $`=`$ $`{\displaystyle \underset{p𝒯_w}{}}z^{n(p)}a^{u(p)}b^{v(p)},`$ (2.3)
where $`n(p),u(p)`$ and $`v(p)`$ are the number of edges in the path $`p`$, the number of vertices in the line $`y=0`$ (excluding the zeroth vertex) and the number of vertices in the line $`y=w`$, respectively.
Let $`_w^n,_w^n`$ and $`𝒯_w^n`$ be the sets of loops, bridges and tails of fixed length $`n`$ in the slit of width $`w`$. The partition function of loops is defined as
$$Z_n^{loop}(w;a,b)=\underset{p_w^n}{}a^{u(p)}b^{v(p)}$$
(2.4)
with the partition functions for bridges and tails defined analogously. Hence the generating functions are related to the partition functions in the standard way, eg for loops we have
$$L_w(z,a,b)=\underset{n}{}z^nZ_n^{loop}(w;a,b).$$
(2.5)
We define the reduced free energy $`\kappa ^{loop}(w;a,b)`$ for loops for fixed finite $`w`$ as
$$\kappa ^{loop}(w;a,b)=\underset{n\mathrm{}}{lim}n^1\mathrm{log}Z_n^{loop}(w;a,b)$$
(2.6)
with the reduced free energies for bridges $`\kappa ^{bridge}(w;a,b)`$ and tails $`\kappa ^{tail}(w;a,b)`$ defined analogously.
Consider the singularities $`z_c^{loop}(w;a,b)`$, $`z_c^{bridge}(w;a,b)`$ and $`z_c^{tail}(w;a,b)`$ of the generating functions $`L_w(z,a,b)`$, $`B_w(z,a,b)`$ and $`T_w(z,a,b)`$, respectively, closest to the origin and on the positive real axis, known as the critical points. Given that the radii of convergence of the generating functions are finite, which we shall demonstrate, and since the partition functions are positive, the critical points exist and are equal in value to the radii of convergence. Hence the free energies exist and one can relate the critical points to the free energies for each $`type`$ of configuration, (that is loops, bridges and tails) as
$`\kappa ^{type}(w;a,b)=\mathrm{log}z_c^{type}(w;a,b).`$ (2.7)
###### Theorem 2.1.
Loops, bridges and tails have the same limiting free energy at every fixed finite $`w`$.
###### Proof.
We sketch the proof: Fix $`w`$ at some finite positive integer, and let $`Z_n^{loop},Z_n^{bridge}`$ and $`Z_n^{tail}`$ be the partition functions of loops, bridges and tails with $`n`$ edges in the strip of width $`w`$. By appending $`w`$ edges to loops and bridges we have the inequalities
$$bZ_n^{loop}Z_{n+w}^{bridge}\text{ and }aZ_n^{bridge}Z_{n+w}^{loop}.$$
(2.8)
Hence $`abZ_n^{loop}aZ_{n+w}^{bridge}Z_{n+2w}^{loop}`$. Taking logarithms, dividing by $`n`$ and taking the limit $`n\mathrm{}`$ shows that the limiting free energies are equal (since the limits exist).
Since every loop is a tail $`Z_n^{loop}Z_n^{tail}`$, and by appending $`w`$ (if $`n`$ and $`w`$ have the same parity) or $`w1`$ edges (if $`n`$ and $`w`$ have different parity) to tails we have
$$aZ_n^{tail}Z_{n+w}^{loop}\text{ or }aZ_n^{tail}Z_{n+w1}^{loop},$$
(2.9)
depending on the parity of $`n+w`$. By a similar sandwiching argument the limiting free energies of loops and tails, and hence bridges, are all equal. ∎
Since the free energies are the same for the three models we define $`\kappa (w;a,b)`$ as
$$\kappa (w;a,b)\kappa ^{loop}(w;a,b)=\kappa ^{bridge}(w;a,b)=\kappa ^{tail}(w;a,b)$$
(2.10)
with $`z_c(w;a,b)`$ being
$$z_c(w;a,b)z_c^{loop}(w;a,b)=z_c^{bridge}(w;a,b)=z_c^{tail}(w;a,b)$$
(2.11)
so that
$$\kappa (w;a,b)=\mathrm{log}z_c(w;a,b).$$
(2.12)
## 3 The half-plane
For loops and tails we first consider the case in which $`w\mathrm{}`$, which reduces the problem to the adsorption of paths to a wall in the half-plane. We begin by noting that once $`w>n`$ for any finite length walk there can no longer be any visits to the top surface so
$$_{w+1}^n=_w^n\text{ and }𝒯_{w+1}^n=𝒯_w^n\text{ for all }w>n.$$
(3.1)
Hence we define the sets $`_{hp}^n`$ and $`𝒯_{hp}^n`$ as
$$_{hp}^n=_{n+1}^n\text{ and }𝒯_{hp}^n=𝒯_{n+1}^n.$$
(3.2)
The limit $`w\mathrm{}`$ can be therefore be taken explicitly. Also, as a consequence of the above, for all $`w>n`$ we have $`v(p)=0`$ for any $`p_w^n`$ and also for any $`p𝒯_w^n`$. Hence the partition function of loops in the half-plane can be defined as
$$Z_n^{loop,hp}(a)=\underset{p_{hp}^n}{}a^{u(p)}$$
(3.3)
with the partition function of tails in the half-plane defined analogously. We define the generating function of loops in the half plane via the partition function as
$$L(z,a)=\underset{n}{}z^nZ_n^{loop,hp}(a)$$
(3.4)
with the generating function of tails $`T(z,a)`$ in the half-plane defined analogously. Note that the above limit does not exist for bridges since for any fixed walk of length $`n`$ we have $`wn`$. In an analogous way to the slit we define the reduced free energy in the half-plane for loops (and tails analogously) as
$$\kappa ^{loop,hp}(w;a)=\underset{n\mathrm{}}{lim}n^1\mathrm{log}Z_n^{loop,hp}(w;a).$$
(3.5)
We note that in defining these free energies for the half-plane the thermodynamic limit $`n\mathrm{}`$ is taken *after* the limit $`w\mathrm{}`$: we shall return to this order of limits later.
While the half-plane solution is well known (Brak *et al* 1998, 2001 or Janse van Rensburg 2000) it is useful to summarise the main results for comparison with our new results for slits. One may factor both loops and tails in the half plane by considering the second last point at which a loop touches the line $`y=0`$ or the last point at which a tail touches the surface. This leads to the functional equations $`L(z,a)=1+az^2L(z,a)L(z,1)`$ and $`T(z,a)=L(z,a)+zL(z,a)T(z,1)`$ for the generating functions for loops and tails. These may be easily solved to give
$`L(z,a)`$ $`=`$ $`{\displaystyle \frac{2}{2a+a\sqrt{14z^2}}}\text{ and}`$ (3.6)
$`T(z,a)`$ $`=`$ $`{\displaystyle \frac{4z}{(a2a\sqrt{14z^2})(12z\sqrt{14z^2})}}.`$ (3.7)
The generating functions are singular when the square root term is zero and when the denominator is zero. These conditions give the locations of the critical points, $`z_c^{loop,hp}(a)`$ and $`z_c^{tail,hp}(a)`$ and so also give the free energies. The values are equal for the two models and so we define $`z_c^{hp}(a)z_c^{loop,hp}(a)=z_c^{tail,hp}(a)`$, and similarly for the free energy we have
$$\kappa ^{hp}(a)\kappa ^{loop,hp}(a)=\kappa ^{tail,hp}(a)=\mathrm{log}z_c^{hp}(a)$$
(3.8)
with
$$\kappa ^{hp}(a)=\{\begin{array}{cc}\mathrm{log}(2)& a2\\ \mathrm{log}\left(\frac{a}{\sqrt{a1}}\right)& a>2.\end{array}$$
(3.9)
The second derivative of $`\kappa ^{hp}(a)`$ is discontinuous at $`a=2`$ and so there is a second-order phase transition at this point. The density of visits in the thermodynamic limit, defined as
$$\rho ^{hp}(a)=\underset{n\mathrm{}}{lim}\frac{v}{n}=\frac{d\kappa ^{hp}(a)}{d\mathrm{log}a}$$
(3.10)
is an order parameter for the transition and has corresponding exponent $`\beta =1`$.
## 4 Exact solution for the generating functions in finite width strips
The solution of the strip problem proceeds in an altogether different manner to the half-plane situation. This begins by using an argument that builds up configurations (uniquely) in a strip of width $`w+1`$ from configurations in a strip of width $`w`$. In this way a recurrence-functional equation is constructed rather than a simple functional equation. Consider configurations of any type — loops, bridges or tails — in a strip of width $`w`$ (see figure 4), and focus on the vertices touching the top wall: call these *top* vertices. These vertices contribute a factor $`b`$ to the Boltzmann weight of the configuration. Now consider a zig-zag path (see figure 4), which is defined as a path of any even length, or one of length zero, in a strip of width 1. The generating function of zig-zag paths is $`1/(1bz^2)`$. Replace each of the *top* vertices in the configuration by any zig-zag path. Since one could choose a single vertex as the zig-zag path all configurations that fit in a strip of width $`w`$ are reproduced. Also, the addition of any non-zero length path at any top vertex will result in a new configuration of width $`w+1`$ and no more. The inverse process is also well defined and so we can write recurrence-functional equations for each of the generating functions.
The generating function $`L_w(z,a,b)`$ satisfies the following functional recurrence:
$`L_1(z,a,b)`$ $`=`$ $`{\displaystyle \frac{1}{1abz^2}};`$ (4.1)
$`L_{w+1}(z,a,b)`$ $`=`$ $`L_w(z,a,{\displaystyle \frac{1}{1bz^2}}).`$ (4.2)
The generating function of bridges, $`B_w(z,a,b)`$, satisfies a similar functional recurrence:
$`B_1(z,a,b)`$ $`=`$ $`{\displaystyle \frac{bz}{1abz^2}};`$ (4.3)
$`B_{w+1}(z,a,b)`$ $`=`$ $`bzB_w(z,a,{\displaystyle \frac{1}{1bz^2}}).`$ (4.4)
We note that the zeroth vertex of the path is weighted $`1`$ and that $`L_w(z,a,b)`$ counts the walk consisting of a single vertex. To compute the generating function of tails we introduce a new variable $`c`$ such that a tail whose last vertex is in the line $`y=w`$ is weighted by $`c`$. Using this variable we arrive at the following functional recurrence:
$`\stackrel{~}{T}_1(z,a,b,c)`$ $`=`$ $`{\displaystyle \frac{(1+cz)}{1abz^2}};`$ (4.5)
$`\stackrel{~}{T}_w(z,a,b,c)`$ $`=`$ $`\stackrel{~}{T}_w(z,a,{\displaystyle \frac{1}{1bz^2}},{\displaystyle \frac{1+cz}{1bz^2}}).`$ (4.6)
The relevant tail generating function is then $`T_w(z,a,b)=\stackrel{~}{T}_w(z,a,b,b)`$.
Using these results it is easy to prove by induction that the generating functions for any finite $`w`$ must be ratios of polynomials in $`z`$. In fact using induction it is possible to prove that the generating functions $`L_w(z,a,b)`$, $`B_w(z,a,b)`$ and $`\stackrel{~}{T}_w(z,a,b,c)`$ have the following forms:
$`L_w(z,a,b)`$ $`=`$ $`{\displaystyle \frac{P_w(z,0,b)}{P_w(z,a,b)}},`$ (4.7)
$`B_w(z,a,b)`$ $`=`$ $`{\displaystyle \frac{bz^w}{P_w(z,a,b)}},`$ (4.8)
$`\stackrel{~}{T}_w(z,a,b,c)`$ $`=`$ $`{\displaystyle \frac{Q_w(z,b,c)}{P_w(z,a,b)}},`$ (4.9)
where the $`P_w(z,a,b)`$ and $`Q_w(z,b,c)`$ are polynomials. By simply iterating the recurrence-functional equations and using, for example the GFUN package of MAPLE<sup>TM</sup> or by combinatorial means (Flajolet 1980 and Viennot 1985), one can find the recurrences $`P_{w+1}=P_wz^2P_{w1}`$ and $`Q_{w+1}=(1+z)Q_wz(1+z)Q_{w1}+z^3Q_{w2}`$. The solution of these, using yet another generating variable $`t`$ over width $`w`$, gives
$$G_p(z,a,b,t)\underset{w1}{}P_w(z,a,b)t^w=t\times \frac{(1abz^2)+(abab)z^2t}{1t+z^2t^2}$$
(4.10)
and
$$G_q(z,a,b,c,t)\underset{w1}{}Q_w(z,b,c)t^w=t\times \frac{(1+cz)zt(c+bz)+cz^3t^2}{(1zt)(1t+z^2t^2)}.$$
(4.11)
Alternately, one can find from substitution into the functional-recurrences (4.2) and (4.6) that
$$P_{w+1}(z,a,b)=(1bz^2)P_w(z,a,\frac{1}{1bz^2})$$
(4.12)
and
$$Q_{w+1}(z,b,c)=(1bz^2)Q_w(z,\frac{1}{1bz^2},\frac{1+cz}{1bz^2}).$$
(4.13)
By multiplying both sides by $`t^w`$ and summing over $`w`$, one obtains functional equations for $`G_p`$ and $`G_q`$, which presumably can be solved directly. We note that $`P_w`$ and $`Q_w`$ are polynomials related to Chebyshev polynomials of the second kind (see for instance Abramowitz and Stegun 1972 or Szegö 1975) as can be seen by comparing generating functions. Other techniques such as transfer matrices (Brak *et al* 1999), constant term (Brak *et al* 1998) and heaps (Bousquet-Mélou and Rechnitzer 2002) would also undoubtedly work on these problems.
## 5 Location of the generating functions singularities
### 5.1 Transformation of the generating variable
We recall that the limiting free energy of a model is determined by the dominant singularities of its generating function. Since the generating functions of these three models are all rational with the same denominator, $`P_w(z,a,b)`$, it is the smallest real positive zero of this polynomial that determines the limiting free energy. Theorem 2.1 shows that this zero (since the only zero of the numerator for bridges is $`0`$ itself) is not cancelled by a zero of the numerator. Below we find the limiting free energy by studying the zeros of $`P_w(z,a,b)`$. Given $`w`$ we would like to find an expression for $`z_c`$ as a function of $`a,b`$.
Since the generating function $`G_p(t)`$ for the $`P_w`$ polynomials is a rational function of $`t`$ we can write it in partial fraction form with respect to $`t`$, and then expand each simple rational piece to find $`P_w`$ as a function of $`w`$. The decomposition of the right-hand side of (4.10) involves expressions containing surds which becomes quite messy. It is easier first to make a substitution that greatly simplifies the partial fraction decomposition and the subsequent analysis. This is chosen so as to simplify the denominator of $`G_p(t)`$. If we set
$$z=\frac{\sqrt{q}}{1+q}$$
(5.1)
the generating function of the polynomials then simplifies to
$`G_p(\sqrt{q}/(1+q),a,b,t)`$ $`=`$ $`abab+{\displaystyle \frac{(1+qaq)(1+qbq)}{(1q)(1+qt)}}`$ (5.2)
$`{\displaystyle \frac{(1+qa)(1+qb)}{(1q)(1+qtq)}}`$
where there are simply linear factors in the denominator. The partial fraction expansion has been taken in (5.2). In our discussion of the singularities of $`P_w`$ we will start with this $`q`$-form. We can convert our results back later by making the inverse substitution
$$q=\frac{12z^2\sqrt{14z^2}}{2z^2}.$$
(5.3)
The above partial fraction form implies that
$$P_w(\sqrt{q}/(1+q),a,b)=\frac{(1+qaq)(1+qbq)}{(1q)(1+q)^{w+1}}\frac{(1+qa)(1+qb)q^w}{(1q)(1+q)^{w+1}},$$
(5.4)
and when $`a=0`$ this simplifies to
$$P_w(\sqrt{q}/(1+q),0,b)=\frac{(1+qbq)}{(1q)(1+q)^w}\frac{(1+qb)q^w}{(1q)(1+q)^w}.$$
(5.5)
We also have
$`Q_w(\sqrt{q}/(1+q),a,b,c)=`$ $`{\displaystyle \frac{q^{\frac{1}{2}+w}\left(1b+q\right)}{\left(1\sqrt{q}\right)\left(1q\right)\left(1+q\right)^w}}`$ (5.6)
$`+{\displaystyle \frac{1+qbq}{\left(1\sqrt{q}\right)\left(1q\right)\left(1+q\right)^w}}`$
$`{\displaystyle \frac{q^{\frac{w}{2}}\left(1c\left(b2c\right)\sqrt{q}+\left(1c\right)q\right)}{\left(1\sqrt{q}\right)^2\left(1+q\right)^w}}.`$
We can therefore find the generating functions for loops, bridges and tails. The generating function $`L_w(\sqrt{q}/(1+q),a,b)`$ can be written as
$$L_w=\frac{(1+q)\left[(1+qbq)(1+qb)q^w\right]}{(1+qaq)(1+qbq)(1+qa)(1+qb)q^w}.$$
(5.7)
Similarly $`B_w`$ and $`T_w`$ can be written as
$$B_w=\frac{bq^{w/2}(1q)(1+q)}{(1+qaq)(1+qbq)(1+qa)(1+qb)q^w}$$
(5.8)
and
$$T_w=\frac{\left(1+q\right)\left(q^{\frac{1+w}{2}}1\right)\left(q^{\frac{w}{2}}\left(1b+q\right)+\left(b1\right)q1\right)}{\left(\sqrt{q}1\right)\left(q^w\left(a1q\right)\left(b1q\right)\left(\left(a1\right)q1\right)\left(\left(b1\right)q1\right)\right)}.$$
(5.9)
### 5.2 Limit $`w\mathrm{}`$ of the generating functions
In the limit that $`w\mathrm{}`$ we have, for $`|q|<1`$,
$`L_w(\sqrt{q}/(1+q),a,b)`$ $`{\displaystyle \frac{(1+q)}{1+qaq}},`$ (5.10)
$`T_w(\sqrt{q}/(1+q),a,b)`$ $`{\displaystyle \frac{(1+q)}{(1+\sqrt{q})(1+qaq)}}`$ (5.11)
and $`B_w(\sqrt{q}/(1+q),a,b)`$ $`0.`$ (5.12)
We note that substituting $`z=z(q)`$ back into the limits of $`L_w`$ and $`T_w`$ gives the generating functions of loops and tails adsorbing in a half-plane (equations 3.6 and 3.7) as required.
### 5.3 Singularities of the generating function of loops
Since each of the models has the same denominator polynomial this confirms that all three models have the same limiting free energy. To determine the free energy we examine the dominant singularity of the denominator polynomial.
The dominant singularity $`q_c`$ is the zero of the denominator of $`L_w(q/(1+q),a,b)`$ that has the minimal value of $`|z|`$, and hence is the solution of
$$q^w=\frac{(1+qaq)(1+qbq)}{(1+qa)(1+qb)}$$
(5.13)
with this property.
Before discussing the cases in which we can find the singularity analytically, we first consider some examples where we locate the singularity numerically. To begin let us consider the symmetric case $`a=b`$ as this was also considered by DiMarzio and Rubin (1971). In figure 5 we plot $`\kappa ^{loop}(20;a,a)`$ with the half-plane thermodynamic limit for comparison.
For larger widths it is clear that the free energy converges to the half-plane curve as one might expect. Let us now consider weighting the vertices in the top wall such that they would strongly attract vertices of the polymer – one may expect this to hamper the convergence to the half-plane curve. Hence we consider $`b=3`$ and again for width $`w=20`$ plot the free energy $`\kappa ^{loop}(20;a,3)`$: this can be found in figure 6.
It can now be seen that for $`a<3`$ the free energy is converging to something close to $`0.75`$ and not the half-plane curve which is $`\mathrm{log}(2)`$ for $`a2`$. Physically, we can understand that a long polymer in a slit with a highly attractive top wall will stick to that wall while the attraction of the bottom wall is lesser in magnitude. However, this points to the fact that this large width limit for infinitely long polymers does *not* demonstrate the same physics as the half-plane case.
## 6 Special parameter values of the generating function
The solutions of the equation (5.13) cannot be written in closed form for all values of $`a`$ and $`b`$. However, there are particular values for which this equation simplifies. When $`a=1`$ or $`2`$ and $`b=1`$ or $`2`$, and when
$$1a=\frac{1}{1b}\text{ equivalently }ab=a+b$$
(6.1)
the right-hand side of equation (5.13) simplifies. In these cases the locations of all the singularities of the generating functions can be expressed in closed form for all widths $`w`$. It is worth considering the generating function directly as some cancellations occur at these special points.
For the special values mentioned the generating functions $`L_w`$ and $`B_w`$ have particularly simple forms, while the forms for $`T_w`$ are slightly more complicated. We give the explicit form for $`L_w`$ at each of the points below:
* when $`a=b=1`$
$$L_w=\frac{(1+q)(1q^{w+1})}{1q^{w+2}};$$
(6.2)
* when $`a=2`$ and $`b=1`$
$$L_w=\frac{(1+q)(1q^{w+1})}{(1q)(1+q^{w+1})};$$
(6.3)
* when $`a=1`$ and $`b=2`$
$$L_w=\frac{(1+q)(1+q^w)}{1+q^{w+1}};$$
(6.4)
* when $`a=b=2`$
$$L_w=\frac{(1+q)(1+q^w)}{(1q)(1q^w)};$$
(6.5)
* on the curve $`ab=a+b`$
$$L_w=\frac{(1+q)((1+qa)+(1+qaq)q^w)}{(1+qa)(1+qaq)(1q^w)}.$$
(6.6)
The singularities in $`z`$ of $`L_w`$, $`B_w`$ and $`T_w`$ may be found from the singularities in $`q`$ by the mapping (5.1). The singularities of $`L_w`$ in $`q`$ occur either for $`|q|=1`$ or for $`q^+`$. The transformation (5.1) maps $`q\{|q|=1\}^+`$ to $`z^+`$. We note that if $`q=e^{i\theta }`$ then $`1/z=2\mathrm{cos}(\theta /2)`$. We concentrate on the dominant singularity $`z_c`$ for which $`|z|`$ is minimal.
Using the expressions for $`L_w`$ given above we find the zeros of the denominators (that do not cancel with a zero of the numerator) and choose the one that minimises $`|z|`$. This gives us $`q_c`$ and hence $`z_c`$. From these we find the free energy $`\kappa `$ at the special points listed above:
* when $`a=b=1`$ the zero is $`q_c=\mathrm{exp}(2\pi i/(w+2))`$ and hence
$$\kappa (w;1,1)=\mathrm{log}\left(2\mathrm{cos}(\pi /(w+2))\right);$$
(6.7)
* when $`a=1`$, $`b=2`$ and when $`a=2`$, $`b=1`$ the zero is $`q_c=\mathrm{exp}(\pi i/(w+1))`$ and hence
$$\kappa (w;2,1)=\kappa (w;1,2)=\mathrm{log}\left(2\mathrm{cos}(\pi /2(w+1))\right);$$
(6.8)
* when $`a=b=2`$ the zero is $`q_c=1`$ and hence
$$\kappa (w;2,2)=\mathrm{log}(2);$$
(6.9)
* when $`ab=a+b`$ and $`ab`$ the zero is $`q_c=(a1)^1`$ while when $`ab=a+b`$ and $`ba`$ the zero $`q_c=(b1)^1=(a1)`$. Hence for $`a>1`$
$$\kappa (w;a,a/(a1))=\mathrm{log}\left(\frac{a}{\sqrt{a1}}\right).$$
(6.10)
In a similar way, one may write down expressions for the locations of all the singularities of the generating functions.
It is instructive to consider the free energy (see equation (6.10)) along the special curve given by $`ab=a+b`$. We first note that it is *independent* of the width of the slit. It is also clearly different at all values of $`a<2`$ to the half-plane free energy (which is $`\mathrm{log}(2)`$). This re-iterates that the large width limit of the slit problem is not necessarily described by the half-plane results.
In order to discuss the effective force between the walls for large $`w`$ we need to calculate the large $`w`$ asymptotics of the free energy. Since we have written down the closed form solution for all $`w`$ of $`\kappa (w;a,b)`$ at some special points it is advantageous to consider the large $`w`$ asymptotics of these first. As noticed above, on the curve $`ab=a+b`$ the $`\kappa `$ is independent of $`w`$ and so this includes the special point $`(a,b)=(2,2)`$. The large $`w`$ asymptotics for $`\kappa `$ for the remaining special points are calculated by expanding in inverse powers of $`w`$. They are
* when $`a=b=1`$
$$\kappa (w,1,1)=\mathrm{log}(2)\frac{\pi ^2}{2w^2}+\frac{2\pi ^2}{w^3}\frac{72\pi ^2+\pi ^4}{12w^4}+O(w^5),$$
(6.11)
* and when $`(a,b)=(1,2)`$ or $`(2,1)`$
$$\kappa (w,1,2)=\kappa (w,2,1)=\mathrm{log}(2)\frac{\pi ^2}{8w^2}+\frac{\pi ^2}{4w^3}\frac{72\pi ^2+\pi ^4}{192w^4}+O(w^5).$$
(6.12)
## 7 Asymptotics for large widths at general parameter values
As we have noticed the large $`w`$ limit is dependent on the order of limits $`w\mathrm{}`$ and $`n\mathrm{}`$. As such we can now consider the large $`w`$ asymptotics of the free energy (which itself is first defined via an $`n\mathrm{}`$ limit). As mentioned above, we are interested in the forces between the walls for large $`w`$ which also involves large $`w`$ asymptotics. However, we are unable to find a closed form solution for $`\kappa (w;a,b)`$ for arbitrary $`a,b`$ and $`w`$. We can, however, find asymptotic expressions for $`\kappa `$ for arbitrary $`a,b`$ and large $`w`$ without the need for such an expression.
For general $`a`$ and $`b`$ the analysis depends on the region of parameter space and the plane naturally breaks up into various regions. These will become different phases when considering the limit $`w\mathrm{}`$. The analysis proceeds in a self-consistent manner perturbing around the solution obtained from the special points in descending powers of $`w`$.
For $`a,b<2`$ we can find the asymptotics of $`q_c`$ as a function of $`w`$,
$$q_c=\mathrm{exp}\left(2\pi i/\left(w+\frac{2(a+bab)}{(2a)(2b)}\right)+O(w^4)\right),$$
(7.1)
which, mapping back to the $`z`$ variable, gives
$$z_c=\frac{1}{2}+\frac{\pi ^2}{4w^2}+\frac{\pi ^2(abab)}{(2a)(2b)w^3}+\left(\frac{5\pi ^4}{48}+\frac{3\pi ^2(abab)^2}{(2a)^2(2b)^2}\right)\frac{1}{w^4}+O(w^5)$$
(7.2)
and so
$$\kappa =\mathrm{log}(2)\frac{\pi ^2}{2w^2}\frac{2\pi ^2(abab)}{(2a)(2b)w^3}\left(\frac{\pi ^4}{12}+\frac{6(abab)^2}{(2a)^2(2b)^2}\right)\frac{1}{w^4}+O(w^5).$$
(7.3)
For $`b=2`$ and $`a<2`$ we have
$$q_c=\mathrm{exp}\left(\pi i/\left(w+\frac{a}{2a}\right)+O(w^4)\right)$$
(7.4)
which gives
$$z_c=\frac{1}{2}+\frac{\pi ^2}{16w^2}\frac{\pi ^2a}{8(2a)w^3}+\left(\frac{5\pi ^4}{768}+\frac{3\pi ^2a^2}{16(2a)^2}\right)\frac{1}{w^4}+O(w^5)$$
(7.5)
and so
$$\kappa =\mathrm{log}(2)\frac{\pi ^2}{8w^2}+\frac{\pi ^2a}{4(2a)w^3}\left(\frac{\pi ^4}{192}+\frac{3\pi ^2a^2}{8(2a)^2}\right)\frac{1}{w^4}+O(w^5).$$
(7.6)
When $`a=2`$ and $`b<2`$ we obtain a similar expression where $`b`$ replaces $`a`$.
When $`a`$ or $`b`$ is greater than $`2`$ the asymptotic form changes. In particular the right-hand side of equation (5.13) is satisfied by either $`q=1/(a1)`$ or $`q=1/(b1)`$ and in the limit as $`w\mathrm{}`$ the left-hand side of the equation is zero. Hence we expand about the point $`q=1/(a1)`$ and find that the next term in the expansion is exponential in $`w`$. More precisely, when $`a>b`$ and $`a>2`$ the asymptotic expansion of $`q_c(w)`$ as $`w\mathrm{}`$ is
$$q_c=\left(\frac{1}{a1}\right)\frac{a(a2)(abab)}{(a1)^2(ab)}\left(\frac{1}{a1}\right)^w+O\left(w\left(\frac{1}{a1}\right)^{2w}\right).$$
(7.7)
Transforming this back to the $`z`$ variable gives
$$z_c=\frac{\sqrt{a1}}{a}\frac{(a2)^2(abab)}{2a(ab)\sqrt{a1}}\left(\frac{1}{a1}\right)^w+O\left(w\left(\frac{1}{a1}\right)^{2w}\right)$$
(7.8)
and so the free energy is
$$\kappa =\mathrm{log}\left(\frac{a}{\sqrt{a1}}\right)+\frac{(a2)^2(abab)}{2(a1)(ab)}\left(\frac{1}{a1}\right)^w+O\left(w\left(\frac{1}{a1}\right)^{2w}\right).$$
(7.9)
For $`b>a`$ and $`b>2`$ one simply interchanges $`a`$ and $`b`$ in the above expressions to obtain the correct results.
Finally, when $`a=b>2`$ the asymptotic expansion of $`q_c(w)`$ is
$$q_c=\left(\frac{1}{a1}\right)\frac{a(a2)}{(a1)^2}\left(\frac{1}{a1}\right)^{w/2}+O\left(w\left(\frac{1}{a1}\right)^w\right)$$
(7.10)
and transforming this back to the $`z`$ variable gives
$$z_c=\frac{\sqrt{a1}}{a}\frac{(a2)^2}{2a\sqrt{a1}}\left(\frac{1}{a1}\right)^{w/2}+O\left(w\left(\frac{1}{a1}\right)^w\right)$$
(7.11)
and
$$\kappa =\mathrm{log}\left(\frac{a}{\sqrt{a1}}\right)+\frac{(a2)^2}{2(a1)}\left(\frac{1}{a1}\right)^{w/2}+O\left(w\left(\frac{1}{a1}\right)^w\right).$$
(7.12)
## 8 The infinitely wide slit
We now come to consider the limit of large slits for infinite length polymers. Even though there isn’t a closed form solution for finite widths of the equation (5.13) we can simply use the large $`w`$ asymptotics above to deduce that
$$\kappa ^{infslit}(a,b)\underset{w\mathrm{}}{lim}\kappa (w;a,b)=\{\begin{array}{cc}\mathrm{log}(2)\hfill & \text{ if }a,b2\hfill \\ \mathrm{log}\left(\frac{a}{\sqrt{a1}}\right)\hfill & \text{ if }a>2\text{ and }a>b\hfill \\ \mathrm{log}\left(\frac{b}{\sqrt{b1}}\right)\hfill & \text{ otherwise.}\hfill \end{array}$$
(8.1)
One immediately sees that the free energy depends on both $`a`$ and $`b`$ rather than only on $`a`$ which reflects the difference from the half-plane result (equation 3.9) noted earlier. In fact, the free energy is symmetric in $`a`$ and $`b`$ which reflects the observation that for infinitely long walks the end-points are irrelevant. For example, as in the case we considered earlier in figure 6 with $`b>2`$ the free energy for $`ab`$ is $`\mathrm{log}(\frac{b}{\sqrt{b1}})`$ (which is approximately $`0.752`$ for $`b=3`$) while it is $`\mathrm{log}(\frac{a}{\sqrt{a1}})`$ for $`a>b`$. The half-plane phase diagram contains a desorbed phase for $`a<2`$ and an adsorbed phase for $`a>2`$. The transition between them is second order with a jump in the specific heat. Using (8.1) one can deduce for the infinite slit that there are 3 phases: one where the polymer is desorbed for $`a,b<2`$; a phase where the polymer is adsorbed onto the bottom wall for $`a>2`$ with $`a>b`$; and a phase where the polymer is adsorbed onto the top wall for $`b>2`$ with $`b>a`$. This is illustrated in figure 7.
The low temperature adsorbed phases are characterised by the order parameter of the thermodynamic density of visits to the bottom and top walls respectively. There are 3 phase transition lines. The first two are given by $`b=2`$ for $`0a2`$ and $`a=2`$ for $`0b2`$. These lines separate the desorbed phase from the adsorbed phases and are lines of second order transitions of the same nature as the one found in the half-plane model. There is also a first order transition for $`a=b>2`$ where the density of visits to each of the walls jumps discontinuously on crossing the boundary non-tangentially.
## 9 Forces between the walls.
Let us define the effective force between the walls induced by the polymer as
$$(w)=\frac{\kappa (w)}{w}=\frac{1}{z_c(w)}\frac{z_c(w)}{w}.$$
(9.1)
Before examining the general case for large $`w`$ it is worth considering the general case for small $`w`$ numerically and at the special points for all $`w`$ where it can be found exactly. To begin, at the special points we can calculate the induced force exactly since we know the value of $`\kappa (w)`$ exactly. We note that since $`\kappa (w)`$ is monotonic at each of these points the induced force has one sign for all $`w`$. In particular,
* when $`a=b=1`$ then $`1/z_c=2\mathrm{cos}(\pi /(w+2))`$ and so
$$=\frac{\pi }{(w+2)^2}\mathrm{tan}(\pi /(w+2))=\frac{\pi ^2}{w^3}\frac{6\pi ^2}{w^4}+O(w^5)$$
(9.2)
which is positive and hence repulsive.
* When $`(a,b)=(2,1),(1,2)`$ then $`1/z_c=2\mathrm{cos}(\pi /2(w+1))`$ and so
$$=\frac{\pi }{2(w+1)^2}\mathrm{tan}(\pi /2(w+1))=\frac{\pi ^2}{4w^3}\frac{3\pi ^2}{4w^4}+O(w^5)$$
(9.3)
which is positive and hence repulsive.
* When $`ab=a+b`$ including $`a=b=2`$ then $`1/z_c=\frac{\sqrt{a1}}{a}=\frac{\sqrt{b1}}{b}`$ and so
$$=0.$$
(9.4)
If one considers more general values of $`a`$ and $`b`$ numerically one observes that in each region of the phase plane $`\kappa (w)`$ is monotonic and hence the force takes on a unique sign at each value of $`a`$ and $`b`$. At general $`(a,b)`$ we are able to find the asymptotic force for large $`w`$ by using the large $`w`$ asymptotic expression for $`\kappa (w;a,b)`$ which we calculated above.
Using the asymptotic expressions for $`\kappa `$ found in section 7 in the $`(a,b)`$-plane we obtain the asymptotics for the force.
* For $`a,b<2`$
$$=\frac{\pi ^2}{w^3}+\frac{6\pi ^2(abab)}{(2a)(2b)w^4}+O(w^5)$$
(9.5)
which is positive and hence repulsive.
* For $`b=2,a<2`$
$$=\frac{\pi ^2}{4w^3}\frac{3\pi ^2a}{4(2a)w^4}+O(w^5)$$
(9.6)
which is positive and hence repulsive.
* For $`a>b`$ and $`a>2`$
$$=\frac{(a2)^2(abab)\mathrm{log}(a1)}{2(a1)(ab)}\left(\frac{1}{a1}\right)^w+O\left(\left(\frac{1}{a1}\right)^{2w}\right).$$
(9.7)
For $`b>a`$ and $`b>2`$ one simply interchanges $`a`$ and $`b`$ in the above expressions to obtain the required result.
* At $`b=a>2`$
$$=\frac{(a2)^2\mathrm{log}(a1)}{4(a1)}\left(\frac{1}{a1}\right)^{w/2}+O\left(\left(\frac{1}{a1}\right)^w\right).$$
(9.8)
This is negative for all $`a>2`$ and so the induced force is attractive.
The regions of the plane which gave different asymptotic expressions for $`\kappa `$ and hence different phases for the infinite slit clearly also give different force behaviours. For the square $`0a,b2`$ the force is repulsive and decays as a power law (ie it is long-ranged) while outside this square the force decays exponentially and so is *short-ranged*. This change coincides with the phase boundary of the infinite slit phase diagram. However, the special curve $`ab=a+b`$ is a line of zero force across which the force, while short-ranged on either side (except at $`(a,b)=(2,2)`$), changes sign. Hence this curve separates regions where the force is attractive (to the right of the curve) and repulsive to the left of the curve. The line $`a=b`$ for $`a>2`$ is also special and, while the force is always short-ranged and attractive, the range of the force on the line is discontinuous and twice the size on this line than close by. All these features lead us to us a *force diagram* that encapsulates these features. This diagram is given in figure 8.
## 10 Discussion
We have solved and analysed in the limit of infinite length three types of polymer configuration, being loops, bridges and tails, in a two-dimensional slit geometry. The types of configurations considered are directed walks based on Dyck paths and the walks interact with both walls of the slit. We have found the exact generating functions at arbitrary width. We have calculated the free energy exactly at various points for arbitrary width, and, importantly, asymptotically (for large widths) for arbitrary $`a`$ and $`b`$. This has allowed us to map out the phase diagram for infinite slits: for finite width the free energy is an analytic function of the Boltzmann weights $`a`$ and $`b`$. This phase diagram is *different* to that obtained in the half-plane even though the generating functions for the half-plane are a formal limit of the finite width generating functions. This arises because the order of the infinite width and infinite polymer length limits (both of which are needed to see a phase transition) are *not* interchangeable.
From the large width asymptotics we have mapped regions of the plane where the free energy decreases or increases with increasing width and, also, whether this happens with an algebraic decay or exponential decay. Using this we have delineated regions where the induced force between the walls is short-ranged or long-ranged and whether it is attractive or repulsive. Three types of behaviour occur: long-ranged repulsive, short-ranged repulsive and short-ranged attractive. Applying such a model to colloidal dispersions implies that the regions where the force is long-ranged and repulsive support steric stabilisation while the regions where the force is short-ranged and attractive promote sensitized flocculation.
Given the curious interchange of limits phenomenon and from the point of view of the application of this theory to colloidal dispersions it would be interesting to investigate the finite polymer length cases both analytically and numerically with a view of searching for a scaling theory for the crossover from half-plane to slit type behaviour.
The model considered here is a directed walk in two dimensions. Although we expect the general features (such as the phase diagrams sketched in Figures 7 and 8) to be the same for self-avoiding walk models in both two and three dimensions, such models are beyond the reach of analytic treatments. It would be interesting to investigate the detailed behaviour of such models by Monte Carlo methods or by exact enumeration coupled with series analysis techniques.
## Acknowledgements
Financial support from the Australian Research Council is gratefully acknowledged by RB, ALO and AR. Financial support from NSERC of Canada is gratefully acknowledged by SGW.
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# Observations of Mkn~421 in 2004 with H.E.S.S. at large zenith angles
## 1 Introduction
Mkn 421 is a ‘BL Lac’ type active galactic nucleus. The broad-band spectral energy distribution is dominated by non-thermal emission that is believed to be produced in a relativistic jet pointing towards the observer. The high energy emission of this object has been studied by previous observations carried out by northern hemisphere ground based Cherenkov telescopes (Punch et al. punch (1992), Aharonian et al. 2002A&A…393…89A (2002), Krennrich et al. 2002ApJ…575L…9K (2002)). Observations of Mkn 421 from the southern hemisphere at large zenith angles benefit from considerable increase of the collection area at higher energies, which results in a better temporal resolution at high energies and a better sampling of the high energy part of the energy spectrum (Okumura et al. 2002ApJ…579L…9O (2002)).
Besides the general interest of understanding the physics of the highly relativistic plasma and its interaction with the ambient medium, the proximity of Mkn 421 (z=0.031) makes it an interesting target to observe the effect of pair-production of gamma-rays with soft (thermal) background photons as part of the extragalactic background light. As the collection area for large zenith angle observations exceeds two square kilometers at energies beyond 10 TeV, the observable energy spectrum could eventually be extended beyond 20 TeV, which is important to probe the mid-to-far infrared part of the extragalactic background light.
## 2 Observations and Data analyses
H.E.S.S. (Hofmann et al. hofmann (2003)) is an imaging atmospheric Cherenkov detector dedicated to the ground based observation of gamma-rays at energies above 100 GeV. Situated in Namibia (2316’S 1630’E), the full array of four telescopes is operational since December 2003. Each telescope has a mirror area of 107 m<sup>2</sup> and is equipped with a camera consisting of 960 photomultiplier tubes (Vincent et al. vincent (2003)). The system has a field of view of 5 and allows to reconstruct the direction of individual showers with a precision of better than 0.1.
The H.E.S.S. observations reported here were carried out for typically 1–2 hours per night from MJD 53107.8 to MJD 53114.9 (April 12–19, 2004) triggered by an increased level of X-ray emission detected by the All-Sky-Monitor (ASM) onboard the RXTE satellite and increased activity detected by the Whipple Cherenkov telescope (Krawczynski private communication). Online analysis of the H.E.S.S. data revealed that the source was also active at TeV energies motivating an extension of the observational campaign for roughly one week. In the beginning of May (MJD 53134.8 corresponding to May 8, 2004), the H.E.S.S. array participated in a multi-wavelength campaign with overlap with pointed X-ray observations with the RXTE satellite.
The observation mode was employing all four telescopes, pointing with 0.5 offset in declination with sign alternating from run to run (28 minutes each run). The runs were selected according to general quality checks such as absolute value of the trigger rate, relative changes in the rate and performance of the cameras (number of operational pixels, homogeneity of the acceptance). In total, 14.7 hours of good quality data were selected for the analyses. The zenith angle of the observations ranges from 60.3 at culmination to 65.4 with an average energy threshold at 1.5 TeV. The average size of images and amplitude detected in individual pixels at the given zenith angles is well below the saturation limit of the H.E.S.S. cameras and read out electronics.
The data were calibrated following the general procedure (Aharonian et al. 2004a ) and standard data reduction (image cleaning, event reconstruction) was applied (Aharonian et al. 2004astro-ph0411582 (2005)). The cuts used for background rejection were identical to the ones used in previous analyses (e.g. Aharonian et al. 2004b , 2004astro-ph0411582 (2005)) with the exception of a loose angular cut $`\theta ^2<0.05(^{})^2`$ (compared to $`\theta ^2<0.02(^{})^2`$ at smaller zenith angles), where $`\theta ^2`$ is the squared angular distance between source and reconstructed shower direction. The relaxed angular cut compensates for the degrading angular resolution at low elevations. Its value was not optimised to result in a maximum significance but rather chosen to ensure a high ($`>80\%`$) and constant acceptance for gamma-rays. The relative energy resolution at the observed zenith angle range expected from simulation is roughly 25 %.
For estimating the background, five off regions identical to the on region in size and separation from the camera center were chosen (Aharonian et al. ahar01 (2001)). The overall signal detected from Mkn 421 amounts to $`N_{\mathrm{on}}=8978`$ and $`N_{\mathrm{off}}=1357.6`$ averaged over the five background regions with a significance of $`S=114\sigma `$ using the likelihood ratio method derived in Li & Ma (1983ApJ…272..317L (1983)).
Finally, in order to reconstruct the energy spectra, collection areas were calculated from simulated air showers taking the range of observed zenith angles and the offset of the source to the camera centre into account. The collection area derived for individual events (using linear interpolation between logarithmic energy bins and linear interpolation between zenith angles) were used as weights to calculate the flux in bins of energy by summing over the events in the on and off regions. The effect of spill-over between adjacent energy bins due to the instrument’s energy resolution was compensated by using collection areas as a function of reconstructed energy (Mohanty et al. 1998APh…..9…15M (1998), Aharonian et al. 2002A&A…393…89A (2002)).
Estimates of systematic uncertainties were derived by changing the atmospheric transparency in the simulations and re-applying the analyses to the data. The small variations observed in the cosmic ray induced event rate during the observations (see also next section) are consistent with smaller variations of the atmospheric transparency than assumed in the evaluation of the systematic uncertainties.
## 3 Energy spectrum from the Crab nebula at large zenith angles
In order to check the procedure for the reconstruction of energy spectra at large zenith angles, a dedicated data-set taken on the Crab nebula for 1.8 hours at zenith angles ranging from 57 to 67 was analysed.
The fit of a power-law ($`dN/dE=N_0(E/\mathrm{TeV})^\mathrm{\Gamma }`$) to the resulting energy spectrum of the Crab nebula yields $`N_0=(2.9\pm 0.4_{\mathrm{stat}}\pm 0.7_{\mathrm{sys}})\times 10^{11}`$ ph/(cm<sup>2</sup> s TeV), $`\mathrm{\Gamma }=2.6\pm 0.1_{\mathrm{stat}}\pm 0.3_{\mathrm{sys}}`$, $`\chi ^2=7(8\mathrm{d}.\mathrm{o}.\mathrm{f}.)`$. The result is in good agreement with the recently published values given by the HEGRA collaboration (Aharonian et al. 2004c ). This gives confidence in the reliability of the instrument’s response and reconstruction technique at large zenith angles.
## 4 Energy spectra and light-curve from Mkn 421
The differential energy spectrum for Mkn 421 was calculated in the same way as described above for the Crab nebula. The entire data-set was used to derive a differential energy spectrum.
The final result for the energy spectrum of Mkn 421 is shown in Fig. 1. As the energy spectrum shows clear evidence for curvature, the collection areas were calculated iteratively until the parameters of the model fit function converged. The differential energy spectrum is well sampled between 1 and 15 TeV. Between 15 and 30 TeV, a signal (corrected for an expected spill-over of 3 events) of 2.2$`\sigma `$ with $`N_{\mathrm{on}}=31`$ and $`N_{\mathrm{off}}=96/5=19.2`$ events is seen. The collection area for this energy bin exceeds 2 km<sup>2</sup> after cuts. Between 30 and 50 TeV an upper limit with a confidence level of 99 % using the method of Feldman & Cousins (fc (1998)) was calculated. All energy bins are centred on the expected average energy for a given bin assuming a spectrum following a power-law with an exponential cutoff.
The result of a fit to the spectrum assuming a power-law with an exponential cutoff (see solid line in Fig. 1), results in $`N_0=(1.55\pm 0.08_{\mathrm{stat}}\pm 0.4_{\mathrm{sys}})\times 10^{10}`$ ph/(cm<sup>2</sup> s TeV), $`\mathrm{\Gamma }=2.1\pm 0.1_{\mathrm{stat}}\pm 0.3_{\mathrm{sys}}`$, $`E_c=3.1(+0.50.4)_{\mathrm{stat}}\pm 0.9_{\mathrm{sys}}`$ TeV, and $`\chi ^2=16(10\mathrm{d}.\mathrm{o}.\mathrm{f}.)`$. Other model fits using more complicated functions which have been commonly used in the past to describe a curved spectrum were tested: a parabola in $`\mathrm{log}(dN/dE)`$ vs. $`\mathrm{log}(E)`$ gives $`\chi ^2=12.6(9\mathrm{d}.\mathrm{o}.\mathrm{f}.)`$ and a super-exponential cutoff $`\chi ^2=11.7(9\mathrm{d}.\mathrm{o}.\mathrm{f})`$. The super-exponential cutoff fit with $`dN/dEE^\mathrm{\Gamma }\times \mathrm{exp}((E/E_c)^\alpha )`$ with $`\alpha =1.6\pm 0.3`$ (systematic uncertainty negligible) is included in Fig. 1 as a dashed line. According to the F-test, the probability of a chance improvement of the $`\chi ^2`$ for the super-exponential cutoff is 6.4 % which is not sufficiently low to give a preference to this fit function (the same holds true for the parabolic fit function).
The cutoff energy $`E_c=3.1(+0.50.4)_{\mathrm{stat}}\pm 0.9_{\mathrm{sys}}`$ TeV found in the observations described here is fully consistent with previous observations in the years 2000 and 2001 by HEGRA of $`3.6(+0.40.3)_{\mathrm{stat}}(+0.90.8)_{\mathrm{sys}}`$ TeV (Aharonian et al. 2002A&A…393…89A (2002)) and VERITAS $`(4.3\pm 0.3)`$TeV (Krennrich et al. 2002ApJ…575L…9K (2002)).
The diurnal integral fluxes above 2 TeV were calculated by taking the sum over the inverse collection area for the events with energies above the chosen value. With the threshold energy raised to 2 TeV, this gives an estimate of the flux which is independent of the variation of the actual energy threshold within a night, at the expense of loosing count statistics. As an additional benefit, the dependence on changes of the spectral shape are negligible as the collection area is not strongly changing with energy for energies well above the threshold. The resulting light-curve is shown in Fig. 2.
The light-curve exhibits night-by-night variability in the first observing week and resumes a lower flux in the later observation. The peak diurnal average flux reaches a value of about 5 times the flux observed from the Crab nebula. The corresponding gamma-ray rate is sufficiently high to probe intra-night variability of the high energy end of the spectrum with unprecedented accuracy.
During the observations of MJD 53113.8–53113.9 (April 18) and 53114.8–53114.9 (April 19), significant variations of the flux within these nights are detected. The hypothesis of a constant flux during these nights results in $`\chi ^2=40.6(7\mathrm{d}.\mathrm{o}.\mathrm{f}.)`$ and $`\chi ^2=28(7\mathrm{d}.\mathrm{o}.\mathrm{f}.)`$ respectively. The intra-night variations as seen during MJD 53113.8–53113.9 (April 18) using bins of 14 min width are shown as an inlay in Fig. 2. In order to exclude variations of the detector’s response or changes in the atmosphere to be responsible for the observed variations, the post-cut (after applying the image shape cuts) cosmic ray rate has been checked for variability which is smaller than 2 % in relative root mean square (RMS) during this night. The observed significant intra-night variability in the light-curve suggests a decay time of less than 1 hour.
The light-curve is not corrected for possible variations of the atmospheric transparency. The atmospheric transmissivity between the position of the air shower and the detector can be probed by the rate of detected cosmic ray events, as the Cherenkov light traverses the same aerosol layer which is believed to dominate possible temporal variations of the atmosphere’s transparency. During the observations, the average cosmic ray rates for individual runs varied between 113 and 143 Hz with an average of 132 Hz and RMS of 13 Hz. The post-cut rate shows even less variation with 8 % relative RMS. The observed variability of Mkn 421 is therefore clearly associated with the source and not a consequence of temporal variations of the detector response or atmospheric transparency. In principle, a correction of the measured gamma-ray flux could be derived and applied to the light-curve. This was not done as the effect is small ($`<10\%`$) in comparison to the observed variability.
In order to study variations of the spectral shape in various flux states, diurnal spectra were calculated. Given the correlation between the cutoff energy $`E_c`$ and the photon index $`\mathrm{\Gamma }`$ derived from the fit of a power-law with an exponential cutoff (correlation coefficient $`0.955`$ between $`E_c^1`$ and photon index $`\mathrm{\Gamma }`$), the power-law index was kept constant ($`\mathrm{\Gamma }=2`$) for the fits applied to spectra obtained for individual nights. The fit was applied to a fixed energy range from 1 to 10 TeV. As a result, a hardening of the spectrum or an increase of the cutoff energy is seen clearly in the correlation of the cutoff energy $`E_c`$ and the integral flux above 2 TeV in Fig. 3. A similar result was obtained when keeping the cutoff energy fixed and letting the photon index vary freely.
## 5 Multiwavelength observations in April 2004
During April 2004, a coordinated multi-wavelength campaign monitored the activity of Mkn 421 in radio, optical, X-ray, and gamma-rays (Cui et al. Cui (2004)). The source was seen to be active in X-rays where observations with the array of proportional chamber units (PCU) onboard the RXTE satellite were performed. These observations were not simultaneous with the observations with the H.E.S.S. array, but by combining the PCU with the ASM data a good temporal coverage overlapping with the H.E.S.S. observations can be achieved. An average of 4 counts/s were detected by the ASM during the first weeks of April. This is sufficiently high to probe the activity of Mkn 421 during individual ASM pointings (dwells).
The Whipple 10 m Cherenkov telescope was observing Mkn 421 simultaneously with the RXTE pointings which were generally starting within a few hours after the H.E.S.S. observations took place. First preliminary results of this campaign have been presented by Cui et al. (Cui (2004)).
In Fig. 4, the different observations are combined such that the observed flux (or count-rate) is normalized to the average flux (or count-rate) during the time between MJD 53107 and 53116. The preliminary Whipple light-curve is derived from the count-rate which is not corrected for different zenith angles of observations. Generally, the Whipple observations were carried out at small zenith angles (Cui et al. Cui (2004)) and according to previous observations, the energy threshold is estimated to be $`400`$ GeV (Krennrich et al. 2002ApJ…575L…9K (2002)). The ASM light-curve and the $`1\sigma `$ uncertainty band is obtained by calculating the sliding average over five dwells.
Even though a detailed correlation analysis is difficult because of a large fraction of the data not being simultaneous, it is interesting to note that the strongest variability occurs at the highest energies, with a maximum relative amplitude of the observed flux of $`F_{\mathrm{max}}/F_{\mathrm{min}}=4.3\pm 0.5`$ as compared to $`1.7\pm 0.2`$ derived from the Whipple count rates and $`2.05\pm 0.02`$ from the RXTE PCU observations. As a measure of the variability the relative RMS of the measurement indicates a variability of $`(17\pm 6)`$ % seen by the Whipple instrument and $`(26\pm 1)`$ % for the X-ray data taken with the PCU detectors, whereas for H.E.S.S. the relative RMS value is ($`51\pm 8)`$ %.
## 6 Conclusion
The results on the time averaged energy spectrum of Mkn 421 presented here confirm previous observations of the existence of a cutoff at $`E_c=3.1(+0.50.4)_{\mathrm{stat}}\pm 0.9_{\mathrm{sys}}`$ TeV which is lower than the cutoff energy of $`6.2\pm 0.4_{\mathrm{stat}}(+2.91.5)_{\mathrm{sys}}`$ TeV observed from Mkn~501 in 1997 (Aharonian et al. ahar99 (1999)). Given the similar red shift of these two objects, this would imply that the cutoff in the Mkn 421 spectrum is intrinsic to the source and not due to absorption. With the observations carried out at large zenith angles and increased collection area, a signal at the level of 2.2 $`\sigma `$ (corrected for the effect of spill-over) was detected between 15 and 30 TeV. This detection in conjunction with the observation of a cutoff in the energy spectrum intrinsic to the source is important for constraining the extragalactic background light at wavelengths beyond 10 $`\mu `$m. Furthermore, the energy spectrum was observed to become harder as the integral flux increases. It is not possible to discern whether the hardening is a consequence of an increase of the cutoff energy from 1.5 to 3.5 TeV or a change in the power-law index as these two parameters are highly correlated.
The relative amplitude and variance of the variability observed above 2 TeV is significantly higher than observed at lower energies during the same activity period by the Whipple telescope and the RXTE pointed X-ray detector. The sparse but similar sampling of the light curve among the different observations (H.E.S.S., Whipple, and RXTE PCU) does not introduce a bias which could explain the observed difference in the variability.
The larger variability seen at multi-TeV energies compared to energies below TeV is consistent with the observed spectral changes, where the relative increase of the flux at higher energies causes a hardening of the observed spectra. Assuming the correlation between integral flux and cut-off energy indicated in Fig. 3 and using the integral flux above 2 TeV as measured with H.E.S.S., the expected light curve above 400 GeV was calculated. As expected, the relative RMS scatter of the lower energy light curve is roughly half of the value observed above 2 TeV well consistent with the observed variability in the Whipple light curve.
###### Acknowledgements.
The support of the Namibian authorities and of the University of Namibia in facilitating the construction and operation of H.E.S.S. is gratefully acknowledged, as is the support by the German Ministry for Education and Research (BMBF), the Max Planck Society, the French Ministry for Research, the CNRS-IN2P3 and the Astroparticle Interdisciplinary Programme of the CNRS, the U.K. Particle Physics and Astronomy Research Council (PPARC), the IPNP of the Charles University, the South African Department of Science and Technology and National Research Foundation, and by the University of Namibia. We appreciate the excellent work of the technical support staff in Berlin, Durham, Hamburg, Heidelberg, Palaiseau, Paris, Saclay, and in Namibia in the construction and operation of the equipment.
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# Strangeness and the Discovery of Quark-Gluon Plasma
## 1 Introduction
The deconfined interacting quark–gluon plasma phase (QGP) is the equilibrium state of matter at high temperature and/or density. It is believed that this state has been present in the early Universe, 10-20$`\mu `$s into its evolution. The question is if, in the short time, $`10^{22}`$$`10^{23}`$ s, available in a laboratory heavy ion collision experiment, the color frozen nuclear phase can melt and turn into the QGP state of matter. There is no valid first principles answer to this question available today, nor as it seems, will a first principles simulation of the dynamic heavy ion environment become available in the foreseeable future. To address this issue we study QGP experimentally, which requires development of laboratory experiments and suitable observables.
To form QGP in the laboratory we perform relativistic heavy ion collisions in which a domain of (space, time) much larger than normal hadron size is formed, in which color-charged quarks and gluons are propagating constrained by external ‘frozen vacuum’, which abhors color . We expect a pronounced boundary in temperature and baryon density between confined and deconfined phases of matter, irrespective of the question if there is, or not, a true phase transition. We search for a boundary between phases considering the size of the interacting region and the magnitude of the reaction energy.
Detailed study of the properties of the deconfined state shows that QGP is rich in entropy and strangeness. The enhancement of entropy $`S`$ arises because the color bonds are broken and gluons can be created. Enhancement of strangeness $`s`$ arises because the mass threshold for strangeness excitation is considerably lower in QGP than in hadron matter. Moreover there are new mechanisms of strangeness formation in QGP involving reactions between (thermal) gluons. Thus $`S`$ and $`s`$ are the two elementary observables which are explored with soft hadronic probes, for further theoretical details and historical developments see our book . The numerical work presented here was carried out with the public package of programs SHARE . This report is a self contained summary of our recent results, see .
Entropy enhancement, observed in terms of enhanced hadron multiplicity per net charge, has been among the first indications of new physics reach of CERN-SPS experimental heavy ion program . The enhancement of strange hadron production both as function of the number of participating baryons, and reaction energy has been explored in several experiments at at BNL-AGS, CERN-SPS and and BNL-RHIC. We refrain from extensive historical survey of these results and present perhaps the latest, STAR-RHIC result in Fig. 1 . In this presentation one sees the yield per participant $`N_{\mathrm{part}}`$ divided by a reference yield obtained in $`pp`$ reactions. We observe that the enhancement rises both with the strangeness content in the hadron and with the size of the reaction region, indicating that the cause of this enhancement is a increased yield of strange quarks in the source, a qualitative expectation we will address in our quantitative analysis below. The gradual increase of the enhancement over the range of $`N_{\mathrm{part}}`$ is an important indicator of the physics mechanisms at work. This behavior agrees with our studies of kinetic strangeness production and strangeness yield increasing with the size of the reaction region. This enhancement of strange antibaryons which demonstrates that a novel strangeness production mechanism is present has been extensively studied at SPS energy range, where it was originally discovered .
Further evidence for parton dynamics prior to final state hadronization is obtained from the study of strange hadron transverse energy spectrum. The identity of hyperon and antihyperon spectra in particular $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}`$ and $`\mathrm{\Xi },\overline{\mathrm{\Xi }}`$ implies that both strange matter and antimatter must have been produced from a common source by the same fundamental mechanism. Furthermore they were not subject to interactions in their passage through the baryon rich hadron gas present at the SPS energy range.
The $`\overline{\mathrm{\Lambda }},\overline{\mathrm{\Xi }}`$ annihilation in baryon-rich hadron gas is strongly momentum dependent. This should deform the shape of the antihyperon spectra as compared to the spectra of hyperons. Thus symmetry of the hyperon-antihyperon spectra also implies that there was no appreciable annihilation of the $`\overline{\mathrm{\Lambda }},\overline{\mathrm{\Xi }}`$ after their formation. The working hypothesis is therefore that hadronization of the QGP deconfined phase formed in high energy nuclear collision is direct, fast (sudden) and occurs without significant sequel interactions. That can be further tested in a study of yields and spectra of unstable resonances .
At RHIC the parton level dynamics is convincingly demonstrated by quark content scaling of azimuthal asymmetry of the collective flow $`v_2`$. Further evidence is derived from the consideration of quark recombination formation of hadrons from QGP. We refer to the recent comprehensive survey of the RHIC result , and presentations at the meeting addressing these very interesting and recent developments. We will assume in this report that the case of quark-parton dynamics prior to hadronization is convincing, and will use our analysis of soft hadron production to find the thresholds of the onset of deconfinement and to determine the properties of the deconfined fireball at the time of its breakup into hadrons.
In next section 2 we introduce the statistical hadronization method of analysis of hadron production. We discuss data analysis as function of impact parameter and energy dependence of soft hadron production in section 3. We will present both the systematics of statistical model parameters, and the associated physical properties. In section 4, we address the physical QGP signatures indicating presence of a phase boundary, giving particular attention to the explanation of the ‘horn’ in the $`\mathrm{K}^+/\pi ^+`$ ratio, and the strangeness to entropy relative yield. We discuss, in section 5, the role these results play in understanding of the phase boundary separating QGP from normal confined matter. We also discuss briefly possible new soft hadron physics at LHC.
## 2 Statistical hadronization model
To describe the yields of particle produced we employ the statistical hadronization model (SHM). SHM is by definition a model of particle production in which the birth process of each particle fully saturates (maximizes) the quantum mechanical probability amplitude, and thus, the relative yields are determined solely by the appropriate integrals of the accessible phase space. For a system subject to global dynamical evolution such as collective flow, this is understood to apply within each local co-moving frame.
When particles are produced in hadronization, we speak of chemical freeze-out. Hadron formation from QGP phase has to absorb the high entropy content of QGP which originates in broken color bonds. The lightest hadron is pion and most entropy per energy is consumed in hadronization by producing these particles abundantly. It is thus important to free the yield of these particles from the chemical equilibrium constraint.
The normalization of the particle yields is, aside of the freeze-out temperature $`T`$, directly controlled by the particle fugacity $`\mathrm{{\rm Y}}_ie^{\sigma _i/T}`$, where $`\sigma _i`$ is particle ‘$`i`$’ chemical potential. Since for each related particle and antiparticle pair, we need two chemical potentials, it has become convenient to choose parameters such that we can control the difference, and sum of these separately. For example, for nucleons $`N`$, and, respectively, antinucleons $`\overline{N}`$ the two chemical factors are:
$$\sigma _N\mu _b+T\mathrm{ln}\gamma _N,\sigma _{\overline{N}}\mu _b+T\mathrm{ln}\gamma _N,$$
(1)
$$\mathrm{{\rm Y}}_N=\gamma _Ne^{\mu _b/T},\mathrm{{\rm Y}}_{\overline{N}}=\gamma _Ne^{\mu _b/T}.$$
(2)
The (baryo)chemical potential $`\mu _b`$, controls the baryon number, arising from the particle difference. $`\gamma _N`$, the phase space occupancy, regulates the number of nucleon–antinucleon pairs present. There are many different hadrons, and in principle, we could assign to each a chemical potential and then look for chemical reactions which relate these chemical potentials. However, a more direct way to accomplish the same objective consists in characterizing each particle by the valance quark content. The relation between quark based fugacity and chemical potentials ($`\lambda _{q,s}=e^{\mu _{q.,s}/T}`$) and the two principal hadron based chemical potentials of baryon number and hadron strangeness $`\mu _i,i=b,\mathrm{S}`$ is:
$$\mu _b=3\mu _q\mu _s=\frac{1}{3}\mu _b\mu _\mathrm{S},\lambda _s=\frac{\lambda _q}{\lambda _\mathrm{S}}.$$
(3)
An important (historical) anomaly is the negative S-strangeness in $`s`$-carrying-baryons. We will in general follow quark flavor and use quark chemical factors to minimize the confusion arising. In the local rest frame, the particle yields are proportional to the momentum integrals of the quantum distribution. As example, for the yield of pions $`\pi `$, nucleons $`N`$ and antinucleons $`\overline{N}`$ we have:
$`\pi `$ $`=`$ $`Vg_\pi {\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{\gamma _q^2e^{E_\pi /T}1}},E_i=\sqrt{m_i^2+p^2},\gamma _q^2<e^{m_\pi /T}`$ (4)
$`N`$ $`=`$ $`Vg_N{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{\gamma _q^3\lambda _q^3e^{E_N/T}+1}},\overline{N}=Vg_N{\displaystyle \frac{d^3p}{(2\pi )^3}\frac{1}{\gamma _q^3\lambda _q^{+3}e^{E_{\overline{N}}/T}+1}}.`$ (5)
There are two types of chemical factors $`\gamma _i`$ and $`\mu _i`$, and thus two types of chemical equilibriums. These are shown in table 1. The absolute equilibrium is reached when the phase space occupancy approaches unity, $`\gamma _i1`$. The distribution of flavor (strangeness) among many hadrons is governed by the relative chemical equilibrium.
In order to arrive at the full particle yield, one has to be sure to include all the hadronic resonances which decay feeding into the yield considered, e.g., the decay $`K^{}K+\pi `$ feeds into $`K`$ and $`\pi `$ yields. The contribution is sensitive to temperature at which these particles are formed. Inclusion of the numerous resonances constitutes a book keeping challenge in study of particle multiplicities, since decays are contributing at the 50% level to practically all particle yields. A public statistical hadronization program, SHARE (Statistical HAdronization with REsonances) has simplified this task considerably .
The resonance decay contribution is dominant for the case for the pion yield. This happens even though each resonance contributes relatively little in the final count. However, the large number of resonances which contribute compensates and the sum of small contributions competes with the direct pion yield. For the more heavy hadrons, generally there is a dominant contribution from just a few, or even from a single resonance. The exception are the $`\mathrm{\Omega },\overline{\mathrm{\Omega }}`$ which have no known low mass resonances, and also $`\varphi `$ – the known resonances are very heavy and very few.
A straightforward test of sudden hadronization and the SHM is that within a particle ‘family’, particle yields with same valance quark content are in relation to each other well described by integrals of relativistic phase space. The relative yield of, e.g., $`K^{}(\overline{s}q)`$ and $`K(\overline{s}q)`$ or $`\mathrm{\Delta }`$ and $`N`$ are controlled by the particle masses $`m_i`$, statistical weights (degeneracy) $`g_i`$ and the hadronization temperature $`T`$. In the Boltzmann limit one has (star denotes the resonance):
$$\frac{N^{}}{N}=\frac{g^{}m^{\mathrm{\hspace{0.17em}2}}K_2(m^{}/T)}{gm^2K_2(m/T)}.$$
(6)
Validity of this relation implies insensitivity of the quantum matrix element governing the coalescence-fragmentation production of particles to intrinsic structure (parity, spin, isospin), and particle mass. The measurement of the relative yield of hadron resonances is a sensitive test of the statistical hadronization hypothesis and lays the foundation to the application of the method in data analysis.
The method available to measure resonance yields depends in its accuracy significantly on the the precise nature of the hadronization process: the resonance yield is derived by reconstruction of the invariant mass of the resonance from decay products energies $`E_i`$ and momenta $`p_i`$. Should the decay products of resonances after formation rescatter on other particles, then often their energies and momenta will change enough for the invariant mass to fail the acceptance conditions set in the experimental analysis. Generally, the rescattering effect depletes more strongly the yields of shorter lived resonances , as a greater fraction of these will decay shortly after formation, when elastic scattering of decay products on other produced particles is possible.
We further hear often the argument that the general scattering process of hadrons in matter can form additional resonance states. In our opinion, the loss of observability (caused by any scattering of any of the decay products is considerably greater than a possible production gain. The loss of resonance yield provides additional valuable information about the freeze-out conditions (temperature and time duration).
## 3 Phase thresholds: volume and energy dependence
### 3.1 Statistical model parameters
In order to explore the properties of the fireball at hadronization as function of the volume at the top RHIC energy $`\sqrt{s_{NN}}=200`$ GeV Au–Au, we study the 11 centrality bins in which the $`\pi ^\pm ,\mathrm{K}^\pm ,p`$ and $`\overline{p}`$ rapidity yields $`dN/dy`$ for $`y_{\mathrm{CM}}=0`$ have been recently presented, see table I and table VIII in Ref. . These 6 particle yields and their ratios change rapidly. On the other hand, the additional two experimental results, the STAR $`\mathrm{K}^{}(892)/\mathrm{K}^{}`$ , and $`\varphi /\mathrm{K}^{}`$ show little centrality dependence.
In addition, three supplemental constraints help to determine the best fit:
A) strangeness conservation, i.e., the (grand canonical) count of $`s`$ quarks in all hadrons equals such $`\overline{s}`$ count for each rapidity unit;
B) the electrical charge to net baryon ratio in the final state is the same as in the initial state;
C) the ratio $`\pi ^+/\pi ^{}=1.\pm 0.02`$, which helps constrain the isospin asymmetry.
This last ratio appears redundant, as we already independently use the yields of $`\pi ^+`$ and $`\pi ^{}`$. These yields have a large systematic error and do not constrain their ratio well, and thus the supplemental constraint is introduced, since SHARE allows for the isospin asymmetry effect. The 7 SHM parameters (volume per unit of rapidity $`dV/dy`$, temperature $`T`$, four chemical parameters $`\lambda _q,\lambda _s,\gamma _q,\gamma _s`$ and the isospin factor $`\lambda _{I3}`$ are in this case studied in a systematic fashion as function of impact parameter using 11 yields and/or ratios and/or constraints, containing one (pion ratio) redundancy.
Although the number of degrees of freedom in such analysis is small, the $`\chi ^2`$ minimization yielding good significance is easily accomplished, showing good consistency of the data sample. The resulting statistical parameters are shown in Fig. 2, as function of participant number. We show on left results for the full non-equilibrium model allowing $`\gamma _q1,\gamma _s1`$ (full circles, blue) and semi-equilibrium setting $`\gamma _s=1`$ (open circles, red). From top to bottom the (chemical) freeze-out temperature $`T`$, the occupancy factors $`\gamma _q`$, $`\gamma _s/\gamma _q`$ and together in the bottom panel the baryon $`\mu _B`$ and hyperon $`\mu _S`$ chemical potentials. On left, in Fig. 2, we present the results for the centrality dependence, as function of participant number, on right as function of energy.
To study the energy dependence, we must assemble several different experimental results from different facilities and experiments . The results of this extensive analysis are shown on right hand side in Fig. 2. We note that when we are able to consider the total particle abundances, the number of participating nucleons is a tacit fit parameter. The lowest energy result is from our AGS study , the SPS data we used are from NA49 energy dependence exploration at the CERN-SPS . These results are for the total particle yields. The highest energy two results are based on studies of RHIC data at 130 and 200 GeV at central rapidity and address the $`dN/dy`$ particle yields, with the highest point corresponding to the results presented at greatest centrality on the left hand side in Fig. 2.
There are several relevant features in Fig. 2. We see on left, that, for $`A>20`$, there is no centrality dependence in freeze-out temperature $`T`$ and chemical potentials $`\mu _{B,S}`$ (up to the variation which can be associated with fluctuation in the data sample). However, there is a change in values of the chemical potentials with reaction energy. This is result of rapidly decreasing baryon density, which due to reduced stopping is distributed over a wider range of rapidity as the reaction energy increases.
The middle sections, in Fig. 2, address the phase space occupancies which were obtained in terms of hadron particle yields. The quark-side occupancy parameters could be considerably different, indeed as model studies show, a factor 2 lower at the discussed conditions : the hadron side phase space size is in general different from the quark-side phase space, since the particle degeneracies, particle spectra are quite different. We see similar behavior of $`\gamma _q`$ as of $`T`$ both for volume (that is, ‘wounded’ participant number $`A`$) and energy dependence seen in Fig. 2: the two lowest energy bins (top AGS and lowest SPS energy) deviate from the behavior seen at all other energies, as do the bins with $`A<20`$. $`\gamma _s/\gamma _q`$ as function of centrality rises steadily indicating longer lifespan of the fireball with increasing size. As function of energy, $`\gamma _s/\gamma _q`$ reaches a plateau at 30 $`A`$GeV, with further rise only seen in the central rapidity results at RHIC. We note that, in our analysis, there is no saturation of the $`\gamma _s`$ as we approach the most central reactions. This is inherent in the data we consider which includes the yields of $`\varphi `$, and $`K^{}`$. This result is consistent with the implication that strangeness is not fully saturated in the QGP source, though it appears over-saturated when measured in the hadron phase space.
The deviation at the most peripheral centrality bins and at lowest reaction energy from trends set by other results could be an indication of the change in the reaction mechanism. As a threshold of centrality and/or energy are crossed a relatively small value of $`\gamma _q0.5`$ grows rapidly to the maximum allowed by pion condensation condition, $`\gamma _qe^{m_\pi /2T}1.6`$. This behavior signals a transformation of a chemically under-saturated phase of matter into something novel, where chemical equilibration is easy, and results in hadronization in a over-saturated phase of matter.
Independent of the chemical (non-)equilibrium assumption, the baryochemical potential $`\mu _B=25\pm 1`$ MeV is seen across 10 centrality bins. Similarly, we find strangeness chemical potential $`\mu _S=5.5\pm 0.5`$ MeV (related to strange quark chemical potential $`\mu _s=\mu _B/3\mu _S`$). The most notable variation, in Fig. 2, is the gradual increase in strangeness phase space occupancy $`\gamma _s/\gamma _q`$ and thus strangeness yield with collision centrality. This effect was predicted and originates in an increasing lifespan of the fireball . The over-saturation of the phase space has been also expected due to both, the dynamics of expansion , and/or reduction in phase space size as a parton based matter turns into HG . This latter effect is also held responsible for the saturation of light quark phase space $`\gamma _qe^{m_\pi /2T}`$. A systematic increase of $`\gamma _s`$ with collision centrality has been reported for several reaction energies .
### 3.2 Hadronization Temperature
Let us now discuss more in depth the magnitude of the hadronization temperature which at high energy/central collisions we find at $`T=140`$ MeV. Some prefer the statistical hadronization to occur at higher temperature, perhaps as high as$`T=175`$ MeV, a point argued at this meeting in great detail by Dr. Peter Braun-Münzinger. In his presentation we heard that he believes that the lattice results will reach to such high temperature near to $`mu_\mathrm{B}=0`$ and this is where one should expect to see hadronization We disagree with both claims. For one,
We note that chemical equilibrium QCD lattice results are mature and yield $`T=163\pm 2`$ MeV when extrapolated to physical quark mass scale. This result is in a very good agreement with the prior work on 2 and 3 flavor QCD which brackets this result by $`T_{n_f=2,3}=T\pm 10`$ MeV . Moreover, the heavy ion collisions present a highly dynamical environment and one has to pay tribute to this especially regarding the value of hadronization temperature. For this reason we do not expect to find that the observed hadronization condition $`T(\mu _\mathrm{B})`$ will line up with the phase boundary curve obtained in study of a statistical system in thermodynamic limit on the lattice:
a) Dependence on parton collective flow:
A widely discussed effect which displaces the hadronization condition from the phase boundary is the expansion dynamics of the fireball. When the collective flow occurs at parton level, the color charge flow, like a wind, pushes out the vacuum , adding to thermal pressure a dynamical component. This can in general lead to supercooling and a sudden breakup of the fireball. We find that this can reduce the effective hadronization temperature by up to 20 MeV .
b) Dependence on quark chemical equilibration:
Lattice result have been discussed for 2-flavor lattice QCD at corresponding to $`\gamma _q=1,\gamma _s=0`$ (called) and for 2+1 flavor, corresponding to $`\gamma _q=\gamma _s=1`$. While the precise nature of the phase limit is still under study it appears that for 2-flavor case the phase boundary temperature rises by about 7-10 MeV compared to the 2+1 case. We refer to the recent review of lattice QCD for further details . Similarly the phase limit in pure gauge case corresponding in loose sense to $`\gamma _q=\gamma _s=0`$ was seen near or even above $`T=200`$ MeV.
These results do suggest that presence and the number of quarks matters regarding the precise location of the phase boundary and its nature. Its importance could be greatly enhanced, should the over-saturation of quark phase space have the same effect as would additional quark degrees of freedom. These are known to cause even for $`\mu _\mathrm{B}=0`$ the conversion of the phase crossover into a 1st-order phase transition which would, with these additional degrees of freedom, be expected at just the temperature we find in the SHM analysis.
We can be nearly sure that the chemical conditions matter and can displace the transition temperature. Because the degree of equilibration in the QGP depends on the collision energy, as does the collective expansion velocity, we cannot at all expect a simple hadronization scheme appropriate for the hadronization of nearly adiabatically expanding Universe.
Leaving this issue we note that, in the data analysis assuming chemical equilibrium, we find $`T=155\pm 8`$ MeV for the chemical equilibrium and strangeness non-equilibrium freeze-out, see Fig. 2. The error is our estimate of the propagation of the systematic data error, combined with the fit uncertainty; the reader should note that the error comparing centrality to centrality is negligible. The freeze-out temperature is for the semi-equilibrium and equilibrium model about 10% greater than the full chemical non-equilibrium freeze-out.
This result for $`T`$, in the equilibrium case, is in mild disagreement (1.5 s.d.) with earlier equilibrium fits . This, we believe, is due to some differences in data sample used, specifically, the hadron resonance production results used provide a very strong constraint for the fitted temperature, and more complete treatment by SHARE of hadron mass spectrum.
Interestingly, it seems that the general consensus about the chemical equilibrium best analysis result is in gross disagreement with the results advanced at this meeting by Dr. Peter Braun-Münzinger.
### 3.3 Physical Properties of the Fireball
We now turn our attention to the physical properties of the hadronizing fireball obtained summing individual contributions made by made of each of the hadronic particles produced. Often particles observed experimentally dominate (e.g. pions dominate pressure, kaon strangeness yield etc). However, it is important to include in this yields of particles predicted in terms of the SHM fit to the available results. Again, on left in Fig. 3. we show the behavior as function of impact parameter and on right as function of energy. From top to bottom we show the pressure $`P`$, energy density $`ϵ`$, entropy density $`\sigma `$, and the dimensionless ratio $`ϵ/T\sigma =E/TS`$. All contributions are evaluated using relativistic expression, see .
When we fitted particle rapidity yields, the global fitted yield normalization factor is $`dV(A)/dy`$, for total particle yields it is $`V`$ which is also a function of centrality trigger. When we consider ratios of two bulk properties e.g. $`E/TS`$, the results are in general smoother indicating cancellation in the error propagating from the fit. The overall error is of the magnitude of the particle yield, for example pressure is dominated by pions and hence its precision is limited by this error. However, on left, the point-to-point error is minimal as the systematic error is common. On right, the absolute error matters as the fluctuations in the results presented show.
We note, in Fig. 3, that as the reaction energy passes the volume threshold $`A=20`$ and even more so, the energy threshold $`6.26\mathrm{GeV}<\sqrt{s_{\mathrm{NN}}}<7.61\mathrm{GeV}`$ the hadronizing fireball becomes much denser. The entropy density jumps by factor 3–4, and the energy and baryon number density by a factor 2–3. The hadron pressure jumps up from $`P=25\mathrm{MeV}/\mathrm{fm}^3`$ initially by factor 2 and ultimately more than factor 3. There is a gradual increase of $`P/ϵ=0.115`$ at low reaction energy to 0.165 at the top available energy. It is important to note that exactly the same behavior of the fireball physical properties arises both as function of reaction volume and reaction energy. This is the case for both the physical properties and the statistical parameters. We believe that this shows a common change in the physical state created as function of energy available and reaction volume.
## 4 Search for an Energy Threshold
### 4.1 Kaon to pion ratio
One of the most interesting questions is if there is an energy threshold for the formation of a new state of matter. An important observable of the deconfined phase of matter is aside of strangeness, the high entropy content, which is arising from broken color bonds. The observable for both is the K$`{}_{}{}^{+}/\pi ^+\overline{s}/\overline{d}`$ yield ratio . This ratio has been studied experimentally and a pronounced horn structure arises. We can describe this structure in our study of the particle yields only within the chemical non-equilibrium model. Although this change is associated a rather sudden modification of chemical conditions in the dense matter fireball, this effect is caused by two distinct phenomena: the rapid rise in strangeness $`\overline{s}`$ production below, and a rise in the antiquark $`\overline{d}`$ yield above a energy reaction threshold.
The measured $`\mathrm{K}^+/\pi ^+`$ ratio by NA49 is shown at top left of Fig. 4, where for comparison we also show the $`pp`$ results. On right top, we present our results reduced to the correct relative scale, both for the total yield ratio for the AGS–RHIC energy range, and for the central rapidity results from RHIC. The solid line guides the eye to the fit results we obtained. To show that the $`K^+/\pi ^+`$ ratio drop is due to a decrease in baryon density which leads to rise in the $`\overline{d}`$ yield, we show in bottom section of Fig. 4, the nearly baryon density independent $`K/\pi `$ double ratio ratio Eq. (7), on left as function of $`\sqrt{s_{\mathrm{NN}}}`$, and on right as function of centrality of the reaction for $`\sqrt{s_{\mathrm{NN}}}=200`$ GeV:
$$\frac{K}{\pi }=\sqrt{\frac{K^+}{\pi ^+}\frac{K^{}}{\pi ^{}}}.$$
(7)
Both upper and lower portion of Fig. 4 are also drawn on same relative scale.
Seen how the horn specifically arises in the one $`K^+/\pi ^+`$, one can wonder if this is really a physical effect and how, in qualitative terms, a parameter $`\gamma _q`$, which controls the light quark yield, can help explain the horn structure seen in top of Fig. 8. We observe that this horn structure in the $`\mathrm{K}^+/\pi ^+`$ ratio traces out the final state valance quark ratio $`\overline{s}/\overline{d}`$, and in language of quark phase space occupancies $`\gamma _i`$ and fugacities $`\lambda _i`$, we have:
$$\frac{\mathrm{K}^+}{\pi ^+}\frac{\overline{s}}{\overline{d}}F(T)\left(\frac{\lambda _s}{\lambda _d}\right)^1\frac{\gamma _s}{\gamma _d}=F(T)\left(\sqrt{\lambda _{I3}}\frac{\lambda _s}{\lambda _q}\right)^1\frac{\gamma _s}{\gamma _q}.$$
(8)
In chemical equilibrium models $`\gamma _s/\gamma _q=1`$, and the $`\mathrm{K}^+/\pi ^+`$ ratio and its horn must arise solely from the variation in the ratio $`\lambda _s/\lambda _q`$ and the change in temperature $`T`$ which both are usually smooth function of reaction energy. The isospin factor $`\lambda _{I3}`$ is insignificant in this consideration.
As collision energy is increased, increased hadron yield leads to a decreasing $`\lambda _q=e^{\mu _\mathrm{B}/3T}`$. We recall the smooth decrease of $`\mu _\mathrm{B}`$ with reaction energy seen in bottom panel in Fig. 2. The two chemical fugacities $`\lambda _s`$ and $`\lambda _q`$ are coupled by the condition that the strangeness is conserved. This leads to a smooth $`\lambda _s/\lambda _q`$. The chemical potential effect is suggesting a smooth increase in the K$`{}_{}{}^{+}/\pi ^+`$ ratio.
For the interesting range of freeze-out temperature, $`F(T)`$ is a smooth function of $`T`$. Normally, one expects that $`T`$ increases with collision energy, hence on this ground as well we expect an monotonic increase in the $`\mathrm{K}^+/\pi ^+`$ ratio as function of reaction energy. With considerable effort, one can arrange the chemical equilibrium fits to bend over to a flat behavior at $`\sqrt{s_{\mathrm{NN}}^{\mathrm{cr}}}`$. It is quasi impossible to generate the horn with chemical equilibrium model. To accomplish this, an additional parameter appears necessary, capable to change rapidly when hadronization conditions change. This is $`\gamma _q`$. It presence also allows $`T`$ to vary in non-monotonic fashion, as is seen in Fig. 2.
### 4.2 QGP degrees of freedom and $`s/S`$ ratio
The full SHM is capable to describe the data, and we now show that it can pinpoint the properties of the phase of matter that was created early on in the reaction. To see this we consider what we learn from the final state data about strangeness and entropy production. For this purpose we consider both the specific per baryon and per entropy yield of strangeness. In addition, we look at the cost in thermal energy to make strangeness. All these quantities are nearly independent of the dynamics of hadronization, since they are related to processes occurring early on, ‘deep’ inside the collision region, and long before hadronization.
In the QGP, the dominant entropy production occurs during the initial glue thermalization, and the thermal strangeness production occurs in parallel and/or just a short time later . The entropy production occurs predominantly early on in the collision during the thermalization phase. Strangeness production by gluon fusion is most effective in the early, high temperature environment, however it continues to act during the evolution of the hot deconfined phase until hadronization . Both strangeness and entropy are nearly conserved in the evolution towards hadronization and thus the final state hadronic yield analysis value for $`s/S`$ is closely related to the thermal processes in the fireball at $`\tau 1`$–2 fm/c. We believe that for reactions in which the system approaches strangeness equilibrium in the QGP phase, one can expect a prescribed ratio of strangeness per entropy, the value is basically the ratio of the QGP degrees of freedom.
We estimate the magnitude of $`s/S`$ deep in the QGP phase, considering the hot stage of the reaction. For an equilibrated non-interacting QGP phase with perturbative properties:
$$\frac{s}{S}\frac{\rho _\mathrm{s}}{\sigma }=\frac{(3/\pi ^2)T^3(m_s/T)^2K_2(m_s/T)}{(32\pi ^2/45)T^3+n_\mathrm{f}[(7\pi ^2/15)T^3+\mu _q^2T]}=\frac{0.027}{1+0.054(\mathrm{ln}\lambda _q)^2}.$$
(9)
Here, we used for the number of flavors $`n_\mathrm{f}=2.5`$ and $`m_s/T=1`$. We see that the result is a slowly changing function of $`\lambda _q`$; for large $`\lambda _q4`$, we find at modest SPS energies, the value of $`s/S`$ is reduced by 10%. Considering the slow dependence on $`x=m_s/T1`$ of $`W(x)=x^2K_2(x)`$ there is minor dependence on the much less variable temperature $`T`$.
The dependence on the degree of chemical equilibration which dominates is easily obtained separating the different degrees of freedom:
$$\frac{s}{S}=0.027\frac{\gamma _s^{\mathrm{QGP}}}{0.38\gamma _\mathrm{G}^{\mathrm{QGP}}+0.12\gamma _s^{\mathrm{QGP}}+0.5\gamma _q^{\mathrm{QGP}}+0.054\gamma _q^{\mathrm{QGP}}(\mathrm{ln}\lambda _q)^2}.$$
(10)
We assume that the interaction effects are at this level of the discussion canceling. Seen Eq. (10) we expect to see a gradual increase in $`s/S`$ as the QGP source of particles approaches chemical equilibrium with increasing collision energy and/or increasing volume.
We repeat that it is important to keep in mind that this ratio $`s/S`$ is established early on in the reaction, the above relations, and the associated chemical conditions we considered, apply to the early hot phase of the fireball. At hadron freeze-out the QGP picture used above does not apply. Gluons are likely to freeze faster than quarks and both are subject to much more complex non-perturbative behavior. However, the value of $`s/S`$ is nearly preserved from the hot QGP to the final state of the reaction.
How does this simple prediction compare to experiment? Given the statistical parameters, we can evaluate the yields of particles not yet measured and obtain the rapidity yields of entropy, net baryon number, net strangeness, and thermal energy, both for the total reaction system and also for the central rapidity condition, also as function of centrality. In passing, we note that, in the most central reaction bin at RHIC-200, $`dB/dy15`$ baryons per unit rapidity interval, implying a rather large baryon stopping in the central rapidity domain.
The rise of strangeness yield with centrality is faster than the rise of baryon number yield: $`(ds/dy)/(dB/dy)s/B`$ is seen in the top left panel in Fig. 5. The solid (blue) lines are for the chemical nonequilibrium central rapidity yields of particles at RHIC-200. Solid (green) lines, on right, are for total hadron yields and thus total yields of the considered quantities, e.g. strangeness, entropy. For the most central head-on reactions, we reach at RHIC-200 $`s/B=9.6\pm 1`$.
In the middle panel in Fig. 5 we compare strangeness with entropy production $`s/S`$, which we just evaluated theoretically. Again, on left, as function of participant number $`A`$ at RHIC-200 and, on right, for the total reaction system for the most central 5-7% reactions. On left, we see a smooth transition from a flat peripheral behavior where $`s/S\stackrel{<}{}0.02`$ to smoothly increasing $`s/S`$ reaching $`s/S0.028`$ in most central reactions. This indicates that even at RHIC-200 for the more central reactions some new mechanisms of strangeness production becomes activated. On right, we see that the change on $`s/S`$ is much more drastic as function of reaction energy. After an initial rapid rise the further increase occurs beyond the threshold energy at slower pace.
In the bottom panel, on left, we see the thermal energy cost $`E_{\mathrm{th}}/s`$ of producing a pair of strange quarks as function of the size of the participating volume (i.e. $`A`$) This quantity shows a smooth drop which can be associated with transfer of thermal energy into collective transverse expansion after strangeness is produced. Thus, it seems that the cost of strangeness production is independent of reaction centrality. The result is different when we consider $`\sqrt{s_{\mathrm{NN}}}`$ dependence of this quantity, see bottom panel on right. There is a very clear change in the energy efficiency of making strangeness at the threshold energy. We will return to discuss possible reaction mechanisms below.
## 5 Final remarks
### 5.1 Phase boundary and hadronization conditions
The chemical freeze-out conditions we have determined presents, in the $`T`$$`\mu _\mathrm{B}`$ plane, a more complex picture than naively expected, see Fig. 6. Considering results shown in Fig. 2, we are able to assign to each point in the $`T`$$`\mu _\mathrm{B}`$ plane the associate value of $`\sqrt{s_{\mathrm{NN}}}`$. The RHIC $`dN/dy`$ results are to outer left. They are followed by RHIC and SPS $`N_{4\pi }`$ results. The dip corresponds to the 30 and 40 $`A`$GeV SPS results. The top right is the lowest 20 $`A`$GeV SPS and top 11.6 $`A`$GeV AGS energy range. To guide the eye, we have added two lines connecting the fit results. We see that the chemical freeze-out temperature $`T`$ rises for the two lowest reaction energies 11.6 and 20 $`A`$ GeV to near the Hagedorn temperature, $`T=160`$ MeV, of boiling hadron matter.
Once the chemical non-equilibrium is allowed for, the data fit turns to be much more precise, and the picture of the phase boundary with smaller ‘measurement’ error reveal a much more complex structure, and contains physics details prior analysis based on a rudimentary model could not uncover. The shape of the hadronization boundary, shown in Fig. 6 in the $`T`$$`\mu _\mathrm{B}`$ plane, is the result of a complex interplay between the dynamics of heavy ion reaction, and the properties of both phases of matter, the inside of the fireball, and the hadron phase we observe. The dynamical effect, capable to shift the location in temperature of the expected phase boundary is due to the expansion dynamics of the fireball which occurs at parton level, and effects of chemical nonequilibrium, see subsection 3.2 for full discussion.
Possibly, not only the location, but also the nature of the phase boundary can be modified by variation of $`\gamma _i`$. We recall that for the 2+1 flavor case, there is a critical point at finite baryochemical potential with $`\mu _\mathrm{B}350`$ MeV . For the case of 3 massless flavors, there can be a 1st order transition at all $`\mu _\mathrm{B}`$ . Considering a classical particle system, one easily sees that an over-saturated phase space, e.g., with $`\gamma _q=1.6,\gamma _s\gamma _q`$ for the purpose of the study of the phase transition acts as being equivalent to a system with 3.2 light quarks and 1.6 massive (strange) quarks present in the confined hadron phase.
Considering the sudden nature of the fireball breakup seen in several observables , we conjecture that the hadronizing fireball leading to $`\gamma _s>\gamma _q=1.6`$ super-cools and experiences a true 1st order phase boundary corresponding also at small $`\mu _\mathrm{B}`$. The system we observe in the final state prior to hadronization is mainly a quark–antiquark system with gluons frozen in prior expansion cooling of the QCD deconfined parton fluid. The quark dominance is necessary to understand e.g. how the azimuthal asymmetry $`v_2`$ varies for different particles .
These quarks and antiquarks have, in principle, at that stage a significant thermal mass. The evidence for this is seen in Fig. 3 in its two bottom panels, showing the dimensionless variable $`E/TS`$. The energy end entropy per particle of non-relativistic and semi-relativistic classical particle gas comprising both quarks and antiquarks is (see section 10, ):
$$\frac{E}{N}m+3/2T+\mathrm{},\frac{S}{N}5/2+m/T+\mathrm{},\frac{E}{TS}\frac{m/T+3/2}{m/T+5/2}.$$
(11)
It is thus possible to interpret the fitted value $`E/TS0.78`$ in terms of a quark matter made of particles with $`maT`$, $`a=2`$ which is close to what is expected based on thermal QCD .
### 5.2 Looking forward to LHC
We expect considerably more violent transverse expansion of the fireball of matter created at LHC. The kinetic energy of this transverse motion must be taken from the thermal energy of the expanding matter, and ultimately this leads to local cooling and thus a reduction in the number of quarks and gluons. The local entropy density decreases, but the expansion assures that the total entropy still increases. Primarily, gluons are feeding the expansion dynamics, while strange quark pair yield, being weaker coupled remain least influenced by this dynamics. Model calculations show that this expansion yields an increase in the final QGP phase strangeness occupancy $`\gamma _s`$ . This mechanism, along with the required depletion of the non-strange degrees of freedom, in the feeding of the expansion, assures an increase in the $`K/\pi `$ ratio given the nearly 30-fold increase of collision energy.
Depending on what we believe to be a valid hadronization temperature for a fast transversely expanding fireball, the possible maximal enhancement in the $`K/\pi `$ ratio may be in the range of a factor 2–3. Perhaps even more interesting than the $`K/\pi `$ ratio enhancement would be the enhancement anomaly in strange (antibaryon) yields. With $`\gamma _s1`$, we find that the more strange baryons and antibaryons are more abundant than the more ‘normal’ species. Specifically of interest would be $`(\mathrm{\Omega }^{}+\overline{\mathrm{\Omega }}^+)/(h^++h^{})`$, $`(\mathrm{\Xi }^{}+\overline{\mathrm{\Xi }}^+)/(h^++h^{})`$, and $`2\varphi /(h^++h^{})`$ which should show a major, up to an order of magnitude shift in relative production strength. Detailed predictions for the yields of these particles require considerable extrapolation of physics conditions from the RHIC to LHC domain and this work is in progress .
Ultimately, strange $`s`$, and $`\overline{s}`$ quarks can exceed in abundance the light quark component, in which case we would need to rethink in much more detail the distribution of global particle yield. The ratios of neutral and charged hadrons would undergo serious change.
Regarding charm, we note that situation will not become as severe. Given the large charm quark mass, we expect that most of charm quark yield is due to first hard interactions of primary partons. For this reason, the yield of strange and light quarks, at time of hadronization, exceeds by about a factor 100 or more that of charm at central rapidity. Thus, even though charm phase space occupancy at hadronization may largely exceed the chemical equilibrium value, seen the low hadronization temperature, e.g., $`m_c/T10`$, it takes a factor $`\gamma _ce^{10}/10^{1.5}=700`$ to compensate hadron yield suppression due to the high charm mass. Said differently, while strange quarks can compete in abundance with light quarks for $`m_s/T1`$, charm (and heavier) flavor(s) will remain suppressed, in absolute yield, at the temperatures we can make presently in laboratory experiments.
### 5.3 Highlights
We have shown that strangeness, and entropy, at SPS and RHIC are well developed tools allowing the detailed study of hot QGP phase. Our detailed discussion of hadronization analysis results has further shown that a systematic study of strange hadrons fingerprints the properties of a new state of matter at point of its breakup into final state hadrons.
We have shown that it is possible to describe the ‘horn’ in the $`K^+/\pi ^+`$ hadron ratio within the chemical non-equilibrium statistical hadronization model. We have shown that appearance of this structure is related to a rapid change in the properties of the hadronizing matter. Of most theoretical relevance and interest are the implications of non-equilibrium hadronization on the possible change in the location and nature of the phase boundary.
In summary, we have presented interpretation of the experimental soft hadron data production and discussed the production of strangeness and entropy that this analysis infers. We have seen, in quantitative way, how the relative strangeness and entropy production in most central high energy heavy ion collisions agrees with quark-gluon degree of freedom counting in hot primordial matter where the values of these quantities have been established.
### Acknowledgments
Work supported in part by a grant from: the U.S. Department of Energy DE-FG02-04ER4131. LPTHE, Univ. Paris 6 et 7 is: Unité mixte de Recherche du CNRS, UMR7589. JR thanks Bikash Sinha and Jan-e Alam and the organizers of the 5th International Conference on Physics and Astrophysics of Quark Gluon Plasma, February 8 — 12, 2005 Salt Lake City, Kolkata, India for their kind hospitality. Dedicated to Professor Bikash Sinha on occasion of his 60th anniversary.
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# 1 Introduction
## 1 Introduction
Over the past year, major progress in the calculation of scattering amplitudes in perturbative Yang-Mills theory has been made. This was triggered by Witten’s discovery that tree-level amplitudes in Yang-Mills can equivalently be derived via a string theory calculation, where the string theory in question is the topological B model with target space a supersymmetric version of Penrose’s twistor space . Witten also observed that tree-level scattering amplitudes, when Fourier transformed to twistor space, have an interesting geometrical structure, namely they have support on algebraic curves; for the simple case of the maximally helicity violating (MHV) amplitude, described by the Parke-Taylor formula, the curve is just a line (for real twistor space). This remarkable observation gives an explanation for the unexpected and previously rather mysterious simplicity of tree-level scattering amplitudes in Yang-Mills such as the Parke-Taylor formula, which is not at all apparent in a calculation performed using standard Feynman rules.
On a different line of development, the simplicity of tree-level scattering amplitudes was linked to the existence of novel recursion relations discovered by Britto, Cachazo and Feng (BCF) , and subsequently proved by the same authors and Witten (BCFW) . The elegant proof of is based on very general properties of amplitudes, such as analyticity and factorisation on multiparticle poles, and hence gave rise to the hope that recursion relations may arise in very different contexts. Indeed, novel recursion relations were also found in general relativity , scalar theory , for the finite rational amplitudes at one-loop in Yang-Mills and massless QCD , and for tree amplitudes involving massive scalars and gluons in Yang-Mills .
The simplicity of tree-level amplitudes in Yang-Mills was exploited by Bern, Dixon, Dunbar and Kosower (BDDK) in order to build one-loop scattering amplitudes . By applying unitarity at the level of amplitudes, rather than Feynman diagrams, these authors were able to construct many one-loop amplitudes in supersymmetric theories, such as the infinite sequence of MHV amplitudes in $`𝒩=4`$ and in $`𝒩=1`$ super Yang-Mills (SYM). The unitarity method of BDDK by-passes the use of Feynman diagrams and its related complications, and generates results of an unexpectedly simple form; for instance, the one-loop MHV amplitude in $`𝒩=4`$ SYM is simply given by the tree-level expression multiplied by a sum of “two-mass easy” box functions, all with coefficient one. As a side remark, we would like to mention that higher-loop amplitudes in $`𝒩=4`$ SYM also display intriguing regularities .
The geometrical structure in twistor space of the amplitudes was also the root of a further important development. In , Cachazo, Svrček and Witten (CSW) proposed a novel perturbative expansion for on-shell amplitudes in Yang-Mills, where the MHV amplitudes are lifted to vertices, joined by simple scalar propagators in order to form amplitudes with an increasing number of negative helicities. Applications at tree level confirmed the validity of the method and led to the derivation of various new amplitudes in gauge theory .
In , a heuristic derivation of the CSW method was given from the twistor string theory. Rather unfortunately, the latter only appears to describe the scattering amplitudes of Yang-Mills at tree level , as at one loop states of conformal supergravity enter the game, and cannot be decoupled in any known limit. The duality between gauge theory and twistor string theory is thus spoiled by quantum corrections. Surprisingly, it was found by three of the present authors that the MHV method at one-loop level does succeed in correctly reproducing the scattering amplitudes of the gauge theory . Furthermore, the twistor space picture of one-loop amplitudes is now in complete agreement with that emerging from the MHV methods, which suggests that the amplitudes at one loop have localisation properties on unions of lines in twistor space; an initial puzzle was indeed clarified and explained in terms of a certain “holomorphic anomaly”, introduced in , and further analysed in . A proof of the MHV method at tree level was finally given in ; at loop level, however, it remains a (well-supported) conjecture.
The initial successful application of the MHV method to $`𝒩=4`$ SYM was followed by calculations of MHV amplitudes in $`𝒩=1`$ SYM , and in pure Yang-Mills , where the four-dimensional cut-constructible part of the infinite sequence of MHV amplitudes was derived. However, amplitudes in non-supersymmetric Yang-Mills theory also have rational terms which escape analyses based on MHV diagrams at one loop or four-dimensional unitarity .
Amplitudes in supersymmetric theories are of course special. They do contain rational terms, but these are uniquely linked to terms which have cuts in four dimensions. In other words, these amplitudes can be reconstructed uniquely from their cuts in four-dimensions – a remarkable result. These cuts are of course four-dimensional tree-level amplitudes, whose simplicity is instrumental in allowing the derivation of analytic, closed-form expressions for the one-loop amplitudes. In non-supersymmetric theories, amplitudes can still be reconstructed from their cuts, but on the condition of working in $`42ϵ`$ dimensions, with $`ϵ0`$ . This is a powerful statement, but it also implies the rather unpleasant fact that one should in principle work with tree-level amplitudes involving gluons continued to $`42ϵ`$ dimensions, which are not simple.
An important simplification is offered by the well-known supersymmetric decomposition of one-loop amplitudes of gluons in pure Yang-Mills. Given a one-loop amplitude $`𝒜_\mathrm{g}`$ with gluons running the loop, one can re-cast it as
$$𝒜_\mathrm{g}=(𝒜_\mathrm{g}+\mathrm{\hspace{0.17em}4}𝒜_\mathrm{f}+\mathrm{\hspace{0.17em}3}𝒜_\mathrm{s})4(𝒜_\mathrm{f}+𝒜_\mathrm{s})+𝒜_\mathrm{s}.$$
(1.1)
Here $`𝒜_\mathrm{f}`$ ($`𝒜_\mathrm{s}`$) is the amplitude with the same external particles as $`𝒜_\mathrm{g}`$ but with a Weyl fermion (complex scalar) in the adjoint of the gauge group running in the loop. This decomposition is useful because the first two terms on the right hand side of (1.1) are contributions coming from an $`𝒩=4`$ multiplet and (minus four times) a chiral $`𝒩=1`$ multiplet, respectively; therefore, these terms are four-dimensional cut-constructible, which simplifies their calculation enormously. The last term in (1.1), $`𝒜_\mathrm{s}`$, is the contribution coming from a scalar running in the loop. The key point here is that the calculation of this term is much easier than that of the original amplitude $`𝒜_\mathrm{g}`$. It is this last contribution which is the focus of this paper.
The root of the simplification lies in the fact that a massless scalar in $`42ϵ`$ dimensions can equivalently be described as a massive scalar in four dimensions . Indeed, if $`L`$ is the $`(42ϵ)`$-dimensional momentum of the massless scalar ($`L^2=0`$), decomposed into a four-dimensional component $`l_{(4)}`$ and a $`2ϵ`$-dimensional component $`l_{(2ϵ)}`$, $`L:=l_{(4)}+l_{(2ϵ)}`$, one has $`L^2:=l_{(4)}^2+l_{(2ϵ)}^2=l_{(4)}^2\mu ^2`$, where $`l_{(2ϵ)}^2:=\mu ^2`$ and the four-dimensional and $`2ϵ`$-dimensional subspaces are taken to be orthogonal. The tree-level amplitudes entering the $`(42ϵ)`$-dimensional cuts of a one-loop amplitude with a scalar in the loop are therefore those involving a pair of massive scalars and gluons. Crucially, these amplitudes have a rather simple form. Some of these amplitudes appear in ; furthermore, a recent paper describes how to efficiently derive such amplitudes using a recursion relation similar to that of BCFW.
Using two-particle cuts in $`42ϵ`$ dimensions, together with the supersymmetric decomposition mentioned above, various amplitudes in pure Yang-Mills were derived in recent years, starting with the pioneering works . In this paper we show that this analysis can be performed with the help of an additional tool: generalised $`(42ϵ)`$-dimensional unitarity.
Generalised four-dimensional unitarity was very efficiently applied in to the calculation of one-loop amplitudes in $`𝒩=4`$ SYM. Amplitudes in this theory can be written as a sum of box functions, multiplied by rational coefficients. To each box function is uniquely associated a (generalised) quadruple cut, so that, schematically, each coefficient of a box function is expressed as a particular quadruple cut of the one-loop amplitude, which is nothing but a product of four tree-level amplitudes. Generalised cuts require the amplitudes to be continued to complexified Minkowski space, which in turn has the consequence that three-point amplitudes no longer vanish, and enter the cut-amplitude in an important way .<sup>1</sup><sup>1</sup>1This circumstance extends to the $`(42ϵ)`$-dimensional three-point scattering amplitudes which will be considered in this paper. The calculation of one-loop amplitudes in $`𝒩=4`$ SYM was in this way turned into an algebraic problem . Using generalised unitarity in four dimensions, the infinite sequence of next-to-MHV amplitudes in $`𝒩=4`$ SYM was determined ; generalised unitarity was also applied to $`𝒩=1`$ SYM, in particular to the calculation of the next-to-MHV amplitude with adjacent negative-helicity gluons . These amplitudes can be expressed solely in terms of triangles, and were efficiently computed in using triple cuts.<sup>2</sup><sup>2</sup>2A new calculation based on localisation in spinor space was also introduced in .
The main point of this paper is the observation that generalised unitarity is actually a useful concept also in $`42ϵ`$ dimensions; in turn this means that generalised $`(42ϵ)`$-dimensional unitarity is relevant for the calculation of non-supersymmetric amplitudes at one loop. In particular in this paper we will be able to compute amplitudes in non-supersymmetric Yang-Mills by using quadruple and triple cuts in $`42ϵ`$ dimensions. This is advantageous for at least three reasons. First of all, working with multiple cuts simplifies considerably the algebra, because several on-shell conditions can be used at the same time; furthermore, for the case of quadruple cuts the integration is actually completely frozen so that the coefficient of the relevant box functions entering the amplitude can be calculated without performing any integration at all. Lastly, the tree-level sub-amplitudes which are sewn together in order to form the multiple cut of the amplitude are simpler than those entering the two-particle cuts of the same amplitude. It seems clear that immediate further progress with this approach will not require major new conceptual advances, and that it will be directly applicable to more complicated and currently unknown amplitudes.
We describe this method in some detail in Section 2, and then move on to present various examples of its application. Specifically, using generalised unitarity in $`42ϵ`$ dimensions we will re-calculate the all-orders in $`ϵ`$ expressions of all one-loop, four gluon scattering amplitudes in non-supersymmetric Yang-Mills, that is $`+`$$`+`$$`+`$$`+`$, $``$$`+`$$`+`$$`+`$, and the two MHV amplitudes $``$$``$$`+`$$`+`$ and $``$$`+`$$``$$`+`$; and finally, the five-gluon all-plus helicity amplitude $`+`$$`+`$$`+`$$`+`$$`+`$. These amplitudes have already been computed to all orders in $`ϵ`$ in , and we find in all cases complete agreement with the results of that paper. The examples we consider are complementary, as they show that this method can be applied to finite amplitudes without infrared divergences, as well as to infrared divergent amplitudes containing both rational and cut-constructible terms. These calculations are described in Section 3 and Section 4. In an Appendix we have collected some useful definitions and formulae.
## 2 Generalised Unitarity in $`D=42ϵ`$ Dimensions
Conventional unitarity and generalised unitarity in four dimensions have been shown to be extremely powerful tools for calculating one-loop and higher-loop scattering amplitudes in supersymmetric gauge theories and gravity. At one-loop, conventional unitarity amounts to reconstructing the full amplitude from the knowledge of the discontinuity or imaginary part of the amplitude. In this process the amplitude is cut into two tree-level, on-shell amplitudes defined in four dimensions, and the two propagators connecting the two sub-amplitudes are replaced by on-shell delta-functions which reduce the loop integration to a phase space integration. In principle this cutting technique is only sensitive to terms in the amplitude that have discontinuities, like logarithms and polylogarithms, and in general any cut-free, rational terms are lost. However, in supersymmetric theories all rational terms turn out to be uniquely linked to terms with discontinuities, and therefore the full amplitudes can be reconstructed in this fashion .
Furthermore, in supersymmetric theories the one-loop amplitudes are known to be linear combinations of scalar box functions, linear triangle functions and linear bubble functions, with the coefficients being rational functions in spinor products. So the task is really to find an efficient way to fix those coefficients with as few manipulations and/or integrations as possible.
The method based on conventional unitarity introduced by BDDK in does not evaluate the phase space integrals explicitly (from which the full amplitude would be obtained by performing a dispersion integral), rather it reconstructs the loop integrand from which one is able to read off the coefficients of the various integral functions. In practice this means that for a given momentum channel the integrand (which is a product of two tree amplitudes) is simplified as much as possible using the condition that the two internal lines are on-shell, and only in the last step the two delta-functions are replaced by the appropriate propagators which turn the integral from a phase space integral back to a fully-fledged loop integral. The resulting integral function will have the correct discontinuities in the particular channel, but, in general, it will also have additional discontinuities in other channels. Nevertheless, working channel by channel one can extract linear equations for the coefficients which allow us in the end to determine the complete amplitude. However, because of the problem of the additional, unwanted discontinuities, this does not provide a diagrammatic method, i.e. one cannot just sum the various integrals for each channel since different discontinuities might be counted with different weights.
It is natural to contemplate if there exist other complementary, or more efficient methods to extract the above mentioned rational coefficients of the various integral functions, and if in particular we can replace more than two propagators by delta functions, so that the loop integration is further restricted - or even completely localised. The procedure of replacing several internal propagators by $`\delta ^{(+)}`$-functions is well known from the study of singularities and discontinuities of Feynman integrals, and goes under the name of generalised unitarity . What turns generalised unitarity into a powerful tool is the fact that generalised cuts of amplitudes can be evaluated with less effort than conventional two-particle cuts.
The most dramatic simplification arises from using quadruple cuts in one-loop amplitudes in $`𝒩=4`$ SYM. In this case it is known that the one-loop amplitudes are simply given by a sum of scalar box functions without triangles or bubbles . Each quadruple cut singles out a unique box function, and because of the presence of the four $`\delta ^{(+)}`$-functions the loop integration is completely frozen; hence, the coefficient of this particular box is simply given by the product of four tree-level scattering amplitudes . An important subtlety arises here because quadruple cuts do not have solutions in real Minkowski space; therefore at intermediate steps one has to work with complexified momenta.
At this point we can push the analogy with the “reconstruction of the Feynman integrand” a step further. Using the on-shell conditions we can pull out the prefactor which is just the product of four tree-level amplitudes in front of the integral, and the integrand of the remaining loop integral becomes just a product of four $`\delta ^{(+)}`$-functions. If we now promote the integral to a Feynman integral by replacing all $`\delta ^{(+)}`$-functions by the corresponding propagators<sup>3</sup><sup>3</sup>3We thank David Kosower for discussions on this point. we arrive at the integral representation of the appropriate box function. Note that no overcounting issue arises, because each quadruple cut selects a unique box function, and the final result is obtained by summing over all quadruple cuts. In some sense, one can really think of this as a true diagrammatic prescription.
As we reduce the amount of supersymmetry to $`𝒩=1`$, life becomes a bit more complicated, since the one-loop amplitudes are linear combinations of scalar box, triangle and bubble integral functions. No ambiguities related to rational terms occur however, thanks to supersymmetry. It is therefore natural to attack the problem in two steps: First, use quadruple cuts to fix all the box coefficients as described in the previous paragraph. Second, use triple cuts to fix triangle and bubble coefficients. Note that the triple cuts also have contributions from the box functions which have been determined in the first step. The three $`\delta ^{(+)}`$-functions are not sufficient to freeze the loop integration completely, and it is advantageous to use again the “reconstruction of the Feynman integrand” method, i.e. use the on-shell conditions to simplify the integrand as much as possible, and lift the integral to a full loop integral by reinstating three propagators. The resulting integrand can be written as a sum of (integrands of) scalar boxes, triangles and bubbles, after standard reduction techniques, like Passarino-Veltman, have been employed.
At this point it is useful to distinguish three types of triple cuts according to the number of external lines attached to each of the three tree-level amplitudes. If $`p`$ of the three amplitudes have more than one external line attached, we call the cut a $`p`$-mass triple cut. Let us start with the 3-mass triple cut. The box terms can be dropped as they have been determined using quadruple cuts, the coefficients of three-mass triangles can be read off directly, and the remaining terms, which are bubbles or triangles with a different triple cut, are dropped as well. Special care is needed for 1-mass and 2-mass triple cuts. First let us note that any bubble can be written as a linear combination of scalar and linear 1-mass triangles or scalar and linear 2-mass triangles depending on whether the bubble depends on a two-particle invariant, $`t_i^{[2]}=(p_i+p_{i+1})^2`$, or on a $`r`$-particle invariant, $`t_i^{[r]}=(p_i+\mathrm{}+p_{i+r1})^2`$, with $`r>2`$. Therefore, what we want to argue is that two-particle cuts are not needed and that 1-mass, 2-mass and bubbles can be determined from the 1-mass and 2-mass triple cuts. Now every 1-mass triple cut is in one-to-one correspondence with a unique two-particle channel $`t_i^{[2]}=(p_i+p_{i+1})^2`$ and allows us to extract the coefficients of 1-mass triangles and bubbles by only keeping terms in the integral depending on that particular $`t_i^{[2]}`$ and dropping all boxes and triangles/bubbles not depending on that particular variable. The 2-mass triple cut is associated with two momentum invariants, say $`P^2`$ and $`Q^2`$, and we only keep 2-mass triangles and bubbles that depend on those two invariants.
In non-supersymmetric theories we have to face the problem that the amplitudes contain additional rational terms that are not linked to terms with discontinuities. This statement is true if we only keep terms in the amplitude up to $`𝒪(ϵ^0)`$. If we work however in $`D=42ϵ`$ dimensions and keep higher orders in $`ϵ`$, even rational terms $`R`$ develop discontinuities of the form $`R(s)^ϵ=Rϵ\mathrm{log}(s)R+𝒪(ϵ^2)`$ and become cut-constructible<sup>4</sup><sup>4</sup>4The idea of using unitarity in $`D=42ϵ`$ dimensions goes back to , and was used in . . In practice, this means that, in our procedure, whenever we cut internal lines by replacing propagators by $`\delta ^{(+)}`$-functions we have to keep the cut lines in $`D`$ dimensions, and in order to proceed we need to know tree-amplitudes with two legs continued to $`D`$ dimensions. Because of the supersymmetric decomposition of one-loop amplitudes in pure Yang-Mills, which was reviewed in the Introduction, we only need to consider the case of a scalar running in the loop. Furthermore, the massless scalar in $`D`$ dimensions can be thought of as a massive scalar in four dimensions $`L^2=l_{(4)}^2+l_{(2ϵ)}^2=l_{(4)}^2\mu ^2`$ whose mass has to be integrated over . Interestingly, a term in the loop integral with the insertion of “mass” term $`(\mu ^2)^m`$ can be mapped to a higher-dimensional loop integral in $`4+2m2ϵ`$ dimensions with a massless scalar . Some of the required tree amplitudes with two massive scalars and all positive helicity gluons have been calculated in using Feynman diagrams and recursive techniques, and more recently all amplitudes with up to four arbitrary helicity gluons and two massive scalars have been presented in .
The comments in the last paragraph make it clear that generalised unitarity techniques can readily be generalised to $`D`$ dimensions and be used to obtain complete amplitudes in pure Yang-Mills and, more generally, in massless, non-supersymmetric gauge theories. The integrands produced by the method described for four dimensional unitarity will now contain terms multiplied by $`(\mu ^2)^m`$ and, therefore, the set of integral functions appearing in the amplitudes includes, in addition to the four-dimensional functions, also higher-dimensional box, triangle and bubble functions (some explicit examples of higher-dimensional integral functions can be found in Appendix A). For example the one-loop $`+`$$`+`$$`+`$$`+`$ gluon amplitude, which vanishes in SYM, is given by a rational function times a box integral with $`\mu ^4`$ inserted, $`I_4[\mu ^4]=(ϵ)(1ϵ)I_4^{82ϵ}=1/6+𝒪(ϵ)`$. Hence this amplitude is a purely rational function in spinor variables.
In the following sections we will describe in detail how this procedure is applied in practice by recalculating all four-gluon scattering amplitudes and the positive helicity five-gluon scattering amplitude in pure Yang-Mills at one-loop level. These examples include the cases of infrared finite amplitudes that are purely rational (and their supersymmetric counterparts vanish), and infrared divergent amplitudes that contain both rational and cut-constructible terms.
## 3 Four-point amplitudes in pure Yang-Mills
In this section we recalculate all the known four-gluon scattering amplitudes, that is $`+`$$`+`$$`+`$$`+`$, $``$$`+`$$`+`$$`+`$, $``$$``$$`+`$$`+`$, and finally $``$$`+`$$``$$`+`$, from quadruple and triple cuts.
### 3.1 The one-loop $`+`$$`+`$$`+`$$`+`$ amplitude
The one-loop $`+`$$`+`$$`+`$$`+`$ amplitude with a complex scalar running in the loop is the simplest of the all-plus gluon amplitudes, and was first derived in using the string-inspired formalism.
The expression in $`42ϵ`$ dimensions, valid to all-orders in $`ϵ`$, is computed in and is given by
$$𝒜_4^{\mathrm{scalar}}(1^+,2^+,3^+,4^+)=\frac{2i}{(4\pi )^{2ϵ}}\frac{[12][34]}{1234}K_4,$$
(3.1)
where<sup>5</sup><sup>5</sup>5Notice also that $`[12][34]/(1234)=s_{12}s_{23}/(12233441)`$.
$$K_4:=I_4[\mu ^4]=ϵ(1ϵ)I_4^{D=82ϵ}=\frac{1}{6}+𝒪(ϵ).$$
(3.2)
In this paper we closely follow the conventions of , with
$$I_n^{D=42ϵ}[f(p,\mu ^2)]:=i()^{n+1}(4\pi )^{2ϵ}\frac{d^4l}{(2\pi )^4}\frac{d^{2ϵ}\mu }{(2\pi )^{2ϵ}}\frac{f(l,\mu ^2)}{(l^2\mu ^2)\mathrm{}[(l_{i=i}^{n1}K_i)^2\mu ^2]},$$
(3.3)
where $`K_i`$ are external momenta (which, in colour-ordered amplitudes, are sums of adjacent null momenta of the external gluons) and $`f(l,\mu ^2)`$ is a generic function of the four-dimensional loop momentum $`l`$ and of $`\mu ^2`$.
The amplitude with four positive helicity gluons is part of the infinite sequence of all-plus helicity gluons, for which a closed expression was conjectured in . The result for all $`n`$ is given by
$$𝒜_n(+,\mathrm{},+)=\frac{i}{48\pi ^2}\underset{1i_1<i_2<i_3<i_4n}{}\frac{i_1i_2[i_2i_3]i_3i_4[i_4i_1]}{1223\mathrm{}n1},$$
(3.4)
or, alternatively,
$$𝒜_n=\frac{i}{96\pi ^2}\underset{1i_1<i_2<i_3<i_4n}{}\frac{s_{i_1i_2}s_{i_3i_4}s_{i_1i_3}s_{i_2i_4}+s_{i_1i_4}s_{i_2i_3}\mathrm{\hspace{0.17em}4}iϵ(i_1i_2i_3i_4)}{1223\mathrm{}n1},$$
(3.5)
where $`ϵ(abcd):=ϵ_{\mu \nu \rho s}a^\mu b^\nu c^\rho d^\sigma `$. As $`ϵ0`$, (3.1) becomes
$$𝒜_4=\frac{i}{48\pi ^2}\frac{s_{12}s_{23}}{12233441}.$$
(3.6)
We see that this amplitude (3.1) consists of purely rational terms, which are cut-free in four dimensions. We now show how to derive (3.1) from quadruple cuts in $`D=42ϵ`$ dimensions.
Consider the quadruple-cut diagram in Figure 1, which is obtained by sewing four three-point scattering amplitudes<sup>6</sup><sup>6</sup>6In the following for the purpose of calculating the (generalised) cuts we drop factors of $`i`$ appearing in the usual definition of tree-amplitudes and propagators. For quadruple and two-particle cuts this does not affect the final result, while for triple cuts this introduces an extra $`(1)`$ factor which we reinstate at the end of every calculation. with one massless gluon and two massive scalars of mass $`\mu ^2`$. From we take the three-point amplitudes for one positive-helicity gluon and two scalars:
$$𝒜(l_1^+,k^+,l_2^{})=𝒜(l_1^{},k^+,l_2^+)=\frac{q|l_1|k]}{qk},$$
(3.7)
where $`l_1+l_2+k=0`$. Here $`|q`$ is an arbitrary reference spinor not proportional to $`|k`$. It is easy to see that (3.7) is actually independent of the choice of $`|q`$.
The $`D`$-dimensional quadruple cut of the amplitude $`+`$$`+`$$`+`$$`+`$ is obtained by combining four three-point tree-level amplitudes,
$$\frac{q_1|l_1|1]}{q_11}\frac{q_2|l_2|2]}{q_22}\frac{q_3|l_3|3]}{q_33}\frac{q_4|l_4|4]}{q_44}.$$
(3.8)
The reference momenta $`q_i,i=1,\mathrm{},4`$ in each of the four ratios in this expression may be chosen arbitrarily. Then, using momentum conservation,
$$l_2=l_1k_2,l_4=l_3k_4,$$
(3.9)
the fact that the external momenta are null, and that the internal momenta square to $`\mu ^2`$, it is easy to see that
$$\frac{q_1|l_1|1]}{q_11}\frac{q_2|l_2|2]}{q_22}=\mu ^2\frac{[12]}{12},$$
(3.10)
and similarly
$$\frac{q_3|l_3|3]}{q_33}\frac{q_4|l_4|4]}{q_44}=\mu ^2\frac{[34]}{34},$$
(3.11)
so that the above expression (3.8) becomes simply
$$\mu ^4\frac{[12][34]}{1234}.$$
(3.12)
Finally, we lift the quadruple-cut box to a box function by reinstating the appropriate Feynman propagators. These propagators then combine with the additional factor of $`\mu ^4`$ in (3.12) to yield the factor $`iK_4/(4\pi )^{2ϵ}`$ which is proportional to the scalar box integral defined in (3.2). Including an additional factor of 2 due to the fact that there is a complex scalar propagating in the loop, the amplitude (3.1) is reproduced correctly.
Next we inspect three-particle cuts. One of the three tree-level amplitudes we sew in the triple-cut amplitude is an amplitude with two positive-helicity gluons and two scalars
$$𝒜(l_1^+,1^+,2^+,l_2^{})=\mu ^2\frac{[12]}{12[(l_1+k_1)^2\mu ^2]}.$$
(3.13)
Consider, for example, the three-particle cut defined by $`1^+,2^+,(3^+,4^+)`$, see Figure 2. Using (3.7) and (3.13), the product of the three tree-level amplitudes gives
$$\frac{q_1|l_1|1]}{q_11}\frac{q_2|l_1|2]}{q_22}\frac{\mu ^2[34]}{34[(l_2k_3)^2\mu ^2]},$$
(3.14)
with $`l_2=l_1k_2`$. As for the quadruple cut, it is easily seen that, on this triple cut,
$$\frac{q_1|l_1|1]}{q_11}\frac{q_2|l_1|2]}{q_22}=\mu ^2\frac{[12]}{12},$$
(3.15)
where we used $`l_1^2=l_2^2=l_4^2=\mu ^2`$. The triple-cut integrand then becomes
$$\frac{[12][34]}{1234}\frac{\mu ^4}{[(l_2k_3)^2\mu ^2]},$$
(3.16)
which, after replacing the three $`\delta ^{(+)}`$ functions by propagators, integrates to (3.1), where we have included an additional $`(1)`$ factor following the comments in footnote 6. The factor of $`2`$ in (3.1) comes from summing over the two “scalar helicities”. The same result comes from evaluating the remaining triple cuts.
We remark that in the case of the quadruple cut we did not even need to insert the solutions of the on-shell conditions for the loop momenta into the expression coming from the cut. This is not true in general; for example, for the five gluon amplitude discussed below the sum over solutions will be essential to obtaining the correct amplitude.
### 3.2 The one-loop $``$$`+`$$`+`$$`+`$ amplitude
The one-loop four gluon scattering amplitude $``$$`+`$$`+`$$`+`$, with a complex scalar running in the loop, is given to all orders in $`ϵ`$ by
$`𝒜_4^{\mathrm{scalar}}(1^{},2^+,3^+,4^+)`$ $`=`$ $`{\displaystyle \frac{2i}{(4\pi )^{2ϵ}}}{\displaystyle \frac{[24]^2}{[12]2334[41]}}{\displaystyle \frac{st}{u}}[{\displaystyle \frac{t(us)}{su}}J_3(s)+{\displaystyle \frac{s(ut)}{tu}}J_3(t)`$
$``$ $`{\displaystyle \frac{tu}{s^2}}J_2(s){\displaystyle \frac{su}{t^2}}J_2(t)+{\displaystyle \frac{st}{2u}}J_4+K_4].`$
We will now show how to derive this result using generalised unitarity cuts.
First consider the quadruple cut (see Figure 3).
The product of tree amplitudes gives
$$\frac{1|l_1|q_1]}{[1q_1]}\frac{q_2|l_2|2]}{q_22}\frac{q_3|l_3|3]}{q_33}\frac{q_4|l_4|4]}{q_44}.$$
(3.18)
It is straightforward to show that, on the quadruple cut,
$`{\displaystyle \frac{q_3|l_3|3]}{q_33}}{\displaystyle \frac{q_4|l_4|4]}{q_44}}`$ $`=`$ $`\mu ^2{\displaystyle \frac{[34]}{34}},`$
$`{\displaystyle \frac{1|l_1|q_1]}{1q_1}}{\displaystyle \frac{q_2|l_2|2]}{q_22}}`$ $`=`$ $`{\displaystyle \frac{[23]}{[31]}}(\mu ^2{\displaystyle \frac{31}{23}}[2|l_1|1),`$
and hence the quadruple cut in Figure 3 gives
$$Q(1^+,2^+,3^+,4^{})=\mu ^2\frac{[\mathrm{3\hspace{0.17em}4}]}{\mathrm{3\hspace{0.17em}4}}\frac{[\mathrm{2\hspace{0.17em}3}]}{[\mathrm{3\hspace{0.17em}1}]}[\mu ^2\frac{\mathrm{3\hspace{0.17em}1}}{\mathrm{2\hspace{0.17em}3}}+[2|l_1|1].$$
(3.19)
In order to compare with (3.2) it is useful to notice that
$$\frac{[34]}{34}\frac{[23]}{[31]}\frac{31}{23}=\frac{[24]^2}{[12]2334[41]}\frac{st}{u}:=𝒩.$$
(3.20)
We conclude that the first term in (3.19) generates
$$\frac{i}{(4\pi )^{2ϵ}}\left(\frac{[24]^2}{[12]2334[41]}\frac{st}{u}\right)K_4,$$
(3.21)
where the prefactor in (3.21) comes from the definition (3.2) and (3.3) for the function $`K_4`$.
The second term in (3.19) corresponds to a linear box integral, which we examine now. We notice that the quadruple cut freezes the loop integration on the solution for the cut. In the linear box term in (3.19) we will then replace $`l_1`$ in $`[2|l_1|1`$ by the solutions of the cut, and sum over the different solutions.
Specifically, in order to solve for the cut-loop momentum $`l_1`$ one has to require
$`l_1^2=l_2^2=l_3^2=l_4^2=\mu ^2,`$
$`l_1=l_4k_1,l_2=l_1k_2,l_3=l_2k_3,l_4=l_3k_4.`$ (3.22)
In order to solve these conditions, it proves useful to use the four linearly independent vectors $`k_1,k_2,k_3`$ and $`K`$, where
$$K_\mu :=ϵ_{\mu \nu \rho \sigma }k_1^\nu k_2^\rho k_3^\sigma .$$
(3.23)
Setting
$$l_1=ak_1+bk_2+ck_3+dK,$$
(3.24)
one finds
$`a={\displaystyle \frac{t}{2u}},b={\displaystyle \frac{1}{2}},c={\displaystyle \frac{s}{2u}},`$ (3.25)
$`d=\pm \sqrt{{\displaystyle \frac{st+4\mu ^2u}{stu^2}}},`$
where
$$s=(k_1+k_2)^2,t=(k_2+k_3)^2,u=(k_1+k_3)^2,$$
(3.26)
and $`s+t+u=0`$.
Then one has
$$[2|l_1|1[2|\frac{l_1^++l_1^{}}{2}|1=c[2|3|1=\frac{s}{2u}[23]31,$$
(3.27)
where $`l_1^\pm `$ denotes the two solutions for the quadruple cut. The square root drops out of the calculation (as it should, given that the amplitude is a rational function). We conclude that the second term in (3.19) gives<sup>7</sup><sup>7</sup>7Recall that in our conventions $`t:=23[32]`$.
$$\frac{i}{(4\pi )^{2ϵ}}\left(\frac{[24]^2}{[12]2334[41]}\frac{st}{u}\right)\frac{st}{2u}J_4,$$
(3.28)
where
$$J_n:=I_n[\mu ^2].$$
(3.29)
Again, the prefactor in (3.28) arises from the definition (3.3).
In total the quadruple cut (3.19) gives
$$\frac{2i}{(4\pi )^{2ϵ}}𝒩\left(K_4+\frac{st}{2u}J_4\right),$$
(3.30)
where we have again included a factor of two for the contribution of a complex scalar. This result matches exactly all the box functions appearing in (3.2).
We now move on to consider triple cuts. We start by considering the triple cut in Figure 4a, which we label as $`(1^{},2^+,(3^+,4^+))`$. It may be shown that this triple cut yields the following expression:
$`TC(1^{},2^+,(3^+,4^+))`$ $`=`$ $`\mu ^2{\displaystyle \frac{[\mathrm{3\hspace{0.17em}4}]}{\mathrm{3\hspace{0.17em}4}}}{\displaystyle \frac{[\mathrm{2\hspace{0.17em}3}]}{[\mathrm{3\hspace{0.17em}1}]}}(\mu ^2{\displaystyle \frac{\mathrm{3\hspace{0.17em}1}}{\mathrm{2\hspace{0.17em}3}}}[2|l_1|1){\displaystyle \frac{1}{(l_2k_3)^2\mu ^2}}`$
$`\mu ^2{\displaystyle \frac{[\mathrm{3\hspace{0.17em}4}]}{\mathrm{3\hspace{0.17em}4}}}{\displaystyle \frac{[2|l_1|1}{\mathrm{2\hspace{0.17em}3}[\mathrm{3\hspace{0.17em}1}]}}.`$ (3.31)
The first line in (3.2) clearly contains the (negative of the) term already studied with quadruple cuts – see (3.19) (for an explanation of the relative minus sign see footnote 6). We now reconsider the linear box term (second term in the first line of (3.2)), and study its Passarino-Veltman (PV) reduction. As we shall see, this box appears also in other triple cuts (see (3.2)).
Let us consider the linear box integral
$$A^\mu :=\frac{d^4l_1}{(2\pi )^4}\frac{d^{2ϵ}\mu }{(2\pi )^{2ϵ}}\frac{\mu ^2l_1^\mu }{(l_1^2\mu ^2)[(l_1k_2)^2\mu ^2][(l_1k_2k_3)^2\mu ^2][(l_1+k_1)^2\mu ^2]}.$$
(3.32)
On general grounds the integral is a linear combination of three of the external momenta,
$$A^\mu =\alpha k_1^\mu +\beta k_2^\mu +\gamma k_3^\mu .$$
(3.33)
For the coefficients we find
$`\alpha `$ $`=`$ $`{\displaystyle \frac{i}{(4\pi )^{2ϵ}}}{\displaystyle \frac{1}{2u}}\left[tJ_4\mathrm{\hspace{0.17em}2}J_3(s)+\mathrm{\hspace{0.17em}2}J_3(t)\right],`$ (3.34)
$`\beta `$ $`=`$ $`{\displaystyle \frac{i}{(4\pi )^{2ϵ}}}{\displaystyle \frac{1}{2}}J_4,`$
$`\gamma `$ $`=`$ $`{\displaystyle \frac{i}{(4\pi )^{2ϵ}}}{\displaystyle \frac{1}{2u}}\left[sJ_4\mathrm{\hspace{0.17em}2}J_3(s)+\mathrm{\hspace{0.17em}2}J_3(t)\right].`$
Taken literally, this means that from the linear box in (3.2) we not only get the $`J_4`$ function but, altogether:
$$\frac{i𝒩}{(4\pi )^{2ϵ}}\left(\frac{st}{2u}J_4\frac{t}{u}J_3(s)+\frac{t}{u}J_3(t)\right).$$
(3.35)
Summarising, the PV reduction of the first line of the triple cut (3.2), lifted to a Feynman integral, gives:
$$\frac{i𝒩}{(4\pi )^{2ϵ}}\left(K_4+\frac{st}{2u}J_4\frac{t}{u}J_3(s)+\frac{t}{u}J_3(t)\right).$$
(3.36)
The last term in (3.36) is clearly spurious – it does not have the right triple cut, and has appeared because we lifted the cut-integral to a Feynman integral; hence we will drop it. In conclusion, the triple cut $`(1^{},2^+,(3^+,4^+))`$ in Figure 4a leads to
$$\frac{i𝒩}{(4\pi )^{2ϵ}}\left(K_4+\frac{st}{2u}J_4\frac{t}{u}J_3(s)\right).$$
(3.37)
We now consider the last term in (3.2), which generates a linear triangle, whose PV reduction we consider now. The linear triangle is proportional to
$$B^\mu :=\frac{d^4l_1}{(2\pi )^4}\frac{d^{2ϵ}\mu }{(2\pi )^{2ϵ}}\frac{\mu ^2l_1^\mu }{(l_1^2\mu ^2)[(l_1k_2)^2\mu ^2][(l_1+k_1)^2\mu ^2]}.$$
(3.38)
On general grounds,
$$B^\mu =\theta k_1^\mu +\tau k_2^\mu ,$$
(3.39)
and hence
$$[2|B|1=0.$$
(3.40)
We conclude that the second line in (3.2) gives a vanishing contribution, so that the content of this triple cut is encoded in (3.37).
Next we consider the triple cut labelled by $`((1^{},2^+),3^+,4^+)`$ and represented in Figure 4b, which gives
$`TC((1^{},2^+),3^+,4^+)`$ $`=`$ $`\mu ^2{\displaystyle \frac{[\mathrm{3\hspace{0.17em}4}]}{\mathrm{3\hspace{0.17em}4}}}{\displaystyle \frac{[\mathrm{2\hspace{0.17em}3}]}{[\mathrm{3\hspace{0.17em}1}]}}[\mu ^2{\displaystyle \frac{\mathrm{3\hspace{0.17em}1}}{\mathrm{2\hspace{0.17em}3}}}+{\displaystyle \frac{\mathrm{1\hspace{0.17em}2}}{\mathrm{2\hspace{0.17em}3}}}3|l_2|2]]{\displaystyle \frac{1}{(l_2+k_2)^2\mu ^2}}`$ (3.41)
$`+\mu ^2{\displaystyle \frac{[\mathrm{3\hspace{0.17em}4}]}{\mathrm{3\hspace{0.17em}4}}}{\displaystyle \frac{1|\mathrm{3\hspace{0.17em}1}l_2\mathrm{2\hspace{0.17em}3}l_2|2]}{\mathrm{1\hspace{0.17em}2}[\mathrm{1\hspace{0.17em}2}]\mathrm{2\hspace{0.17em}3}[\mathrm{3\hspace{0.17em}1}]}}.`$
The first term of (3.41) clearly corresponds to the function $`K_4`$ already fixed using quadruple cuts. The second term can be rewritten as follows. Introducing $`l_1:=l_2+k_2`$, we have
$$\frac{12}{23}3|l_2|2]=[2|l_2|1+\frac{13}{23}[(l_2+k_2)^2\mu ^2],$$
(3.42)
therefore we can rewrite (3.41) as
$`TC((1^{},2^+),3^+,4^+)`$ $`=`$ $`\mu ^2{\displaystyle \frac{[\mathrm{3\hspace{0.17em}4}]}{\mathrm{3\hspace{0.17em}4}}}{\displaystyle \frac{[\mathrm{2\hspace{0.17em}3}]}{[\mathrm{3\hspace{0.17em}1}]}}(\mu ^2{\displaystyle \frac{\mathrm{3\hspace{0.17em}1}}{\mathrm{2\hspace{0.17em}3}}}[2|l_1|1){\displaystyle \frac{1}{(l_2+k_2)^2\mu ^2}}`$
$`+`$ $`\mu ^2{\displaystyle \frac{[\mathrm{3\hspace{0.17em}4}]}{\mathrm{3\hspace{0.17em}4}}}\left({\displaystyle \frac{1|\mathrm{3\hspace{0.17em}1}l_2\mathrm{2\hspace{0.17em}3}l_2|2]}{\mathrm{1\hspace{0.17em}2}[\mathrm{1\hspace{0.17em}2}]\mathrm{2\hspace{0.17em}3}[\mathrm{3\hspace{0.17em}1}]}}{\displaystyle \frac{[23]}{[31]}}{\displaystyle \frac{31}{23}}\right).`$ (3.43)
We know already that the PV reduction of the first line of (3.2) corresponds to (3.36) – with the term containing $`J_3(t)`$ removed – so we now study the second line, which will give new contributions.
The second term in the second line corresponds to a scalar triangle, more precisely it gives a contribution
$$\frac{i𝒩}{(4\pi )^{2ϵ}}J_3(s).$$
(3.44)
The first term corresponds to a linear triangle, and now we perform its PV reduction. The relevant integral is
$$C^\mu :=\frac{d^4l_2}{(2\pi )^4}\frac{d^{2ϵ}\mu }{(2\pi )^{2ϵ}}\frac{\mu ^2l_2^\mu }{(l_2^2\mu ^2)[(l_2k_3)^2\mu ^2][(l_2+k_1+k_2)^2\mu ^2]}.$$
(3.45)
On general grounds,
$$C^\mu =\lambda k_3^\mu +\kappa (k_1+k_2)^\mu .$$
(3.46)
A quick calculation shows that
$$\lambda =\frac{i}{(4\pi )^{2ϵ}}\left[J_3(s)\frac{2}{s}J_2(s)\right],\kappa =\frac{i}{(4\pi )^{2ϵ}}\frac{1}{s}J_2(s).$$
(3.47)
The first term in the second line of (3.2) gives then
$$\frac{i𝒩}{(4\pi )^{2ϵ}}\left(\frac{u}{s}J_3(s)+\frac{ut}{s}J_2(s)\right),$$
(3.48)
where $`𝒩`$ is defined in (3.20). Altogether, the second line of (3.2) gives
$$\frac{i𝒩}{(4\pi )^{2ϵ}}\left(\left(1+\frac{u}{s}\right)J_3(s)+\frac{ut}{s}J_2(s)\right),$$
(3.49)
whereas from the first line of the same equation we get
$$\frac{i𝒩}{(4\pi )^{2ϵ}}\left(K_4+\frac{st}{2u}J_4\frac{t}{u}J_3(s)\right),$$
(3.50)
where we have dropped the term $`J_3(t)`$ for reasons explained earlier.
We conclude that the function which incorporates all the right cuts in the channels considered so far is equal to the sum of (3.49) and (3.50), which gives
$$\frac{i𝒩}{(4\pi )^{2ϵ}}\left(K_4+\frac{st}{2u}J_4\frac{t}{u}J_3(s)\left(1+\frac{u}{s}\right)J_3(s)+\frac{ut}{s^2}J_2(s)\right).$$
(3.51)
Using $`t/u1u/s=s/uu/s`$, (3.51) becomes
$$\frac{i𝒩}{(4\pi )^{2ϵ}}\left(K_4+\frac{st}{2u}J_4+\left(\frac{s}{u}\frac{u}{s}\right)J_3(s)+\frac{ut}{s^2}J_2(s)\right).$$
(3.52)
To finish the calculation one has to consider the two remaining triple cuts, that is $`(4^+,1^{},(2^+,3^+))`$ and $`((4^+,1^{}),2^+,3^+)`$. These cuts can be obtained from the previously considered cuts by exchanging $`s`$ with $`t`$.
Our conclusion is therefore that the function (including the usual factor of 2) with the correct quadruple and triple cuts is:
$`{\displaystyle \frac{2i𝒩}{(4\pi )^{2ϵ}}}(K_4+{\displaystyle \frac{st}{2u}}J_4+({\displaystyle \frac{s}{u}}{\displaystyle \frac{u}{s}})J_3(s)+{\displaystyle \frac{ut}{s^2}}J_2(s)`$ (3.53)
$`+({\displaystyle \frac{t}{u}}{\displaystyle \frac{u}{t}})J_3(t)+{\displaystyle \frac{us}{t^2}}J_2(t)).`$
This agrees precisely with (3.1) using the identities
$$\frac{t(us)}{su}=\frac{s}{u}\frac{u}{s},\frac{s(ut)}{tu}=\frac{t}{u}\frac{u}{t}.$$
(3.54)
### 3.3 The one-loop $``$$``$$`+`$$`+`$ amplitude
We now turn our attention to the one-loop four point amplitudes with two negative helicity gluons. We start by considering the one-loop amplitude $`𝒜_4^{\mathrm{scalar}}(1^{},2^{},3^+,4^+)`$, which is given by <sup>8</sup><sup>8</sup>8Here for simplicity we drop the functions $`I_1`$ and $`I_2(0)`$, which are zero in the massless case . We also include a factor of two as we are considering complex scalars.
$$𝒜_4^{\mathrm{scalar}}(1^{},2^{},3^+,4^+)=2\frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}\left(\frac{t}{s}K_4+\frac{1}{s}J_2(t)+\frac{1}{t}I_2^{62ϵ}(t)\right).$$
(3.55)
To begin with, we consider the quadruple cut of the amplitude, represented in Figure 5. It is given by
$$\frac{1|l_1|q_1]}{[1q_1]}\frac{2|l_1|q_2]}{[2q_2]}\frac{q_3|l_3|3]}{q_3\mathrm{\hspace{0.17em}3}}\frac{q_4|l_4|4]}{q_4\mathrm{\hspace{0.17em}4}}.$$
(3.56)
By choosing $`q_1=2`$, $`q_2=1`$, $`q_3=4`$, $`q_4=3`$, (3.56) can be rewritten as
$$i\frac{t}{s}𝒜_4^{\mathrm{tree}}\mu ^4,$$
(3.57)
where
$$𝒜_4^{\mathrm{tree}}=i\frac{\mathrm{1\hspace{0.17em}2}^3}{\mathrm{2\hspace{0.17em}3}\mathrm{3\hspace{0.17em}4}\mathrm{4\hspace{0.17em}1}}.$$
(3.58)
Reinstating the four cut propagators and integrating over the loop momentum, (3.57) gives
$$\frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}\left(\frac{t}{s}K_4\right),$$
(3.59)
where $`K_4`$ is defined in (3.2).
Next we consider triple cuts. We begin our analysis with the triple cut in Figure 6a. This yields
$$\frac{\mu ^2[\mathrm{3\hspace{0.17em}4}]}{\mathrm{3\hspace{0.17em}4}2(l_23)}\frac{1|l_1|q_1]}{[1q_1]}\frac{2|l_1|q_2]}{[2q_2]}=\mu ^4\frac{\mathrm{1\hspace{0.17em}2}[\mathrm{3\hspace{0.17em}4}]}{[\mathrm{1\hspace{0.17em}2}]\mathrm{3\hspace{0.17em}4}}\frac{1}{2(l_23)},$$
(3.60)
which, upon reinstating the cut propagators and performing the loop momentum integration gives
$$\frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}\left(\frac{t}{s}K_4\right).$$
(3.61)
This function had already been detected with the quadruple cut, as discussed earlier.
Next we move on to consider the triple cut in Figure 6b. This yields
$$\frac{1|l_3|4]^2}{2t(l_34)}\frac{2|l_1|q_1]}{[2q_1]}\frac{q_2|l_2|3]}{q_2\mathrm{\hspace{0.17em}3}}.$$
(3.62)
We can re-cast (3.62) as follows. Firstly, we write
$$\frac{1|l_3|4]q_2|l_3|3]}{q_2\mathrm{\hspace{0.17em}3}}=\mu ^2\frac{1|4|3]}{\mathrm{3\hspace{0.17em}4}}\frac{2(l_34)1|l_3|3]}{\mathrm{3\hspace{0.17em}4}},$$
(3.63)
and secondly
$$\frac{1|l_3|4]2|l_1|q_1]}{[2q_1]}=\mu ^2\frac{2|1|4]}{[\mathrm{1\hspace{0.17em}2}]}\frac{2(l_34)2|1|4]}{[\mathrm{1\hspace{0.17em}2}]}+\frac{2(l_34)2|l_3|4]}{[\mathrm{1\hspace{0.17em}2}]}.$$
(3.64)
The expression (3.62) becomes a sum of six terms $`T_i`$, $`i=1,\mathrm{},6`$, where
$`T_1`$ $`=`$ $`{\displaystyle \frac{1|4|3]2|1|4]\mu ^4}{t\mathrm{3\hspace{0.17em}4}[\mathrm{1\hspace{0.17em}2}]2(l_34)}},`$
$`T_2`$ $`=`$ $`{\displaystyle \frac{1|4|3]2|1|4]\mu ^2}{t\mathrm{3\hspace{0.17em}4}[\mathrm{1\hspace{0.17em}2}]}},`$
$`T_3`$ $`=`$ $`{\displaystyle \frac{1|4|3]2|l_3|4]\mu ^2}{t\mathrm{3\hspace{0.17em}4}[\mathrm{1\hspace{0.17em}2}]}},`$
$`T_4`$ $`=`$ $`{\displaystyle \frac{2|1|4]1|l_3|3]\mu ^2}{t\mathrm{3\hspace{0.17em}4}[\mathrm{1\hspace{0.17em}2}]}},`$
$`T_5`$ $`=`$ $`{\displaystyle \frac{2|1|4]1|l_3|3]2(l_34)}{t\mathrm{3\hspace{0.17em}4}[\mathrm{1\hspace{0.17em}2}]}},`$
$`T_6`$ $`=`$ $`{\displaystyle \frac{1|l_3|3]2|l_3|4]2(l_34)}{t\mathrm{3\hspace{0.17em}4}[\mathrm{1\hspace{0.17em}2}]}}.`$ (3.65)
Next we replace the delta functions with propagators, and integrate over the loop momentum. To evaluate the integrals, we use the linear, quadratic and cubic triangle integrals in $`42ϵ`$ dimensions listed in the Appendix. The integration of the expressions gives
$`T_1`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{t}{s}}K_4\right),`$
$`T_2`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{t}{s}}J_3(t)\right),`$
$`T_3`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{t}{s}}J_3(t){\displaystyle \frac{1}{s}}J_2(t)\right),`$
$`T_4`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{1}{s}}J_2(t)\right),`$
$`T_5`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{t}{2s}}I_2(t)+{\displaystyle \frac{u}{s}}I_3^{62ϵ}(t)\right),`$
$`T_6`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{t}{4s}}I_2(t)\left({\displaystyle \frac{3}{2s}}+{\displaystyle \frac{1}{t}}\right)I_2^{62ϵ}(t){\displaystyle \frac{u}{s}}I_3^{62ϵ}(t)\right).`$ (3.66)
We now use (A.26) in relating $`J_2(t)`$ to $`I_2(t)`$ and $`I_2^{62ϵ}(t)`$, and get
$$T_5+T_6𝒜_4^{\mathrm{tree}}\left(\frac{1}{s}J_2(t)\frac{1}{t}I_2^{62ϵ}(t)\right).$$
(3.67)
Adding up the six $`T_i`$ terms, and including the usual factor of two, we obtain
$$\frac{2𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}\left(\frac{t}{s}K_4\frac{1}{s}J_2(t)\frac{1}{t}I_2^{62ϵ}(t)\right),$$
(3.68)
which precisely agrees with (3.55).
### 3.4 The one-loop $``$$`+`$$``$$`+`$ amplitude
Now we consider the one-loop amplitude with a complex scalar in the loop, $`𝒜_4^{\mathrm{scalar}}(1^{},2^+,3^{},4^+)`$, which is given by
$`𝒜_4^{\mathrm{scalar}}(1^{},2^+,3^{},4^+)`$ $`=`$ $`2{\displaystyle \frac{1}{(4\pi )^{2ϵ}}}𝒜_4^{\mathrm{tree}}({\displaystyle \frac{st}{u^2}}K_4{\displaystyle \frac{s^2t^2}{u^3}}I_4^{62ϵ}+{\displaystyle \frac{st}{u^2}}I_3^{62ϵ}(t)`$
$`+`$ $`{\displaystyle \frac{st}{u^2}}I_3^{62ϵ}(s){\displaystyle \frac{st(st)}{u^3}}J_3(t){\displaystyle \frac{st(ts)}{u^3}}J_3(s)+{\displaystyle \frac{s}{u^2}}J_2(t)+{\displaystyle \frac{t}{u^2}}J_2(s)`$
$`+`$ $`{\displaystyle \frac{s}{tu}}I_2^{62ϵ}(t)+{\displaystyle \frac{t}{su}}I_2^{62ϵ}(s)+{\displaystyle \frac{ts^2}{u^3}}I_2(t)+{\displaystyle \frac{st^2}{u^3}}I_2(s)).`$
The relevant quadruple cut is represented in Figure 7, and gives:
$`{\displaystyle \frac{1|l_1|q_1]}{[1q_1]}}{\displaystyle \frac{q_2|l_2|2]}{q_2\mathrm{\hspace{0.17em}2}}}{\displaystyle \frac{3|l_3|q_3]}{[3q_3]}}{\displaystyle \frac{q_4|l_4|q_4]}{q_4\mathrm{\hspace{0.17em}4}}}`$
$`={\displaystyle \frac{1}{[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}}}(\mathrm{1\hspace{0.17em}3}\mu ^2+\mathrm{1\hspace{0.17em}2}3|l_1|2])([\mathrm{2\hspace{0.17em}4}]\mu ^2[\mathrm{3\hspace{0.17em}4}]3|l_1|2])`$
$`=i𝒜_4^{\mathrm{tree}}\left({\displaystyle \frac{st\mu ^4}{u^2}}+{\displaystyle \frac{2s^2t|l_1|2]\mu ^2}{u^23|1|2]}}+{\displaystyle \frac{s^3t3|l_1|2]^2}{u^23|1|2]^2}}\right),`$ (3.70)
where
$$𝒜_4^{\mathrm{tree}}=i\frac{13^4}{12233441}.$$
(3.71)
Averaging over the two solutions of the quadruple cut we obtain the following expression:
$$i𝒜_4^{\mathrm{tree}}\left(\frac{st}{u^2}\mu ^4+\frac{2s^2t^2}{u^3}\mu ^2+\frac{s^3t^3}{2u^4}\right).$$
(3.72)
After reinstating the four cut propagators and integrating over the loop momentum, (3.72) gives
$$\frac{1}{(4\pi )^{2ϵ}}𝒜_4^{\mathrm{tree}}\left(\frac{st}{u^2}K_4\frac{2s^2t^2}{u^3}J_4\frac{s^3t^3}{2u^4}I_4\right).$$
(3.73)
We now use the identity (A.26) in ignoring functions that do not have a quadruple cut to write this as
$$\frac{1}{(4\pi )^{2ϵ}}𝒜_4^{\mathrm{tree}}\left(\frac{st}{u^2}K_4+\frac{s^2t^2}{u^3}I_4^{62ϵ}\right).$$
(3.74)
We now consider triple cuts. There is only one independent triple cut, and we consider, for instance, the triple cut in Figure 8, which gives
$$\frac{1|l_3|4]^2}{2t(l_34)}\frac{3|l_3|q_2]}{[3q_2]}\frac{q_1|l_1|2]}{q_1\mathrm{\hspace{0.17em}2}}.$$
(3.75)
Using straightforward spinor manipulations, and taking into account properties of the cut momenta, one finds that the above expression may be expanded as a product of two sets of terms. The first is
$$\frac{1|l_3|4]3|l_3|q_2]}{[3q_2]}=\frac{\mu ^23|1|4]}{[\mathrm{1\hspace{0.17em}3}]}\frac{t3|l_3|4]}{[\mathrm{1\hspace{0.17em}3}]}+\frac{2(l_34)3|l_3|4]}{[\mathrm{1\hspace{0.17em}3}]},$$
(3.76)
whereas the second is
$$\frac{1|l_3|4]q_1|l_1|2]}{q_1\mathrm{\hspace{0.17em}2}}=\frac{\mu ^21|4|2]}{\mathrm{2\hspace{0.17em}4}}+\frac{4|1|2]1|l_3|4]}{\mathrm{2\hspace{0.17em}4}}\frac{2(l_3.4)1|l_3|2]}{\mathrm{2\hspace{0.17em}4}}.$$
(3.77)
The expression (3.75) becomes then a sum of nine terms $`R_i`$, $`i=1,\mathrm{},9`$, where
$`R_1`$ $`=`$ $`{\displaystyle \frac{1|4|2]3|1|4]\mu ^4}{t[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}2(l_34)}},`$
$`R_2`$ $`=`$ $`{\displaystyle \frac{4|1|2]3|1|4]1|l_3|4]\mu ^2}{t[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}2(l_34)}},`$
$`R_3`$ $`=`$ $`{\displaystyle \frac{3|1|4]1|l_3|2]\mu ^2}{t[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}}},`$
$`R_4`$ $`=`$ $`{\displaystyle \frac{1|4|2]3|l_3|4]\mu ^2}{[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}2(l_34)}},`$
$`R_5`$ $`=`$ $`{\displaystyle \frac{4|1|2]3|l_3|4]1|l_3|4]}{[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}2(l_34)}},`$
$`R_6`$ $`=`$ $`{\displaystyle \frac{3|l_3|4]1|l_3|2]}{[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}}},`$
$`R_7`$ $`=`$ $`{\displaystyle \frac{1|4|2]3|l_3|4]\mu ^2}{t[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}}},`$
$`R_8`$ $`=`$ $`{\displaystyle \frac{4|1|2]3|l_3|4]1|l_3|4]}{t[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}}},`$
$`R_9`$ $`=`$ $`{\displaystyle \frac{3|l_3|4]1|l_3|2]2(l_34)}{t[\mathrm{1\hspace{0.17em}3}]\mathrm{2\hspace{0.17em}4}}}.`$ (3.78)
The term $`R_5`$ becomes a quadratic box integral when the three delta functions are replaced with propagators. We can use the properties of the cut momenta to re-write $`R_5`$ as a sum of terms which will give a box integral, a linear box integral and a linear triangle integral as follows,
$$R_5=\frac{4|1|2][4|\mathrm{3\hspace{0.17em}1}|4]\mu ^2}{[\mathrm{1\hspace{0.17em}3}]^2\mathrm{2\hspace{0.17em}4}2(l_34)}+\frac{t4|1|2][4|3l_3|4]}{[\mathrm{1\hspace{0.17em}3}]^2\mathrm{2\hspace{0.17em}4}2(l_34)}\frac{4|1|2][4|3l_3|4]}{[\mathrm{1\hspace{0.17em}3}]^2\mathrm{2\hspace{0.17em}4}}.$$
(3.79)
We now replace the delta functions with propagators and integrate over the cut momenta. Note that one must drop any terms without cuts in the $`t`$-channel. This must be used for all the linear box integrals that appear above. Using the results for the linear box and the linear, quadratic and cubic triangle integrals in $`42ϵ`$ dimensions listed in the Appendix gives
$`R_1`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{st}{u^2}}K_4\right),`$
$`R_2`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{s^2t^2}{2u^3}}J_4{\displaystyle \frac{s^2t}{u^3}}J_3(t)\right),`$
$`R_3`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{st}{u^2}}J_3(t)+{\displaystyle \frac{s}{u^2}}J_2(t)\right),`$
$`R_4`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{s^2t^2}{2u^3}}J_4+{\displaystyle \frac{st^2}{u^3}}J_3(t)\right),`$
$`R_5`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{s^2t^2}{u^3}}J_4+{\displaystyle \frac{s^3t^3}{2u^4}}I_4+{\displaystyle \frac{s^2t^3}{u^4}}I_3(t)+{\displaystyle \frac{s^2t}{u^3}}I_2(t)\right),`$
$`R_6`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{st}{2u^2}}I_2(t)\right),`$
$`R_7`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{s}{u^2}}J_2(t)\right),`$
$`R_8`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{s^2}{u^2}}I_3^{62ϵ}(t)\right),`$
$`R_9`$ $``$ $`{\displaystyle \frac{𝒜_4^{\mathrm{tree}}}{(4\pi )^{2ϵ}}}\left({\displaystyle \frac{st}{4u^2}}I_2(t)+\left({\displaystyle \frac{s}{2u^2}}{\displaystyle \frac{s^2}{u^2t}}\right)I_2^{62ϵ}(t)+{\displaystyle \frac{s^2}{u^2}}I_3^{62ϵ}(t)\right).`$ (3.80)
Now using (A.26) in , and ignoring all terms without cuts in the $`t`$-channel, it is easy to show that the sum of these nine terms leads to the result
$`𝒜^{\mathrm{t}\mathrm{cut}}(1^{},2^+,3^{},4^+)`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^{2ϵ}}}𝒜_4^{\mathrm{tree}}({\displaystyle \frac{st}{u^2}}K_4{\displaystyle \frac{s^2t^2}{u^3}}I_4^{62ϵ}+{\displaystyle \frac{st}{u^2}}I_3^{62ϵ}(t)`$
$``$ $`{\displaystyle \frac{st(st)}{u^3}}J_3(t)+{\displaystyle \frac{s}{u^2}}J_2(t)+{\displaystyle \frac{s}{tu}}I_2^{62ϵ}(t)+{\displaystyle \frac{ts^2}{u^3}}I_2(t)).`$
Next, one must also include the corresponding terms coming from the $`s`$-channel version of the of triple cut in Figure 8. This just yields (3.4) with $`t`$ replaced by $`s`$. Combining these two expressions, without double-counting the box contributions (which appear in both cuts), and including the usual factor of two, one precisely reproduces the amplitude for this process (3.4)
## 4 The $`+`$$`+`$$`+`$$`+`$$`+`$ amplitude
The five-gluon all-plus one loop amplitude, with a scalar in the loop, is given by
$$𝒜_5(1^+,2^+,3^+,4^+,5^+)=\frac{i}{96\pi ^2C_5}\left[s_{12}s_{23}+s_{23}s_{34}+s_{34}s_{45}+s_{45}s_{51}+s_{51}s_{12}+4iϵ(1234)\right],$$
(4.1)
where $`C_5:=1223344551`$ and $`ϵ(abcd):=ϵ_{\mu \nu \rho \sigma }a^\mu b^\nu c^\rho d^\sigma `$.
An expression for the five-gluon amplitude valid to all orders in $`ϵ`$ appears in ,
$`𝒜_{5;1}^{\mathrm{scalar}}(1^+,2^+,3^+,4^+,5^+)={\displaystyle \frac{i}{C_5}}{\displaystyle \frac{ϵ(1ϵ)}{(4\pi )^{2ϵ}}}[s_{23}s_{34}I_4^{(1),82ϵ}+s_{34}s_{45}I_4^{(2),82ϵ}`$
$`+s_{45}s_{51}I_4^{(3),82ϵ}+s_{51}s_{12}I_4^{(4),82ϵ}+s_{12}s_{23}I_4^{(5),82ϵ}`$ (4.2)
(4.3)
$`+\mathrm{\hspace{0.17em}4}i(42ϵ)ϵ(1234)I_5^{102ϵ}].`$ (4.4)
The result (4.1) is obtained from (4.4) by taking the $`ϵ0`$ limit, where
$$ϵ(1ϵ)I_4^{82ϵ}\frac{1}{6},ϵ(1ϵ)I_5^{102ϵ}\frac{1}{24},ϵ(1ϵ)I_6^{102ϵ}0.$$
(4.5)
Here we will find that we can reproduce the full amplitude using only quadruple cuts in $`42ϵ`$ dimensions.
Let us start by considering the diagram in Figure 9, which represents the quadruple cut where gluons $`4`$ and $`5`$ enter the same tree amplitude. The momentum constraints on this quadruple cut are given by
$`l_1^2=l_2^2=l_3^2=l_4^2=\mu ^2,`$
$`l_1=l_4k_1,l_2=l_1k_2,l_3=l_2k_3,l_4=l_3k_4k_5.`$ (4.6)
It will prove convenient to solve for the momentum $`l_3`$, which we expand in the basis of vectors $`k_1,k_2,k_3`$ and $`K`$, where $`K`$ is defined in (3.23). One finds that the solution of (4) is given by<sup>9</sup><sup>9</sup>9We notice that, had we solved for $`l_1`$, the solution would have taken the form (3.24) with the same coefficients $`a`$, $`b`$, $`c`$, $`d`$ of (3.25) - but with $`u`$ defined by $`u=st(k_4+k_5)^2.`$
$$l_3=ak_1+bk_2+ck_3+dK,$$
(4.7)
with
$`a={\displaystyle \frac{t}{2u}},b={\displaystyle \frac{1}{2}},c=1{\displaystyle \frac{s}{2u}},`$ (4.8)
$`d=\pm \sqrt{{\displaystyle \frac{st+4\mu ^2u}{stu^2}}},`$
where the kinematical invariants $`s`$, $`t`$, $`u`$ are again defined by (3.26), but now $`s+t+u=(k_4+k_5)^2`$.
Considering the diagram in Figure 9, the product of tree-level amplitudes entering the quadruple cut can be written as
$$\frac{q_1|l_1|1]}{q_11}\frac{q_2|l_2|2]}{q_22}\frac{q_3|l_3|3]}{q_33}\frac{\mu ^2[45]}{45[(l_3k_4)^2\mu ^2]}.$$
(4.9)
Using (3.10), and choosing $`q_3=2`$, (4.9) can be recast as
$``$ $`\mu ^4{\displaystyle \frac{[12]}{12}}{\displaystyle \frac{[45]}{45}}{\displaystyle \frac{1}{23}}{\displaystyle \frac{2|l_3|3]}{(l_3k_4)^2\mu ^2}}={\displaystyle \frac{\mu ^4}{1223344551}}{\displaystyle \frac{\mathrm{Tr}_{}(5123l_34)}{(l_3k_4)^2\mu ^2}}`$ (4.10)
$`=`$ $`{\displaystyle \frac{\mu ^4}{1223344551}}{\displaystyle \frac{\mathrm{Tr}_+(123l_343)+\mathrm{Tr}_+(123l_342)}{(l_3k_4)^2\mu ^2}}.`$
Using momentum conservation, and
$$\mathrm{Tr}_+(abcd)=2\left[(ab)(cd)(ac)(bd)+(ad)(bc)+iϵ(abcd)\right],$$
(4.11)
it is easy to see that
$$\frac{\mathrm{Tr}_+(123l_343)+\mathrm{Tr}_+(123l_342)}{(l_3k_4)}=4(12)(23)\mathrm{\hspace{0.17em}4}i\frac{(34)ϵ(12l_33)(12)ϵ(234l_3)}{(l_3k_4)}.$$
(4.12)
We set
$$V(l_3)=iϵ(12l_33)(34)iϵ(234l_3)(12).$$
(4.13)
Now we wish to sum the expression (4.10) over the solutions (4.8), including a factor of $`1/2`$. Writing these solutions as $`l_3^\pm =x\pm y`$, where $`y`$ contains the term involving the momentum $`K`$, it is straightforward to show that
$$\frac{1}{2}\underset{l_3^\pm }{}\frac{\mathrm{Tr}_+(123l_343)+\mathrm{Tr}_+(123l_342)}{(l_3k_4)}=4(12)(23)\mathrm{\hspace{0.17em}4}\frac{V(x)(x4)V(y)(y4)}{(x4)^2(y4)^2},$$
(4.14)
and
$$\frac{V(x)(x4)V(y)(y4)}{(x4)^2(y4)^2}=\frac{i}{2}\mu ^2ϵ(1234)\left[\frac{1}{(l_3^+4)}+\frac{1}{(l_3^{}4)}\right].$$
(4.15)
Summarising, we have found that
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{l_3^\pm }{}}{\displaystyle \frac{\mathrm{Tr}_+(123l_343)+\mathrm{Tr}_+(123l_342)}{(l_3k_4)}}=4(12)(23)+\mathrm{\hspace{0.17em}2}i\mu ^2ϵ(1234)\left[{\displaystyle \frac{1}{(l_3^+4)}}+{\displaystyle \frac{1}{(l_3^{}4)}}\right]`$
$`=s_{12}s_{23}\mathrm{\hspace{0.17em}4}i\mu ^2ϵ(1234)\left[{\displaystyle \frac{1}{(l_3^+k_4)^2\mu ^2}}+{\displaystyle \frac{1}{(l_3^{}k_4)^2\mu ^2}}\right].`$ (4.16)
From (4.10), we see that the full amplitude in the quadruple cut is obtained by multiplying (4.16) by $`\mu ^4/C_5`$. Next, we lift the cut integral to a full Feynman integral, and get
$`\mathrm{\hspace{0.17em}2}{\displaystyle \frac{\mu ^4}{C_5}}\left[s_{12}s_{23}\mathrm{\hspace{0.17em}4}i\mu ^2ϵ(1234)\left({\displaystyle \frac{1}{(l_3^+k_4)^2\mu ^2}}+{\displaystyle \frac{1}{(l_3^{}k_4)^2\mu ^2}}\right)\right]`$
$`{\displaystyle \frac{i}{C_5(4\pi )^{2ϵ}}}\left[I_4^{(5),42ϵ}[\mu ^4]s_{12}s_{23}+\mathrm{\hspace{0.17em}8}iI_5^{42ϵ}[\mu ^6]ϵ(1234)\right]`$
$`={\displaystyle \frac{i}{C_5}}{\displaystyle \frac{ϵ(1ϵ)}{(4\pi )^{2ϵ}}}\left[s_{12}s_{23}I_4^{(5),82ϵ}+\mathrm{\hspace{0.17em}4}i(42ϵ)ϵ(1234)I_5^{102ϵ}\right],`$ (4.17)
where the factor of $`2`$ in the first line of (4.17) comes from adding, as usual, the two possible quadruple cuts of the amplitude (which are equal, since they are obtained one from the other by simply flipping all the internal “scalar helicities”).
Let us now discuss the result we have found. The first term in the last line of (4.17) gives the $`s_{12}s_{23}`$ term in (4.4). The other quadruple cut diagrams, which come from cyclic relabelling of the external legs, will similarly generate the other $`ϵ(1234)`$-independent terms in (4.4). Finally, the $`ϵ(1234)`$ term in (4.17) – a pentagon integral term – matches the $`ϵ(1234)`$ term in (4.4).
Thus we have shown that the five gluon amplitude $`+`$$`+`$$`+`$$`+`$$`+`$ may be reconstructed directly using quadruple cuts in $`42ϵ`$ dimensions.
## Acknowledgements
This work was partially supported by a Particle Physics and Astronomy Research Council award M Theory, string theory and duality, and a European Union Framework 6 Marie Curie Research Training Network grant Superstrings. AB would like to thank the Albert-Einstein-Institute in Golm and the School of Physics and Astronomy at Tel-Aviv University for hospitality during various stages of this work. GT would like to thank Marco Matone and the Physics Department at the University of Padova for hospitality during the initial stage of this work. We would like to thank James Bedford, Emil Bjerrum-Bohr, Stefano Catani, Chong-Sun Chu, Dave Dunbar, Nigel Glover, Massimiliano Grazzini, Harald Ita, Valya Khoze, David Kosower, Marco Matone and Sanjaye Ramgoolam for stimulating conversations.
Appendix A: Tensor Integrals
In this section we summarise the tensor bubble, tensor triangle and tensor box integrals used in this paper.
The scalar $`n`$-point integral functions in $`D=4+2m2ϵ`$ dimensions are defined as
$`I_n^DI_n^D[1]`$ $`=`$ $`i(1)^{n+1}(4\pi )^{D/2}{\displaystyle \frac{d^DL}{(2\pi )^D}\frac{1}{L^2(Lp_1)^2\mathrm{}(L_{i=1}^{n1}p_i)^2}}`$
$`=`$ $`{\displaystyle \frac{i(1)^{n+1}}{\pi ^{2+mϵ}}}{\displaystyle \frac{d^{4+2m}ld^{2ϵ}\mu }{(l^2\mu ^2)((lp_1)^2\mu ^2)\mathrm{}((l_{i=1}^{n1}p_i)^2\mu ^2)}}.`$
The higher dimensional integral functions are related to $`42ϵ`$ dimensional integrals with a factor $`\mu ^{2m}`$ inserted in the integrand. For $`m=1,2`$ one finds
$$I_n[\mu ^2]J_n=(ϵ)I_n^{62ϵ},\mathrm{and}I_n[\mu ^4]K_n=(ϵ)(1ϵ)I_n^{82ϵ}.$$
(A.2)
In our paper we encounter bubble functions with $`m=0,1`$, triangles with one massive external line and $`m=0,1`$, and boxes with four massless external lines and $`m=0,1,2`$:
$`I_2(P^2)={\displaystyle \frac{r_\mathrm{\Gamma }}{ϵ(12ϵ)}}(P^2)^ϵ,I_2^{62ϵ}(P^2)={\displaystyle \frac{r_\mathrm{\Gamma }}{2ϵ(12ϵ)(32ϵ)}}(P^2)^{1ϵ},`$
$`I_3(P^2)={\displaystyle \frac{r_\mathrm{\Gamma }}{ϵ^2}}(P^2)^{1ϵ},I_3^{62ϵ}(P^2)={\displaystyle \frac{r_\mathrm{\Gamma }}{2ϵ(1ϵ)(12ϵ)}}(P^2)^ϵ,`$
$`I_4={\displaystyle \frac{r_\mathrm{\Gamma }}{st}}\left\{{\displaystyle \frac{1}{ϵ^2}}\left[(s)^ϵ+(t)^ϵ\right]+{\displaystyle \frac{1}{2}}\mathrm{log}^2\left({\displaystyle \frac{s}{t}}\right)+{\displaystyle \frac{\pi ^2}{2}}\right\}+𝒪(ϵ),`$
$`(ϵ)I_4^{62ϵ}=0+𝒪(ϵ),(ϵ)(1ϵ)I_4^{82ϵ}={\displaystyle \frac{1}{6}}+𝒪(ϵ).`$ (A.3)
Note that the expressions for the bubbles and triangles are valid to all orders in $`ϵ`$, whereas for the box functions we have only kept the leading terms which contribute up to $`𝒪(ϵ^0)`$ in the amplitudes.
We now move on to present the result of the PV reduction for various tensor integrals which are relevant for this paper. Note that the expressions are presented in terms of scalar $`n`$-point integral functions $`I_n^𝒟`$ in various dimensions $`𝒟`$, specifically in terms of $`I_n`$, $`I_n^{62ϵ}`$ and $`I_n^{82ϵ}`$ in $`42ϵ`$, $`62ϵ`$ and $`82ϵ`$ dimensions, respectively. The expressions are valid to all orders in $`ϵ`$, if $`I_n`$, $`I_n^{62ϵ}`$ and $`I_n^{82ϵ}`$ are evaluated to all orders, and the PV reductions have been performed in a fashion that naturally leads to coefficients without explicit $`ϵ`$ dependence (the reader may consult for more details on this particular variant of PV reductions).
For the linear and two-tensor bubbles we have (see Figure 10a):
$`I_2\left[L_3^\mu \right]={\displaystyle \frac{1}{2}}I_2(p_2+p_3)^\mu ,`$ (A.4)
$`I_2\left[L_3^\mu L_3^\nu \right]={\displaystyle \frac{1}{2}}I_2^{62ϵ}\delta _{[42ϵ]}^{\mu \nu }+\left({\displaystyle \frac{1}{4}}I_2+{\displaystyle \frac{1}{2t}}I_2^{62ϵ}\right)(p_2+p_3)^\mu (p_2+p_3)^\nu .`$ (A.5)
For the linear, two- and three-tensor triangles (see Figure 10b):
$`I_3\left[L_3^\mu \right]={\displaystyle \frac{1}{t}}I_2p_2^\mu +\left(I_3+{\displaystyle \frac{1}{t}}I_2\right)p_3^\mu ,`$ (A.6)
$`I_3\left[L_3^\mu L_3^\nu \right]={\displaystyle \frac{1}{2t}}I_2p_2^\mu p_2^\nu +\left({\displaystyle \frac{1}{t}}I_3^{62ϵ}+{\displaystyle \frac{1}{2t}}I_2\right)\left(p_2^\mu p_3^\nu +p_2^\nu p_3^\mu \right)`$
$`+\left({\displaystyle \frac{3}{2t}}I_2+I_3\right)p_3^\mu p_3^\nu {\displaystyle \frac{1}{2}}I_3^{62ϵ}\delta _{[42ϵ]}^{\mu \nu },`$ (A.7)
$`I_3\left[L_3^\mu L_3^\nu L_3^\rho \right]=\left({\displaystyle \frac{1}{4t}}I_2+{\displaystyle \frac{1}{2t^2}}I_2^{62ϵ}\right)\left(p_2^\mu p_2^\nu p_2^\rho \right)`$
$`\left({\displaystyle \frac{1}{4t}}I_2+{\displaystyle \frac{3}{2t^2}}I_2^{62ϵ}\right)\left(p_2^\mu p_2^\nu p_3^\rho +p_2^\mu p_3^\nu p_2^\rho +p_3^\mu p_2^\nu p_2^\rho \right)`$
$`+\left({\displaystyle \frac{1}{4t}}I_2+{\displaystyle \frac{3}{2t^2}}I_2^{62ϵ}{\displaystyle \frac{2}{t}}I_3^{62ϵ}\right)\left(p_2^\mu p_3^\nu p_3^\rho +p_3^\mu p_3^\nu p_2^\rho +p_3^\mu p_2^\nu p_3^\rho \right)`$
$`+\left({\displaystyle \frac{7}{4t}}I_2+{\displaystyle \frac{1}{2t^2}}I_2^{62ϵ}I_3\right)\left(p_3^\mu p_3^\nu p_3^\rho \right)+{\displaystyle \frac{1}{2t}}I_2^{62ϵ}\left(\delta ^{\mu \nu }p_2^\rho +\delta ^{\mu \rho }p_2^\nu +\delta ^{\rho \nu }p_2^\mu \right)`$
$`+\left({\displaystyle \frac{1}{2t}}I_2^{62ϵ}+{\displaystyle \frac{1}{2}}I_3^{62ϵ}\right)\left(\delta ^{\mu \nu }p_3^\rho +\delta ^{\mu \rho }p_3^\nu +\delta ^{\rho \nu }p_3^\mu \right).`$ (A.8)
Finally, for the linear box:
$`I_4\left[L_3^\mu \right]`$ $`=`$ $`\left({\displaystyle \frac{t}{2u}}I_4{\displaystyle \frac{1}{u}}\left(I_3(t)I_3(s)\right)\right)p_1^\mu {\displaystyle \frac{1}{2}}I_4p_2^\mu `$ (A.9)
$`+`$ $`\left({\displaystyle \frac{tu}{2u}}I_4{\displaystyle \frac{1}{u}}\left(I_3(t)I_3(s)\right)\right)p_3^\mu ,`$
where, as usual, $`I_n^D`$ denote $`D`$-dimensional scalar $`n`$-point integral functions, $`s:=(p_1+p_2)^2`$, $`t:=(p_2+p_3)^2`$, $`u:=(p_1+p_3)^2`$, and $`I_n`$ is an abbreviation for the $`(42ϵ)`$-dimensional integral functions.
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# State Specific Kohn–Sham Density Functional Theory
## I Introduction
The Kohn-Sham version of density functional theory plays a major role in both quantum chemistry and condensed matter physics Dreizler and E. K. U. Gross (1990); Parr and Yang (1989); Springborg (1997); Ellis (1995); E. K. U. Gross and Dreizler (1994); Seminario and Politzer (1995); Handy (1997). The local density approximation Kohn and Sham (1965) has been widely used for the solid state. While for molecules, by far, the most successful functional, a hybrid one Becke (1993); Burke et al. (1997); Perdew et al. (1996); Ernzerhof (1996), is known as B3LYP Becke (1993); P. J. Stephens et al. (1994).
The Kohn–Sham approach, however, does have well known shortcomings. For example, a constraint search definition Levy (1978, 1979, 1982); Levy and Perdew (1985) is required to treat the $`v`$–representability problem that arises in the original Kohn–Sham method Kohn and Sham (1965). Unfortunately, this formal definition is difficult to consider when deriving approximate functionals. Furthermore, in contrast to wave function based methods, the exchange-correlation functional is an unknown, implicit functional, and there is no systematic method to improve approximations. In addition, there are well known errors arising from Coulomb self-interactions that appears when using approximate functionals Parr and Yang (1989); Dreizler and E. K. U. Gross (1990); Koch and Holthausen (2000). Also, the most widely used approximate functional for molecular systems, the B3LYP functional, includes a component of the exact exchange-potential, even though the Kohn–Sham approach requires the noninteracting state to come from a local potential. The optimized potential method Fiolhais et al. (2003); Sharp and Horton (1953); Talman and Shadwick (1976); Li et al. (1993); Shaginyan (1994); Görling and Levy (1994); Grabo et al. (2000) is an approach to convert a nonlocal operator into a local potential. Unfortunately, this method leads to potentials that are not invariant to a unitary transformation of orbitals and depend explicitly on the individual orbitals and orbital energies.
The formalism presented below uses an electronic-energy functional containing a correlation energy functional $`_{\text{co}}`$ that depends on the external potential $`v`$ and on the one-particle density matrix $`\rho _1`$ of determinantal states. Since the $`v`$–representability problem does not appear, a constrain search definition is not needed. Also, since the approach uses the exact exchange-potential, errors from Coulomb self-interactions do not occur. The energy functionals, however, contains multiple solutions, since any one-particle density matrix $`\rho _1`$ delivering the density from the interacting state yields a solution. In order to obtain the Kohn–Sham solution, the nonlocal operators are converted into local ones using an approach developed by Sala and Görling Sala and Görling (2001). In contrast to the optimized potential method Fiolhais et al. (2003); Sharp and Horton (1953); Talman and Shadwick (1976); Li et al. (1993); Shaginyan (1994); Görling and Levy (1994); Grabo et al. (2000), the energy functionals and local potentials are invariant to a unitary transformation of orbitals and do not depend on the individual orbital or the orbital energies. A density functional formalism is also derived that assumes that the one-particle density matrices of interest have $`v`$-representable (non-interacting) densities and that these density matrices can be written as an explicit functional of the electron density.
Previously we have shown that the correlation energy from many body perturbation theory Lindgren and Morrison (1986); Harris et al. (1992); Raimes (1972) can be written as an explicit functional of $`v`$ and $`\rho _1`$ Finley (2003). In a similar manner, but using less restrictive energy denominators, the correlation energy functionals presented below can be shown to be an explicit functional of $`v`$ and $`\rho _1`$ Finley (2005). Hence, in contrast to the Kohn–Sham method, it maybe possible to derive approximate functionals that can be improved in a systematic manner. For simplicity, we only consider noninteracting closed-shell states and target states that are nondegenerate, singlet ground-states.
## II The energy functionals and trial wave functions
Our interest is in finding the ground-state eigenvalue of the Hamiltonian operator,
$`\widehat{H}_{Nv}=\widehat{T}+\widehat{V}_{\text{ee}}+\widehat{V}_v,`$ (1)
where
$`\widehat{T}`$ $`=`$ $`{\displaystyle \underset{i}{\overset{N}{}}}(\frac{1}{2}_i^2),`$ (2)
$`\widehat{V}_{\text{ee}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{\overset{N}{}}}r_{ij}^1,`$ (3)
$`\widehat{V}_v`$ $`=`$ $`{\displaystyle \underset{i}{\overset{N}{}}}v(i),`$ (4)
and $`v`$ is the external potential; $`N`$ is the number of electrons. Since the Hamiltonian $`\widehat{H}_{Nv}`$ is determined by $`N`$ and $`v`$, so are the ground state wave functions $`|\mathrm{\Psi }_{Nv}`$ that satisfy the Schrödinger equation:
$`\widehat{H}_{Nv}|\mathrm{\Psi }_{Nv}=_{Nv}|\mathrm{\Psi }_{Nv},`$ (5)
where, for simplicity, we only consider wave functions that are nondegenerate, singlet ground-states.
Using a second quantization approach, our spin-free Hamiltonian does not depend on $`N`$, and it can be expressed by
$`\widehat{H}_v={\displaystyle \underset{ij}{}}(i|(\frac{1}{2}^2)|j)\widehat{E}_{ij}+{\displaystyle \underset{ij}{}}(i|v|j)\widehat{E}_{ij}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ijkl}{}}(ij|kl)\widehat{E}_{ijkl},`$ (6)
where the symmetry-adapted excitation operators are given by
$`\widehat{E}_{ij}`$ $`=`$ $`{\displaystyle \underset{\sigma }{}}a_{i\sigma }^{}a_{j\sigma },`$ (7)
$`\widehat{E}_{ijkl}`$ $`=`$ $`{\displaystyle \underset{\sigma \lambda }{}}a_{i\sigma }^{}a_{k\lambda }^{}a_{l\lambda }a_{j\sigma },`$ (8)
and the one- and two electrons integrals are spin-free integrals written in chemist’s notation Szabo and N. S. Ostlund (1982) using a spatial orbital set, say $`\{\chi \}`$; this set has the following form:
$`\psi _{j\sigma }(𝐱)=\chi _j(𝐫)\sigma (\omega );\sigma =\alpha ,\beta ,`$ (9)
where the spatial and spin coordinates, $`𝐫`$ and $`\omega `$, are denoted collectively by $`𝐱`$.
Wave function-based methods including perturbation theory, configuration interaction, and coupled cluster theory, use one or more reference states to express $`\mathrm{\Psi }`$ and $``$. For closed-shell ground-state wave functions, a single determinant can be used, where closed-shell determinantal, or noninteracting, states can be constructed from a set of doubly occupied spatial-orbitals; these occupied orbitals also determine the spin-less one-particle density-matrix of the noninteracting state, given by Parr and Yang (1989); McWeeny (1960)
$`\rho _1(𝐫_1,𝐫_2)=2{\displaystyle \underset{w\{\chi _o\}}{}}\chi _w(𝐫_1)\chi _w^{}(𝐫_2),`$ (10)
where the sum is over the occupied orbitals; this set of orbitals is denoted by $`\{\chi _o\}`$.
For later use, we also mention that for a complete basis set we have
$`2\delta (𝐫_1𝐫_2)=\rho _1(𝐫_1,𝐫_2)+\kappa _{\rho _1}(𝐫_1,𝐫_2),`$ (11)
where $`\kappa _{\rho _1}`$ is determined by the excited orbitals,
$`\kappa _{\rho _1}(𝐫_1,𝐫_2)=2{\displaystyle \underset{r\{\chi _u\}}{}}\chi _r(𝐫_1)\chi _r^{}(𝐫_2),`$ (12)
and $`\{\chi _u\}`$ denotes the set of orbitals orthogonal to the occupied set $`\{\chi _o\}`$. The operator form of Eq. (11) is
$`2\widehat{I}=\widehat{\rho }_1+\widehat{\kappa }_{\rho _1},`$ (13)
where $`\widehat{I}`$ is the identity operator; so, the kernels of the three operators within Eq. (13) are given by the corresponding terms within Eq. (11).
It is well known that there is a one-to-one mapping between determinantal states and their one-particle density matrices Parr and Yang (1989); J.-P. Blaizot and Ripka (1986), say $`\gamma `$, where for a closed-shell state described by the orbitals given by Eq. (9), we have P. A. M. Dirac (1930, 1931); P. -O. Löwdin (1955a, b)
$`\gamma (𝐱_1,𝐱_2)={\displaystyle \underset{w\{\chi _o\}}{}}{\displaystyle \underset{\sigma }{}}\chi _w(𝐫_1)\chi _w^{}(𝐫_2)\sigma (\omega _1)\sigma ^{}(\omega _2),`$ (14)
and by using Eq. (10), we obtain
$`\gamma (𝐱_1,𝐱_2)={\displaystyle \frac{1}{2}}\rho _1(𝐫_1,𝐫_2)\delta _{\omega _1\omega _2}.`$ (15)
Since our closed-shell determinantal states are determined by $`\rho _1`$, we denote these kets by $`|\rho _1`$.
According to the Hohenberg-Kohn theorem Hohenberg and Kohn (1964); Parr and Yang (1989); Dreizler and E. K. U. Gross (1990), the external potential $`v`$ is determined by the density, and the density also determines $`N`$. So, in principle, we can replace the variables $`N`$ and $`v`$ by the electronic density $`n`$ and, at least for nondegenerate ground-states, write
$`\widehat{H}_v|\mathrm{\Psi }_n=_n|\mathrm{\Psi }_n;nN,v,`$ (16)
where these functions serve as density-dependent trial-wave functions for the Kohn-Sham approach. Notice we have omitted the $`N`$ subscript on the Hamiltonian operator, since $`\widehat{H}_v`$ is independent of $`N`$ when this operator is expressed in second quantization.
As an alternative to a density-dependent wave function, we consider trial wave functions, say $`|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}`$, that are determined by the one-body external potential $`v`$ and, in addition, by the spin-less one-particle density-matrix $`\rho _1`$ of a noninteracting state, and, as mentioned previously, these noninteracting states are denoted by $`|\rho _1`$.
By definition, our trial wave function $`|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}`$ yields the exact ground-state wave function $`|\mathrm{\Psi }_n`$ when the noninteracting density $`\rho _s`$, i.e., the density of $`|\rho _1`$, equals the exact density $`n`$ of the interacting state $`|\mathrm{\Psi }_n`$, where $`n`$ also determines the $`v`$ and $`N`$. This state of affairs can be represented by the following:
$`|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}=|\mathrm{\Psi }_n;\rho _1\rho _s=n,nN,v.`$ (17)
In other words, $`\rho _1`$ determines $`\rho _s`$, and when $`\rho _s=n`$, $`|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}`$ yields $`|\mathrm{\Psi }_n`$. Letting $`\varrho _1`$ denote the one-particle density matrix of interest, we can write
$`|\stackrel{~}{\mathrm{\Psi }}_{v\varrho _1}=|\mathrm{\Psi }_n;\varrho _1n,nN,v.`$ (18)
For later use, we also mention that the density $`n`$ of an interacting state can be partitioned as
$`n=\rho _s+\rho _c,`$ (19)
where the correlation density is given by
$`\rho _c(𝐫)={\displaystyle \frac{\mathrm{\Psi }_n|\widehat{\mathrm{\Gamma }}(𝐫)|\mathrm{\Psi }_n}{\mathrm{\Psi }_n|\mathrm{\Psi }_n}}\rho _s(𝐫),`$ (20)
and $`\widehat{\mathrm{\Gamma }}`$ is the density operator, given by Eq. (87).
Using our trial wave function, we introduce a variational energy functional:
$`E_v[\rho _1]={\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{H}_v|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}.`$ (21)
Our trial wave functions $`|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}`$ and energy functionals $`E_v[\rho _1]`$ are assumed to be explicit functionals of $`\rho _1`$ and $`v`$. However, two different one-particle density matrices, say $`\rho _1`$ and $`\rho _{1}^{}{}_{}{}^{}`$, that yield the same density $`\rho _s`$, i.e., $`\rho _1\rho _s`$ and $`\rho _{1}^{}{}_{}{}^{}\rho _s`$, yield the same $`|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}`$ and $`E_v[\rho _1]`$, so these functions are implicit functionals of $`\rho _s`$, and, therefore, we can write $`|\stackrel{~}{\mathrm{\Psi }}_{v\rho _s}`$ and $`E_v[\rho _s]`$. However, we will continue to consider them as functionals of their explicit variable $`\rho _1`$.
Using Eqs. (16) and (18), we observe that our energy functional $`E_v`$, given by Eq. (21), delivers the exact energy $`_n`$ when the one-particle density matrix determines the exact density $`n`$:
$`E_v[\varrho _1]=_n,\varrho _1n,nN,v,`$ (22)
and for an arbitrary density we get
$`E_v[\rho _1]_n,\rho _1\rho _sN,`$ (23)
where the density $`\rho _s`$ from the noninteracting state $`|\rho _1`$ is not necessarily $`v`$-representable.
## III Trial Hamiltonians
Our trial wave function is a ground-state eigenfunction of a Hamiltonian operator that depend explicitly on the one-particle density of a noninteracting state:
$`\widehat{H}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}=\stackrel{~}{}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}.`$ (24)
As in our trial wave functions $`|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}`$ and energy functionals $`E_v[\rho _1]`$, the trial Hamiltonians $`\widehat{H}_{v\rho _1}`$ are explicit functionals of $`\rho _1`$, but implicit functionals of $`\rho _s`$. So two trial Hamiltonians, say $`\widehat{H}_{v\rho _1}`$ and $`\widehat{H}_{v\rho _{1}^{}{}_{}{}^{}}`$, are equal if both $`\rho _1`$ and $`\rho _{1}^{}{}_{}{}^{}`$ yield the same density, i.e., $`\rho _1,\rho _{1}^{}{}_{}{}^{}\rho _s`$.
Our trial Hamiltonians must be chosen so that Eq. (18) is satisfied, indicating the following identity:
$`\widehat{H}_{v\varrho _1}=\widehat{H}_v,\varrho _1n,nN,v.`$ (25)
There are many ways to obtain a trial Hamiltonian that satisfies Eq. (25). Consider the following trial Hamiltonian obtained by adding a term to the Hamiltonian:
$`\widehat{H}_{v\rho _1}=\widehat{H}_v+\lambda {\displaystyle 𝑑𝐫g_{\rho _c}(𝐫)\left(\widehat{\mathrm{\Gamma }}(𝐫)\rho _s(𝐫)\right)},\rho _1\rho _s,`$ (26)
where $`\widehat{\mathrm{\Gamma }}(𝐫)`$ is the density operator, given by Eq. (87); $`(\widehat{\mathrm{\Gamma }}(𝐫)\rho _s(𝐫))`$ is the one-body portion of $`\widehat{\mathrm{\Gamma }}(𝐫)`$ when this operator is written in normal-ordered form Čížek (1966, 1969); Lindgren and Morrison (1986); Paldus and Čížek (1975), given by Eq. (86). Furthermore, $`\lambda `$ is an arbitrary constant, and the functional $`g`$ is also arbitrary, except that it vanishes when the correlation density $`\rho _c`$ vanishes
$`\underset{\rho _c0}{lim}g_{\rho _c}(𝐫)=0,`$ (27)
where $`\rho _c`$ is defined by Eqs. (19) and (20).
Since $`(\widehat{\mathrm{\Gamma }}(𝐫)\rho _s(𝐫))`$ is normal-ordered, we have
$`\rho _1|\left(\widehat{\mathrm{\Gamma }}(𝐫)\rho _s(𝐫)\right)|\rho _1=0.`$ (28)
Therefore, the added term appearing in Eq. (26) can be considered a sort of correlation term, since it does not contribute in first order. Hence, we have
$`\rho _1|\widehat{H}_{v\rho _1}|\rho _1=\rho _1|\widehat{H}_v|\rho _1.`$ (29)
One possible choice for $`g_{\rho _c}`$, and presented in Appendix A, is given by
$`g_{\rho _c}(𝐫_1)={\displaystyle 𝑑𝐫_2r_{12}^1\rho _c(𝐫_2)}.`$ (30)
## IV A generalization of the Kohn-Sham formalism
We now obtain a generalization of the Kohn-Sham formalism. Substituting Eq. (1) into Eq. (21) gives
$`E_v[\rho _1]={\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{T}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}+{\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{V}_{\text{ee}}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}+{\displaystyle 𝑑𝐫v(𝐫)\rho _s(𝐫)}+{\displaystyle 𝑑𝐫v(𝐫)\stackrel{~}{\rho }_c(𝐫)},`$ (31)
where $`\stackrel{~}{\rho }_c`$ is the correlation density of the trial wave function, i.e, as in Eq. (20), we have
$`\stackrel{~}{\rho }_c(𝐫)={\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{\mathrm{\Gamma }}(𝐫)|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}\rho _s(𝐫)=\stackrel{~}{n}\rho _s(𝐫),\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}\stackrel{~}{n},\rho _1\rho _s,`$ (32)
and $`\stackrel{~}{n}`$ is the density of $`\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}`$.
Through the first-order, the kinetic energy and electron-electron repulsion energy are given, respectively, by
$`\rho _1|\widehat{T}|\rho _1`$ $`=`$ $`{\displaystyle 𝑑𝐫_1\left[\frac{1}{2}_1^2\rho _1(𝐫_1,𝐫_2)\right]_{𝐫_2=𝐫_1}},`$ (33)
$`\rho _1|\widehat{V}_{\text{ee}}|\rho _1`$ $`=`$ $`E_J[\rho _s]+E_\text{x}[\rho _1],`$ (34)
where the Coulomb and exchange energies are
$`E_J[\rho _s]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle r_{12}^1𝑑𝐫_1𝑑𝐫_2\rho (𝐫_1)\rho (𝐫_2)},`$ (35)
$`E_\text{x}[\rho _1]`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle r_{12}^1𝑑𝐫_1𝑑𝐫_2\rho _1(𝐫_1,𝐫_2)\rho _1(𝐫_2,𝐫_1)}.`$ (36)
Adding and subtracting $`\rho _1|\widehat{T}|\rho _1`$ and $`\rho _1|\widehat{V}_{\text{ee}}|\rho _1`$, Eq. (31) can be written as
$`E_v[\rho _1]={\displaystyle 𝑑𝐫_1\left[\frac{1}{2}_1^2\rho _1(𝐫_1,𝐫_2)\right]_{𝐫_2=𝐫_1}}+{\displaystyle 𝑑𝐫v(𝐫)\rho _s(𝐫)}`$
$`+E_J[\rho _s]+E_\text{x}[\rho _1]+E_{\text{co}}[\rho _1,v]+{\displaystyle 𝑑𝐫v(𝐫)\stackrel{~}{\rho }_c(𝐫)},`$ (37)
where the correlation-energy functional is given by
$`E_{\text{co}}[\rho _1,v]={\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{T}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}\rho _1|\widehat{T}|\rho _1+{\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{V}_{\text{ee}}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}\rho _1|\widehat{V}_{\text{ee}}|\rho _1.`$ (38)
Recognizing the first four terms from Eq. (IV) as the energy through the first order, $`_1`$, we can write
$`E_v[\rho _1]=_1[\rho _1,v]+E_{\text{co}}[\rho _1,v]+{\displaystyle 𝑑𝐫v(𝐫)\stackrel{~}{\rho }_c(𝐫)},`$ (39)
where
$`_1[\rho _1,v]=\rho _1|H_v|\rho _1={\displaystyle 𝑑𝐫_1\left[\frac{1}{2}_1^2\rho _1(𝐫_1,𝐫_2)\right]_{𝐫_2=𝐫_1}}`$ (40)
$`+{\displaystyle 𝑑𝐫_1v(𝐫_1)\rho (𝐫_1)}+{\displaystyle \frac{1}{2}}{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho (𝐫_1)\rho (𝐫_2)}{\displaystyle \frac{1}{4}}{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho _1(𝐫_1,𝐫_2)\rho _1(𝐫_2,𝐫_1)}.`$
Now consider the correlation energy that is obtained by wave function methods. Using the notation from Eq. (5), and a reference state $`|\rho _1`$, the correlation energy is given by
$`_{\text{co}}[\rho _1,v]={\displaystyle \frac{\mathrm{\Psi }_{Nv}|\widehat{H}_v|\mathrm{\Psi }_{Nv}}{\mathrm{\Psi }_{Nv}|\mathrm{\Psi }_{Nv}}}_1[\rho _1,v],`$ (41)
where previously we have shown that $`_{\text{co}}`$ can be written as an explicit functional of $`v`$ and $`\rho _1`$ Finley (2003). In a similar manner, but using less restrictive energy denominators, our correlation energy functional $`E_{\text{co}}`$, given by Eq. (38), can be shown to be an explicit functional of $`v`$ and $`\rho _1`$ Finley (2005). Therefore, by requiring the last term within Eq. (39) to be an explicit functional of $`v`$ and $`\rho _1`$, $`E_v`$ can also be written as an explicit functional of $`v`$ and $`\rho _1`$ Finley (2005).
We now focus our attention on minimizing the energy functional $`E_v`$, subject to the constraint that the spin-less one-particle density-matrix $`\rho _1`$ comes from a closed-shell single-determinantal state. For the more general case of a determinantal state, say $`|\gamma `$, with the (spin-dependent) one-particle density matrix $`\gamma `$, as in Eq. (14), the two necessary conditions for $`\gamma `$ to satisfy are given by the following J.-P. Blaizot and Ripka (1986); Parr and Yang (1989):
$`{\displaystyle \gamma (𝐱_3,𝐱_4)\delta (𝐱_3𝐱_4)𝑑𝐱_3𝑑𝐱_4}=N,`$ (42)
$`{\displaystyle \gamma (𝐱_3,𝐱_5)\gamma (𝐱_5,𝐱_4)𝑑𝐱_5}=\gamma (𝐱_3,𝐱_4),`$ (43)
where the first relation indicates that the electron density yields the number of electrons $`N`$; the second relation indicates that $`\gamma `$ is indempotent. For our special closed-shell case, we substitute Eq. (15) into the above constrains, yielding the following conditions:
$`{\displaystyle \rho _1(𝐫_3,𝐫_4)\delta (𝐫_3𝐫_4)𝑑𝐫_3𝑑𝐫_4}=N,`$ (44)
$`{\displaystyle \rho _1(𝐫_3,𝐫_5)\rho _1(𝐫_5,𝐫_4)𝑑𝐫_5}=2\rho _1(𝐫_3,𝐫_4).`$ (45)
It is well know that the functional derivative of $`_1`$ with respect to the $`\gamma `$ yields the kernel of the Fock operator Parr and Yang (1989). For the closed-shell case, we have
$`F(𝐫_1,𝐫_2)={\displaystyle \frac{\delta _1[\rho _1,v]}{\delta \rho _1(𝐫_2,𝐫_1)}},`$ (46)
where, using Eq. (40), the Fock kernel is given by
$`F_{\rho _1}(𝐫_1,𝐫_2)=\delta (𝐫_1𝐫_2)\left(\frac{1}{2}_2^2+v(𝐫_2)+{\displaystyle 𝑑𝐫_3r_{23}^1\rho (𝐫_3)}\right)+v_\text{x}^{\rho _1}(𝐫_1,𝐫_2),`$ (47)
and the exchange operator, say $`\widehat{v}_\text{x}^{\rho _1}`$, has the following kernel:
$`v_\text{x}^{\rho _1}(𝐫_1,𝐫_2)={\displaystyle \frac{1}{2}}r_{12}^1\rho _1(𝐫_1,𝐫_2).`$ (48)
By generalizing Eq. (46), we define a generalized, or exact, Fock operator $`\widehat{}`$, where the kernel of this operator is
$`_{\rho _1}(𝐫_1,𝐫_2)={\displaystyle \frac{\delta E_v[\rho _1]}{\delta \rho _1(𝐫_2,𝐫_1)}}=F_{\rho _1}(𝐫_1,𝐫_2)+v_{\text{co}}^{\rho _1}(𝐫_1,𝐫_2)+v_{\text{ec}}^{\rho _1}(𝐫_1,𝐫_2),`$ (49)
and the correlation operator $`\widehat{v}_{\text{co}}^{\rho _1}`$ and external-correlation operator $`\widehat{v}_{\text{ec}}^{\rho _1}`$ are defined by their kernels:
$`v_{\text{co}}^{\rho _1}(𝐫_1,𝐫_2)`$ $`=`$ $`{\displaystyle \frac{\delta E_{\text{co}}[\rho _1,v]}{\delta \rho _1(𝐫_2,𝐫_1)}},`$ (50)
$`v_{\text{ec}}^{\rho _1}(𝐫_1,𝐫_2)`$ $`=`$ $`{\displaystyle \frac{\delta \left(𝑑𝐫_3v(𝐫_3)\stackrel{~}{\rho }_c(𝐫_3)\right)}{\delta \rho _1(𝐫_2,𝐫_1)}}.`$ (51)
Minimizing the functional $`E_v`$, given by Eq. (39), subject to the constraints given by Eqs. (44) and (45), is very similar to the corresponding Hartree–Fock derivation Parr and Yang (1989) and the derivation for reference-state one-particle density matrix theory Finley (2004a, 2003, b). The only difference being that the spin variable has been eliminated, and we have a factor of two appearing in Eq. (45). Therefore, we only state the main results, i.e., this minimization yields the exact electronic energy $`_n`$ for the interacting state, as given by Eq. (22), where the one-particle density-matrix $`\varrho _1`$ satisfies the following conditions:
$`\widehat{\kappa }_{\varrho _1}\widehat{}_{\varrho _1}\widehat{\varrho }_1`$ $`=`$ $`0,`$ (52)
$`\widehat{\varrho }_1\widehat{}_{\varrho _1}\widehat{\kappa }_{\varrho _1}`$ $`=`$ $`0,`$ (53)
and the kernels of the operators $`\widehat{\rho }_1`$ and $`\widehat{\kappa }_{\rho _1}`$ are given by the terms on the right side of Eq. (11); also, as mentioned previously, $`\varrho _1`$ yields the exact density $`n`$ of the interacting state $`\mathrm{\Psi }_n`$. Using Eqs. (52) and (53), it is readily shown that $`\widehat{}_{\varrho _1}`$ and $`\widehat{\varrho }_1`$ commute:
$`[\widehat{}_{\varrho _1},\widehat{\varrho }_1]=0,`$ (54)
and the occupied orbitals satisfy a generalized Hartree–Fock equation:
$`\widehat{}_{\varrho _1}\chi _w={\displaystyle \underset{x\varrho _1}{}}\epsilon _{xw}\chi _x,`$ (55)
where the notation $`x\varrho _1`$ indicates a summation over the occupied orbitals from the determinantal state $`|\varrho _1`$; $`\chi _w`$ is also an occupied orbital from $`|\varrho _1`$. Furthermore, we can choose orbitals that diagonalize the matrix $`\epsilon _{xw}`$, yielding exact, canonical Hartree–Fock equations:
$`\left(\frac{1}{2}^2+v+v_j^n+\widehat{v}_\text{x}^{\varrho _1}+\widehat{v}_{\text{co}}^{\varrho _1}+\widehat{v}_{\text{ec}}^{\varrho _1}\right)\chi _w=\epsilon _w\chi _w,\chi _w\varrho _1,`$ (56)
where the Coulomb operator is defined by
$`v_j^\rho (𝐫_1)\chi (𝐫_1)={\displaystyle 𝑑𝐫_2r_{12}^1\rho (𝐫_2)\chi (𝐫_1)},`$ (57)
and we have
$`\varrho _1(𝐫,𝐫)=n(𝐫).`$ (58)
Equation (56) is also satisfied by the canonical excited orbitals.
For later use, we also mention that the determinantal states $`|\varrho _1`$ satisfy the following noninteracting Schrödinger equation:
$`{\displaystyle \underset{i=1}{\overset{N}{}}}\widehat{}_{\varrho _1}(𝐫_i)|\varrho _1=2\left({\displaystyle \underset{w}{}}\epsilon _w\right)|\varrho _1.`$ (59)
Appendix B presents an alternative way of partitioning the energy functional that differs from Eq. (39).
## V Conversion of the nonlocal potential into a local one
As mentioned previously, our energy functionals $`E_v`$ are implicit functionals of the noninteracting density $`\rho _s`$. Hence, any one-particle density-matrix that yields the interacting density minimizes our energy functional, i.e., we have
$`_n=E_v[\varrho _1]=E_v[\varrho _1^{}]=E_v[\varrho _1^{\prime \prime }]\mathrm{},`$ (60)
where
$`n(𝐫)=\varrho _1(𝐫,𝐫)=\varrho _1^{}(𝐫,𝐫)=\varrho _1^{\prime \prime }(𝐫,𝐫)\mathrm{},`$ (61)
and there are other solutions besides Eq. (56), e.g,
$`\widehat{}_{\varrho _1^{}}\chi _w=\left(\frac{1}{2}^2+v+v_j^n+\widehat{w}_{\varrho _1^{}}\right)\chi _w=\epsilon _w^{}\chi _w,\chi _w\varrho _1^{},`$ (62)
where the nonlocal potential $`\widehat{w}_{\rho _1}`$ is given by
$`\widehat{w}_{\rho _1}=\widehat{v}_\text{x}^{\rho _1}+\widehat{v}_{\text{co}}^{\rho _1}+\widehat{v}_{\text{ec}}^{\rho _1}.`$ (63)
Assuming $`n`$ is a noninteracting $`v`$-representable density, there exist a noninteracting state, say $`|\phi _1`$, that has $`n`$ as its density:
$`n(𝐫)=\phi _1(𝐫,𝐫),`$ (64)
and this determinant—assuming it is a closed-shell determinant—is the ground-state solution of the following noninteracting Schrödinger equation:
$`{\displaystyle \underset{i=1}{\overset{N}{}}}\widehat{f}(𝐫_i)|\phi _1=2\left({\displaystyle \underset{w}{}}ϵ_w\right)|\phi _1,`$ (65)
where
$`\widehat{f}=\frac{1}{2}^2+v_s,`$ (66)
and $`v_s`$ is a local potential. Therefore, the canonical occupied orbitals from $`|\phi _1`$ satisfy the following one-particle Schrödinger equation:
$`\widehat{f}\varphi _w=\left(\frac{1}{2}^2+v+v_j^n+v_{\text{xc}}\right)\varphi _w=ϵ_w\varphi _w,\varphi _w\phi _1,`$ (67)
where with no loss of generality, we have required $`v_s`$ to be defined by
$`v_s=v+v_j^n+v_{\text{xc}}.`$ (68)
By definition, or using Eqs. (60), (61), and (64), $`\phi _1`$ is a one-particle density matrix that minimizes our energy functional:
$`_n=E_v[\phi _1],`$ (69)
and, therefore, $`\phi _1`$ also satisfies Eq. (59):
$`{\displaystyle \underset{i=1}{\overset{N}{}}}\widehat{}_{\phi _1}(𝐫_i)|\phi _1=2\left({\displaystyle \underset{w}{}}ϵ_w\right)|\phi _1.`$ (70)
Hence, it follows from Eqs. (65) and (70) that $`|\phi _1`$ is an eigenstate of two different noninteracting Hamiltonians. By comparing Eq. (62) and (67) with $`\varrho _1^{}=\phi _1`$, we see that the two operators, $`\widehat{}_{\phi _1}`$ and $`\widehat{f}`$, are identical, except that $`\widehat{}_{\phi _1}`$ contains the nonlocal operator $`\widehat{w}_{\phi _1}`$ and $`\widehat{f}`$ contains the local potential $`v_{\text{xc}}`$. Furthermore, the occupied orbitals from Eq. (62) and (67) with $`\varrho _1^{}=\phi _1`$ may differ by a unitary transformation, but they yield the same one-particle density matrix:
$`\phi _1(𝐫_1,𝐫_2)=2{\displaystyle \underset{w\phi _1}{}}\chi _w(𝐫_1)\chi _w^{}(𝐫_2)=2{\displaystyle \underset{w\phi _1}{}}\varphi _w(𝐫_1)\varphi _w^{}(𝐫_2).`$ (71)
Using the approach by Sala and Görling Sala and Görling (2001), and Eqs. (65), (70), (62) and (67), but permitting the orbitals to be complex, it is readily demonstrated that $`v_{\text{xc}}`$ is given by
$`v_{\text{xc}}(𝐫)={\displaystyle \frac{1}{2n(𝐫)}}{\displaystyle }d𝐫_1[2w(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1)\phi _1(𝐫,𝐫_1){\displaystyle }d𝐫_2\phi _1(𝐫_2,𝐫)w(𝐫_1,𝐫_2)`$ (72)
$`+\phi _1(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1)v_{\text{xc}}(𝐫_1)].`$
By substituting $`v_{\text{xc}}`$ repeatedly on the right side we can obtain an expansion for $`v_{\text{xc}}`$:
$`v_{\text{xc}}(𝐫)={\displaystyle \frac{1}{2n(𝐫)}}[2w(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1)\phi _1(𝐫,𝐫_1)\phi _1(𝐫_2,𝐫)w(𝐫_1,𝐫_2)`$
$`+\phi _1(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1){\displaystyle \frac{1}{n(𝐫_1)}}\{w(𝐫_2,𝐫_1)\phi _1(𝐫_1,𝐫_2){\displaystyle \frac{1}{2}}\phi _1(𝐫_1,𝐫_2)\phi _1(𝐫_3,𝐫_1)w(𝐫_2,𝐫_3)\}`$
$`+\phi _1(𝐫_1,𝐫)\phi _1(𝐫,𝐫_1){\displaystyle \frac{1}{2n(𝐫_1)}}\phi _1(𝐫_2,𝐫_1)\phi _1(𝐫_1,𝐫_2){\displaystyle \frac{1}{n(𝐫_2)}}w(𝐫_3,𝐫_2)\phi _1(𝐫_2,𝐫_3)+\mathrm{}],`$ (73)
where there are integrations over the dummy variables $`𝐫_1`$, $`𝐫_2`$ and $`𝐫_3`$. The leading term of Eq. (73) is the Slater potential Slater (1951); Harbola and Sahni (1993); Hirata et al. (2001); this term also appears within the Krieger–Li–Iafrate (KLI) approximation of the optimized potential method Fiolhais et al. (2003); Krieger et al. (1992); Li et al. (1993); Hirata et al. (2001).
The orbitals $`\varphi _w`$ satisfying Eq. (67) are the Kohn–Sham orbitals Kohn and Sham (1965); $`|\phi _1`$ is the Kohn–Sham noninteracting state. However, $`\widehat{f}`$ differs from the Kohn–Sham operator, since, in addition to depending explicitly $`\phi _1`$, instead of $`n`$, $`\widehat{f}`$ depends explicitly on the external potential $`v`$ from the interacting Hamiltonian $`\widehat{H}_v`$. Furthermore, the external-correlation operator $`\widehat{v}_{\text{ec}}^{\rho _1}`$ does not appear in Kohn–Sham formalism. In addition, unlike the original Kohn–Sham approach Kohn and Sham (1965), the $`N`$-representability problem does not arise, nor the need to introduce a constraint-search definition Levy (1978, 1979, 1982); Levy and Perdew (1985) to avoid this problem.
In our derivation we have assumed that $`|\phi _1`$ is a ground state solution of Eq. (65). However, the results may also be valid if $`|\phi _1`$ is an excited state solution, since the Sala and Görling approach may also be valid in this case.
## VI Conversion of the one-particle density-matrix functionals into density functionals
For noninteracting states, the wave function is determined by the one-particle density matrix. For certain closed-shell determinantal states, we can write $`\rho _1[\rho _s]`$, where this functional includes all densities that are noninteracting $`v`$-representable, but it is also defined for all $`N`$-representable densities. Using the constraint search approach Levy (1978, 1979, 1982); Levy and Perdew (1985), for a given density, say $`\rho ^{}`$, the functional $`\rho _1[\rho ^{}]`$ yields the one-particle density matrix that minimizes the expectation value of the kinetic energy:
$`\text{Min}\text{}\rho _1|\widehat{T}|\rho _1=\rho _1\left[\rho ^{}\right]|\widehat{T}|\rho _1\left[\rho ^{}\right],`$ (74)
where the search is over all determinantal states that have a density of $`\rho ^{}`$.
Substituting $`\rho _1[\rho ]`$ into $`E_{\text{co}}`$ of Eq. (IV) gives
$`E_v[\rho _1]={\displaystyle 𝑑𝐫_1\left[\frac{1}{2}_1^2\rho _1(𝐫_1,𝐫_2)\right]_{𝐫_2=𝐫_1}}+{\displaystyle 𝑑𝐫v(𝐫)\rho _s(𝐫)}`$
$`+E_J[\rho _s]+E_\text{x}[\rho ]+E_{\text{co}}[\rho ,v]+{\displaystyle 𝑑𝐫v(𝐫)\stackrel{~}{\rho }_c(𝐫)},\rho \rho _1,`$ (75)
where, using $`\rho _1[\rho ]`$, the last term is also a functional of $`v`$ and $`\rho `$. This equation differs from the Kohn–Shan density functional, since the correlation-energy functional depends on the external potential $`v`$, and the last term does not appear in the Kohn–Sham approach. However, mathematically speaking, the minimization of Eq. (VI) follows the same procedure as in the Kohn–Sham method, yielding
$`\widehat{f}\varphi _w=\left(\frac{1}{2}^2+v+v_j^n+v_\text{x}^n+v_{\text{co}}^n+v_{\text{ec}}^n\right)\varphi _w=ϵ_w\varphi _w,\varphi _w\phi _1,`$ (76)
where the local potentials are given by
$`v_\text{x}^\rho (𝐫)`$ $`=`$ $`{\displaystyle \frac{\delta E_\text{x}[\rho ,v]}{\delta \rho (𝐫)}},`$ (77)
$`v_{\text{co}}^\rho (𝐫)`$ $`=`$ $`{\displaystyle \frac{\delta E_{\text{co}}[\rho ,v]}{\delta \rho (𝐫)}},`$ (78)
$`v_{\text{ec}}^\rho (𝐫)`$ $`=`$ $`{\displaystyle \frac{\delta \left(𝑑𝐫_1v(𝐫_1)\stackrel{~}{\rho }_c(𝐫_1)\right)}{\delta \rho (𝐫)}}.`$ (79)
Assuming the density $`n`$ from the interacting state is noninteracting $`v`$-representable, we have
$`E_v[n]=_n,n\text{ is noninteracting }v\text{-representable}.`$ (80)
Note that Eq. (VI) is a valid energy functional only when the one-particle density matrix that enters the first term is the same one generated by the functional $`\rho _1[\rho ]`$; this is the case, at least when $`\rho `$ is non-interacting $`v`$-representable.
## Appendix A A possible choice for $`g_{\rho _c}`$
The electron-electron repulsion operator is spin-free and can be written as
$`\widehat{V}_{\text{ee}}={\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}(ij|r_{12}^1|kl)\widehat{E}_{ijkl},`$ (81)
where the two-electron integral is written in chemist’s notation Szabo and N. S. Ostlund (1982) and the two-electron spin-adapted excitation-operator is given by Eq. (8). This operator can also be written as
$`\widehat{V}_{\text{ee}}={\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\widehat{\mathrm{\Gamma }}_2(𝐫_2,𝐫_1)},`$ (82)
where the pair-function operator is given by
$`\widehat{\mathrm{\Gamma }}_2(𝐫_2,𝐫_1)={\displaystyle \frac{1}{2}}{\displaystyle \underset{ijkl}{}}\chi _j(𝐫_1)\chi _i^{}(𝐫_1)\chi _l(𝐫_2)\chi _k^{}(𝐫_2)\widehat{E}_{ijkl},`$ (83)
and this operator yields the diagonal elements of the spinless two-particle density matrix as the expectation value. Writing this operator in normal-ordered form Čížek (1966, 1969); Lindgren and Morrison (1986); Paldus and Čížek (1975) with respect to the vacuum state $`|\rho _1`$, we have
$`\widehat{V}_{\text{ee}}={\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho _2(𝐫_2,𝐫_1)_{\rho _1}}+{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho _s(𝐫_2)\widehat{\mathrm{\Gamma }}(𝐫_1)_{\rho _s}}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho _1(𝐫_2,𝐫_1)\widehat{\mathrm{\Gamma }}(𝐫_1,𝐫_2)_{\rho _1}}+{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\widehat{\mathrm{\Gamma }}_2(𝐫_2,𝐫_1)_{\rho _1}},`$ (84)
where, examining each term in turn, from the first term we have
$`\rho _2(𝐫_2,𝐫_1)_{\rho _1}={\displaystyle \frac{1}{2}}\rho _s(𝐫_2)\rho _s(𝐫_1){\displaystyle \frac{1}{4}}\rho _1(𝐫_2,𝐫_1)\rho _1(𝐫_1,𝐫_2),`$ (85)
and this function is the diagonal elements of the spinless second-order density matrix of the determinantal state $`|\rho _1`$. From the second term, we have
$`\widehat{\mathrm{\Gamma }}(𝐫)_{\rho _s}={\displaystyle \underset{ij}{}}\chi _j(𝐫)\chi _i^{}(𝐫)\{\widehat{E}_{ij}\}_{\rho _1},\rho _1\rho _s,`$ (86)
and this operator is the one-body portion of the density operator, where the density operator is given by
$`\widehat{\mathrm{\Gamma }}(𝐫)={\displaystyle \underset{ij}{}}\chi _j(𝐫)\chi _i^{}(𝐫)\widehat{E}_{ij}.`$ (87)
Note that we can write
$`\widehat{\mathrm{\Gamma }}(𝐫)_{\rho _s}=\widehat{\mathrm{\Gamma }}(𝐫)\rho _s(𝐫),`$ (88)
indicating that $`\widehat{\mathrm{\Gamma }}(𝐫)_{\rho _s}`$ is determined by $`\rho _s`$ and not by $`\rho _1`$; two different one-particle density matrices that yield the same density have the same $`\widehat{\mathrm{\Gamma }}(𝐫)_{\rho _s}`$.
Returning to Eq. (A), from the third term we have
$`\widehat{\mathrm{\Gamma }}(𝐫_1,𝐫_2)_{\rho _1}={\displaystyle \underset{ij}{}}\chi _j(𝐫_1)\chi _i^{}(𝐫_2)\{\widehat{E}_{ij}\}_{\rho _1},`$ (89)
and this operator is the one-body portion of the one-particle density-matrix operator, given by
$`\widehat{\mathrm{\Gamma }}(𝐫_1,𝐫_2)={\displaystyle \underset{ij}{}}\chi _j(𝐫_1)\chi _i^{}(𝐫_2)\widehat{E}_{ij}=\rho _1(𝐫_1,𝐫_2)+\widehat{\mathrm{\Gamma }}(𝐫_1,𝐫_2)_{\rho _1}.`$ (90)
And from the last term, we have
$`\widehat{\mathrm{\Gamma }}_2(𝐫_2,𝐫_1)_{\rho _1}={\displaystyle \frac{1}{2}}{\displaystyle \underset{ijkl}{}}\chi _j(𝐫_1)\chi _i^{}(𝐫_1)\chi _l(𝐫_2)\chi _k^{}(𝐫_2)\{\widehat{E}_{ijkl}\}_{\rho _1},`$ (91)
and this operator is the two-body portion of the pair-function operator, Eq. (83).
To obtain a slight modification of $`\widehat{V}_{\text{ee}}`$, we replace the determinantal state density $`\rho _s`$, that appears in Eq. (A), with the exact density $`n`$, giving
$`\widehat{V}_{\text{ee}}^{\rho _1}={\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho _2(𝐫_2,𝐫_1)_{\rho _1}}+{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1n(𝐫_2)\widehat{\mathrm{\Gamma }}(𝐫_1)_{\rho _s}}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho _1(𝐫_2,𝐫_1)\widehat{\mathrm{\Gamma }}(𝐫_1,𝐫_2)_{\rho _1}}+{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\widehat{\mathrm{\Gamma }}_2(𝐫_2,𝐫_1)_{\rho _1}},`$ (92)
and this operator can also be written as
$`\widehat{V}_{\text{ee}}^{\rho _1}=\widehat{V}_{\text{ee}}+{\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho _c(𝐫_2)\left(\widehat{\mathrm{\Gamma }}(𝐫_1)\rho _s(𝐫_1)\right)},`$ (93)
Replacing $`\widehat{V}_{\text{ee}}`$ by $`\widehat{V}_{\text{ee}}^{\rho _1}`$ within the Hamiltonian operator, we have obtain a trial Hamiltonian:
$`\widehat{H}_{v\rho _1}=\widehat{H}_v+\lambda {\displaystyle 𝑑𝐫_1𝑑𝐫_2r_{12}^1\rho _c(𝐫_2)\left(\widehat{\mathrm{\Gamma }}(𝐫_1)\rho _s(𝐫_1)\right)},`$ (94)
where $`\lambda `$ is unity, but it can be permitted to be any constant value. Comparing this equation with Eq. (26) yields Eq. (30).
## Appendix B Energy Functional using Intermediate Normalization
Using Eq. (26), our energy functional $`E_v`$, Eq. (21), can be also be written as
$`E_v[\rho _1]={\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\widehat{H}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}{\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}}}\lambda {\displaystyle 𝑑𝐫g_{\rho _c}(𝐫)\left(\widehat{\mathrm{\Gamma }}(𝐫)\rho _s(𝐫)\right)}.`$ (95)
By requiring our trial wave functions to satisfy intermediate normalization,
$`\rho _1|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}=1,`$ (96)
we have
$`E_v[\rho _1]=\rho _1|\widehat{H}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}\lambda {\displaystyle 𝑑𝐫g_{\rho _c}(𝐫)\left(\widehat{\mathrm{\Gamma }}(𝐫)\rho _s(𝐫)\right)}.`$ (97)
This form suggest the following partitioning:
$`E_v[\rho _1]=_1[\rho _1,v]+\stackrel{~}{E}_{\text{co}}[\rho _1,v]\lambda {\displaystyle 𝑑𝐫g_{\rho _c}(𝐫)\left(\widehat{\mathrm{\Gamma }}(𝐫)\rho _s(𝐫)\right)},`$ (98)
where $`\stackrel{~}{E}_{\text{co}}`$ is the correlation-energy (functional) of the trial wave function:
$`\stackrel{~}{E}_{\text{co}}[\rho _1,v]=\rho _1|\widehat{H}_{v\rho _1}|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}^Q,`$ (99)
and the correlation function $`\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}^Q`$ is defined by
$`|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}=|\rho _1+|\stackrel{~}{\mathrm{\Psi }}_{v\rho _1}^Q.`$ (100)
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# Dependence of the superfluidity criterion on the resonance between one-particle (Bogoliubov) and two-particle series
## Abstract
We study how the superfluidity depends on the resonance between one- and two-particle series. The frequency of the spectrum of two-particle solutions in an interval is calculated.
We consider a system of $`N`$ identical bosons of mass $`m`$ that lie on a three-dimensional torus $`𝐓`$ with the side lengths $`L_1`$, $`L_2`$, and $`L_2`$. We assume that the bosons mutually interact and that the interaction potential has the form
$$V\left(N^{1/3}(xy)\right),$$
(1)
where $`V(\xi )`$ is an even smooth function rapidly decreasing at infinity and $`x,y`$ are the boson coordinates on $`𝐓`$. The interaction potential (1) depends on $`N`$ such that the interaction radius decreases as the particle number $`N`$ increases, and the average number of particles with which a single particle interacts remains constant in this case.
In the ultrasecond quantization over pairs, an ultrasecond-quantized operator $`\overline{\widehat{H}}`$, whose explicit form was previously presented, corresponds to the boson system in question. As stated, this ultrasecond-quantized operator satisfies identity
$$\overline{\widehat{H}}=\overline{\widehat{E}}\widehat{A},$$
(2)
where $`\overline{\widehat{E}}`$ is the ultrasecond-quantized unit operator and $`\widehat{A}`$ is an operator nonuniquely chosen in the ultrasecond quantization space. Further, we assume that the operator $`\widehat{A}`$ has the form
$`\widehat{A}={\displaystyle 𝑑x𝑑y\widehat{B}^+(x,y)\left(\frac{\mathrm{}^2}{2m}(\mathrm{\Delta }_x+\mathrm{\Delta }_y)+V\left(N^{1/3}(xy)\right)\right)\widehat{B}^{}(x,y)}+`$
$`+2{\displaystyle 𝑑x𝑑y𝑑x^{}𝑑y^{}V\left(N^{1/3}(xy)\right)\widehat{B}^+(x,y)\widehat{B}^+(x^{},y^{})\widehat{B}^{}(x,x^{})\widehat{B}^{}(y,y^{})},`$ (3)
where $`\widehat{B}^+(x,y)`$ and $`\widehat{B}^{}(x,y)`$ are the corresponding bosonic creation and annihilation operators for a particle pair in the Fock ultrasecond quantization space. By identity (2), an asymptotic expression for operator (3) can be found for establishing an asymptotic representation for the spectrum of the boson system under consideration in the limit as $`N\mathrm{}`$.
Because the function (1) times $`N`$ in a sense converges to the Dirac function in the limit as $`N\mathrm{}`$, the operator (3) involves the small parameter $`1/N`$ in the second term for this limiting case. This means that to find the asymptotic expressions for eigenfunctions and eigenvalues of the operator $`\widehat{A}`$, we can use the semiclassical methods developed in . The asymptotic expressions for eigenfunctions and eigenvalues are determined by the symbol of the operator (3), and this symbol is called the true symbol for the ultrasecond-quantized problem. The true symbol corresponding to the operator (3) is the following functional defined for a pair of functions $`\mathrm{\Phi }^+(x,y)`$ and $`\mathrm{\Phi }(x,y)`$:
$`[\mathrm{\Phi }^+(),\mathrm{\Phi }()]={\displaystyle 𝑑x𝑑y\mathrm{\Phi }^+(x,y)\left(\frac{\mathrm{}^2}{2m}(\mathrm{\Delta }_x+\mathrm{\Delta }_y)\right)\mathrm{\Phi }(x,y)}+`$
$`+2{\displaystyle 𝑑x𝑑y𝑑x^{}𝑑y^{}\left(NV\left(N^{1/3}(xy)\right)\right)\mathrm{\Phi }^+(x,y)\mathrm{\Phi }^+(x^{},y^{})\mathrm{\Phi }(x,x^{})\mathrm{\Phi }(y,y^{})}.`$ (4)
The conservation of the number of particles in the system for the functions $`\mathrm{\Phi }^+(x,y)`$ and $`\mathrm{\Phi }(x,y)`$ implies the condition
$$𝑑x𝑑y\mathrm{\Phi }^+(x,y)\mathrm{\Phi }(x,y)=\frac{1}{2}.$$
(5)
According to the asymptotic methods in , in the limit as $`N\mathrm{}`$, an asymptotic series of eigenfunctions and eigenvalues of the operator (3) corresponds to each solution of the equations
$$\mathrm{\Omega }\mathrm{\Phi }(x,y)=\frac{\delta }{\delta \mathrm{\Phi }^+(x,y)},\mathrm{\Omega }\mathrm{\Phi }^+(x,y)=\frac{\delta }{\delta \mathrm{\Phi }(x,y)}$$
(6)
that also satisfy condition (5). It follows from the explicit form of the true symbol (4) that system (6) can be written as
$`\mathrm{\Omega }\mathrm{\Phi }(x,y)={\displaystyle \frac{\mathrm{}^2}{2m}}(\mathrm{\Delta }_x+\mathrm{\Delta }_y)\mathrm{\Phi }(x,y)+`$
$`+{\displaystyle 𝑑x^{}𝑑y^{}\left(NV\left(N^{1/3}(xy)\right)+NV\left(N^{1/3}(x^{}y^{})\right)\right)\mathrm{\Phi }^+(x^{},y^{})\mathrm{\Phi }(x,x^{})\mathrm{\Phi }(y,y^{})},`$
$`\mathrm{\Omega }\mathrm{\Phi }^+(x,y)={\displaystyle \frac{\mathrm{}^2}{2m}}(\mathrm{\Delta }_x+\mathrm{\Delta }_y)\mathrm{\Phi }^+(x,y)+`$
$`+{\displaystyle 𝑑x^{}𝑑y^{}\left(NV\left(N^{1/3}(xx^{})\right)+NV\left(N^{1/3}(yy^{})\right)\right)\mathrm{\Phi }(x^{},y^{})\mathrm{\Phi }^+(x,x^{})\mathrm{\Phi }^+(y,y^{})}.`$ (7)
In the limit as $`N\mathrm{}`$, system (7) supplemented with the condition (5) has the family of solutions
$`\mathrm{\Phi }_k^+(x,y)={\displaystyle \frac{1}{L_1L_2^2}}\mathrm{cos}\left(k(xy)\right),`$
$`\mathrm{\Phi }_k(x,y)={\displaystyle \frac{1}{L_1L_2^2}}{\displaystyle \underset{l}{}}\phi _{k,l}\mathrm{exp}\left(il(xy)\right),`$ (8)
where $`k`$ and $`l`$ are three-dimensional vectors of the form
$$2\pi (\frac{n_1}{L_1},\frac{n_2}{L_2},\frac{n_3}{L_2}),$$
the numbers $`n_1`$, $`n_2`$, and $`n_3`$ are integers, the functions $`\phi _{k,l}`$ in formula (2) have the form
$$\phi _{k,l}=\frac{1}{2V_0}\left(\frac{\mathrm{}^2}{2m}(k^2l^2)+V_0\pm \sqrt{\left(\frac{\mathrm{}^2}{2m}(k^2l^2)+V_0\right)^2V_0^2}\right),$$
(9)
with the plus sign for $`l^2>k^2`$ and the minus sign for $`l^2<k^2`$, and $`V_0`$ here and hereafter denotes the expression
$$V_0=\frac{1}{L_1L_2^2}𝑑xV(x),$$
(10)
in which the integral is taken over the space $`𝐑^3`$. The vector $`k`$ in (8) plays the role of a parameter indexing different solutions of system (7) with conditions (5). The solutions (8) are standing waves associated with the series without fluidity.
The leading asymptotic term for the eigenvalues in the series corresponding to solution (8) is equal to the value of the product of the symbol (4) on the functions (8) times $`N`$,
$$E_k=N\left(\frac{\mathrm{}^2k^2}{2m}+\frac{V_0}{2}\right).$$
(11)
The asymptotic expressions for the eigenvalues and eigenfunctions, in particular, the asymptotic terms following (11) are determined by the solutions of system (7) and also by the solutions of the variational system of equations corresponding to (7). The detailed study of solutions of the variational system is here omitted. In what follows, we present the results of this study obtained, in particular, by using the Maple program.
If $`k=0`$, then the asymptotic series related to this solution (8) of system (7), (5) is the Bogoliubov series corresponding to the ground state without fluidity. The quasiparticle spectrum for this series is expressed by the well-known formula
$$\lambda _l=\sqrt{\left(\frac{\mathrm{}^2l^2}{2m}+V_0\right)^2V_0^2}.$$
(12)
If $`k0`$, then the variational system equations corresponding to (7) gives the quasiparticle spectrum
$`\lambda _{1,k,l}=`$
$`=\pm {\displaystyle \frac{\mathrm{}^2}{2m}}\sqrt{k^4+{\displaystyle \frac{l^2}{2}}+{\displaystyle \frac{l_1^2}{2}}k^2l^2k^2l_1^2+{\displaystyle \frac{1}{2}}(l^2+l_1^22k^2)\sqrt{(l^2l_1^2)^2+\left({\displaystyle \frac{4mV_0}{\mathrm{}^2}}\right)^2}},`$
$`\lambda _{2,k,l}=`$
$`=\pm {\displaystyle \frac{\mathrm{}^2}{2m}}\sqrt{k^4+{\displaystyle \frac{l^2}{2}}+{\displaystyle \frac{l_1^2}{2}}k^2l^2k^2l_1^2{\displaystyle \frac{1}{2}}(l^2+l_1^22k^2)\sqrt{(l^2l_1^2)^2+\left({\displaystyle \frac{4mV_0}{\mathrm{}^2}}\right)^2}},`$ (13)
where $`l_1=l+2k`$ and $`lk`$. In the formula for $`\lambda _{1,k,l}`$ the plus sign is chosen for $`l^2>k^2`$ and the minus sign for $`l^2<k^2`$; in the formula for $`\lambda _{2,k,l}`$, the plus sign is chosen for $`l_1^2>k^2`$ and the minus sign for $`l_1^2<k^2`$.
The explicit form (13) implies that this quasiparticle spectrum contains some negative values, which means that the series corresponding to the solution (8) with $`k0`$ is nonmetastable. In this connection, we suggest the following explanation for the dependence of the critical velocity on capillary width. In what follows, we assume that $`L_1L_2`$. We consider the Bogoliubov series corresponding to the flow with the velocity $`\mathrm{}k_0/m`$ along the capillary, where $`k_0=2\pi (n_1/L_1,0,0)`$. For the boson system in question, the leading asymptotic term for the eigenvalues in this series is
$$N\left(\frac{\mathrm{}^2k_0^2}{2m}+\frac{V_0}{2}\right).$$
(14)
We now assume that there is a relation between $`L_1`$ and $`L_2`$ such that there exists a vector $`k=2\pi (0,n_2/L_2,n_3/L_2)`$ for which the corresponding value of the term in (11) is exactly equal to (14). This means that there can be resonance between the current states of the Bogoliubov series and the states of the nonmetastable series corresponding to (8). If the value of $`L_1`$ is very large, then resonance is also possible for the case where the value of the term in (11) is close to expression (14) but need not coincide with it. The existence of such a resonance indicates the possibility of a transition from the current state to a nonmetastable state from which the system falls to the lowest energy level, which indicates the loss of superfluidity.
The following versions of series of two-particle solutions are determined by the relation:
$$Q(k_1,k_2,m_1,m_2,m_3,l_1,l_2,l_3)=\frac{l_3^2m_3^2}{m_1^2+m_2^2+k_1^2+k_2^2l_1^2l_2^2},$$
(15)
where $`k_1,k_2,l_1,l_2,l_3,m_1,m_2,m_3`$ are integers. We impose the following conditions on the numbers $`k_1,k_2,l_1,l_2,l_3,m_1,m_2,m_3`$:
$$Q>0,k_1^2+k_2^2>0,$$
(16)
where $`Q`$ is an irreducible fraction, i.e., its numerator and denominator do not have common divisors.
Professor A. A. Karatsuba kindly calculated for me the frequency of the spectrum of two-particle solutions in the interval $`[a,a+h]`$. The asymptotic expression in the limit as $`L\mathrm{}`$ has the following form.
We are given the functions (15) and the function
$$V=\sqrt{l_1^2+l_2^2+\frac{l_3^2}{Q}},$$
(17)
where $`k_1,k_2,m_1,m_2,m_3,l_1,l_2,l_3`$ are nonnegative integers that do not exceed $`L`$; moreover, $`Q>0`$ and $`k_1^2+k_2^2>0`$.
We are also given the numbers $`a1`$, $`0<h1`$.
Let $`K`$ be the number of sets $`k_1,k_2,m_1,m_2m_3,l_1,l_2,l_3`$ satisfying the inequality
$$a<Va+h.$$
(18)
Then $`K`$ is determined by the approximate formula
$$K=\left(\frac{1}{5}+\frac{1}{9}\right)\frac{\pi ^3}{128}L^2((a+h)^6a^6)=\frac{7\pi ^3}{2880}L^2((a+h)^6a^6).$$
(19)
If the distance between the spectrum points is of the order of $`\frac{1}{N}`$, i.e., if it is determined with the desired accuracy, then the Bogoliubov series is in resonance.
This readily shows that $`a4`$, and hence resonance occurs approximately at the points $`8\pi h/mL_2`$. This resonance is braking the fluid, and the superfluidity disappears.
We here present A. A. Karatsuba’s proof of formula (19).
We introduce the new notation
$$r=m_1^2+m_2^2+k_1^2+k_2^2,l=l_1^2+l_2^2,\xi =\frac{m_3}{l_3}0.$$
In this notation, we have
$$Q=l_3^2\frac{1\xi ^2}{rl}>0.$$
Two cases, (I) and (II), are possible:
$$\text{(I)}0\xi <1,l<r;\text{(II)}\xi >1,l>r.$$
The function $`V`$ can be rewritten as
$$V=\sqrt{l+\frac{rl}{1\xi ^2}}=\sqrt{\frac{rl\xi ^2}{1\xi ^2}}$$
and hence (19) becomes
$$a^2<V^2=\frac{rl\xi ^2}{1\xi ^2}(a+h)^2.$$
(20)
Case (I). From (20) we obtain
$`a^2(1\xi ^2)+l\xi ^2<r(a+h)^2(1\xi ^2)+l\xi ^2,`$
$`l<r,0\xi <1.`$ (21)
If in (21) we have the relation
$$a^2(1\xi ^2)+l\xi ^2l,\text{i.e.},a^2l,$$
then system (21) is equivalent to
$`l<r(a+h)^2(1\xi ^2)+l\xi ^2,`$
$`a^2l<(a+h)^2,`$
$`0\xi <1.`$ (22)
But if in (21) we have
$$a^2(1\xi ^2)+l\xi ^2>l,\text{i.e.},l<a^2,$$
then system (21) is equivalent to
$`a^2(1\xi ^2)+l\xi ^2<r(a+h)(1\xi ^2)+l\xi ^2,`$
$`l<a^2,`$
$`0\xi <1.`$ (23)
Let $`K_1`$ be the number of sets satisfying either (22) or (23).
It follows from systems (22) and (23) that
$$l<a^2\text{or}l<(a+h)^2,$$
i.e., we always have $`l_1,l_2<a+ha+1`$, since $`h1`$.
Precisely in the same way we obtain
$$r<(a+h)^2,\text{i.e.},m_1,m_2,k_1,k_2<a+ha+1.$$
Therefore, the six variables $`l_1,l_2,m_1,m_2mk_1,k_2`$ are bounded by the value $`a+1`$, and they play an unimportant role for small $`a`$. The main parameter is $`L`$, $`L+\mathrm{}`$, and this parameter is related only to the variable $`\xi `$,
$$\xi =\frac{m_3}{l_3}<1,0m_3<l_3L.$$
The ranges of these variables are determined automatically by the first two inequalities in systems (22) and (23). In what follows, we present two asymptotic formulas (in the limit as $`x\mathrm{}`$).
The number of sets $`m_1,m_2,k_1,k_2`$ for which $`r=m_1^2+m_2^2+k_1^2+k_2^2x`$ is equal to
$$\frac{1}{24}v_4(\sqrt{x}),$$
where $`v_4(\sqrt{x})`$ is the volume of a four-dimensional sphere of radius $`\sqrt{x}`$, i.e.,
$$v_4(\sqrt{x})=(\sqrt{x})^4\frac{\pi ^{4/2}}{\mathrm{\Gamma }(4/2+1)}=x^2\frac{\pi ^2}{2};$$
(24)
the number of sets $`l_1,l_2`$ for which $`l_1^2+l_2^2x`$ is equal to
$$\frac{1}{4}\pi x.$$
(25)
These are approximate formulas; the larger $`x`$, the more precise they are.
From (24) and (22), (23) we obtain
$`K_1`$ $`={\displaystyle \underset{0\xi <1}{}}\underset{a^2l<(a+h)^2}{{\displaystyle ^{}}}{\displaystyle \frac{\pi ^2}{32}}\{((a+h)^2(1\xi ^2)+l\xi ^2)^2(l)^2\}`$
$`+{\displaystyle \underset{0\xi <1}{}}\underset{l<a^2}{{\displaystyle ^{}}}{\displaystyle \frac{\pi ^2}{32}}\{((a+h)^2(1\xi ^2)+l\xi ^2)^2(a^2(1\xi ^2)+l\xi ^2)^2)\}.`$
The prime on the sum over $`l`$ means that the numbers $`l`$ are regarded with multiplicity taken into account (the number of solutions of the equation is $`l=l_1^2+l_2^2`$, $`l_10`$, $`l_20`$); the multiplicity is equal to $`\pi /4`$ (averaged over $`l`$), i.e.,
$`K_1`$ $`={\displaystyle \underset{0\xi <1}{}}{\displaystyle \underset{a^2l<(a+h)^2}{}}{\displaystyle \frac{\pi ^2}{128}}\{((a+h)^2(1\xi ^2)+l\xi ^2)^2l^2\}`$
$`+{\displaystyle \underset{0\xi <1}{}}{\displaystyle \underset{l<a^2}{}}{\displaystyle \frac{\pi ^2}{128}}\{((a+h)^2(1\xi ^2)+l\xi ^2)^2(a^2(1\xi ^2)+l\xi ^2)^2)\}`$
The sum over $`l`$ (with a good accuracy) is equal to the integral
$`K_1`$ $`={\displaystyle \frac{\pi ^3}{128}}{\displaystyle \underset{0\xi <1}{}}{\displaystyle _{a^2}^{(a+h)^2}}\{((a+h)^2(1\xi ^2)+l\xi ^2)^2l^2\}𝑑l`$
$`+{\displaystyle \frac{\pi ^3}{128}}{\displaystyle \underset{0\xi <1}{}}{\displaystyle _0^{a^2}}\{((a+h)^2(1\xi ^2)+l\xi ^2)^2(a^2(1\xi ^2)+l\xi ^2)^2\}𝑑l`$
$`={\displaystyle \frac{\pi ^3}{128}}((a+h)^6a^6){\displaystyle \underset{0\xi <1}{}}\left({\displaystyle \frac{2}{3}}\xi ^2+{\displaystyle \frac{1}{3}}\xi ^4\right).`$
Since
$$\xi =\frac{m_3}{l_3},0m_3<l_3L,$$
the sum over $`\xi `$ (with a good accuracy) is equal to the integral
$$_0^L_0^l\left(\frac{2}{3}\left(\frac{m}{l}\right)^2+\frac{1}{3}\left(\frac{m}{l}\right)^4\right)𝑑m𝑑l=\frac{L^2}{2}\left(\frac{2}{3}\frac{1}{3}+\frac{1}{15}\right)=\frac{L^2}{5}.$$
Hence we obtain
$$K_1=\frac{\pi ^3}{128}((a+h)^6a^6)\frac{L^2}{5}.$$
Case (II).
$$\xi >1,l>r,\text{i.e.},l_3<m_2,r<l.$$
Let
$$0<\xi _1=\frac{1}{\xi }=\frac{l_3}{m_3}<1,r<l.$$
Then we have the equality
$$V=\sqrt{\frac{rl\xi ^2}{1\xi ^2}}=\sqrt{\frac{lr\xi _1^2}{1\xi _1^2}}$$
and the condition $`a<V(a+h)`$ implies
$$a^2<\frac{lr\xi _1^2}{1\xi _1^2}(a+h)^2$$
or
$`a^2(1\xi _1^2)+r\xi _1^2<l(a+h)^2(1\xi _1^2)+r\xi _1^2,`$
$`r<l,0<\xi _1<1.`$ (26)
It follows from (26) that $`r<(a+h)^2`$. If $`r>a^2(1\xi _1^2)+r\xi _1^2`$, i.e., if $`r>a^2`$, then (26) is replaced by the system
$`rl<l(a+h)^2(1\xi _1^2)+r\xi _1^2,`$
$`a^2<r<(a+h)^2,0<\xi _1<1.`$ (27)
If $`r\mathrm{`}a^2(1\xi _1^2)+r\xi _1^2`$, i.e., if $`ra^2`$, then (26) is replaced by (28):
$`a^2(1\xi _1^2)+r\xi _1^2<l(a+h)^2(1\xi _1^2)+r\xi _1^2,`$
$`ra^2,0<\xi _1<1.`$ (28)
It follows from Case (I) and systems (22) and (23) that $`\xi `$ is replaced by $`\xi _1`$, while $`r`$ and $`l`$ interchange.
Therefore, we use formulas (25), (27), and (28) to find $`K_2`$, which is the number of sets satisfying (27) and (28):
$`K_2`$ $`={\displaystyle \frac{1}{4}}\pi {\displaystyle \underset{0<\xi _1<1}{}}\underset{a^2<r<(a+h)^2}{{\displaystyle ^{}}}((a+h)^2(1\xi _1^2)+r\xi _1^2r)`$
$`+{\displaystyle \frac{1}{4}}\pi {\displaystyle \underset{0<\xi _1<1}{}}\underset{ra^2}{{\displaystyle ^{}}}((a+h)^2a^2)(1\xi _1^2).`$
The prime on the sum over $`r`$ means that the number of solutions $`r`$ is regarded with multiplicity taken into account (the number of solutions of the equation is $`r=m_1^2+m_2^2+k_1^2+k_2^2`$).
The number of solutions of the inequality
$$m_1^2+m_2^2+k_1^2+k_2^2x,m_1,m_2,k_1,k_20,.$$
is approximately equal to
$$\frac{\pi ^2}{32}x^2.$$
We obtain
$`\underset{a^2<r<(a+h)^2}{{\displaystyle ^{}}}r=\underset{a^2<r<(a+h)^2}{{\displaystyle ^{}}}{\displaystyle \underset{0t<r}{}}1={\displaystyle \underset{0t<(a+h)^2}{}}\underset{\mathrm{max}(a^2,t)<r<(a+h)^2}{{\displaystyle ^{}}}1`$
$`={\displaystyle \underset{0ta^2}{}}\underset{a^2<r<(a+h)^2}{{\displaystyle ^{}}}1+{\displaystyle \underset{a^2<t<(a+h)^2}{}}\underset{t<r<(a+h)^2}{{\displaystyle ^{}}}1`$
$`={\displaystyle \frac{\pi ^2}{32}}\left({\displaystyle \underset{0t(a+h)^2}{}}((a+h)^4a^4)+{\displaystyle \underset{a^2<t<(a+h)^2}{}}((a+h)^4t^2)\right)`$
$`={\displaystyle \frac{\pi ^2}{32}}\left((a+h)^4a^2a^6+(a+h)^4((a+h)^2a^2){\displaystyle \frac{1}{3}}((a+h)^6a^6)\right)`$
$`={\displaystyle \frac{\pi ^2}{316}}((a+h)^6a^6);`$
We derive the following formula for $`K_2`$:
$`K_2`$ $`={\displaystyle \frac{\pi }{4}}{\displaystyle \underset{0<\xi _1<1}{}}\{(a+h)^2(1\xi _1^2){\displaystyle \frac{\pi ^2}{32}}(()^4a^4)`$
$`(1\xi _1^2){\displaystyle \frac{\pi ^2}{316}}((a+h)^6a^6)+((a+h)^2a^2)(1\xi _1^2){\displaystyle \frac{\pi ^2}{32}}a^4\}`$
$`={\displaystyle \frac{\pi ^3}{1232}}((a+h)^6a^6){\displaystyle \underset{0<\xi _1<1}{}}(1\xi _1^2).`$
Since $`\xi _1=\frac{l}{m}`$, $`0<l<mL`$, we have
$$\underset{0<\xi _1<1}{}(1\xi _1^2)=_0^L_0^m\left(1\left(\frac{l}{m}\right)^2\right)𝑑l𝑑m=\frac{L^2}{3}$$
and
$$K_2=\frac{\pi ^3}{3236}L^2((a+h)^6a^6).$$
The fact that $`K=K_1+K_2`$ implies the relation
$$K=\left(\frac{1}{5}+\frac{1}{9}\right)\frac{1}{128}\pi ^3L^2((a+h)^6a^6)=\frac{7\pi ^3}{2880}L^2((a+h)^6a^6).$$
Remarks.
1. The formula for $`K`$ is an approximate formula.
2. This formula is sufficiently exact if $`L`$ is large and $`a`$ is not too small but significantly less than $`L`$; for example, $`a=\sqrt{L}`$, $`L+\mathrm{}`$.
3. If $`a`$ is chosen to be not large, for example, $`a=4`$, $`h=\frac{1}{2}`$, or $`a=3`$, $`h=\frac{1}{2}`$, or $`a=5`$, $`h=\frac{1}{2}`$, then we can obtain sufficiently exact formulas for $`L8`$, since, in this case, we can exactly calculate the multiplicities of $`l`$ and the $`l`$ themselves (as sums of two squared numbers). But, in this case, we must perform specific calculations.
Example of calculations for small $`a`$.
Let $`a=4`$, $`h=1`$; $`L>5`$.
In system (22), $`l`$ satisfy the condition
$$16l=l_1^2+l_2^2<25.$$
The pairs $`(l_1,l_2)`$ satisfying this condition are:
$$(0,4),(1,4),(2,4),(3,3),(4,1),(4,1),(4,0),$$
i.e., from (22) we obtain two systems for $`l=16`$, $`l=17`$, and $`l=20`$ and one system for $`l=18`$:
$`16<r259\xi ^2,`$
$`0\xi <1,`$ (29)
$`17<r258\xi ^2,`$
$`0\xi <1,`$ (30)
$`20<r255\xi ^2,`$
$`0\xi <1,`$ (31)
$`18<r257\xi ^2,`$
$`0\xi <1.`$ (32)
To calculate the number of sets $`(m_1,m_2,k_1,k_2)`$ for which $`r=m_1^2+m_2^2+k_1^2`$ satisfies systems (29)–(32), we can either use the computer software (in this case, we must divide $`\xi `$ into intervals, where $`9\xi ^2`$ lies between two neighboring integers from $`0`$ to $`9`$, the same for $`8\xi ^2`$, $`5\xi ^2`$, and $`7\xi ^2`$, and calculate the number of sets directly) or use the formula for the number of integer points in the four-dimensional ball and thus find the points with positive coordinates. This is an approximate formula and if
$$m_1^2+m_2^2+k_1^2x,m_i,k_i0,$$
then this number is equal to
$$\frac{\pi ^2}{32}x^2.$$
(33)
Therefore, for system (22), we obtain the number of sets
$`{\displaystyle \frac{\pi ^2}{32}}{\displaystyle \underset{0\xi <1}{}}\{((259\xi ^2)^216^2)2+((258\xi ^2)^217^2)2`$
$`+((255\xi ^2)^220^2)2+((257\xi ^2)^218^2)\}.`$ (34)
Since the number of such $`l`$ is equal to $`2`$, we use the factor $`2`$. In system (23), we have
$$l=l_1^2+l_2^2<16,$$
and the pairs $`(l_1,l_2)`$ are: $`\underset{¯}{(0,3)}`$, $`\underset{¯}{(0,2)}`$, $`\underset{¯}{(0,1)}`$, $`\underset{¯}{(0,0)}`$, $`\underset{¯}{(1,2)}`$, $`\underset{¯}{(1,3)}`$, $`\underset{¯}{(1,0)}`$, $`\underset{¯}{(2,0)}`$, $`\underset{¯}{(2,1)}`$, $`(2,2)`$, $`(2,3)`$, $`\underset{¯}{(3,0)}`$, $`\underset{¯}{(3,2)}`$, $`(3,2)`$.
| l= | 0 | 1 | 4 | 9 | 5 | 10 | 8 | 13 |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| multiplicity | 1 | 2 | 2 | 2 | 2 | 2 | 1 | 2 |
System (23) for each $`l`$ has the form
$$16(1\xi ^2)+l\xi ^2<r25(1\xi ^2)+l\xi ^2,0\xi <1.$$
Using (33), we obtain the following formula (with multiplicity of $`l`$ taken into account) for the number of sets for system (23):
$`{\displaystyle \frac{\pi ^2}{32}}{\displaystyle \underset{0\xi <1}{}}\{(25(1\xi ^2))^2(16(1\xi ^2))^2`$
$`+((25(1\xi ^2)+\xi ^2)^2(16(1\xi ^2)+\xi ^2)^2)2`$
$`+((25(1\xi ^2)+4\xi ^2)^2(16(1\xi ^2)+4\xi ^2)^2)2`$
$`+((25(1\xi ^2)+9\xi ^2)^2(16(1\xi ^2)+9\xi ^2)^2)2`$
$`+((25(1\xi ^2)+5\xi ^2)^2(16(1\xi ^2)+5\xi ^2)^2)2`$
$`+((25(1\xi ^2)+10\xi ^2)^2(16(1\xi ^2)+10\xi ^2)^2)2`$
$`+((25(1\xi ^2)+8\xi ^2)^2(16(1\xi ^2)+8\xi ^2)^2)1`$
$`+((25(1\xi ^2)+13\xi ^2)^2(16(1\xi ^2)+13\xi ^2)^2)2\}`$
$`={\displaystyle \frac{\pi ^2}{32}}{\displaystyle \underset{0\xi <1}{}}\{(25(1\xi ^2))^2(16(1\xi ^2))^2`$
$`+((2524\xi ^2)^2(1615\xi ^2)^2)2`$
$`+((2521\xi ^2)^2(1612\xi ^2)^2)2`$
$`+((2516\xi ^2)^2(167\xi ^2)^2)2`$
$`+((2520\xi ^2)^2(1611\xi ^2)^2)2`$
$`+((2515\xi ^2)^2(166\xi ^2)^2)2`$
$`+((2517\xi ^2)^2(168\xi ^2)^2)1`$
$`+((2512\xi ^2)^2(163\xi ^2)^2)2\}.`$ (35)
Since $`\xi =\frac{m_3}{l_3}`$, $`0m_3l_3`$, $`1l_3L`$, the number of sets corresponding to Case (I) is either calculated directly according to (34)–(35) or is approximately replaced by an integral of the form
$$_0^L𝑑l_0^l𝑑m\mathrm{}$$
equal to
$$\frac{\pi ^2}{64}L^2\mu ,$$
where
$`\mu `$ $`={\displaystyle _0^1}\{2((259u^2)^216^2)+2((258u^2)^217^2)`$
$`+2((255u^2)^220^2)+((257u^2)^218^2)`$
$`+(25(1u^2))^2(16(1u^2))^2+2((2524u^2)^2(1615u^2)^2)`$
$`+2((2521u^2)^2(1612u^2)^2)+2((2516u^2)^2(167u^2)^2)`$
$`+2((2520u^2)^2(1611u^2)^2)+2((2515u^2)^2(166u^2)^2)`$
$`+1((2517u^2)^2(168u^2)^2)+2((2512u^2)^2(163u^2)^2)\}du.`$
Thus, we have
$$K_1=\frac{\pi ^2}{64}L^2\mu .$$
Case (II), i.e., the calculation of $`K_2`$ is the same; only the multiplicity of $`r`$, $`16<r<25`$, $`r=m_1^2+m_2^2+k_1^2+k_2^2`$, is calculated either by an exhaustive search of the sets $`(m_1,m_2,k_1,k_2)`$ (also for $`r<16`$) or by the formula for the number of integer points in the four-dimensional ball; the number of sets for $`lr`$ from (27) and (28) is calculated approximately by the formula
$$\frac{\pi }{4}x.$$
Therefore, if $`\kappa (r)`$ is the multiplicity of $`r`$, then
$`K_2`$ $`={\displaystyle \underset{0<\xi _1<1}{}}\left\{{\displaystyle \underset{16<r<25}{}}\kappa (r){\displaystyle \frac{\pi }{4}}(25r)(1\xi _1^2)+{\displaystyle \underset{r<16}{}}{\displaystyle \frac{\pi }{4}}\kappa (r)(2516)(1\xi _1^2)\right\}=`$
$`={\displaystyle \frac{L^2}{2}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\pi }{4}}\left({\displaystyle \underset{16<r<25}{}}\kappa (r)(25r)+9{\displaystyle \underset{r<16}{}}\kappa (r)\right).`$
But we have already calculated that
$`{\displaystyle \underset{r<16}{}}\kappa (r)`$ $`={\displaystyle \frac{\pi ^2}{32}}4^4;`$
$`{\displaystyle \underset{16<r<25}{}}\kappa (r)`$ $`={\displaystyle \frac{\pi ^2}{32}}(5^44^4);`$
$`{\displaystyle \underset{16<r<25}{}}\kappa (r)r`$ $`={\displaystyle \frac{\pi ^2}{48}}(5^64^6).`$
Hence we have
$`K_2`$ $`={\displaystyle \frac{\pi }{12}}L^2\left(25{\displaystyle \frac{\pi ^2}{32}}(5^44^4){\displaystyle \frac{\pi ^2}{48}}(5^64^6)+9{\displaystyle \frac{\pi ^2}{32}}4^4\right),`$
$`K`$ $`=K_1+K_2=cL^2.`$
All these formulas are approximate.
The number $`\mu =_0^1\mathrm{}𝑑m`$ and the constant $`c`$ in the formula for $`K`$ have not been calculated, but these calculations are very simple.
The author wishes to express his deep gratitude to Professor A. A. Karatsuba and to G. V. Koval’ for their help in preparing this paper.
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# Ground State Baryons in 𝑈̃(12) Scheme
## 1 Introduction
The spectroscopy of light-quark $`qqq`$ baryons is longstanding problem of hadron physics. The non-relativistic quark model(NRQM) successfully explain the properties of ground state 56-multiplet of $`SU(6)`$ spin-flavor symmetry. On the other hand, NRQM predicts the negative-parity states as the next low-lying states from orbital excitation, while the experiments show the clear evidence of the Roper resonance $`N(1440)`$ of the second nucleon state with positive-parity. The situation is similar for $`\mathrm{\Delta }`$, $`\mathrm{\Lambda }`$ and $`\mathrm{\Sigma }`$ systems. The positive-parity $`\mathrm{\Delta }(1600)`$, $`\mathrm{\Lambda }(1600)`$ and $`\mathrm{\Sigma }(1660)`$ have too light masses to be naturally assigned as radially excited states in NRQM. The mass of the negative parity $`\mathrm{\Lambda }(1405)`$ is also too light to be assigned as the first excited 70-multiplet in NRQM.
Recently, we have proposed a covariant level-classification scheme of hadrons based on $`\stackrel{~}{U}(12)`$ group. In this scheme, the squared-mass spectra of hadrons including light constituent quarks are classified as the representation of $`\stackrel{~}{U}(12)`$ spin-flavor group. In this scheme we have introduced the expansion bases of spinor wave functions(WF) of composite hadrons. Each spinor index corresponding to light quark freedom is expanded by free Dirac spinors $`u_{r,s}(v_\mu )`$.
$`u_{+,s}(v_\mu )`$ $`=`$ $`\left(\begin{array}{c}ch\theta \chi ^{(s)}\\ sh\theta 𝒏\mathbf{}𝝈\chi ^{(s)}\end{array}\right),u_{,s}(v_\mu )=\left(\begin{array}{c}sh\theta 𝒏\mathbf{}𝝈\chi ^{(s)}\\ ch\theta \chi ^{(s)}\end{array}\right),`$ (5)
where $`ch\theta ,sh\theta =\sqrt{(\omega \pm 1)/2}`$ and the $`v_\mu P_\mu /M=(𝒏\omega _3;i\omega )_\mu `$ is the four velocity of the relevant hadron. This method is invented for keeping the Lorentz-covariance of the composite system. (Concerning the $`\stackrel{~}{U}(12)`$ scheme and its group theoretical arguments, see our ref..) We should note both $`u_+`$ and $`u_{}`$ is necessary for expansion bases of quark spinor index. They are called ur-citon spinors for historical reason.
The $`u_{+,s}`$ corresponds to the ordinary spinor freedom appearing in NRQM, while the $`u_{,s}`$ represents the relativistic effect. In the $`\stackrel{~}{U}(12)`$-classification scheme, the $`u_{,s}`$ is supposed to appear as new degrees of freedom for light constituents, being independent of $`u_{+,s}`$. The freedom corresponding $`r`$ index of $`u_{r,s}`$ is called $`\rho `$-spin, while the ordinary Pauli-spin freedom described by $`\chi ^{(s)}`$ is called $`\sigma `$-spin, where $`\rho \times \sigma `$ corresponds to the $`\rho \sigma `$ decomposition of Dirac $`\gamma `$ matrices. The index $`s`$ of $`u_{r,s}`$ represents the eigen value of $`\sigma _3`$, while the $`r`$ does the eigen value of $`\rho _3`$ at the hadron rest frame, where $`v_\mu =v_{0\mu }=(\mathrm{𝟎};i)_\mu `$.
Because of this extra $`SU(2)`$ spin freedom, the $`SU(6)_{SF}`$ is extended to $`\stackrel{~}{U}(12)`$, or more precisely $`U(12)_{\mathrm{stat}}`$ at the hadron rest frame, as $`U(12)_{\mathrm{stat}}SU(2)_\rho \times SU(2)_\sigma \times SU(3)_F`$ (see, ref.). The ground-state $`qqq`$ baryons and anti-baryons are assigned as the completely symmetric $`(\mathrm{𝟏𝟐}\times \mathrm{𝟏𝟐}\times \mathrm{𝟏𝟐})_{sym}=\mathrm{𝟑𝟔𝟒}`$ representation of $`U(12)_{\mathrm{stat}}`$ or $`\stackrel{~}{U}(12)`$. The corresponding flavor-spinor WFs $`\mathrm{\Phi }_{ABC}(v_\mu )`$ are represented by the direct product of flavor and spinor WFs as
$`\mathrm{\Phi }_{ABC}(v)`$ $``$ $`|F_{abc}u_\alpha (v)u_\beta (v)u_\gamma (v),`$ (6)
where $`A=(a,\alpha )`$ etc. denote the (flavor,spinor) indices. $`|F(u(v))`$ represents the flavor(spinor) WF. At the rest frame of hadron, the $`u_\alpha (v)`$’s are decomposed as $`\rho `$ and $`\sigma `$ spin WFs denoted as $`|\rho ,\rho _3/2`$ and $`|\sigma _3/2`$, respectively. The ground-state baryons in $`\stackrel{~}{U}(12)`$ are classified as $`\mathrm{𝟑𝟔𝟒}/2=\mathrm{𝟓𝟔}_E+\mathrm{𝟕𝟎}_G+\mathrm{𝟓𝟔}_F`$ in terms of $`SU(6)_{SF}`$. The explicit forms of WFs are given in Table 1.
The $`\mathrm{𝟓𝟔}_E`$ has all the three spinor indices with positive $`\rho _3`$ and $`\rho `$-spin WF is given by the completely symmetric $`|\rho ,\frac{3}{2}_S`$. This freedom corresponds to the ones appearing in NRQM, and the corresponding states are called Pauli-states(or Paulons) . The $`\mathrm{𝟓𝟔}_E`$ includes the $`N(939)`$-octet and $`\mathrm{\Delta }(1232)`$-decouplet. While $`\mathrm{𝟕𝟎}_G(\mathrm{𝟓𝟔}_F)`$ have one index(two indices) with negative $`\rho _3`$ and are described by the $`\rho `$-spin WFs $`|\rho ,\frac{1}{2}_{\alpha (\beta )}`$ ($`|\rho ,\frac{1}{2}_S`$). These states are out of the conventional non-relativistic $`SU(6)_{SF}`$ framework, and are called chiral states(or chiralons), since the $`u_{+,s}`$ changes $`u_{,s}`$ by multiplying the chiral matrix $`\gamma _5=\rho _1`$ as $`\rho _1u_{+,s}=u_{,s}`$. We predict, in addition to the ordinary $`\mathrm{𝟓𝟔}_E`$, the existence of the extra $`\mathrm{𝟓𝟔}_F`$ with positive-parity and $`\mathrm{𝟕𝟎}_G`$ with negative-parity in lower mass region since they are $`S`$-wave states. There is an interesting possibility that all the problematic resonances mentioned above are explained as chiral states in $`\stackrel{~}{U}(12)`$ scheme. In order to check this possibility, we consider the mass spectra and decay properties of $`\mathrm{𝟓𝟔}_F`$ and $`\mathrm{𝟕𝟎}_G`$.
## 2 Mass Spectra
In NRQM the mass spectra of $`\mathrm{𝟓𝟔}_E`$ multiplet are given by the Hamiltonian
$`H_{NR}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}m_i^C+{\displaystyle \underset{i<j}{}}{\displaystyle \frac{a}{m_im_j}}𝝈^{\mathbf{(}𝒊\mathbf{)}}\mathbf{}𝝈^{\mathbf{(}𝒋\mathbf{)}}.`$ (7)
We extend this $`H_{NR}`$ to the form applicable to the chiral $`\rho `$-spin space. Intuitively, we take the form,
$`H`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}m_i^C+{\displaystyle \underset{i<j}{}}{\displaystyle \frac{(b_1+b_2\rho _3^{(i)})(b_1+b_2\rho _3^{(j)})a}{m_im_j}}𝝈^{\mathbf{(}𝒊\mathbf{)}}\mathbf{}𝝈^{\mathbf{(}𝒋\mathbf{)}}+C_\chi .`$ (8)
The first two terms reduce to $`H_{NR}`$ in non-relativistic $`|\rho ,\frac{3}{2}_S=|_\rho `$ space with constraint $`b_1+b_2=1`$. The $`C_\chi `$ term is phenomenologically introduced to describe the overall mass splittings between $`\mathrm{𝟓𝟔}_F`$ and $`\mathrm{𝟓𝟔}_E`$ and between $`\mathrm{𝟕𝟎}_G`$ and $`\mathrm{𝟓𝟔}_E`$. The $`C_\chi `$ for $`\mathrm{𝟓𝟔}_E`$ is taken to be 0 as $`C_E=0`$, while two $`C_\chi `$’s, $`C_F`$for $`\mathrm{𝟓𝟔}_F`$ and $`C_G`$ for $`\mathrm{𝟕𝟎}_G`$, are taken as independent parameters.
By taking the expectation value $`F\rho \sigma |H|F\rho \sigma `$, by WF given in Table 1, we obtain the mass formula of ground-state baryon systems. We use the seven parameters, which are determined as: $`m_n^C,m_s^C,A`$ are from the masses of $`N(939)`$, $`\mathrm{\Lambda }(1116)`$ and $`\mathrm{\Delta }(1232)`$. $`r`$ is from the mass of $`\mathrm{\Omega }(1672)`$. $`b_2,C_E`$ are from the masses of $`N(1440)`$ and $`\mathrm{\Sigma }(1660)`$. $`C_G`$ is from the mass of $`\mathrm{\Lambda }(1406)`$. For $`\mathrm{\Lambda }(1406)`$, the singlet-octet mixing in $`\mathrm{𝟕𝟎}_G`$ coming from $`T_{33}`$ breaking of spin-spin interaction is taken into account. We can predict the masses of all the other ground-state baryons in Table 2.
## 3 Strong pionic and kaonic decays
Next we consider the strong decays with one $`\pi `$ or $`K`$ emission. The relevant decay amplitudes $`T`$ are obtained by using the low energy theorem of Nambu-Goldstone bosons as $`T=J_{A\mu }^i\frac{1}{f_{\varphi ^i}}_\mu \varphi ^i`$, where $`J_{A_\mu }^i`$ is the baryonic axial-vector current, $`\varphi ^i`$ denotes the pseudoscalar octet, and $`f_\varphi `$ are their decay constants. The $`T`$ is given by the quark effective interaction,
$`L_{\mathrm{eff}}^{\mathrm{quark}}`$ $`=`$ $`\overline{q}i\gamma _5\gamma _\mu {\displaystyle \frac{\lambda ^i}{2}}q{\displaystyle \frac{1}{f_{\varphi ^i}}}_\mu \varphi ^i\overline{u}_q(p^{})i\gamma _5\gamma _\mu iq_\mu {\displaystyle \frac{\lambda ^i}{2f_{\varphi ^i}}}u_q(p)\varphi ^i(q),`$ (9)
$`L_{\mathrm{eff}}^{\mathrm{quark}}`$ $`=`$ $`g_\pi \overline{q}(si\gamma _5\varphi )qg_\pi \overline{u}_q(p^{})i\gamma _5{\displaystyle \frac{\lambda ^i}{\sqrt{2}}}u_q(p)\varphi ^i(q),(\varphi {\displaystyle \frac{\varphi ^i\lambda ^i}{\sqrt{2}}}),`$ (10)
where we also give the momentum representation. The $`L_{\mathrm{eff}}^{\mathrm{quark}}(L_{\mathrm{eff}}^{\mathrm{quark}})`$ is derivative (non-derivative) type interaction. We may use either interaction or both, since the $`L_{\mathrm{eff}}^{\mathrm{quark}}`$ is reduced to $`L_{\mathrm{eff}}^{\mathrm{quark}}`$ by using Dirac equation of constituent quark spinor as $`g_\pi \overline{u}_q(p^{})i\gamma _5u_q(p)=(g_\pi /2m_q)\overline{u}_q(p^{})i\gamma _5\gamma _\mu iq_\mu u_q(p)`$. Thus, the $`L_{\mathrm{eff}}^{\mathrm{quark}}`$ is equivalent to $`L_{\mathrm{eff}}^{\mathrm{quark}}`$. However, in $`\stackrel{~}{U}(12)`$-scheme, we cannot use the Dirac equation since there are no constituent quark spinors. We have only urciton spinors with velocity of the hadron itself. Thus, in the effective interaction of urciton spinors, the non-derivative type coupling shows very different effects from the derivative type coupling. We introduce both interaction with independent parameters,
$`L_{\mathrm{eff}}^{urciton}`$ $`=`$ $`g_{ND}/\sqrt{3}\overline{u}_{r^{},s^{}}(v_\mu ^{})(si\gamma _5\varphi )u_{r,s}(v_\mu )`$ (11)
$`+g_D/\sqrt{3}\overline{u}_{r^{},s^{}}(v_\mu ^{})(iv^{}\gamma )(s+i\gamma _5\varphi )(iv^\varphi \gamma )u_{r,s}(v_\mu ),`$
where $`v_\mu ^\varphi q_\mu /m_\varphi `$ and the factor $`1/\sqrt{3}`$ comes from the overlapping of color WFs. The $`(iv^{}\gamma )`$ factor takes the +1 for the Pauli states with $`r^{}=1`$. It is introduced for keeping the invariance of linear chiral transformation: $`u(v)exp\{i\beta ^i\lambda ^i\gamma _5/2\}u(v)`$ and $`(s+i\gamma _5\varphi )exp\{i\beta ^i\lambda ^i\gamma _5/2\}(si\gamma _5\varphi )exp\{i\beta ^i\lambda ^i\gamma _5/2\}`$. The matrix elements between the spinors in Pauli states($`r=+1`$) are expected to take the strength consistent with the low energy theorem, Eq. (9): For pion emission,
$`{\displaystyle \frac{m_\pi g_{ND}}{M+M^{}}}g_D{\displaystyle \frac{\sqrt{3}m_\pi }{\sqrt{2}f_\pi }},`$ (12)
where $`M(M^{})`$ is the mass of initial(final) baryon. We will check later this relation is satisfied approximately in our choice of parameters. For kaon emission we replace $`g_D`$ by $`g_D\frac{m_Kf_\pi }{m_\pi f_K}`$.
For the decays of flavor octet and singlet, we must consider the effect of $`U_A(1)`$ breaking six-point interaction, $`L_{\mathrm{int}}=G_D\{\mathrm{det}\overline{q}(1+\gamma _5)q+\mathrm{det}\overline{q}(1\gamma _5)q\}`$. By contracting the one $`\overline{q},q`$ pair with the spinor WF of pseudoscalar meson, we obtain the effective urciton interaction with the parameter $`g_{det}`$
$`L_{\mathrm{det}}^{\mathrm{urciton}}`$ $`=`$ $`ig_{det}ϵ_{a^{}b^{}c^{}}ϵ^{abc}\varphi _a{}_{}{}^{a^{}}\overline{u}(v^{})^b^{}u(v)_b\overline{u}(v^{})^c^{}\gamma _5u(v)_c,`$ (13)
where $`ϵ^{abc}`$ is the antisymmetric tensor of flavor $`SU(3)`$, which vanishes for decouplet states.
Next we consider the way to take the overlapping of spectator-quark indices. In the relativistic quantum mechanics, the overlapping leading to the probability amplitude is given by $`u^{}(\mathrm{𝟎})(\mathrm{})u(\mathrm{𝟎})`$, where we have the hydrogen-like composite system in mind, and 0 represents the velocity of the whole system. The system is non-covariant as a whole entity, and exactly speaking the overlapping is defined in the frame of both initial and final hadrons being at rest.
By using the four velocity $`v_0=(\mathrm{𝟎},i)`$ at rest, $`u^{}(\mathrm{𝟎})(\mathrm{})u(\mathrm{𝟎})=u^{}(v_0)(\mathrm{})u(v_0)=\overline{u}(v_0)(iv_0\gamma )(\mathrm{})u(v_0)`$, which is simply extended to the case with the final hadron with velocity $`v(v_0)`$ as $`\overline{u}(v)(iv\gamma )(\mathrm{})u(v_0)`$. We use this form for the spectator-quark transition where $`(\mathrm{})=1`$.
In the rest frame of initial baryon, the overlapping of WF is given by
$`3F\rho \sigma (v_\mu )|`$ $`(\mathrm{vertex})^{(1)}`$ $`(iv\gamma )^{(2)}(iv\gamma )^{(3)}|F\rho \sigma (v_{0\mu }),`$ (14)
where the interaction vertex, denoted as (vertex$`)^{(1)}`$, is given by Eqs. (11) and (13). We consider the total symmetricity of WF except for colors and the vertex is inserted only to the first spinor-index with the factor 3. The formula of the decay widths between $`\mathrm{𝟓𝟔}_E`$ and $`\mathrm{𝟓𝟔}_F`$ are summarized in Table 3.
In the actual calculation, the urciton spinor of final baryon with $`v_\mu `$ is obtained from the one with $`v_{0\mu }`$ by acting the Lorentz booster $`\overline{B}(v)`$ as
$`\overline{u}(v)(iv\gamma )=u^{}(v_0)\overline{B}(v),`$ $`\overline{u}(v)=u^{}(v_0)\rho _3\overline{B}(v),`$ (15)
where $`\overline{B}(v)=(ch\theta \rho _1𝒏\mathbf{}𝝈sh\theta )`$. The $`sh\theta `$ appears in the amplitude for the process with $`\rho _3`$-flipping in $`\overline{B}(v)`$. It is related with momentum of the final baryon through $`2ch\theta sh\theta =𝒑^{\mathbf{}}/M^{}`$, which is suppressed by baryon mass $`M^{}`$, not by constituent quark mass $`m_q`$. This suppression behavior is considered as a kind of conservation law in $`\stackrel{~}{U}(12)`$ scheme, named $`\rho _3`$-line rule.
The vertex of pseudoscalar emission is proportional to $`i\gamma _5=i\rho _1`$, which is the $`\rho _3`$-flip interaction. Thus, the transitions within $`\mathrm{𝟓𝟔}_E`$(which has $`|\rho ,_S`$) and $`\mathrm{𝟓𝟔}_F`$(which has $`|\rho ,_S`$) requires at least one $`\rho _1𝒏\mathbf{}𝝈sh\theta `$ factor in $`\overline{B}(v)`$. The relevant processes are $`\rho _3`$-violationg $`P`$-wave decays.
On the other hand, the decays of $`\mathrm{𝟕𝟎}_G`$(which has $`|\rho ,_{\alpha (\beta )}`$) to $`N_E\varphi `$ or $`D_E\varphi `$ are $`\rho _3`$-conserving $`S`$-wave, and their decay amplitudes are proportional to $`ch^3\theta `$, not suppressed by $`sh\theta `$. Their decay widths $`\mathrm{\Gamma }`$ generally becomes $``$a few GeV, which are quite large. The negative-parity $`\mathrm{𝟕𝟎}_G`$ baryons are expected not to be observed as resonances, but as backgrounds. The decay of $`\mathrm{\Lambda }(1406)`$ is the only exception. Its $`\mathrm{\Sigma }\pi `$ decay width is given by $`\mathrm{\Gamma }_{\mathrm{\Sigma }\pi }=\frac{M^{}𝒑^{\mathbf{}}}{4\pi M}H^2`$ with $`H=(c_ts_t)(g_{ND}+g_D\omega ^\pi )(c_t\frac{s_t}{2})\frac{g_{det}}{\sqrt{3}}`$, where $`c_t`$ and $`s_t`$ are the coefficients of singlet-octet mixing, given in the caption of Table 2. Only $`\mathrm{\Lambda }(1406)`$ has small width in $`\mathrm{𝟕𝟎}_G`$ because of its small decay phase space, and partly because of the cancellations between $`c_t`$ and $`s_t`$ terms and between $`g_{ND}`$ and $`g_{det}`$ terms.
The parameters $`g_{ND}`$, $`g_D`$ and $`g_{det}`$ are determined from the experimental decay widths of $`\mathrm{\Delta }(1232)N\pi `$, $`N(1440)N\pi `$ and $`\mathrm{\Lambda }(1406)\mathrm{\Sigma }\pi `$, as $`(g_{ND},g_D,g_{det})=(37.65,0.761,22.39)`$. We can see the relation Eq. (12) is approximately satisfied. Now we can predict the widths of all the other decay channels. The results are shown in Table 4.
As can be seen in Table 4, the predicted values of decay widths seem to be consistent with the present experimental data although they have large uncertainty. We use $`\mathrm{\Gamma }(N(1440)N\pi )`$ as input. However, its magnitude is predicted through the pion low energy theorem and our assignment of $`N(1440)`$ as ground $`N_F`$ state. If the $`N(1440)`$ is assigned as a radially excited $`2S`$ state $`N(2S)`$ as in the conventional quark model, the decay width becomes much smaller than the $`\mathrm{\Gamma }_{\mathrm{exp}}`$ because of the small overlapping factor of $`1S`$ and $`2S`$ WFs. The $`\mathrm{\Lambda }(1617)`$ and $`\mathrm{\Sigma }(1660)`$ also show the plausible properties as octet chiral states in $`\mathrm{𝟓𝟔}_F`$. Their rather small widths compared with $`N(1440)`$ are explained by the cancellation due to determinant-type interaction.
On the other hand the total decay widths of $`D_F`$ multiplet are predited with the values more than $``$500MeV, much larger than the other decay modes. Naively, in Table 3 the factor $`\frac{10}{9}`$ of $`D_FD_E\varphi `$ decay is an order of magnitude larger than the other channels. This factor comes from the overlapping of $`\sigma `$-spin WFs. For instance, $`\mathrm{\Delta }(1608)`$ have $`\mathrm{\Delta }\pi `$ widths as $``$1.7GeV, which is much larger than the width of the observed resonance $`\mathrm{\Delta }(1600)`$.
We can successfully explain the property of $`N(1440)`$, which is assigned in $`\stackrel{~}{U}(12)`$ scheme as a $`S`$-wave $`N_F`$ state out of the conventional $`SU(6)_{SF}`$ framework. Our new assignment raises a question at the same moment, where the genuine radially-excited $`N(2S)`$ state exists. Experimentally there is no plausible candidate of $`N(2S)`$ observed below 2 GeV.
We can predict the mass of $`N(2S)`$ by using the covariant oscillator quark model (COQM) with one-gluon-exchange potential. The resulting $`M_{N(2S)}`$ is $`1.74`$GeV.<sup>1</sup><sup>1</sup>1 In COQM the WFs are determined by the spin-independent 4-dimensional oscillator potential, reproducing the linearly-rising Regge trajectories. By using this WF, we can predict the magnitudes of Coulomb $`V_C`$, spin-spin $`V_{SS}`$, orbit-orbit $`V_{OO}`$, tensor $`V_T`$ and spin-orbit $`V_{SO}`$ potentials. The orbitally excited $`L=2`$ states are around $`M_{1D}1.95`$GeV, where $`\mathrm{\Delta }_{\frac{7}{2}^+}(1950)`$ and $`N_{\frac{5}{2}^+}(2000)`$ are observed. This $`M_{1D}`$ is used as input and by using the values of $`V_C`$, $`V_{SS}`$ and $`V_{OO}`$, we obtain the $`M_{N(2S)}`$ as $`1.74`$GeV. By using this $`M_{N(2S)}`$ and the decay coupling constants determined in this section, the decay widths of $`N(2S)`$ can be predicted. We obtain small $`\mathrm{\Gamma }_{N(2S)N\pi }60`$MeV because this is $`\rho _3`$-violating process and partly because of the small overlapping factor from space-time WFs. The small $`\mathrm{\Gamma }`$ is also obtained in the other quark model. However, in $`\stackrel{~}{U}(12)`$ scheme, we must consider the decays to the negative-parity $`\mathrm{𝟕𝟎}_G`$ baryons, that is, $`N(2S)N_{\frac{1}{2}}(1249)\pi `$ and $`\mathrm{\Delta }_{\frac{1}{2}}(1284)\pi `$ (See Table 2). These processes are $`\rho _3`$-conserving and the large decay widths are predicted, with 360(290)MeV for the former(latter). There is another decay mode such as $`N(2S)N\sigma N\pi \pi `$. Totally the $`N(2S)`$ is predicted to have decay width $`>0.7`$GeV, and not to be observed as a resonance, but as a $`N\pi \pi `$ background. The other octet $`2S`$ baryons are also not expected to be observed for the same reason in our scheme. This prediction seems to be consistent with the present experimental situation.
## 4 Radiative decays of $`\mathrm{\Delta }(1232)`$, $`\mathrm{\Lambda }(1406)`$ and $`N(1440)`$
Finally we consider the radiative decays of $`\mathrm{\Delta }(1232)`$, $`\mathrm{\Lambda }(1406)`$ and $`N(1440)`$, observed experimentally. The spin-current interaction is given in the covariant form by
$`L_{spin}^{\mathrm{urciton}}`$ $`=`$ $`\overline{u}(v)(iv\gamma )A_\mu (q){\displaystyle \frac{Q^{(i)}e}{2m_i}}i\sigma _{\mu \nu }iq_\nu (b_1^{}+b_2^{}(iv_0\gamma ))u(v_0).`$ (16)
This form aside from the factor $`(b_1^{}+b_2^{}(iv_0\gamma ))`$ is used in the analyses of radiative decays in COQM and in the analyses of $`D_s`$ system. Analyses of magnetic moments of the nucleon give $`m_n(m_s)=336(510)`$MeV. In order to keep this result we require a constraint $`b_1^{}+b_2^{}=1`$. As a result, the interaction vertex between $`u^{}(\mathrm{𝟎})`$ and $`u(\mathrm{𝟎})`$ is given by
$`\frac{e|𝒒|}{\sqrt{2}m_n}\overline{B}(v)Q_M\sigma _{}(1\pm \rho _1)(b_1^{}+b_2^{}\rho _3)`$, where the sign corresponds to the helicity $`1`$ of the photon emitted to the $`z`$ direction. There is no effect of $`b_2^{}`$ term for the transition between Pauli-states, while it gives sizable contribution for the decays of chiral states with $`\rho _3`$-flipping.
The parameter $`b_2^{}`$ is taken to be $`b_2^{}=0.136`$, which is determined by the width of $`\mathrm{\Lambda }(1406)\mathrm{\Lambda }\gamma `$ as input. Now we can predict the radiaitve decay amplitudes $`A_{\frac{1}{2}}`$ or $`A_{\frac{3}{2}}`$, where $`A`$ is related with the ordinary Lorentz-invariant amplitude $`T`$ as $`A=T/\sqrt{2|𝒒|}`$. We show only the numerical results in Table 5.
It is remarkable that the amplitude of $`\mathrm{\Delta }(1232)N\gamma `$ is greatly improved from the prediction by NRQM, where $`A_{\frac{3}{2}}=\sqrt{|𝒒|/6}e/m_n=0.187`$GeV$`^{\frac{1}{2}}`$ leading $`\mathrm{\Gamma }=`$377keV, which is much smaller than the experimental value. In our covariant $`\stackrel{~}{U}(12)`$ scheme we obtain the relativistic recoil factor $`ch^3\theta +ch^2\theta sh\theta =1.17`$ for the amplitude, leading to the great improvement, $`\mathrm{\Gamma }=516`$keV. There is no effect from $`b_2^{}`$ for $`\mathrm{\Delta }`$ decay, while there are sizable cancellation between $`b_1^{}`$ and $`b_2^{}`$ terms for $`\mathrm{\Lambda }(1406)`$ decays. If we fix $`b_2^{}=0`$ we obtain $`\mathrm{\Lambda }\gamma (\mathrm{\Sigma }^0\gamma )`$ width as 52(64)keV, which are not so much different from the experimental values.
On the other hand, for $`N(1440)N\gamma `$ decay, the amplitude vanishes until the order of $`ch^2\theta sh\theta `$. There is a contribution of order $`ch\theta sh^2\theta `$, however, because of the suppression by $`sh\theta `$, the predicted $`A_{\frac{1}{2}}`$ is of an order smaller than the experimental value. This problem may be possibly resolved by considering the small mixing of $`N_E`$ WF to the $`N_F`$ WF in $`N(1440)`$, of which effect is proportional to $`ch^3\theta `$. We consider in this paper only ideal case that $`N(939)`$($`N(1440)`$) is purely $`N_E`$ Pauli state($`N_F`$ chiral state). This problem will be considered elsewhere.
## 5 Concluding Remarks
In this letter we study the spectra and decay properties of light-quark baryon systems in $`\stackrel{~}{U}(12)`$ classification scheme of hadrons. The $`N(1440)`$, $`\mathrm{\Lambda }(1600)`$ and $`\mathrm{\Sigma }(1660)`$ have the property of the octet chiral states in the extra ground-state $`\mathrm{𝟓𝟔}_F`$ multiplet. The existence of $`\mathrm{\Xi }(1820)`$ in the same multiplet is predicted with the decay width 40 MeV to $`\mathrm{\Xi }^{}\pi `$ channel, and the negligible decay width to $`\mathrm{\Lambda }K`$ channel.
The negative parity $`\mathrm{𝟕𝟎}_G`$ multiplet has generally very wide widths, and are expected to be not observed as resonances, except for $`\mathrm{\Lambda }(1406)`$, where the cancellation due to $`U_A(1)`$ breaking interaction diminishes its decay width to $`\mathrm{\Sigma }\pi `$. These results suggest the $`\stackrel{~}{U}(12)`$ representation is actually realized in baryon spectroscopy of light quarks. Especially the parity-inversion problem of excited nucleon spectra is naturally resolved in $`\stackrel{~}{U}(12)`$ scheme, since the mass and strong decay property of $`N(1440)`$ are explained as the chiral $`S`$-wave $`N_F`$ state. Ordinary radially excited $`N(2S)`$ state is not observed because of its large predicted widths of the decays to chiral $`\mathrm{𝟕𝟎}_G`$ baryons.
The property of $`\mathrm{\Delta }(1600)`$ cannot be explained as the ground-state chiralon. However, we expect the existence of $`P`$-wave chiral state in this energy region, and the width of $`\mathrm{\Delta }(1600)`$ may be explained from the possible mixing effect to this $`P`$-wave chiral state. This problem will be considered elsewhere.
The author would like to express his sincere gratitude to Professor S. Ishida and Professor K. Yamada for their useful comments.
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# Understanding transverse coherence properties of X-ray beams in third generation Synchrotron Radiation sources
## 1 Introduction
In recent years, continuous evolution of third generation light sources has allowed dramatic increase of brilliance with respect to older designs, which has triggered a number of new techniques and experiments unthinkable before. Among the most exciting properties of today third generation facilities is the high flux of coherent X-rays provided. The availability of intense coherent X-ray beams has fostered the development of new coherence-based techniques like fluctuation correlation dynamics, phase imaging, coherent X-ray diffraction (CXD) and X-ray holography. In this context, understanding the evolution of transverse coherence properties of Synchrotron Radiation (SR) along the beam line is of fundamental importance.
In general, when dealing with this problem, one should account for the fact that Synchrotron Radiation is a random statistical process. Therefore, the evolution of transverse coherence properties should be treated in terms of probabilistic statements: the shot noise in the electron beam causes fluctuations of the beam density which are random in time and space. As a result, the radiation produced by such a beam has random amplitudes and phases.
Statistical Optics GOOD , MAND , NEIL affords convenient tools to deal with fluctuating electromagnetic fields in an appropriate way. Among the most important quantities needed to describe coherent phenomena in the framework of Statistical Optics is the correlation function of the electric field. In any interference experiment one needs to know the system (second order) correlation function of the signal at a certain time and position with the signal at another time and position. Alternatively, and equivalently, one can describe the same experiment in frequency domain. In this case one is interested in the correlation function of the Fourier transform of the time domain signal at a certain frequency and position with the Fourier transform of the time domain signal at another frequency and position. The signal one is interested to study is, indeed, the Fourier transform of the original signal in time domain. In SR experiments the analysis in frequency domain is much more natural than that in the time domain. In fact, up-to-date detectors are limited to about $`100`$ ps time resolution and they are by no means able to resolve a single X-ray pulse in time domain. They work, instead, by counting the number of photons at a certain frequency over an integration time longer than the radiation pulse. Therefore, in this paper we will deal with signals in the frequency domain and we will often refer to the ”Fourier transform of the electric field” simply as ”the field”.
For some particular experiment one may be interested in higher order correlation functions (for instance, in the correlation between the intensities) which, in general, must be calculated separately. In the particular case when the field fluctuations can be described as a Gaussian process, the field is often said to obey Gaussian statistics. In this case, with the help of the Moment Theorem GOOD one can recover correlation functions of any order from the knowledge of the second order one: this constitutes a great simplification to the task of describing coherence properties of light. A practical example of a field obeying Gaussian statistics is constituted by the case of polarized thermal light. This is more than a simple example: in fact, Statistical Optics has largely developed in connection with problems involving optical sources emitting thermal light like the sun, other stars, or incandescent lamps. As a consequence, Gaussian statistics is often taken for granted. Anyway, it is not a priori clear wether Synchrotron Radiation fields obey it or not; our analysis will show that Synchrotron Radiation is indeed a Gaussian random process. Therefore, as is also the case for polarized thermal light and any other signal obeying Gaussian statistics, when we deal with Synchrotron Radiation the basic quantity to consider is the second order correlation function of the field. Moreover, as already discussed, in Synchrotron Radiation experiments it is natural to work in the space-frequency domain, so that we will focus, in particular, on the second order correlation function in the space-frequency domain.
Besides obeying Gaussian statistics, polarized thermal light has two other specific properties allowing simplifications of the theory: the first is stationarity<sup>1</sup><sup>1</sup>1Here we do not distinguish between different kind of stationarity because, under the assumption of a Gaussian process, these concepts simply coincide. and the second is quasi-homogeneity. Exactly as the property of Gaussian statistics, also stationarity and quasi-homogeneity of the source are usually taken for granted in Statistical Optics problems but, unlike it, they do not belong, in general, to Synchrotron Radiation fields. In fact, as we will show, Synchrotron radiation fields are intrinsically non-stationary and not always quasi-homogeneous. Nevertheless, up to now it has been a widespread practice to assume that undulator sources are completely incoherent (i.e. homogeneous) and to apply the well known van Cittert-Zernike (VCZ) theorem for calculating the degree of transverse coherence in the far-field approximation PETR . Using the VCZ theorem, the electric field cross-correlation function in the far field is usually calculated (aside for a geometrical phase factor) as a Fourier transformation of the intensity distribution of the source, customarily located at the exit of the undulator.
Although the VCZ theorem only deals with completely incoherent sources, there exists an analogous generalized version of it which allows to extend the treat the case of quasi-homogeneous sources as well. Actually there is no unambiguous choice of terminology in literature regarding the scope of the VCZ theorem. For instance, a very well-known textbook MAND reports of a ”Zernike-propagation” equation dealing with any distance from the source. Also, sometimes GOOD , the generalized VCZ theorem is referred to as Schell’s theorem (and also used in some paper CHAN ). In this paper we will refer to the VCZ theorem and its generalized version only in the limit for a large distance from the source and for, respectively, homogeneous and quasi-homogeneous sources. However, irrespectively of different denominations, the fundamental fact holds, that once a cross-correlation function is known on a given source plane it can be propagated through the beamline at any distance from the source. It should be noted that, from this viewpoint, a source simply denotes an initial plane down the beamline from which the cross-correlation function is propagated further. Then, the position of the source down the beamline is suggested only on the ground of opportunity. On the contrary, when dealing with the VCZ theorem, the source must be (quasi)-homogeneous which explains the customary location at the exit of the undulator.
In some cases, the VCZ theorem or its generalized version may provide a convenient method for calculating the degree of transverse coherence in various parts of the beamline once the transverse coherence properties of the photon beam are specified at the exit of the undulator, that is at the source plane. In most SR applications though, such treatment is questionable. First, the source (even at the exit of the undulator) may not be quasi-homogeneous. Second, even for specific sets of problem parameters where the quasi-homogeneous model is accurate, the specification of the far-field zone depends not only on the electron beam sizes, but also on the electron beam divergencies (in both direction) and on the intrinsic divergence of the radiation connected with the undulator device. At the time being, widespread and a-critical use of the VCZ theorem and its generalization shows that there is no understanding of transverse coherence properties of X-ray beams in third generation Synchrotron Radiation sources.
If, on the one hand, the definition of the far-zone and the possible non quasi-homogeneity of SR sources constitute serious problems in the description of the coherence properties of Synchrotron light, on the other hand the intrinsic non-stationarity of the SR process does not play a very important role. In particular, as we will show, assumption of a minimal undulator bandwidth much larger than the characteristic inverse bunch duration (which is always verified in practice) allows to separate the correlation function in space-frequency domain in the product of two functions. The first function is a spectral correlation describing correlation in frequency. The second function describes correlation in space and is well-known also in the case of stationary processes as the cross-spectral density of the process. Then, the cross-spectral density can be studied independently at any given frequency giving information on the spatial correlation of the field. Subsequently, the knowledge of the spectral correlation function brings back the full expression for the space-frequency correlation.
In this paper we aim at the development of a theory of transverse coherence capable of providing very specific predictions, relevant to practice, regarding the cross-spectral density of undulator radiation at various positions along the beam-line. A fully general study of undulator sources is not a trivial one. Difficulties arise when one tries to include simultaneously the effect of intrinsic divergence of the radiation due to the presence of the undulator, of electron beam size and electron beam divergence into the insertion device. The full problem, including all effects, poses an unsolvable analytical challenge, and numerical calculations are to be preferred. Generally, the cross-spectral density of the undulator radiation is controlled by nine physical parameters which model both the electron beam and the undulator: the horizontal and vertical geometrical emittances of the electron beam $`ϵ_{x,y}`$, the horizontal and vertical minimal betatron functions $`\beta _{x,y}^o`$, the observation distance down the beamline $`z_o`$, the observation frequency $`\omega `$, the undulator resonant frequency $`\omega _o`$, the undulator length $`L_w`$ and the length of the undulator period $`\lambda _w`$.
We will make a consistent use of dimensional analysis. Dimensional analysis of any problem, performed prior to analytical or numerical investigations, not only reduces the number of independent terms, but also allows one to classify the grouping of dimensional variables in a way that is most suitable for subsequent study. The algorithm for calculating the cross-spectral density can be formulated as a relation between dimensionless quantities. After appropriate normalization, the radiation cross-spectral density from an undulator device is described by six dimensionless quantities: the normalized emittances $`\widehat{ϵ}_{x,y}=\omega _oϵ_{x,y}/c`$, the normalized betatron functions $`\widehat{\beta }_{x,y}=\beta _{x,y}^o/L_w`$, the normalized observation distance $`\widehat{z}_o=z_o/L_w`$ and the normalized detuning parameter $`\widehat{C}=2\pi N_w(\omega \omega _o)/\omega _o`$, where $`N_w=L_w/\lambda _w`$ is the number of undulator periods.
At some point in this work we will find it convenient to pose $`\widehat{C}=0`$. In other words we will assume that parameters are tuned at perfect resonance. It is relevant to note that even under this simplifying assumptions, conditions for the undulator source to be quasi-homogeneous still include four parameters $`\widehat{ϵ}_{x,y}`$ and $`\widehat{\beta }_{x,y}`$. For storage rings that are in operation or planned in the $`\mathrm{\AA }`$ngstrom wavelength range, the parameter variation of $`\widehat{ϵ}_x1010^3`$, $`\widehat{ϵ}_y10^110`$, $`\widehat{\beta }_{x,y}10^110`$ are possible: these include many practical situations in which the assumption of quasi-homogeneous sources and, therefore the (generalized) VCZ theorem, is not accurate.
In this paper we will first deal with the most general case of non-homogeneous sources. In fact, from a practical viewpoint, it is important to determine the cross-spectral density as a function of $`\widehat{ϵ}_{x,y}`$, $`\widehat{\beta }_{x,y}`$, $`\widehat{C}`$ and $`\widehat{z}_o`$. Once a general expression for the cross-spectral density is found, it can be used as a basis for numerical calculations. A second goal of this work is to find the region of applicability of the quasi-homogeneous source model (i.e. of the generalized VCZ theorem) which will arise automatically from the dimensional analysis of the problem. Finally, we will derive analytical expressions for the cross-spectral density at $`\widehat{C}=0`$ in various parts of the beamline.
Results may also be obtained using numerical techniques alone, starting from the Lienard-Wiechert expressions for the electromagnetic field and applying the definition of the field correlation function without any analytical manipulation. Yet, computer codes can calculate properties for a given set of parameters, but can hardly improve physical understanding, which is particularly important in the stage of planning experiments: understanding of correct approximations and their region of applicability with the help of a consistent use of dimensional analysis can simplify many tasks a lot, including practical and non-trivial ones. Moreover, at the time being, no code capable to deal with transverse coherence problems has been developed at all.
It should be noted that some theoretical attempt to follow this path has been proposed in TAKA . Among the results of that paper is the fact that van Cittert-Zernike theorem could not be applied unless the electron beam divergence is much smaller than the diffraction angle, which is never verified in practice in the horizontal plane. We will show that this conclusion is incorrect.
We organize our work as follows. After this Introduction, in Section 2 we present a second-order theory of coherence for fields generated by Synchrotron Radiation sources. In Section 3 we give a derivation of the cross-spectral density for undulator-based sources in reduced units. Subsequently, we analyze the evolution of the cross-spectral density function through the beamline in the limit for $`\widehat{C}=0`$. A particular case of quasi-homogeneous sources and its applicability region is treated under several simplifying assumptions in Section 4. Effects of the vertical emittance on the cross-spectral density are discussed in detail in Section 5, while a treatment of some non quasi-homogeneous source is given in the following Section 6. Obtained results include approximate design formula capable of describing in very simple terms the evolution of the coherence length along the beamline in many situation of practical interest. A good physical insight is useful to identify possible applications of given phenomena. In particular in Section 7 we selected one practical application to exploit the power of our approach. We show that, by means of a simple vertical slit, it is possible to manipulate transverse coherence properties of an X-ray beam to obtain a convenient coherent spot-size on the sample. This invention was devised almost entirely on the basis of theoretical ideas of rather complex and abstract nature which have been described in this paper. Finally, in Section 8, we come to conclusions.
## 2 Second-order coherence theory of fields generated by Synchrotron Radiation sources
### 2.1 Thermal light and Synchrotron Radiation: some concepts and definitions
A great majority of optical sources emits thermal light. Such is the case of the sun and the other stars, as well as of incandescent lamps. This kind of radiation consists of a large number of independent contributions (radiating atoms) and is characterized by random amplitudes and phases in space and time. The electromagnetic fields can be then conveniently described in terms of Statistical Optics, a branch of Physics that has been intensively developed during the last few decades. Today one can take advantage of a lot of existing experience and theoretical basis for the descriptions of fluctuating electromagnetic fields GOOD , MAND .
Consider the light emitted by a thermal source passing through a polarization analyzer (see Fig. 1). Properties of polarized thermal light are well-known in Statistical Optics, and are referred to as properties of completely chaotic, polarized light GOOD , MAND . Thermal light is a statistical random process and statements about such process are probabilistic statements. Statistical processes are handled using the concept of statistical ensemble, drawn from Statistical Mechanics, and statistical averages are performed indeed, over many ensembles, or realizations, or outcomes of the statistical process under study.
Polarized thermal light is a very particular kind of random process in that it is Gaussian, stationary and ergodic. Let us discuss these characteristics in more detail.
The properties of Gaussian random processes are well-known in Statistical Optics. For instance, the real and imaginary part of the complex amplitudes of the electric field from a polarized thermal source have Gaussian distribution, while the instantaneous radiation power fluctuates in accordance with the negative exponential distribution. Gaussian statistics alone, guarantees that higher-order correlation functions can be expressed in terms of second-order correlation functions. Moreover, it can be shown GOOD that a linearly filtered Gaussian process is also a Gaussian random process. As a result, the presence of a spectral filter (monochromator) and a spatial filter as in the system depicted in Fig. 1 do not change the statistics of the signal, because they simply act as linear filters.
Stationarity is a subtle concept. There are different kinds of stationarity. Strict-stationarity means that all ensemble averages are independent on time. Wide-sense stationarity means that the signal average is independent on time and that the second order correlation function in time depends only on the difference of the observation times. However, for Gaussian processes strict and wide-sense stationarity coincide GOOD , MAND . As a consequence of the definition of stationarity, necessary condition for a certain process to be stationary is that the signal last forever. Yet, if a signal lasts much longer than its coherence time $`\tau _c`$ (which fixes the short-scale duration of the field fluctuations) and it is observed for a time much shorter than its duration $`\sigma _T`$, but much longer than its coherence time it can be reasonably considered as everlasting and it has a chance to be stationary as well, as in the case of thermal light.
Ergodicity is a stronger requirement than stationarity. Qualitatively, we may state that if, for a given random process all ensemble averages can be substituted by time averages, the process under study is said to be ergodic: all the statistical properties of the process can be derived from one single realization. A process must be strictly stationary in order to be ergodic. There exist stationary processes which are not ergodic. One may consider, for instance, the random constant process: this is trivially strictly stationary, but not ergodic because a single (constant) realization of the process does not allow one to characterize the process from a statistical viewpoint. However, this is a pathologic case when both the coherence time $`\tau _c`$ and the duration time of the signal $`\sigma _T`$ are infinite. On the contrary, a stationary process like the radiation from an incandescent lamp driven by a constant current has, virtually, infinite duration. In this case different ensembles are simply different observations, for given time intervals, of the same, statistically identical phenomenon: then, the concept of ensemble average and time average are equivalent and the process is also ergodic.
Statistical Optics was developed starting with signals characterized by Gaussian statistics, stationarity and ergodicity. Let us consider any Synchrotron radiation source. Like thermal light, also Synchrotron Radiation is a random process. In fact, relativistic electrons in a storage ring emit Synchrotron Radiation passing through bending magnets or undulators. The electron beam shot noise causes fluctuations of the beam density which are random in time and space from bunch to bunch. As a result, the radiation produced has random amplitudes and phases. As already declared in the Introduction we will demonstrate that the SR field obeys Gaussian statistics. In contrast with thermal light though, Synchrotron Radiation is intrinsically non-stationary (and, therefore, non-ergodic) because even if its short pulse duration cannot be resolved by detectors working in the time domain, it can nonetheless be resolved by detectors working in frequency domain. For this reason, in what follows the averaging brackets $`\mathrm{}`$ will always indicate the ensemble average over bunches.
In spite of differences with respect to the simpler case of thermal light, as we will see in this paper, also Synchrotron Radiation fields can be described in terms of Statistical Optics. Fig. 2 shows the geometry of the experiment under consideration. The problem is to describe the statistical properties of Synchrotron Radiation at the detector installed after the spatial and spectral filters. Radiation at the detector consists of a carrier modulation of frequency $`\omega `$ subjected to random amplitude and phase modulation. The Fourier decomposition of the radiation contains frequencies spread about the monochromator bandwidth $`\mathrm{\Delta }\omega _\mathrm{m}`$: it is not possible, in practice, to resolve the oscillations of the radiation fields which occur at the frequency of the carrier modulation. It is therefore appropriate, for comparison with experimental results, to average the theoretical results over a cycle of oscillations of the carrier modulation.
Fig. 3 gives a qualitative illustration of the type of fluctuations that occur in cycle-averaged Synchrotron Radiation beam intensity. Within some characteristic time, a given random function appears to be smooth, but when observed at larger scales the same random function exhibits ”rough” variations. The time scale of random fluctuations is the coherence time $`\tau _c`$. When $`\tau _c\sigma _T`$ the radiation beyond the monochromator is partially coherent. This case is shown in Fig. 3: there, we can estimate $`\tau _c\mathrm{\Delta }\omega _\mathrm{m}^1`$. If the radiation beyond the monochromator is partially coherent, a spiky spectrum is to be expected. The nature of the spikes is easily described in terms of Fourier transform theory. We can expect that the typical width of the spectrum envelope should be of order of $`\mathrm{\Delta }\omega /\omega (\tau _c\omega )^1`$. Also, the spectrum of the radiation from a bunch with typical duration $`\sigma _T`$ at the source plane should contain spikes with characteristic width $`\mathrm{\Delta }\omega /\omega (\omega \sigma _T)^1`$, as a consequence of the reciprocal width relations of Fourier transform pairs (see, again, Fig. 3).
### 2.2 Second-order correlations in space-frequency domain
We start our discussion in the most generic way possible, considering a fixed polarization component of the Fourier transform at frequency $`\omega `$ of the electric field produced at location $`(z_o,\stackrel{}{r}_o)`$, in some cartesian coordinate system, by a given collection of sources. We will denote it with $`\overline{E}_{}(z_o,\stackrel{}{r}_o,\omega )`$ and it will be linked to the time domain field $`E_{}(z_o,\stackrel{}{r}_o,t)`$ through the Fourier transform
$$\overline{E}_{}(\omega )=_{\mathrm{}}^{\mathrm{}}𝑑tE_{}(t)e^{i\omega t},$$
(3)
so that
$$E_{}(t)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑\omega \overline{E}_{}(\omega )e^{i\omega t}.$$
(4)
This very general collection of sources includes the case of an ultra relativistic electron beam going through a certain magnetic system and in particular an undulator, which is our case of interest. In this case $`z_o`$ is simply the observation distance along the optical axis of the system and $`\stackrel{}{r}_o`$ are the transverse coordinates of the observer on the observation plane. The contribution of the $`k`$-th electron to the field Fourier transform at the observation point depends on the transverse offset $`(l_{xk},l_{yk})`$ and deflection angles $`(\eta _{xk},\eta _{yk})`$ that the electron has at the entrance of the system with respect to the optical axis. Moreover, an arrival time $`t_k`$ at the system entrance has the effect of multiplying the field Fourier transform by a phase factor $`\mathrm{exp}(i\omega t_k)`$ (that is, in time domain the electric field is retarded by a time $`t_k`$). At this point we do not need to specify explicitly the dependence on offset and deflection. The total field Fourier transform can be written as
$$\overline{E}_{}(z_o,\stackrel{}{r}_o,\omega )=\underset{k=1}{\overset{N}{}}\overline{E}_s(\stackrel{}{\eta }_k,\stackrel{}{l}_k,z_o,\stackrel{}{r}_o,\omega )\mathrm{exp}(i\omega t_k),$$
(5)
where $`\stackrel{}{\eta }_k,\stackrel{}{l}_k`$ and $`t_k`$ are random variables and $`N`$ is the number of electrons in the beam. It follows from Eq. (5) that the Fourier transform of the Synchrotron Radiation pulse at a fixed frequency and a fixed point in space is a sum of a great many independent contributions, one for each electron, of the form $`\overline{E}_s(\stackrel{}{\eta }_k,\stackrel{}{l}_k,z_o,\stackrel{}{r}_o,\omega )\mathrm{exp}(i\omega t_k)`$. For simplicity we make three assumptions about the statistical properties of elementary phasors composing the sum, which are generally satisfied in Synchrotron Radiation problems of interest.
1) We assume that for a beam circulating in a storage ring random variables $`t_n`$ are independent from $`\stackrel{}{\eta }_n`$ and $`\stackrel{}{l}_n`$. This is always verified, because the random arrival times of electrons, due to shot noise, do not depend on the electrons offset and deflection with respect to the $`z`$-direction. Eq. (5) states that the $`k`$-th elementary contribution to the total $`\overline{E}_{}`$ can be written as a product of the complex phasors $`\mathrm{exp}(i\omega t_k)`$, and $`\overline{E}_s`$ that, in its turn, can be written as a product of modulus and phase as $`\overline{E}_s(\stackrel{}{\eta }_k,\stackrel{}{l}_k,z_o,\stackrel{}{r}_o,\omega )=\overline{E}_{sk}\mathrm{exp}(i\varphi _k)`$. Under the assumption of statistical independence of $`t_n`$ from $`\stackrel{}{\eta }_n`$ and $`\stackrel{}{l}_n`$ the complex phasors $`\mathrm{exp}(i\omega t_k)`$, and $`\overline{E}_s`$ are statistically independent of each other and of all the other elementary phasors for different values of $`k`$. The ensemble average of a given function $`f`$ of random variables $`\stackrel{}{\eta }_n`$, $`\stackrel{}{l}_n`$ and $`t_n`$ is by definition:
$`f(\stackrel{}{\eta }_n,\stackrel{}{l}_n,t_n)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta _{xn}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\eta _{yn}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑l_{xn}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑l_{yn}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_n`$ (6)
$`\times f(\stackrel{}{\eta }_n,\stackrel{}{l}_n,t_n)P(\stackrel{}{\eta }_n,\stackrel{}{l}_n,t_n),`$ (7)
where $`P(\stackrel{}{\eta }_n,\stackrel{}{l}_n,t_n)`$ is the probability density distribution in the joint random variables $`\stackrel{}{\eta }_n`$, $`\stackrel{}{l}_n`$, $`t_n`$. Independence of $`t_n`$ from $`\stackrel{}{\eta }_n`$ and $`\stackrel{}{l}_n`$ allows us to write
$$P(\stackrel{}{\eta }_n,\stackrel{}{l}_n,t_n)=F_{\eta _x,l_x}(\eta _{xn},l_{xn})F_{\eta _y,l_y}(\eta _{yn},l_{yn})F_t(t_n),$$
(8)
where we also assumed that the distribution in the horizontal and vertical planes are not correlated. Since electrons arrival times are completely uncorrelated from transverse coordinates and offsets, the shapes of $`F_{\eta _x,l_x}`$, $`F_{\eta _y,l_y}`$ and $`F_t`$ are the same for all electrons.
2) We assume that the random variables $`\overline{E}_{sk}`$ (at fixed frequency $`\omega `$), are identically distributed for all values of $`k`$, with a finite mean $`\overline{E}_{sk}`$ and a finite second moment $`\overline{E}_{sk}^2`$. This is always the case in practice because electrons are identical particles.
3) We assume that the electron bunch duration $`\sigma _T`$ is large enough so that $`\omega \sigma _T1`$: under this assumption the phases $`\omega t_k`$ can be regarded as uniformly distributed on the interval $`(0,2\pi )`$. The assumption $`\omega \sigma _T1`$ is justified by the fact that $`\omega `$ is the undulator resonant frequency, which is high enough to guarantee that $`\omega \sigma _T1`$ for any practical choice of $`\sigma _T`$.
The formal summation of phasors with random lengths and phases is illustrated in Fig. 4. Under the three previously discussed assumptions we can use the central limit theorem to conclude that the real and the imaginary part of $`\overline{E}_{}`$ are distributed in accordance to a Gaussian law. Detailed proof of this fact is given in Appendix A. As a result, Synchrotron Radiation is a Gaussian random process and second-order field correlation function is all we need in order to specify the field statistical properties. In fact, as already remarked, higher-order correlation functions can be expressed in terms of second-order correlation functions.
In Synchrotron Radiation experiments with third generation light sources detectors are limited to about $`100`$ ps time resolution and are by no means able to resolve a single X-ray pulse in time domain: they work, instead, by counting the number of photons at a certain frequency over an integration time longer than the pulse. Therefore, for Synchrotron Radiation related issues the frequency domain is much more natural a choice than the time domain, and we will deal with signals in the frequency domain throughout this paper. The knowledge of the second-order field correlation function in frequency domain
$$\mathrm{\Gamma }_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega ,\omega ^{})=\overline{E}_{}(z_o,\stackrel{}{r}_{o1},\omega )\overline{E}_{}^{}(z_o,\stackrel{}{r}_{o2},\omega ^{}),$$
(9)
is all we need to completely characterize the signal from a statistical viewpoint. For the sake of completeness it is nonetheless interesting to remark that it is possible (and often done, in Statistical Optics) to give equivalent descriptions of the process in time domain as well. First, note that the time domain process $`E_{}(t)`$ is linked to $`\overline{E}_{}(\omega )`$ by Fourier transform, and that a linearly filtered Gaussian process is also a Gaussian process (see GOOD 3.6.2). As a result, $`E_{}(t)`$ is a Gaussian process as well. Second, the operation of ensemble average is linear with respect to Fourier transform integration. This guarantees, that the knowledge of $`\mathrm{\Gamma }_\omega `$ in frequency domain is completely equivalent to the knowledge of the second-order correlation function between $`E_{}(z_o,\stackrel{}{r}_{o1},t_1)`$ and $`E_{}(z_o,\stackrel{}{r}_{o2},t_2)`$. The latter is usually known as mutual coherence function and was first introduced in WOLF :
$$\mathrm{\Gamma }_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1,t_2)=E_{}(z_o,\stackrel{}{r}_{o1},t_1)E_{}^{}(z_o,\stackrel{}{r}_{o2},t_2).$$
(10)
For the rest of this paper we will abandon almost entirely any reference to the time domain and work consistently in frequency domain with the help of Eq. (9) because, as has already been said, this is a natural choice for Synchrotron Radiation applications. In particular, as has already been anticipated, under non-restrictive assumptions on characteristic bandwidths of the process, it is possible to break the correlation function $`\mathrm{\Gamma }_\omega (\omega ,\omega ^{})`$ in space-frequency domain in the product of two factors, the spectral correlation function $`F_\omega (\omega \omega ^{})`$ , and the cross-spectral density of the process $`G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$ MAND . The cross-spectral density can be studied independently at any given frequency giving information on the spatial correlation of the field. Subsequently, the knowledge of the spectral correlation function brings back the full expression for $`\mathrm{\Gamma }_\omega `$.
Substituting Eq. (5) in Eq. (9) one has
$`\mathrm{\Gamma }_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega ,\omega ^{})={\displaystyle \underset{m=1}{\overset{N}{}}}\overline{E}_s(\stackrel{}{\eta }_m,\stackrel{}{l}_m,z_o,\stackrel{}{r}_{o1},\omega )`$ (11)
$`\times {\displaystyle \underset{n=1}{\overset{N}{}}}\overline{E}_s^{}(\stackrel{}{\eta }_n,\stackrel{}{l}_n,z_o,\stackrel{}{r}_{o2},\omega ^{})\mathrm{exp}[i(\omega t_m\omega ^{}t_n)].`$ (12)
Expanding Eq. (12) one has
$`\mathrm{\Gamma }_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega ,\omega ^{})={\displaystyle \underset{m=1}{\overset{N}{}}}\overline{E}_s(\stackrel{}{\eta }_m,\stackrel{}{l}_m,z_o,\stackrel{}{r}_{o1},\omega )`$ (13)
$`\times \overline{E}_s^{}(\stackrel{}{\eta }_m,\stackrel{}{l}_m,z_o,\stackrel{}{r}_{o1},\omega ^{})\mathrm{exp}[i(\omega \omega ^{})t_m]`$ (14)
$`+{\displaystyle \underset{mn}{}}\overline{E}_s(\stackrel{}{\eta }_m,\stackrel{}{l}_m,z_o,\stackrel{}{r}_{o1},\omega )\mathrm{exp}(i\omega t_m)`$ (15)
$`\times \overline{E}_s^{}(\stackrel{}{\eta }_n,\stackrel{}{l}_n,z_o,\stackrel{}{r}_{o2},\omega ^{})\mathrm{exp}(i\omega ^{}t_n).`$ (16)
With the help of Eq. (7) and Eq. (8), the ensemble average $`\mathrm{exp}(i\omega t_k)_t`$ can be written as the Fourier transform of the bunch longitudinal profile function $`F_t(t_k)`$, that is
$$\mathrm{exp}(i\omega t_k)_t=_{\mathrm{}}^{\mathrm{}}𝑑t_kF_t(t_k)e^{i\omega t_k}=F_\omega (\omega ).$$
(17)
Using Eq. (17), Eq. (16) can be written as
$`\mathrm{\Gamma }_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega ,\omega ^{})={\displaystyle \underset{m=1}{\overset{N}{}}}F_\omega (\omega \omega ^{})\overline{E}_s(\stackrel{}{\eta }_m,\stackrel{}{l}_m,z_o,\stackrel{}{r}_{o1},\omega )`$ (18)
$`\times \overline{E}_s^{}(\stackrel{}{\eta }_m,\stackrel{}{l}_m,z_o,\stackrel{}{r}_{o2},\omega ^{})_{\stackrel{}{\eta },\stackrel{}{l}}+{\displaystyle }_{mn}F_\omega (\omega )F_\omega (\omega ^{})`$ (19)
$`\times \overline{E}_s(\stackrel{}{\eta }_m,\stackrel{}{l}_m,z_o,\stackrel{}{r}_{o1},\omega )_{\stackrel{}{\eta },\stackrel{}{l}}\overline{E}_s^{}(\stackrel{}{\eta }_n,\stackrel{}{l}_n,z_o,\stackrel{}{r}_{o2},\omega ^{})_{\stackrel{}{\eta },\stackrel{}{l}},`$ (20)
where $`F_\omega ^{}(\omega ^{})=F_\omega (\omega ^{})`$ because $`F_t`$ is a real function. When the radiation wavelengths of interest are much shorter than the bunch length we can safely neglect the second term on the right hand side of Eq. (20) since the form factor product $`F_\omega (\omega )F_\omega (\omega ^{})`$ goes rapidly to zero for frequencies larger than the characteristic frequency associated with the bunch length: think for instance, at a centimeter long bunch compared with radiation in the Angstrom wavelength range. It should be noted, however, that when the radiation wavelength of interested is longer than the bunch length the second term in Eq. (20) is dominant with respect to the first, because it scales with the number of particles squared: in this case, analysis of the second term leads to a treatment of Coherent Synchrotron Radiation phenomena (CSR). In this paper we will not be concerned with CSR and we will neglect the second term in Eq. (20), assuming that the radiation wavelength of interest is shorter than the bunch length: then, it should be noted that $`F_\omega (\omega \omega ^{})`$ depends on the difference between $`\omega `$ and $`\omega ^{}`$, and the first term cannot be neglected. We can therefore write
$`\mathrm{\Gamma }_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega ,\omega ^{})={\displaystyle \underset{m=1}{\overset{N}{}}}F_\omega (\omega \omega ^{})`$ (21)
$`\times \overline{E}_s(\stackrel{}{\eta }_m,\stackrel{}{l}_m,z_o,\stackrel{}{r}_{o1},\omega )\overline{E}_s^{}(\stackrel{}{\eta }_m,\stackrel{}{l}_m,z_o,\stackrel{}{r}_{o2},\omega ^{})_{\stackrel{}{\eta },\stackrel{}{l}}`$ (22)
$`=NF_\omega (\omega \omega ^{})\overline{E}_s(\stackrel{}{\eta },\stackrel{}{l},z_o,\stackrel{}{r}_{o1},\omega )\overline{E}_s^{}(\stackrel{}{\eta },\stackrel{}{l},z_o,\stackrel{}{r}_{o2},\omega ^{})_{\stackrel{}{\eta },\stackrel{}{l}}.`$ (23)
As one can see from Eq. (23) each electron is correlated just with itself: cross-correlation terms between different electrons was, in fact, included in the second term on the right hand side of Eq. (20), which has been dropped. It is important to note that if the dependence of $`\overline{E}_s`$ on $`\omega `$ and $`\omega ^{}`$ is slow enough, so that $`\overline{E}_s`$ does not vary appreciably on the characteristic scale of $`F_\omega `$, we can substitute $`\overline{E}_s^{}(\stackrel{}{\eta },\stackrel{}{l},z_o,\stackrel{}{r}_{o2},\omega ^{})`$ with $`\overline{E}_s^{}(\stackrel{}{\eta },\stackrel{}{l},z_o,\stackrel{}{r}_{o2},\omega )`$ in Eq. (23).
The situation is depicted in Fig. 5. On the one hand, the characteristic scale of $`F_\omega `$ is given by $`1/\sigma _T`$, where $`\sigma _T`$ is the characteristic bunch duration. On the other hand, the bandwidth of single particle undulator radiation at resonance is given by $`\omega _o/N_w`$, where $`\omega _o`$ is the resonant frequency and $`N_w`$ is the number of undulator periods (of order $`10^210^3`$). In the case of an electron beam the undulator spectrum will exhibit a longer tail, as has been shown in OURS , which guarantees that $`\omega _o/N_w`$ is, indeed, a minimum for the radiation bandwidth, and is the right quantity to be compared with $`1/\sigma _T`$. As an example, for wavelengths of order $`1\AA `$, $`N_w10^3`$ and $`\sigma _T30`$ ps (see PETR ), $`\omega _o/N_w210^{16}`$ Hz which is much larger than $`1/\sigma _T310^{10}`$ Hz. As a result we can simplify Eq. (23) to
$`\mathrm{\Gamma }_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega ,\omega ^{})=NF_\omega (\omega \omega ^{})G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$ (24)
where
$$G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )=\overline{E}_s(\stackrel{}{\eta },\stackrel{}{l},z_o,\stackrel{}{r}_{o1},\omega )\overline{E}_s^{}(\stackrel{}{\eta },\stackrel{}{l},z_o,\stackrel{}{r}_{o2},\omega )_{\stackrel{}{\eta },\stackrel{}{l}}.$$
(25)
Eq. (24) fully characterizes the system under study from a statistical viewpoint. However, in practical situations, the observation plane is behind a monochromator or, equivalently, the detector itself is capable of analyzing the energy of the photons. The presence of a monochromator simply modifies the right hand side Eq. (24) for a factor $`T(\omega )T^{}(\omega ^{})`$, where $`T`$ is the monochromator transfer function:
$`\mathrm{\Gamma }_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega ,\omega ^{})=NF_\omega (\omega \omega ^{})T(\omega )T^{}(\omega ^{})G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega ).`$ (26)
Independently on the characteristics (and even on the presence) of the monochromator, it should be noted that both in Eq. (24) and Eq. (26), correlation in frequency and space are expressed by two separate factors. In particular, in both these equation, spatial correlation is expressed by the cross-spectral density function $`G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$. In other words, we are able to deal separately with spatial and spectral part of the correlation function in space-frequency domain with the only non-restrictive assumption that $`\omega _o/N_w1/\sigma _T`$. From now on we will be concerned with the calculation of the correlation function $`G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$, independently on the shape of the remaining factors on the right hand side of Eq. (26) which can have a simple or a complicated structure, accounting for the characteristics of the monochromator.
Before proceeding with the analysis of $`G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$ though, let us spend some words on these remaining factors; the presence of a monochromator introduces another bandwidth of interest. If we indicate the bandwidth of the monochromator with $`\mathrm{\Delta }\omega _\mathrm{m}`$ and the central frequency of interest at which the monochromator is tuned with $`\omega _o`$ (typically, the undulator resonant frequency), then $`T`$ is peaked around $`\omega _o`$ and goes rapidly to zero as we move out of the range $`(\omega _o\mathrm{\Delta }\omega _\mathrm{m}/2,\omega _o+\mathrm{\Delta }\omega _\mathrm{m}/2)`$. Now, if the characteristic bandwidth of the monochromator, $`\mathrm{\Delta }\omega _\mathrm{m}`$, is large enough so that $`T`$ does not vary appreciably on the characteristic scale of $`F_\omega `$, i.e. $`\mathrm{\Delta }\omega _\mathrm{m}1/\sigma _T`$, then $`F_\omega (\omega \omega ^{})`$ is peaked at $`\omega =\omega ^{}`$. In this case the process resembles more and more a stationary process, although it will be still intrinsically non-stationary. Consider a signal observed for a time much shorter than its duration, but much longer than its coherence time, and such that the stationary model applies to it. Now imagine that we extend the observation time to a duration which is still much shorter than the signal duration, but long enough that we need to account for the intrinsical non-stationarity of the process due to finite signal duration. In this case the stationary model does not apply anymore strictly. To describe this situation, we can define a property weaker than stationarity, but nonetheless very interesting from a physical standpoint: quasi-stationarity. The time domain correlation function (that is, the mutual coherence function) can be written as
$`\mathrm{\Gamma }_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1,t_2)={\displaystyle \frac{N}{(2\pi )^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega {\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega ^{}F_\omega (\omega \omega ^{})T(\omega )T^{}(\omega ^{})`$ (27)
$`\times G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )\mathrm{exp}(i\omega t_1)\mathrm{exp}(i\omega ^{}t_2).`$ (28)
When $`\mathrm{\Delta }\omega _\mathrm{m}1/\sigma _T`$, and with the help of new variables $`\mathrm{\Delta }\omega =\omega \omega ^{}`$ and $`\omega `$, we can simplify Eq. (28) accounting for the fact that $`F_\omega (\omega \omega ^{})`$ is strongly peaked around $`\mathrm{\Delta }\omega =0`$. In fact we can consider $`T(\omega )T^{}(\omega ^{})G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )|T(\omega )|^2G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$, so that we can integrate separately in $`\mathrm{\Delta }\omega `$ and $`\omega `$ to obtain
$`\mathrm{\Gamma }_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1,t_2)={\displaystyle \frac{N}{(2\pi )^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mathrm{\Delta }\omega F_\omega (\mathrm{\Delta }\omega )\mathrm{exp}\left(i\mathrm{\Delta }\omega t_2\right)`$ (29)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\omega |T(\omega )|^2G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )\mathrm{exp}[i\omega (t_1t_2)]`$ (30)
$`=F_t\left(t_2\right)G_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1t_2).`$ (31)
In other words, in the quasi-stationary case, $`\mathrm{\Gamma }_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1,t_2)`$ is split on the product of two factors, a ”reduced mutual coherence function”, that is $`G_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1t_2)`$, and an intensity profile, that is $`F(t_2)`$.
If we now assume $`\mathrm{\Delta }\omega _\mathrm{m}N_w/\omega _o1`$ (that is usually true), $`G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$ $`=G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega _o)`$ is a constant function of frequency within the monochromator line. In this case, $`G_\omega `$ contains all the information about spatial correlations between different point and is, in fact, the quantity of central interest in our study, but it is independent on the frequency $`\omega `$. As a result we have
$$\mathrm{\Gamma }_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1,t_2)=Ng_t(t_1t_2)F_t\left(t_2\right)G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega _o),$$
(32)
which means that the mutual coherence function $`\mathrm{\Gamma }_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1,t_2)`$ is reducible, in the sense that it can be factorized as a product of two factors, the first $`Ng_t(t_1t_2)F_t(t_2)`$, characterizing the temporal coherence and the second $`G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega _o)`$ describing the spatial coherence of the system<sup>2</sup><sup>2</sup>2It is interesting, for the sake of completeness, to discuss the relation between $`G_\omega `$ and the mutual intensity function as usually defined in textbooks GOOD , MAND in quasimonochromatic conditions. The assumption $`\mathrm{\Delta }\omega _\mathrm{m}1/\sigma _T`$ in the limit $`\sigma _T\mathrm{}`$ describes a stationary process. Now letting $`\mathrm{\Delta }\omega _m0`$ slowly enough so that $`\mathrm{\Delta }\omega _\mathrm{m}1/\sigma _T`$, Eq. (31) remains valid while both $`F_\omega `$ and $`|T(\omega )|^2`$ become approximated better and better by Dirac $`\delta `$-functions, $`\delta (\mathrm{\Delta }\omega )`$ and $`\delta (\omega \omega _o)`$, respectively. Then $`\mathrm{\Gamma }_tG_\omega \mathrm{exp}[i\omega _o(t_1t_2)]`$. Aside for an unessential factor, depending on the normalization of $`F_\omega `$ and $`|T(\omega )|^2`$, this relation between $`\mathrm{\Gamma }_t`$ and $`G_\omega `$ allows identification of the mutual intensity function with $`G_\omega `$ as in GOOD , MAND .. This case is of practical importance. In fact for $`1\AA `$ radiation we typically have $`\mathrm{\Delta }\omega _\mathrm{m}/\omega _o10^4÷10^5`$ and $`N_w10^2÷10^3`$, i.e. $`\mathrm{\Delta }\omega _\mathrm{m}N_w/\omega _o1`$. It should be noted that, although Eq. (32) describes the case when $`\mathrm{\Delta }\omega _\mathrm{m}N_w/\omega _o1`$ and $`\mathrm{\Delta }\omega _\mathrm{m}1/\sigma _T`$ , only the former assumption is important for the mutual coherence function to be reducible. In fact, if the former is satisfied but the latter is not, from Eq. (28) one would simply have:
$$\mathrm{\Gamma }_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1,t_2)=N\stackrel{~}{g}_t(t_1,t_2)G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega _o),$$
(33)
that is still reducible.
Eq. (31) contains two important facts:
(a) The temporal correlation function, that is $`G_t(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},t_1t_2)`$, and the spectral density distribution of the source, that is $`H(z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )=N|T(\omega )|^2G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$, form a Fourier pair.
(b) The intensity distribution of the radiation pulse $`F_t(t_2)`$ and the spectral correlation function $`F_\omega (\mathrm{\Delta }\omega )`$ form a Fourier pair.
The statement (a) can be regarded as an analogue, for quasi-stationary sources, of the well-known Wiener-Khinchin theorem, which applies to stationary sources and states that the temporal correlation function and the spectral density are a Fourier pair. Since there is symmetry between time and frequency domains, a ”anti” Wiener-Khinchin theorem must also hold, and can be obtained by the usual Wiener-Khinchin theorem by exchanging frequencies and times. This is simply constituted by the statement (b).
The assumption of quasi-stationarity is not vital for the following of this work, since the cross-spectral density can be studied in any case as a function of frequency. In this respect it should be noted that, although in the large majority of the cases monochromator characteristics are not good enough to allow resolution of the non-stationary process, there are cases when it is not allowed to treat the process as if it were quasi-stationary. For instance, in NOST a particular monochromator is described with a relative resolution of $`10^8`$ at wavelengths of about $`1\AA `$, or $`\omega _o210^{19}`$ Hz. Let us consider, as in NOST , the case of radiation pulses of $`32`$ ps duration. Under the already accepted assumption $`1/\sigma _T\omega _o/N_w`$, we can identify the radiation pulse duration with $`\sigma _T`$. Then we have $`\mathrm{\Delta }\omega _\mathrm{m}210^{11}`$ Hz which is of order of $`2\pi /\sigma _T210^{11}`$ Hz: this means that the monochromator has the capability of resolving the non-stationary processes in the frequency domain. On the contrary, also in this case, the mutual coherence function is reducible, in the sense specified before, because $`\mathrm{\Delta }\omega _\mathrm{m}/\omega _o10^8`$ and $`N_w10^2÷10^3`$, i.e. $`\mathrm{\Delta }\omega _\mathrm{m}N_w/\omega _o1`$. However, such accuracy level is not usual in Synchrotron Radiation experiments.
To sum up, condition $`\omega _o/N_w1/\sigma _T`$ alone allows separate treatment of transverse coherence properties at a given frequency through the function $`G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$. The condition for the monochromator bandwidth $`\mathrm{\Delta }\omega _\mathrm{m}1/\sigma _T`$ defines a quasi-stationary process. If the monochromator bandwidth is such that $`\mathrm{\Delta }\omega _m>\omega _o/N_w`$, the ratio $`N_w/(\omega _o\sigma _T)`$ gives us a direct measure of the accuracy of the stationary approximation which does not depend, in this case, on the presence of the monochromator. When a monochromator is present, with a bandwidth $`\mathrm{\Delta }\omega _m<\omega _o/N_w`$, it is the ratio $`1/(\mathrm{\Delta }\omega _\mathrm{m}\sigma _T)`$ which gives such a measure. The condition $`\mathrm{\Delta }\omega _\mathrm{m}\omega _o/N_w`$ ensures, instead, that the mutual coherence function of the signal is reducible, in the sense specified by Eq. (33). In the following we will need only the first of these conditions, $`\omega _o/N_w1/\sigma _T`$. In fact, this is all we need in order to separate transverse and temporal coherence effects, as shown in Eq. (26).
Once transverse and temporal coherence effects are separated one can focus on the study of transverse coherence through the function $`G_\omega (z_o,\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )`$.
There exists an important class of sources, called quasi-homogeneous. As we will see, quasi-homogeneity is the spatial analogue of quasi-stationarity.
In general, quasi-homogeneous sources are a particular class of Schell’s model sources. Schell’s model sources are defined by the condition that their cross-spectral density at the source plane (that is for a particular value of $`z_o`$) is of the form
$`G_\omega (\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )=\left|E(\stackrel{}{r}_{o1},\omega )\right|^2^{1/2}\left|E(\stackrel{}{r}_{o2},\omega )\right|^2^{1/2}`$ (34)
$`\times g(\stackrel{}{r}_{o2}\stackrel{}{r}_{o1},\omega ),`$ (35)
where $`g(\stackrel{}{r}_{o2}\stackrel{}{r}_{o1},\omega )`$ is the spectral degree of coherence (that is normalized to unity by definition, i.e. $`g(0,\omega )=1`$)<sup>3</sup><sup>3</sup>3Sometimes, loosely speaking, we will refer to $`g`$ as to ”the cross-spectral density”, or to ”the field correlation function” the difference being just a normalization factor.. Equivalently one may simply define Schell’s model sources using the condition that the spectral degree of coherence depends on the positions across the source only through the difference $`\stackrel{}{r}_{o2}\stackrel{}{r}_{o1}`$ (see MAND , 5.2.2) from which Eq. (35) follows.
Quasi-homogeneous sources are Schell’s sources obeying the following extra-assumption: the spectral density $`G_\omega (\stackrel{}{r}_o,\stackrel{}{r}_o,\omega )`$ at the source plane, considered as a function of $`\stackrel{}{r}_o`$, varies so slowly with the position that it is approximatively constant over distances across the source, which are of the order of the correlation length $`\mathrm{\Delta }`$ (that is the effective width of $`|G_\omega (\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )|`$, see Fig. 6 for a qualitative illustration in one dimension).
Since for quasi-homogeneous sources the spectral density is assumed to vary slowly we are allowed to make the approximation
$`G_\omega (\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )=I(\stackrel{}{r}_{o1},\omega )g(\stackrel{}{r}_{o2}\stackrel{}{r}_{o1},\omega ),`$ (36)
where
$$I(\stackrel{}{r}_{o1},\omega )=\left|E(\stackrel{}{r}_{o1},\omega )\right|^2$$
(37)
is the field intensity distribution. A schematic illustration of the measurement of the cross-spectral density of a undulator source is given in Fig. 7.
We reported definitions and differences between Schell’s sources and quasi-homogeneous sources as treated in MAND in order to review some conventional language. However, in our paper we will study and classify sources with the help of parameters from dimensional analysis of the problem. In particular, in the following we will introduce quantities that model, in dimensionless units, the electron beam dimension and divergence in the $`x`$ and $`y`$ directions ($`\widehat{N}_{x,y}`$ and $`\widehat{D}_{x,y}`$, respectively). When some of these parameters are much larger or much smaller than unity, in certain particular combinations discussed in the following part of this work, we will be able to point out simplifications of analytical expressions. These simplifications do not depend, in general, on the fact that the source is quasi-homogeneous or not, but simply on the fact that some of the above parameters are large or small. It will be possible to describe some parameter combination in terms of Schell’s or quasi-homogeneous model, but this will not be the case, in general. In order to link the physical properties of certain kind of sources to a certain range of parameters we will find it convenient to extend the concept of quasi-homogeneity to the concept of ”weak quasi-homogeneity”. It is better to familiarize with our new definition already here: a given wavefront at fixed position $`\widehat{z}_o`$ will be said to be weakly quasi-homogeneous, by definition, when the modulus of the spectral degree of coherence $`|g|`$ depends on the position across the source only through the difference $`\stackrel{}{r}_{o2}\stackrel{}{r}_{o1}`$, i.e. $`|g|=w(\stackrel{}{r}_{o2}\stackrel{}{r}_{o1},\omega )`$. This is equivalent to generalize Eq. (36) to
$`|G_\omega (\stackrel{}{r}_{o1},\stackrel{}{r}_{o2},\omega )|=I(\stackrel{}{r}_{o1},\omega )w(\stackrel{}{r}_{o2}\stackrel{}{r}_{o1},\omega ).`$ (38)
As a remark to both the definitions of quasi-homogeneity and weak quasi-homogeneity, it should be noted that they only involve conditions on the cross-spectral densities through Eq. (36) or Eq. (38): one can apply these definitions to any wavefront at any position $`z_o`$. This is consistent with what has been remarked in the Introduction: the choice of a source plane down the beamline is just a convention. Therefore, our definition of weak quasi-homogeneity is completely separated from the concept of source plane. It may seem, at first glance, that the definition of ”weak quasi-homogeneity” is somehow artificial but it is, on the contrary, very convenient from a practical viewpoint. In fact, in any coherent experiment, the specimen is illuminated by coherent light from some kind of aperture, or diaphragm. Think, for instance, to the usual process of selection of transversely coherent light through a spatial filter, where a diaphragm is placed downstream a pinhole. Physically, when the modulus of the spectral degree of coherence depends only on the coordinate difference $`\stackrel{}{r}_{o2}\stackrel{}{r}_{o1}`$ the coherence properties of the beam do not depend on the position of the diaphragm with respect to the transverse coordinate of the center of the pinhole, which we may imagine on the $`z`$ axis: for instance, the coherence length will depend only on $`\stackrel{}{r}_{o2}\stackrel{}{r}_{o1}`$ and not on the average position $`(\stackrel{}{r}_{o2}+\stackrel{}{r}_{o1})/2`$. Mathematically, as has already been said, the definition of weak quasi-homogeneity is linked with particular combination of small and large parameters which will lead to simplifications of equations and to analytical treatment of several interesting cases.
After having introduced the definition of ”weak quasi-homogeneity” we should go back to the concept of usual quasi-homogeneity to describe a vary particular feature of it: quasi-homogeneity can be regarded as the spatial equivalent of quasi-stationarity, as anticipated before. Exactly as the time domain has a reciprocal description in terms of frequency, the space domain has a reciprocal description in terms of transverse (two-dimensional) wave vectors. However, since the frequency is fixed, the ratio between the horizontal or vertical component of the wave vector and the longitudinal wave number is representative of the propagation angle of a plane wave at fixed frequency. Therefore any given signal on a two-dimensional plane can be represented in terms of superposition of plane waves with the same frequency and different angles of propagation, which goes under the name of angular spectrum. This defines an angular domain which is the reciprocal of the space domain. Intuitively the angular spectrum representation constitutes a picture of the effects of propagation in the far zone whereas the near zone is described by the space domain. In this Section, an analogous of the Wiener-Khinchin theorem and its reciprocal form for quasi-stationary processes has been discussed. Substitution of times with position vectors and frequencies with angular vectors allows to derive similar statements for the near (space domain) and far (angular domain) zone MAND . If we use $`\stackrel{}{\theta }=\stackrel{}{r}_o/z_o`$ (and $`\stackrel{}{\theta }_{1,2}=\stackrel{}{r}_{o1,2}/z_o`$) as variables to describe radiant intensity and cross-spectral density in the far field, and if we identify points on the source plane with $`\stackrel{}{r}_{}`$ (and $`\stackrel{}{r}_{1,2}`$), we have:
(a’) The cross-spectral density of the field at the source plane $`g(\stackrel{}{r}_2\stackrel{}{r}_1)`$ and the angular distribution of the radiant intensity $`I(\stackrel{}{\theta })`$ are a Fourier Pair.
(b’) The cross-spectral density of the far field $`g(\stackrel{}{\theta }_2\stackrel{}{\theta }_1)`$ and the source-intensity distribution $`I(\stackrel{}{r}_{})`$ are, apart for a simple geometrical phase factor, a Fourier Pair.
The statement (b’) can be regarded as an analogue, for quasi-homogeneous sources, of the far-zone form of the van Cittert-Zernike theorem. The statement (a’) instead, is due to the symmetry between space and angle domains, and can be seen as an ”anti” VCZ theorem. This discussion underlines the link between the VCZ theorem and the Wiener-Khinchin theorem.
Many undulator radiation sources are quasi-homogeneous sources in the usual way. In the case of a quasi-homogeneous source of typical linear dimension $`d`$, the angular spectrum at distance $`z_od\omega \mathrm{\Delta }/c`$ is expected to exhibit speckles with typical linear dimension $`z_o(d\omega /c)^1`$, as illustrated in Fig. 8, which shows an intuitive picture of the propagation of transverse coherence.
They generate fields which are relatively simple to analyze mathematically and still rich in physical features. Although both VCZ and ”anti” VCZ theorem are based on the assumption of usual quasi-homogeneity we will see that these are often applicable, at least in some sense, also in the case of ”weakly quasi-homogeneous” wavefronts.
## 3 Evolution of the cross-spectral density function through the undulator beamline
In our work OURS we presented an expression for the reduced field $`\stackrel{~}{E}_{}(\omega )`$ $`=\overline{E}_{}(\omega )\mathrm{exp}(i\omega z_o/c)`$ of a single particle with offset $`\stackrel{}{l}`$ and deflection $`\stackrel{}{\eta }`$ with respect to the optical axis $`z`$ in an undulator. In order to derive our result, we used a Green’s function approach to solve the paraxial Maxwell equations for the Fourier transform of the electric field and we took advantage of a consistent use of the resonance approximation. The field $`\stackrel{~}{E}_{}`$ differs from $`\overline{E}_{}`$ for a phase factor which depends on the variable $`z_o`$ and on the frequency $`\omega `$ only: therefore, the use of one expression instead of the other in the equation for $`G`$ does not change the result. In OURS , we presented results in normalized units in the far field zone for a particle with offset and deflection. Based on that work we can calculate the field in normalized units for a particle with offset and deflection at any distance from the exit of the undulator, where the center of the undulator is taken at $`z=0`$, as specified in Fig. 9:
$`\widehat{E}_s=\widehat{z}_o{\displaystyle _{1/2}^{1/2}}𝑑\widehat{z}^{}{\displaystyle \frac{1}{\widehat{z}_o\widehat{z}^{}}}\mathrm{exp}\left\{i\left[\left(\widehat{C}+{\displaystyle \frac{\stackrel{}{\widehat{\eta }}^2}{2}}\right)\widehat{z}^{}+{\displaystyle \frac{\left(\stackrel{}{\widehat{r}}_o\stackrel{}{\widehat{l}}\stackrel{}{\widehat{\eta }}\widehat{z}^{}\right)^2}{2(\widehat{z}_o\widehat{z}^{})}}\right]\right\}.`$ (39)
Eq. (39) is valid for the system tuned at resonance with the fundamental harmonic $`\omega _o`$. This means that we are considering a large number of undulator periods $`N_w1`$ and that we are looking at frequencies near the fundamental and at angles within the main lobe of the directivity diagram. In this situation one can neglect the vertical $`y`$-polarization of the field with an accuracy $`(4\pi N_w)^1`$. This constitutes a great simplification of the problem since, at any position of the observer, we may consider the electric field Fourier transform, $`\widehat{E}_s`$, as a complex scalar quantity corresponding to the surviving $`x`$-polarization component of the original vector quantity. Normalized units were defined as
$`\widehat{E}_s={\displaystyle \frac{c^2z_o\gamma \stackrel{~}{E}_s}{K\omega eL_wA_{JJ}}},`$ (40)
$`\stackrel{}{\widehat{\eta }}=\stackrel{}{\eta }\sqrt{{\displaystyle \frac{\omega L_w}{c}}},`$ (41)
$`\widehat{C}=L_wC=2\pi N_w{\displaystyle \frac{\omega \omega _o}{\omega _o}},`$ (42)
$`\stackrel{}{\widehat{r}}_o=\stackrel{}{r}_o\sqrt{{\displaystyle \frac{\omega }{L_wc}}},`$ (43)
$`\stackrel{}{\widehat{l}}=\stackrel{}{l}\sqrt{{\displaystyle \frac{\omega }{L_wc}}},`$ (44)
$`\widehat{z}={\displaystyle \frac{z}{L_w}}.`$ (45)
$`K`$ being the deflection parameter, $`L_w`$ being the undulator length,
$$A_{JJ}=J_0\left(\frac{K^2}{4+2K^2}\right)J_1\left(\frac{K^2}{4+2K^2}\right),$$
(46)
$$\omega _o=\frac{4\pi c\gamma ^2}{\lambda _w\left(1+K^2/2\right)}$$
(47)
being the resonant frequency, $`J_n`$ the Bessel function of the first kind of order $`n`$, $`\lambda _w`$ the undulator period, $`(e)`$ the electron charge and $`\gamma `$ the relativistic Lorentz factor. $`\stackrel{}{\widehat{l}}`$ is the normalized offset in the center of the undulator. Finally, the parameter $`\widehat{C}`$ represents the normalized detuning, which accounts for small deviation in frequency from resonance.
As it is shown in Appendix B, Eq. (39) can be rewritten as
$`\widehat{E}_s={\displaystyle _{1/2}^{1/2}}{\displaystyle \frac{\widehat{z}_od\widehat{z}^{}}{\widehat{z}_o\widehat{z}^{}}}\mathrm{exp}\left\{i\left[\mathrm{\Phi }_U+\widehat{C}\widehat{z}^{}+{\displaystyle \frac{\widehat{z}_o\widehat{z}^{}}{2(\widehat{z}_o\widehat{z}^{})}}\left(\stackrel{}{\widehat{\theta }}{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }}\right)^2\right]\right\}`$ (48)
where
$`\stackrel{}{\widehat{\theta }}={\displaystyle \frac{\stackrel{}{\widehat{r}}_o}{\widehat{z}_o}}`$ (49)
represents the observation angle and $`\mathrm{\Phi }_U`$ is given by
$$\mathrm{\Phi }_U=\left(\stackrel{}{\widehat{\theta }}\frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}\right)^2\frac{\widehat{z}_o}{2}.$$
(50)
Eq. (48) is of the form
$$\widehat{E}_s(\widehat{C},\widehat{z}_o,\stackrel{}{\widehat{\theta }}\frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}\stackrel{}{\widehat{\eta }})=\mathrm{exp}(i\mathrm{\Phi }_U)S[\widehat{C},\widehat{z}_o,\left(\stackrel{}{\widehat{\theta }}\frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}\stackrel{}{\widehat{\eta }}\right)^2].$$
(51)
Starting from the next Section we will restrict our attention to the case $`\widehat{C}=0`$ for simplicity. Therefore, it may be interesting to note that in the particular case $`\widehat{C}=0`$, the function $`S`$ can be represented in terms of the exponential integral function Ei as:
$$S(0,\widehat{z}_o,\zeta ^2)=\mathrm{exp}(i\widehat{z}_o\zeta ^2/2)\widehat{z}_o\left[\mathrm{Ei}\left(\frac{i\widehat{z}_o^2\zeta ^2}{1+2\widehat{z}_o}\right)\mathrm{Ei}\left(\frac{i\widehat{z}_o^2\zeta ^2}{1+2\widehat{z}_o}\right)\right]$$
(52)
It is easy to show that the expression for the function $`S()`$ reduces to a $`\mathrm{sinc}()`$ function as $`\widehat{z}_o1`$ . In fact, in this limiting case, the expression for the electric field from a single particle, given in Eq. (39) is simplified to
$`\widehat{E}_s=\mathrm{exp}(i\mathrm{\Phi }_U){\displaystyle _{1/2}^{1/2}}𝑑\widehat{z}^{}\mathrm{exp}\left\{i\widehat{z}^{}\left[\widehat{C}+{\displaystyle \frac{1}{2}}\left(\stackrel{}{\widehat{\theta }}{\displaystyle \frac{\stackrel{}{\widehat{l}_x}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }}\right)^2\right]\right\},`$ (53)
Eq. (53) can be integrated analytically giving
$$\widehat{E}_s=\mathrm{exp}(i\mathrm{\Phi }_U)\mathrm{sinc}\left(\frac{\widehat{C}}{2}+\frac{\zeta ^2}{4}\right),$$
(54)
where
$$\zeta =\stackrel{}{\widehat{\theta }}\frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}\stackrel{}{\widehat{\eta }}.$$
(55)
A comparison between $`\mathrm{sinc}(\zeta ^2/4)`$ and the real and imaginary parts of $`S(0,\widehat{z}_o,\zeta ^2)`$ for $`\widehat{z}_o=1`$ is given in Fig. 10.
Let us now go back to the general case for $`\widehat{z}_o1/2`$ and use Eq. (48) to calculate the cross-spectral density. The cross-spectral density $`G_\omega `$ is given Eq. (25) in dimensional units and as a function of dimensional variables. Since the field in Eq. (48) is given in normalized units and as a function of normalized variables $`\widehat{z}_o`$, $`\stackrel{}{\widehat{\theta }}_{x,y}`$ and $`\widehat{C}`$, it is convenient to introduce a version of $`G_\omega `$ defined by means of the field in normalized units:
$`\widehat{G}(\widehat{z}_o,\stackrel{}{\widehat{\theta }}_1,\stackrel{}{\widehat{\theta }}_2,\widehat{C})=\widehat{E}_s(\widehat{C},\widehat{z}_o,\stackrel{}{\widehat{\theta }}_1{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }})\widehat{E}_s^{}(\widehat{C},\widehat{z}_o,\stackrel{}{\widehat{\theta }}_2{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }})_{\stackrel{}{\eta },\stackrel{}{l}}.`$ (56)
(57)
Transformation of $`G_\omega `$ in Eq. (25) to $`\widehat{G}`$ (and viceversa) can be easily performed shifting from dimensional to normalized variables and multiplying $`G_\omega `$ by an inessential factor:
$$\widehat{G}=\left(\frac{c^2z_o\gamma }{K\omega eL_wA_{JJ}}\right)^2G_\omega .$$
(58)
As a result, we can always use $`\widehat{G}`$ in stance of $`G_\omega `$. Substituting Eq. (48) in Eq. (57) we obtain:
$`\widehat{G}(\widehat{z}_o,\stackrel{}{\widehat{\theta }}_1,\stackrel{}{\widehat{\theta }}_2,\widehat{C})=S[\widehat{C},\widehat{z}_o,(\stackrel{}{\widehat{\theta }}_1{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }})^2]S^{}[\widehat{C},\widehat{z}_o,(\stackrel{}{\widehat{\theta }}_2{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }})^2]`$ (59)
$`\times \mathrm{exp}\left\{i[(\stackrel{}{\widehat{\theta }}_1{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}})^2(\stackrel{}{\widehat{\theta }}_2{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}})^2]{\displaystyle \frac{\widehat{z}_o}{2}}\right\}_{\stackrel{}{\eta },\stackrel{}{l}}`$ (60)
Expanding the exponent in the exponential factor in the right hand side of Eq. (60), it is easy to see that terms in $`\widehat{l}_{x,y}^2`$ cancel out. Terms in $`\widehat{\theta }_{x,y}^2`$ contribute for a common factor, and only linear terms in $`\widehat{l}_{x,y}`$ remain inside the ensemble average sign. Substitution of the ensemble average with integration over the beam distribution function leads to
$`\widehat{G}(\widehat{z}_o,\stackrel{}{\widehat{\theta }}_1,\stackrel{}{\widehat{\theta }}_2,\widehat{C})=\mathrm{exp}\left[i\left(\stackrel{}{\widehat{\theta }}_1^2\stackrel{}{\widehat{\theta }}_2^2\right){\displaystyle \frac{\widehat{z}_o}{2}}\right]{\displaystyle 𝑑\stackrel{}{\widehat{\eta }}𝑑\stackrel{}{\widehat{l}}F_{\stackrel{}{\widehat{\eta }},\stackrel{}{\widehat{l}}}(\stackrel{}{\widehat{\eta }},\stackrel{}{\widehat{l}})\mathrm{exp}\left[i(\stackrel{}{\widehat{\theta }}_2\stackrel{}{\widehat{\theta }_1})\stackrel{}{\widehat{l}}\right]}`$ (61)
$`\times S[\widehat{C},\widehat{z}_o,\left(\stackrel{}{\widehat{\theta }}_1{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }}\right)^2]S^{}[\widehat{C},\widehat{z}_o,\left(\stackrel{}{\widehat{\theta }}_2{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }}\right)^2].`$ (62)
Here integrals $`d\stackrel{}{\widehat{\eta }}`$ and in $`d\stackrel{}{\widehat{l}}`$ are to be intended as integrals over the entire plane spanned by the $`\stackrel{}{\widehat{\eta }}`$ and $`\stackrel{}{\widehat{l}}`$ vectors. Eq. (62) is very general and can be used as a starting point for computer simulations.
We already assumed that the distribution in the horizontal and vertical planes are not correlated, so that $`F_{\stackrel{}{\widehat{\eta }},\stackrel{}{\widehat{l}}}=F_{\widehat{\eta }_x,\widehat{l}_x}F_{\widehat{\eta }_y,\widehat{l}_y}`$. If the transverse phase space is specified at position $`\widehat{z}_o=0`$ corresponding to the minimal values of the $`\beta `$-functions, we can write $`F_{\widehat{\eta }_x,\widehat{l}_x}=F_{\widehat{\eta }_x}F_{\widehat{l}_x}`$ and $`F_{\widehat{\eta }_y,\widehat{l}_y}=F_{\widehat{\eta }_y}F_{\widehat{l}_y}`$ with
$`F_{\widehat{\eta }_x}(\widehat{\eta }_x)={\displaystyle \frac{1}{\sqrt{2\pi \widehat{D}_x}}}\mathrm{exp}\left({\displaystyle \frac{\widehat{\eta }_x^2}{2\widehat{D}_x}}\right),`$ (63)
$`F_{\widehat{\eta }_y}(\widehat{\eta }_y)={\displaystyle \frac{1}{\sqrt{2\pi \widehat{D}_y}}}\mathrm{exp}\left({\displaystyle \frac{\widehat{\eta }_y^2}{2\widehat{D}_y}}\right),`$ (64)
$`F_{\widehat{l}_x}(\widehat{l}_x)={\displaystyle \frac{1}{\sqrt{2\pi \widehat{N}_x}}}\mathrm{exp}\left({\displaystyle \frac{\widehat{l}_x^2}{2\widehat{N}_x}}\right),`$ (65)
$`F_{\widehat{l}_y}(\widehat{l}_y)={\displaystyle \frac{1}{\sqrt{2\pi \widehat{N}_y}}}\mathrm{exp}\left({\displaystyle \frac{\widehat{l}_y^2}{2\widehat{N}_y}}\right).`$ (66)
From Eq. (45) and Eq. (49) it is easy to see that
$$\widehat{D}_{x,y}=\sigma _{x^{},y^{}}^2\frac{\omega L_w}{c}$$
(67)
$$\widehat{N}_{x,y}=\sigma _{x,y}^2\frac{\omega }{cL_w}$$
(68)
where $`\sigma _{x,y}`$ and $`\sigma _{x^{},y^{}}`$ are the rms transverse bunch dimension and angular spread. Parameters $`\widehat{N}_{x,y}`$ will be indicated as the beam diffraction parameters and are, in fact, analogous to Fresnel numbers and correspond to the normalized square of the electron beam sizes, whereas $`\widehat{D}_{x,y}`$ represent the normalized square of the electron beam divergences.
Substitution of relations (66) in Eq. (62) yields, at perfect resonance ($`\widehat{C}=0`$):
$`\widehat{G}(\widehat{z}_o,\stackrel{}{\widehat{\theta }}_1,\stackrel{}{\widehat{\theta }}_2)={\displaystyle \frac{\mathrm{exp}\left[i\left(\stackrel{}{\widehat{\theta }}_1^2\stackrel{}{\widehat{\theta }}_2^2\right)\widehat{z}_o/2\right]}{4\pi ^2\sqrt{\widehat{D}_x\widehat{D}_y\widehat{N}_x\widehat{N}_y}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\eta }_x\mathrm{exp}\left({\displaystyle \frac{\widehat{\eta }_x^2}{2\widehat{D}_x}}\right)`$ (69)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\eta }_y\mathrm{exp}({\displaystyle \frac{\widehat{\eta }_y^2}{2\widehat{D}_y}}){\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{l}_x\mathrm{exp}({\displaystyle \frac{\widehat{l}_x^2}{2\widehat{N}_x}}){\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{l}_y\mathrm{exp}({\displaystyle \frac{\widehat{l}_y^2}{2\widehat{N}_y}})`$ (70)
$`S[\widehat{z}_o,\left(\stackrel{}{\widehat{\theta }}_1{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }}\right)^2]S^{}[\widehat{z}_o,\left(\stackrel{}{\widehat{\theta }}_2{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }}\right)^2]\mathrm{exp}\left[i(\stackrel{}{\widehat{\theta }}_2\stackrel{}{\widehat{\theta }_1})\stackrel{}{\widehat{l}}\right].`$ (71)
For notational simplicity, in Eq. (71) we have substituted the proper notation $`\widehat{G}(\widehat{z}_o,\stackrel{}{\widehat{\theta }}_1,\stackrel{}{\widehat{\theta }}_2,\widehat{C})`$ with the simplified dependence $`\widehat{G}(\widehat{z}_o,\stackrel{}{\widehat{\theta }}_1,\stackrel{}{\widehat{\theta }}_2)`$ because we will be treating the case $`\widehat{C}=0`$ only. Consistently, also $`S[\widehat{z}_o,(\stackrel{}{\widehat{\theta }}\stackrel{}{\widehat{l}}/\widehat{z}_o\stackrel{}{\widehat{\eta }})^2]`$ is to be understood as a shortcut notation for $`S[\widehat{C},\widehat{z}_o,(\stackrel{}{\widehat{\theta }}\stackrel{}{\widehat{l}}/\widehat{z}_o\stackrel{}{\widehat{\eta }})^2]`$ calculated at $`\widehat{C}=0`$.
## 4 Undulator radiation as a quasi-homogeneous source
When describing physical principles it is always important to find a model which provides the possibility of an analytical description without loss of essential information about the feature of the random process.
In order to get a feeling for some realistic magnitude of parameters we start noting that the geometrical emittances of the electron beam are simply given by $`ϵ_{x,y}=\sigma _{x,y}\sigma _{x^{},y^{}}`$. Here they will be normalized as $`\widehat{ϵ}_{x,y}=2\pi ϵ_{x,y}/\lambda `$. Then $`\sigma _{x,y}^2=\beta _{x,y}^oϵ_{x,y}`$, where $`\beta _{x,y}^o`$ are the minimal values of the horizontal and vertical betatron functions. In this paper we will assume that the betatron functions will have their minimal value at the undulator center. Therefore we have $`ϵ_{x,y}=\sigma _{x^{},y^{}}^2\beta _{x,y}^o`$ or, in normalized units, $`\widehat{ϵ}_{x,y}=\widehat{D}_{x,y}\widehat{\beta }_{x,y}`$, where $`\widehat{\beta }_{x,y}=\beta _{x,y}^o/L_w`$. Equivalently we can write $`\widehat{ϵ}_{x,y}=\sqrt{\widehat{D}_{x,y}\widehat{N}_{x,y}}`$. It follows that $`\widehat{N}_{x,y}=\widehat{ϵ}_{x,y}\widehat{\beta }_{x,y}`$. Now taking $`\lambda =1\AA `$, $`ϵ_x=1÷3`$ nm, $`ϵ_y=10^2ϵ_x`$ and $`\beta _x^o=10^1÷10L_w`$ one obtains, in normalized units, $`\widehat{ϵ}_x=10^2÷310^2`$, $`\widehat{ϵ}_y=1÷3`$ and $`\widehat{\beta }_x=10^1÷10`$: therefore $`\widehat{D}_x1`$, $`\widehat{N}_x1`$. This is always the case in situations of practical interest, with $`\widehat{N}_x`$ which may range from values much smaller to much larger than $`\widehat{D}_x`$.
Assuming $`\widehat{D}_x1`$ and $`\widehat{N}_x1`$, independently on the values of $`\widehat{D}_y`$ and $`\widehat{N}_y`$, introduces simplifications in the expression for the cross-correlation function and allows further analytical investigations. As we will see, in particular, a model of the electron beam based on these assumptions contains (but it is not limited to) the class of quasi-homogeneous sources discussed in Section 2.2.
In the next Sections 4.1 and 4.2 and later in Section 5, we will see what are the conditions in terms of the dimensionless parameters $`\widehat{N}_{x,y}`$ and $`\widehat{D}_{x,y}`$ for some undulator radiation wavefront at position $`\widehat{z}_o`$, to be quasi-homogeneous in the usual and in the weak sense (according to the definition in Section 2.2), we will justify the introduction of the concept of weak quasi-homogeneity itself and we will discuss the applicability regions of the VCZ (and ”anti” VCZ) theorem. Then, in Section 6, we will also discuss some case characterized by non weakly quasi-homogenous fields.
### 4.1 A simple model
To provide a first analysis of the problem we adopt some simplifying assumptions that are only occasionally met in practice.
As already assumed vertical emittance is much smaller than horizontal emittance. For notational simplicity we will make the assumptions $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$. This means that we theoretically assume $`\widehat{\eta }_y1`$ and $`\widehat{l}_y1`$. As a result, the terms in $`\widehat{\eta }_y`$ and $`\widehat{l}_y`$ can be neglected in the $`S()`$ term on the right hand side of Eq. (71). Although this model includes obvious schematization it is still close to reality in many situations, and it is only to be considered as a provisory model for physical understanding to be followed, below, by more comprehensive generalizations. In this Section we will restrict our attention at the correlation function for $`\widehat{\theta }_{y1}=\widehat{\theta }_{y2}`$ that is on any horizontal plane. Here again, for notational simplicity, we will substitute the proper notation $`\widehat{G}(\widehat{z}_o,\widehat{\theta }_{x1},0,\widehat{\theta }_{x2},0)`$ with $`\widehat{G}(\widehat{z}_o,\widehat{\theta }_{x1},\widehat{\theta }_{x2})`$. Eq. (71) can be greatly simplified leading to
$`\widehat{G}(\widehat{z}_o,\widehat{\theta }_{x1},\widehat{\theta }_{x2})={\displaystyle \frac{1}{2\pi \sqrt{\widehat{D}_x\widehat{N}_x}}}\mathrm{exp}\left[i\left(\widehat{\theta }_{x1}^2\widehat{\theta }_{x2}^2\right)\widehat{z}_o/2\right]`$ (72)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\eta }_x\mathrm{exp}({\displaystyle \frac{\widehat{\eta }_x^2}{2\widehat{D}_x}}){\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{l}_x\mathrm{exp}({\displaystyle \frac{\widehat{l}_x^2}{2\widehat{N}_x}})\mathrm{exp}\left[i(\widehat{\theta }_{x2}\widehat{\theta }_{x1})\widehat{l}_x\right]`$ (73)
$`\times \left\{S[\widehat{z}_o,\left(\widehat{\theta }_{x1}\widehat{l}_x/\widehat{z}_o\widehat{\eta }_x\right)^2]S^{}[\widehat{z}_o,\left(\widehat{\theta }_{x2}\widehat{l}_x/\widehat{z}_o\widehat{\eta }_x\right)^2]\right\}.`$ (74)
Let us now introduce
$$\mathrm{\Delta }\widehat{\theta }=\frac{\widehat{\theta }_{x1}\widehat{\theta }_{x2}}{2}$$
(75)
$$\overline{\theta }=\frac{\widehat{\theta }_{x1}+\widehat{\theta }_{x2}}{2}$$
(76)
With this variables redefinition we obtain
$`\widehat{G}(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })={\displaystyle \frac{1}{2\pi \sqrt{\widehat{D}_x\widehat{N}_x}}}\mathrm{exp}\left(i2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\right){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\eta }_x\mathrm{exp}\left({\displaystyle \frac{\widehat{\eta }_x^2}{2\widehat{D}_x}}\right)`$ (77)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{l}_x\mathrm{exp}({\displaystyle \frac{\widehat{l}_x^2}{2\widehat{N}_x}})\mathrm{exp}[2i\mathrm{\Delta }\widehat{\theta }\widehat{l}_x]\{S[\widehat{z}_o,(\overline{\theta }+\mathrm{\Delta }\widehat{\theta }\widehat{l}_x/\widehat{z}_o\widehat{\eta }_x)^2]`$ (78)
$`\times S^{}[\widehat{z}_o,(\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{l}_x/\widehat{z}_o\widehat{\eta }_x)^2]\}.`$ (79)
A double change of variables $`\widehat{\eta }_x\widehat{\eta }+\overline{\theta }`$ followed by $`\widehat{l}_x/\widehat{z}_o\widehat{\varphi }\widehat{\eta }`$ yields
$`\widehat{G}(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })={\displaystyle \frac{\mathrm{exp}\left(i2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\right)}{2\pi \sqrt{\widehat{D}\widehat{N}/\widehat{z}_o^2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\eta }\mathrm{exp}\left({\displaystyle \frac{(\widehat{\eta }+\overline{\theta })^2}{2\widehat{D}}}+2i\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\widehat{\eta }\right)`$ (80)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }\mathrm{exp}({\displaystyle \frac{(\widehat{\varphi }\widehat{\eta })^2}{2\widehat{N}/\widehat{z}_o^2}})\mathrm{exp}(2i\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\widehat{\varphi })`$ (81)
$`\times S^{}[\widehat{z}_o,(\widehat{\varphi }\mathrm{\Delta }\widehat{\theta })^2]S[\widehat{z}_o,(\widehat{\varphi }+\mathrm{\Delta }\widehat{\theta })^2].`$ (82)
where we have posed $`\widehat{D}=\widehat{D}_x`$ and $`\widehat{N}=\widehat{N}_x`$ for notational simplicity. Eq. (82) can be also written as
$`\widehat{G}(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })={\displaystyle \frac{\mathrm{exp}\left(i2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\right)}{2\pi \sqrt{\widehat{D}\widehat{N}/\widehat{z}_o^2}}}\mathrm{exp}({\displaystyle \frac{\overline{\theta }^2}{2\widehat{D}}}){\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }[\mathrm{exp}({\displaystyle \frac{\widehat{\varphi }^2}{2\widehat{N}/\widehat{z}_o^2}})`$ (83)
$`\times \mathrm{exp}\left(2i\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\widehat{\varphi }\right)S^{}[\widehat{z}_o,(\widehat{\varphi }\mathrm{\Delta }\widehat{\theta })^2]S[\widehat{z}_o,(\widehat{\varphi }+\mathrm{\Delta }\widehat{\theta })^2]`$ (84)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\eta }\mathrm{exp}({\displaystyle \frac{\widehat{N}/\widehat{z}_o^2+\widehat{D}}{2\widehat{D}\widehat{N}/\widehat{z}_o^2}}\widehat{\eta }^2+{\displaystyle \frac{\widehat{\varphi }}{\widehat{N}/\widehat{z}_o^2}}\widehat{\eta }{\displaystyle \frac{\overline{\theta }}{\widehat{D}}}\widehat{\eta }+2i\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\widehat{\eta })].`$ (85)
The integral in $`\widehat{\eta }`$ can be performed analytically thus leading to
$`\widehat{G}(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })={\displaystyle \frac{\mathrm{exp}\left(i2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\right)}{\sqrt{2\pi (\widehat{N}/\widehat{z}_o^2+\widehat{D})}}}`$ (86)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }^2+4\widehat{N}\mathrm{\Delta }\widehat{\theta }^2\widehat{D}+4i(\widehat{N}/\widehat{z}_o)\overline{\theta }\mathrm{\Delta }\widehat{\theta }}{2(\widehat{N}/\widehat{z}_o^2+\widehat{D})}}\right]`$ (87)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }\mathrm{exp}[{\displaystyle \frac{\widehat{\varphi }^2+2\widehat{\varphi }\left(\overline{\theta }+2i(\widehat{N}/\widehat{z}_o)\mathrm{\Delta }\widehat{\theta }\right)}{2(\widehat{N}/\widehat{z}_o^2+\widehat{D})}}]S^{}[\widehat{z}_o,(\widehat{\varphi }\mathrm{\Delta }\widehat{\theta })^2]`$ (88)
$`\times S[\widehat{z}_o,(\widehat{\varphi }+\mathrm{\Delta }\widehat{\theta })^2].`$ (89)
It is important to remember again that an asymptotic formula for $`\widehat{z}_o1`$ can be obtained from Eq. (89) simply substituting $`S[\widehat{z}_o,(\widehat{\varphi }\pm \mathrm{\Delta }\widehat{\theta })^2]`$ with $`\mathrm{sinc}[(\widehat{\varphi }\pm \mathrm{\Delta }\widehat{\theta })^2/4]`$. Then, it is easy to understand that $`S`$ is bound to go to zero for values of $`(\widehat{\varphi }\pm \mathrm{\Delta }\widehat{\theta })^2`$ larger than unity, exactly as the asymptotic terms in $`\mathrm{sinc}()`$ would do. In fact, once $`\widehat{C}`$ is set to zero, $`S`$ depends parametrically on the normalized distance $`\widehat{z}_o`$ alone, that is $`S=S[\widehat{z}_o,(\widehat{\varphi }\pm \mathrm{\Delta }\widehat{\theta })^2]`$, and gives the previously found asymptotic expression of $`\mathrm{sinc}[(\widehat{\varphi }\pm \mathrm{\Delta }\widehat{\theta })^2/4]`$ in the limit for $`\widehat{z}_o1`$. Since here $`\widehat{z}_o`$ is supposed to be at least of order unity ($`\widehat{z}_o>1/2`$), we can conclude that $`S`$ must be different from zero only for values of $`(\widehat{\varphi }\pm \mathrm{\Delta }\widehat{\theta })^2`$ of order unity (to be more precise, for values $`(\widehat{\varphi }\pm \mathrm{\Delta }\widehat{\theta })^24`$, $`(\widehat{\varphi }\pm \mathrm{\Delta }\widehat{\theta })^2/4`$ being the arguments of the sinc function) as it can be seen, for instance, from Fig. 10 for a particular case. Thus Eq. (89) and its asymptotic equivalent for $`\widehat{z}_o1`$ share the same mathematical structure.
Let us now introduce the non-restrictive assumptions:
$`\widehat{N}1,\widehat{D}1`$ (90)
and define
$$\widehat{A}=\frac{\widehat{N}}{\widehat{z}_o^2}.$$
(91)
The physical interpretation of $`\widehat{A}`$ follows from that of $`\widehat{\sigma }/\widehat{z}_o`$: $`\widehat{A}`$ is the dimensionless square of the apparent angular size of the source at the observer point position, calculated as if the source was positioned at $`\widehat{z}_o=0`$. If $`\widehat{N}1`$ and $`\widehat{D}1`$ we have $`(2\widehat{A}\widehat{z}_o^2\widehat{D})/(\widehat{A}+\widehat{D})1`$ for any value of $`\widehat{z}_o`$ and any choice of $`\widehat{N}`$ and $`\widehat{D}`$. As a result, from the exponential factor $`\mathrm{exp}[2\widehat{A}\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }^2\widehat{D}/(\widehat{A}+\widehat{D})]`$ outside the integral sign in Eq. (89) we have that $`G_\omega (\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })`$ is different from zero only for $`\mathrm{\Delta }\widehat{\theta }1`$. Then we can neglect terms in $`\mathrm{\Delta }\widehat{\theta }`$ in the factors $`S()`$ within the integral sign thus getting
$`\widehat{G}(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })={\displaystyle \frac{\mathrm{exp}\left(i2\overline{\theta }\widehat{z}_o\mathrm{\Delta }\widehat{\theta }\right)}{\sqrt{2\pi (\widehat{A}+\widehat{D})}}}\mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }^2+4\widehat{A}\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }^2\widehat{D}+4i\widehat{A}\overline{\theta }\widehat{z}_o\mathrm{\Delta }\widehat{\theta }}{2(\widehat{A}+\widehat{D})}}\right]`$ (92)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }\mathrm{exp}[{\displaystyle \frac{\widehat{\varphi }^2+2\widehat{\varphi }\overline{\theta }}{2(\widehat{A}+\widehat{D})}}]\mathrm{exp}[i{\displaystyle \frac{2\widehat{\varphi }\widehat{A}\widehat{z}_o\mathrm{\Delta }\widehat{\theta }}{\widehat{A}+\widehat{D}}}]|S[\widehat{z}_o,\widehat{\varphi }^2]|^2.`$ (93)
The maximal value of $`\overline{\theta }`$ is related with the width of the exponential function $`\mathrm{exp}[\overline{\theta }^2`$ $`/(2\widehat{A}`$ $`+2\widehat{D})]`$ outside the integral sign in Eq. (93). It follows that in the limit for $`\widehat{D}1`$ we can neglect the exponential factor $`\mathrm{exp}[(\widehat{\varphi }^2+2\widehat{\varphi }\overline{\theta })/(2\widehat{A}+2\widehat{D})]`$ within the integral sign: in fact, its argument assumes values of order unity for $`\widehat{\varphi }1`$, but the factor $`S[\widehat{z}_o,\widehat{\varphi }^2]^2`$ cuts off the integrand for $`\widehat{\varphi }1`$. Therefore Eq. (93) can be simplified as follows:
$`\widehat{G}(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })={\displaystyle \frac{\mathrm{exp}\left(i2\overline{\theta }\widehat{z}_o\mathrm{\Delta }\widehat{\theta }\right)}{\sqrt{2\pi (\widehat{A}+\widehat{D})}}}\mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }^2}{2(\widehat{A}+\widehat{D})}}\right]\mathrm{exp}\left[{\displaystyle \frac{2i\widehat{A}\overline{\theta }\widehat{z}_o\mathrm{\Delta }\widehat{\theta }}{\widehat{A}+\widehat{D}}}\right]`$ (94)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}\widehat{D}\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }^2}{\widehat{A}+\widehat{D}}}\right]{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }\mathrm{exp}\left[i\left({\displaystyle \frac{2\widehat{A}}{\widehat{A}+\widehat{D}}}\widehat{z}_o\mathrm{\Delta }\widehat{\theta }\right)\widehat{\varphi }\right]\left|S[\widehat{z}_o,\widehat{\varphi }^2]\right|^2.`$ (95)
The integral in Eq. (95) is simply the Fourier transform of the function $`f(\widehat{\varphi })=S[\widehat{z}_o,\widehat{\varphi }^2]^2`$ with respect to the variable $`2\widehat{A}\widehat{z}_o\mathrm{\Delta }\widehat{\theta }/(\widehat{A}+\widehat{D})`$. Since the function $`f(\widehat{\varphi })`$ has values sensibly different from zero only as $`\widehat{\varphi }`$ is of order unity or smaller, its Fourier Transform will also be suppressed for values of $`2\widehat{A}\widehat{z}_o|\mathrm{\Delta }\widehat{\theta }|/(\widehat{A}+\widehat{D})`$ larger than unity, by virtue of the Bandwidth Theorem. This means that the integral in Eq. (95) gives non-negligible contributions only up to some maximal value of $`|\mathrm{\Delta }\widehat{\theta }|`$:
$$|\mathrm{\Delta }\widehat{\theta }|_{\mathrm{max}}\frac{1}{2\widehat{z}_o}\left(1+\frac{\widehat{D}}{\widehat{A}}\right).$$
(96)
On the other hand, the exponential factor outside the integral in Eq. (95) will cut off the function $`\widehat{G}`$ around some other value
$$|\mathrm{\Delta }\widehat{\theta }|_{\mathrm{max2}}\frac{1}{2\widehat{z}_o}\left(\frac{1}{\widehat{D}}+\frac{1}{\widehat{A}}\right)^{1/2}.$$
(97)
It is easy to see that, for any value of $`\widehat{z}_o`$, $`|\mathrm{\Delta }\widehat{\theta }|_{\mathrm{max}}|\mathrm{\Delta }\widehat{\theta }|_{\mathrm{max2}}`$. In fact we have
$$\frac{|\mathrm{\Delta }\widehat{\theta }|_{\mathrm{max}}}{|\mathrm{\Delta }\widehat{\theta }|_{\mathrm{max2}}}\sqrt{\widehat{D}}\sqrt{1+(\widehat{D}/\widehat{A})}>\sqrt{\widehat{D}}1,$$
(98)
in the limit for $`\widehat{D}1`$. As a result the Fourier transform in Eq. (95) is significant only for values of the variable $`2\widehat{A}\widehat{z}_o\mathrm{\Delta }\widehat{\theta }/(\widehat{A}+\widehat{D})`$ near to zero and contributes to $`\widehat{G}`$ only by the inessential factor
$$_{\mathrm{}}^{\mathrm{}}𝑑\widehat{\varphi }\left|S[\widehat{z_o},\varphi ^2]\right|^2=\mathrm{constant}.$$
(99)
In order to use the correlation function $`\widehat{G}`$ for calculation of coherence length and other statistical properties, one has to use the spectral degree of coherence $`g`$, which can be presented as a function of $`\overline{\theta }`$ and $`\mathrm{\Delta }\widehat{\theta }`$ instead of $`x_{o2}`$ and $`x_{o1}`$:
$$g(\overline{\theta },\mathrm{\Delta }\widehat{\theta })=\frac{\widehat{G}(\overline{\theta },\mathrm{\Delta }\widehat{\theta })}{\left|\widehat{E}_s\left(\overline{\theta }+\mathrm{\Delta }\widehat{\theta }\right)\right|^2^{1/2}\left|E\left(\overline{\theta }\mathrm{\Delta }\widehat{\theta }\right)\right|^2^{1/2}}.$$
(100)
From Eq. (95) we obtain :
$`g(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })=\mathrm{exp}\left(i2\overline{\theta }\widehat{z}_o\mathrm{\Delta }\widehat{\theta }\right)\mathrm{exp}\left[{\displaystyle \frac{2i\widehat{A}\overline{\theta }\widehat{z}_o\mathrm{\Delta }\widehat{\theta }}{\widehat{A}+\widehat{D}}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}\widehat{D}\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }^2}{\widehat{A}+\widehat{D}}}\right].`$ (101)
In the asymptotic limit for a large value of $`\widehat{z}_o`$, $`\widehat{A}1`$, Eq. (101) can be simplified to
$`g(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })=\mathrm{exp}\left(i2\overline{\theta }\widehat{z}_o\mathrm{\Delta }\widehat{\theta }\right)\mathrm{exp}\left[{\displaystyle \frac{2i\widehat{A}\overline{\theta }\widehat{z}_o\mathrm{\Delta }\widehat{\theta }}{\widehat{D}}}\right]\mathrm{exp}\left[2\widehat{A}\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }^2\right].`$ (102)
From Eq. (95) it is easy to see that the region of interest where the field intensity is not negligible is when $`\overline{\theta }\sqrt{\widehat{D}}`$. Therefore, Eq. (102) can be further approximated to
$`g(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })=\mathrm{exp}\left(i2\overline{\theta }\widehat{z}_o\mathrm{\Delta }\widehat{\theta }\right)\mathrm{exp}\left[2\widehat{A}\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }^2\right].`$ (103)
It is interesting to calculate the transverse coherence length $`\widehat{\xi }_c`$ as a function of the observation distance $`\widehat{z}_o`$. For any experiment, complete information on the coherence properties of light are given by the function $`g`$. When calculating the coherence length one applies a certain algorithm to $`g`$ thus extracting a single number. This number does not include all information about the coherence properties of light, and the algorithm applied to $`g`$ is simply a convenient definition. Then, in order to calculate a coherence length one has, first, to choose a definition among all the possible convenient ones. In this paper we will simply follow the approach by Mandel, originally developed for the time domain, but trivially extensible to any domain of interest, in our case the angular domain. The coherence length, naturally normalized to the diffraction length $`\sqrt{L_wc/\omega }`$ is defined as
$$\widehat{\xi }_c(\widehat{z}_o)=2_{\mathrm{}}^{\mathrm{}}|g(\mathrm{\Delta }\widehat{\theta })|^2d(\widehat{z}_o\mathrm{\Delta }\widehat{\theta }),$$
(104)
where the factor $`2`$ in front of the integral on the right hand side is due to the fact that we chose Mandel’s approach and that our definition of $`\mathrm{\Delta }\widehat{\theta }`$ differs of a factor $`1/2`$ from his definition. Performing the integration in Eq. (104) with the help of Eq. (101) yields:
$$\widehat{\xi }_c(\widehat{z}_o)=\sqrt{\pi }\left(\frac{1}{\widehat{A}}+\frac{1}{\widehat{D}}\right)^{1/2}.$$
(105)
### 4.2 Discussion
The coherence length in Eq. (105) exhibits linear dependence on $`\widehat{z}_o`$, that is $`\widehat{\xi }_c\sqrt{\pi /\widehat{N}}\widehat{z}_o`$ while for $`\widehat{z}_o1/2`$ that is at the end of the undulator, it converges to a constant $`\widehat{\xi }_c[\pi /(4\widehat{N})+\pi /\widehat{D}]^{1/2}`$. Eq. (105) and its asymptotes are presented in Fig. 11 and Fig. 12 for the case $`\widehat{N}=10^3`$, $`\widehat{D}=10`$. It is evident that at the exit of the undulator, $`\widehat{\xi }_c1/\sqrt{\widehat{D}}`$, because $`\widehat{N}\widehat{D}`$. On the other hand, horizontal dimension of the light spot is simply proportional to $`\sqrt{\widehat{N}}`$ as it is evident from Eq. (101). This means that the horizontal dimension of the light spot is determined by the electron beam size, as is intuitive, while the beam angular distribution is printed in the fine structures of the intensity function, that are of the dimension of the coherence length. In the limit for $`\widehat{z}_o1`$ the situation is reversed. The radiation field at the source can be presented as a superposition of plane waves, all at the same frequency $`\omega _o`$, but with different propagation angles with respect to the $`z`$-direction. Since the radiation at the exit of the undulator is partially coherent, a spiky angular spectrum is to be expected. The nature of the spikes is easily described in terms of Fourier transform theory, in perfect analogy with what has been said about the frequency spectrum in Section 2.1. From Fourier transform theorem or, directly, from Eq. (95) or from geometrical optics arguments we can expect an angular spectrum envelope with Gaussian distribution and rms width of $`\sqrt{\widehat{D}}`$.
Also, the angular spectrum should contain spikes with characteristic width $`1/\sqrt{\widehat{N}}`$, as a consequence of the reciprocal width relations of Fourier transform pairs (see Fig. 13). This can be seen realized in mathematical form from the expression for the cross-spectral density, Eq. (101) and from the equation for the coherence length, Eq. (105). Since $`\widehat{N}1`$, the horizontal width of the coherence spot is much smaller than the vertical one.
It is also important to remark that the asymptotic behavior for $`\widehat{A}1`$ of $`g`$ in Eq. (103) and $`\widehat{\xi }_c`$
$$\widehat{\xi }_c\sqrt{\frac{\pi }{\widehat{N}}}\widehat{z}_o$$
(106)
are direct application of van Cittert-Zernike theorem. In fact, the last exponential factor on the right hand side of Eq. (102) is simply linked with the Fourier transform of $`F_{\widehat{l}_x}(\widehat{l}_x)`$. We derived Eq. (102) for $`\widehat{N}1`$ and $`\widehat{D}1`$, with $`\widehat{z}_o^2\widehat{D}\widehat{N}`$: in non-normalized units these conditions mean that the VCZ theorem is applicable when the electron beam divergence is much larger than the diffraction angle, i.e. $`\sigma _x^{}^2\lambda /(2\pi L_w)`$, the electron beam dimensions are much larger than the diffraction size, i.e. $`\sigma _x^2\lambda L_w/2\pi `$, and $`(\sigma _x^{}z_o)^2\sigma _x^2`$. On the contrary, authors of TAKA state that, in order for the van Cittert-Zernike theorem to be applicable, ”the electron-beam divergence must be much smaller than the photon divergence”, that is our diffraction angle, i.e. $`\sigma _x^{}\sqrt{\lambda /(2\pi L_w)}`$ (reference TAKA , page 571, Eq. (57)). Our derivation shows that this conclusion is incorrect.
In GOOD (paragraph 5.6.4) a rule of thumb is given for the applicability region of the generalization of the VCZ theorem to quasi-homogeneous sources. The rule of thumb requires $`z_o>2d\mathrm{\Delta }/\lambda `$ where $`d`$ is ”the maximum linear dimension of the source”, that is the diameter of a source with uniform intensity and $`\mathrm{\Delta }`$ ”represents the maximum linear dimension of a coherence area of the source”. In our case $`d2\sigma _x`$, since $`\sigma _x`$ is the rms source dimension, and from Eq. (105) we have $`\mathrm{\Delta }=\xi _c\lambda /(2\sqrt{\pi }\sigma _x^{})`$. The rule of thumb then requires $`z_o>2\sigma _x/(\sqrt{\pi }\sigma _x^{})`$: in dimensionless this reads $`\widehat{z}_o\sqrt{\widehat{N}/\widehat{D}}`$. This is parametrically in agreement with our limiting condition $`\widehat{z}_o^2\widehat{D}\widehat{N}`$, even though these two conditions are obviously different when it come to actual estimations: our condition is, in fact, only an asymptotic one. To see how well condition $`\widehat{z}_o\sqrt{\widehat{N}/\widehat{D}}`$ works in reality we might consider the plot in Fig. 11. There $`\widehat{N}=10^3`$ and $`\widehat{D}=10`$ so that, following GOOD we may conclude that a good condition for the applicability of the VCZ theorem should be $`\widehat{z}_o10`$. However as it is seen from the figure, the linear asymptotic behavior is not yet a good approximation at $`\widehat{z}_o10`$. This may be ascribed to the fact that the derivation in GOOD is not generally valid, but has been carried out for sources which drop to zero very rapidly outside the maximum linear dimension $`d`$ and whose correlation function also drops rapidly to zero very rapidly outside maximum linear dimension $`\mathrm{\Delta }`$.
However, at least parametrically, the applicability of the VCZ theorem in the asymptotic limit $`\widehat{z}_o^2\widehat{D}\widehat{N}`$ can be also expected from the condition $`z_o>2d\mathrm{\Delta }/\lambda `$ in GOOD . In other words, with the help of our approach we were able to specify an asymptotic region where the VCZ theorem holds. Such a region overlaps with predictions from Statistical Optics. Statistical Optics can describe propagation of the cross-spectral density only once it is known at some source plane position. Our treatment allows us to specify the cross-spectral density at the exit of the undulator, but it should be noted that we do not need to use customary results of Statistical Optics and propagate the cross-spectral density from the exit of the undulator in order to obtain the cross-spectral density at some distance along the beamline. In fact our approach, which consists in taking advantage of the system Green’s function in paraxial approximation and, subsequently, of the resonance approximation, allows us to calculate the cross-spectral density directly at any distance from the exit of the undulator.
Let us now consider the structure of Eq. (101) and discuss the meaning of the phase terms in $`\overline{\theta }\mathrm{\Delta }\widehat{\theta }`$. These are important in relation with the condition for quasi-homogeneous source: their presence couples the two variables $`\overline{\theta }`$ and $`\mathrm{\Delta }\widehat{\theta }`$ and prevents the source to be quasi-homogeneous at any given value of $`\widehat{z}_o`$<sup>4</sup><sup>4</sup>4Note, again, that the definition of ”source plane” is just conventional. One may define the source plane at the exit of the undulator, that is at $`\widehat{z}_o=1/2`$, but there is no fundamental reason for such a definition: one may pick any value of $`\widehat{z}_o`$ as the source position., unless they compensate each other in some parameter region.
Let us discuss the limit, $`\widehat{A}1`$. We may consider two subcases. First, consider $`\widehat{A}\widehat{D}1`$. In this case, inspection of Eq. (101) shows that the two phase terms compensate and the source is quasi-homogeneous, because the cross-spectral density is factorized in a function of $`\overline{\theta }`$ and a function of $`\mathrm{\Delta }\widehat{\theta }`$. It should be noted that if condition $`\widehat{A}1`$ is not satisfied at the exit of the undulator, where $`\widehat{z}_o1`$, then it is never satisfied. If $`\widehat{N}\widehat{D}1`$ we have a quasi-homogeneous source at the exit of the undulator.
Second, consider $`\widehat{D}\widehat{N}1`$. This correspond to a situation with a low value of the normalized betatron function in the horizontal direction. Figure 14 shows a numerical example with $`\widehat{ϵ}_x=100`$ and $`\widehat{\beta }_x=0.3`$ that is $`\widehat{N}=30`$ and $`\widehat{D}=300`$: the value for the horizontal betatron function is similar to the low-$`\beta `$ case reported at page 12, Table 2.2.2 in PETR , where $`\beta _x=1.3`$ m for a $`5`$m-long insertion device. The value $`\widehat{ϵ}_x=100`$ corresponds to a wavelength of about $`0.6\AA `$ for the PETRA III case. When $`\widehat{D}\widehat{N}1`$ no compensation of the phase terms in Eq. (101) is possible, not even at the exit of the undulator. In this case, whatever the value of $`\widehat{z}_o`$ we can never have a quasi-homogeneous wavefront. This constitutes no problem. Simply, the wavefront is not-quasi-homogeneous in this case. However, we may interpret the situation by saying that a ”virtual” quasi-homogeneous source placed in the center of the undulator would result in the non-homogeneous source described by Eq. (101) at the exit of the undulator. Although it physically makes no sense to discuss about Eq. (101) inside the undulator, the ”virtual” source analogy is suggested by the fact that setting $`\widehat{z}_o=0`$ in Eq. (101), both phase terms become zero.
Note that in general, whatever the values of $`\widehat{N}1`$ and $`\widehat{D}1`$, one never has quasi-homogeneous sources in the limit for $`\widehat{A}1`$. In fact, in the asymptotic $`\widehat{A}1`$ only the phase factor $`\mathrm{exp}(i2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o)`$ contributes, which couples $`\overline{\theta }`$ and $`\mathrm{\Delta }\widehat{\theta }`$. Such a factor is connected with phase of the field from a single electron in an undulator in the far zone, $`\omega (x_o^2+y_o^2)/(2cz_o)`$, which represents, in paraxial approximation, the phase difference between the point $`(x_o,y_o,z_o)`$ and the point $`(0,0,z_o)`$: in the asymptotic for large values of $`\widehat{z}_o`$, the electric field generated by a single electron with offset and deflection in an undulator has a spherical wavefront (see OURS ). When one calculates the field correlation function at two different points, he ends up with a contribution equal to the difference (due to complex conjugation) between $`\omega (x_{o2}^2+y_{o2}^2)/(2cz_o)`$ and $`\omega (x_{o1}^2+y_{o1}^2)/(2cz_o)`$, that for the vertical (and separately, the horizontal) direction gives exactly the shift $`2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o`$ in normalized units. It should be noted that such a reasoning is not limited to Synchrotron Radiation sources, but it is quite general since, as already discussed, it relies on the fact that the wavefront of a single radiator (in our case, an electron) produces a spherical wavefront in the far field. This is, for instance, the case of thermal sources as well. In other words, if the far field radiation of a quasi-homogeneous source is taken as a new source, that new source will never be quasi-homogeneous.
A common property of all situations with $`\widehat{N}1`$ and $`\widehat{D}1`$ is that, for any value of $`\widehat{z}_o`$, the modulus of $`g`$, i.e. $`|g|`$, is always independent on $`\overline{\theta }`$. Moreover, it is always possible to apply the VCZ theorem starting either from a virtual quasi-homogeneous source placed at center of the undulator when $`\widehat{D}\widehat{N}1`$, or otherwise from a real one placed at the exit of the undulator (or at any other position $`\widehat{z}_o`$ close enough to the exit of the undulator to guarantee a quasi-homogeneous wavefront). These observations suggest to extend the concept of quasi-homogeneity, and introduce the new concept of ”weak quasi-homogeneity” as discussed before. With respect to the new coordinates $`\overline{\theta }`$ and $`\mathrm{\Delta }\widehat{\theta }`$, a given wavefront at fixed position $`\widehat{z}_o`$ is said to be weakly quasi-homogeneous when $`|g|`$ is independent of $`\overline{\theta }`$. With this new definition at hand we can restate some of our conclusions in a slightly different language. We have seen that in the far field, when the VCZ theorem holds, the wavefronts are weakly quasi-homogeneous, but never quasi-homogeneous in the usual sense. In the case $`\widehat{N}\widehat{D}1`$ we pass from quasi-homogeneous wavefronts (in the usual sense) in the near field to weakly-quasi homogeneous wavefronts (but not quasi-homogenous in the usual sense) in the far field. Note that the wavefronts are always weakly quasi-homogeneous, even during the transition from near to far zone. In the case $`\widehat{D}\widehat{N}1`$ instead, the VCZ theorem is applicable already from the exit of the undulator, as it can be seen from Fig. 14, and the wavefront is not quasi-homogeneous in the usual sense, but still weakly quasi-homogeneous from the very beginning.
The weak quasi-homogeneity of the wavefronts at any value of $`\widehat{z}_o`$, i.e. the fact that $`|g|`$ is independent of $`\overline{\theta }`$ for any value of $`\widehat{z}_o`$ guarantees that the plot in Fig. 14 is universal. It should be noted that this fact depends on the choice of large parameters $`\widehat{N}1`$ and $`\widehat{D}1`$, but it is also strictly related with the Gaussian nature of the electron distribution in angles and offsets, that is a well-established fact for storage-ring beams. If angles or offsets were obeying different distribution laws, in general, one could not perform the integral in $`\widehat{\eta }`$ in Eq. (82) and, in general, $`|g|`$ would have shown a dependence on $`\overline{\theta }`$: our noticeable result is linked with the properties of the exponential elementary function. However, it should be clear that even in the case when angles or offsets were obeying different distribution laws, i.e when the plot in Fig. 14 is not universal, we could have situations when wavefronts are quasi-homogeneous in the usual sense near the exit of the undulator and are weakly quasi-homogeneous in the far field limit, but not along the transition between these two zones. A more detailed discussion of this issue will be given in Section 5, where we will be discussing conditions for the source to be quasi-homogeneous.
Another remark to be made pertains the applicability of the VCZ theorem. As we deal with a quasi-homogenous source (in the usual sense) the knowledge of $`I(\overline{\theta })`$ and $`g(\mathrm{\Delta }\widehat{\theta })`$ in the far zone allow, respectively, the calculation at the source plane of $`g(\mathrm{\Delta }\widehat{x})`$ through the ”anti” VCZ theorem and of $`I(\overline{x})`$ through the VCZ theorem (here we consider only one dimension, the horizontal one $`x`$). Viceversa the knowledge of $`I`$ and $`g`$ at the source allow calculation of $`I`$ and $`g`$ in the far field. In terms of intensity, all information regarding wavefront evolution (assuming a quasi-homogeneous source, in the usual sense) is included in $`I(\overline{\theta })`$ and $`I(\overline{x})`$. For instance, the knowledge of $`I(\overline{\theta })`$ allows calculation of $`g(\mathrm{\Delta }\widehat{x})`$ at the source plane through the ”anti” VCZ theorem. Then, the knowledge of $`g(\mathrm{\Delta }\widehat{x})`$ and $`I(\overline{x})`$ allow the calculation of the cross-spectral density, which can be propagated at any distance, and allow to recover $`g(\mathrm{\Delta }\widehat{\theta })`$. So, complete characterization of the undulator source is given when $`I(\overline{\theta })`$ and $`I(\overline{x})`$, when the source is assumed quasi-homogenous.
Yet we have seen that, when the electron distribution in angles and offsets are Gaussian and $`\widehat{N}1`$ and $`\widehat{D}1`$, the VCZ theorem holds also in the case $`\widehat{D}\widehat{N}1`$, when the source is non quasi-homogeneous in the usual sense. We have seen that the situation can be equivalently described with the help of a virtual quasi-homogeneous source in the middle of the undulator. However, such an interpretation is only valid a posteriori. For the case $`\widehat{N}1`$ and $`\widehat{D}1`$ there is a non quasi-homogeneous wavefront at the undulator exit; before our approach was presented one would have concluded that the VCZ theorem cannot be applied, since the spectral degree of coherence does not form a Fourier pair with the intensity distribution at the undulator exit. Our approach is based on the simplification of mathematical results through the use of small and large parameters and subsequent understanding and interpretation of these simplified results: in our analysis we were never limited to the treatment of quasi-homogenous cases alone.
As a closing remark about the coherence length we like to draw the reader’s attention on the fact that the dimensional form $`\xi _c`$ of the coherence length, given in normalized units by $`\widehat{\xi }_c`$ in Eq. (105), does not include the undulator length. This is to be expected since, in the limit $`\widehat{N}1`$ and $`\widehat{D}1`$, the typical size and divergence of the electron beam are much larger than the diffraction size $`\sqrt{cL_w/\omega }`$ and angle $`\sqrt{c/(\omega L_w)}`$, which are intrinsic properties of the undulator radiation. As a result, in this limit, the evolution of the radiation beam is a function of the electron beam parameters only, and does not depend on the undulator length. In the following Section 5, where we will extend our model to a two-dimensional case, we will see that the quasi-homogeneous approximation is valid in many practical situations, but we will have to account for diffraction of undulator radiation in the vertical direction. In this case the dimensional coherence $`\xi _c`$ length will be a function of the undulator length $`L_w`$ as well.
## 5 Effect of the vertical emittance on the cross-spectral density
Up to now we were dealing with the field correlation function $`g`$ within the framework of a one-dimensional model.
In fact we considered the limit for $`\widehat{ϵ}_y\widehat{\beta }_y1`$ and $`\widehat{ϵ}_y/\widehat{\beta }_y1`$ and we calculated $`g`$ for $`\widehat{\theta }_{y1}=\widehat{\theta }_{y2}=0`$ and $`\widehat{C}=0`$, so that our attention was focused on coherent effects in the horizontal direction. We will now extend our considerations to a two-dimensional model always for $`\widehat{C}=0`$. This can be done by a straightforward generalization of Eq. (89) which can be obtained from Eq. (71) following the same steps which lead to Eq. (89), but this time without assumptions on $`\widehat{N}_y`$, $`\widehat{D}_y`$, $`\widehat{\theta }_{y1}`$ and $`\widehat{\theta }_{y2}`$. Finally, at the end of calculations, our final expression for $`\widehat{G}`$ should be normalized to
$$\widehat{W}=\left|\widehat{E}_s\left(\stackrel{}{\widehat{\theta }}_1\right)\right|^2^{1/2}\left|\widehat{E}_s\left(\stackrel{}{\widehat{\theta }}_2\right)\right|^2^{1/2}.$$
(107)
As has already been seen in Section 4, after normalization to $`\widehat{W}`$ we will obtain the spectral degree of coherence $`g`$. With this in mind we will neglect, step after step, unnecessary multiplicative factors that, in any case, would be finally disposed after normalization of the final result. Retaining indexes $`x`$ and $`y`$ in our notation we obtain
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (108)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (109)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_y^2+4\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+4i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{2(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (110)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{\widehat{\varphi }_x^2+2\widehat{\varphi }_x\left(\overline{\theta }_x+2i\widehat{A}_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x\right)}{2(\widehat{A}_x+\widehat{D}_x)}}]`$ (111)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\widehat{\varphi }_y^2+2\widehat{\varphi }_y\left(\overline{\theta }_y+2i\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)}{2(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (112)
$`\times S^{}[\widehat{z}_o,(\widehat{\varphi }_x\mathrm{\Delta }\widehat{\theta }_x)^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,(\widehat{\varphi }_x+\mathrm{\Delta }\widehat{\theta }_x)^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2].`$ (113)
We will still assume $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$. This allows to factorize the right hand side of Eq. (113) in the product of contribution depending on horizontal ($`\overline{\theta }_x,\mathrm{\Delta }\widehat{\theta }_x`$) coordinates only with a second depending on vertical ($`\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_y`$) coordinates only. In fact, from the exponential factor outside the integral sign in Eq. (113) it is possible to see that the maximum value of $`\mathrm{\Delta }\widehat{\theta }_x^2`$ is of order $`(\widehat{A}_x+\widehat{D}_x)/(\widehat{A}_x\widehat{D}_x\widehat{z}_o^2)1`$. As a result, $`\mathrm{\Delta }\widehat{\theta }_x`$ can be neglected inside the $`S`$ functions in Eq. (113). Moreover, since $`\widehat{D}_x1`$ one can also neglect the exponential factor in $`\widehat{\varphi }_x^2+2\widehat{\varphi }_x\overline{\theta }_x`$ inside the integral sign. This leads to
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (114)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (115)
$`\mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_y^2+4\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+4i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{2(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (116)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x\mathrm{exp}\left[i\widehat{\varphi }_x{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{\widehat{A}_x+\widehat{D}_x}}\right]{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{\widehat{\varphi }_y^2+2\widehat{\varphi }_y\left(\overline{\theta }_y+2i\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)}{2(\widehat{A}_y+\widehat{D}_y)}}]`$ (117)
$`\times S^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2].`$ (118)
Based on the same reasoning in Section 4.1, that we repeat here for completeness, we can also neglect the phase factor in $`\widehat{\varphi }_x`$ under the integral in $`d\widehat{\varphi }_x`$ in Eq. (118). Such integral in $`d\widehat{\varphi }_x`$ is simply the Fourier transform of the function
$`f(\widehat{\varphi }_x)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[{\displaystyle \frac{\widehat{\varphi }_y^2+2\widehat{\varphi }_y\left(\overline{\theta }_y+2i\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)}{2(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (119)
$`\times S^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2].`$ (120)
with respect to the variable $`2\widehat{A}_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x/(\widehat{A}_x+\widehat{D}_x)`$. In the argument of the $`S`$ functions on the right hand side of Eq. (120), $`\widehat{\varphi }_x^2`$ is always summed to positively defined quantities. This remark allows one to conclude that $`f(\widehat{\varphi }_x)`$ has values sensibly different from zero only as $`\widehat{\varphi }_x`$ is of order unity or smaller. Therefore, its Fourier Transform will also be suppressed for values of $`2\widehat{A}_x\widehat{z}_o|\mathrm{\Delta }\widehat{\theta }_x|/(\widehat{A}_x+\widehat{D}_x)`$ larger than unity, by virtue of the Bandwidth Theorem. This means that the integral in $`d\widehat{\varphi }_x`$ in Eq. (118) gives non-negligible contributions only up to some maximal value of $`|\mathrm{\Delta }\widehat{\theta }_x|`$:
$$|\mathrm{\Delta }\widehat{\theta }|_{x\mathrm{max}}\frac{1}{2\widehat{z}_o}\left(1+\frac{\widehat{D}_x}{\widehat{A}_x}\right).$$
(121)
On the other hand, the exponential factor outside the integral in Eq. (118) will cut off the function $`\widehat{G}`$ around some other value
$$|\mathrm{\Delta }\widehat{\theta }|_{x\mathrm{max2}}\frac{1}{2\widehat{z}_o}\left(\frac{1}{\widehat{D}_x}+\frac{1}{\widehat{A}_x}\right)^{1/2}.$$
(122)
It is easy to see that, for any value of $`\widehat{z}_o`$, $`|\mathrm{\Delta }\widehat{\theta }|_{x\mathrm{max}}|\mathrm{\Delta }\widehat{\theta }|_{x\mathrm{max2}}`$. In fact we have
$$\frac{|\mathrm{\Delta }\widehat{\theta }|_{x\mathrm{max}}}{|\mathrm{\Delta }\widehat{\theta }|_{x\mathrm{max2}}}\sqrt{\widehat{D}_x}\sqrt{1+(\widehat{D}_x/\widehat{A}_x)}>\sqrt{\widehat{D}_x}1,$$
(123)
in the limit for $`\widehat{D}_x1`$. As a result the Fourier transform in Eq. (118) is significant only for values of the variable $`2\widehat{A}_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x/(\widehat{A}_x+\widehat{D}_x)`$ near to zero. As a result we obtain the following equation for $`\widehat{G}`$:
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (124)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (125)
$`\mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_y^2+4\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+4i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{2(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (126)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{\widehat{\varphi }_y^2+2\widehat{\varphi }_y\left(\overline{\theta }_y+2i\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)}{2(\widehat{A}_y+\widehat{D}_y)}}]`$ (127)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_xS^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2],`$ (128)
where horizontal and vertical coordinates are obviously factorized.
Eq. (128) has been derived assuming $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$. Note that assuming setting $`\mathrm{\Delta }\widehat{\theta }_y=\overline{\theta }_y=0`$ one can obtain Eq. (101) from Eq. (128). This proves that Eq. (101) has a wider range of validity than that for $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$ (as the reader will remember, these assumptions were made just for notational simplicity). In fact, as $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$ horizontal and vertical direction factorize and the horizontal factor is always that in Eq. (101), independently on $`\widehat{N}_y`$ and $`\widehat{D}_y`$.
Under one of the two extra assumptions $`\widehat{N}_y1`$ or $`\widehat{D}_y1`$, Eq. (128) can be further simplified and often describes a weakly quasi-homogeneous wavefront according to the definition given at the end of Section 4.1. At the end of the Section we will show that as $`\widehat{D}_y1`$ we always have a weakly quasi-homogeneous wavefront, for any value of $`\widehat{N}_y`$. It will also be seen that the same applies when $`\widehat{N}_y1`$ and $`\widehat{A}_y1`$ (far field) or $`\widehat{A}_y1`$ (near field) for any value of $`\widehat{D}_y`$. However, as $`\widehat{N}_y1`$, $`\widehat{D}_y1`$ and $`\widehat{A}_y1`$ we have an intermediate region between the near and far region were, in general, wavefronts are not quasi-homogeneous, not even in the weak case.
For simplicity of discussion we will set $`\overline{\theta }_x=\overline{\theta }_y=0`$ thus obtaining from Eq. (128)
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{(\widehat{A}_x+\widehat{D}_x)}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y}{(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (129)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{\widehat{\varphi }_y^2+2\widehat{\varphi }_y\left(2i\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)}{2(\widehat{A}_y+\widehat{D}_y)}}]`$ (130)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_xS^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2],`$ (131)
which is easier to manipulate. It should be reminded that, if one is interested in ascertaining the weak quasi-homogeneity of a wavefront, one has to deal with the full Eq. (128). Moreover, it should be noted that on the one hand, within the weak quasi-homogeneous case, Eq. (131) is quite general and we can extract from it all important information on the transverse coherence independently on the values of $`\overline{\theta }_x`$ and $`\overline{\theta }_y`$ . On the other hand though, in the case the weakly quasi-homogeneous assumption fails, Eq. (131), e.g. when $`\widehat{N}_y1`$, $`\widehat{D}_y1`$ and $`\widehat{A}_y1`$ as we will see, $`|g|`$ depends on $`\overline{\theta }_y`$ and the study of Eq. (131) has a more restricted range of validity, namely for the particular value of $`\overline{\theta }_y=0`$.
In Section 5.1 we will assume $`\widehat{N}_y1`$ and arbitrary $`\widehat{D}_y`$, while in Section 5.2 we will study the case with arbitrary $`\widehat{N}_y`$ and $`\widehat{D}_y1`$.
In general, the coherence length in the $`\widehat{y}`$ direction (calculated at $`\mathrm{\Delta }\widehat{\theta }_x=0`$, but trivially extendible to the case $`\mathrm{\Delta }\widehat{\theta }_x0`$), $`\widehat{\xi }_{cy}`$ is a function of $`\widehat{D}_y`$, $`\widehat{N}_y`$ and $`\widehat{z}_o`$, as it can be concluded by inspection of the general expression for $`\widehat{G}`$ in Eq. (113).
As $`\widehat{N}_y1`$ and $`\widehat{D}_y`$ is arbitrary we will demonstrate that $`\widehat{\xi }_{cy}`$ can be approximated as
$$\widehat{\xi }_{cy}=\mathrm{\Phi }[\widehat{D}_y,\widehat{A}_y],$$
(132)
where $`\mathrm{\Phi }`$ is a universal function of the dimensionless parameters $`\widehat{D}_y`$ and $`\widehat{A}_y`$. We will first study the asymptotic cases $`\widehat{D}_y1`$ and $`\widehat{D}_y1`$, which are useful for physical understanding. Then we will generalize our results accounting for the influence of a finite divergence of the electron beam on the cross-spectral density. An analytical approximation for the function $`\mathrm{\Phi }`$ will be proposed. This will be chosen to match in a very simple way both asymptotic expressions for $`\widehat{D}_y1`$ and $`\widehat{D}_y1`$: we will demonstrate that for any value of $`\widehat{D}_y`$, discrepancies between the approximated and the actual (numerically calculated) expression are less than $`10\%`$, though there is no theoretical reason to assume, a priori, this relatively good accuracy. Also, at $`\widehat{A}_y1`$ (and $`\widehat{N}_y1`$) we will see that the VCZ theorem always hold.
As $`\widehat{N}_y`$ is arbitrary and $`\widehat{D}_y1`$ we will show, instead, that the coherence length (calculated again at $`\mathrm{\Delta }\widehat{\theta }_x=0`$, but trivially extendible to the case $`\mathrm{\Delta }\widehat{\theta }_x0`$) can be approximated as
$$\widehat{\xi }_{cy}=\mathrm{\Psi }[\widehat{D}_y,\widehat{z}_o^2/\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)],$$
(133)
$`\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)`$ being a universal function of the Fresnel number $`\widehat{N}_y`$. $`\mathrm{\Psi }`$ is a universal function of the dimensionless parameters $`\widehat{D}_y`$, $`\widehat{N}_y`$ and $`\widehat{z}_o`$. As usual in this paper, we will first study the asymptotic cases $`\widehat{N}_y1`$ and $`\widehat{N}_y1`$. Then we will generalize our results accounting for any value of $`\widehat{N}_y`$. As before, an analytical approximation for $`\mathrm{\Psi }`$ will be proposed. This will be chosen to match in a very simple way both asymptotic expressions for $`\widehat{N}_y1`$ and $`\widehat{N}_y1`$: again, we will demonstrate that for any value of $`\widehat{N}_y`$, discrepancies between the approximated and the actual (numerically calculated) expression are less than $`13\%`$ though there is no theoretical reason to assume, a priori, this relatively good accuracy.
The case for a large Fresnel number with a finite electron beam divergence, or viceversa of a large beam divergence and a finite Fresnel number, is of practical importance. To give a numerical example, let us put $`\lambda =1\AA `$, and consider a typical vertical emittance (for third generation sources in operation) $`ϵ_y310^{11}`$ m, that is $`\widehat{ϵ}_y2`$. On the one hand, if $`\widehat{\beta }_y=3`$, we have $`\widehat{D}_y0.6`$ and $`\widehat{N}_y6`$. On the other hand, if $`\widehat{\beta }_y=0.3`$ we have, viceversa, $`\widehat{D}_y6`$ and $`\widehat{N}_y0.6`$.
### 5.1 Very large Fresnel number $`\widehat{N}_y1`$, arbitrary divergence parameter $`\widehat{D}_y`$
As a matter of fact, the only important assumptions used to derive results Section 4.1 were that $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$ that are valid here as well. When these assumptions are granted, results can be factorized as a product of factors dependent separately on the $`x`$ and on the $`y`$ directions. Then the spectral degree of coherence in the horizontal direction will always be the same as in Section 4.1. Differences will arise here, of course, due to $`\widehat{\theta }_{y1}\widehat{\theta }_{y2}`$.
#### 5.1.1 Case with divergence parameter $`\widehat{D}_y1`$.
This case is the easiest to analyze, because one can follow step by step the derivation for the one dimensional model given in the previous Section 4.1. Calculations in the vertical direction $`y`$ simply follow the derivation for the horizontal direction $`x`$. As a result, one can start from Eq. (113) and perform, separately for the $`x`$ and the $`y`$ directions, the same simplifications which hold for the one-dimensional model. By comparison with Eq. (101), one obtains directly the result
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{D}_x}{\widehat{A}_x+\widehat{D}_x}}\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{D}_y}{\widehat{A}_y+\widehat{D}_y}}\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right].`$ (134)
Normalization of Eq. (134) according to Eq.(107) has been obtained simply setting $`\widehat{G}(\widehat{z}_o,0,0)=1`$. In this case as well, two-dimensional Fourier Transform of $`|S[\widehat{z}_o,\widehat{\varphi }_x^2+\widehat{\varphi }_y^2]|^2`$ calculated with respect the variables $`2\widehat{A}_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x`$ $`/(\widehat{A}_x+\widehat{D}_x)`$ and $`2\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y`$ $`/(\widehat{A}_y+\widehat{D}_y)`$ gives, in analogy with Eq. (99), an unessential factor
$$_{\mathrm{}}^{\mathrm{}}𝑑\widehat{\varphi }_x_{\mathrm{}}^{\mathrm{}}𝑑\widehat{\varphi }_y\left|S[\widehat{z}_o,\widehat{\varphi }_x^2+\widehat{\varphi }_y^2]\right|^2=\mathrm{constant},$$
(135)
which has been included in the normalization. Again, similarly as before one can calculate the coherence area $`\widehat{\mathrm{\Omega }}_c(\widehat{z}_o)`$ defined in analogy with $`\widehat{\xi }_c(\widehat{z}_o)`$ as
$$\widehat{\mathrm{\Omega }}_c(\widehat{z}_o)=\widehat{\xi }_{cx}(\widehat{z}_o)\widehat{\xi }_{cy}(\widehat{z}_o)$$
(136)
Performing the integration yields:
$$\widehat{\mathrm{\Omega }}_c=\pi \left(\frac{1}{\widehat{A}_x}+\frac{1}{\widehat{D}_x}\right)^{1/2}\left(\frac{1}{\widehat{A}_y}+\frac{1}{\widehat{D}_y}\right)^{1/2}$$
(137)
In the limit for a large value of $`\widehat{z}_o`$ the coherence area exhibits quadratic dependence on $`\widehat{z}_o`$, that is $`\widehat{\mathrm{\Omega }}_c\pi \widehat{z}_o^2/(\widehat{N}_x\widehat{N}_y)^{1/2}`$ while for $`\widehat{z}_o1/2`$, that is at the end of the undulator, it converges to the constant value $`\widehat{\mathrm{\Omega }}_c`$ $`\pi \left[1/(4\widehat{N}_x)+1/\widehat{D}_x\right]^{1/2}\left[1/(4\widehat{N}_y)+1/\widehat{D}_y\right]^{1/2}`$.
It should be noted as before that the asymptotic behavior for $`\widehat{z}_o1`$ of $`g`$
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\right]\mathrm{exp}\left[2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right]`$ (138)
and $`\widehat{\mathrm{\Omega }}_c`$
$$\widehat{\mathrm{\Omega }}_c=\frac{\pi \widehat{z}_o^2}{\sqrt{\widehat{N}_x\widehat{N}_y}}$$
(139)
are direct application of van Cittert-Zernike theorem, as it must be since $`\widehat{A}_{x,y}1`$. In fact, Eq. (138) is simply linked with the two-dimensional Fourier transform of $`F_{\widehat{l}_x}(\widehat{l}_x)F_{\widehat{l}_y}(\widehat{l}_y)`$.
#### 5.1.2 Case with divergence parameter $`\widehat{D}_y1`$.
With respect to the situation treated in Paragraph 5.1.1, this case requires a more careful analysis of the relations between different parameters. In fact, on the one hand $`\widehat{D}_y1`$ implies that the electron beam divergence drops out of the problem parameters, but on the other hand in this case the divergence of the radiation is described by the intrinsic divergence of undulator radiation, that is described by a more complicate mathematical function, compared with a Gaussian. In relation with this, it should be noted that simplifications in Section 4.1 pr 5.1.1 were based on the very specific properties of the Gaussian function representing the electron distributions in offset and deflection. Luckily, this is a realistic description in storage ring beam physics.
Let us consider Eq.(131). In order to derive it we only used the assumptions $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$, for $`\overline{\theta }_x=\overline{\theta }_y=0`$. The result of operations on the right hand side of Eq. (131) depend on how $`\widehat{N}_y`$ scales with respect to $`\widehat{z}_o^2`$. The cases for $`\widehat{A}_y>1`$ cannot be dealt with fully analytically. In the following we will analyze the asymptotic situation $`\widehat{A}_y1`$ and then we will treat semi-analytically the generic situation for all values of $`\widehat{A}_y`$. As we will see, as soon as $`\widehat{A}_y<1`$ we start to be in the applicability region of the VCZ theorem.
It should be noted that the dependence in $`\mathrm{\Delta }\widehat{\theta }_x`$ and $`\mathrm{\Delta }\widehat{\theta }_y`$ in Eq. (131) are already separated. Therefore, what has been said in Section 4.1 regarding the behavior of coherence properties in the horizontal direction hold independently of the behavior of coherence properties in the vertical direction.
(A) Far zone case: $`\widehat{A}_y1`$. — Since $`\widehat{N}_y1`$ it must be $`\widehat{z}_o^2\widehat{N}_y1`$ in order to allow for $`\widehat{A}_y1`$. Since we are working in quasi-homogeneous source condition ($`\widehat{N}_{x,y}1`$) in the limiting situation $`\widehat{A}_y1`$ we should recover VCZ theorem: this is the far field case. Even in the presence of the extra parameter $`\widehat{D}_y`$ we can treat the generic case $`\widehat{A}_y1`$ independently on how $`\widehat{A}_y`$ compares with respect to $`\widehat{D}_y`$.
Eq. (131) can be simplified on the assumptions $`\widehat{A}_y1`$ and $`\widehat{D}_y1`$. In fact, the Gaussian exponential factor inside the integral in $`d\widehat{\varphi }_y`$ in Eq. (131) imposes a maximal value $`\widehat{\varphi }_y^2\widehat{A}_y+\widehat{D}_y1`$.
Simultaneously, from the oscillating factor, always inside the integral in $`d\widehat{\varphi }_y`$, we have a condition for the maximal value of $`\mathrm{\Delta }\widehat{\theta }_y^21/\widehat{z}_o^21`$: in fact, if this condition is not fulfilled the integrand will be highly oscillatory. Alternatively, we can obtain a similar condition from the Gaussian exponential factor in $`\mathrm{\Delta }\widehat{\theta }_y`$ outside the integral, since $`\mathrm{\Delta }\widehat{\theta }_y^21/(\widehat{A}_y\widehat{z}_o^2)1`$.
As a result, the dependence of $`S`$ on $`(\widehat{\varphi }+\mathrm{\Delta }\widehat{\theta }_y)^2`$ can be dropped giving
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{(\widehat{A}_x+\widehat{D}_x)}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y}{(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (140)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x|S[\widehat{z}_o,\widehat{\varphi }_x^2]|^2`$ (141)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{\widehat{\varphi }_y^2+2\widehat{\varphi }_y\left(2i\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)}{2(\widehat{A}_y+\widehat{D}_y)}}].`$ (142)
The integral in $`d\widehat{\varphi }_y`$ can be performed giving
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{\widehat{A}_x+\widehat{D}_x}}\right]`$ (143)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y}{\widehat{A}_y+\widehat{D}_y}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y^2\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (144)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x|S[\widehat{z}_o,\widehat{\varphi }_x^2]|^2.`$ (145)
Normalizing $`\widehat{G}`$ in such a way that $`\widehat{G}(\widehat{z}_o,0,0)=1`$ we obtain the following expression for the spectral degree of coherence $`g`$:
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{\widehat{A}_x+\widehat{D}_x}}\right]`$ (146)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y}{\widehat{A}_y+\widehat{D}_y}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y^2\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2}{\widehat{A}_y+\widehat{D}_y}}\right].`$ (147)
Finally, combination of the second and the third exponential functions yields the result
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{(\widehat{A}_x+\widehat{D}_x)}}\right]\mathrm{exp}\left[2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right],`$ (148)
that is, again, a direct application of the van Cittert-Zernike theorem. For $`\mathrm{\Delta }\widehat{\theta }_x=0`$ we have
$$\xi _{cy}=\left(\frac{\pi }{\widehat{A}_y}\right)^{1/2}.$$
(149)
(B) Case $`\widehat{A}_y1`$. — This case encompasses situations with $`\widehat{z}_o1`$ as well as situations with $`\widehat{z}_o1`$.
Let us first consider $`\widehat{z}_o1`$. Looking at the integral in $`\widehat{\varphi }_y`$ in Eq. (131) it is easy to recognize that its integrand is highly oscillatory in $`\widehat{\varphi }_y`$ when $`2\widehat{\varphi }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y1`$, since $`\widehat{A}_y/(\widehat{A}_y+\widehat{D}_y)<1`$. Therefore, the integrand will contribute to the integral significatively only up to values $`2\widehat{\varphi }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y1`$. On the other hand, the terms in $`S`$ give non negligible contributions only for values of $`\widehat{\varphi }_y`$ up to order unity. As a result, it must be $`2\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y1`$. As $`\widehat{z}_o1`$ the width of $`g`$ in $`\mathrm{\Delta }\widehat{\theta }_y`$ is then of order unity. When $`\widehat{z}_o`$ becomes larger than unity, the angular width $`\mathrm{\Delta }\widehat{\theta }_y`$ will decrease and asymptotically, as $`\widehat{z}_o1`$, one will have a rapidly oscillating integrand for $`\mathrm{\Delta }\widehat{\theta }_y1/\widehat{z}_o1`$. However note that $`\mathrm{\Delta }\widehat{\theta }_y`$ must be multiplied by $`\widehat{z}_o`$ in order to obtain the correlation length which means that this remains constant and comparable with the diffraction length $`\sqrt{cL_w/\omega }`$ as $`\widehat{z}_o`$ increases. As a result, under the assumption $`\widehat{z}_o1`$, terms in $`\mathrm{\Delta }\widehat{\theta }_y`$ can be dropped in the functions $`S()`$ of Eq. (131), and the functions $`S`$ can be substituted with their limiting form $`\mathrm{sinc}`$. Moreover, in the limit for $`\widehat{A}_y1`$ and $`\widehat{D}_y1`$ we have
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{D}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2}{\widehat{A}_x+\widehat{D}_x}}\right]\mathrm{exp}\left[2\widehat{D}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right]`$ (150)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[i\left(2\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\widehat{\varphi }_y]{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x\mathrm{sinc}^2[(\widehat{\varphi }_x^2+\widehat{\varphi }_y^2)/4],`$ (151)
where the simplification in the phase under the integral in $`d\widehat{\varphi }_y`$ is possible for $`\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\widehat{\varphi }_y\widehat{D}_y/\widehat{A}_y1`$ and the exponential function $`\mathrm{exp}[\widehat{\varphi }_y^2/2(\widehat{A}_y+\widehat{D}_y)]`$ under the integral in $`d\widehat{\varphi }_y`$ can be neglected because the $`\mathrm{sinc}()`$ function has characteristic length in $`\widehat{\varphi }_y`$ of order unity. Eq. (151) is therefore valid in the limit for $`\widehat{A}_y1`$, $`\widehat{D}_y1`$ and $`\widehat{z}_o1`$. The integral in $`d\widehat{\varphi }_y`$ in Eq. (151) is simply the Fourier transform of the universal function
$$I_S(\widehat{\varphi }_y)=_{\mathrm{}}^{\mathrm{}}𝑑\widehat{\varphi }_x\mathrm{sinc}^2\left[\left(\widehat{\varphi }_x^2+\widehat{\varphi }_y^2\right)/4\right].$$
(152)
done with respect to the variable $`2\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y`$, conjugate to $`\widehat{\varphi }_y`$, that is
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{D}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2}{\widehat{A}_x+\widehat{D}_x}}\right]\mathrm{exp}\left[2\widehat{D}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right]`$ (153)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[i\left(2\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\widehat{\varphi }_y]I_S(\widehat{\varphi }_y).`$ (154)
It is not difficult to see that the width of this Fourier transform in $`\mathrm{\Delta }\widehat{\theta }_y\widehat{z}_o`$ is much smaller than the characteristic width imposed by the Gaussian exponentials outside the integration sign, their ratio being of order $`\widehat{D}_y1`$. This means that the Gaussian in $`\mathrm{\Delta }\widehat{\theta }_y^2`$ outside the integral sign is almost constant with respect to the behavior of the Fourier transform, and can be neglected, to obtain
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{D}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2}{\widehat{A}_x+\widehat{D}_x}}\right]{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[i\left(2\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\widehat{\varphi }_y\right]I_S(\widehat{\varphi }_y).`$ (155)
(156)
The assumption $`\widehat{z}_o1`$ was vital for the derivation of Eq. (156). However, for all values of $`\widehat{z}_o`$ such that $`\widehat{A}_y1`$ (and, therefore, up to the exit of the undulator at $`\widehat{z}_o=1/2`$), it is easy to see from Eq. (113) that the source is quasi-homogeneous, so that the ”anti” VCZ theorem applies and the cross-spectral density of the field at the source plane forms a Fourier couple with the angular distribution of the radiant intensity, which is simply given by $`I_S`$ as defined in Eq. (152).
It is interesting to justify the fact that $`I_S`$ is in fact the angular distribution of the radiant intensity. To this purpose, it is sufficient to note that $`I_S`$ is simply, normalization factors aside, Eq. (176) of OURS , which represents the intensity from a beam with $`\widehat{ϵ}_x\mathrm{}`$ and $`\widehat{ϵ}_y0`$. A representation of $`I_S(\widehat{\varphi }_y)`$ is given in Fig. 15.
As a result, one may therefore conclude that Eq. (156) is valid in general, for any value of $`\widehat{z}_o`$ such that $`\widehat{A}_y1`$, in the asymptotic limit $`\widehat{D}_y1`$.
It is important to note that Eq. (156) can be written as
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{D}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2}{(\widehat{A}_x+\widehat{D}_x)}}\right]\gamma (\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y),`$ (157)
where $`\gamma (\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y)`$, given by
$`\gamma (\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[i\left(2\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\widehat{\varphi }_y\right]I_S(\widehat{\varphi }_y),`$ (158)
is a universal function normalized to unity. It is possible to calculate Eq. (158) analytically. To this purpose, it is sufficient to note that the Fourier Transform
$`\gamma _1(\xi ,\eta )={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[i(\xi \widehat{\varphi }_x+\eta \widehat{\varphi }_y)\right]\mathrm{sinc}^2\left({\displaystyle \frac{\widehat{\varphi }_x^2+\widehat{\varphi }_y^2}{4}}\right)`$ (159)
can be evaluated with the help of the Bessel-Fourier formula as
$`\gamma _1(\lambda )=2\pi {\displaystyle _0^{\mathrm{}}}𝑑\varphi \varphi J_0\left(\varphi \lambda \right)\mathrm{sinc}^2\left({\displaystyle \frac{\varphi ^2}{4}}\right)`$ (160)
$`=2\pi \left[\pi +\lambda ^2\mathrm{Ci}\left({\displaystyle \frac{\lambda ^2}{2}}\right)2\mathrm{sin}\left({\displaystyle \frac{\lambda ^2}{2}}\right)2\mathrm{S}\mathrm{i}\left({\displaystyle \frac{\lambda ^2}{2}}\right)\right]`$ (161)
where $`\lambda ^2=\xi ^2+\eta ^2`$, $`\varphi ^2=\varphi _x^2+\varphi _y^2`$, $`\mathrm{Si}()`$ is the sine integral function and $`\mathrm{Ci}()`$ is the cosine integral function. As a result one has
$`\gamma (\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle \frac{2}{\pi }}\left[{\displaystyle \frac{\pi }{2}}+2\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\mathrm{Ci}\left(2\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right)\mathrm{sin}\left(2\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right)\mathrm{Si}\left(2\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right)\right].`$ (162)
(163)
The function $`\gamma (\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y)`$ is illustrated in Fig. 17.
It should be noted that $`\gamma (\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y)`$ is the spectral degree of coherence $`g`$ calculated for $`\mathrm{\Delta }\widehat{\theta }_x=0`$ in the limit for $`\widehat{A}_y1`$ and $`\widehat{D}_y1`$, in agreement with Eq. (157).
(C) Case $`\widehat{A}_y1.`$ In case (B) we have shown that terms in $`\mathrm{\Delta }\widehat{\theta }_y`$ can be dropped in the functions $`S()`$ of Eq. (131), for every value of $`\widehat{z}_o`$, and the functions $`S`$ can be substituted with their limiting form $`\mathrm{sinc}`$. In the case of $`\widehat{A}_y1`$ and $`\widehat{D}_y1`$ we have
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{D}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2}{\widehat{A}_x+\widehat{D}_x}}\right]\mathrm{exp}\left[2\widehat{D}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right]`$ (164)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[i2\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\widehat{\varphi }_y]\mathrm{exp}[{\displaystyle \frac{\widehat{\varphi }_y^2}{2\widehat{A}_y}}]`$ (165)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x\mathrm{sinc}^2[(\widehat{\varphi }_x^2+\widehat{\varphi }_y^2)/4].`$ (166)
The integral in $`d\widehat{\varphi }_y`$ in Eq. (166) is the Fourier transform of
$$F_S(\widehat{\varphi }_y)=I_S(\widehat{\varphi }_y)\mathrm{exp}\left[\frac{\widehat{\varphi }_y^2}{2\widehat{A}_y}\right]$$
(167)
done with respect to the variable $`2\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y`$, conjugate to $`\widehat{\varphi }_y`$. Similarly as before, it is not difficult to see that the width of this Fourier transform in $`\mathrm{\Delta }\widehat{\theta }_y\widehat{z}_o`$ is much smaller than the characteristic width imposed by the Gaussian exponentials outside the integration sign, their ratio being of order $`\widehat{D}_y1`$. This means that the Gaussian in $`\mathrm{\Delta }\widehat{\theta }_y^2`$ outside the integral sign is almost constant with respect to the behavior of the Fourier transform, and can be neglected. Using the definition of $`F_S`$ in Eq. (167) we have
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{D}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2}{\widehat{A}_x+\widehat{D}_x}}\right]{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[i2\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\widehat{\varphi }_y\right]F_S(\widehat{\varphi }_y).`$ (168)
(169)
At this point, no simplification is possible and numerical analysis of the problem should be undertaken in order to calculate $`\widehat{G}`$, followed by normalization according to
$`\widehat{G}(\widehat{z}_o,0,0)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_yF_S(\widehat{\varphi }_y),`$ (170)
as already said, in order to find an expression for the complex degree of coherence $`g`$. This discussion underlines again how the correlation angle $`\mathrm{\Delta }\widehat{\theta }_y`$ in the $`y`$ direction behaves. It starts from a constant value equal to the diffraction angle $`\sqrt{c/(\omega L_w)}`$ when $`\widehat{z}_o1`$, which corresponds to the maximal possible value of $`\widehat{A}_y`$ once $`\widehat{N}_y`$ is fixed. Then it decreases as $`\widehat{z}_o`$ grows. Asymptotically in limit for $`\widehat{z}_o1`$ (but still such that $`\widehat{z}_o^2\widehat{N}_y`$), it behaves as $`1/\widehat{z}_o`$, as it is clear from the fact that the function $`F_S`$, which is Fourier Transformed in Eq. (169), does not depend on any parameter. However, it should be noted that $`\mathrm{\Delta }\widehat{\theta }_y`$ must be multiplied by $`\widehat{z}_o`$ in order to obtain the correlation length, which means that this remains constant and comparable with the diffraction length $`\sqrt{cL_w/\omega }`$ as $`\widehat{z}_o`$ increases. Correlation length in the $`x`$ direction is governed instead by the Gaussian exponential function in $`\mathrm{\Delta }\widehat{\theta }_x`$ outside the integral sign in Eq. (169), exactly as in the simplified model treated in Section 4.1.
We performed some numerical calculation with the aim of giving the reader an exemplification. We set $`\mathrm{\Delta }\widehat{\theta }_x=0`$ and $`\widehat{N}_y=10`$. Assuming that we are in the asymptotic limit $`\widehat{D}1`$ we can rely on what has been said in this paragraph to calculate $`g`$ and on Eq. (104) to calculate the correlation length $`\xi _{cy}`$ in the vertical direction for any value of $`\widehat{z}_o`$. These numerical results must agree with the Van Cittert-Zernike limit for $`\widehat{A}_y1`$ treated in paragraph 5.1.1 (A): in fact, Eq. (148) and Eq. (104) yield immediately a linear dependence of $`\xi _{cy}`$ on $`\widehat{z}_o`$ as $`\widehat{A}_y1`$. Then, using paragraph 5.1.1 (B) we can extend the function $`\xi _{cy}(\widehat{z}_o)`$ to all values of $`\widehat{z}_o`$. In this way we obtain $`\xi _{cy}(\widehat{z}_o)`$ for every value of $`\widehat{z}_o`$. In Fig. 16 we plot Eq. (167) for two particular values of $`\widehat{A}_y`$: $`\widehat{A}_y=10`$ (solid line) and $`\widehat{A}_y=1`$ (dashed line). Note that the solid line Fig. 16 is still similar to Fig. 15, although the tails are changing and the width is already smaller.
The following step is to calculate the Fourier Transform of $`F_S`$ according to Eq. (169) after normalization procedure according to Eq. (170). This gives $`g`$ calculated for $`\mathrm{\Delta }\widehat{\theta }_x=0`$. Fig. 18 illustrates $`g`$ for $`\widehat{z}_o=1/2`$ as a function of $`\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y`$ at $`\mathrm{\Delta }\widehat{\theta }_x=0`$. Note that for this parameter choice we are still in the limit for $`\widehat{A}_y1`$, which explains why Fig. 18 is practically identical to the universal plot Fig. 17.
One can calculate the coherence length $`\widehat{\xi }_{cy}(\widehat{z}_o)`$ straightforwardly by means of Eq. (104). The curve obtained can be then compared, in the limit for $`\widehat{A}_y1`$, with the van Cittert-Zernike behavior illustrated in Paragraph 5.1.2 (A). Fig. 19 shows our results. The black circles represent actual numerical calculations. The asymptotic limit for $`\widehat{A}_y1`$ (VCZ theorem) is shown with a dotted line. When $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$, the coherence length $`\widehat{\xi }_{cy}(\widehat{z}_o)`$ can be calculated with the approximated formula
$$\widehat{\xi }_{cy}\left[a^2+\frac{\pi }{\widehat{A}_y}\right]^{1/2}.$$
(171)
Under the approximation of negligibly small electron beam divergence in the vertical direction, the normalized coherence length is thus a universal function of one dimensionless parameter, $`\widehat{A}_y`$. On the one hand, Eq. (171) accounts for the asymptotic behavior as $`\widehat{A}_y1`$ (VCZ theorem). On the other hand, the value of $`\widehat{\xi }_{cy}`$ in Eq. (171) approaches the constant value $`\widehat{\xi }_{cy}a`$ for asymptotically large values of $`\widehat{A}_y`$. Beside accurately reproducing the asymptotes for small and large values of $`\widehat{A}_y`$, Eq. (171) provides an accuracy of several per cent with respect to the result of numerical calculations (when $`\widehat{A}_y`$ is within the limits $`(0,\mathrm{})`$), in the whole range of the parameter $`\widehat{A}_y`$. The solid line in Fig. 19 is calculated with the approximated formula (171), where calculated $`a`$ numerically using Eq. (156) to calculate $`\widehat{\xi }_{cy}`$ and obtaining $`a1.12`$. Therefore, according to Eq. (171), when $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$, we have $`\widehat{\xi }_{cy}=const=1.12`$ at the undulator exit (with accuracy $`\widehat{D}_y1`$ and $`1/\widehat{N}_y1`$) and, in dimensional units, $`\xi _{cy}=1.12\sqrt{L_wc/\omega }`$. In this case the coherence length is a function of undulator length and wavelength due to the intrinsic divergence of the undulator radiation and $`a=1.12`$ is a universal constant.
#### 5.1.3 Case with finite divergence parameter $`\widehat{D}_y`$.
We will first discuss, in Paragraph 5.1.3 (A), the limit for $`\widehat{A}_y1`$: in this case, whatever the value of $`\widehat{D}_y`$, we will recover the VCZ theorem. Note that, although we are discussing the limit $`\widehat{N}_y1`$ there always be values of $`\widehat{z}_o`$ large enough so that $`\widehat{A}_y1`$. Further on, in Paragraph 5.1.3 (B), we will discuss the case $`\widehat{A}_y1`$: note that, since we are discussing the limit for $`\widehat{N}_y1`$, this will always be the case near the exit of the undulator at $`\widehat{z}_o=1/2`$. We will be particularly interested to this situation; in fact, the study of the case $`\widehat{A}_y1`$ near the exit of the undulator will allow us to give an explicit representation of Eq. (132).
(A) Far zone case $`\widehat{A}_y1`$. — Eq. (131) is still valid and can be calculated numerically, in principle, for any value of $`\widehat{z}_o`$. In the case $`\widehat{A}_y1`$ it can be further simplified. In fact, in this situation, the second exponential function of the right hand side of Eq. (131) limits the possible values of $`\mathrm{\Delta }\widehat{\theta }_y`$ to $`\mathrm{\Delta }\widehat{\theta }_y1`$. Moreover $`\widehat{z}_o1`$ so that Eq. (166) is valid in this case. In particular, remembering the definition of $`F_S(\widehat{\varphi }_y)`$ in Eq. (167) we have
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{(\widehat{A}_x+\widehat{D}_x)}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y}{(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (172)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[i{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\widehat{\varphi }_y]F_S(\widehat{\varphi }_y).`$ (173)
It is not difficult to see that the ratio between the characteristic width in $`\mathrm{\Delta }\widehat{\theta }_y`$ of the exponential function outside the integral sign and the exponential function in inside the integral sign in Eq. (173) is of order $`\sqrt{\widehat{A}_y/\widehat{D}_y}`$ and it is always much smaller than unity unless $`\widehat{D}_y1`$: such a case has already been treated before in Paragraph 5.1.1 (A), and has been shown to obey the VCZ theorem <sup>5</sup><sup>5</sup>5Alternatively one may note directly that $`\widehat{\varphi }_y`$ can only range over values much smaller than unity. As a result, the dependence of $`I_S`$ on $`\widehat{\varphi }_y`$ can be dropped, giving an extra normalization constant to be disposed of. Then, the integral in $`d\widehat{\varphi }`$ performed giving, as in Paragraph 5.1.1 (A), $`\mathrm{exp}[2\widehat{A}_y^2\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2/(\widehat{A}_y+\widehat{D}_y)]`$ to be combined with the exponential function in $`\mathrm{\Delta }\widehat{\theta }_y`$ outside the integral sign, giving $`\mathrm{exp}[2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2]`$.. For all other values of $`\widehat{D}_y`$ we have, automatically, $`\widehat{D}_y\widehat{A}_y`$, so that we can neglect the integral in $`d\widehat{\varphi }_y`$ and we get back once more the Van Cittert-Zernike regime. To sum up, we obtain the following expression for $`g`$, which is valid for $`\widehat{A}_y1`$ with no restrictions on $`\widehat{D}_y`$:
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{(\widehat{A}_x+\widehat{D}_x)}}\right]\mathrm{exp}\left[2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right].`$ (174)
Calculation of the coherence length from Eq. (174) at $`\mathrm{\Delta }\widehat{\theta }_x=0`$ gives once more the behavior
$$\xi _{cy}=\left(\frac{\pi }{\widehat{A}_y}\right)^{1/2},$$
(175)
which is consistent with the partial result in Paragraph 5.1.1 (A).
(B) Near zone, $`\widehat{A}_y1`$. — Equation (166) is still valid in this case and following the same line of reasoning as paragraph 5.1.2 (B) one gets
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{D}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2}{(\widehat{A}_x+\widehat{D}_x)}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{D}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2}{(\widehat{A}_y+\widehat{D}_y)}}\right]\gamma (\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y),`$ (176)
(177)
valid for $`\widehat{A}_y1`$ and arbitrary $`\widehat{D}_y`$. Note that since we are working in the limit for $`\widehat{N}_y1`$, for $`\widehat{z}_o=1/2`$ we have $`\widehat{A}_y1`$. We can see from Eq. (177) that, for finite values of $`\widehat{D}_y`$, the cross-spectral density $`g`$, evaluated at the exit of the undulator for $`\mathrm{\Delta }\widehat{\theta }_x=0`$, is given by the product of the exponential function $`\mathrm{exp}[2\widehat{D}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2]`$ with the function illustrated in Fig. 18. This remark is intuitively sound. Since $`\widehat{N}_y1`$ in fact we have weakly quasi-homogeneous wavefronts near the exit of the undulator, and we can use the ”anti” VCZ theorem to conclude that $`g`$ must for a Fourier couple with the intensity distribution in the far zone. This will simply be, for any arbitrary $`\widehat{D}_y`$, a convolution between a Gaussian distribution with rms width equal to $`\sqrt{\widehat{D}_y}`$ and $`I_S`$, which is the angular distribution of radiant intensity for a beam with $`\widehat{ϵ}_x\mathrm{}`$ and $`\widehat{ϵ}_y0`$. Finally, the Fourier transform of such a convolution between two function is simply given by the product of the Fourier transforms of the two functions.
(C) Approximate formula. — With in mind Eq. (137), Eq. (171) and Eq. (175) we make the working hypothesis that Eq. (132) has the form
$$\widehat{\xi }_{cy}=\left(\frac{\pi }{\widehat{D}_{\mathrm{eff}}(\widehat{D}_y)}+\frac{\pi }{\widehat{A}_y}\right)^{1/2}.$$
(178)
Within the assumption $`\widehat{N}_y1`$, we have seen that if $`\widehat{D}_y1`$ Eq. (137) is valid with relative accuracy $`1/\widehat{D}_y`$. This means that, in this limit, $`\widehat{D}_{\mathrm{eff}}(\widehat{D}_y)=\widehat{D}_y`$. Then we have seen that if $`\widehat{D}_y1`$ Eq. (171) holds with accuracy $`1/\widehat{N}_y`$. This means that $`\widehat{D}_{\mathrm{eff}}(0)=2.50`$. Moreover we have seen that in the far zone case, in the limit for $`\widehat{A}_y1`$, the VCZ theorem holds, in agreement with Eq. (175).
We are now in position to calculate $`\widehat{D}_{\mathrm{eff}}=f(\widehat{D}_y)`$ for any value of $`\widehat{D}_y`$. Evaluation of Eq. (177) at $`\widehat{z}_o=1/2`$, followed by normalization according to $`\widehat{G}(\widehat{z}_o,0,0)=1`$ gives the function $`g`$. Further integration of $`|g|^2`$ to calculate the correlation function allows to recover $`\widehat{D}_{\mathrm{eff}}=f(\widehat{D}_y)`$ as plotted in Fig. 20.
One may choose to calculate $`\widehat{D}_{\mathrm{eff}}=f(\widehat{D}_y)`$ numerically, but it is also possible to use, with reasonable accuracy, the following analytical interpolation of $`\widehat{D}_{\mathrm{eff}}`$:
$$\widehat{D}_{\mathrm{eff}}\widehat{D}_{\mathrm{eff}}(0)+\widehat{D}_y=2.50+\widehat{D}_y.$$
(179)
There is of some interest to compare the exact and interpolated $`\widehat{D}_{\mathrm{eff}}(\widehat{D}_y)`$ functions. Fig. 21 shows the function
$$\mathrm{\Delta }(\widehat{D}_y)=\left|1\frac{\widehat{D}_{\mathrm{eff}}(\widehat{D}_y)}{\widehat{D}_{\mathrm{eff}}(0)+\widehat{D}_y}\right|.$$
(180)
There is seen to be good agreement between the interpolated and exact $`\widehat{D}_{\mathrm{eff}}`$ functions for small and large value of $`\widehat{D}_y`$. Noticeable discrepancies for $`\widehat{D}_y`$ close to unity are, anyway, less than $`10\%`$.
Our conclusive result is, therefore, the following: when $`\widehat{N}_y1`$ and $`\widehat{D}_y`$ assumes arbitrary values we have:
$$\widehat{\xi }_{cy}\left(\frac{\pi }{2.50+\widehat{D}_y}+\frac{\pi }{\widehat{A}_y}\right)^{1/2}.$$
(181)
Also, it is important to remember that under conditions $`\widehat{N}_y1`$ and arbitrary $`\widehat{D}_y`$ the spectral degree of coherence is given by Eq. (177) in the near zone, and by Eq. (174) in the far zone.
The final step is to check that our main work hypothesis, i.e. that the coherence length has the form in Eq. (178) is correct. This can be done comparing Eq. (181) with numerical calculations for any given value of $`\widehat{N}_y1`$ and a finite $`\widehat{D}_y`$, which give a good agreement.
### 5.2 Very large divergence $`\widehat{D}_y1`$, arbitrary Fresnel number $`\widehat{N}_y`$
We now move on to treat the case with arbitrary beam transverse size compared with the diffraction size (i.e. arbitrary Fresnel number $`\widehat{N}_y`$) and large divergence compared with the diffraction angle (i.e. $`\widehat{D}_y1`$). The particular case for $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$ overlaps with the previous Section 5.1 and has already been treated in Section 5.1.1. The conclusion was that $`\xi _{cy}=(\pi /\widehat{D}_y+\pi /\widehat{A}_y)^{1/2}`$. In all the other remaining cases Eq. (131) can be simplified as follows:
$`\widehat{G}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{\widehat{A}_x+\widehat{D}_x}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (182)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_xS^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2],`$ (183)
that is
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{\widehat{A}_x+\widehat{D}_x}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y}{\widehat{A}_y+\widehat{D}_y}}\right]\stackrel{~}{f}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_y),`$ (184)
(185)
where
$`\stackrel{~}{f}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_xS^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2],`$ (186)
As is shown in Appendix C, having defined
$`\beta (\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right],`$ (187)
we have the important result
$$\stackrel{~}{f}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_y)=\beta (\mathrm{\Delta }\widehat{\theta }_y)$$
(188)
for every choice of $`\widehat{z}_o`$. $`\beta `$ is defined in such a way to be normalized to unity. If we account for Eq. (188) we obtain the following expression for the spectral degree of coherence $`g`$:
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x}{\widehat{A}_x+\widehat{D}_x}}\right]\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y}{\widehat{A}_y+\widehat{D}_y}}\right]\beta (\mathrm{\Delta }\widehat{\theta }_y).`$ (189)
(190)
It should be noted that, as $`\widehat{N}_y1`$, the width of the gaussian function in $`\mathrm{\Delta }\widehat{\theta }_y`$ in Eq. (190) becomes much smaller than unity and the function $`\beta `$ can be considered constant and drops out of the normalized expression for $`g`$. So, even if we did not analyze here the limit for $`\widehat{N}_y1`$ (we did it in Paragraph 5.1.1), we see that the limit of Eq. (190) for $`\widehat{N}_y1`$ restitutes the results found in Paragraph 5.1.1. Therefore we conclude that Eq. (190) is valid for any value of $`\widehat{N}_y`$.
As is shown in Appendix C, the function $`\beta (\mathrm{\Delta }\widehat{\theta }_y)`$ can also be calculated as
$$\beta (\mathrm{\Delta }\widehat{\theta }_y)=\frac{1}{\pi }_0^{\mathrm{}}𝑑\alpha \alpha J_0\left(\alpha \frac{\mathrm{\Delta }\widehat{\theta }_y}{2}\right)\left[\pi 2\mathrm{S}\mathrm{i}(\alpha ^2)\right]^2.$$
(191)
where $`\mathrm{Si}`$ indicates the sine integral function. The representation of $`\beta `$ in Eq. (191) is easier to deal with numerically, because it involves a one-dimensional integration only. Performing the integral, one can tabulate $`|\beta |`$ to obtain the plot in Fig. 22. This is an universal plot. It should be noted that $`\beta (\mathrm{\Delta }\widehat{\theta }_y)`$ is the spectral degree of coherence $`g`$ calculated for $`\mathrm{\Delta }\widehat{\theta }_x=0`$ in the limit for $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$. Its generalization for arbitrary $`\widehat{N}_y`$ is given by Eq. (190). Using Eq. (190) and the tabulated values for the universal function $`\beta `$ we can therefore calculate $`g`$ numerically for any choice of $`\widehat{N}_y`$ and subsequently, we can calculate the coherence length $`\xi _{cy}(\widehat{z}_o)`$. For instance, in the particular case $`\widehat{N}_y=10`$ and $`\mathrm{\Delta }\widehat{\theta }_x=0`$ we obtain the simple linear behavior
$$\xi _{cy}(\widehat{z}_o)=0.54\widehat{z}_o$$
(192)
For a generic value of $`\widehat{N}_y`$, one can introduce an effective function $`\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)`$ so that
$$\xi _{cy}(\widehat{z}_o)=\sqrt{\frac{\pi }{\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)}}\widehat{z}_o$$
(193)
On the one hand, the function $`\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)`$ can be computed numerically, similarly as we did for the particular case $`\widehat{N}_y=10`$. $`\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)`$ is represented by the solid line in Fig. 23.
On the other hand, one may also use an interpolation for $`\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)`$. First, numerical calculations tell us that, in the particular case $`\widehat{N}_y0`$, we have
$$\xi _{cy}(\widehat{z}_o)=\sqrt{\frac{\pi }{0.35}}\widehat{z}_o.$$
(194)
Second, as $`\widehat{D}_y1`$ and $`\widehat{N}_y\widehat{D}_y`$ we have
$$\xi _{cy}(\widehat{z}_o)\sqrt{\frac{\pi }{\widehat{N}_y}}\widehat{z}_o.$$
(195)
The simpler interpolated formula which satisfies both asymptotes is therefore:
$$\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)\widehat{N}_y+0.35,$$
(196)
The interpolation of $`\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)`$ is represented by black circles line in Fig. 23.
There is of some interest to compare the exact and interpolated $`\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)`$ functions. Fig. 24 shows the function
$$\mathrm{\Delta }(\widehat{N}_y)=\left|1\frac{\widehat{N}_{\mathrm{eff}}(\widehat{N}_y)}{\widehat{N}_{\mathrm{eff}}(0)+\widehat{N}_y}\right|.$$
(197)
There is seen to be good agreement between the interpolated and exact $`\widehat{N}_{\mathrm{eff}}`$ functions for small and large value of $`\widehat{N}_y`$. Noticeable discrepancies for $`\widehat{N}_y`$ close to unity are, anyway, less than $`13\%`$.
Since in the case $`\widehat{D}_y1`$ and $`\widehat{N}_y1`$ we concluded that $`\xi _{cy}=(\pi /\widehat{D}_y+\pi /\widehat{A}_y)^{1/2}`$, we can formulate the hypothesis that, for $`\widehat{D}_y1`$ and generic value of $`\widehat{N}_y`$ one has
$$\xi _{cy}(\widehat{z}_o)=\left(\frac{\pi }{\widehat{D}_y}+\frac{\pi }{0.35+\widehat{N}_y}\widehat{z}_o^2\right)^{1/2}.$$
(198)
The final step is to check that such hypothesis is correct. This can be done comparing Eq. (198) with numerical calculations for any given value of $`\widehat{D}_y1`$ and finite $`\widehat{N}_y`$, which give a good agreement.
### 5.3 Conditions for the source to be quasi-homogeneous
Up to this moment we discussed, for simplicity, the case $`\overline{\theta }_x=\overline{\theta }_y=0`$. We will now treat the generic case with arbitrary $`\overline{\theta }_x`$ and $`\overline{\theta }_y`$. This discussion will reduce to the relation between the weakly quasi-homogeneous condition and our assumptions $`\widehat{N}_x1`$, $`\widehat{D}_x1`$ and either $`\widehat{N}_y1`$ or $`\widehat{D}_y1`$. When we deal with weakly quasi-homogeneous wavefront, the results found for $`\overline{\theta }_x=\overline{\theta }_y=0`$ have extended validity for generic values of $`\overline{\theta }_x`$ and $`\overline{\theta }_y`$. As already said before though, we will find that the wavefronts are not always weakly quasi-homogeneous in the vertical $`y`$ direction. In this case, previously found results are only valid in the particular case $`\overline{\theta }_y=0`$. Let us consider the situation in more detail. We will start with Eq. (128), that may also be written as
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (199)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (200)
$`\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (201)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{(\overline{\theta }_y+\widehat{\varphi }_y)^2+4i\widehat{\varphi }_y\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{2(\widehat{A}_y+\widehat{D}_y)}}]`$ (202)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_xS^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2].`$ (203)
#### 5.3.1 Very large divergence parameter $`\widehat{D}_y1`$
Because of the properties of the $`S`$ function $`\widehat{\varphi }_y`$ can only change of a quantity $`\mathrm{\Delta }\widehat{\varphi }_y1`$, otherwise the $`S`$ functions will drop to zero. Then, since $`\widehat{D}_y1`$ we have from Eq. (203)
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (204)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (205)
$`\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (206)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_y^2}{2(\widehat{A}_y+\widehat{D}_y)}}\right]{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[{\displaystyle \frac{2i\widehat{\varphi }_y\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (207)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_xS^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2].`$ (208)
that is obviously weakly quasi-homogeneous.
#### 5.3.2 Very large Fresnel number $`\widehat{N}_y1`$.
(A) Case $`\widehat{A}_y1`$.— In the case $`\widehat{A}_y1`$ we can follow the same line of reasoning in Section 5.3.1 simply replacing the roles of $`\widehat{D}_y`$ with $`\widehat{A}_y`$, obtaining
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (209)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (210)
$`\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (211)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_y^2}{2(\widehat{A}_y+\widehat{D}_y)}}\right]{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[{\displaystyle \frac{2i\widehat{\varphi }_y\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{(\widehat{A}_y+\widehat{D}_y)}}\right]`$ (212)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_xS^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2].`$ (213)
that is obviously weakly quasi-homogeneous.
(B) Case $`\widehat{A}_y1.`$— In this situation it must be $`\widehat{z}_o1`$. Therefore we can substitute all $`S`$ functions with $`\mathrm{sinc}`$ functions. The case $`\widehat{D}_y1`$ has been already treated. Let us, therefore, first assume $`\widehat{D}_y1`$. Using the fact that $`\mathrm{\Delta }\widehat{\theta }_y^21/\widehat{z}_o^21`$, Eq. (203) gives directly
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (214)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]\mathrm{exp}\left[2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right]`$ (215)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{(\widehat{\varphi }_y+\overline{\theta }_y)^2}{2\widehat{D}_y}}]I_S(\widehat{\varphi }_y),`$ (216)
that is obviously weakly quasi-homogeneous.
Now let us consider the case $`\widehat{D}_y1`$. Eq. (203) can be simplified on the assumptions $`\widehat{A}_y1`$ and $`\widehat{D}_y1`$. In fact, the Gaussian exponential factor inside the integral in $`d\widehat{\varphi }_y`$ in Eq. (203) imposes a maximal value $`(\overline{\theta }_y+\widehat{\varphi }_y)^2\widehat{A}_y+\widehat{D}_y1`$.
Simultaneously, from the Gaussian exponential factor in $`\mathrm{\Delta }\widehat{\theta }_y`$ outside the integral, we have a condition for the maximal value of $`\mathrm{\Delta }\widehat{\theta }_y^21/\widehat{z}_o^21`$ since $`\mathrm{\Delta }\widehat{\theta }_y^21/(\widehat{A}_y\widehat{z}_o^2)1`$.
As a result, the dependence of $`S`$ on $`(\widehat{\varphi }+\mathrm{\Delta }\widehat{\theta }_y)^2`$ can be substituted by a dependence on $`\overline{\theta }_y`$ and, as has already been said, the $`S`$ functions can be substituted with $`\mathrm{sinc}`$ functions, giving
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (217)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (218)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]I_S(\overline{\theta }_y)`$ (219)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{(\overline{\theta }_y+\widehat{\varphi }_y)^2+2\widehat{\varphi }_y\left(2i\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)}{2(\widehat{A}_y+\widehat{D}_y)}}].`$ (220)
The integral in $`d\widehat{\varphi }_y`$ can be performed giving
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (221)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (222)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (223)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y^2\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2}{\widehat{A}_y+\widehat{D}_y}}\right]\mathrm{exp}\left[{\displaystyle \frac{2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]I_S(\overline{\theta }_y).`$ (224)
Normalizing $`\widehat{G}`$ in such a way that $`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,0,0)=1`$ we obtain the following expression for the spectral degree of coherence $`g`$:
$`g(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (225)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]\mathrm{exp}\left[2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right],`$ (226)
which generalizes Eq. (147) for any value of $`\overline{\theta }_x`$ and $`\overline{\theta }_y`$, and shows weak quasi-homogeneity of the wavefronts.
(C) Case $`\widehat{A}_y1.`$— Since $`\widehat{N}_y1`$ and $`\widehat{z}_o1`$, the exponential function in $`\mathrm{\Delta }\widehat{\theta }_y`$ outside the integral sign of Eq. (128) impose $`\mathrm{\Delta }\widehat{\theta }_y1`$ so that
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (227)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (228)
$`\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (229)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{(\overline{\theta }_y+\widehat{\varphi }_y)^2+4i\widehat{\varphi }_y\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{2(\widehat{A}_y+\widehat{D}_y)}}]I_S(\widehat{\varphi }_y).`$ (230)
The intensity distribution can be found from Eq. (230) setting $`\mathrm{\Delta }\widehat{\theta }_x=\mathrm{\Delta }\widehat{\theta }_y=0`$ thus obtaining
$`I(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y)=\mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2}{2(\widehat{A}_x+\widehat{D}_x)}}\right]{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[{\displaystyle \frac{(\overline{\theta }_y+\widehat{\varphi }_y)^2}{2(\widehat{A}_y+\widehat{D}_y)}}\right]I_S(\widehat{\varphi }_y).`$ (231)
(232)
It is evident by inspection that we cannot factorize Eq. (230) to obtain $`|\widehat{G}|=I(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y)w(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)`$ (where have put $`|g|=w(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)`$). As a result we conclude that, in this case, the wavefront is not quasi-homogeneous, not even in the weak sense.
#### 5.3.3 Discussion.
Let us discuss the results obtained in this analysis of quasi-homogeneity.
In Section 4.1 we have seen that weakly quasi-homogeneous wavefronts which are not quasi-homogeneous in the usual sense are present in the far field, when the VCZ theorem holds. From our analysis in the one-dimensional framework, the notion of far zone arises in the $`x`$ direction, when the apparent angular dimension $`\widehat{A}_x`$ of the source is much smaller than the divergence of the radiation beam that, in this case, can be identified with the electron beam divergence $`\widehat{D}_x1`$. Both $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$. If $`\widehat{A}_x\widehat{D}_x`$ the wavefront is quasi-homogenous but only in the weak sense, one is in the far zone and the VCZ theorem applies. If $`\widehat{A}_x\widehat{D}_x`$ one is in the near field zone and the wavefront is quasi-homogeneous in the usual sense. Transition from the near to the far zone always involves, in this case, weakly quasi-homogeneous wavefront.
In the two-dimensional framework studied in the present Section, as for the one-dimensional model, both $`\widehat{N}_x`$ and $`\widehat{D}_x1`$. The $`x`$ and the $`y`$ coordinates appear factorized and the far zone applies now separately to both $`x`$ and $`y`$ directions, meaning that $`\widehat{A}_{x,y}`$ of the source is much smaller than the divergence of the radiation beam. In other words, we are in the far zone as soon as the (square of the) beam size at a given $`\widehat{z}_o`$ begins to be much larger than the (square of the) initial size of radiation (i.e. much larger than the Fresnel numbers $`\widehat{N}_{x,y}`$).
In the case $`\widehat{D}_y1`$, independently on the value of $`\widehat{N}_y`$, we will always have weakly quasi-homogeneous (but not always quasi-homogeneous in the usual sense!) wavefronts at any distance $`\widehat{z}_o`$. This is exactly the situation discussed for the $`x`$ direction, where transitions from the near ($`\widehat{A}_y\widehat{D}_y`$) to the far field ($`\widehat{A}_y\widehat{D}_y`$) involve only weakly quasi-homogeneous wavefronts. Moreover, in the far field zone, the VCZ theorem is valid.
In the case with $`\widehat{D}_y1`$, $`\widehat{N}_y1`$ we will have quasi-homogeneous wavefronts in the usual sense at $`\widehat{z}_o1`$ (near zone) and weak quasi-homogeneous wavefronts in the far zone, at $`\widehat{z}_o1`$.
Finally, if $`\widehat{N}_y1`$, $`\widehat{D}_y1`$ and $`\widehat{A}_y1`$, i.e. for distances $`\widehat{z}_o\sqrt{\widehat{N}_y}`$, we have seen that the wavefront is non quasi-homogenous, not even in the weak sense.
These results can be seen in terms of convolution between the Gaussian distribution of intensity associated to the electron beam emittance and the distribution of intensity due to intrinsic properties of undulator radiation. The only case, among those treated up to now, when such a convolution does not simplify into the product of separate factors is when $`\widehat{N}_y1`$, $`\widehat{D}_y`$ not much larger than unity and $`\widehat{A}_y1`$: in this case we do not have quasi-homogeneous wavefronts, not even in the weak sense.
## 6 Radiation from some non-homogeneous undulator sources
In Section 4.1 we treated a simplified situation with $`\widehat{N}_x1`$, $`\widehat{D}_x1`$ and $`\widehat{\theta }_{y1}=\widehat{\theta }_{y2}`$. Moreover, for simplicity of calculations we assumed $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$. In Section 5, instead, we treated the case of electron beams with $`\widehat{N}_x1`$, $`\widehat{D}_x1`$ and either $`\widehat{N}_y1`$ or $`\widehat{D}_y1`$ (or both). It is important to note that, in the $`x`$ direction, the results obtained in Section 4.1 are the same as the one in Section 5. In fact as $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$, the cross-spectral density factorizes in the product of two contributions depending separately on the $`x`$ and $`y`$ coordinates, and under $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$ the derivation of the $`x`$-dependent factor is always the same.
We have seen that in some of the cases discussed in Section 5.3, the assumptions $`\widehat{N}_x1`$, $`\widehat{D}_x1`$ and either $`\widehat{N}_y1`$ or $`\widehat{D}_y1`$ (or both) were enough to guarantee that the wavefront is weakly quasi-homogeneous in the sense specified by Eq. (38). In this Section we will extend our analytical investigations to some cases outside the range of parameters treated before, where the weakly quasi-homogeneous assumption is not fulfilled in the far zone. In particular, we will demonstrate that, under conditions $`\widehat{N}_x1`$, $`\widehat{D}_x1`$ and both $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$, wavefronts are not weakly quasi-homogenous in the vertical $`y`$ direction (although they are in the horizontal $`x`$ direction).
First, in Section 6.1 we will analyze the case $`\widehat{N}_x1`$, $`\widehat{D}_x1`$. Assuming a vertical emittance of the electron beam much smaller than the horizontal emittance $`\widehat{ϵ}_y\widehat{ϵ}_x`$, we have, automatically $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$. This corresponds to a practically important situation. For instance, consider a VUV beamline at a third generation light source with $`\lambda =30`$ nm, $`ϵ_x=310^9`$ m, and $`ϵ_y=0.0310^9`$ m, i.e. $`\widehat{ϵ}_x=0.6`$ and $`\widehat{ϵ}_y=610^3`$. If $`\widehat{\beta }_x=3`$ we would have $`\widehat{N}_x=2`$ and $`\widehat{D}_x=0.2`$. Second, in Section 6.2 we will study the situation $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$ with $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$, that will give us back also the limiting case for $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$ with $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$ (already discussed in Section 4.1). The situation with finite vertical Fresnel number and negligible vertical divergence (compared with the diffraction angle) is a very practical one: for instance, given a third generation light source with $`\lambda =1\AA `$ and $`ϵ_y=10^{11}`$ m, i.e. $`\widehat{ϵ}_y=0.6`$, a value $`\widehat{\beta }_y=6`$ corresponds to $`\widehat{D}_y=0.1`$ and $`\widehat{N}_y=3.6`$.
Although in these two cases, the weakly quasi-homogeneous assumption is not fulfilled we will see that the choice $`\widehat{z}_o1`$ will allow us to treat these situations in analogy with respect to some weakly quasi-homogeneous case we already dealt with.
### 6.1 Case $`\widehat{N}_x1`$, $`\widehat{D}_x1`$.
Assuming a vertical emittance of the ring much smaller than the horizontal emittance $`\widehat{ϵ}_y\widehat{ϵ}_x`$, we have, automatically, $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$. We start with Eq. (113), that can be specialized to an equation dependent on the $`x`$ coordinates only and, in the case for $`\widehat{z}_o1`$, can be written as
$`\widehat{G}(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })={\displaystyle \frac{\mathrm{exp}\left(i2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\right)}{\sqrt{2\pi (\widehat{N}/\widehat{z}_o^2+\widehat{D})}}}\mathrm{exp}\left[{\displaystyle \frac{4\widehat{N}\mathrm{\Delta }\widehat{\theta }^2\widehat{D}+4i(\widehat{N}/\widehat{z}_o)\overline{\theta }\mathrm{\Delta }\widehat{\theta }}{2(\widehat{N}/\widehat{z}_o^2+\widehat{D})}}\right]`$ (233)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }\mathrm{exp}[{\displaystyle \frac{\left(\widehat{\varphi }+\overline{\theta }\right)^2+4i\widehat{\varphi }(\widehat{N}/\widehat{z}_o)\mathrm{\Delta }\widehat{\theta }}{2(\widehat{N}/\widehat{z}_o^2+\widehat{D})}}]`$ (234)
$`\times \mathrm{sinc}\left[(\widehat{\varphi }\mathrm{\Delta }\widehat{\theta })^2/4\right]\mathrm{sinc}\left[(\widehat{\varphi }+\mathrm{\Delta }\widehat{\theta })^2/4\right],`$ (235)
where we systematically omitted $`x`$ subscripts. Since $`\widehat{A}+\widehat{D}1`$, the exponential factor in $`(\varphi +\overline{\theta })^2`$ inside the integral sign in Eq. (235) behaves like a $`\delta `$-Dirac function with respect to the $`\mathrm{sinc}()`$ functions inside the same integral: as a result we can substitute $`\varphi `$ with $`\overline{\theta }`$ in the $`\mathrm{sinc}()`$ functions, which drop out of the integral sign. Then, the integral in $`d\widehat{\varphi }`$ can be calculated analytically so that we have:
$`\widehat{G}(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })=\mathrm{exp}\left(i2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\right)\mathrm{exp}\left[{\displaystyle \frac{2\widehat{N}\mathrm{\Delta }\widehat{\theta }^2\widehat{D}}{\widehat{N}/\widehat{z}_o^2+\widehat{D}}}\right]`$ (236)
$`\times \mathrm{sinc}\left[(\overline{\theta }\mathrm{\Delta }\widehat{\theta })^2/4\right]\mathrm{sinc}\left[(\overline{\theta }+\mathrm{\Delta }\widehat{\theta })^2/4\right]`$ (237)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2(\widehat{N}^2/\widehat{z}_o^2)\mathrm{\Delta }\widehat{\theta }^2}{\widehat{N}/\widehat{z}_o^2+\widehat{D}}}\right].`$ (238)
Combination of the second and the third exponential function yields
$`\widehat{G}(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })=\mathrm{exp}\left(i2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\right)\mathrm{exp}\left[2\widehat{N}\mathrm{\Delta }\widehat{\theta }^2\right]`$ (239)
$`\times \mathrm{sinc}\left[(\overline{\theta }\mathrm{\Delta }\widehat{\theta })^2/4\right]\mathrm{sinc}\left[(\overline{\theta }+\mathrm{\Delta }\widehat{\theta })^2/4\right].`$ (240)
Finally, in order to obtain the degree of spectral coherence $`g`$ we should normalize $`\widehat{G}`$ according to Eq. (107) <sup>6</sup><sup>6</sup>6Note that, in this case, we are dealing with non quasi-homogeneous wavefronts and, as has already been said, normalizing according to $`\widehat{G}(\widehat{z}_o,\overline{\theta },0)=1`$ is not the same of normalizing according to Eq. (107). thus obtaining
$`g(\widehat{z}_o,\overline{\theta },\mathrm{\Delta }\widehat{\theta })=\mathrm{exp}\left(i2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o\right)\mathrm{exp}\left[2\widehat{N}\mathrm{\Delta }\widehat{\theta }^2\right].`$ (241)
According to Eq. (241) the spectral degree of coherence $`g`$ is such that $`|g|`$ is only function of $`\mathrm{\Delta }\widehat{\theta }`$. Moreover, the dependence of $`g`$ on the phase $`2\overline{\theta }\mathrm{\Delta }\widehat{\theta }\widehat{z}_o`$ is a feature for radiation from Schell’s model sources in the far field, that we have already encountered many times in the study of weakly quasi-homogenous cases. As a result we can conclude that the radiation of the undulator source at $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$ represents the far field radiation of a Schell’s model source. Finally, it should be noted that in a two-pinhole experiment, for any vertical position of the pinholes, the fringe visibility, i.e. the modulus of the spectral degree of coherence depends only on the separation along the horizontal $`x`$ direction. As a result, while in Section 4.1 we put $`\widehat{\theta }_{y1}=\widehat{\theta }_{y2}`$, thus selecting from the very beginning a horizontal plane, the present case can be fully described by a one-dimensional model, independently on the choice of transverse coordinates of the pinholes.
### 6.2 Case $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$ with $`\widehat{N}_x1`$, $`\widehat{D}_x1`$
In this situation we go back to a two-dimensional model. When $`\widehat{z}_o1`$ we have $`\widehat{A}_y1`$, which is the limiting case treated in Paragraph 5.1.2 (A): the difference is that, now $`\widehat{N}_y1`$. Let us start with Eq. (128) written as
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (242)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (243)
$`\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (244)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{(\overline{\theta }_y+\widehat{\varphi }_y)^2+4i\widehat{\varphi }_y\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{2(\widehat{A}_y+\widehat{D}_y)}}]`$ (245)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_xS^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2],`$ (246)
$`\widehat{A}_y1`$ and $`\widehat{D}_y1`$ impose a maximal value of $`(\overline{\theta }_y+\widehat{\varphi }_y)^2\widehat{A}_y+\widehat{D}_y1`$. Moreover the $`S`$ functions can be substituted with $`\mathrm{sinc}`$ functions since we are working in the limit for $`\widehat{z}_o1`$. Then, from Eq. (246) we obtain
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (247)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (248)
$`\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (249)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]`$ (250)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_y\mathrm{exp}[{\displaystyle \frac{(\overline{\theta }_y+\widehat{\varphi }_y)^2+4i\widehat{\varphi }_y\widehat{A}_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{2(\widehat{A}_y+\widehat{D}_y)}}],`$ (251)
As done before, the integral in $`d\widehat{\varphi }_y`$ can be performed giving
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (252)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]`$ (253)
$`\mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\widehat{D}_y+2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (254)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_y^2\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2}{\widehat{A}_y+\widehat{D}_y}}\right]\mathrm{exp}\left[{\displaystyle \frac{2i\widehat{A}_y\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y}{\widehat{A}_y+\widehat{D}_y}}\right]`$ (255)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right],`$ (256)
that is
$`\widehat{G}(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (257)
$`\times \mathrm{exp}\left[{\displaystyle \frac{\overline{\theta }_x^2+4\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+4i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{2(\widehat{A}_x+\widehat{D}_x)}}\right]\mathrm{exp}\left[2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right]`$ (258)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right].`$ (259)
Finally, normalization according to Eq. (107) yields
$`g(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (260)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+2i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{\widehat{A}_x+\widehat{D}_x}}\right]\mathrm{exp}\left[2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2\right]`$ (261)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\widehat{\varphi }_x\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]`$ (262)
$`\times \left[{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x\mathrm{sinc}^2\left\{{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right\}\right]^{1/2}`$ (263)
$`\times \left[{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x\mathrm{sinc}^2\left\{{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right\}\right]^{1/2}.`$ (264)
Obviously $`|g|`$ is a function of both $`\overline{\theta }_y`$ and $`\mathrm{\Delta }\theta `$, so that in this case we have neither weak quasi-homogeneity, neither wavefronts which can be described by Schell’s model. However it should be noted that the integrals in Eq. (264) do not contain any parametric dependence.
It is interesting to have some final comment on Eq. (264). In the limit for $`\widehat{N}_y1`$, we recover the case discussed in Section 4.1, the only difference being that we did not set $`\widehat{\theta }_{y1}=\widehat{\theta }_{y2}=0`$: in Section 4.1 we chose to deal with a one-dimensional model putting ourselves on the horizontal plane. Now, Eq. (264) allows to study the full two-dimensional situation in the limit for $`\widehat{N}_y1`$. In this case we see that the exponential function $`\mathrm{exp}[2\widehat{A}_y\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_y^2]`$ can be neglected because we have a maximum value of $`\mathrm{\Delta }\widehat{\theta }_y1`$. As a result we obtain:
$`g(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x\right]`$ (265)
$`\times \mathrm{exp}\left[{\displaystyle \frac{2\widehat{A}_x\widehat{z}_o^2\mathrm{\Delta }\widehat{\theta }_x^2\widehat{D}_x+2i\widehat{A}_x\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x}{\widehat{A}_x+\widehat{D}_x}}\right]\mathrm{exp}\left[i2\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right]\chi (\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_y),`$ (266)
where
$`\chi (\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]`$ (267)
$`\times \left[{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x\mathrm{sinc}^2\left\{{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right\}\right]^{1/2}`$ (268)
$`\times \left[{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x\mathrm{sinc}^2\left\{{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right\}\right]^{1/2}.`$ (269)
It should be noted that Eq. (269) does not depend on parameters and is, in fact, a universal function. Besides a geometrical factor $`\mathrm{exp}[i2\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y]`$, the function $`\chi `$ represents the spectral degree of coherence in the vertical direction, once the horizontal coordinates are fixed. The fact that it is a universal function means that even in the case of zero vertical emittance we never have full coherence in the vertical direction. On the one hand, this phenomenon can be seen to be an influence of the presence of horizontal emittance on the vertical coherence properties of the photon beam, as the integral in $`d\widehat{\varphi }_x`$ in $`\chi `$ comes from an integration over the horizontal electron beam distribution. On the other hand, being $`\chi `$ a universal function, the influence of the horizontal emittance on the vertical coherence does not depend, in the limit for $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$, on the actual values of $`\widehat{N}_x`$ and $`\widehat{D}_x`$. It is straightforward to see that $`\chi `$ is symmetric with respect to $`\mathrm{\Delta }\widehat{\theta }_y`$ and with respect to the exchange of $`\mathrm{\Delta }\widehat{\theta }_y`$ with $`\overline{\theta }_y`$. When $`\overline{\theta }_y=0`$, i.e. $`\widehat{\theta }_{y1}=\widehat{\theta }_{y2}`$, we obviously obtain $`\chi (0,\mathrm{\Delta }\widehat{\theta }_y)=1`$ that corresponds to complete coherence. In Fig. 25 we plot the three-dimensional representation of $`\chi (\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_y)`$. In order to get a feeling for the behavior of $`\chi `$ we also plot, in Fig. 26 and Fig. 27, two cuts of Fig. 25 illustrating, respectively, the behavior of $`\chi `$ for a fixed $`\overline{\theta }_y=0.5`$ (or fixed $`\mathrm{\Delta }\widehat{\theta }_y=0.5`$) and at $`\overline{\theta }_y=\mathrm{\Delta }\widehat{\theta }_y`$.
As it is evident from Fig. 25, $`\chi `$ exhibits, for any fixed value of $`\mathrm{\Delta }\widehat{\theta }_y`$, many different zeros in $`\overline{\theta }_y`$. In Fig. 28 we illustrate some of these zeros as a function of $`\mathrm{\Delta }\widehat{\theta }_y`$, $`\overline{\theta }_{y,Z}(\mathrm{\Delta }\widehat{\theta }_y)`$. The interest of this plot is that, once a certain distance $`\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y`$ between two pinholes is fixed, it illustrates at what position of the pinhole system, $`\overline{\theta }_{y,Z}`$, the spectral degree of coherence drops from unity to zero for the first time.
It is interesting to compare Fig. 28, with the directivity diagram of the radiant intensity $`I_S(\overline{\theta }_y)`$. This comparison is shown in Fig. 29. In the limit for $`\widehat{N}_y1`$ and $`\widehat{D}_y1`$, one may increase the degree of coherence of the beam by spatially filtering the radiation in the far field. If a vertical slit is used with aperture $`d`$ much larger than the horizontal coherence length, i.e. $`d\widehat{\xi }_{cx}`$, one would have poor coherence. Decreasing the aperture of the slit will increase the coherence of the X-ray beam up to some value $`d`$ smaller than $`\widehat{\xi }_c`$. Within our assumption $`\widehat{N}_x1`$ one has the far field approximation $`\widehat{\xi }_{cx}=(\pi /\widehat{N}_x)^{1/2}\widehat{z}_o`$. When $`d`$ becomes smaller and smaller with respect to $`(\pi /\widehat{N}_x)^{1/2}\widehat{z}_o`$ the spectral degree of coherence $`g`$ can be identified with the universal function $`\chi `$, as once can see by inspecting Eq. (266). As a result, as $`d`$ becomes smaller one loses photons, but the X-ray beam transverse coherence ceases to improve because, as is seen in Fig. 29, the transverse degree of coherence $`g=\chi `$ drops to zero along the vertical radiation pattern of the filtered X-ray beam: for instance, from Fig. 29 one can see that $`\chi `$ drops to zero for the first time at $`\mathrm{\Delta }\widehat{\theta }1`$ $`\overline{\theta }_y2`$, where the X-ray flux is still intense. This behavior of the degree of coherence should be taken into account at the stage of planning experiments. To give an example, after spatial filtering, one may conduct a two-pinhole experiment (like the one illustrated in Fig. 7) and find, surprisingly, that for some vertical position $`\overline{\theta }_y`$ of the pinholes (at fixed $`\mathrm{\Delta }\widehat{\theta }_y`$) well within the radiation pattern diagram he will have no fringes, but for some other vertical position he can find perfect visibility. So, without the knowledge of the function $`\chi `$ a user would not even have the possibility to predict the outcomes of a simple two-pinhole experiment.
From the definitions of $`\chi `$, $`\beta `$, $`\gamma `$ and $`I_S`$ it can be seen that all universal functions introduced in this work are partial cases of the more generic
$`M(\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right].`$ (270)
In fact
$`\chi (\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle \frac{M(\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_y)}{\left[M(\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y,0)\right]^{1/2}\left[M(\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y,0)\right]^{1/2}}}`$ (271)
$`I_S(\overline{\theta }_y)=M(\overline{\theta }_y,0)`$ (272)
$`\beta (\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\xi M(\xi ,\mathrm{\Delta }\widehat{\theta }_y)`$ (273)
$`\gamma (x)={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\xi \mathrm{exp}[i(2x)\xi ]M(\xi ,0).`$ (274)
The knowledge of $`M`$ is all one needs to calculate coherence properties out of many experimental setups, in very practical situations. It is therefore worth to tabulate $`M`$. We present a 3D plot of $`M`$, obtained from such tabulation, in Fig. 30.
Finally, it is interesting to sum up and compare results for the far field region obtained in this Section (non-homogeneous undulator source) with results obtained in Section 5. Many users performing coherent experiments with X-ray beams are interested in the beam coherence properties in the far field. We have seen that in the most general situation for third generation light sources one is interested in the case $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$, which guarantees factorization of results in the $`x`$ and $`y`$ direction, with arbitrary $`\widehat{N}_y`$ and $`\widehat{D}_y`$. In this case we will not have, in general, weakly quasi-homogeneous radiation in the vertical direction. The spectral degree of coherence can be found by simplifying Eq. (128) in the mathematical limit $`\widehat{z}_o\mathrm{}`$:
$`g(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (275)
$`\times \mathrm{exp}\left[2\widehat{N}_x\mathrm{\Delta }\widehat{\theta }_x^2\right]\mathrm{exp}\left[2\widehat{N}_y\mathrm{\Delta }\widehat{\theta }_y^2\right]{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[{\displaystyle \frac{(\overline{\theta }_y+\widehat{\varphi }_y)^2}{2\widehat{D}_y}}\right]M(\widehat{\varphi }_y,\mathrm{\Delta }\widehat{\theta }_y)`$ (276)
$`\times \left\{{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[{\displaystyle \frac{(\widehat{\varphi }_y+\overline{\theta }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{2\widehat{D}_y}}\right]I_S(\widehat{\varphi }_y)\right\}^{1/2}`$ (277)
$`\times \left\{{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}\left[{\displaystyle \frac{(\widehat{\varphi }_y+\overline{\theta }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{2\widehat{D}_y}}\right]I_S(\widehat{\varphi }_y)\right\}^{1/2}.`$ (278)
We can see that for any value of $`\widehat{N}_y`$ and $`\widehat{D}_y`$, in the far field limit we obtain a contribution to the cross-spectral density for the $`x`$ and for the $`y`$ direction. The contribution for the $`y`$ direction can be expressed in terms of the product of an exponential function and convolutions between the (Gaussian) electron beam divergence and universal functions.
On the one hand, as $`\widehat{D}_y1`$ we have
$`g(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (279)
$`\times \mathrm{exp}\left[2\widehat{N}_x\mathrm{\Delta }\widehat{\theta }_x^2\right]\mathrm{exp}\left[2\widehat{N}_y\mathrm{\Delta }\widehat{\theta }_y^2\right]\chi (\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_y).`$ (280)
On the other hand, as $`\widehat{D}_y1`$ we have weakly quasi-homogeneous wavefronts and
$`g(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (281)
$`\times \mathrm{exp}\left[2\widehat{N}_x\mathrm{\Delta }\widehat{\theta }_x^2\right]\mathrm{exp}\left[2\widehat{N}_y\mathrm{\Delta }\widehat{\theta }_y^2\right]\beta (\mathrm{\Delta }\widehat{\theta }_y)`$ (282)
where $`\beta (\mathrm{\Delta }\widehat{\theta }_y)`$ is given in Eq. (191).
Moreover, as $`\widehat{N}_y1`$, and for arbitrary $`\widehat{D}_y`$ we have:
$`g(\widehat{z}_o,\overline{\theta }_x,\overline{\theta }_y,\mathrm{\Delta }\widehat{\theta }_x,\mathrm{\Delta }\widehat{\theta }_y)=\mathrm{exp}\left[i2\left(\overline{\theta }_x\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_x+\overline{\theta }_y\widehat{z}_o\mathrm{\Delta }\widehat{\theta }_y\right)\right]`$ (283)
$`\times \mathrm{exp}\left[2\widehat{N}_x\mathrm{\Delta }\widehat{\theta }_x^2\right]\mathrm{exp}\left[2\widehat{N}_y\mathrm{\Delta }\widehat{\theta }_y^2\right].`$ (284)
It should be noted that Eq. (284) is simply a consequence of the application of the VCZ theorem in both horizontal and vertical directions.
## 7 Application: Coherent X-ray beam expander scheme
In this Section we show how transverse coherence properties of an X-ray beam can be manipulated to obtain a larger coherent spot-size on a sample.
The idea of increasing the horizontal width of the coherence spot is based on the use of a downstream slit for selection of the transversely coherent fraction of undulator radiation. Imagine a slit very close to the exit of the undulator with an aperture $`d`$ comparable with the coherent length of the radiation at the exit of the undulator $`\widehat{\xi }_{cx}=\sqrt{\pi /\widehat{D}_x}`$, as illustrated in Fig. 31. The new radiation source after the slit is now coherent and characterized by a horizontal dimension of the light spot equal to $`\sqrt{\pi /\widehat{D}_x}`$. In the far field one can take advantage of the reciprocal width relations of Fourier transform pairs or, equivalently, the expression for the Fraunhoffer diffraction pattern from a slit, i.e. a $`\mathrm{sinc}`$ function, to calculate the magnitude of the coherence spot. There is, of course, some arbitrary convention to agree upon when it comes to the definition of the width of the $`\mathrm{sinc}`$ function but, numerical factors aside, this reasoning shows qualitatively that the coherence spot is of order $`\sqrt{\widehat{D}_x}\widehat{z}_o`$ which is $`\widehat{ϵ}_x`$ times larger than the spot size dimension in the case of free-space propagation, of order $`\widehat{z}_o/\sqrt{\widehat{N}_x}`$. The radiation beyond the slit must be then spectrally filtered by a monochromator (not shown in Fig. 31) to further narrow the spectral bandwidth. Here we assume that the radiation frequency $`\omega `$ is equal to the fundamental frequency $`\omega _o`$. The radiation beyond the slit is transversely coherent when the aperture $`d`$ is equal (at most) to the coherence length $`\xi _{cx}`$.
Let us present a numerical example illustrating the improvement of the horizontal coherence length obtained by slit application. Let us consider the case when the electron horizontal emittance is large $`\widehat{ϵ}_x1`$ and the vertical emittance is small $`\widehat{ϵ}_y1`$, with $`\widehat{N}_x\widehat{D}_x`$. Since $`\widehat{N}_x=\widehat{ϵ}_x\widehat{\beta }_x`$ and $`\widehat{D}_x=\widehat{ϵ}_x/\widehat{\beta }_x`$, this is the case, for instance when $`\widehat{\beta }_x10`$. Then, for $`\widehat{ϵ}_x=100`$ we have $`\widehat{N}_x=1000`$ and $`\widehat{D}_x=10`$. This particular numerical example has been considered in Section 4.1 to illustrate, in free space, the behavior of the coherence length as a function of the position along the beamline, given in Eq. (105). Suppose we install the slit at $`\widehat{z}_o=2`$. From Eq. (105) we have $`\widehat{\xi }_{cx\widehat{z}_o=2}\sqrt{\pi /\widehat{D}_x}0.6`$. At $`\widehat{z}_o=12`$ we have a situation close to the asymptotic behavior where the field diffracted by the slit can be treated in the Fraunhofer approximation and $`\widehat{\xi }_{cx\widehat{z}_o=12}\sqrt{\widehat{D}_x}\widehat{z}_o30`$. Comparison of the asymptotic behaviors after spatial filtering with respect to the free-space propagation case is given in Fig. 32. Note that only the asymptotic behaviors near the slit and at large distance $`\widehat{z}_o1`$ are plotted for the spatially filtered radiation. What is important is, in fact, the comparison between the coherent distance at large values of $`\widehat{z}_o`$ with and without spatial filtering in the near zone.
This is an example in which evolution of transverse coherence through the beam line plays an important role. In fact, the ability of spatially filter radiation by a slit requires the knowledge of the transverse coherence length variation along the beamline.
Let us calculate the number of coherent photons observed beyond the aperture. In the region of parameters where $`\widehat{N}_x1`$ and $`\widehat{D}_x1`$, the number of transversely coherent photons into the slit aperture $`d=\xi _{cx}`$ can be calculated as
$$(N_{\mathrm{ph}})_{\mathrm{coh}}=\frac{dN_{\mathrm{ph}}}{dx}\xi _{cx}.$$
(285)
In the near-zone limit the slit is positioned at a position down the beamline $`z_s\beta _x`$ so that, from Eq. (105) and Eq. (89) we have:
$$\frac{dN_{\mathrm{ph}}}{dx}=\frac{N_{\mathrm{ph}}}{\sqrt{2\pi \sigma _x^2}}$$
(286)
and
$$\xi _{cx}=\sqrt{\pi }\frac{\lambda }{2\pi \sigma _x^{}}.$$
(287)
Therefore we can write
$$\frac{dN_{\mathrm{ph}}}{dx}\xi _{cx}=\sqrt{\pi }\frac{\lambda N_{\mathrm{ph}}}{\sqrt{(2\pi )^3\sigma _x^2\sigma _x^{}^2}}.$$
(288)
At the opposite extreme, with the distance $`z_s`$ much larger than $`\beta _x`$ we find that, always from Eq. (105) and Eq. (89)
$$\frac{dN_{\mathrm{ph}}}{dx}=\frac{N_{\mathrm{ph}}}{\sqrt{2\pi z_s^2\sigma _x^{}^2}}$$
(289)
and
$$\xi _{cx}=\sqrt{\pi }\frac{\lambda z_s}{2\pi \sigma _x}.$$
(290)
Therefore we can write
$$\frac{dN_{\mathrm{ph}}}{dx}\xi _{cx}=\sqrt{\pi }\frac{\lambda N_{\mathrm{ph}}}{\sqrt{(2\pi )^3\sigma _x^2\sigma _x^{}^2}}.$$
(291)
Thus, the number of transversely coherent photons into the slit aperture $`d=\xi _{cx}`$ will be independent of the distance $`z_s`$. This means that operation of spatial filtering in the near field region, as proposed by us, will not diminish the number of coherent photons with respect to the usual practice in which spatial filtering to obtain coherent radiation is performed in the far field: its effect will only be that of increasing the horizontal dimension of the coherent spot size. This possibility to create transversely coherent radiation with large divergence in the horizontal direction is important for many experiments. The distance of the slit from the exit of the undulator sets a limit to achievable linear dimension of coherence area. If we build a coherent X-ray beam line, we want to have large linear dimension of coherence area at the specimen position. Therefore, in order to have a largest linear dimension of coherence area at the specimen position $`\xi _{cx}(z_o)\sigma _x^{}z_o`$ we must have a slit aperture of at most of the size $`d\lambda /(2\pi \sigma _x^{})`$ installed in the near zone at $`z_s\beta _x`$. In order not to loose coherent photons instead, the slit aperture must be at least of the size $`d\lambda /(2\pi \sigma _x^{})`$. The right compromise is thus a slit aperture of the size $`d\lambda /(2\pi \sigma _x^{})`$ installed in the near zone at $`z_s\beta _x`$.
## 8 Conclusions
Before this work, no satisfactory theory describing spatial coherence from undulator radiation sources has been built. In this paper we developed such a theory of transverse coherence dealing with X-ray beams, with particular attention to third generation light sources.
First we studied Synchrotron Radiation as a random statistical process using the language of Statistical Optics. Statistical Optics developed around Gaussian, stationary processes characterized by quasi-homogeneous sources; under these assumptions, the characterization of statistical properties of the process are greatly simplified and the van Cittert-Zernike theorem (or its generalized version) can be used in order to describe the X-ray beam partial coherence properties in the far field region. However, for Synchrotron Radiation, there is no a priori reason to hold these assumptions satisfied.
We showed that Synchrotron Radiation is a Gaussian random process. As a result, statistical properties of Synchrotron Radiation are described satisfactory by second-order field correlation functions. We used a frequency domain analysis to describe them from a mathematical viewpoint. This choice is very natural. In fact, up-to-date detectors are limited to about $`100`$ ps time resolution: therefore, in real-life experiments with third generation light sources, detectors are by no means able to resolve a single X-ray pulse in time domain and work, instead, by counting the number of photons at a certain frequency over an integration time longer than the radiation pulse. As a consequence of the frequency domain analysis we could study the spatial correlation for a given frequency content using the cross-spectral density of the system which, independently of the spectral correlation function, can be used to extract useful information even if the process is not stationary.
We gave an expression for the process cross-spectral density dependent on six dimensionless parameters. Subsequently we tuned parameters at perfect resonance, thus obtaining a simplified expression.
First we studied the limit of applicability of the quasi-homogeneous model from an analytical viewpoint, within the framework of simplifying assumptions, namely in the limit of small electron beam divergence and Fresnel number in the vertical direction, of large electron beam divergence and Fresnel number in the horizontal direction, and performing calculations for the cross-spectral density on the horizontal plane only. This simplified study allowed us to introduce the concept of weakly quasi-homogeneous radiation and virtual quasi-homogeneous source while discussing the applicability region of the VCZ theorem.
Second, we studied the effect of the vertical emittance on the cross-spectral density. This study led us to analyze both cases in which the sources are weakly quasi-homogeneous and cases when they are not quasi-homogeneous at all. In the limit for large horizontal beam divergence and Fresnel number, which is always satisfied for third generation light sources, we found that the spectral degree of coherence factorizes in the product of factors depending separately on the horizontal and on the vertical coordinates. In the far field limit the vertical part of the spectral degree of coherence can be expressed in terms of the product of an exponential function (which, alone, would simply satisfy the VCZ theorem) and convolutions between the electron beam divergence in the vertical direction and a universal functions, that we introduced in our work. The universality of such a function implied that even for zero vertical emittance we never have full coherence in the vertical direction. This unexpected result is due to the influence of the horizontal emittance on the vertical coherence properties of the photon beam. Because of this, the degree of coherence changes between zero and unity within the diffraction angle. We also studied the near field zone. When one is interested in the evolution of the degree of coherence along the beamline back up to the exit of the undulator the situation becomes much more complicated with respect to the far zone case, as the observation distance is one of the problem parameters. There are many more asymptotic situations which can be studied, and a large part of our paper is devoted to the calculation of these asymptotic situations. We provided approximate estimations for the vertical coherence length that are valid from the far zone and back, up to the exit of the undulator in the case when either the vertical Fresnel number or the vertical electron beam divergence are much larger than unity. These can be used at the stage of planning experiments.
It should be noted that, throughout this work, we did not discuss the accuracy of the approximation of small and large parameters. In order to do so, one needs to develop a perturbation theory for each asymptotic case studied here, which would considerably increase the size of this paper. As a result we leave this issue for future work.
Finally, we selected an application to show the power of our approach. We discussed how the transverse coherence properties of an X-ray beam can be manipulated to obtain a convenient coherent spot-size on the sample with the help of a simple vertical slit; this invention was predicted almost entirely on the basis of theoretical ideas of rather complex and abstract nature discussed in the previous parts of the paper.
## Appendix A: Random phasor sum
The field of thermal light can be regarded as a sum of a great many independent contributions. The complex envelope of polarized thermal light at fixed time and a fixed point in space is a sum of a very large number of complex phasors
$$E(\stackrel{}{r},t)=\underset{k=1}{\overset{N}{}}\alpha _ke^{i\psi _k},$$
(292)
where $`N`$ is the number of radiating atoms. Statistical properties of elementary phasors that are generally satisfied in thermal light problems of interest are as follows:
a) The amplitudes $`\alpha _k=Re(\alpha _k)`$ and the phases $`\psi _k`$ are statistically independent of each other and of the amplitudes and phases of all other elementary phases for different values of $`k`$.
b) The random variables $`\alpha _k`$ are identically distributed for all $`k`$ with mean value $`<\alpha >`$ and second moment $`<\alpha ^2>`$.
c) The phases $`\psi _k`$ are uniformly distributed over the interval $`(0,2\pi )`$.
The reader can find in GOOD that when assumptions from a) to c) are satisfied, the real ($`Re(E)`$) and imaginary ($`Im(E)`$) parts of the field are distributed in accordance with the Gaussian law in the limit for $`N\mathrm{}`$, so that
$`Re(E)=Im(E)=0,`$ (293)
$`[Re(E)]^2=[Im(E)]^2={\displaystyle \frac{\alpha ^2}{2}}N=\sigma ^2,`$ (294)
$`Re(E)Im(E)=0,`$ (295)
$`p(Re(E),Im(E))={\displaystyle \frac{1}{2\pi \sigma ^2}}\mathrm{exp}\left[{\displaystyle \frac{[Re(E)]^2+[Im(E)]^2}{2\sigma ^2}}\right],`$ (296)
where $`p(Re(E),Im(E))`$ is the joint probability density function.
In Section 2.1 we discussed statistical properties of Synchrotron Radiation and we were led to assumptions 1), 2) and 3) which are weaker than a), b) and c). Here we will demonstrate that assumptions from a) to c) can be relaxed to assumptions from 1) to 3) without changes in results. We will derive results valid when the amplitudes $`\alpha _k`$ are complex $`\alpha _k=|\alpha _k|\mathrm{exp}(i\varphi _k)`$.
After denoting with $`r`$ and $`i`$ the real and imaginary parts of the fields and after substituting notation $`Q`$ with $`\overline{Q}`$ we first demonstrate that $`\overline{r}=\overline{i}=0`$. We have straightforwardly
$`\overline{r}={\displaystyle \frac{1}{N}}{\displaystyle \underset{k=1}{\overset{N}{}}}\left(\overline{\alpha _k\mathrm{cos}\varphi _k}\overline{\mathrm{cos}\psi _k}\overline{\alpha _k\mathrm{sin}\varphi _k}\overline{\mathrm{sin}\psi _k}\right)=0,`$ (297)
$`\overline{i}={\displaystyle \frac{1}{N}}{\displaystyle \underset{k=1}{\overset{N}{}}}\left(\overline{\alpha _k\mathrm{cos}\varphi _k}\overline{\mathrm{sin}\psi _k}+\overline{\alpha _k\mathrm{sin}\varphi _k}\overline{\mathrm{cos}\psi _k}\right)=0,`$ (298)
because all averages over trigonometric functions are zero.
Second, we demonstrate that $`\overline{r^2}=\overline{i^2}=\overline{\alpha ^2}/2`$. Again, direct calculation shows
$`\overline{r^2}={\displaystyle \frac{1}{N}}{\displaystyle \underset{k,n=1}{\overset{N}{}}}(\overline{\alpha _k\alpha _n\mathrm{cos}\varphi _k\mathrm{cos}\varphi _n}\overline{\mathrm{cos}\psi _k\mathrm{cos}\psi _n}`$ (299)
$`+\overline{\alpha _k\alpha _n\mathrm{sin}\varphi _k\mathrm{sin}\varphi _n}\overline{\mathrm{sin}\psi _k\mathrm{sin}\psi _n}`$ (300)
$`\overline{\alpha _k\alpha _n\mathrm{cos}\varphi _k\mathrm{sin}\varphi _n}\overline{\mathrm{cos}\psi _k\mathrm{sin}\psi _n}`$ (301)
$`\overline{\alpha _k\alpha _n\mathrm{sin}\varphi _k\mathrm{cos}\varphi _n}\overline{\mathrm{sin}\psi _k\mathrm{sin}\psi _n})`$ (302)
$`={\displaystyle \frac{1}{2N}}{\displaystyle \underset{k=1}{\overset{N}{}}}\overline{\alpha _k^2\left(\mathrm{cos}^2\psi _k+\mathrm{sin}^2\psi _k\right)}={\displaystyle \frac{\overline{\alpha ^2}}{2}}.`$ (303)
Moreover it is easy to see that $`\overline{i^2}=\overline{r^2}`$.
Finally, we show that $`\overline{ri}=0`$. In fact
$`\overline{ri}={\displaystyle \frac{1}{N}}{\displaystyle \underset{k,n=1}{\overset{N}{}}}(\overline{\alpha _k\alpha _n\mathrm{cos}\varphi _k\mathrm{cos}\varphi _n}\overline{\mathrm{cos}\psi _k\mathrm{sin}\psi _n}`$ (304)
$`+\overline{\alpha _k\alpha _n\mathrm{cos}\varphi _k\mathrm{sin}\varphi _n}\overline{\mathrm{cos}\psi _k\mathrm{cos}\psi _n}`$ (305)
$`\overline{\alpha _k\alpha _n\mathrm{sin}\varphi _k\mathrm{cos}\varphi _n}\overline{\mathrm{sin}\psi _k\mathrm{sin}\psi _n}`$ (306)
$`\overline{\alpha _k\alpha _n\mathrm{sin}\varphi _k\mathrm{sin}\varphi _n}\overline{\mathrm{sin}\psi _k\mathrm{cos}\psi _n})`$ (307)
$`={\displaystyle \frac{1}{2N}}{\displaystyle \underset{k}{}}\overline{|\alpha _k|^2\left(\mathrm{cos}\varphi _k\mathrm{sin}\varphi _k\mathrm{sin}\varphi _k\mathrm{cos}\varphi _k\right)}=0.`$ (308)
As a result we have that real and imaginary parts have zero means, equal variances and are uncorrelated. Use of the central limit theorem allows to conclude that the resulting phasor sum is a circular complex Gaussian random variable.
## Appendix B: A useful transformation of the expression for the undulator radiation field
We start reporting here, for convenience, Eq. (39), that represents the field (in normalized units) produced by a particle with offset and deflection at any distance $`\widehat{z}_o1/2`$ from the exit of the undulator, where the undulator center is taken at $`\widehat{z}_o=0`$:
$`\widehat{E}_s=\widehat{z}_o{\displaystyle _{1/2}^{1/2}}𝑑\widehat{z}^{}{\displaystyle \frac{1}{\widehat{z}_o\widehat{z}^{}}}\mathrm{exp}\left\{i\left[\left(\widehat{C}+{\displaystyle \frac{\stackrel{}{\widehat{\eta }}^2}{2}}\right)\widehat{z}^{}+{\displaystyle \frac{\left(\stackrel{}{\widehat{r}}_o\stackrel{}{\widehat{l}}\stackrel{}{\widehat{\eta }}\widehat{z}^{}\right)^2}{2(\widehat{z}_o\widehat{z}^{})}}\right]\right\}.`$ (309)
In this Appendix we show that $`\widehat{E}_s`$ as reported in Eq. (309) may be described as
$`\widehat{E}_s={\displaystyle _{1/2}^{1/2}}{\displaystyle \frac{\widehat{z}_od\widehat{z}^{}}{\widehat{z}_o\widehat{z}^{}}}\mathrm{exp}\left\{i\left[\mathrm{\Phi }_U+\widehat{C}\widehat{z}^{}+{\displaystyle \frac{\widehat{z}_o\widehat{z}^{}}{2(\widehat{z}_o\widehat{z}^{})}}\left(\stackrel{}{\widehat{\theta }}{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }}\right)^2\right]\right\}`$ (310)
where $`\mathrm{\Phi }_U`$ is given by
$$\mathrm{\Phi }_U=\left[\left(\widehat{\theta }_x\frac{\widehat{l}_x}{\widehat{z}_o}\right)^2+\left(\widehat{\theta }_y\frac{\widehat{l}_y}{\widehat{z}_o}\right)^2\right]\frac{\widehat{z}_o}{2}.$$
(311)
This means that $`\widehat{E}_s`$ is of the form
$$\widehat{E}_s(\widehat{C},\widehat{z}_o,\stackrel{}{\widehat{\theta }}(\stackrel{}{\widehat{l}}/\widehat{z}_o)\stackrel{}{\widehat{\eta }})=\mathrm{exp}(i\mathrm{\Phi }_U)S(\widehat{C},\widehat{z}_o,\stackrel{}{\widehat{\theta }}(\stackrel{}{\widehat{l}}/\widehat{z}_o)\stackrel{}{\widehat{\eta }}).$$
(312)
Let us introduce in this Appendix, for simplicity of notation, $`\stackrel{}{\widehat{\xi }}=\stackrel{}{\widehat{r}}_o\stackrel{}{\widehat{l}}`$ and $`\stackrel{}{\widehat{\varphi }}=\stackrel{}{\xi }/\widehat{z}_o`$. It is easy to rewrite the phase in the integrand of Eq. (309), which we denote with $`\mathrm{\Phi }_T`$ as
$`\mathrm{\Phi }_T=\widehat{C}\widehat{z}^{}+{\displaystyle \frac{\stackrel{}{\widehat{\eta }}^2}{2}}\widehat{z}^{}+{\displaystyle \frac{1}{2}}\left\{\left[\stackrel{}{\widehat{\varphi }}^2\widehat{z}_o2\stackrel{}{\widehat{\varphi }}\stackrel{}{\widehat{\eta }}\widehat{z}^{}+{\displaystyle \frac{\stackrel{}{\widehat{\eta }}^2\widehat{z}^2}{\widehat{z}_o}}\right]\left[1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\widehat{z}^{}}{\widehat{z}_o}}\right)^n\right]\right\}.`$ (313)
Further algebraic manipulation of Eq. (313) yields
$`\mathrm{\Phi }_T=\widehat{C}\widehat{z}^{}+{\displaystyle \frac{\stackrel{}{\widehat{\varphi }}^2\widehat{z}_o}{2}}+{\displaystyle \frac{\widehat{z}^{}}{2}}(\stackrel{}{\widehat{\varphi }}\stackrel{}{\widehat{\eta }})^2+{\displaystyle \frac{1}{2}}\{{\displaystyle \frac{\stackrel{}{\widehat{\eta }}^2\widehat{z}^2}{\widehat{z}_o}}+{\displaystyle \frac{\stackrel{}{\widehat{\eta }}^2\widehat{z}^2}{\widehat{z}_o}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\widehat{z}^{}}{\widehat{z}_o}}\right)^n`$ (314)
$`+\stackrel{}{\widehat{\varphi }}^2\widehat{z}_o{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\widehat{z}^{}}{\widehat{z}_o}}\right)^n2\stackrel{}{\widehat{\varphi }}\stackrel{}{\widehat{\eta }}\widehat{z}^{}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\widehat{z}^{}}{\widehat{z}_o}}\right)^n\}.`$ (315)
It is easy to recognize $`\mathrm{\Phi }_U`$ in the second term on the right hand side of Eq. (315). Furthermore, the fourth term on the right hand side of Eq. (315) can be further manipulated, leading to
$`\mathrm{\Phi }_T=\widehat{C}\widehat{z}^{}+\mathrm{\Phi }_U+{\displaystyle \frac{\widehat{z}^{}}{2}}\left(\stackrel{}{\widehat{\varphi }}\stackrel{}{\widehat{\eta }}\right)^2`$ (316)
$`+{\displaystyle \frac{1}{2}}\left\{\widehat{z}_o\left(\stackrel{}{\widehat{\varphi }}^2+\stackrel{}{\widehat{\eta }}^2\right){\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\widehat{z}^{}}{\widehat{z}_o}}\right)^n2\stackrel{}{\widehat{\varphi }}\stackrel{}{\widehat{\eta }}\widehat{z}_o{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\widehat{z}^{}}{\widehat{z}_o}}\right)^n\right\},`$ (317)
that is
$`\mathrm{\Phi }_T=\widehat{C}\widehat{z}^{}+\mathrm{\Phi }_U+{\displaystyle \frac{\widehat{z}^{}}{2}}\left(\stackrel{}{\widehat{\varphi }}\stackrel{}{\widehat{\eta }}\right)^2+{\displaystyle \frac{\widehat{z}_o}{2}}\left\{\left(\stackrel{}{\widehat{\varphi }}\stackrel{}{\widehat{\eta }}\right)^2{\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\widehat{z}^{}}{\widehat{z}_o}}\right)^n\right\}`$ (318)
or
$`\mathrm{\Phi }_T=\widehat{C}\widehat{z}^{}+\mathrm{\Phi }_U+{\displaystyle \frac{\widehat{z}_o\widehat{z}^{}}{2(\widehat{z}_o\widehat{z}^{})}}\left(\stackrel{}{\widehat{\varphi }}\stackrel{}{\widehat{\eta }}\right)^2`$ (319)
Therefore, since
$$\stackrel{}{\widehat{\varphi }}\stackrel{}{\widehat{\eta }}=\stackrel{}{\widehat{\theta }}\frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}\stackrel{}{\widehat{\eta }}$$
(320)
we have
$`\widehat{E}_s={\displaystyle _{1/2}^{1/2}}{\displaystyle \frac{\widehat{z}_od\widehat{z}^{}}{\widehat{z}_o\widehat{z}^{}}}\mathrm{exp}\left\{i\left[\mathrm{\Phi }_U+\widehat{C}\widehat{z}^{}+{\displaystyle \frac{\widehat{z}_o\widehat{z}^{}}{2(\widehat{z}_o\widehat{z}^{})}}\left(\stackrel{}{\widehat{\theta }}{\displaystyle \frac{\stackrel{}{\widehat{l}}}{\widehat{z}_o}}\stackrel{}{\widehat{\eta }}\right)^2\right]\right\}`$ (321)
that is Eq. (310), quantum erat demonstrandum.
## Appendix C: Autocorrelation function for undulator sources
In this Appendix we demonstrate the validity of Eq. (188). Having defined
$`\stackrel{~}{f}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_xS^{}[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2]S[\widehat{z}_o,\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2],`$ (322)
and
$`\beta (\mathrm{\Delta }\widehat{\theta }_y)={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right]\mathrm{sinc}\left[{\displaystyle \frac{\widehat{\varphi }_x^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y)^2}{4}}\right],`$ (323)
we have want to demonstrate that Eq. (188) holds, that is
$$\stackrel{~}{f}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_y)=\beta (\mathrm{\Delta }\widehat{\theta }_y).$$
(324)
The proof is based on the autocorrelation theorem, which states that if the (two-dimensional) Fourier Transform of a function $`w(x,y)`$ with respect to variables $`\alpha _x`$ and $`\alpha _y`$ is indicated by $`\overline{w}(\alpha _x,\alpha _y)`$, then the Fourier transform of the two-dimensional autocorrelation function of $`w(x,y)`$ with respect to the same variables $`\alpha _x`$ and $`\alpha _y`$ is given by $`|\overline{w}(\alpha _x,\alpha _y)|^2`$. In formulas, after definition of the autocorrelation function
$$𝒜[w](x,y)=_{\mathrm{}}^{\mathrm{}}𝑑\eta _{\mathrm{}}^{\mathrm{}}𝑑\xi w(\eta +x,\xi +y)w^{}(\eta ,\xi ),$$
(325)
which is equivalent to
$$𝒜[w](x,y)=_{\mathrm{}}^{\mathrm{}}𝑑\eta _{\mathrm{}}^{\mathrm{}}𝑑\xi w(\eta +x/2,\xi +y/2)w^{}(\eta x/2,\xi y/2),$$
(326)
the autocorrelation theorem states that
$$_{\mathrm{}}^{\mathrm{}}𝑑x_{\mathrm{}}^{\mathrm{}}𝑑y\mathrm{exp}[i(\alpha _xx+\alpha _yy)]𝒜[w](x,y)=|\overline{w}(\alpha _x,\alpha _y)|^2.$$
(327)
First we extend the definition of $`\stackrel{~}{f}`$
$`\stackrel{~}{f}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x^{},\mathrm{\Delta }\widehat{\theta }_y^{})={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_xS^{}[\widehat{z}_o,(\widehat{\varphi }_x\mathrm{\Delta }\widehat{\theta }_x^{}/2)^2+(\widehat{\varphi }_y\mathrm{\Delta }\widehat{\theta }_y^{}/2)^2]`$ (328)
$`\times S[\widehat{z}_o,(\widehat{\varphi }_x+\mathrm{\Delta }\widehat{\theta }_x^{}/2)^2+(\widehat{\varphi }_y+\mathrm{\Delta }\widehat{\theta }_y^{}/2)^2],`$ (329)
where we changed variables from $`\mathrm{\Delta }\widehat{\theta }_{x,y}`$ to $`\mathrm{\Delta }\widehat{\theta }_{x,y}^{}=2\mathrm{\Delta }\widehat{\theta }_{x,y}`$ Then we can apply the autocorrelation theorem in Eq. (327) to the function $`\stackrel{~}{f}`$ thus obtaining the following relation:
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mathrm{\Delta }\widehat{\theta }_x^{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mathrm{\Delta }\widehat{\theta }_y^{}\mathrm{exp}[i(\alpha _x\mathrm{\Delta }\widehat{\theta }_x^{}+\alpha _y\mathrm{\Delta }\widehat{\theta }_y^{})]\stackrel{~}{f}(\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x^{},\mathrm{\Delta }\widehat{\theta }_y^{})`$ (330)
$`={\displaystyle \frac{1}{2\pi ^2}}\left|{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}[i(\alpha _x\widehat{\varphi }_x+\alpha _y\widehat{\varphi }_y)]S[\widehat{z}_o,\widehat{\varphi }_x^2+\widehat{\varphi }_y^2]\right|^2,`$ (331)
where $`\alpha _{x,y}`$ are now conjugated variables with respect to the angles $`\widehat{\varphi }_{x,y}`$ on which $`S`$ depends. We will denote with $`\overline{S}`$ the two-dimensional Fourier Transform of $`S`$, that is:
$`\overline{S}(\alpha _x,\alpha _y)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}[i(\alpha _x\widehat{\varphi }_x+\alpha _y\widehat{\varphi }_y)]S[\widehat{z}_o,\widehat{\varphi }_x^2+\widehat{\varphi }_y^2].`$ (332)
The relation between the function $`S`$ and the undulator field is given by Eq. (51), and one has
$`\overline{S}(\alpha _x,\alpha _y)={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_x{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{\varphi }_y\mathrm{exp}[i(\alpha _x\widehat{\varphi }_x+\alpha _y\widehat{\varphi }_y)]\mathrm{exp}[i(\widehat{\varphi }_x^2+\widehat{\varphi }_y^2)\widehat{z}_o/2]`$ (333)
$`\times \widehat{E}_s[\widehat{z}_o,\widehat{\varphi }_x^2+\widehat{\varphi }_y^2].`$ (334)
After definition of $`\overline{\alpha }_{x,y}=\alpha _{x,y}/\widehat{z}_o`$ one can write $`\overline{S}`$, as a function of $`\overline{\alpha }_{x,y}`$ instead of $`\alpha _{x,y}`$. Then, one can switch to the new integration variables $`\widehat{x}=\widehat{\varphi }_x\widehat{z}_o`$ and $`\widehat{y}=\widehat{\varphi }_y\widehat{z}_o`$ to obtain:
$`\overline{S}(\overline{\alpha }_x,\overline{\alpha }_y)={\displaystyle \frac{1}{\widehat{z}_o^2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{x}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{y}\mathrm{exp}[i(\overline{\alpha }_x\widehat{x}+\overline{\alpha }_y\widehat{y})]\mathrm{exp}[i(\widehat{x}^2+\widehat{y}^2)/(2\widehat{z}_o)]`$ (335)
$`\times \widehat{E}_s[\widehat{z}_o,\widehat{x}^2+\widehat{y}^2].`$ (336)
(337)
where the expression for $`\widehat{E}_s[\widehat{z}_o,\widehat{x}^2+\widehat{y}^2]`$ is given in Eq. (39). Now we have to calculate the Fourier transform of the product of two factors: $`\mathrm{exp}[i(\widehat{x}^2+\widehat{y}^2)/(2\widehat{z}_o)]`$ and $`\widehat{E}_s`$. Let us look for the Fourier transform of each factor.
A direct calculation shows that
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{x}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{y}\mathrm{exp}[i(\overline{\alpha }_x\widehat{x}+\overline{\alpha }_y\widehat{y})]\widehat{E}_s(\widehat{z}_o,\widehat{x}^2+\widehat{y}^2)`$ (338)
$`=2i\pi \widehat{z}_o\mathrm{exp}\left[i{\displaystyle \frac{(\overline{\alpha }_x^2+\overline{\alpha }_y^2)\widehat{z}_o}{2}}\right]\mathrm{sinc}\left[{\displaystyle \frac{\overline{\alpha }_x^2+\overline{\alpha }_y^2}{4}}\right].`$ (339)
Second, let us deal with the Fourier transform of $`\mathrm{exp}[i(\widehat{x}^2+\widehat{y}^2)/(2\widehat{z}_o)]`$:
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{x}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\widehat{y}\mathrm{exp}[i(\overline{\alpha }_x\widehat{x}+\overline{\alpha }_y\widehat{y})]\mathrm{exp}[i(\widehat{x}^2+\widehat{y}^2)/(2\widehat{z}_o)]`$ (340)
$`=4i\widehat{z}_o\mathrm{exp}[i(\overline{\alpha }_x^2+\overline{\alpha }_y^2)\widehat{z}_o/2].`$ (341)
Since the Fourier transform of a product is equal to the convolution of the Fourier transforms of the factors we have
$`\overline{S}(\overline{\alpha }_x,\overline{\alpha }_y)=8\pi {\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑u{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑w\mathrm{exp}\{i[(\overline{\alpha }_xu)^2+(\overline{\alpha }_yw)^2]\widehat{z}_o/2\}`$ (342)
$`\times \mathrm{exp}[i(u^2+w^2)\widehat{z}_o/2]\mathrm{sinc}[(u^2+w^2)/4].`$ (343)
and therefore we have
$`|\overline{S}(\overline{\alpha }_x,\overline{\alpha }_y)|=8\pi \left|{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑u{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑w\mathrm{exp}[i(\overline{\alpha }_xu+\overline{\alpha }_yw)\widehat{z}_o]\mathrm{sinc}[(u^2+w^2)/4]\right|.`$ (344)
Going back to old variables $`\alpha _{x,y}`$ we obtain
$`|\overline{S}(\alpha _x,\alpha _y)|=8\pi \left|{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑u{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑w\mathrm{exp}[i(\alpha _xu+\alpha _yw)]\mathrm{sinc}[(u^2+w^2)/4]\right|.`$ (345)
which is independent of $`\widehat{z}_o`$. As a result Eq. (331) is also independent on $`\widehat{z}_o`$, i.e. the Fourier transform of $`\stackrel{~}{f}`$ is independent of $`\widehat{z}_o`$.
Now, on the one hand in the limit for $`\widehat{z}_o\mathrm{}`$ the function $`\stackrel{~}{f}`$ transforms into $`\beta `$, because the $`S`$ functions in $`\stackrel{~}{f}`$ tend asymptotically to the $`\mathrm{sinc}`$ functions in $`\beta `$. On the other hand, if the Fourier transform of $`\stackrel{~}{f}`$ is independent of $`\widehat{z}_o`$, also $`\stackrel{~}{f}`$ is independent of $`\widehat{z}_o`$. As a result it can only be $`\stackrel{~}{f}(\mathrm{\Delta }\widehat{\theta }_y^{})=\beta (\mathrm{\Delta }\widehat{\theta }_y^{})`$, and $`\stackrel{~}{f}(\mathrm{\Delta }\widehat{\theta }_y)=\beta (\mathrm{\Delta }\widehat{\theta }_y)`$ that is Eq. (324) holds, quantum erat demonstrandum.
Note that, based on the autocorrelation theorem, it is also possible to give an analytic expression for the Fourier transform of $`\beta `$. After definition of $`\beta =\beta (\widehat{z}_o,\mathrm{\Delta }\widehat{\theta }_x^{},\mathrm{\Delta }\widehat{\theta }_y^{})`$ as in Eq. (329), application of the autocorrelation theorem simply states that the two-dimensional Fourier Transform of $`\beta `$, that will be indicated with $`\overline{\beta }`$ can be written as
$$\overline{\beta }(\alpha _x,\alpha _y)=\frac{1}{2\pi ^2}\left|_{\mathrm{}}^{\mathrm{}}𝑑\widehat{\varphi }_x_{\mathrm{}}^{\mathrm{}}𝑑\widehat{\varphi }_y\mathrm{exp}[i(\alpha _x\widehat{\varphi }_x+\alpha _y\widehat{\varphi }_y)]\mathrm{sinc}\left(\frac{\widehat{\varphi }_x^2+\widehat{\varphi }_y^2}{4}\right)\right|^2.$$
(346)
Introducing $`\alpha ^2=\alpha _x^2+\alpha _y^2`$ and representing the two-dimensional Fourier Transform of $`\beta `$ in terms of Fourier-Bessel transform we obtain
$$\overline{\beta }(\alpha )=\frac{1}{2\pi ^2}\left|2\pi _0^{\mathrm{}}𝑑rrJ_0(r\alpha )\mathrm{sinc}\left(\frac{r^2}{4}\right)\right|^2=2\left[\pi 2\mathrm{S}\mathrm{i}(\alpha ^2)\right]^2,$$
(347)
where $`\mathrm{Si}`$ indicates the sine integral function.
Finally, one can get back a simpler representation of the function $`\beta `$ in terms of a one-dimensional integration simply performing an anti Fourier-Bessel transform:
$$\beta (\mathrm{\Delta }\widehat{\theta }^{})=\frac{1}{2\pi }_0^{\mathrm{}}𝑑\alpha \alpha J_0(\alpha \mathrm{\Delta }\widehat{\theta }^{})\overline{\beta }(\alpha ),$$
(348)
where $`\mathrm{\Delta }\widehat{\theta }^{}_{}{}^{}2=\mathrm{\Delta }\widehat{\theta }_x^{}_{}{}^{}2+\mathrm{\Delta }\widehat{\theta }_y^{}_{}{}^{}2`$. For $`\mathrm{\Delta }\widehat{\theta }_x=0`$ we obtain
$$\beta (\mathrm{\Delta }\widehat{\theta }_y)=\frac{1}{\pi }_0^{\mathrm{}}𝑑\alpha \alpha J_0\left(\alpha \frac{\mathrm{\Delta }\widehat{\theta }_y}{2}\right)\left[\pi 2\mathrm{S}\mathrm{i}(\alpha ^2)\right]^2.$$
(349)
## Acknowledgements
The authors wish to thank Hermann Franz, Petr Ilinski and Ivan Vartanyants for many useful discussions, Jochen Schneider and Edgar Weckert for their interest in this work.
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# Finite-curvature scaling in optical lattice systems
## Acknowledgement
We thank M. Foulkes, N. I. Gidopoulos and A. J. Schofield for useful discussions. JQ thanks the University of Birmingham for hospitality while some of this work was carried out. We acknowledge financial support from EPSRC (CH), the Leverhulme Trust and CCLRC (JQ).
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# Microwave emission from a crystal of molecular magnets – The role of a resonant cavity
## I Introduction
Magnetic complex molecules have attracted a great deal of attention in recent years, because they have remarkable properties, related to their high magnetic anisotropy and large value of spinBarbara et al. (1999); Wernsdorfer and Sessoli (1999); Villain (2003): for Mn<sub>12</sub>Ac as well as for Fe<sub>8</sub>O, the quantum number of the total spin is $`S=10.`$ Accordingly, the eigenvalues of $`S_z`$, (the spin component in the direction of the easy axis), can take $`21`$ different values: $`m=10,\mathrm{},10.`$ One of the most interesting properties of these molecules is that in a slowly changing external magnetic field the magnetization of the crystal consisting of such molecules exhibit series of steps at sufficiently low temperatures.Friedman et al. (1996) The effect can be explained by assuming that the energy levels of the molecules become doubly degenerate at the corresponding values of the magnetic field and this resonance condition increases the possibility of the transition between the degenerate states with different values of $`m`$ and $`m^{}`$. This kind of quantum tunneling between spin levels leads to a sudden change in the magnetic moment of the crystal and is therefore of fundamental importance as being a macroscopically observable quantum effect.
In an important theoretical work Chudnovsky and GaraninChudnovsky and Garanin (2002) proposed that resonant magnetic tunneling could be accompanied by the emission of electromagnetic radiation, the possibility of superradiance from magnetic molecules has been considered further in Ref. \[Henner and Kaganov, 2003\] and bursts of microwave pulses have actually been detected in recent experiments.Tejada et al. (2004); Vanacken et al. (2004); Hernandez-Minguez et al. (2005) According to Refs. \[Chudnovsky and Garanin, 2002; Henner and Kaganov, 2003; Joseph et al., 2004\], the possible physical mechanism responsible for this phenomenon can be superradiance (SR)Dicke (1954); Benedict et al. (1996); Gross and Haroche (1982) which is an interesting collective effect predicted first by Dicke in 1954, and has been experimentally observed in several physical systems since then.Benedict et al. (1996) However, when the radiation emitted by magnetic molecules was detected, the sample was placed in a container, which acted as a waveguide. This cavity changes the mode structure of the electromagnetic field surrounding the sample, which is known to have crucial effects on the dynamics of the emitted radiation. Studies in SR with other physical systems like the ensemble of proton spins Benedict et al. (1996); Kiselev et al. (1988); Bazhanov et al. (1990); Yukalov and Yukalova (2004) in the MHz, and with Rydberg atoms in the GHz domainHaroche and Raimond (1985) show that the presence of a resonant cavity may enhance the collectivity of the radiating individual dipoles, as first proposed by Bloembergen and Pound Bloembergen and Pound (1954), and which seems to be necessary to obtain radiation in the case of molecular magnets, as well.Yukalov and Yukalova (2005a, b) Additionally, it has also been demonstrated that external resonators such as Fabry-Perot mirrors can enhance the relaxation of a crystal of molecular nanomagnets.Tejada et al. (2003) Inspired by these facts, we investigate in this paper the interplay between the radiation and the changes of the magnetization of a macroscopic sample of molecules Mn<sub>12</sub>-Ac inside a nearly resonant cavity. We also note that the interesting proposal of Ref. \[Leuenberger and Loss, 2001\] to use these molecules for implementing a quantum algorithm gives another motivation to study their radiative properties.
The present paper is organized as follows: In Sec. II we investigate the relevant magnetic level structure and describe a method that allows us to reduce the problem to a set of level pairs. The interaction of the molecules with the cavity field is considered in Sec. III. In Sec. IV we discuss the approximate analytical consequences of our model and present numerical results as well. Finally we summarize and draw the conclusions (Sec. V).
## II Magnetic level structure
Experiments including magnetization measurements, Friedman et al. (1996); Mertes et al. (2001) neutronMirebeau et al. (1999) and EPRBarra et al. (1997); Hill et al. (1998, 2003) studies on crystals of Mn<sub>12</sub>Ac and Fe<sub>8</sub>O suggest that the Hamiltonian responsible for the magnetic properties can be written as a sum of two terms:
$$H_S=H_0+H_1.$$
(1)
Here $`H_0`$ is diagonal in the eigenbasis $`\{|m\}`$ of the (dimensionless) $`z`$ component of the spin operator, $`S_z`$:
$$H_0=DS_z^2FS_z^4\stackrel{~}{\mu }B_0S_z,$$
(2)
where the last term describes the coupling to an external magnetic field applied in the $`z`$ direction (easy axis): $`𝐁_0=(0,0,B_0)`$ with $`\stackrel{~}{\mu }=g\mu _B`$. On the other hand, $`H_1`$ consists of termsMertes et al. (2001); Barra et al. (1997) that do not commute with $`S_z`$:
$$H_1=C(S_+^4+S_{}^4)+E(S_+^2+S_{}^2)/2+K(S_++S_{})/2.$$
(3)
As most of the experiments where microwave radiation emitted by magnetic molecules was detected have been performed on Mn<sub>12</sub>-Ac, from now on, we shall consider this molecule as the representative example. In this case the values of the parameters in $`H_0`$ are $`D/k_B=0.56K,`$ and $`F/k_B=1.110^3K`$. The coefficients in $`H_1`$ do not have unanimously accepted values, but $`H_1`$ can be considered as a small correction to $`H_0`$. However, as the transitions between levels with different $`m`$ and $`m^{}`$ are induced by terms that do not commute with $`S_z`$, the importance of $`H_1`$ is fundamental from this point of view. Tetragonal symmetry would allow only the quartic term, but there is experimental evidenceMertes et al. (2001) showing the presence of weak quadratic and linear terms in $`H_1`$. We shall return to the determination of the coefficients in $`H_1`$ at the end of this section.
In a typical experimental situation the external field $`B_0`$ slowly changes in time, and consequently so does $`H_0`$. Considering the total Hamiltonian (1) as the generator of the time evolution (which means that relaxation effects are not taken into account), the corresponding time dependent Schrödinger equation governs the dynamics. However, it turns out that its solution is not feasible, because the time independent part of $`H_0`$ forces a much faster evolution than the slow variation due to the change of the magnetic field. The fundamental frequencies $`\omega _{mm^{}}(D/\mathrm{})(m^2m^2)`$ are in the range $`10^{10}10^{11}s^1`$ being very fast compared with the time dependence of the magnetic field $`B_0`$ that among ordinary circumstancesTejada et al. (2004) cause a change on the scale $`10s^1`$. Therefore there is a significant variation in the molecular state as a consequence of the first term in (2), while nothing happens due to $`\stackrel{~}{\mu }B_0S_z`$, which appears in the last term of $`H_0`$. On the other hand, we know that an appreciable change in the state occurs only if two energy levels become close to each other, as stipulated by a simple time dependent perturbation calculation, where the energy difference between the levels appears in the denominator of the transition probability. As the dominant term in the Hamiltonian (1) is $`H_0`$, we can approximately find the points where two energy levels are close to each other by calculating the eigenvalues of $`H_0`$, which are obtained by the mere substitution $`S_zm`$. Simple algebra shows that these eigenvalues become doubly degenerate with given $`m`$ and $`m^{}`$ at the following values of $`B_0`$:
$$\stackrel{~}{\mu }B_0=D(m+m^{})\left(1+\frac{F}{D}\left[m^2+m^2\right]\right).$$
(4)
A part of the level scheme of the total Hamiltonian $`H_S`$ is shown in Fig. 1 as function of $`B_0`$. The special values of $`B_0`$ given by Eq. (4) – where the levels of $`H_0`$ cross – can be clearly identified in this figure. However, as it is known, the presence of $`H_1`$ perturbs the eigenvalues leading to a splitting of the levels instead of crossing as shown by the inset. The resonance condition implies that appreciable changes in the population of the levels is expected around these avoided crossings (sometimes also called anticrossings). Thus the system can be efficiently approximated by a set of level pairs, each of which is to be considered as an effective two-level system, similarly to the figure shown in the inset.
Reducing the problem to a set of level pairs means technically that one applies degenerate perturbation theory around each avoided level crossing (determined by the values of $`B_0`$ in Eq. (4)) following a technique proposed independently by van Vleckvan Vleck (1929), Des Cloiseauxdes Cloizeaux (1960) and other authors as summarized in Ref. \[Klein, 1974\]. In the vicinity of a given avoided level crossing we find a unitary transformation that block diagonalizes $`H_S`$ on the zero order two-dimensional eigensubspaces and generate an effective Hamiltonian $`H_e`$ that has the same eigenvalues as $`H_S`$:
$$H_e=UH_SU^{}$$
(5)
with
$$U=\underset{a}{}(P_a^0P_aP_a^0)^{1/2}P_a^0P_a,$$
(6)
where $`a`$ labels the various eigensubspaces with quasi-degenerate eigenvalues, and $`P_a^0`$ are orthogonal projections on the degenerate eigensubspaces of $`H_0`$. For a given value of $`a,`$ corresponding to the avoided crossing of levels $`m`$ and $`m^{}`$, we have $`P^0=|mm|+|m^{}m^{}|.`$ As the perturbation is turned on, the $`\{|m,|m^{}\}`$ eigenstates of $`H_0`$ evolve into $`\{|\varphi _m,|\varphi _m^{}\}`$ and we consider $`P=|\varphi _m\varphi _m|+|\varphi _m^{}\varphi _m^{}|`$ which projects onto the space arising from the zero order subspace of interest. By appropriate expansionsGaranin (1991); Leuenberger and Loss (2000); Yoo and Park (2005) of this unitary operator $`U`$, one can obtain a perturbation series, leading to useful analytical approximations of $`H_e.`$ Alternatively, the numerically exact projections $`P_a`$ can be found by the diagonalization of $`H_S`$, and then Eqs. (5) and (6) provide the operator $`H_e`$. The latter method is followed in this paper and we obtain in each two-dimensional subspace spanned by $`\{|m,|m^{}\}`$ the following matrix for the effective Hamiltonian:
$$H_e(t)=\left(\begin{array}{cc}\epsilon _0+w/2& \mathrm{\Delta }_0/2\\ \mathrm{\Delta }_0/2& \epsilon _0w/2\end{array}\right)_{m,m^{}}$$
(7)
with time dependent elements, and of course the values depend also on the pair $`\{|m,|m^{}\}`$. Here $`\epsilon _0`$ is the energy where the given crossing would occur, $`w`$ is proportional to the time dependent external field in the $`z`$ direction, while the offdiagonal element $`\mathrm{\Delta }_0`$ is the level splitting responsible for the effective coupling between the levels. $`w`$ will be assumed to be linear in time with constant $`\dot{B}_0`$, yielding $`w(t)=\stackrel{~}{\mu }\dot{B}_0(tt_0)(mm^{})`$ with $`t_0`$ being the time instant when the crossing would occur. Note that this is a reasonable approximation even for a time scale much longer than the expected duration of the transition to be described (see e.g. Fig. 1. in Ref. \[Vanacken et al., 2004\]). This linear approximation corresponds to the usual Landau-Zener-Stückelberg (LZS) model Landau (1932); Zener (1932); Stückelberg (1932), by the aid of which we can calculate the probability of a given $`mm^{}`$ transition: $`P_{mm^{}}=1\mathrm{exp}(\pi \mathrm{\Delta }_0^2/2\mathrm{}w)`$. Note that in this expression both $`w`$ and $`\mathrm{\Delta }_0`$ depend on the labels $`m,m^{}`$. For a given pair of levels and sweep rate $`\dot{B}_0`$, $`P_{m,m^{}}`$ is determined by the magnitude of the level splitting $`\mathrm{\Delta }_0`$, i.e., essentially by the parameters in $`H_1`$. If we assume that initially the system is in thermal equilibrium, we can consider a series of transitions at the values of $`B_0`$ given by Eq. (4). Calculating the expectation value of the operator $`S_z`$ (which is proportional to the magnetization) we obtain a staircase-like hysteresis loop that can be compared with the experimental curvesMertes et al. (2001) at a given temperature and sweep rate. As our results depend on the coefficients in $`H_1`$, the minimization of the difference between the steps in the calculated hysteresis curve and the experimental plots gives the desired parameter values. We have found the best agreement for $`K=0.025\stackrel{~}{\mu }B_0`$, $`E/k_B=4.48`$ $`10^3`$ $`K`$, $`C/k_B=1.36`$ $`10^5`$ $`K`$, therefore these parameter values will be used in the following.
The method summarized in this section, first of all, gives us information about the magnitude of the terms in the spin Hamiltonian, and describes how to obtain the level splittings $`\mathrm{\Delta }_0`$. As we shall see in the next section, the coupling of the molecular system to the cavity field at a given $`mm^{}`$ transition can also be determined in this way.
## III Interaction with the resonant cavity field
In this section we describe the interaction of an ensemble of magnetic molecules with a quasi-resonant cavity field. The time dependence of the level structure, considered in the previous section, will bring a certain level pair into resonance with the cavity field at a given value of the external magnetic field $`𝐁_0.`$ Usually there are avoided level crossings before the resonance, where for low lying $`m`$ and $`m^{}`$ the LZS transition probabilities $`P_{m,m^{}}`$ are small, leading to an almost complete inversion. We shall consider a dynamical equation for the density operator of the molecules and assume that a dipole moment is generated during a certain transition between the split levels $`m`$ and $`m^{}`$ of a given molecule which in turn serves as a source of microwave radiation influencing the transitions in all other molecules. To take into account this radiation mediated interaction, we have to add a new term to the effective two-level Hamiltonian describing the interaction with the magnetic dipole field of the cavity, $`\stackrel{}{}`$, which is an additional field beyond the stronger but almost static magnetic field $`𝐁_0`$ creating the inversion between the levels. The fast dynamics of the magnetic dipoles will be a forced oscillation generated by the interaction with the external field that is characterized by the interaction Hamiltonian $`H_I=\stackrel{~}{\mu }\stackrel{}{}\stackrel{~}{𝐒}_e.`$ Here $`\stackrel{~}{𝐒}_e=USU^{}=(\stackrel{~}{S}_x,\stackrel{~}{S}_y,\stackrel{~}{S}_z)`$ denotes the spin operator one obtains after the unitary transformation (5) described in the previous section, by restricting it to the actual two-dimensional subspace we consider. The matrix elements of $`\stackrel{~}{𝐒}_e`$ usually differ from those of $`𝐒.`$ As for the field strength $`\stackrel{}{},`$ a detailed model should take into account the mode structure of the cavity. However, the characteristic features of the dynamics due to the cavity field can be captured by a simpler model to be discussed here. We consider the microwave field as a single transverse (TM) mode being perpendicular to the $`z`$ (easy) axis and having a frequency $`\mathrm{\Omega }`$ and equal amplitudes in the $`x`$ and $`y`$ directions: $`H_I=\stackrel{~}{\mu }(\stackrel{~}{S}_x+\stackrel{~}{S}_y)/\sqrt{2}`$. We note that the choice of the polarization does not have essential influence on the results presented here. Now the equations describing the dynamics of the two-level system (without relaxation) can be written as
$$\frac{\varrho }{t}=\frac{i}{\mathrm{}}[H^{},\varrho ],$$
(8)
where $`\varrho `$ is the density operator of the effective two-level system and
$$H^{}=H_e+H_I=\left(\begin{array}{cc}\epsilon _0+\mathrm{}\omega /2& \mathrm{\Delta }/2\\ \mathrm{\Delta }^{}/2& \epsilon _0\mathrm{}\omega /2\end{array}\right)_{mm^{}}.$$
(9)
Here $`\mathrm{}\omega (t)=w(t)2\stackrel{~}{\mu }s^{}`$ and $`\mathrm{\Delta }=\mathrm{\Delta }_02\stackrel{~}{\mu }s`$ with $`s^{}`$ and $`s`$ being the diagonal and offdiagonal elements resulting from the coupling operator $`(\stackrel{~}{S}_x+\stackrel{~}{S}_y)/\sqrt{2}`$. This leads to the following equations for the population differences and the coherences between the states:
$`{\displaystyle \frac{d}{dt}}(\varrho _{mm}\varrho _{m^{}m^{}})`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}(\mathrm{\Delta }^{}\varrho _{mm^{}}\mathrm{\Delta }\varrho _{m^{}m})`$ (10)
$`{\displaystyle \frac{d}{dt}}\varrho _{mm^{}}`$ $`=`$ $`i\omega (t)\varrho _{mm^{}}+{\displaystyle \frac{i}{2\mathrm{}}}\mathrm{\Delta }(\varrho _{mm}\varrho _{m^{}m^{}}).`$ (11)
As $`\mathrm{\Omega }`$ is in the terahertz domain, we can separate a slowly varying amplitude of the time varying field and write
$$=(\frac{1}{2}B(t)e^{i\mathrm{\Omega }t}+c.c.)u_k(z),$$
(12)
where $`u_k(z)`$ is the corresponding mode function of the cavity, and $`\left|\dot{B}(t)\right|\mathrm{\Omega }\left|B\right|`$. This makes straightforward a similar separation for the offdiagonal elements of the density matrix: $`\varrho _{mm^{}}=R_{mm^{}}e^{i\mathrm{\Omega }t}`$ where $`R_{mm^{}}`$ is again assumed to vary slowly compared with $`e^{i\mathrm{\Omega }t}.`$ Substituting into Eq. (10) we can neglect terms oscillating with frequency $`2\mathrm{\Omega }`$ as they do not contribute essentially to the evolution of the state, i.e, this is the standard rotating wave approximation (RWA)Sargent et al. (1974). Introducing the notation $`Z_{mm^{}}=\varrho _{mm}\varrho _{m^{}m^{}}`$ for the inversion between levels $`m`$ and $`m^{}`$, we obtain:
$`{\displaystyle \frac{d}{dt}}Z_{mm^{}}`$ $`=`$ $`{\displaystyle \frac{i}{2\mathrm{}}}\left\{\left(R_{mm^{}}\mathrm{\Delta }_0^{}e^{i\mathrm{\Omega }t}R_{mm^{}}^{}\mathrm{\Delta }_0e^{i\mathrm{\Omega }t}\right)\stackrel{~}{\mu }(B^{}s_{mm^{}}^{}R_{mm^{}}Bs_{mm^{}}R_{mm^{}}^{})u_k^2(z)\right\}`$ (13)
$`{\displaystyle \frac{d}{dt}}R_{mm^{}}`$ $`=`$ $`i(\omega (t)\mathrm{\Omega })R_{mm^{}}+{\displaystyle \frac{i}{\mathrm{}}}(\mathrm{\Delta }_0e^{i\mathrm{\Omega }t}\stackrel{~}{\mu }Bs_{mm^{}})Z_{mm^{}}\text{.}`$ (14)
At a given time instant only one of the level pairs get into resonance with the cavity, therefore from now on we shall omit the indices $`m,m^{}`$. We shall come back to this point, and discuss the mechanism which selects the actual level pair. We can also average out the equations over a time period of a cycle of the oscillation that eliminates the terms varying with frequency $`\mathrm{\Omega }`$. A similar procedure can be performed in space over the wavelength $`\lambda =2\pi /k`$, and we shall make use of $`\frac{1}{\lambda }_0^\lambda u_k^2(z)=1/2`$.
We also have to describe the effects caused by other degrees of freedom. These additional interactions – among which the strongest one is the spin phonon coupling, i.e, the oscillation of the atoms in the lattice – are not taken into account by the Hamiltonian (1), but can significantly influence the dynamics. The effects due to the reservoir of phonons (i) can be dissipative, by simply taking up the energy from the spin system and (ii) dephasing, by randomly disturbing the relative phases of the magnetic states. Dissipative terms lead to a decay of the diagonal elements while dephasing reduce the off-diagonal elements of $`\varrho `$. The diagonal terms relax generally much slower than the offdiagonal ones, therefore we consider only this (so called transversal) relaxation. As usual, it will be taken into account by assuming a simple exponential decay with a time constant $`T_2`$, the order of magnitude of which can be estimated between $`10^5`$ and $`10^7`$ *s* in the temperature range we are interested inLeuenberger and Loss (2000). With this term we have:
$`{\displaystyle \frac{d}{dt}}Z`$ $`=`$ $`{\displaystyle \frac{i\stackrel{~}{\mu }}{2\mathrm{}}}(B^{}s^{}RBsR^{}),`$ (15)
$`{\displaystyle \frac{d}{dt}}R`$ $`=`$ $`i(\omega (t)\mathrm{\Omega })R{\displaystyle \frac{i}{2\mathrm{}}}\stackrel{~}{\mu }BsZR/T_2.`$ (16)
These equations are familiar from the theory of magnetic resonance, and as it is known, the effect of the field $`B`$ is essential when the frequency originating from the slowly varying longitudinal field gets close to the cavity frequency: $`\omega (t)\mathrm{\Omega }`$.
We also treat the mode amplitude of the cavity as a dynamical quantity. Following the usual semiclassical approach of radiation-matter interaction theorySargent et al. (1974), the time varying field resulting from the magnetic molecules of the crystal will be described here as the field of a sample with time dependent magnetic dipole moment density $``$ in the $`xy`$ plane. The appropriate component of the transverse $``$ field originating from $``$ as a source, obeys the damped inhomogeneous wave equation
$$\mathrm{\Delta }\dot{}/(c^2T_c)\ddot{}/c^2=\ddot{}/c^2,$$
(17)
where $`T_c`$ is the cavity lifetime. Within the cavity we expand the field into modes, and in accordance with Eq. (12), we also write:
$$=(\frac{1}{2}He^{i\mathrm{\Omega }t}+c.c.)u_k(z),=(\frac{1}{2}Me^{i\mathrm{\Omega }t}+c.c.)f(z)$$
(18)
where $`f(z)`$ is nonzero only within the sample, where it can be taken equal to $`u_k(z)`$. If one substitutes into Eq. (17), and makes an approximation exploiting that the amplitudes, $`H`$ and $`M`$ are slowly varyingSargent et al. (1974) with respect to $`e^{i\mathrm{\Omega }t}`$, one obtains the equation
$$\frac{dH}{dt}=i\mathrm{\Omega }\eta \frac{M}{2}\frac{1}{2T_c}H.$$
(19)
The filling factor
$$\eta =_𝒞u_k(z)f(z)𝑑z/_𝒞u_k^2(z)𝑑zl/L$$
(20)
arises when we project the resulting equation on the mode in question, by integrating over the volume of the cavity. Here $`l/L`$ is the ratio of the lengths of the sample and the cavity, corresponding to the geometries reported in the experimentsTejada et al. (2004); Vanacken et al. (2004).
The corresponding component of the transverse magnetization of the sample is given by
$$=N_0\stackrel{~}{\mu }Tr\left[\varrho (\stackrel{~}{S}_x+\stackrel{~}{S}_y)/\sqrt{2}\right]=N_0\stackrel{~}{\mu }(s^{}\varrho _{mm^{}}+s\varrho _{m^{}m}),$$
(21)
where $`N_0`$ is the number density of the molecules participating in the transition $`mm^{}`$. We note that the static part of $``$ containing $`s^{},`$ does not give rise to radiation. Eq. (21) connects the microscopic dynamical variables with the macroscopic ones. Recalling that $`\varrho _{mm^{}}=R_{mm^{}}e^{i\mathrm{\Omega }t}`$, we see that $`M=2N_0\stackrel{~}{\mu }s^{}R_{mm^{}}.`$ The slowly varying magnetic induction field acting on the molecules is given by $`B=\mu _0(H+\beta M)`$, where $`\mu _0`$ is the vacuum permeability and $`\beta `$ may differ from unity giving account of a local field correction resulting from the near field of the dipolesJackson (1998). It is natural to measure the time variable in units of the characteristic time
$$T_0=\left(\frac{2\mathrm{}}{\eta N_0\mathrm{\Omega }\mu _0\stackrel{~}{\mu }^2\left|s\right|^2}\right)^{1/2}.$$
(22)
As we shall see, the relation of $`T_0`$ and the rate of relaxation characterizes the dynamics: in case of $`T_0/T_21`$ the phase correlation of the individual emitters is conserved during the process and a superradiant pulse (or a sequence of pulses) can be emitted. On the other hand, $`T_0/T_2>1`$ indicates that relaxation effects are too strong to allow SR to occur.
To perform the calculations it is straightforward to introduce the dimensionless magnetic field strength and induction amplitudes:
$$h=\left(\frac{\mu _0}{N_0\mathrm{}\mathrm{\Omega }}\right)^{1/2}H,\text{ }b=\left(\frac{1}{\mu _0N_0\mathrm{}\mathrm{\Omega }}\right)^{1/2}B.$$
(23)
The field intensity can be measured as the energy density averaged out over the time and space period. The dimensionless field intensity $`I=\left|h\right|^2/2`$ gives the number of emitted photons of energy $`\mathrm{}\mathrm{\Omega }`$ per number of molecules participating in the given transition. Outside the sample $`h=b,`$ while within the sample one has
$$b=h+2\beta \frac{1}{\eta T_0\mathrm{\Omega }}Re^{i\psi },$$
(24)
where $`\psi `$ is the phase of the offdiagonal coupling constant: $`s=|s|e^{i\psi }`$. The dynamical equation for the magnetic field in dimensionless form reads:
$$\frac{dh}{d\tau }=\frac{\kappa }{2}h+iRe^{i\psi },$$
(25)
where $`\tau =t/T_0`$ and $`\kappa =(T_0/T_c),`$ is the damping coefficient of the cavity. These equations are to be solved together with Eqs. (15,16), which take the dimensionless form:
$`{\displaystyle \frac{d}{d\tau }}Z`$ $`=`$ $`i(b^{}Re^{i\psi }bR^{}e^{i\psi }),`$ (26)
$`{\displaystyle \frac{d}{d\tau }}R`$ $`=`$ $`iT_0(\omega (\tau )\mathrm{\Omega })Ribe^{i\psi }Z\gamma R.`$ (27)
with $`\gamma =T_0/T_2`$. Note that numerically (using SI units) $`T_0|s|^1\times 10^8s`$, where we substituted $`N_0\eta =10^{23}m^3`$, corresponding to the values reported in the experimentTejada et al. (2004). Depending on the transition $`mm^{}`$, usually the magnitude of the dimensionless matrix element $`|s|`$ is much less than unity, thus $`10ns`$ (obtained with $`|s|=1`$) is basically a lower bound for $`T_0`$. This value – at least at low temperatures – is less than the time scale of the relaxation, $`T_2`$, but clearly by orders of magnitude larger than the period of the microwave radiation $`1/\mathrm{\Omega }`$, thus our rotating wave approximation leading to Eqs. (15,16) is valid. Additionally, as $`\mathrm{\Omega }`$ is around $`10^{11}`$ $`s^1`$, Eq. (24) shows that it is a very good assumption to take $`h=b`$ within the sample as well.
## IV Results and Discussion
The generic scheme for the emission from the ensemble of magnetic molecules considered in this paper starts with inverted two-level systems that come into resonance with the cavity field at a certain value of the external magnetic field $`B_0`$. Besides resonance, an additional requirement for the transverse radiation to begin is that the wavelength corresponding to the transition frequency should be comparable or smaller than the size of the sample, otherwise the non-transverse near-field of the sources would dominate the emitted field at the location of the other molecules. This explains why in the experiments reported in Ref. \[Tejada et al., 2004\] the emission is in the $`mm`$ range wavelength.
Now we shall analyze if the observed radiation can be considered as superradiance demanding $`\gamma 1`$, or is it rather a maser effect, where the absence of phase relaxation is not crucial. Therefore we first assume that $`\gamma =0,`$ and see that equations (26) and (27) admit a simple constant of motion:
$$Z^2+2\left|R\right|^2=Z_0^2.$$
(28)
If $`\omega (\tau )`$ is changing sufficiently slowly, the condition of resonance $`\omega (\tau )\mathrm{\Omega }_c=0`$ is sustained during the dynamics of the emission. Then writing $`Z=Z_0\mathrm{cos}\theta (\tau )`$, and $`\left|R\right|=(Z_0/\sqrt{2})\mathrm{sin}\theta (\tau )`$, a simple equation yielding essential physical insight into the nature of the problem can be obtained. With the assumptions that $`R`$ is real, $`\psi =\pi /2`$ and using Eqs. (26,27), one has $`b=\frac{d}{d\tau }\theta (\tau )/\sqrt{2}\dot{\theta }/\sqrt{2}`$, and from Eq. (25) we obtain
$$\ddot{\theta }(\tau )+\kappa \dot{\theta }(\tau )/2Z_0\mathrm{sin}\theta (\tau )=0,$$
(29)
which is the equation of a damped pendulum ($`\theta `$ measured from the inverted position), being often discussed in coherent atom-field interactions. The physically realistic initial condition for $`\theta _0`$ is a small ($`\theta _00`$) value, as initially we expect the offdiagonal element of the density matrix to be small. This comes from the small initial polarization as a remnant of the coherence between the levels during the magnetic tunneling transition. Putting in such an initial condition takes into account the rapidly varying terms omitted in Eqs. (14) and (13) which would also lead to a small but nonzero $`R`$ initially. Experimental results reported in Ref. \[Tejada et al., 2004\] suggest that the radiation appears after a crossing with small tunneling probability $`P_{mm^{}}`$, thus there is a significant inversion present in the system, meaning $`Z_01`$ in the beginning of the radiation process.
If the cavity lifetime is short, $`\kappa `$ is large, the second term dominates over the first in Eq. (29). Neglecting $`\ddot{\theta },`$ the equation of the overdamped pendulum admits an analytical solution. For the emitted intensity one obtains
$$\left|b\right|^2=\dot{\theta }^2(\tau )/2=\frac{1}{2\tau _R^2}\mathrm{sech}^2[(\tau \tau _d)/\tau _R],$$
(30)
where $`\tau _R=\frac{\kappa }{2Z_0}`$, and $`\tau _d=\tau _R(\mathrm{ln}\theta _0/2)`$. In this bad cavity limit $`T_c`$ can be estimated as $`L/c,`$the time needed for a photon to leave the cavity of length $`L`$. Assuming $`Z_0=1`$, the characteristic time of the emission in usual units is given by $`T_R=T_0\tau _R=T_0^2\frac{c}{2L}=\frac{\mathrm{}c}{\eta LN_0\mathrm{\Omega }_c\mu _0\stackrel{~}{\mu }^2\left|s\right|^2}`$. We see that this time constant is inversely proportional to the number density of the molecules, $`N_0`$, while according to Eq. (30), the intensity of the radiation is proportional to $`N_0^2`$: these are the characteristic features of superradiance Benedict et al. (1996). In addition, there is a delay time $`T_d=\tau _dT_0`$ necessary for the appearance of the pulse described in Eq. (30).
In the case the cavity losses do not dominate the process, the energy of the field is fed back into the crystal and according to Eq. (29) this leads to a damped periodic process. Assuming a perfect cavity one has a kind of Rabi oscillations with a time dependent field: energy is exchanged periodically between the crystal and the field.
These considerations based on the analytic solutions, however, become only valid approximately, as they do not take into account the factor $`\omega (\tau )\mathrm{\Omega }_c`$ in Eq. (27). As we assume a constant $`\dot{B}_0`$, we can introduce a constant dimensionless external field sweep rate $`v`$ via
$$T_0(\omega (\tau )\mathrm{\Omega }_c)=v\tau ,$$
(31)
where the origin of the time axis is chosen so that $`\tau =0`$ corresponds to exact resonance: $`\omega (0)=\mathrm{\Omega }_c`$. As an example, at the $`m=10m^{}=8`$ transition with $`\dot{B}_0=30`$ mT/s and $`T_0=10^6s`$ we have $`v=T_0^2\stackrel{~}{\mu }\dot{B}_0(mm^{})/\mathrm{}0.1`$, leading to a dynamics significantly different from the analytical solution.
Qualitatively, we expect that around $`\tau =0`$ (crossing point) the coupling begins to act, and creates a superposition of the levels. The coherence of the levels begin to increase accompanied by a finite transition probability: $`Z`$ will differ substantially from $`Z(\mathrm{}).`$ Then the oscillation of the pendulum, i.e., the radiation starts, but as the levels separate, their energy difference and therefore the oscillation frequency becomes larger. Taking relaxation into account, the amplitude of these oscillations diminishes, the molecules do not emit radiation any longer and simultaneously $`Z`$ will reach a stationary value $`Z(\mathrm{}),`$ analogously to what is usually called quantum tunneling of magnetization, because a different $`Z`$ means different value of the expectation value of $`S_z.`$ In this sense the inclusion of this time dependent detuning leads to a similar effect as discussed in the problem of tunneling, but the resonance condition is ensured by the inclusion of the time dependent cavity field, therefore the dynamics is more complicated than in LZS theory.
Quantitatively, Fig. 2 shows the inversion and the intensity of the emitted radiation as a function of time for different strengths of the cavity decay. The effect of the resonance is clear, appreciable radiation and change in $`Z`$ is seen after $`\tau =0`$. As the cavity damping becomes stronger we have less oscillations in the emitted intensity and the process starts later. For small values of $`\kappa `$ the final inversion $`Z(\mathrm{})`$ is determined by the sweep rate, but when cavity losses become significant, the energy of the molecular system is lost via the cavity field during the process leading to $`Z(\mathrm{})1`$. A bad cavity ($`\kappa 1`$) overdamps the pendulum and the system leaves the vicinity of the resonance before observable emission occurs. We note that the inversion $`Z(\mathrm{})`$ is not necessarily closely related to the final magnetization of the sample, because following the photon emission, cascade (not purely two-level) transitions related to a given side of the two-well potential can also have a considerable probability.
So far it has been assumed that the phase memory of the system is conserved, $`\gamma =T_0/T_2`$ is small and accordingly the damping term $`\gamma R`$ was neglected in Eq. (27). This was the assumption that led to the coherent behavior of the molecules resulting in superradiant emission. However, in reality there are at least two reasons to consider nonzero $`\gamma `$. One of them is the spin-phonon coupling which is temperature dependent, thus can be reduced by cooling the sample. Additionally, if the size of the system is smaller than the wavelength the near field dipole-dipole coupling between the molecules becomes important and can be shown to lead to an effective phase relaxation.Gross and Haroche (1982); Friedberg et al. (1973) Microscopic studies of the latter effect can be found in Refs. \[Benedict et al., 1996; Zaitsev et al., 1983\], for a detailed recent work see Ref. \[Davis et al., 2005\]. At very low temperatures this effect can be even stronger than the homogeneous broadening mechanism caused by elastic collisions with the phonons. At temperatures around $`2`$K, however, where the experiments observing the radiation have been performed, the dephasing is predominantly due to spin-phonon interactions instead of dipole-dipole coupling.Sessoli (1995); Leuenberger and Loss (1999) In the present work all these relaxation mechanisms are incorporated effectively by an appropriately chosen damping coefficient $`\gamma `$.
The consequences of phase relaxation is shown in Fig. 3, where a moderate constant cavity decay ($`\kappa =0.2`$) is also taken into account. For weak dephasing, we have similar oscillations in $`Z`$ and pulse structure as shown in Fig. 2. Increasing the value of $`\gamma `$ the coherent Rabi oscillations disappear. Additionally, the final inversion $`Z(\mathrm{})`$ is a monotonically increasing function of $`\gamma `$, and this can be considered as a remarkable difference between the two decay mechanisms. Note that this saturation effect can be responsible for the additional steps in the hysteresis curve published in Ref. \[Tejada et al., 2004\] following the most pronounced one which is accompanied by microwave radiation: The nonzero population that remains on the upper level after the transition can lead to an observable change of the magnetization of the sample at a next avoided level crossing.
In the case of strong dephasing, the time derivative of $`R`$ can be neglected with respect of $`\gamma R,`$ and from Eq. (27) we obtain that $`R`$ follows adiabatically the time dependence of $`b`$. Substituting back into Eq. (26) we obtain the following rate equation description of the process:
$`{\displaystyle \frac{dZ}{d\tau }}`$ $`=`$ $`Z\left|b\right|^2{\displaystyle \frac{2\gamma }{\gamma ^2+v^2\tau ^2}},`$
$`{\displaystyle \frac{db}{d\tau }}`$ $`=`$ $`{\displaystyle \frac{\kappa }{2}}b+{\displaystyle \frac{bZ}{\gamma +iv\tau }}.`$ (32)
Here atomic coherence does not play any role, thus the process cannot be termed as superradiance, it is rather a maser, operating on the inverted magnetic levels.
An important experimental result is that the position of the peaks in $`dM/dB_0`$ corresponding to the radiation process does depend on the external field sweep rate $`v`$. In our model this is related to the time spent by the sample around the resonance. For a slowly changing $`B_0`$, the dynamics is similar to the case of constant detuning, where an analytical solution is known, while increasing the value of $`v`$, an appropriate numerical solution of the dynamics is needed. In the case of superradiance, Fig. 4 shows $`dM/dB_0`$ as a function of $`v\tau `$, i.e., the dimensionless external field $`B_0`$. As we can see, for larger values of the sweep rate $`v`$, the height of the emission peak decreases and its position is shifted towards higher field values. The shift is in agreement with the experimental findings, while as a consequence of the coherent interaction, $`dM/dB_0`$ exhibits oscillations with sign changes.
If we assume that the maser effect is responsible for the radiation and use the rate equations (32) to calculate the dynamics of the system, somewhat different results are obtained. As Fig. 5 shows, $`dM/dB_0`$ accompanied with maser radiation, scales similarly with $`v`$ as in the case of weak dephasing: larger sweep rates correspond to peaks at higher $`B_0`$ fields, thus taking the cavity effects into account, this scaling property is not characteristic for SR. However, the oscillations seen in the superradiant case are absent in Fig. 5.
We shall analyze now from the point of view of transversal relaxation if the observed radiation could be superradiance. The reduction procedure summarized in Sec. II allows us to calculate the matrix element $`|s|`$ and thus the characteristic time (22) of the emission process for any transition. As the experimentally observed radiation peaks were around $`1.4`$ T, we focus on this value of the external magnetic field. The level structure of the Hamiltonian (1) provides the resonant transition frequency for a given transition $`mm^{}`$, as well as the population of the upper level according to the Boltzmann factor. In this way we can calculate $`T_0`$ for any transition as a function of the temperature of the sample. As for SR to occur the dephasing rate $`\gamma =T_0/T_2`$ must be small, so we should look for the transition with the minimal value of $`T_0`$. We found that below approx. $`0.8`$ K the transition $`m=10m^{}=8`$ provides the shortest $`T_0`$, while above this temperature the transition from $`m=6`$ to $`m^{}=4`$ yields the minimal characteristic time. As Fig. 6 shows, for low temperatures, i.e., ground state tunneling, the minimal $`T_0`$ is of the order of seconds. This is a consequence of the very small coupling coefficient $`s`$, corresponding to this transition. For higher temperatures when the transition $`m=6m^{}=4`$ provides the shortest characteristic time, $`T_0`$ significantly decreases as a function of temperature. This is a consequence of the relation $`T_01/(|s|\sqrt{N_0})`$ (see Eq. (22)), where $`N_0`$ is temperature dependent. That is, above $`0.8`$ K, the most favorable conditions for superradiance might be realized in the case of the transition $`m=6m^{}=4,`$ with a still strong temperature dependence of $`T_0`$.
However, the energy emitted during a transition process is not necessarily the highest for the lowest $`T_0`$. In fact, the population of the $`m=6`$ level – which determines the maximum number of the active molecules – is not large enough to explain the magnitude of the emitted energy observed in a recent experiment. Ref. \[Hernandez-Minguez et al., 2005\] reports on radiative bursts of duration of a few milliseconds, where (at 2 K) the total energy emitted by the sample was detected to be around $`3`$ *nJ*. Using the parameters of the experiment,Hernandez-Minguez et al. (2005) we investigated all the possible transitions and found the best agreement with the experimental data for the transition $`m=8m^{}=6`$, giving a value of $`T_0`$ in the *ms* range and a total emitted energy to be around 1.5 *nJ*. (Note that for initial states below $`m=8`$ the time scale of the process turns out to be too long, while for $`m>8`$ the number of active molecules is too small.) Thus our model predicts that the process having the most important role in producing the observed radiation is the transition $`m=8m^{}=6.`$
As we have seen, the character of the emission depends on the ratio $`\gamma =T_0/T_2,`$ where $`T_2`$ is decreasing with increasing temperature. According to Fig. 6, even for the shortest possible $`T_0`$ and relatively weak dephasing, with $`T_2`$ around $`10^5`$$`10^6`$ *s* (Ref. \[Leuenberger and Loss, 1999\]), we have $`\gamma >1`$ at $`2`$ K. The millisecond time scale obtained here and observed also in the experiments is clearly longer than the relaxation time $`T_2`$. Thus – unless a yet unknown effect decreases the disturbance caused by phase relaxation – the process responsible for the experimentally observedTejada et al. (2004) bursts of radiation seems to be rather a maser effect than superradiance.
The fact that no emission was seen for external fields lower than 1.4 T (where there can be a resonance as well) can be explained by the strength of the coupling to the transversal mode: $`|s|`$ is generally at least an order of magnitude larger for the possibly relevant transitions at $`1.4`$ T than at previous resonances. However, in a good resonator – like the Fabry-Perot mirrors in Ref. \[Tejada et al., 2003\] – one may expect that several resonances have observable consequences, radiation and – as it has been detected – enhanced magnetic relaxation rates. Additionally, we note that at high sweep ratesVanacken et al. (2004) the system passes not only a single cavity resonance during the emission process, and consecutive resonances can broaden the peaks in the $`dM/dB_0(B_0)`$ plots.
## V Summary and Conclusions
In this paper we developed a model for the interaction of a crystal of molecular magnets with the magnetic field of a surrounding cavity. The sample itself generates this transversal field $`B`$, while it also acts back on the molecules. The most important point of our treatment is that the cavity mode with fixed frequency $`\mathrm{\Omega }`$ comes to resonance with a magnetic transition at a given value of the external longitudinal magnetic field. Around this resonance the interaction of the molecules with the mode significantly increases leading to an observable burst of electromagnetic radiation as well as a change in the magnetization of the sample. Our model can describe different mechanisms of this radiation, in fact, there is a continuous transition from superradiance to maser-like effects. The crucial parameter here is the ratio of two time scales, the characteristic time of the process and the dephasing time $`\gamma =T_0/T_2`$. For small values of $`\gamma `$ the time evolution of the molecules is coherent allowing for the strong collective effect of superradiance in a cavity. In the case of strong dephasing, the sample still can emit electromagnetic radiation, but now the coherence of the molecules plays no role, the maser rate equations with time dependent detuning can describe the process. For moderate values of $`\gamma `$ we have a transition between the two processes. By calculating the intensity of the emitted radiation we have shown that with increasing the sweep rate of the external magnetic field, the emission peaks are shifted towards higher field values in accordance with the experimental results. This statement holds for both emission mechanisms, but the detailed functional dependencies are different for SR and maser emission.
Based on realistic approximations for $`T_0`$ and $`T_2`$, the process responsible for the experimentally observed bursts of electromagnetic radiation is most probably not superradiance, but rather a maser effect. The comparison of time resolved experiments on the emitted radiation with our theoretical results would provide the necessary information in order to settle this question. We expect that at very low temperatures, when spin-phonon relaxation is weaker, the collective features of the radiation may become dominant. While this is an interesting problem on its own, it is expected that the analysis of the radiated field can yield additional information on the process of quantum tunneling, as well as on the detailed properties of the interaction of these crystals and the field.
###### Acknowledgements.
We thank O. Kálmán for her valuable comments. This work was supported by the Flemish-Hungarian Bilateral Programme, the Flemish Science Foundation (FWO-Vl), the Belgian Science Policy and the Hungarian Scientific Research Fund (OTKA) under Contracts Nos. T48888, D46043, M36803, M045596.
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# Solitary waves in elongated clouds of strongly-interacting bosons
## Abstract
We examine the propagation of solitary waves in elongated clouds of trapped bosonic atoms as the confinement, the strength of the interatomic interaction, and the atom density are varied. We identify three different physical regimes and develop a general formalism that allows us to interpolate between them. Finally we pay special attention to the transition to the Tonks-Girardeau limit of strongly-interacting bosons.
By appropriate manipulation of the trapping potential, confined ultracold atoms can achieve conditions of reduced dimensionality. Taking advantage of this property, two recent experiments described in Refs. Paredes ; Science have managed to create conditions such that bosonic atoms approached the so-called Tonks-Girardeau phase TGIR predicted for strongly-interacting bosons in one dimension. In the experiment of Ref. Science , bosonic atoms were confined in elongated traps, and the gas was observed to approach the Tonks-Girardeau limit as the transverse confinement was increased. Evidence for this transition was obtained from measurements of the energy and of the axial size of the gas. In the experiment of Ref. Paredes the momentum distribution provided evidence for this transition.
Motivated by these experiments, we will consider here how solitary waves emerge in strongly-interacting, quasi-one-dimensional atomic gases. We first present a general description of the problem, identifying three physically distinct regimes Liebnew . We then develop a formalism that allows us to calculate the density profile, energy, and momentum associated with grey/dark solitary waves Lieb ; Tsuzuki ; JKP ; Kol3 ; GW ; KPa ; JK ; KPa2 . This formalism can be applied to any quasi-one dimensional system given only knowledge of the energy per unit length as function of the density per unit length. We find interesting changes in the properties of solitary waves as a function of density which suggest that the study of solitary waves in these elongated clouds could provide confirmation of the transition to the Tonks-Girardeau limit Burger ; Den . Our results also suggest eventual technological applications, e.g., the transmission of signals in atomic waveguides.
Here, we will neglect the effects of trapping along the $`z`$-axis of weak confinement and will for simplicity consider a cylindrical trap, $`V=M\omega _{}^2(x^2+y^2)/2`$, where $`M`$ is the atom mass, and $`\omega _{}`$ is the frequency of the trapping potential in the transverse direction. We approximate a short-ranged atom-atom interaction as $`V_{\mathrm{int}}(𝐫𝐫^{})=U_0\delta (𝐫𝐫^{})`$ where $`U_0=4\pi \mathrm{}^2a/M`$ and $`a`$ is the s-wave scattering length for elastic atom collisions. The Gross-Pitaevskii equation for the order parameter, $`\mathrm{\Psi }`$, then has the form
$`i\mathrm{}_t\mathrm{\Psi }=(\mathrm{}^2^2/2M+U_0|\mathrm{\Psi }|^2+V)\mathrm{\Psi }.`$ (1)
Following Ref. KP , we assume that the transverse dimension of the cloud is sufficiently small and the corresponding time scale sufficiently short that the transverse profile of the particle density can adjust to the equilibrium form appropriate for the local atomic density. With this approximation, the problem becomes one-dimensional, and the solitary pulse can be described by a local velocity, $`v(z)`$, and a local density of particles per unit length, $`\sigma (z)=|\mathrm{\Psi }(x,y,z)|^2𝑑x𝑑y`$ KP . The order parameter can then be written as a simple product JKP with $`\mathrm{\Psi }(𝐫,t)=f(z,t)g(x,y,\sigma )`$, where $`g`$ is the equilibrium wavefunction for the transverse profile. If $`g`$ is chosen to be normalized, $`|g|^2𝑑x𝑑y=1`$, the equations above imply that $`|f|^2=\sigma `$.
In this problem there are three qualitatively distinct physical regimes and two corresponding transitions separating them. The first of these transitions involves a change in the transverse profile of the cloud. For weak transverse confinement, $`|g|^2`$ can be calculated using the Thomas-Fermi approximation and is simply a parabola. For transverse confinement of a strength sufficient to ensure that the chemical potential is much smaller than $`\mathrm{}\omega _{}`$, $`|g|^2`$ becomes Gaussian. The second transition involves a change in the longitudinal profile of solitary waves. For weak transverse confinement, the atoms form a condensate. In the limit of strong interactions, strong transverse confinement and small linear densities, however, bosonic atoms behave in a certain sense like non-interacting fermions. This transition involves no change in the (Gaussian) transverse profile since the chemical potential of the gas is much smaller than $`\mathrm{}\omega _{}`$ in both cases.
We first consider the transition in the transverse direction. The critical value of $`\sigma a`$ for this transition is determined by the condition that the kinetic energy associated with transverse motion, $`\mathrm{}^2/MR^2`$, is comparable to the typical interaction energy, $`nU_0`$, where $`R`$ is the transverse width of the cloud and $`n`$ is the typical (three-dimensional) density. Since $`nN/(Z\pi R^2)`$, where $`N/Z=\sigma `$ is the density per unit length, one sees that the value of $`\sigma `$ at the crossover is $`\sigma _{c,1}1/a`$. It is interesting to note that both energy scales vary as $`1/R^2`$ with the consequence that $`\sigma _{c,1}`$ is independent of the oscillator length $`a_{}=(\mathrm{}/M\omega _{})^{1/2}`$.
Under typical experimental conditions (i.e., where a single trap rather than a series of tubes is used), $`\sigma a1`$. The cloud is thus in the Thomas-Fermi regime. The situation changes, however, with greater transverse confinement. To see this, we recall that the cloud expands along the long axis of the trap when it is squeezed transversely. Since $`Z`$ increases, $`\sigma =N/Z`$ decreases.
In the limit $`\sigma a1`$, $`|g|^2`$ has a Gaussian form, $`|g|^2=\mathrm{e}^{(\rho /a_{})^2}/(\pi a_{}^2)`$. As shown in Ref. JKP , $`f`$ then satisfies the equation
$`i\mathrm{}_tf=(\mathrm{}^2/2M)_z^2f+\mathrm{}\omega _{}(1+2a|f|^2)f.`$ (2)
We see from this equation that $`f\mathrm{e}^{i\omega _{}(1+2a\sigma _0)t}`$ as $`|z|\mathrm{}`$, where $`\sigma _0`$ is the background linear density. Rewriting Eq. (2) using the variable $`w=fe^{i\omega _{}(1+2a\sigma _0)t}`$, we obtain
$`i\mathrm{}_tw=(\mathrm{}^2/2M)_z^2w+\mathrm{}\omega _{}2a(|w|^2\sigma _0)w.`$ (3)
Equation (3) includes a familiar (i.e., quadratic) nonlinear term and leads to a speed of sound, $`c_1`$, which satisfies the equation $`Mc_1^2=2\mathrm{}\omega _{}\sigma _0a`$. Since $`\sigma _0=n_0\pi a_{}^2`$, we see that $`Mc_1^2=n_0U_0/2`$ JKP .
As shown in Refs. JKP ; JK , the density associated with a solitary wave of velocity $`u`$ has the form
$`\sigma (z)/\sigma _01=\mathrm{cos}^2\theta /\mathrm{cosh}^2(z\mathrm{cos}\theta /\zeta ),`$ (4)
where $`\theta =\mathrm{sin}^1(u/c_1)`$ and $`\zeta =2\xi (n_0)`$. Here $`\xi (n_0)`$ is the coherence length for $`n_0=\sigma _0/(\pi a_{}^2)`$, so that $`\zeta =a_{}/(2\sigma _0a)^{1/2}`$. Further, the dispersion relation connecting the energy and momentum of a solitary wave, $`(𝒫)`$, is given parametrically as $`/_0=(4\sqrt{2}/3)\mathrm{cos}^3\theta `$, where $`_0=\mathrm{}\omega _{}(\sigma _0a)^{1/2}\sigma _0a_{}`$, and $`𝒫/𝒫_0=\pi u/|u|2\theta \mathrm{sin}2\theta `$, where $`𝒫_0=\sigma _0\mathrm{}`$.
We now consider the second, longitudinal transition. As noted above, the gas approaches the so-called Tonks-Girardeau limit, in which the bosons to some extent behave like non-interacting fermions, as the transverse confinement is increased. The motion of the atoms in this limit is effectively one-dimensional since transverse degrees of freedom are frozen out when $`\mathrm{}\omega _{}`$ is much larger than the chemical potential. In the crossover region, each atom occupies a length of order $`1/\sigma _0`$ along the axis of the trap. The corresponding kinetic energy is on the order of $`\mathrm{}^2\sigma _0^2/M`$. The typical interaction energy, on the other hand, is on the order of $`n_0U_0`$, where $`n_0\sigma _0/a_{}^2`$ so that $`n_0U_0\mathrm{}^2\sigma _0a/(Ma_{}^2)`$. Thus, the ratio between the interaction energy and the kinetic energy is $`a/(\sigma _0a_{}^2)[\sigma _0\xi (n_0)]^2`$. Given the assumption of strong interactions, this quantity is much larger than unity for small values of $`\sigma _0`$ (i.e., for low densities), small values of $`a_{}`$ (i.e., for strong transverse confinement), or large values of $`a`$ (i.e., for strong interactions). The transition between the two regimes takes place for $`\sigma _{c,2}a/a_{}^2`$.
For values of $`\sigma _0`$ much smaller than $`\sigma _{c,2}`$ Kol2 ,
$`i\mathrm{}_tf=(\mathrm{}^2/2M)(_z^2f+\pi ^2|f|^4f).`$ (5)
In this case we see that $`fe^{i(\pi ^2\sigma _0^2\mathrm{}/2M)t}`$ as $`|z|\mathrm{}`$. We thus rewrite Eq. (5) using the variable $`w=f\mathrm{e}^{i(\pi ^2\sigma _0^2\mathrm{}/2M)t}`$ to obtain
$`i\mathrm{}_tw=(\mathrm{}^2/2M)[_z^2w+\pi ^2(|w|^4\sigma _0^2)w].`$ (6)
The speed of sound $`c_2`$ is now given as $`Mc_2^2=\pi ^2\mathrm{}^2\sigma _0^2/M`$ or $`c_2=\pi \mathrm{}\sigma _0/M`$. This is precisely the Fermi velocity in one dimension with $`k_F=\sigma _0\pi `$.
As shown in Ref. Kol3 , the density of the cloud associated with the solitary wave is now
$$\sigma (z)/\sigma _01=\frac{3\mathrm{cos}^2\theta }{2+(1+3\mathrm{sin}^2\theta )^{1/2}\mathrm{cosh}(2\pi \sigma _0z\mathrm{cos}\theta )},$$
(7)
where $`\theta =\mathrm{sin}^1(u/c_2)`$. The dispersion relation can again be expressed parametrically using
$`/_0={\displaystyle \frac{\sqrt{3}\pi }{2}}\mathrm{cos}^2\theta \mathrm{ln}\left[{\displaystyle \frac{2+3\mathrm{cos}^2\theta }{(1+3\mathrm{sin}^2\theta )^{1/2}}}\right],`$ (8)
with $`_0=\mathrm{}^2\sigma _0^2/M`$, and
$$𝒫/𝒫_0=\pi (/_0)\mathrm{tan}\theta +\mathrm{cos}^1\left[\frac{3\mathrm{sin}^2\theta 1}{(1+3\mathrm{sin}^2\theta )^{1/2}}\right].$$
(9)
It is natural to ask how one actually proceeds from one region to another. The transverse transition was studied in some detail in Refs. JKP ; KPa2 . Here, we focus on the second, longitudinal transition using a formalism of more general applicability. The basic ingredient required is the equation of state, i.e., the energy of the gas per unit length, $`ϵ(\sigma )`$, as function of $`\sigma `$. This can be calculated from the Lieb-Liniger model LL and can be written as $`ϵ(\sigma )=(\mathrm{}^2\sigma ^3/2M)\stackrel{~}{e}(\gamma )`$, where $`\stackrel{~}{e}(\gamma )`$ is a numerically known function and $`\gamma =2a/(\sigma a_{}^2)`$.
Given $`ϵ(\sigma )`$, one can immediately determine the sound velocity as $`Mc^2=\sigma _0^2ϵ/\sigma ^2(\sigma _0)`$. The energy of the solitary wave is
$`={\displaystyle \left(\frac{\mathrm{}^2}{2M}\frac{w^{}}{z}\frac{w}{z}+ϵ(\sigma )\sigma ϵ^{}(\sigma _0)+C\right)𝑑z}`$ (10)
with $`ϵ^{}=ϵ/\sigma `$. The final term in this equation represents the energy of the background density of atoms and ensures convergence of the integral. Thus, $`C=ϵ(\sigma _0)+\sigma _0ϵ^{}(\sigma _0)`$.
Equation (10) implies that $`w`$ satisfies the equation
$`i\mathrm{}_tw=(\mathrm{}^2/2M)_z^2w+[ϵ^{}(\sigma )ϵ^{}(\sigma _0)]w.`$ (11)
Writing $`w=\sqrt{\sigma }\mathrm{e}^{i\varphi }`$ and separating the real and imaginary parts of Eq. (11), we obtain the two hydrodynamic equations
$`{\displaystyle \frac{\mathrm{}^2}{2M}}\left({\displaystyle \frac{\sqrt{\sigma }}{z}}\right)^2=ϵ(\sigma )ϵ(\sigma _0)(\sigma \sigma _0)ϵ^{}(\sigma _0)`$
$`Mu^2{\displaystyle \frac{(\sigma \sigma _0)^2}{2\sigma }}`$ (12)
and $`v=(\mathrm{}/M)\varphi /z=u(1\sigma _0/\sigma )`$. Here, we have imposed the boundary condition that $`v0`$ as $`\sigma \sigma _0`$. Eq. (10) can be written as
$`={\displaystyle }[{\displaystyle \frac{\mathrm{}^2}{2M}}\left({\displaystyle \frac{\sqrt{\sigma }}{z}}\right)^2+{\displaystyle \frac{\mathrm{}^2\sigma }{2M}}\left({\displaystyle \frac{\varphi }{z}}\right)^2+ϵ(\sigma )ϵ(\sigma _0)`$
$`(\sigma \sigma _0)ϵ^{}(\sigma _0)]dz.`$ (13)
Combining the above equations, we obtain
$$=2\left[ϵ(\sigma )ϵ(\sigma _0)(\sigma \sigma _0)ϵ^{}(\sigma _0)\right]𝑑z.$$
(14)
Finally, the momentum of the solitary wave is given as
$$𝒫=M(\sigma \sigma _0)v(z)𝑑z=Mu\frac{(\sigma \sigma _0)^2}{\sigma }𝑑z.$$
(15)
We thus see that, given knowledge of the energy per unit length $`ϵ(\sigma )`$, Eq. (12) allows us to determine the shape of the solitary wave $`\sigma (z)`$ for a given velocity, $`u`$. Equations (14) and (15) then give $`(u)`$ and $`𝒫(u)`$, which can be combined to establish the dispersion relation, $`=(𝒫)`$. While the profiles of the solitary wave and the corresponding dispersion relations are known in all three limits examined earlier Lieb ; Tsuzuki ; Kol3 ; KPa ; JK ; KPa2 , Eqs. (11) – (15) allow us to interpolate between them. As indicated, our approach is quite general and merely requires knowledge of $`ϵ(\sigma )`$.
Let us now consider a specific example for the case $`\gamma _0=2a/(\sigma _0a_{}^2)=1`$. It is convenient to express $`ϵ`$ explicitly in terms of $`\sigma `$ with the result that $`ϵ(\sigma )=[\mathrm{}^2\sigma _0^3/2M]\stackrel{~}{e}(\sigma /\sigma _0)`$, where a reliable interpolation formula for $`\stackrel{~}{e}`$ is given as $`\stackrel{~}{e}(y)=(\pi ^2\gamma _0^3y^3/3+\kappa \gamma _0y^5)/(\gamma _0^3+4\gamma _0^2y+4\gamma _0y^2+\kappa y^3)`$ with $`y=\sigma (z)/\sigma _0`$ and $`\kappa 6.879`$.
For $`\gamma _0=1`$, Eq. (12) then has the form
$`\left({\displaystyle \frac{\sqrt{y}}{\stackrel{~}{z}}}\right)^2={\displaystyle \frac{\pi ^2y^3/3+\kappa y^5}{1+4y+4y^2+\kappa y^3}}{\displaystyle \frac{\pi ^2/3+\kappa }{9+\kappa }}`$
$`(y1)\stackrel{~}{ϵ}/y|_{y=1}\left({\displaystyle \frac{Mc}{\mathrm{}\sigma _0}}\right)^2{\displaystyle \frac{u^2}{c^2}}{\displaystyle \frac{(y1)^2}{y}},`$ (16)
where $`\stackrel{~}{z}=\sigma _0z`$. Given that $`\stackrel{~}{e}(y)/y1.471`$ and $`^2\stackrel{~}{e}(y)/y^21.872`$ for $`y=1`$, we find that the sound velocity is given as $`(Mc/\mathrm{}\sigma _0)^20.936`$. For a given value of $`u/c`$, Eq. (16) gives the profile of the solitary wave $`\sigma (z)/\sigma _0`$. Figure 1 shows such solutions for $`u/c=0.1`$ (bottom curve at $`z=0`$), 0.4, 0.7, and 0.9 (top).
The energy and the momentum can then be calculated from Eqs. (14) and (15):
$`(u)={\displaystyle \frac{\mathrm{}^2\sigma _0^2}{M}}{\displaystyle _{\mathrm{}}^+\mathrm{}}[{\displaystyle \frac{\pi ^2y^3/3+\kappa y^5}{1+4y+4y^2+\kappa y^3}}{\displaystyle \frac{\pi ^2/3+\kappa }{9+\kappa }}`$
$`1.471(y1)]d\stackrel{~}{z},`$ (17)
and
$`𝒫(u)=Mu{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\stackrel{~}{z}(y1)^2/y.`$ (18)
The dispersion relation $`=(𝒫)`$, initially obtained by Lieb, is shown as the continuous (lower) line in Fig. 2. The maximum value of $`𝒫`$ corresponds to a dark solitary wave with $`u=0`$. For long wavelengths, the solitary waves become ordinary sound waves with the usual linear dispersion relation Tsuzuki ; KPa ; JK ; KPa2 . From Eq. (11) one can also calculate the dispersion relation appropriate for the Bogoliubov mode. This has the usual form
$`=\sqrt{\left(𝒫^2/2M\right)^2+(c𝒫)^2}.`$ (19)
The Bogoliubov dispersion relation, $`=(𝒫)`$, is shown as the dashed (higher) line in Fig. 2.
The two conditions derived above establishing the typical values of $`\sigma `$ that characterize the three regimes, $`\sigma _{c,1}a1`$ and $`\sigma _{c,2}a_{}^2a`$, give values of $`\sigma `$ which differ by roughly one order of magnitude when the transverse confinement is strong as in Ref. Science . For example, for $`a100`$ Å and $`\omega _{}=2\pi \times 70.7`$ kHz (which implies that $`a_{}0.04\mu `$m), we see that $`\sigma _{c,1}10^6`$/cm and $`\sigma _{c,2}10^5`$/cm.
Another remarkable feature of this problem is that it is independent of the scattering length in the Tonks-Girardeau limit. Specifically, in the two limiting cases (i.e., the one-dimensional Gross-Pitaevskii equation and the Tonks-Girardeau equation), the characteristic widths of solitary waves is on the order of the coherence length $`\xi a_{}/(\sigma _0a)^{1/2}`$ and on the order of the interparticle spacing $`1/\sigma _0`$, respectively. This allows us to understand the qualitative features of the profiles shown in Fig. 3. For the largest value of $`\sigma _0=10\sigma _c`$, which explores the Gross-Pitaevskii limit, the characteristic width of the disturbance (in units of $`\sigma _0^1`$) is on the order of $`(\sigma _0a_{}^2/a)^{1/2}1/\sqrt{\gamma _0}`$, which is larger than unity. In the opposite Tonks-Girardeau limit with $`\sigma =\sigma _c/10`$, the width is of order unity in units of $`\sigma _0^1`$. Therefore, as one moves from one regime to the other, the width of the solitary wave in the dimensionless units of Fig. 3 decreases and eventually reaches a constant value of order unity. This fact can be used as an experimental signature of the transition from one limit to the other.
The present model is valid provided that the solitary waves have a width which is larger than the atom-atom spacing. While our predictions are thus expected to be reliable for waves whose size is larger than $`1/\sigma _0`$ (e.g., sound waves), they are pushed to the limits of their validity in the extreme cases of the Tonks-Girardeau limit ($`\gamma 1`$) and for narrow (and thus slow) pulses. Predictions in these regions are expected to be qualitatively, but not quantitatively, correct Kol3 . The Boboliubov spectrum, on the other hand, agrees with the exact result Lieb for $`𝒫0`$ and $`𝒫\mathrm{}`$ for all values of $`\gamma `$.
Finally, as seen in Fig. 2, the dispersion relations for the two modes converge in the limit of long wavelength since both correspond to small amplitude sound waves in this limit. For shorter wavelengths, however, they diverge. The characteristic momentum, $`𝒫_c`$, for which differences appear is on the order of $`\mathrm{}/\xi `$. Thus, $`𝒫_c/𝒫_0\sqrt{a/(\sigma _0a_{}^2)}\sqrt{\gamma _0}`$, which is on the order of unity at the transition and larger than unity as one approaches the Tonks-Girardeau regime. This behavior is quite different from that between the three-dimensional and the one-dimensional problems studied in Refs. KPa ; JK ; KPa2 , where the two modes have different energies and momenta. In this respect the present problem more closely resembles the homogeneous problem initially studied by Lieb Lieb ; Tsuzuki .
We thank Nikos Papanicolaou and Stephanie Reimann for useful discussions. M.Ö. thanks Peter Drummond for presenting the Lieb-Liniger model to him. We also acknowledge financial support from the European Community project ULTRA-1D (NMP4-CT-2003-505457).
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# 2005 International Linear Collider Workshop - Stanford, U.S.A. SLAC-PUB-11295 Supersymmetry Parameter Analysis with Fittino
## I INTRODUCTION
If low-energy Supersymmetry (SUSY) susy is realized in nature, the next generation of colliders, the Large Hadron Collider (LHC) and the International Linear Collider (ILC) are likely to copiously produce SUSY particles and will allow for precise measurements of their properties. Once SUSY is established experimentally, it is the main task to explore the unknown mechanism of SUSY breaking (SSB). The most general way to do so is to reconstruct the parameters of the general Minimal Supersymmetric Standard Model (MSSM) mssm parameter space, which is significantly more ambitious than a test of specific models of SUSY breaking, like e. g. minimal supergravity (mSUGRA) msugra can be tested against the data in a relative straight-forward manner (see e. g. msugrafits ) due to their small number of parameters.
The general MSSM Lagrangian contains more than 100 new parameters which are only related to each other through the unknown SSB mechanism. Many of them (like complex phases) are already limited to some extent by measurements, e.g. by the absence of neutron and electron electric dipole moments edm . Furthermore universality of first and second generation appears to be a reasonable approximation while large Yukawa couplings in the third generation lead to significant differences. Applying these constraints, the number of free parameters is reduced to 19, including the Standard Model top quark mass as a parameter to account for parametric uncertainties. It is of high interest how well this 19-dimensional parameter space can be restricted with future measurements from LHC and ILC.
The experimental collaborations ATLAS atlas and CMS CMS at the LHC have performed detailed simulations of the possibilities to extract mass information from their future data. At a future electron positron collider for collision energies up to 1 TeV, the International Linear Collider (ILC), the kinematically accessible part of the superpartner spectrum can be studied in great detail due to favorable background conditions and the well-known initial state ref:LHCILC . The attempt of an evaluation of the MSSM Lagrangian parameters and their associated errors is only useful if the experimental errors of the future measurements are known and under control. The by far best-studied SUSY scenario to serve as a basis for such an evaluation is an mSUGRA-inspired benchmark scenario, the SPS1a scenario ref:SPS . For this scenario with a relatively light superpartner spectrum, a wealth of experimental simulations exists and has recently been compiled in the framework of the international LHC/ILC study group ref:LHCILC . In this paper, we perform a global fit of the 19 MSSM parameters to those expected measurements augmented by possible measurements of topological cross-sections (i. e. cross sections times branching fractions) that will be possible at the ILC with polarized beams. Since a detailed experimental simulation for some measurements is lacking, we estimate their uncertainties conservatively from the predicted cross-sections and transferring experimental efficiency from well-studied cases.
As theoretical basis for the global fit presented in this paper we use the loop-level calculations as implemented into the program SPheno ref:SPheno . In SPheno, the masses and decay branching fractions of the superpartners are calculated as well as production cross sections in $`\text{e}^+\text{e}^{}`$ collisions.
In this paper we present the result of a global fit of the MSSM Lagrangian parameters at the electro-weak scale for a slightly modified SPS1a benchmark scenario using the program Fittino ref:Fittino ; ref:FittinoProgram . Previous evaluations of the errors of those parameters msugrafits did not attempt to develop a strategy to extract those parameters from data without a priori knowledge. Within Fittino, special attention is given precisely to this task, i. e. to find the parameter set which is most consistent with the data in a $`\chi ^2`$ minimization before a careful evaluation of errors and correlations is performed. A similar code sfitter exploits a different strategy for this task.
This paper is organized as follows. In Section II we shortly describe the approach used in Fittino. In Subsection II.1, the input observables are explained. The results of the fit and the error evaluation method are summarized in Subsection II.2 and conclusions are drawn in Section III.
## II SPS1a’ FIT
From the numerous fitting options provided by Fittino we have chosen the following fit procedure to extract the low-energy Lagrangian parameters. First, in order to be independent of human bias, start values for the parameters are calculated using tree-level relations between parameters and a few observables ref:FittinoProgram . For fits with many parameters these values are not good enough to allow a fitting tool like MINUIT ref:Minuit to find the global minimum due to the amount of loop-level induced cross-dependencies between the individual sectors of the MSSM. Therefore, in a second step, the parameters are refined, using a simulated annealing approach ref:FittinoProgram ; SimAnn ; SimAnn2 . As a result, the parameter values are close to the global minimum so that a global fit using MINUIT can find the exact minimum in a third step. To determine the parameter uncertainties and correlations many individual fits with input values randomly smeared within their uncertainty range are carried out. The parameter uncertainties and the correlation matrix are derived from the spread of the fitted parameter values.
### II.1 Input Observables from LHC and ILC
A number of anticipated LHC and ILC measurements serve as input observables to the fit. For the ILC, running at center-of-mass energies of 400 GeV, 500 GeV and 1 TeV is considered with 80 % electron and 60 % positron polarization. The predicted values of the observables are calculated using the following prescriptions:
* Masses:
The experimental uncertainties of the mass measurements are taken from ref:LHCILC .
* Cross-sections:
Only $`\text{e}^+\text{e}^{}`$ cross-sections are included in the fit. However, the measurement of absolute cross-sections is impossible for many channels, in which only a fraction of the final states can be reconstructed. Therefore, absolute cross-section measurements are only used for the Higgs-strahlung production of the light Higgs boson, which is studied in detail in ref:HiggsBr .
* Cross-sections times branching fractions:
Since no comprehensive study of the precision of cross-section times branching fraction measurements is available, the uncertainty is assumed to be the error of a counting experiment with the following assumptions:
+ The selection efficiency amounts to 50 %.
+ 80 % polarization of the electron beam and 60 % polarization of the positron can be achieved.
+ 500 fb<sup>-1</sup> per center-of-mass energy and polarization is collected.
+ The relative precision is not allowed to be better than 1 % and the absolute accuracy is at most 0.1 fb to account for systematic uncertainties.
All production processes and decays of SUSY particles and Higgs bosons are used which have a cross-section times branching ratio value of more than 1 fb in one of the $`\text{e}^+\text{e}^{}`$ polarization states LL, RR, LR, RL at $`\sqrt{s}=400`$ and $`500`$ GeV and LR or RL at $`\sqrt{s}=1000`$ GeV.
* Branching fractions:
The three largest branching fractions of the lightest Higgs boson are included. The uncertainties on the Higgs branching fractions are taken from ref:HiggsBr .
* Standard Model parameters:
The present uncertainties for $`m_\text{W}`$ and $`m_\text{Z}`$ are used as conservative estimates. The experimental uncertainty for the top mass $`m_\text{t}`$ is assumed to be 50 MeV ref:HiggsBr .
In order to check the influence of theoretical uncertainties on the fit results, the fit has been performed twice, once with experimental uncertainties only and a second time with experimental and theoretical uncertainties. Theoretical uncertainties for the mass predictions are scale uncertainties taken from ref:SPA . The experimental and theoretical contributions are added in quadrature. The assumed experimental and theoretical uncertainties for the masses are given in ref:LHCILC . For the cross-section and cross-section times branching fraction measurements, the smallest allowed relative precision is raised to 2 % for the fit including theoretical uncertainties (as opposed to 1 % for the fit without theoretical uncertainties). The full list of observables used in the fit and their uncertainties can be obtained from ref:Fittino .
### II.2 Fit Results
The input observables described in Section II.1 are used to determine the SUSY Lagrangian parameters in a global fit under the assumptions mentioned in Section I. In total 18 SUSY parameters remain to be fitted. In addition to those, the top mass $`m_\text{t}`$ is fitted, since it strongly influences parts of the MSSM observables. Thus 19 Lagrangian parameters are simultaneously determined in this fit.
As shown in Table 1 all parameters are perfectly reconstructed at their input values. Due to the fact that the input observables are unsmeared in this fit, the final $`\chi ^2`$ is close to zero at $`\chi ^2=5.3\times 10^5`$.
Most parameters are reconstructed to a precision better than or around 1 %. For the U(1) gaugino mass parameter $`M_1`$, an accurary below the per mil level is achieved for the fit including experimental uncertainties only. However, its precision is worse by a factor of 2 once theoretical uncertainties are included. The uncertainties on other parameters increase by up to a factor of 5 if theoretical uncertainties are taken into account, such as $`X_\text{t}`$. In order to fully benefit from the precision data provided in particular by the ILC, work must be invested to reduce uncertainties of theoretical predications.
After a successful convergence of the simulated annealing algorithm to the input parameter values, the fit uncertainties are evaluated by carrying out many individual fits with input values randomly smeared within their uncertainty range using a Gaussian probability density. The complete covariance matrix and the correlation matrix are derived from the spread of the fitted parameter values. For the fit without theoretical uncertainties 1002 such fits were made. The corresponding number for the case including theory uncertainties amounts to 993. For each of those fits, the parameter set is minimized. For a large and complex parameter space with large correlations among the parameters this method has turned out to be more robust than a MINOS error analysis. An example of the outcome of this procedure for the fit (with experimental uncertainties only) is shown in Figure 1 for the parameters $`\mathrm{tan}\beta `$ and $`M_1`$ for the 1002 independent fits.
Figure 2 shows the distribution of the $`\chi ^2`$ values for these 1002 independent fits for the case without theoretical uncertainties. The mean $`\chi ^2`$ obtained from a fit of the $`\chi ^2`$ distribution to the observed distribution is 124.6, in agreement with the expectation of 126 $`\pm `$ 0.64. This shows that the fits converge well to the true minimum of the $`\chi ^2`$ surface for the smeared observables, implying that the uncertainties extracted from the toy fit value distributions are correct.
### II.3 Relative Importance of Observables
In order to determine automatically the most important input observables influencing the precision with which a certain parameter, the contribution of each observable to $`\mathrm{\Delta }\chi ^2=\chi _{1\sigma }^2\chi _{\text{min}}^2`$ is calculated for the parameter of interest. $`\chi _{\text{min}}^2`$ is the $`\chi ^2`$ value at the minimum of the fit and $`\chi _{1\sigma }^2`$ is the corresponding value if the parameter is varied within $`\pm 1\sigma `$. A large $`\mathrm{\Delta }\chi ^2`$ indicates a strong correlation of the parameter with other parameters. A large contribution $`\mathrm{\Delta }\chi _i^2`$ of an observable $`i`$ to the total $`\mathrm{\Delta }\chi ^2`$ with respect to the contributions by the other observables means that the parameter is mainly determined by the measurement of observable $`i`$.
Generally the precision of most parameters is determined by ILC data, as shown as an example for $`\mathrm{tan}\beta `$ and $`\mu `$ in Table 2. As a surprise it can be seen that even for the uncertainty of the squark mass parameters the most constraining input is provided by ILC slepton mass measurements. It shows that the small contribution of the squark mass parameter $`M_{\stackrel{~}{\text{q}}_R}`$ to the slepton mass $`m_{\stackrel{~}{e}_R}`$ has a larger effect compared to the good precision of the measurement of $`m_{\stackrel{~}{e}_R}`$, than the large contribution of the squark mass parameter to the much less precisely measured squark mass.
### II.4 LHC and ILC Synergy
In order to test the dependence of the fit results on the synergy of LHC and ILC, a fit has been performed with all ILC-influenced observables removed. While the LHC observables may allow to determine a further reduced set of parameters (e.g. mSUGRA, if a simplified model described by fewer parameters still fits the data), the complete set of parameters studied here is not constrained by the LHC mass observables which currently enter the fit. For most parameters, the resulting parameter uncertainties exceed the LHC+ILC uncertainties by orders of magnitude. Only for 3 of the parameters studied here, namely $`M_{\stackrel{~}{\mathrm{q}}_L}`$, $`M_{\stackrel{~}{\mathrm{q}}_R}`$ and $`M_3`$, the uncertainty is in the same order of magnitude. This exemplifies the synergy effects of the LHC and the ILC.
## III CONCLUSIONS
The Lagrangian parameters of the MSSM assuming universality for the first and second generation and real parameters but without assumptions on the SUSY breaking mechanism are correctly reconstructed without usage of a priori information. This has been achieved using precision measurements at the LHC and ILC as input to a global fit exploiting the techniques implemented in the program Fittino.
Most of the Lagrangian parameters can be determined to a precision around the percent level. For some parameters an accuracy of better than 1 per mil is achievable if theoretical uncertainties are neglected. A fit including them revealed that theory uncertainties can significantly deteriorate the precision of the Lagrangian parameter determination. More work is needed to improve the accuracy of theoretical predictions in order to fully benefit from the experimental precision. The only parameter which cannot be strongly constrained by the observables used in the presented fit is $`X_\text{b}`$ resulting in a precision of only 30 % for this parameter.
###### Acknowledgements.
The authors are grateful to Gudrid Moortgat-Pick, Werner Porod, Sven Heinemeyer and the whole SPA working group for very fruitful discussions.
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# Travelling kinks in discrete ϕ⁴ models
## 1 Introduction
Spatially discretized partial differential equations (or, equivalently, chains of coupled ordinary differential equations) have attracted considerable attention recently. One of the issues that has been vigorously debated and that will concern us in this paper, is whether discrete systems can support solitary waves travelling without losing energy to resonant radiation and decelerating as a result. We address this issue for one of the prototype models of nonlinear physics, the $`\varphi ^4`$-theory:
$$u_{tt}=u_{xx}+\frac{1}{2}u(1u^2).$$
(1.1)
The $`\varphi ^4`$-equation (1.1) is Lorentz-invariant, and so the existence of the travelling kink
$$u(x,t)=\mathrm{tanh}\frac{xcts}{2\sqrt{1c^2}},$$
(1.2)
where $`|c|<1`$ and $`s`$, is an immediate consequence of the existence of the stationary kink for $`c=0`$. On the other hand, if we discretize equation (1.1) in $`x`$,
$$\ddot{u}_n=\frac{u_{n+1}2u_n+u_{n1}}{h^2}+f(u_{n1},u_n,u_{n+1}),$$
(1.3)
the translation and Lorentz invariances are lost, and the existence of the travelling kink (and even of an arbitrarily centred stationary one) becomes a nontrivial matter. In equation (1.3), $`u_n`$, $`n`$, $`t`$, $`h`$ is the lattice spacing, and the nonlinearity $`f(u_{n1},u_n,u_{n+1})`$ satisfies the continuity condition
$$f(u,u,u)=\frac{1}{2}u(1u^2).$$
(1.4)
We restrict ourselves to symmetric discretizations, i.e.
$$f(u_{n1},u_n,u_{n+1})=f(u_{n+1},u_n,u_{n1}).$$
(1.5)
Equation (1.1) results from (1.3) in the continuum limit, where $`u_n(t)=u(x_n,t)`$, $`x_n=nh`$ and $`h0`$. In this limit, the truncation error of the Taylor series is $`𝒪(h^2)`$. We shall be concerned with monotonic kink solutions of (1.3): $`u_{n+1}(t)u_n(t)`$ for all $`n`$. As $`h0`$, such monotonic discrete kinks approach the continuous kink (1.2).
The most common, one-site discretization of the nonlinearity function is given by
$$f(u_{n1},u_n,u_{n+1})=\frac{1}{2}u_n(1u_n^2).$$
(1.6)
It is a well established fact \[L88\], however, that the discrete Klein-Gordon equation (1.3)+(1.6) admits only a countable set of stationary monotonic kinks with the boundary conditions
$$\underset{n\mathrm{}}{lim}u_n(t)=1,\underset{n+\mathrm{}}{lim}u_n(t)=+1.$$
(1.7)
Physically, this fact is related to the presence of the Peierls-Nabarro barrier, an effective potential periodic with the spacing of the lattice. Half of the stationary kinks are centred at the minima (the on-site kinks) and the other half (the off-site kinks) at the maxima of the Peierls-Nabarro potential. There are no continuous families of stationary discrete kinks of the form $`u_n=u(ns)`$, with $`s`$ a free parameter, which would interpolate between the two solutions. In an abuse of terminology, we will be calling such families “translationally invariant kinks” — although, in the first place, translation invariance is a property of an equation rather than a solution, and in the second, all lattice equations are of course *not* translationally invariant. As for propagating waves, of special importance are kinks moving at constant speed and without the emission of radiation. We will be referring to such kinks, i.e. solutions of the form $`u_n=u(ncts)`$ where $`u(\xi )`$ is a monotonically growing function satisfying the boundary conditions (1.7), as sliding kinks, to emphasise the fact that they do not experience any radiative friction. Being an obstacle to the “translationally invariant” kinks, the Peierls-Nabarro barrier is also detrimental to the existence of sliding kinks — at least for small $`c`$ (see reviews in \[S03\] and \[IJ05\]).
In an attempt to find a discrete model with “translationally invariant” and sliding kinks, Speight and Ward \[SW94, S97\] considered a hamiltonian discretization of the form
$`f(u_{n1},u_n,u_{n+1})`$ $`=`$ $`{\displaystyle \frac{1}{12}}(2u_n+u_{n+1})\left(1{\displaystyle \frac{u_n^2+u_nu_{n+1}+u_{n+1}^2}{3}}\right)`$ (1.8)
$`+`$ $`{\displaystyle \frac{1}{12}}(2u_n+u_{n1})\left(1{\displaystyle \frac{u_n^2+u_nu_{n1}+u_{n1}^2}{3}}\right).`$
In the static limit, the corresponding energy admits a topological lower bound which is saturated by a first- (rather than second-) order difference equation. This equation is readily shown to have a one-parameter continuous family of stationary kink solutions $`u_n=u(ns)`$ for $`0h2`$ (see Proposition 1 in \[S97\]). The parameter $`s`$ of the family defines the position of the kink relative to the lattice. Since all members of the family have the same (lowest attainable) energy, the stationary kink experiences no Peierls-Nabarro barrier. As for travelling kinks, Speight and Ward’s numerical simulations revealed that although moving kinks in this model do lose energy to Cherenkov radiation and decelerate as a result, this happens at a slower rate than a similar process in equation (1.6) (see figures 4 and 5 in \[S97\]).
Another line of attack was chosen by Bender and Tovbis \[BT97\] who proposed a different discretization supporting a continuous family of arbitrarily centred stationary kinks:
$$f(u_{n1},u_n,u_{n+1})=\frac{1}{4}\left(u_{n+1}+u_{n1}\right)\left(1u_n^2\right).$$
(1.9)
In this case, the family arises due to the suppression of the stationary kink’s resonant radiation. In fact, the family of stationary kinks can be found explicitly as
$$u_n(t)=\mathrm{tanh}\left[a(ns)\right],$$
(1.10)
where $`a=\mathrm{arcsinh}\left(h/2\right)`$ for all $`h`$. (The solution (1.10) coincides with the stationary dark soliton of the repulsive Ablowitz-Ladik equation \[HA93\].)
Finally, the nonlinearity
$$f(u_{n1},u_n,u_{n+1})=\frac{1}{8}\left(u_{n+1}+u_{n1}\right)\left(2u_{n+1}^2u_{n1}^2\right).$$
(1.11)
was introduced by Kevrekidis \[K03\], who demonstrated the existence of a two-point invariant and hence a first-order difference equation associated with the stationary equation. Consequently, the discretization (1.11) also supports a continuous family of stationary kinks for all $`h[0,h_0]`$ with some $`h_0>0`$. (For general discussion, see \[S99\], \[BOP05\] and \[DKY05a\].) A relevant property of the model (1.11), which is related to the existence of a two-point invariant \[K03\] and indicates some additional underlying symmetry, is the conservation of momentum. (See also \[DKY05b\].)
Since the reasons for the nonexistence of “translationally invariant” kinks and of sliding kinks are apparently related (the breaking of symmetries of the underlying continuum theory or, speaking physically, the presence of the Peierls-Nabarro barrier), the availability of “translation-invariant” stationary kinks in the models (1.8), (1.9) and (1.11) suggests that they might have sliding kinks as well. It is the purpose of the present study to find out whether this is indeed the case. We shall analyse the persistence of continuous families of stationary kinks $`u_n=u(ns)`$ for nonzero velocities; in other words, examine the existence of solutions of the form $`u(ncts)`$ where $`u(z)`$ is a monotonically growing function satisfying (1.7), and $`c0`$. We develop an accurate numerical test in the limit $`h0`$ which shows whether or not standing and travelling kinks of the discrete $`\varphi ^4`$ model (1.3) bifurcate from the exact kink solutions (1.2) of its continuous counterpart (1.1). The analysis of this bifurcation poses a singular problem in perturbation theory which can be analysed using two (inner and outer) matched asymptotic scales on the complex plane \[TTJ98, T00a\]. In particular, the nonvanishing of the Stokes constant in the inner asymptotic equation serves as a sufficient condition for the non-existence of continuous solutions of the difference equations \[TTJ98\].
Our test will be based on computing the Stokes constant for the differential-difference equation underlying the lattice system. We will examine all four discretizations of the $`\varphi ^4`$ theory mentioned above, i.e. equations (1.6), (1.8), (1.9) and (1.11). Since translationally invariant stationary kinks $`u_n=u(ns)`$ do exist for the three exceptional nonlinearities (1.8), (1.9) and (1.11), the Stokes constant is *a priori* vanishing for $`c=0`$ in these three cases. However, we will show that in all three cases the Stokes constant acquires a nonzero value as soon as $`c`$ deviates from zero. It remains nonzero for all $`c`$ except a few isolated values which define the particular velocities of the sliding kinks in the corresponding model. There is one such isolated zero of the Stokes constant for the nonlinearity (1.8) and three sliding velocites for the discretization (1.11). Consequently, the main conclusion of this work is that the sliding kinks, i.e. kinks travelling at a constant speed without the emission of radiation, can occur only at particular values of the velocity. The sliding velocities are, of course, functions of the discretization spacing $`h`$, so that sliding kinks arise along continuous curves on the $`(c,h)`$-plane.
We conclude this introduction with a remark on a convention adopted in the remainder of this paper — namely, that the linear part of the function $`f(u_{n1},u_n,u_{n+1})`$ in (1.3) can always be fixed to $`\frac{1}{2}u_n`$ without loss of generality. Indeed, the most general function satisfying (1.4) and (1.5) is $`f=\left(\frac{1}{2}2a\right)u_n+a\left(u_{n+1}+u_{n1}\right)+\text{cubic terms}`$, where $`a`$ is arbitrary. Since $`h^2`$ in (1.3) is also a free parameter, we can always make a replacement $`h\stackrel{~}{h}`$ such that $`1/h^2+a=1/\stackrel{~}{h}^2`$. This gives
$$f(u_{n1},u_n,u_{n+1})=\frac{1}{2}u_nQ(u_{n1},u_n,u_{n+1}),$$
(1.12)
where $`Q`$ is a homogeneous polynomial of degree 3 which is independent of the parameter $`h`$.
The plan of this paper is as follows. In the next section (section 2) we review the construction of the outer and inner asymptotic solutions in the limit $`h0`$. Section 3 contains details of the numerical computation of the Stokes constants while the last section (section 4) summarises the results of our work.
## 2 Inner and outer asymptotic expansions in the limit $`h0`$
We are looking for a sliding-kink solution of the discrete $`\varphi ^4`$ models (1.3) in the form
$$u_n(t)=\varphi (z),z=h(ns)ct,$$
(2.1)
where $`\varphi (z)`$ is assumed to be a twice differentiable function of $`z`$, that satisfies the differential advance-delay equation
$$c^2\varphi ^{\prime \prime }(z)=\frac{\varphi (z+h)2\varphi (z)+\varphi (zh)}{h^2}+\frac{1}{2}\varphi (z)Q(\varphi (zh),\varphi (z),\varphi (z+h)),$$
(2.2)
with the boundary conditions $`\varphi (z)\pm 1`$ as $`z\pm \mathrm{}`$. The velocity $`c`$ is assumed to be smaller than 1 in modulus. If a solution to this boundary-value problem (i.e. a heteroclinic orbit) exists, then the parameter $`s`$ is arbitrary due to the translation invariance of the advance-delay equation (2.2). The scaling parameter $`h`$ (which stands for the lattice step-size) can be used to reduce equation (2.2) to a singularly perturbed differential equation as $`h0`$ \[TTJ98\]. Formal asymptotic solutions of the problem (2.2) can be constructed at the inner and outer asymptotic scales. The formal series represent convergent asymptotic solutions of the singular perturbation problem only if the Stokes constants are all zero \[T00a\].
Asymptotic analysis beyond all orders of perturbation theory was pioneered by Kruskal and Segur \[KS91\] and has been utilised by many authors. It was extended by Pomeau et. al. \[PRG88\] to allow the computation of radiation coefficients from Borel summation of series rather than from the numerical solution of differential equations. Essentially the same method has been applied to different problems by Grimshaw and Joshi \[GJ95, G95\] and Tovbis and collaborators \[TTJ98, T00a, T00b, TP05\]. In this paper, we shall work with formal inner and outer asymptotic series for the problem (2.2) without attempting rigorous analysis of their asymptoticity.
### 2.1 Outer asymptotic series
Assuming that the solution $`\varphi (z)`$ is a real analytic function of $`z`$, we consider the Taylor series for the second difference in a strip $`𝒟_\delta =\{z:|\mathrm{Im}z|<\delta \}`$, where $`\delta >0`$:
$$\varphi (z+h)2\varphi (z)+\varphi (zh)=h^2\varphi ^{\prime \prime }(z)+\underset{n=2}{\overset{\mathrm{}}{}}h^{2n}\frac{2}{(2n)!}\varphi ^{(2n)}(z).$$
(2.3)
Since the cubic polynomial $`Q(u_{n1},u_n,u_{n+1})`$ satisfies the continuity and symmetry relations (1.4) and (1.5), the nonlinearity of (2.2) can also be expanded in a Taylor series in the same strip:
$$Q(\varphi (zh),\varphi (z),\varphi (z+h))=\frac{1}{2}\varphi ^3(z)+\underset{n=1}{\overset{\mathrm{}}{}}h^{2n}Q_{2n}(\varphi ,(\varphi ^{})^2,\mathrm{},\varphi ^{(2n)}),$$
(2.4)
where the coefficients $`Q_{2n}`$ depend on even derivatives and even powers of odd derivatives of $`\varphi (z)`$, and also $`Q_{2n}(\varphi ,0,\mathrm{},0)=0`$. The differential advance-delay equation (2.2) can thus be written as
$`(1c^2)\varphi ^{\prime \prime }+{\displaystyle \frac{1}{2}}\varphi (1\varphi ^2)+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}h^{2n}\left({\displaystyle \frac{2}{(2n+2)!}}\varphi ^{(2n+2)}Q_{2n}(\varphi ,(\varphi ^{})^2,\mathrm{},\varphi ^{(2n)})\right)=0.`$
For $`h=0`$, equation (2.1) becomes the travelling wave reduction of the continuous model (1.1), with the explicit solution
$$\varphi _0(z)=\mathrm{tanh}\xi ;\xi =\frac{z}{2\sqrt{1c^2}},|c|<1.$$
(2.6)
We will search for solutions of equation (2.1) of the form
$$\widehat{\varphi }(z)=\varphi _0(z)+\underset{n=1}{\overset{\mathrm{}}{}}h^{2n}\varphi _{2n}(z).$$
(2.7)
Substituting the expansion (2.7) into (2.1) we get, at order $`h^{2n}`$,
$$\varphi _{2n}=H_{2n},$$
where the linearised operator $``$ is given by
$$=\frac{d^2}{d\xi ^2}+46\mathrm{sech}^2\xi ,$$
and $`H_{2n}`$ are polynomials in $`\varphi _0,\varphi _2,\mathrm{},\varphi _{2n2}`$ and their derivatives. The kernel of $``$ is one-dimensional, and spanned by an even eigenfunction $`y_0=\mathrm{sech}^2\xi `$. The rest of the spectrum of $``$ is positive. It is not difficult to prove by induction that if $`\varphi _{2k}(z)`$ are all odd in $`z`$ for $`k=0,1,\mathrm{},n1`$, the nonhomogeneous term $`H_{2n}`$ is also odd in $`z`$ and hence, by the Fredholm alternative, there exists a unique odd bounded solution $`\varphi _{2n}(z)`$ for $`z`$. Moreover, since $`H_{2n}`$ decays to zero exponentially fast as $`|z|\mathrm{}`$, the function $`\varphi _{2n}(z)`$ is also exponentially decaying for any $`n1`$. The perturbation $`\varphi _2(z)`$ in particular satisfies the nonhomogeneous equation
$$\varphi _2=\frac{1}{3}\left(\varphi _0^{(\mathrm{iv})}+\alpha \varphi _0^{\prime \prime }+\beta \varphi _0^2\varphi _0^{\prime \prime }+\gamma \varphi _0(\varphi _0^{})^2\right),$$
(2.8)
where the numerical coefficients depend on whether the nonlinearity function $`f`$ is given by (1.6), (1.8), (1.9) or (1.11):
$`\text{One-site (}\text{1.6}\text{)}:`$ $`\alpha =\beta =\gamma =0;`$
$`\text{Speight-Ward (}\text{1.8}\text{)}:`$ $`\alpha =1,\beta =\gamma =4;`$
$`\text{Bender-Tovbis (}\text{1.9}\text{)}:`$ $`\alpha =3,\beta =3,\gamma =0;`$
$`\text{Kevrekidis (}\text{1.11}\text{)}:`$ $`\alpha =3,\beta =9,\gamma =6.`$
The odd bounded solution $`\varphi _2(z)`$ of the nonhomogeneous equation (2.8) is:
$$\varphi _2(z)=𝒜\mathrm{tanh}\xi \mathrm{sech}^2\xi +\xi \mathrm{sech}^2\xi ,$$
(2.9)
where
$`𝒜={\displaystyle \frac{(1c^2)(\gamma +2\beta )+6}{72(1c^2)^2}},={\displaystyle \frac{(1c^2)(\alpha +\beta )+1}{24(1c^2)^2}}.`$
The hat in the series (2.7) indicates that the series is formal, i.e. it may or may not converge \[TTJ98, T00a\], depending on the choice of $`c`$ and $`Q`$ in equation (2.2). We shall be referring to (2.7) as the outer asymptotic expansion.
### 2.2 Inner asymptotic series
The leading-order term (2.6) of the outer expansion (2.1) has poles at $`\xi =\frac{\pi \mathrm{}}{2}(1+2n)`$, where $`n`$. We apply the scaling transformation
$$z=h\zeta +\mathrm{}\pi \sqrt{1c^2},\varphi (z)=\frac{1}{h}\psi (\zeta )$$
(2.10)
to equation (2.2) in order to study the convergence of the formal asymptotic solution (2.7) near the pole $`\xi =\frac{\pi \mathrm{}}{2}`$ (see \[TTJ98, T00a\]). This yields the following differential advance-delay equation for $`\psi (\zeta )`$:
$`c^2\psi ^{\prime \prime }(\zeta )=\psi (\zeta +1)2\psi (\zeta )+\psi (\zeta 1)Q(\psi (\zeta 1),\psi (\zeta ),\psi (\zeta +1))+{\displaystyle \frac{h^2}{2}}\psi (\zeta ).`$ (2.11)
The following are the cubic functions $`Q`$ for each of the four discretizations that we deal with in this paper:
$`\text{One-site (}\text{1.6}\text{)}:`$ $`Q`$ $`={\displaystyle \frac{1}{2}}\psi ^3(\zeta );`$
$`\text{Speight-Ward (}\text{1.8}\text{)}:`$ $`Q`$ $`={\displaystyle \frac{1}{36}}[\psi ^3(\zeta +1)+3\psi ^2(\zeta +1)\psi (\zeta )+3\psi (\zeta +1)\psi ^2(\zeta )`$
$`+4\psi ^3(\zeta )+3\psi (\zeta 1)\psi ^2(\zeta )+3\psi ^2(\zeta 1)\psi (\zeta )+\psi ^3(\zeta 1)];`$
$`\text{Bender-Tovbis (}\text{1.9}\text{)}:`$ $`Q`$ $`={\displaystyle \frac{1}{4}}\psi ^2(\zeta )\left[\psi (\zeta +1)+\psi (\zeta 1)\right];`$
$`\text{Kevrekidis (}\text{1.11}\text{)}:`$ $`Q`$ $`={\displaystyle \frac{1}{8}}[\psi ^3(\zeta +1)+\psi ^2(\zeta +1)\psi (\zeta 1)+\psi (\zeta +1)\psi ^2(\zeta 1)`$
$`+\psi ^3(\zeta 1)].`$
We note that the heteroclinic orbit becomes small as $`h0`$ under the normalization (2.10): if $`\varphi (z)\pm 1`$ as $`z\pm \mathrm{}`$, then $`\psi (\zeta )\pm h`$ as $`\mathrm{Re}\zeta \pm \mathrm{}`$. The formal asymptotic series (2.7) in the new variables (2.10) becomes a new formal series
$$\widehat{\psi }(\zeta )=\widehat{\psi }_0(\zeta )+\underset{n=1}{\overset{\mathrm{}}{}}h^n\widehat{\psi }_n(\zeta ),$$
(2.12)
where each term $`\widehat{\psi }_n(\zeta )`$ can be expanded in a formal series in descending powers of $`\zeta `$. In particular, the leading-order function $`\widehat{\psi }_0(\zeta )`$ has the general form
$$\widehat{\psi }_0(\zeta )=\underset{n=0}{\overset{\mathrm{}}{}}\frac{a_{2n}}{\zeta ^{2n+1}}.$$
(2.13)
By comparing the series (2.12) and (2.13) with the solutions (2.6) and (2.9) in variables (2.10), we note the correspondence:
$$a_0=2\sqrt{1c^2},a_2=\frac{(1c^2)(\gamma +2\beta )+6}{9\sqrt{1c^2}}.$$
We shall be referring to (2.12) as the inner asymptotic expansion. The odd powers of $`h`$ in the inner asymptotic expansion (2.12) appear due to the matching conditions with the outer asymptotic expansion (2.7) under the scaling (2.10), as well as due to the non-zero boundary conditions for the heteroclinic orbits $`\psi (\zeta )\pm h`$ as $`\mathrm{Re}\zeta \pm \mathrm{}`$.
### 2.3 Leading-order problem for an inner solution
Convergence of the formal inverse-power series (2.13) for the leading-order solution $`\widehat{\psi }_0(\zeta )`$ depends on the values of the Stokes constants \[T00a\]. Computation of the Stokes constants is based on Borel–Laplace transforms of the inner equation (2.11) \[TTJ98\]. Assuming continuity in $`h`$, we study the leading-order solution $`\psi _0(\zeta )=lim_{h0}\psi (\zeta )`$ of the truncated inner equation
$$c^2\psi _0^{\prime \prime }(\zeta )=\psi _0(\zeta +1)2\psi _0(\zeta )+\psi _0(\zeta 1)Q(\psi _0(\zeta 1),\psi _0(\zeta ),\psi _0(\zeta +1)).$$
(2.14)
By substituting the series (2.13) into equation (2.14), one can derive a recurrence relation between the coefficients in the set $`\{a_n\}_{n=0}^{\mathrm{}}`$. The Stokes constants can be computed from the asymptotic behavior of the coefficients $`a_n`$ for large $`n`$. Alternatively, the leading-order solution $`\psi _0(\zeta )`$ and the Stokes constants can be defined by using the Borel–Laplace transform:
$$\psi _0(\zeta )=_\gamma V_0(p)\mathrm{}^{p\zeta }p.$$
(2.15)
The choice of the contour of integration $`\gamma `$ determines the domain of $`\psi _0(\zeta )`$ in the complex $`\zeta `$-plane. We define two solutions $`\psi _0^{(s)}(\zeta )`$ and $`\psi _0^{(u)}(\zeta )`$, which lie on the stable and unstable manifolds respectively, such that
$$\underset{\mathrm{Re}\zeta +\mathrm{}}{lim}\psi _0^{(s)}(\zeta )=0,\underset{\mathrm{Re}\zeta \mathrm{}}{lim}\psi _0^{(u)}(\zeta )=0.$$
(2.16)
We note that the stable and unstable solutions tend to the stationary point at the origin, since the heteroclinic orbits connect the stationary points at $`\psi =\pm h`$ which move to the origin as $`h0`$. The three stationary points coalesce to become a degenerate stationary point at the origin within the truncated inner equation (2.14).
The Borel–Laplace transform (2.15) produces the stable solution $`\psi _0^{(s)}(\zeta )`$ when the contour of integration $`\gamma _s`$ lies in the first quadrant of the complex $`p`$-plane and extends from $`p=0`$ to $`p=\mathrm{}`$. Similarly, it produces the unstable solution $`\psi _0^{(u)}(\zeta )`$ when the contour of integration $`\gamma _u`$ lies in the second quadrant. We choose the integration contours in such a way that $`\mathrm{arg}p\pi /2`$ as $`p\mathrm{}`$, so that the solutions $`\psi _0^{(s)}(\zeta )`$ and $`\psi _0^{(u)}(\zeta )`$ are defined by (2.15) for all complex $`\zeta `$ with $`\mathrm{Im}\zeta <0`$.
The Borel transform $`V_0(p)`$ satisfies the following integral equation, which follows from (2.14) and (2.15):
$$\left(4\mathrm{sinh}^2\frac{p}{2}c^2p^2\right)V_0(p)=\widehat{Q}\left[V_0(p)\right].$$
(2.17)
Here, $`\widehat{Q}\left[V(p)\right]`$ denotes a double convolution of $`V(p)`$ with itself (in this case, the hat is used to denote an operator). We list below the convolutions $`\widehat{Q}\left[V(p)\right]`$ for each of the four models under consideration:
$`\text{One-site (}\text{1.6}\text{)}:`$ $`2\widehat{Q}`$ $`=V(p)V(p)V(p);`$
$`\text{Speight-Ward (}\text{1.8}\text{)}:`$ $`18\widehat{Q}`$ $`=\mathrm{cosh}p\left[V(p)V(p)V(p)\right]+3\left\{\mathrm{cosh}p\left[V(p)V(p)\right]\right\}`$
$`V`$ $`(p)+3\left[\mathrm{cosh}pV(p)\right]V(p)V(p)+2V(p)V(p)V(p);`$
$`\text{Bender-Tovbis (}\text{1.9}\text{)}:`$ $`2\widehat{Q}`$ $`=\left[\mathrm{cosh}pV(p)\right]V(p)V(p);`$
$`\text{Kevrekidis (}\text{1.11}\text{)}:`$ $`4\widehat{Q}`$ $`=\mathrm{cosh}p\left[V(p)V(p)V(p)\right]`$
$`+\left[\mathrm{cosh}pV(p)\right]\left[\mathrm{}^pV(p)\right]\left[\mathrm{}^pV(p)\right],`$
where the asterisk $``$ denotes the convolution integral for the Borel–Laplace transform:
$$V(p)W(p)=_0^pV(pp_1)W(p_1)p_1,$$
and the integration is performed from the origin to the point $`p`$ on the complex plane, along the contour $`\gamma `$. The inverse power series (2.13) for the limiting solution $`\psi _0(\zeta )`$ becomes the following power series for the Borel transform $`V_0(p)`$:
$$\widehat{V}_0(p)=\underset{n=0}{\overset{\mathrm{}}{}}v_{2n}p^{2n},v_{2n}=\frac{a_{2n}}{(2n)!},$$
(2.18)
where $`v_0=a_0=2\sqrt{1c^2}`$. The hat denotes a formal series which might only converge for some values of $`p`$. The virtue of the integral form (2.17) is that the limiting behavior of $`v_n`$ for large $`n`$ can be related to singularities of $`V_0(p)`$, which in turn correspond to the oscillatory tails of $`\psi _0(\zeta )`$.
If the sliding kink exists, the inverse-power series for $`\psi _0(\zeta )`$ will converge for all $`\zeta `$ such that $`\mathrm{Im}\zeta <0`$. This implies that the stable and unstable solutions $`\psi _0^{(s)}(\zeta )`$ and $`\psi _0^{(u)}(\zeta )`$ coincide, i.e. that the contour $`\gamma _s`$ in the right half of the complex $`p`$-plane can be continously deformed to the contour $`\gamma _u`$ in the left half-plane (see figure 1). If, however, there are any singularities between the two contours, then a continuous deformation is possible only if the residues are zero. The residues are proportional to the values of the Stokes constants. When the Stokes constants are nonzero, the formal power series (2.18) for the solution $`V_0(p)`$ of the integral equation (2.17) diverges for some values of $`p`$ in the sector between the contours $`\gamma _s`$ and $`\gamma _u`$.
### 2.4 Stokes constants
The Borel transform $`V_0(p)`$ is singular near the points in the $`p`$-plane where the coefficient in front of $`V_0(p)`$ on the left-hand side of the integral equation (2.17) vanishes \[T00a\], except for the point $`p=0`$ where the right hand side is also zero. That is, singularities occur when $`(2/p)\mathrm{sinh}(p/2)=\pm c`$. The location of these singularities is important because the stable and unstable solutions are not, in fact, uniquely defined by (2.16); different solutions are generated depending on where the contours lie relative to the singularities of $`V_0(p)`$ with $`\mathrm{Re}p0`$. Exploiting this nonuniqueness, we wish to choose the contours $`\gamma _s`$ and $`\gamma _u`$ to lie above all the singularities with nonzero real part; this will minimise the number of singularities between the stable and unstable solutions.
It is not difficult to show that the contour $`\gamma _s`$ extending from 0 to $`\mathrm{}`$ can be chosen in such a way, i.e. so that that there are no singularities between it and the imaginary axis. Indeed, assume, for definiteness, that $`c>0`$. Let $`(2n_c1)`$ be the number of positive roots of the equation $`\mathrm{sin}q=cq`$ and denote the real and imaginary parts of $`p/2`$ by $`\kappa `$ and $`q`$: $`p/2=\kappa +iq`$. In the $`(\kappa ,q)`$-plane, consider a rectangular region $`𝒟`$ bounded by the horizontal segments $`q=2\pi n`$ and $`q=ϵ`$ at the top and bottom, and vertical segments $`\kappa =ϵ`$ and $`\kappa =ϵ`$ on the left and right. Here $`n`$ is any positive integer greater than $`n_c`$ and $`ϵ>0`$ is taken to be small. Using the argument principle, we can count the number of (complex) roots of the equation $`\mathrm{sinh}(p/2)=c(p/2)`$ in the region $`𝒟`$. We have
$$\mathrm{tan}\mathrm{arg}\frac{p}{2}=\frac{\mathrm{cosh}\kappa \mathrm{sin}qcq}{\mathrm{sinh}\kappa \mathrm{cos}qc\kappa }.$$
On the right lateral side, where $`\kappa =ϵ`$, this becomes
$$\mathrm{tan}\mathrm{arg}\frac{p}{2}\frac{1}{ϵ}\frac{\mathrm{sin}qcq}{\mathrm{cos}qc}.$$
(2.19)
As we move from $`q=ϵ`$ to $`q=2\pi n`$, the numerator in (2.19) will change sign $`(2n_c1)`$ times. In a similar way, moving down along the left side there will be $`(2n_c1)`$ more zero crossings, while no zero crossings will occur along the horizontal segments. This means that the argument can change by no more than $`(4n_c2)\pi `$ and hence there are at most $`(2n_c1)`$ roots in the region $`𝒟`$, no matter how large $`n`$ is. Similarly, we can show that the equation $`\mathrm{sinh}(p/2)=c(p/2)`$ has no more than $`2n_c`$ roots in the region $`𝒟`$, if $`2n_c`$ is the number of positive roots of $`\mathrm{sin}q=cq`$. The upshot is that for any finite $`c`$, there are only a finite number of singularities with small real parts; the singularities cannot accumulate to the imaginary axis. For $`c0`$, the singularities with nonzero real parts lie on the curves
$$q=\pm \sqrt{\frac{1}{c^2}\mathrm{cosh}^2\kappa \kappa ^2\mathrm{coth}^2\kappa }\pm \frac{1}{c}\mathrm{cosh}\kappa \text{as }|\kappa |\mathrm{}.$$
Accordingly, in order for the integration contours $`\gamma _s`$ and $`\gamma _u`$ to lie above these singularities, they must be curvilinear (and not just rays) as shown in figure 1.
Having chosen the contours $`\gamma _s`$ and $`\gamma _u`$ to lie above the singularities in the first and second quadrants respectively, the only singularities of $`V_0(p)`$ that determine whether the stable solution $`\psi _0^{(s)}(\zeta )`$ can be continuously transformed into the unstable solution $`\psi _0^{(u)}(\zeta )`$ are those at non-zero pure imaginary values of $`p`$. We will be referring to these values as resonances. The set of resonances $`_c`$ is defined by the transcendental equation
$$_c=\{p=ik,k_+:\frac{2}{k}\mathrm{sin}\frac{k}{2}=\pm c\}.$$
(2.20)
When $`c=0`$, the set $`_0`$ is infinite-dimensional and can be described explicitly:
$$_0=\left\{p=2\pi n\mathrm{},n\right\}.$$
Let $`p_1=\mathrm{}k_1`$ be the smallest imaginary root in the set $`_c`$. It is clear from (2.20) that $`0<k_1<2\pi `$ for $`c(0,1)`$, so that $`k_12\pi `$ as $`c0^+`$ and $`k_10`$ as $`c1^{}`$. The set of resonances $`_c`$ is finite-dimensional for $`c(0,1)`$ and it consists of only one root $`p_1=\mathrm{}k_1`$ for $`c(c_1,1)`$, where $`c_10.22`$.
Due to the resonances, a function $`\psi _0(\zeta )`$ that satisfies the truncated inner equation (2.14) may have oscillatory tails as $`|\mathrm{Re}\zeta |\mathrm{}`$. Adding the solutions of equation (2.14) linearised about $`\widehat{\psi }_0(\zeta )`$, the general bounded solution of (2.14) in the limit $`|\zeta |\mathrm{}`$ can be represented as \[TTJ98\]:
$`\psi _0(\zeta )=\widehat{\psi }_0(\zeta )+{\displaystyle \underset{\mathrm{}k_n_c}{}}\alpha _n\widehat{\phi }_n(\zeta )\mathrm{}^{\mathrm{}k_n\zeta }+\text{multiple harmonics}.`$ (2.21)
Here, $`\widehat{\psi }_0(\zeta )`$ is given by the power series (2.13); $`\alpha _n`$ are coefficients which we will be referring to as amplitudes in what follows; $`k_n>0`$ are roots of (2.20) for $`p=\mathrm{}k_n`$, and the functions $`\widehat{\phi }_n(\zeta )\mathrm{}^{\mathrm{}k_n\zeta }`$, $`n1`$, satisfy the linearised truncated inner equation (2.14). In particular, the equation for the leading-order term $`\widehat{\phi }_1(\zeta )`$ is
$`\mathrm{}^{\mathrm{}k_1}\widehat{\phi }_1(\zeta +1)+(c^2k_1^22)\widehat{\phi }_1(\zeta )+\mathrm{}^{\mathrm{}k_1}\widehat{\phi }_1(\zeta 1)+2\mathrm{}c^2k_1\widehat{\phi }_1^{}(\zeta )c^2\widehat{\phi }_1^{\prime \prime }(\zeta )`$
$`=D_1Q\widehat{\phi }_1(\zeta 1)+D_2Q\widehat{\phi }_1(\zeta )+D_3Q\widehat{\phi }_1(\zeta +1),`$ (2.22)
where $`D_{1,2,3}Q`$ are the partial derivatives of $`Q(\psi _0(\zeta 1),\psi _0(\zeta ),\psi _0(\zeta +1))`$ with respect to its first, second and third argument respectively, evaluated at $`\psi _0=\widehat{\psi }_0(\zeta )`$.
If the amplitude $`\alpha _n`$ is nonzero for some $`n`$, the formal power series (2.13) does not converge because the solution (2.21) does not decay as $`|\mathrm{Re}\zeta |\mathrm{}`$. The amplitudes $`\alpha _n`$ are proportional to the Stokes constants computed for the formal power series (2.13). Each oscillatory term in the sum (2.21) becomes exponentially small in $`h`$ when we transform from $`\zeta `$ to $`z`$ using the transformation (2.10). Since $`p_1=\mathrm{}k_1`$ is the element of $`_c`$ with the smallest imaginary part, it follows that the $`n=1`$ term dominates the sum in (2.21) when the transformation (2.10) is made (unless $`\alpha _1=0`$). Furthermore, when $`c(c_1,1)`$, where $`c_10.22`$, it is the *only* term in the sum since the resonant set $`_c`$ consists of just the one root $`p_1=\mathrm{}k_1`$. We shall, therefore, only be concerned with the leading-order Stokes constant, which multiplies the function $`\widehat{\phi }_1(\zeta )`$.
If $`\widehat{\psi }_0(\zeta )`$ is given by the power series (2.13), the solution of the linearized equation (2.22) can also be represented by a formal power series:
$$\widehat{\phi }_1(\zeta )=\zeta ^r\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}b_{\mathrm{}}\zeta ^{\mathrm{}},$$
(2.23)
where we can set $`b_0=1`$ due to the linearity of (2.22). Substituting (2.13) and (2.23) into (2.22) and using (2.20), the coefficient $`b_1`$ can be determined from
$`2\mathrm{}r\zeta ^{r1}(c^2k_1\mathrm{sin}k_1)+\zeta ^{r2}[r(r1)(\mathrm{cos}k_1c^2)`$
$`+2\mathrm{}b_1(r1)(c^2k_1\mathrm{sin}k_1)6(1c^2)]+𝒪(\zeta ^{r3})=0.`$
In this equation, the coefficient of each power of $`\zeta `$ should be set to zero. In order to set the coefficent in front of the first term to zero in the situation where $`c0`$, we must choose $`r=0`$. The second term then gives
$$b_1=\frac{3\mathrm{}(1c^2)}{c^2k_1\mathrm{sin}k_1},$$
after which all the other coefficients $`b_2`$, $`b_3`$, …, can be computed recursively. On the other hand, in the situation with $`c=0`$, we have $`k_1=2\pi `$ and the coefficient in front of $`\zeta ^{r1}`$ is zero regardless of the value of $`r`$. Setting the coefficient in front of $`\zeta ^{r2}`$ to zero requires that we choose either $`r=3`$ or $`r=2`$, and hence we have two different descending-power series, one starting with $`\zeta ^3`$ and the other one with $`\zeta ^2`$. We shall focus on the former as it dominates the latter in the limit $`\zeta \mathrm{}`$. Again, the succeeding terms in (2.23) are determined recursively.
Thus, we have established that the leading-order oscillatory term in the expansion (2.21) behaves as
$$\{\begin{array}{cc}\alpha _1\left[1+\frac{b_1}{\zeta }+𝒪\left(\frac{1}{\zeta ^2}\right)\right]\mathrm{}^{\mathrm{}k_1\zeta }\hfill & \text{for }c0\text{ and}\hfill \\ \alpha _1\left[\zeta ^3+b_1\zeta ^2+𝒪\left(\zeta ^1\right)\right]\mathrm{}^{2\pi \mathrm{}\zeta }\hfill & \text{for }c=0\text{.}\hfill \end{array}$$
(2.24)
For $`c0`$, the two leading order terms in the expression above are generated by, respectively, a simple pole and a logarithmic singularity of the Borel transform $`V_0(p)`$ at $`p=\mathrm{}k_1`$. For $`c=0`$ they are generated by a quadruple pole of $`V_0(p)`$ at $`p=2\pi \mathrm{}`$. From the fact that $`V_0`$ is an even function of $`p`$, we deduce the structure of this function near the poles:
$$V_0(p)\{\begin{array}{cc}\frac{k_1^2K_1(c)}{p^2+k_1^2}\sigma (c)\mathrm{ln}(p^2+k_1^2)+\mathrm{}\hfill & \text{(for }c0\text{)}\hfill \\ \frac{6(2\pi )^8S_1}{\left(p^2+4\pi ^2\right)^4}+\frac{2(2\pi )^6\rho }{\left(p^2+4\pi ^2\right)^3}+\mathrm{}\hfill & \text{(for }c=0\text{)}\hfill \end{array}$$
(2.25)
as $`p\pm \mathrm{}k_1`$. Here $`K_1(c)`$ and $`S_1`$ are the leading-order Stokes constants for $`c0`$ and $`c=0`$, respectively; $`\sigma (c)`$ and $`\rho `$ are independent of $`p`$, and $`\mathrm{}`$ stands for terms with even slower growth as $`pik_1`$.
To show that these singularities do indeed give rise to the oscillatory tails in (2.24), we compare the two integrals $`\psi _0^{(s)}`$ and $`\psi _0^{(u)}`$ for a given value of $`\zeta `$. To this end, we deform the paths of integration $`\gamma _s`$ and $`\gamma _u`$ to $`\gamma _s^{}`$ and $`\gamma _u^{}`$ respectively, without crossing any singularities. This is illustrated in figure 1. There are two contributions to the difference $`\psi _0^{(s)}(\zeta )\psi _0^{(u)}(\zeta )`$. The first comes from integrating around the pole at $`p=\mathrm{}k_1`$, and is equal to $`2\pi \mathrm{}`$ times the residue of the function $`V_0(p)\mathrm{}^{p\zeta }`$ at $`p=\mathrm{}k_1`$, determined from (2.25). The second contribution (manifest only in the $`c0`$ case) arises because the *integrand* increases as the singularity is encircled, since it is a branch point of the logarithm. Since $`\mathrm{ln}z`$ can be written as $`\mathrm{ln}|z|+\mathrm{}\mathrm{arg}z`$, where $`z=p\mathrm{}k_1`$, we see that $`V_0(p)`$ increases by $`2\pi \mathrm{}\sigma (c)`$ as the branch point $`p=\mathrm{}k_1`$ is encircled in the $`c0`$ case. Therefore, the difference in the integrand of (2.15) along the paths $`\gamma _s^{}`$ and $`\gamma _u^{}`$ is $`2\pi \mathrm{}\sigma (c)\mathrm{}^{p\zeta }`$, which must be integrated along the path of integration from $`p=\mathrm{}k_1`$ to infinity, to give $`2\pi \mathrm{}\sigma (c)\mathrm{}^{\mathrm{}k_1\zeta }/\zeta `$. (We have considered the integration on a Riemann surface in order to account for branch points.) Adding together the two contributions discussed above, we have
$$\psi _0^{(s)}(\zeta )\psi _0^{(u)}(\zeta )=\{\begin{array}{cc}\left[\pi k_1K_1(c)\frac{2\pi \mathrm{}\sigma (c)}{\zeta }+𝒪\left(\frac{1}{\zeta ^2}\right)\right]\mathrm{}^{\mathrm{}k_1\zeta }\hfill & \text{for }c0\hfill \\ \frac{1}{128}[16\pi ^3\mathrm{}S_1\zeta ^3\hfill & \\ +(192\pi ^4S_1+\rho )\zeta ^2+𝒪(\zeta ^1)]\mathrm{}^{2\pi \mathrm{}\zeta }\hfill & \text{for }c=0\text{.}\hfill \end{array}$$
(2.26)
If we take the limit $`\mathrm{Re}\zeta \mathrm{}`$, then the unstable solution $`\psi _0^{(u)}(\zeta )`$ decays to zero as a power law, according to the expansion (2.13). Thus, the stable solution $`\psi _0^{(s)}(\zeta )`$ has the oscillatory tail given by the representation (2.21) with the amplitude factor
$$\alpha _1=\{\begin{array}{cc}\pi k_1K_1(c)\hfill & \text{for }c0\hfill \\ \frac{\mathrm{}\pi ^3}{8}S_1\hfill & \text{for }c=0\text{.}\hfill \end{array}$$
(2.27)
Similarly, if we take the limit $`\mathrm{Re}\zeta +\mathrm{}`$, then the stable solution $`\psi _0^{(s)}(\zeta )`$ decays to zero, while the unstable solution $`\psi _0^{(u)}(\zeta )`$ has the representation (2.21) with the amplitude factor given by the negative of expression (2.27). By comparing the other terms on the right-hand side of (2.26) to the corresponding terms in (2.24), $`\sigma (c)`$ and $`\rho `$ can be uniquely determined.
We now match the leading-order singular behaviour of $`V_0(p)`$ near $`p=\pm \mathrm{}k_1`$, given by (2.25), to the formal power series (2.18). Expanding the expressions in (2.25) as power series gives us
$$V_0(p)\{\begin{array}{cc}K_1(c)\sigma (c)\mathrm{ln}(k_1^2)+\mathrm{}\hfill & \\ +_{n=1}^{\mathrm{}}(1)^nk_1^{2n}\left(K_1(c)+\frac{\sigma (c)}{n}+\mathrm{}\right)p^{2n}\hfill & \text{for }c0\hfill \\ _{n=0}^{\mathrm{}}\frac{(1)^n(n+2)(n+1)}{(2\pi )^{2n}}\left[(n+3)S_1+\rho +\mathrm{}\right]p^{2n}\hfill & \text{for }c=0\text{,}\hfill \end{array}$$
(2.28)
as $`p\pm \mathrm{}k_1`$. These series converge for all $`|p|<k_1`$; in particular, they are valid for $`p\pm \mathrm{}k_1`$, provided $`|p|<k_1`$. Hence we can replace (2.25) with (2.28) in this neighbourhood. In (2.28), the ellipses stand for coefficients of the expansion of terms with a slower growth as $`p\pm \mathrm{}k_1`$ which were dropped in (2.25). The discarded terms would modify the coefficients of the power series (2.28); however, there are terms which would not be affected by these modifications, namely terms with large $`n`$. For example, the coefficients proportional to $`\sigma (c)`$ and $`\rho `$ in (2.28) are a factor of $`n`$ smaller than those proportional to $`K_1(c)`$ and $`S_1`$; the discarded coefficients would be even smaller. Therefore the leading singular behaviour of $`V_0(p)`$ as $`p\pm \mathrm{}k_1`$ is determined just by the large-$`n`$ coefficients of the power series (2.28), and hence only the large-$`n`$ coefficients should be matched to the coefficients of the expansion (2.18). This gives the Stokes constant as a limit of the coefficients $`v_{2n}`$ of the series (2.18):
$$\{\begin{array}{cc}K_1(c)=\underset{n\mathrm{}}{lim}(1)^nk_1^{2n}v_{2n}\hfill & \text{for }c0\hfill \\ S_1=\underset{n\mathrm{}}{lim}\frac{(1)^n(2\pi )^{2n}v_{2n}}{(n+3)(n+2)(n+1)}\hfill & \text{for }c=0\text{.}\hfill \end{array}$$
(2.29)
This formula is used in the next section for numerical computations of the leading-order Stokes constant $`K_1(c)`$ for $`c0`$.
Note that, since (2.18) matches (2.28) in the limit $`n\mathrm{}`$, the formal power series $`\widehat{V}_0(p)`$ also has radius of convergence $`k_1`$. However, the formal inverse-power series $`\widehat{\psi }_0(\zeta )`$ diverges for all $`\zeta `$ unless $`\widehat{V}_0(p)`$ converges everywhere (which requires that all the Stokes constants be zero).
Next, we note that as $`c0`$, the Stokes constant $`K_1(c)`$ does not tend to $`S_1`$, its value at $`c=0`$. This discontinuity is due to the fact that, as $`c0`$, pairs of simple roots in the resonant set $`_c`$ coalesce. (E.g. $`\mathrm{}k_1`$ coalesces with $`\mathrm{}k_2`$ at $`2\pi \mathrm{}`$, and so on.) As a result, all roots are double and the representation of $`\widehat{\phi }_1(\zeta )`$ is discontinuous at $`c=0`$, with the power degree $`r`$ of the prefactor in (2.23) jumping from $`r=0`$ for $`c0`$ to $`r=3`$ for $`c=0`$. In particular, in exceptional models, i.e., discrete models with continuous families of stationary kinks (like (1.8), (1.9) and (1.11)) the constant $`S_1`$ is *a priori* zero while the limit of $`K_1(c)`$ as $`c0`$ may be nonvanishing. In fact, numerical computations of the top limit in (2.29) indicate that the Stokes constant blows up as $`c0`$. Renormalisation of $`K_1(c)`$ for small $`c`$ is, however, a nontrivial asymptotic problem which is beyond the scope of our current investigation.
For $`c(c_1,1)`$, where $`c_10.22`$, the resonant set $`_c`$ contains only one root $`\mathrm{}k_1`$ and, therefore, there is just one Stokes constant $`K_1(c)`$, which completely determines the convergence of the formal power series for $`\widehat{\psi }_0(\zeta )`$. If $`K_1(c_0)=0`$ at some point $`c_0(c_1,1)`$, the stable and unstable solutions $`\psi _0^{(s)}(\zeta )`$ and $`\psi _0^{(u)}(\zeta )`$ coincide to leading order. Arguments based on the implicit function theorem (see \[TP05\]) reveal a heteroclinic bifurcation which occurs on crossing a smooth curve $`c=c_{}(h)`$ on the $`(c,h)`$-plane, with $`c_{}(0)=c_0`$.
On the other hand, for $`c(0,c_1)`$ the resonant set $`_c`$ contains more than one root. If $`K_1(c_0)=0`$ for some $`c_0(0,c_1)`$, this alone is not sufficient for the convergence of the formal power series $`\widehat{\psi }_0(\zeta )`$. The higher-order Stokes constants $`K_2(c)`$, $`K_3(c)`$, …, must be introduced and computed from the asymptotic behavior of the power series $`\widehat{V}_0(p)`$.
As we shall show in the next section, the function $`K_1(c)`$ does have zeros in the case of the discretizations (1.8) and (1.11). All these zeroes are “safe”; that is, all $`c_0`$ values lie in the interval $`(c_1,1)`$, so that the higher-order Stokes constants do not have to be computed.
## 3 Numerical computations of the Stokes constant
In this section, we report on the numerical computation of the Stokes constants $`K_1(c)`$ for the four different discretizations of the $`\varphi ^4`$ model (1.3) under consideration. Our numerical method utilises the expression (2.29) of the Stokes constant in terms of the coefficients of the formal power series solution (2.18). First, we obtain the recurrence relation for the coefficients in the set $`\{v_n\}_{n=0}^{\mathrm{}}`$ by substituting the power series expansion (2.18) into the limiting integral equation (2.17), and using the convolution formula
$$p^np^m=\frac{n!m!}{(n+m+1)!}p^{n+m+1}.$$
(3.1)
After that, we compute the asymptotic behavior of these coefficients as $`n\mathrm{}`$ and evaluate the limit (2.29) numerically for a fixed value of $`c0`$.
In order to calculate the Stokes constant for the four models in a uniform way, we write a general symmetric homogeneous cubic polynomial $`Q(u_{n1},u_n,u_{n+1})`$ as
$$Q=\underset{\alpha =1}{\overset{1}{}}\underset{\beta =\alpha }{\overset{1}{}}\underset{\gamma =\beta }{\overset{1}{}}a_{\alpha ,\beta ,\gamma }u_{n+\alpha }u_{n+\beta }u_{n+\gamma },$$
(3.2)
where $`a_{\alpha ,\beta ,\gamma }`$ are numerical coefficients, with $`\alpha ,\beta ,\gamma \{1,0,1\}`$ and $`\alpha \beta \gamma `$. The symmetry implies that $`a_{\alpha ,\beta ,\gamma }=a_{\gamma ,\beta ,\alpha }`$ and therefore it is sufficient to specify just six coefficients. The values of these coefficients for the four nonlinearities in question are given in Table 1.
By applying the Borel–Laplace transform (2.15) to equation (2.14) with $`Q`$ as in (3.2), we obtain the corresponding cubic convolution function $`\widehat{Q}[V(p)]`$ on the right hand side of the integral equation (2.17):
$$\widehat{Q}\left[V(p)\right]=\underset{\alpha =1}{\overset{1}{}}\underset{\beta =\alpha }{\overset{1}{}}\underset{\gamma =\beta }{\overset{1}{}}a_{\alpha ,\beta ,\gamma }\mathrm{}^{\alpha p}V(p)\mathrm{}^{\beta p}V(p)\mathrm{}^{\gamma p}V(p).$$
(3.3)
To derive the recurrence formula for the coefficents $`v_{2n}`$ in (2.18), it will be more convenient to consider the power series expansion which consists of both even and odd powers of $`p`$:
$$\widehat{V}_0(p)=\underset{n=0}{\overset{\mathrm{}}{}}v_np^n.$$
(3.4)
We now substitute the series (3.4) into (2.17) with $`\widehat{Q}\left[V(p)\right]`$ given by (3.3) and use the convolution formula (3.1). Equating the coefficients of $`p^{n+2}`$ where $`n=0,1,2,\mathrm{}`$, in the resulting equation, we find that
$`{\displaystyle \underset{i=0}{\overset{[n/2]}{}}}{\displaystyle \frac{2}{(2i+2)!}}v_{n2i}c^2v_n={\displaystyle \underset{\alpha =1}{\overset{1}{}}}{\displaystyle \underset{\beta =\alpha }{\overset{1}{}}}{\displaystyle \underset{\gamma =\beta }{\overset{1}{}}}{\displaystyle \frac{a_{\alpha ,\beta ,\gamma }}{(n+2)(n+1)}}\{{\displaystyle \underset{i=0}{\overset{n}{}}}\left({\displaystyle \underset{k=0}{\overset{ni}{}}}{\displaystyle \frac{\alpha ^k}{k!}}v_{nik}\right)`$
$`\times \left[{\displaystyle \underset{j=0}{\overset{i}{}}}\left({\displaystyle \underset{l=0}{\overset{j}{}}}{\displaystyle \frac{\beta ^l}{l!}}v_{jl}\right)\left({\displaystyle \underset{m=0}{\overset{ij}{}}}{\displaystyle \frac{\gamma ^m}{m!}}v_{ijm}\right){\displaystyle \frac{j!(ij)!}{i!}}\right]{\displaystyle \frac{i!(ni)!}{n!}}\},`$ (3.5)
where $`[n/2]`$ is the integer part of $`n/2`$ and $`0^0=1`$. Equation (3) is a recurrence relation between the coefficients $`\{v_n\}_{n=0}^{\mathrm{}}`$. Solving equation (3) with $`n=0`$, we get $`v_0=2\sqrt{1c^2}`$. Note that this result is independent of the choice of $`a_{\alpha ,\beta ,\gamma }`$, i.e. independent of the model. Letting $`v_1=0`$ and making use of the symmetry of $`Q`$, one can show by induction that the coefficients of all odd powers in (3.4) are zero (as we concluded previously on the basis that the outer expansion is odd).
To prevent overflow or underflow when evaluating the recurrence relation numerically, we shall work with the normalised coefficients
$$w_n=(1)^nk_1^{2n}v_{2n}$$
so that the Stokes constant (2.29) for $`c0`$ is given by
$$K_1(c)=\underset{n\mathrm{}}{lim}w_n.$$
(3.6)
Reformulating (3) in terms of $`w_n`$, we use the relation (3.6) to compute $`w_n`$ numerically. We truncate the sums involving $`1/(2i+2)!`$, $`1/l!`$ and $`1/m!`$ when these factors become smaller than $`10^{50}`$, and evaluate the sums involving the combinatorial factors in two halves. In the first, the summation index increases from zero to the halfway point, and in the second it decreases from its maximum. This ensures that the combinatorial factors are always decreasing from one step to the next so that they can be accurately determined recursively. We also truncate these sums when the combinatorial factors fall below $`10^{50}`$.
These expedients result in a numerical routine fast enough to allow for evaluation of the recurrence relation up to very large $`n`$; this is essential given the slow convergence of $`w_n`$ to a constant. Matching (2.18) to (2.28) yields
$$v_{2n}(1)^nk_1^{2n}[K_1(c)+\sigma (c)/n]\text{as }n\mathrm{};$$
therefore, the rate at which $`w_n`$ converges to $`K_1(c)`$ is of order $`1/n`$:
$$w_n=K_1(c)+\frac{\sigma (c)}{n}+\frac{\stackrel{~}{\sigma }(c)}{n^2}+𝒪\left(\frac{1}{n^3}\right).$$
(3.7)
Although the convergence of $`w_n`$ to $`K_1(c)`$ is extremely slow, we can accelerate the process by using (3.7). Defining
$$\stackrel{~}{w}_nw_n+n(w_nw_{n1}),$$
we get
$$\stackrel{~}{w}_n=K_1(c)\frac{\sigma (c)+\stackrel{~}{\sigma }(c)}{n^2}+𝒪\left(\frac{1}{n^3}\right).$$
The convergence of the sequence $`\stackrel{~}{w}_n`$ is much faster than that of $`w_n`$; see Figure 2. The relative error
$$E(n)=\frac{\sigma (c)+\stackrel{~}{\sigma }(c)}{n^2}\frac{1}{K_1(c)}$$
can be written as
$$E(n)=\frac{n}{2}\frac{\stackrel{~}{w}_n\stackrel{~}{w}_{n1}}{\stackrel{~}{w}_n}$$
plus terms of order $`1/n^4`$. This gives an empirical criterion for the termination of the process. We continued our computations until $`E(n)`$ reached a value smaller than $`10^3`$, i.e. until the percentage error dropped below 0.1%. For $`c>0.5`$, the value of $`n`$ to which we have to compute in order to achieve this accuracy is less than $`100`$, increasing for smaller values of $`c`$ to approximately $`\mathrm{5\hspace{0.17em}000}`$ for $`c=0.005`$. Consequently, the above numerical algorithm is not suited to the study of the $`c0`$ limit, and would have to be modified for that purpose.
Figure 3 displays the Stokes constant computed using the above numerical procedure, for the four models of Table 1. We see that the Stokes constant $`K_1(c)`$ vanishes almost nowhere in $`c0`$ in all four models. There are, however, several isolated zeros: $`K_1(c)=0`$ for $`c_00.45`$ in the case of the Speight-Ward nonlinearity (1.8) and for $`c_00.37`$, $`0.63`$ and $`0.83`$ in the case of the Kevrekidis discretization (1.11). Importantly, all of these lie in the region $`(c_1,1)`$ where the resonance set (2.20) consists of only one value, $`p_1=\mathrm{}k_1`$. (Here $`c_10.22`$.) Therefore, there is a sliding kink in the $`h0`$ limit for each of these isolated values of velocity. Furthermore, strong parallels between our current setting and that of solitons of the fifth-order KdV equation \[TP05\] suggest that sliding kinks should exist along a curve on the $`(c,h)`$ plane emanating from each of the points $`(c_0,0)`$. In other words, we conjecture that there is a radiationless kink travelling with a certain particular speed $`c_{}(h)`$ for each $`h`$ in the case of the Speight-Ward nonlinearity, and that there are three such velocities (for each $`h`$) in the case of the Kevrekidis model. For small $`h`$, $`c_{}(h)`$ should be close to the above values $`c_0`$.
In order to verify the existence of kinks sliding at these isolated velocities by an independent method, we solved the differential advance-delay equation (2.2) numerically. The infinite line was approximated by an interval of length $`2L=200`$, with the antiperiodic boundary conditions $`\varphi (L)=\varphi (L)`$. We utilised Newton’s iteration with an eighth-order finite-difference approximation of the second derivative; the step size was chosen to be $`h/10`$. The continuum solution (1.2) was used as an initial guess.
If we find a solution to the advance-delay equation with $`\varphi (z)`$ decaying to a constant for large positive and negative $`z`$, then we regard this solution as (a numerical approximation to) a radiationless travelling kink. We were able to tune $`c`$ for a fixed value of $`h`$ so that the radiation was reduced to the order of $`10^{12}`$, whereupon the finite accuracy of our numerical scheme prevented any further reduction. To make sure that the radiation does vanish rather than reaching a local minimum but remaining nonzero, we plot the average magnitude of the radiation near the ends of the interval as a function of $`c`$, for fixed $`h`$. This is defined as the average of $`\left[\varphi (z)\overline{\varphi }\right]^2`$ over the last 20 units of the interval, where $`\overline{\varphi }`$ is the average value of $`\varphi (z)`$ over these last 20 units. The results are shown in figure 4. Note the straight-line behavior of the graphs near the isolated values of $`c`$; this indicates that the coefficient of the sinusoid superimposed over the kink’s flat asymptote crosses through zero (rather than attaining a small but nonzero minimum). The supression of radiation at the isolated points is thereby verified.
Finally, the last question that we would like to address here is whether the intensity of the radiation from the moving discrete kink depends on the type of discretization. More specifically, we would like to know whether the choice of one of the exceptional discretizations (which, by definition, support translationally invariant stationary kinks) serves to reduce the radiation from the moving kinks. Speight and Ward have already given an affirmative answer for their exceptional discretization; here we consider the one-parameter nonlinearity
$$Q=\frac{(1\mu )}{2}u_n^3+\frac{\mu }{4}u_n^2(u_{n+1}+u_{n1}),$$
(3.8)
which interpolates between the one-site nonlinearity (1.6) (for which $`\mu =0`$) and the exceptional discretization (1.9) of Bender and Tovbis (for which $`\mu =1`$). Figure 5 shows the Stokes constant for the model (3.8), as $`\mu `$ changes from 0 to 1 for fixed values of $`c`$. The Stokes constant is indeed seen to be drastically reduced as $`\mu `$ approaches $`1`$ — that is, in the limit of the exceptional discretization. (It nonetheless remains nonzero, of course, unless $`c=0`$.)
## 4 Concluding remarks and conclusions
In this paper we have investigated the existence of sliding kinks — i.e. discrete kinks travelling at a constant velocity over a flat background, without emitting any radiation — in four discrete versions of the quartic-coupling theory. One of these models is the most common, one-site, discretization. As the overwhelming majority of discrete $`\varphi ^4`$-equations, it supports travelling kinks, but these kinks do radiate and decelerate as a result. The other three discretizations we considered are all exceptional in the sense that they all support one-parameter continuous families of stationary kinks where the free parameter defines the position of the kink relative to the lattice. This property is clearly nongeneric; the translation invariance of the continuous $`\varphi ^4`$-theory is broken by the discretization and hence in generic disretizations kinks may only be centred at a site or midway between two sites. Since the nonexistence of “translationally-invariant” and of sliding kinks in the generic models can be ascribed to similar factors, viz., the breaking of the translation and Lorentz invariances, it was hoped that the exceptional discretizations might turn out to be equally exceptional from the point of view of sliding kinks. Our approach was based on the computation of the Stokes constants associated with the putative sliding kink in a given equation.
The main conclusion of our work is that the sliding kinks do exist in the discrete $`\varphi ^4`$ theories, but only with special, isolated, velocities (which of course depend on $`h`$). There is one such velocity in the exceptional model of Speight and Ward, and three different sliding velocities in the discretization of Kevrekidis. It is natural to expect that the sliding kinks should play the role of attractors similarly to the fronts moving with “stable velocities” in dissipative systems; that is, radiating travelling kinks should evolve into kinks travelling with the sliding velocities — if there are such velocities in the system. Not every discretization supports sliding kinks, of course; in particular, no sliding velocities arise for the generic, one-site, nonlinearity and even for the exceptional discretization of Bender and Tovbis.
One natural way of trying to construct the sliding kinks is via power series expansions in powers of $`c^2`$; for the exceptional discretizations, this construction can be carried out to any order. This approach was pursued in the recent work of Ablowitz and Musslimani \[AM03\]. Our results indicate, however, that these power series will not converge and exponentially small terms (terms lying beyond all orders of the power expansion) emerge because of the singular behaviour of the Stokes constant $`K_1(c)`$ as $`c0`$. Detailed studies of this singular limit will be presented elsewhere.
The exceptional discretizations have richer underlying symmetries than generic nonlinearities but the “translation invariance” of the stationary kink alone does not automatically guarantee the existence of the sliding velocities. The exact relation between the “translational invariance” and mobility of kinks is still to be clarified; at this stage it is worth mentioning that the Stokes constant associated with (and hence the intensity of radiation from) a moving kink is several orders of magnitude smaller in exceptional models than in generic discretizations.
Finally, it is instructive to point out some parallels with an earlier work of Flach, Zolotaryuk and Kladko \[FZK99\] who also studied the phenomenon of kink sliding in Klein-Gordon lattices. In the scheme of \[FZK99\], one postulates an analytic expression for the sliding kink, $`u_n(t)=\varphi (ncts)`$, with some explicit function $`\varphi (z)`$, and then reconstructs the Klein-Gordon nonlinearity for which this is an exact solution. Our present conclusions are in agreement with the results of these authors who observed that for a given $`h`$, the kink may only slide at a particular, isolated, velocity. The two approaches, ours and that of \[FZK99\], are reciprocal; while we examine the existence of sliding kinks for particular discretizations of the $`\varphi ^4`$-theory, with fixed parameters independent of the kink’s velocity, in the “inverse method” of \[FZK99\] one assumes an explicit solution of a particular form but does not have any control over the resulting nonlinearities. Consequently, the discrete Klein-Gordon models generated by the “inverse method” are not discretizations of the $`\varphi ^4`$-theory and do include explicit dependence on the velocity of the sliding kink.
O.O. was supported by funds provided by the South African government and the University of Cape Town. D.P. thanks the Department of Mathematics at UCT for hospitality during his visit and the NRF of South Africa for financial support which made the visit possible. I.B. is a Harry Oppenheimer Fellow; also supported by the NRF under grant 2053723.
## References
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# Tradeoff Between Source and Channel Coding for Erasure Channels
## I Introduction
In Shannon presented his celebrated result on the asymptotic optimality of separable source and channel coding. However, for finite block length systems, the importance and superior performance of joint source channel coders has been well recognized and is an area of active research (see for example ). Specific research thrusts have included investigating source coders that incorporate channel information in the design, channel coders that provide unequal error protection to various source bits, and iterative source-channel decoders.
An important issue in joint source channel coding is the tradeoff between source and channel coding rates. For a fixed source vector dimension and channel capacity, there is a tradeoff between the source and channel coding rates. A high rate channel code implies more bits for the source coder which results in a high quality representation at the source but has a higher probability of being received in error. Similarly, a low rate channel code results in fewer bits for the source coder; consequently, the representation is of lower quality at the source but there is a higher probability of being received without error at the receiver. This tradeoff has been quantified for binary symmetric channels (BSC) and Gaussian channels .
This paper addresses the problem of optimal allocation of rate between a source encoder and a channel encoder for transmission over erasure channels. The system under investigation is a concatenation of a vector quantizer with a channel coder and the objective is to minimize the end-to-end distortion. Upper and lower bounds on the channel coding rate are constructed that minimizes the end-to-end distortion. The upper bound on channel coding rate is derived using the sphere packing and straight line exponents as a bound for performance of the channel code. Similarly, the lower bound on rate is derived based on the expurgated error exponent for the erasure channel. The proposed bounds suggest that the optimal channel coding rate is substantially smaller than the channel capacity. Asymptotically, as the erasure probability $`ϵ0`$, the optimal channel coding rate equals 1. The resulting upper and lower bounds are then adapted to obtain the optimal coding rate for packet erasure channels.
The closed form approximations for the optimal coding rate are derived under the assumption of asymptotically small erasure probabilities. Also, a high rate quantization regime is considered and hence, distortion achieved asymptotically with large k equals the rate distortion bound. The proposed bounds are independent of the source distribution for sources with fixed dimensionality and finite support.
The rest of the paper is organized as follows. In Section 2, we define the system and present the notation and assumptions made in this paper. In Section 3 we evaluate the upper and lower bounds on the rate for erasure channels using the expurgated, sphere packing and straight line bounds. In Section 4 we present some numerical results and conclude in Section 5.
A model of the communication system under investigation is given in Fig. 1. Consider a random vector $`\text{X}\mathrm{}^k`$ that has a probability density function f over support set $`A`$, a closed bounded subset of $`\mathrm{}^k`$ with nonempty interior. Let X be quantized by a vector quantizer $`Q:\mathrm{}^kC`$ where $`C=[\text{y}_1,\text{y}_2,\mathrm{},\text{y}_M]`$ is the codebook of the vector quantizer with $`m=\mathrm{log}M`$ bits per source symbol. All logarithms are to base 2. Consequently, the quantizer can be modeled as ,
$$Q\left(\text{x}\right)=\underset{i=1}{\overset{M}{}}\text{y}_i\text{1}_{S_i}\left(\text{x}\right)$$
(1)
where $`\{S_i\}_{i=1}^M`$ is a partition of $`\mathrm{}^k`$ into disjoint regions, each of which is represented by $`\text{y}_i`$ and $`\text{1}_{S_i}(.)`$ is the indicator function which equals 1 if x lies in the $`i^{th}`$ cell of the partition. The average distortion using this quantizer is given by,
$`D_m\left(Q\right)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{M}{}}}{\displaystyle _{S_i}}\text{x}\text{y}_i^pf\left(x\right)𝑑x`$ (2)
$`=`$ $`2^{pRr+O\left(1\right)},`$ (3)
where (3) follows from Zador’s distortion formula and $`p`$ is the power of the distortion measure. Traditional quantization theory has worked on computing optimal quantizers that achieve the infimum of $`D_m\left(Q\right)`$. The quantizer design has been expanded to include the effect of channel errors; however, the problem is extremely challenging and little analytical results are known in such cases.
In our system, we consider that the $`m`$-bit source codewords are first randomly permuted using a mapping $`\pi `$ and then passed to a channel encoder of rate $`r=m/n`$ before transmission over a binary erasure channel with erasure probability $`ϵ`$. For simplicity, we have included the index assignment $`\pi `$ as part of the source encoder. The results are independent of the index assignment. The channel encoder generates a unique $`n`$-bit channel codeword for each of the $`m`$-bit source codewords. The added redundancy $`nm`$ is used to protect the source code word from channel impairments. The transmission rate per source component is $`R=n/k`$, and the quantization rate is $`R_s=m/k=Rr`$. Following the notations used in , we denote $`a_i=O\left(b_i\right)`$ if $`|a_i|/b_ic`$ for some $`c>0`$ and $`i`$ sufficiently large. We denote $`a_i=\mathrm{\Omega }\left(b_i\right)`$ if $`|a_i|/b_ic`$ for some $`c>0`$ and sufficiently large $`i`$. Finally, $`a_i=o\left(b_i\right)`$ if $`lim_i\mathrm{}a_i/b_i=0`$.
## II Binary Erasure Channel
We now consider obtaining bounds on the coding rate for a binary erasure channel (BEC). The end-to-end distortion for the system in Fig. 1, is readily given by
$$D_R(Q,ϵ)=\underset{i,j=1}{\overset{M}{}}q\left(j|i\right)_{S_i}\text{x}\text{y}_i^pf\left(x\right)𝑑x$$
(4)
where $`ϵ`$ is the bit erasure probability and $`q\left(j|i\right)`$ is the conditional probability that the channel decoder decides in favor of the $`j^{th}`$ channel codeword when the $`i^{th}`$ codeword was transmitted.
### II-A Lower bound on channel coding rate
The lower bound on the channel coding rate is obtained by upper bounding the distortion at the decoder. Assuming small bit erasure probability and following , the total distortion may be upper bounded as,
$$D_R(Q,ϵ)D_m\left(Q\right)+O(1)\underset{1iM}{\mathrm{max}}P_{e|i}$$
(5)
In (5), the total distortion is the sum of the distortion due to the vector quantizer ( $`D_m\left(Q\right)`$), and the distortion due to the errors in transmission. The positive O(1) term is due to the fact that $`f`$ has support $`A`$, and $`y_j`$ is contained in $`A`$ for all $`j`$ . The problem of interest is posed as follows: “Given a binary erasure channel with $`R`$ and source $`\text{X}\mathrm{}^k`$, find the optimal rate $`r`$ that minimizes the distortion $`D_R(Q,ϵ,)`$
For an arbitrary binary discrete memoryless channel, Shannon’s channel coding theorem guarantees that for channel code rates $`r`$ below capacity, the probability of error is upper bounded by
$$\underset{1iM}{\mathrm{max}}P_{e|i}2^{nE_{ex}(r)+o(r)},$$
(6)
where, $`E_{ex}\left(r\right)`$ is the expurgated error exponent and is an exponentially decreasing function of the rate. The dependence of $`P_{e|i}`$ on $`n`$ indicates that the decoding error probability can be decreased by increasing the length of the channel codewords. The expurgated error exponent is given by ,
$$E_{ex}\left(r\right)=\underset{\rho 1}{sup}\left[\rho r+\underset{\text{q}}{\mathrm{max}}E_x(\rho ,\text{q})\right]$$
(7)
where,
$$E_x(\rho ,\text{q})=$$
$$\rho \mathrm{log}\underset{k=0}{\overset{K1}{}}\underset{i=0}{\overset{K1}{}}q\left(k\right)q\left(i\right)\left[\underset{j=0}{\overset{J1}{}}\sqrt{P\left(j|k\right)P\left(j|i\right)}\right]^{1/\rho }.$$
(8)
Note that $`𝐪=[q(0)q(1)\mathrm{}q(K1)]`$ represents the probability of the input channel alphabets and $`P(j|i)`$ is the probability of receiving output symbol $`j`$ when input symbol $`i`$ is transmitted. In $`\left(\text{8}\right)`$, $`K`$ and $`J`$ represent, respectively, the cardinality of the input and output alphabets of the channel. For a binary erasure channel, $`J=3`$ and $`K=2`$. For the binary erasure channel,the transition probability matrix is given by
For a symmetric channel, the q that maximizes the error exponent is the uniform probability assignment . Thus for the binary erasure channel, $`\text{q}=\left[q\left(0\right)q\left(1\right)\right]=\left[\mathrm{0.5\hspace{0.25em}0.5}\right]`$. Substituting the upper bound for the probability of error into (5), we obtain the end-to-end distortion as
$$D_R(Q,ϵ,\pi )2^{pRr+O\left(1\right)}+2^{kRE_{ex}(r)+o(r)}$$
(9)
Consider the case of large $`R`$: To ensure that neither of the two terms on the right hand side of (9) dominates the distortion upper bound, we choose the exponents of the two terms to be within o(1) of each other , . Hence, we set
$$E_{ex}\left(r\right)=\frac{p}{k}r_{ex}+o\left(1\right),$$
(10)
to obtain the channel coding rate that optimizes the end-to-end distortion at the decoder. This optimal rate is characterized by Theorem 1, which is similar to Theorem 1 in .
Theorem 1: The upper bound on the minimum $`p^{th}`$ power distortion, averaged over all index assignments of a $`k`$-dimensional cascaded good vector-quantizer and channel encoder that transmits over a binary erasure channel with bit erasure probability $`ϵ`$, is achieved with a channel code rate $`r_{ex}`$ satisfying
$$r_{ex}=12^{c_ϵ}\left(\frac{\mathrm{log}\mathrm{log}\left(1/ϵ\right)+\mathrm{log}e+c_ϵ}{\mathrm{log}\left(1/ϵ\right)}\right)$$
$$+O\left(\frac{\mathrm{log}\mathrm{log}\left(1/ϵ\right)}{\mathrm{log}^2\left(1/ϵ\right)}\right)+o\left(1\right),$$
(11)
where, $`c_ϵ`$ satisfies
$$\frac{p}{k}2^{c_ϵ}\frac{\left(p/k\right)\left(\mathrm{log}\mathrm{log}\left(1/ϵ\right)+\mathrm{log}e+c_ϵ\right)2^{c_ϵ}}{log\left(1/ϵ\right)}1=0.$$
(12)
Proof:
For a BEC, evaluating (8), we obtain
$$\underset{\text{q}}{\mathrm{max}}E_x(\rho ,\text{q})=\rho \left[1\mathrm{log}\left(1+ϵ^{1/\rho }\right)\right]$$
(13)
and thus the expurgated error exponent becomes
$$E_{ex}\left(r\right)=\underset{\rho 1}{sup}\left\{\rho \left[1r\mathrm{log}\left(1+ϵ^{1/\rho }\right)\right]\right\}$$
(14)
The $`\rho `$ which maximizes the error exponent and also satisfies (10) is given by
$$\rho =\frac{\mathrm{log}\left(1/ϵ\right)}{\mathrm{log}\mathrm{log}\left(1/ϵ\right)+c_ϵ}$$
(15)
Substituting for $`\rho `$ and $`c_ϵ`$ into (10) we obtain the optimal rate (11) and hence the theorem is proved. $`\mathrm{}`$
Note that the expression for the expurgated error joint source channel rate is similar to the BSC case with the difference being the argument of the $`\mathrm{log}`$ term. Appendix 1 in provides details on the derivation for $`\rho `$ and $`c_ϵ`$. A further simplification in the expression for the rate can be obtained by neglecting the $`O(1)`$ and $`o(r)`$ terms and equating (14) to the exponent of the source coding distortion yielding
$$r_{ex}=\frac{\rho }{\frac{p}{k}+\rho }\left[1\mathrm{log}\left(1+ϵ^{1/\rho }\right)\right]$$
(16)
Numerical values of $`r_{ex}`$ is given in Figure 2 and are explained in Section IV.
### II-B Upper bound on channel coding rate
Following the analysis in , the upper bound on the average distortion minimized over all channel code rates for large $`R`$ and small bit error probability for the binary erasure channel can be obtained as
$$D_R(Q,ϵ,\pi )D_m\left(Q\right)\left(1P_e\right)+\mathrm{\Omega }\left(1\right)\frac{1}{M}\underset{k=1}{\overset{M}{}}P_{e|k}$$
$$=2^{pRr+O\left(1\right)}\left(1P_e\right)+\mathrm{\Omega }\left(1\right)P_e$$
(17)
where $`P_e`$ is the probability of error occurring in the channel. A lower bound on this probability of error is given by,
$$P_e2^{nE_{sl}(r)+o(n)}=2^{kRE_{sl}\left(r\right)+o\left(R\right)}$$
(18)
where $`E_{sl}`$ is the straight line exponent. The straight line exponent $`E_{sl}\left(r\right)`$ is a linear function of $`r`$ which is tangent to the sphere packing exponent $`E_{sp}\left(r\right)`$ and also satisfies $`E_{sl}\left(0\right)=E_{ex}\left(0\right)`$. The sphere packing exponent is given by,
$$E_{sp}\left(r\right)=\underset{\rho 0}{sup}\left[\rho r+\underset{\text{q}}{\mathrm{max}}E_o(\rho ,\text{q})\right]$$
(19)
and,
$$\underset{\text{q}}{\mathrm{max}}E_o(\rho ,\text{q})=\mathrm{log}\underset{j=0}{\overset{J}{}}\left[\underset{k=0}{\overset{K}{}}q(k)P(j|k)^{1/\left(1+\rho \right)}\right]^{1+\rho }$$
(20)
The straight line exponent can be written as,
$$E_{sl}\left(r_{sl}\right)=E_{ex}\left(0\right)+r_{sl}\frac{\left[E_{sp}\left(r^{}\right)E_{ex}\left(0\right)\right]}{r^{}}$$
(21)
where, $`r^{}`$ is the rate at which the straight line exponent meets the sphere packing exponent tangentially. The straight line exponent has also been characterized in for a binary erasure channel. Now, the end-to-end distortion is thus bounded as
$$D_R(Q,ϵ,\pi )2^{pRr+O(1)}+2^{kRE_{sl}(r)+o(R)}$$
(22)
The channel coding rate that minimizes this bound is now characterized in Theorem 2.
Theorem 2: An upper bound on the channel code rate $`r`$ that minimizes the $`p^{th}`$ power distortion averaged over all random index assignments of a $`k`$-dimensional cascaded good vector quantizer for a binary erasure channel with small effective bit erasure probability $`ϵ`$ and large $`R`$ is given by
$$r_{sl}=\frac{E_{ex}\left(0\right)}{\frac{p}{k}\frac{E_{sp}\left(r^{}\right)E_{ex}\left(0\right)}{r^{}}}$$
(23)
Proof: As in the earlier case, for large $`R`$, to prevent either of the terms in the distortion bound (22) from dominating the other, we set the straight line exponent to be linearly proportional to the exponent term of the noiseless-optimal distortion within $`o\left(1\right)`$ of each other. Thus,
$$E_{sl}\left(r\right)=\frac{p}{k}r_{sl}+o\left(1\right)$$
(24)
Substituting (21) in (24), the theorem is proved $`\mathrm{}`$.
Note that to completely characterize $`r_{sl}`$ we need to explicitly evaluate the sphere packing exponent $`E_{sp}(r)`$. It is easily seen that a uniform probability assignment for the input states to the channel $`q(.)`$ maximizes $`E_o(\rho ,\text{q})`$ and thus $`E_{sp}(r)`$ can be evaluated as,
$$E_{sp}\left(r\right)=\underset{\rho 0}{sup}\left\{\rho \left(1r\right)\mathrm{log}\left[\left(1ϵ\right)+ϵ2^\rho \right]\right\}$$
(25)
Note that $`(\text{25})`$ is a concave function of $`\rho `$ and hence the supremum can be replaced by the max operator. The $`\rho `$ which maximizes (25) satisfies
$$r=\frac{\left(1ϵ\right)}{\left(1ϵ\right)+2^\rho ϵ}$$
(26)
We can use this relation between the rate and $`\rho `$ to express the sphere packing exponent in terms of the channel encoding rate for a given erasure channel as,
$$E_{sp}\left(r\right)=r\mathrm{log}r+(1r)\mathrm{log}(1r)r\mathrm{log}\left(\frac{1ϵ}{ϵ}\right)\mathrm{log}ϵ$$
(27)
At $`r^{}`$, the slope of the sphere packing exponent equals the slope of the straight line exponent. Thus,
$$\frac{E_{sp}}{r}_{r=r^{}}=\frac{E_{sp}\left(r^{}\right)E_{ex}\left(0\right)}{r^{}}$$
(28)
Differentiating $`(\text{27})`$ and substituting in (28), we get
$$r^{}=12^{E_{ex}(0)\mathrm{log}\left(1/ϵ\right)}$$
(29)
It turns out that $`E_{sp}(r^{})`$ is nearly 0 for small values of $`ϵ`$.
## III Numerical Results
The bounds derived above for the erasure channel can be easily extended to the case of packet erasures. We use a simplified model for the packet erasure channel and assume that a packet erasure occurs if any of the bits within the packet suffers an erasure. Although this assumption simplifies the packet erasure channel model, it is useful in obtaining closed form bounds on the coding rate over such channels. For a packet of size $`P`$ bits, the probability of a packet erasure $`\delta `$ is given by $`\delta =1\left(1ϵ\right)^P`$, where as before $`ϵ`$ denotes the probability of bit erasure. The error exponent for a binary erasure channel with erasure probability $`ϵ`$ and a $`2^P`$-ary erasure channel with erasure probability $`\delta `$ is the same. Hence, given the packet erasure probability $`\delta `$, we consider an equivalent binary erasure channel with bit erasure probability $`ϵ=1\left(1\delta \right)^{1/P}`$ and find the bounds on the rate and distortion for the corresponding BEC.
The plot of the upper and lower bound on channel coding rate as a function of the erasure channel probability for k=4 and squared distortion measure is given in Fig. 2. The bounds for various packet sizes $`P=1`$, 10 and 100 are shown in Fig. 2 . It is observed that for a given packet size, as the erasure probability increases, the channel coding rate decreases indicating that more bits need to be invested on channel coding to combat a hostile channel. Further, for a given packet erasure probability, as the packet size increases, the channel coding rate increases implying that more bits can be allocated for source coding with larger packet size. Fig. 3 offers a different perspective on the results. From the bounds on the channel coding rate, we can get the bounds on the distortion due to channel coding. By virtue of our optimal joint source-channel coding criterion, the total distortion will be twice the distortion due to channel coding. Hence, given an end-to-end limit on the distortion, we can get the minimum packet length to be chosen from Fig. 3. The squared distortion metric with $`k=4`$, $`R=10`$ and packet erasure probability $`\delta =10^3`$ was chosen. The $`o\left(r\right)`$ and $`o\left(R\right)`$ terms were neglected in the terms for distortion due to noisy channel decoding in $`(\text{9})`$ and $`(\text{22})`$.The asymptotic nature of the curve indicates that large packet size is not required for packet erasure channels with small erasure probabilities.
## IV Conclusion
The results presented in this paper provide a mechanism for optimal concatenation of source and channel coders. Analytic results are provided for lower and upper bounds for a binary erasure channel and for packet erasure channels. The results on packet erasure channel enable us to obtain the bounds on the packet size for a specified bound on the distortion and given packet erasure probability. Alternately, for a given packet erasure probability, we can find bounds on the channel encoding rate for various packet lengths. By studying the optimal rate allocation for a bit and packet erasure channel, one can apply these results for transmission in a wide range of scenarios, including wireline channels with congestion. In future work, these bounds should be expanded to include transmission over more sophisticated channel models.
## Acknowledgment
This work has been supported in part by Nokia Inc.
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# Resolutions of ideals of any six fat points in P2
## 1 Introduction
We begin by describing the problem we solve here, using terminology familiar to experts. Those readers not already familiar with the jargon can rest easy, since we will recall what the terms mean in section 2.
Given general points $`p_1,\mathrm{},p_n`$ of P<sup>2</sup> and arbitrary positive integers $`m_i`$, it is an open problem to determine the graded Betti numbers for the minimal free resolution of the ideal $`I(Z)`$ of the fat point subscheme $`Z=m_1p_1+\mathrm{}+m_np_n`$. It is even an open problem to determine just the Hilbert function of $`I(Z)`$. Partly because of the difficulty of these problems and because a standard approach to them involves considering special configurations of points, and partly because of the intrinsic interest, there has been growing interest in these problems not only for general points but also when the points need not be general, both in the plane and in higher dimensions (see, for example, \[BGV1\], \[BGV2\], \[C\], \[FHL\], \[Fr\], \[GMS\], \[GV1\], \[GV2\], \[H1\], \[H2\], \[H3\], \[H4\], \[HR\]).
In particular, \[GMS\] raises the question of finding all Hilbert functions and graded Betti numbers for ideals of double point subschemes of the plane; i.e., for $`2p_1+\mathrm{}+2p_n\text{P}\text{2}`$, for all possible configurations of the points $`p_i`$. As \[GMS\] discusses, the Hilbert functions which occur for simple point subschemes $`Z=p_1+\mathrm{}+p_n\text{P}\text{2}`$ are known for all possible configurations of the points $`p_i`$; one goal that \[GMS\] works toward is to find all Hilbert functions occurring for double point subschemes $`2Z=2p_1+\mathrm{}+2p_n\text{P}\text{2}`$ such that the support scheme $`Z=p_1+\mathrm{}+p_n\text{P}\text{2}`$ has given Hilbert function. While \[GMS\] shows that for each Hilbert function of simple points there is a Hilbert function which in each degree has minimal value, it leaves unsolved the problem of how to actually find this minimal Hilbert function, even for small values of $`n`$ (such as $`n=6`$), and it raises the question of whether there is also a maximal Hilbert function. (It is worth mentioning that while we talk about the Hilbert function of the ideal $`I(Z)`$, \[GMS\] talks about the Hilbert function of the quotient ring $`R/I(Z)`$, where $`R`$ is the homogeneous coordinate ring of P<sup>2</sup>. Thus what is for us a maximal Hilbert function is for \[GMS\] a minimal Hilbert function.)
We answer all of these questions for the case of 6 points of P<sup>2</sup> (see section 3). Moreso, we give a general approach for answering any problems of the kinds raised in \[GMS\], for any fat point subschemes of P<sup>2</sup> with support at 6 points, regardless of the multiplicities $`m_i`$. More precisely, define a configuration type of $`n`$ points by requiring that sets $`\{p_1,\mathrm{},p_n\}\text{P}\text{2}`$ and $`\{p_1^{},\mathrm{},p_n^{}\}\text{P}\text{2}`$ of distinct points have the same configuration type if and only if, after reordering the points $`p_i^{}`$ if need be, the ideals of $`Z=m_1p_1+\mathrm{}+m_np_n`$ and $`Z^{}=m_1p_1^{}+\mathrm{}+m_np_n^{}`$ have the same Hilbert function for every choice of the nonnegative integers $`m_i`$. We show not only that the set of all configurations of 6 points of P<sup>2</sup> fall into only 11 different types (see Corollary 2.3 and section 3), but that ideals of any two subschemes $`Z=m_1p_1+\mathrm{}+m_6p_6`$ and $`Z^{}=m_1p_1^{}+\mathrm{}+m_6p_6^{}`$ whose points have the same type also have the same graded Betti numbers (see Theorem 3.1 and Example 3.2). Our method also allows us to write down the Hilbert function and graded Betti numbers for any $`Z=m_1p_1+\mathrm{}+m_6p_6`$, given only the coefficients $`m_i`$ and given the configuration type of the points with respect to a specific ordering of the points. (Figure 1 shows the 11 different configuration types of 6 points. Thus type 1 consists of 6 general points; for type 2, three of the points are collinear, etc. Type 11 has all six points on an irreducible conic.)
What is new here is the explicit enumeration of the 11 types (this is easy), and the determination of the graded Betti numbers (this is where most of the effort of this paper lies). It follows from our main result, Theorem 3.1, that numerical Bezout considerations (as discussed in Remark 2.4 and demonstrated in Example 3.2) suffice to determine the graded Betti numbers of a fat point ideal supported at any 6 distinct points of P<sup>2</sup>. (By numerical Bezout considerations we are referring to the version of Bezout’s theorem that tells us that two effective divisors $`C`$ and $`D`$ on an algebraic surface must have a common component if their intersection $`CD`$ is negative. We give a procedure for computing the graded Betti numbers that depends only on computing intersections of divisor classes on a blow up of P<sup>2</sup>, which amounts to taking dot products of integer entry vectors. This procedure is easy to carry out by hand, as shown by Example 3.2. An awk script implementing it can be run over the web by visiting
http://www.math.unl.edu/$``$bharbour/6ptres/6reswebsite.html .)
The facts that, for any $`n8`$ points of the plane, numerical Bezout considerations determine the Hilbert function of any fat point subscheme supported at those points, and that there are only finitely many different configuration types of $`n8`$ points, follow from the main result of \[H2\]. However, these facts seem not to be widely recognized (the authors of \[GMS\], for example, were quite interested when we mentioned this to them), perhaps because the finite set of configurations has never been explicitly written down.
Enumerating these finitely many types for $`7n8`$ takes considerably more effort than doing so for $`n=6`$; in a not yet written preprint, Geramita, Harbourne and Migliore find 29 types of distinct points for $`n=7`$ and 146 for $`n=8`$. Determining how the graded Betti numbers behave will be much more difficult, both because of the many cases that need to be considered, and because the behavior of the graded Betti numbers is more subtle (see \[H5\] and \[FHH\], which work out the graded Betti numbers for 7 and 8 general points respectively, versus \[F1\], which works out the case of 6 general points). Moreover, for $`n>8`$ points, the number of types is infinite. (For example, just by taking points in various configurations on a smooth non-supersingular plane cubic curve, for any positive integer $`r`$ one can by Proposition 1.2 of \[H1\] arrange for the Hilbert function of $`I(mp_1+\mathrm{}+mp_9)`$ in degree $`t=3m`$ to be $`m/r+1`$. Thus the number of Hilbert functions increases with $`m`$, so for $`m`$ large enough no given finite set of types will be sufficient to encompass all of them.)
We now briefly discuss additional background for our work in this paper. The Hilbert function for ideals $`I(Z)`$ of fat point subschemes $`Z\text{P}\text{2}`$ supported at $`n9`$ general points is well known; see, for example, Nagata \[N\], or, for $`n=6`$, Giuffrida \[Gf\]. For $`n>9`$ general points, the problem of finding the Hilbert function of $`I(Z)`$ has been solved only in special cases. As mentioned above, the problem of finding the Hilbert function of $`I(Z)`$ as long as $`Z`$ has support at $`n8`$ points, even possibly infinitely near, was solved, in principle, in \[H2\], without however classifying the possible configuration types.
Type 1 Type 2 Type 3
Type 4 Type 5 Type 6
Type 7 Type 8 Type 9
Type 10 Type 11
Figure 1
A logical next step is to determine the graded Betti numbers for minimal free resolutions of ideals of fat point subschemes in P<sup>2</sup> supported at any configuration of points. Previous results have been given in various cases. The first results are due to Catalisano \[C\], who determined the minimal free resolutions for fat point subschemes supported at distinct points on an irreducible plane conic. The case that the conic is not irreducible or the points are possibly infinitely near was handled in \[H4\]. (Since a connected curve of degree at most 2 in any projective space lies in a plane in that space, by applying \[FHL\] the results of \[C\] and \[H4\] actually also give the Hilbert function and graded Betti numbers for fat points in projective space of any dimension, as long as the support of the points is contained in a connected curve of degree at most 2.) Various cases in which the points of the support are contained in complete intersections in P<sup>2</sup> are studied in \[BGV1\], \[BGV2\], \[GV1\] and \[GV2\]. Additional special configurations are handled in \[GMS\], but only in case of points of multiplicity 2.
Since any five points lie on a smooth conic, Catalisano’s result handles the case of fat point subschemes supported at five general points. The case of 6 general points was worked out by Fitchett \[F1\]. For the case of seven general points, see \[H5\], and for eight general points, see \[FHH\]. Numerous special cases for 9 or more general points have been done (for $`n9`$ general points of multiplicity 1, 2 or 3, see \[GGR\], \[I\] and \[GI\], respectively; for $`n`$ general points of multiplicity $`m`$ when $`m`$ is not too small and $`n`$ is an even square, in light of \[E\], see \[HHF\]; additional cases are handled by \[HR\]). The problem for general points is otherwise open. There is a conjecture for the Hilbert function of the ideal of any fat point subscheme of P<sup>2</sup> supported at general points (see \[H7\] for a discussion), and there are conjectures in special cases for resolutions (see \[H6\] and \[HHF\]), but so far no general conjecture for the resolution has been posed.
In this paper, we extend \[C\] and \[H4\] to the case of any 6 distinct points of P<sup>2</sup>. Our approach involves a case by case analysis for the different configuration types of 6 points in P<sup>2</sup>, depending on finding sets of generators of the cone of nef divisor classes on the surface $`X`$ obtained by blowing up the 6 points. At first glance verifying our result even for a single configuration of points would seem to require checking an infinite number of cases, since there are infinitely many nef divisor classes. The fact that our methods make the problem tractable is of interest in its own right.
## 2 Background
We begin by discussing our methods in more detail. So let $`p_1,\mathrm{},p_n`$ be distinct points of P<sup>2</sup>. Given nonnegative integers $`m_i`$, the fat point subscheme $`Z=m_1p_1+\mathrm{}+m_np_n\text{P}\text{2}`$ is, by definition, defined by the ideal $`I(Z)=I(p_1)^{m_1}\mathrm{}I(p_n)^{m_n}`$, where $`I(p_i)R=k[\text{P}\text{2}]`$ is the ideal generated by all forms (in the polynomial ring $`R`$ in three variables over the field $`k`$) vanishing at $`p_i`$. The support of $`Z`$ consists of the points $`p_i`$ for which $`m_i`$ is positive.
The minimal free resolution of $`I(Z)`$ is an exact sequence of the form
$$0F_1F_0I(Z)0$$
where each $`F_i`$ is a free graded $`R`$-module, where the grading is with respect to the usual grading of $`R`$ by degree, and all entries of the matrix defining the homomorphism $`F_1F_0`$ are homogeneous polynomials in $`R`$ of degree at least 1. To determine $`F_0`$ up to graded isomorphism, it is enough to determine the dimensions of the cokernels of the multiplication maps $`\mu _{Z,i}:I(Z)_iR_1I(Z)_{i+1}`$ for each $`i0`$, where, given a graded $`R`$-module $`M`$, $`M_t`$ denotes the graded component of degree $`t`$. If we denote $`\text{dim cok}(\mu _{Z,i1})`$ by $`t_i`$, then $`F_0=_{i>0}R[i]^{t_i}`$, where $`R[i]`$ is the free graded $`R`$-module of rank 1 with a shift in degrees given by $`R[i]_j=R_{ji}`$. The Hilbert functions of $`I(Z)`$ and $`F_0`$ then determine $`F_1`$ up to graded isomorphism. In fact, if we denote the Hilbert function of $`Z`$ by $`h_Z`$ (i.e., $`h_Z(i)=\text{dim }I(Z)_i`$), and if $`\mathrm{\Delta }`$ denotes the difference operator (i.e., $`\mathrm{\Delta }h_Z(i)=h_Z(i)h_Z(i1)`$), then $`F_1=_{i>0}R[i]^{s_i}`$, where $`s_i=t_i(\mathrm{\Delta }^3h_Z)(i)`$ (see \[FHH\], p. 685).
Thus to determine $`F_0`$ and $`F_1`$ it is enough to determine the Hilbert function of $`I(Z)`$ and the rank of $`\mu _{Z,i}`$ for each $`i`$. The Hilbert function of $`I(Z)`$ can be obtained by applying the result of \[H2\]. It follows from Theorem 3.1 that the ranks of the $`\mu `$ can be found by a maximal rank criterion, as we now explain.
Given $`Z`$, let $`\alpha (Z)`$ be the least degree $`j`$ such that $`h_Z(j)>0`$; i.e., such that $`I(Z)_j0`$. For each $`t\alpha (Z)`$, let $`\gamma (Z,t)`$ be the gcd of $`I(Z)_t`$. Thus $`\gamma (Z,t)`$ is a homogeneous form of some degree $`d_{Z,t}`$. If $`d_{Z,t}=0`$, it is convenient to set $`\gamma (Z,t)=1`$, but if $`d_{Z,t}>0`$, then $`\gamma (Z,t)`$ defines a plane curve $`C=C_{Z,t}`$ of degree $`d_{Z,t}`$. Let $`m_i^{}`$ be the multiplicity $`\text{mult}_{p_i}(C)`$ of the curve at the point $`p_i`$. Thus we get a fat points subscheme $`Z_t^{}=m_1^{}p_1+\mathrm{}+m_n^{}p_n`$. Let $`Z_t^+=(m_1m_1^{})_+p_1+\mathrm{}+(m_nm_n^{})_+p_n`$, where for any integer $`m`$, $`m_+=\text{max }(0,m)`$. Then clearly $`I(Z)_t=\gamma (Z,t)I(Z_t^+)_{td_{Z,t}}`$.
For $`n8`$ and $`t\alpha (Z)`$, it is known that
$$\text{dim}(I(Z_t^+)_{td_{Z,t}})=\left(\genfrac{}{}{0pt}{}{td_{Z,t}+2}{2}\right)\underset{i}{}\left(\genfrac{}{}{0pt}{}{(m_im_i^{})_++1}{2}\right),$$
as a consequence of the fact that a nef divisor $`F`$ on a blow up $`X`$ of P<sup>2</sup> at $`n8`$ points has $`h^1(X,𝒪_X(F))=0=h^2(X,𝒪_X(F))`$ \[H2\]. For $`n8`$, as we discuss in more detail below in Remark 2.4, one can determine $`\alpha (Z)`$ using purely numerical Bezout considerations, and for each $`t\alpha (Z)`$, one can also determine $`Z_t^{}=m_1^{}p_1+\mathrm{}+m_n^{}p_n`$ and $`d_{Z,t}`$ purely numerically, from Bezout considerations. (In order to determine these quantities in the case of $`n=6`$ distinct points, in addition to having the coefficients $`m_i`$, one needs to know only the configuration type with respect to a specific ordering of the points; i.e., one needs to know only whenever there is a line going through three or more of the points $`p_i`$, and which points those are, and if there is a conic going through all 6 points.)
Given that we can determine the Hilbert function of the ideal $`I(Z)`$, to determine the graded Betti numbers $`t_i`$ and $`s_i`$ of the resolution, therefore, it is enough to determine $`t_i`$ for each $`i`$. Since we know the Hilbert function, we know $`\alpha (Z)`$ and clearly, $`t_i=0`$ for $`i<\alpha (Z)`$, and $`t_i=h_Z(\alpha (Z))`$ for $`i=\alpha (Z)`$. If $`i`$ is large enough, the Hilbert function and Hilbert polynomial coincide; i.e., we will have $`\text{dim}(I(Z)_i)=\left(\genfrac{}{}{0pt}{}{i+2}{2}\right)_i\left(\genfrac{}{}{0pt}{}{m_i+1}{2}\right)`$. Let $`\tau (Z)`$ be the least $`i`$ such that this holds, and let $`\sigma (Z)=\tau (Z)+1`$. Regularity considerations \[DGM\] then imply that $`t_i=0`$ for $`i>\sigma (Z)`$.
So assume $`\alpha (Z)i<\sigma (Z)`$. Since $`I(Z)_i=\gamma (Z,i)I(Z_i^+)_{id_{Z,i}}`$ for $`i\alpha (Z)`$, multiplying by $`\gamma (Z,i)`$ gives an inclusion $`I(Z_i^+)_{id_{Z,i}+1}I(Z)_{i+1}`$ and a vector space isomorphism between the images of $`\mu _{Z_i^+,id_{Z,i}}`$ and $`\mu _{Z,i}`$. From the inclusions $`\text{Im}(\mu _{Z_i^+,id_{Z,i}})I(Z_i^+)_{id_{Z,i}+1}I(Z)_{i+1}`$ it now follows that
$$t_{i+1}=\text{dim cok}(\mu _{Z,i})=\text{dim cok}(\mu _{Z_i^+,id_{Z,i}})+(h_Z(i+1)h_{Z_i^+}(id_{Z,i}+1)).$$
Since $`\gamma (Z_i^+,id_{Z,i})=1`$, and assuming that we can determine Hilbert functions, this reduces the problem of computing $`\text{dim cok}(\mu _{Z,i})`$ for an arbitrary $`Z`$ in degrees $`i\alpha (Z)`$ to the problem of computing $`\text{dim cok}(\mu _{Z,i})`$ for an arbitrary $`Z`$ but only in degrees $`i\alpha (Z)`$ such that $`\gamma (Z,i)=1`$. This is what we do. Our main result, Theorem 3.1, essentially says that if $`Z`$ has support at any 6 distinct points of P<sup>2</sup>, and if $`i\alpha (Z)`$ is such that $`\gamma (Z,i)=1`$, then $`\mu _{Z,i}`$ has maximal rank (meaning that $`\mu _{Z,i}`$ is either injective or surjective and hence $`t_{i+1}`$ is either $`h_Z(i+1)3h_Z(i)`$ or 0, respectively). Since $`\gamma (Z_i^+,id_{Z,i})=1`$, it follows that $`\mu _{Z_i^+,id_{Z,i}}`$ has maximal rank, and hence everything on the right hand side of the displayed formula above is in terms of Hilbert functions of fat points supported at the given 6 points. Computing those Hilbert functions thus computes $`\text{dim cok}(\mu _{Z,i})`$.
In order to compute the graded Betti numbers for the minimal free resolution of fat point subschemes $`Z`$ with support at 6 points, we thus need to determine their Hilbert functions and, for each degree $`i`$, we need to determine $`Z_i^+`$ and the degree of $`\gamma (Z,i)`$. The easiest context in which this can be done involves the intersection theory on the surface obtained by blowing up the points. This will also be the context we use to study the rank of $`\mu _{Z,i}`$.
Let $`\pi :X\text{P}\text{2}`$ be the birational morphism obtained by blowing up distinct points $`p_1,\mathrm{},p_n`$ of P<sup>2</sup>. Let $`\text{Cl}(X)`$ be the divisor class group of $`X`$. Let $`E_0`$ be the pullback to $`X`$ of the class of a line on P<sup>2</sup>, and let $`E_1,\mathrm{},E_n`$ be the classes of the exceptional divisors of the blow ups of $`p_1,\mathrm{},p_n`$. Then $`\text{Cl}(X)`$ is formally just a free abelian group with a preferred orthogonal basis $`E_0,\mathrm{},E_n`$. This basis is called an exceptional configuration. (The bilinear form on $`\text{Cl}(X)`$ is given by $`E_iE_j=0`$ for all $`ij`$, $`E_0^2=1`$ and $`E_i^2=1`$ for $`i>0`$.) We are mainly interested in the case that $`n=6`$; hereafter, we will often but not always assume that $`n=6`$.
Problems involving fat points with support at points $`p_1,\mathrm{},p_n`$ on P<sup>2</sup> can be translated to problems involving divisors on $`X`$. Given $`Z`$ and $`t`$, the vector space $`I(Z)_t`$ is a vector subspace of the space of sections $`H^0(\text{P}\text{2},𝒪_{\text{P}\text{2}}(t))`$. The latter is referred to as a complete linear system; $`I(Z)_t`$ is typically a proper subspace, in which case it is referred to as an incomplete linear system. However, we can associate to $`Z=m_1p_1+\mathrm{}+m_np_n`$ and $`t`$ the divisor class $`F(Z,t)=tE_0m_1E_1\mathrm{}m_nE_n`$ on $`X`$, in which case $`I(Z)_t`$ can be canonically identified (as a vector space) with the complete linear system $`H^0(X,𝒪_X(F(Z,t)))`$.
Given a divisor or divisor class $`F`$ on $`X`$, it will be convenient to write $`h^i(X,F)`$ in place of $`h^i(X,𝒪_X(F))`$, and we will refer to a divisor class $`F`$ as effective if $`h^0(X,F)>0`$; i.e., if it is the class of an effective divisor. In particular, $`\text{dim }I(Z)_t=h^0(X,F(Z,t))`$ for all $`Z`$ and $`t`$, and the ranks of $`\mu _{Z,t}`$ and $`\mu _{F(Z,t)}`$ are equal, where
$$\mu _{F(Z,t)}:H^0(X,F(Z,t))H^0(X,E_0)H^0(X,F(Z,t)+E_0)$$
is the natural map given by multiplication.
Whenever $`N`$ is a prime divisor (i.e., a reduced irreducible curve) such that $`F(Z,t)N<0`$, we have $`h^0(X,F(Z,t))=h^0(X,M)`$, where $`M=F(Z,t)N`$. Moreover, clearly the kernels of $`\mu _{F(Z,t)}`$ and $`\mu _M`$ have the same dimension, so if we can compute $`h^0`$ for arbitrary divisors on $`X`$, finding the rank of $`\mu _{F(Z,t)}`$ is equivalent to doing so for $`\mu _M`$. If we have a complete list of prime divisors $`N`$ of negative self-intersection, then whenever $`F(Z,t)`$ is effective, we can subtract off prime divisors of negative self-intersection to obtain an effective class $`M`$ which is nef (meaning that $`MD0`$ for every effective divisor $`D`$), in which case $`h^0(X,F(Z,t))=h^0(X,M)`$ and the kernels of $`\mu _{F(Z,t)}`$ and $`\mu _M`$ have the same dimension, thereby reducing the problem to the case of computing $`h^0(X,M)`$ and ranks of $`\mu _M`$ only when $`M`$ is nef.
This is very helpful, since for $`n8`$, $`h^1(X,M)=0=h^2(X,M)`$ whenever $`M`$ is nef (\[H2\]) and hence $`h^0(X,M)=(M^2K_XM)/2+1`$ by Riemann-Roch. Thus for $`n8`$, the Hilbert function of $`I(m_1p_1+\mathrm{}+m_np_n)`$ is completely determined by the coefficients $`m_i`$ and by the set of classes of prime divisors of negative self-intersection on the surface $`X`$ obtained by blowing up the points $`p_i`$. (For $`n9`$, this is no longer true. This is because $`h^1(X,𝒪_X(F))=0`$ can fail for nef divisors when $`n9`$, as shown by considering a general pencil of cubics.)
But whereas $`\mu _M`$ is always surjective for nef divisors $`M`$ for any $`n5`$ distinct (or even possibly infinitely near) points \[H4\], $`\mu _M`$ can fail to have maximal rank for nef divisors when $`n7`$ \[H5\], even for $`n`$ general points. However, for $`n=6`$ general points, $`\mu _M`$ always at least has maximal rank when $`M`$ is nef \[F1\]. This leaves open the question of whether $`\mu _M`$ may fail to have maximal rank for some nef $`M`$ for some particular choice of $`n=6`$ distinct points; we show that $`\mu _M`$ has maximal rank for any nef $`M`$ for all choices of the points $`p_i`$.
We begin by determining the subset $`\text{NEG}(X)\text{Cl}(X)`$ of divisor classes of effective reduced irreducible divisors of negative self-intersection. Among all 6 point blow ups $`X`$ of P<sup>2</sup>, it turns out there are only finitely many possible subsets $`\text{NEG}(X)`$, and $`\text{NEG}(X)`$ is itself always finite. (By Corollary 2.3, up to reordering the points, the possible subsets $`\text{NEG}(X)`$ correspond bijectively with the configuration types of Figure 1.) We can then obtain our result by an analysis for each possible subset $`\text{NEG}(X)`$. As a practical matter, it is easier to consider the subset
$$\text{neg}(X)=\{C\text{NEG}(X)\text{ : }C^2<1\},$$
since $`\text{neg}(X)`$ is a proper (and usually substantially smaller) subset of $`\text{NEG}(X)`$, but $`\text{neg}(X)`$ determines $`\text{NEG}(X)`$, by Remark 2.2. In fact, the elements of $`\text{neg}(X)`$ correspond to the curves displayed in Figure 1. (For example, for configuration type 1, $`\text{neg}(X)`$ is empty, for configuration type 2, $`\text{neg}(X)`$ consists of the divisor class of the proper transform of the line through the three collinear points, etc.)
While $`\text{NEG}(X)`$ and $`\text{neg}(X)`$ depend on the particular points $`p_i`$, we now define a fixed finite subset of the divisor class group $`\text{Cl}(X)`$ which contains them. Consider $`𝒬`$, where $`=\{E_i:i>0\}`$ ($``$ here is for blow up of a point), $`=\{E_0E_{i_1}\mathrm{}E_{i_r}:r2,0<i_1<\mathrm{}<i_r6\}`$ ($``$ here is for points on a line), and $`𝒬=\{2E_0E_{i_1}\mathrm{}E_{i_r}:r5,0<i_1<\mathrm{}<i_r6\}`$ ($`𝒬`$ here is for points on a conic, defined by a quadratic equation).
The next result, which is well known but hard to cite in the form we need, shows that there are only finitely many possibilities for $`\text{NEG}(X)`$, since it is a subset of $`𝒬`$. (The finiteness remains true as long as $`n<9`$ but can fail for $`n9`$. In addition, more possibilities occur than the ones listed here if $`n`$ is 7 or 8.)
###### Lemma 2.1
Let $`X`$ be obtained by blowing up 6 distinct points of P<sup>2</sup>. Then the following hold:
$`\text{NEG}(X)𝒬`$, and every class in $`\text{NEG}(X)`$ is the class of a smooth rational curve;
for any nef $`F\text{Cl}(X)`$, $`F`$ is effective (hence $`h^2(X,F)=0`$), $`|F|`$ is base point free, $`h^0(X,F)=(F^2K_XF)/2+1`$ and $`h^1(X,F)=0`$;
$`\text{NEG}(X)`$ generates the subsemigroup $`\text{EFF}(X)\text{Cl}(X)`$ of classes of effective divisors; and
any class $`F`$ is nef if and only if $`FC0`$ for all $`C\text{NEG}(X)`$.
Proof. Riemann-Roch for a smooth rational surface $`X`$ states that $`h^0(X,A)h^1(X,A)+h^2(X,A)=(A^2K_XA)/2+1`$ holds for any divisor class $`A`$. Also, $`K_X=3E_0E_1\mathrm{}E_6`$, so $`K_XE_0=3`$. If $`F`$ is effective, then $`FE_00`$, since $`E_0`$ is nef. (The reason $`E_0`$ is nef is that it is the class of an irreducible divisor of nonnegative self-intersection, hence any effective divisor meets it nonnegatively. More generally, any effective divisor which meets each of its components nonnegatively is nef.) By duality, $`h^2(X,F)=h^0(X,K_XF)`$, and $`h^0(X,K_XF)=0`$ since $`K_XE_0=3`$, hence $`(K_XF)E_0<0`$. This verifies the parenthetical remark in part (b). Similarly, $`h^2(X,K_X)=0`$, so we have $`h^0(X,K_X)=K_X^2+1+h^1(X,K_X)`$, but for us $`K_X^2=3`$, so $`h^0(X,K_X)=4+h^1(X,K_X)`$. Thus $`K_X`$ is the class of an effective divisor, say $`D`$. Moreover, the subgroup $`K_X^{}\text{Cl}(X)`$ of all classes orthogonal to $`K_X`$ is negative definite. This is easy to see since the classes $`E_1E_2`$, $`E_1+E_22E_3`$, $`E_1+E_2+E_33E_4`$, $`E_1+E_2+E_3+E_44E_5`$ and $`2E_0E_1\mathrm{}E_6`$ have negative self-intersection but are linearly independent and pairwise orthogonal, hence give an orthogonal basis of $`K_X^{}`$ over the rationals. On the other hand, it is not hard to check that $`E_0E_1E_2E_3`$, $`E_1E_2`$, $`\mathrm{}`$, $`E_5E_6`$ give a $`𝐙`$-basis for $`K_X^{}`$, and since each basis element has self-intersection $`2`$, it follows that $`A^2`$ is even for every $`AK_X^{}`$. I.e., $`K_X^{}`$ is even and negative definite.
To justify (a), let $`C`$ be the class of a reduced irreducible divisor on $`X`$, with $`C^2<0`$. Since $`E_0`$ is nef, we know $`CE_00`$. If $`CE_0=0`$, then $`C`$ must be a component of one of the $`E_i`$, hence $`C`$, since each $`E_i`$ is reduced and irreducible. If $`CE_0=1`$, then $`C`$ is the proper transform of a line in P<sup>2</sup>, so $`C`$. If $`CE_0=2`$, then $`C`$ is the proper transform of a smooth conic in P<sup>2</sup>, so $`C𝒬`$. By explicitly applying adjunction $`C^2+CK_X=2g2`$, where $`g`$ is the (a priori arithmetic) genus of $`C`$, any $`C𝒬`$ which is the class of a prime divisor has $`g=0`$ and so is the class of a smooth rational curve.
Now it suffices to show that we cannot have $`CE_0>2`$. If $`CD<0`$, then $`C`$ is the class of an irreducible component of $`D`$, hence $`E_0(K_XC)0`$, so $`E_0C3`$. If $`CE_0=3`$, then $`C`$ is the proper transform of an irreducible plane cubic. But an irreducible plane cubic has at most one singular point, which must be of multiplicity 2. Thus its proper transform is either $`3E_0E_{i_1}\mathrm{}E_{i_r}`$, with $`0<i_1<\mathrm{}<i_r6`$, or $`3E_02E_{i_1}E_{i_2}\mathrm{}E_{i_r}`$, with $`0<i_2<\mathrm{}<i_r6`$ and $`0<i_1<6`$ such that $`i_1i_j`$ for $`j>1`$. But in neither case would we have $`C^2<0`$, so $`CE_03`$ cannot happen.
Now say $`CD0`$. From adjunction, since $`0CD=K_XC`$, we have $`1C^22`$, with $`g=0`$ in any case, hence $`C`$ is a smooth rational curve. If $`CD=0`$, then $`CK_X^{}`$ and adjunction gives $`C^2=2`$, but since $`K_X^{}`$ is negative definite, it has only finitely many classes of self-intersection $`2`$. One can show that the only classes in $`K_X^{}`$ of self-intersection $`2`$ are $`\pm E_i\pm E_j`$, $`0<i<j6`$, $`\pm E_0\pm E_i\pm E_j\pm E_k`$, $`0<i<j<k6`$, and $`\pm 2E_0\pm E_1\pm \mathrm{}\pm E_6`$. (To see this, assume that $`A=aE_0b_1E_1\mathrm{}b_6E_6K_X^{}`$. Thus $`3a=b_1+\mathrm{}+b_6`$. Working over $`\text{Cl}(X)_𝐙𝐐`$, let $`m=(b_1+\mathrm{}+b_6)/6`$, so $`a=2m`$, and define $`B=aE_0m(E_1+\mathrm{}+E_6)`$. Then $`BK_X=0`$, but $`A^2B^2=2m^2`$. If $`A^2=2`$, then we must have $`a2`$, in order to have $`m1`$. Thus $`a`$ is either 0, 1 or 2, and now it is easy to enumerate solutions $`A^2=2`$.) Among these classes, only those in $`𝒬`$ can be classes of prime divisors. (This is because a prime divisor must, first, meet $`E_0`$ nonnegatively, and second, when expressed as a linear combination $`a_0E_0a_iE_i`$, if $`a_j<0`$ for some $`j>0`$, then it must be a component of $`E_j`$ and thus must be in $``$.)
If $`CD>0`$, then $`C^2=1=K_XC`$. Let $`YX`$ be obtained by blowing up a seventh, general point $`p_7`$. This morphism induces an inclusion $`\text{Cl}(X)\text{Cl}(Y)`$. Then, arguing as above, $`K_Y^{}`$ is even and negative definite, and the only solutions to $`A^2=2`$ for $`AK_Y^{}`$ are of the form $`\pm E_i\pm E_j`$, $`0<i<j7`$, $`\pm E_0\pm E_i\pm E_j\pm E_k`$, $`0<i<j<k7`$, and $`\pm 2E_0\pm E_{i_1}\pm \mathrm{}\pm E_{i_6}`$, $`0<i_1<\mathrm{}<i_67`$. Thus $`CE_7`$ is in $`K_Y^{}`$, with $`(CE_7)^2=2`$, since $`CE_7=0`$, and (keeping in mind that $`C`$ is a prime divisor also on $`Y`$ and that $`CK_X=1`$) it follows that $`C`$ is among $`E_i`$, $`0<i6`$, $`E_0E_iE_j`$, $`0<i<j6`$, and $`2E_0E_{i_1}\mathrm{}E_{i_5}`$, $`0<i_1<\mathrm{}<i_56`$. This finishes the proof of (a).
To prove (b), we have $`h^1(X,F)=0`$ and $`h^2(X,F)=0`$ by Theorem 8, \[H2\]. Thus $`h^0(X,F)=(F^2K_XF)/2+1`$ follows by Riemann-Roch. But $`F^20`$ holds for nef divisors (Proposition 4, \[H2\]), so $`F`$ is effective. To see that $`|F|`$ is base point free, note that a nef divisor in $`K_X^{}`$ must be 0. Now apply Theorem III.1(a,b) of \[H3\] to see that $`|F|`$ is fixed component free, and has a base point only if $`K_XF=1`$, in which case, using $`Y`$ as above, we see that $`FE_7`$ must be effective, but $`FE_7K_Y^{}`$, so $`F^21=(FE_7)^20`$. But $`(FE_7)^2=0`$ implies $`FE_7=0`$, which is impossible since then $`0=FE_7=E_7^2=1`$. Thus $`0>(FE_7)^2=F^21`$, so $`F^2=0`$. However, we also have $`K_XF=1`$, which contradicts $`h^0(X,F)=(F^2K_XF)/2+1`$, since $`h^0(X,F)`$ must be an integer. Thus we cannot have $`K_XF=1`$ if $`F`$ is nef.
Consider (c). Let $`G`$ be the class of an effective divisor. We can write $`G=N+F`$, where $`N`$ is the fixed part of $`|G|`$, and $`F`$ is nef. Note that no component of $`N`$ can be nef, since nef divisors (in our situation) are base point free, whereas components of $`N`$ are fixed. Thus the class of every component of $`N`$ is in $`\text{NEG}(X)`$. Now, if a class $`F=a_0E_0a_1E_1\mathrm{}a_6E_6`$ is nef for a particular set of distinct points $`p_i`$, then it remains nef when the points $`p_i`$ are general, and if $`F`$ is effective when the points are general, it was effective to begin with. (This is because by semicontinuity the effective subsemigroup can never get smaller as the points are specialized, so the nef cone can never enlarge.) And if the points $`p_i`$ are general, then $`\text{NEG}(X)`$ consists of the exceptional classes; i.e., the classes $`E_i`$, $`i>0`$, $`E_0E_iE_j`$, $`0<i<j6`$, and $`2E_0E_{i_1}\mathrm{}E_{i_5}`$, $`0<i_1<\mathrm{}<i_56`$. It follows from \[H1\], that the class of every effective divisor is a nonnegative sum of exceptional classes. (The results of \[H1\] show that it is enough to show that $`E_0`$, $`E_0E_1`$, $`2E_0E_1E_2`$, and $`3E_0E_1\mathrm{}E_j`$, $`3j6`$ are, but this is easy; for example, $`E_0=(E_0E_1E_2)+E_1+E_2`$.) Thus given a class $`F`$ which is nef for a given set of points $`p_i`$, $`FE`$ is effective for some $`E`$ among the classes $`E_i`$, $`i>0`$, $`E_0E_iE_j`$, $`0<i<j6`$, and $`2E_0E_{i_1}\mathrm{}E_{i_5}`$, $`0<i_1<\mathrm{}<i_56`$. If $`E`$ is a prime divisor, then $`E\text{NEG}(X)`$. If not, then $`EN^{}<0`$ for some $`N^{}\text{NEG}(X)`$ (otherwise, $`E`$ is nef, hence $`h^0(X,E)=(E^2K_XE)/2+1=1`$, but also $`|E|`$ must be base point free, hence $`h^0(X,E)>1`$).
Thus either way there is an $`N^{}\text{NEG}(X)`$ such that $`FN^{}`$ is effective. By replacing $`F`$ by $`FN^{}`$ and repeating the process, we eventually reach the case that $`F=0`$, hence any effective divisor is a sum of elements of $`\text{NEG}(X)`$.
Finally, we prove (d). To show $`F`$ is nef, we just need to show that $`FC0`$ for each class $`C`$ of an effective divisor. But each such $`C`$ is a nonnegative sum of classes in $`\text{NEG}(X)`$ and any class in $`\text{NEG}(X)`$ is the class of an effective divisor. It follows that $`FC0`$ for the class $`C`$ of an effective divisor if and only if $`FC0`$ for every $`C\text{NEG}(X)`$.
###### Remark 2.2
We now show how $`\text{neg}(X)`$ determines $`\text{NEG}(X)`$. In fact,
$$\text{NEG}(X)=\text{neg}(X)\{C𝒬\text{ }|\text{ }C^2=1,CD0D\text{neg}(X)\}.$$
The forward inclusion follows from Lemma 2.1(a). For the reverse, say $`C^2=1`$ for some $`C𝒬`$. It is easy to check case by case that each such $`C`$ is effective, hence $`CC^{}<0`$ for some $`C^{}\text{NEG}(X)`$. Given that $`CD0`$ for all $`D\text{neg}(X)`$, then $`C^{}\text{NEG}(X)\text{neg}(X)`$. But any two distinct elements of $`𝒬`$ of self-intersection $`1`$ meet nonnegatively, hence $`C=C^{}\text{NEG}(X)`$.
By Lemma 2.1 and Remark 2.2 it follows that specifying $`\text{neg}(X)`$ as a subset of $`𝒬`$ is equivalent to specifying the configuration type of the six points blown up to obtain $`X`$:
###### Corollary 2.3
Let $`A`$ and $`A^{}`$ be sets of six distinct points of P<sup>2</sup>. Then $`A`$ and $`A^{}`$ have the same configuration type if and only if, for some orderings $`A=\{p_1,\mathrm{},p_6\}`$ and $`A^{}=\{p_1^{},\mathrm{},p_6^{}\}`$, we have $`f(\text{neg}(X))=\text{neg}(X^{})`$, where $`X`$ is the surface obtained by blowing up the points $`p_i`$, $`X^{}`$ is the surface obtained by blowing up the points $`p_i^{}`$, $`E_0,\mathrm{},E_6`$ and $`E_0^{},\mathrm{},E_6^{}`$ are the corresponding exceptional configurations and $`f:\text{Cl}(X)\text{Cl}(X^{})`$ is the map defined by $`f(E_i)=E_i^{}`$ for all $`i`$.
Proof. If $`A`$ and $`A^{}`$ have the same configuration type, then $`f(\text{EFF}(X))=\text{EFF}(X^{})`$, hence $`f(\text{NEG}(X))=\text{NEG}(X^{})`$ (since NEG is the set of all $`C`$ in EFF such that $`C^2<0`$ but $`C`$ is not the sum of two nontrivial elements of EFF), so $`f(\text{neg}(X))=\text{neg}(X^{})`$. Conversely, by Remark 2.2, $`\text{neg}(X)`$ determines $`\text{NEG}(X)`$, and, by Lemma 2.1 (and the proof of Lemma 2.1(c)), $`\text{NEG}(X)`$ determines $`h^0(X,G)`$ for any class $`G`$. I.e., if $`f(\text{neg}(X))=\text{neg}(X^{})`$, then $`h^0(X,G)=h^0(X^{},f(G))`$ for every class $`G`$, hence $`A`$ and $`A^{}`$ have the same configuration type.
The next remark shows explicitly how to determine Hilbert functions, given $`\text{NEG}(X)`$ (or, equivalently by Remark 2.2, given $`\text{neg}(X)`$).
###### Remark 2.4
Given a fat points subscheme $`Z=m_1p_1+\mathrm{}+m_6p_6`$ with support at 6 distinct points, for each $`t`$ consider the class $`F=F(Z,t)=tE_0m_1E_1\mathrm{}m_6E_6`$. For each $`C\text{NEG}(X)`$, check $`FC`$. If $`FC<0`$, then $`h_Z(t)=h^0(X,F)=h^0(X,FC)`$, so we can replace $`F`$ by $`FC`$ while preserving $`h^0`$. Continue replacing the current $`F`$ by $`FC`$ whenever the current $`F`$ meets some $`C\text{NEG}(X)`$ negatively. Eventually we obtain an $`F`$ such that either $`FE_0<0`$, in which case $`0=h^0(X,F)=h^0(X,F(Z,t))`$, or $`FC0`$ for all $`C\text{NEG}(X)`$, in which case $`F`$ is nef and hence $`h^0(X,F(Z,t))=h^0(X,F)`$ is given by Lemma 2.1(b). This procedure thus gives us a way to determine the value $`h_Z(t)`$ of the Hilbert function $`h_Z`$ for every $`t`$. Note that determining $`h_Z(t)`$ involves nothing more than integer arithmetic and addition and subtraction in the rank 7 free abelian group $`\text{Cl}(X)`$. It requires only that we know $`\text{NEG}(X)`$ (or even just $`\text{neg}(X)`$) and the multiplicities $`m_i`$ of the points of support of $`Z`$. We do not need to know the points $`p_i`$ themselves.
When $`t\alpha (Z)`$, we also want to know the multiplicity $`m_i^{}=\text{mult}_{p_i}(C_{Z,t})`$ and degree $`d_{Z,t}`$ of the curve $`C_{Z,t}`$ defined by $`\gamma (Z,t)`$, whenever $`\gamma (Z,t)`$ has positive degree. But $`\gamma (Z,t)`$ by Lemma 2.1 just defines the fixed component of the linear system $`I(Z)_t=H^0(X,F(Z,t))`$, and hence if $`F`$ is the nef divisor class obtained by successively subtracting from $`F(Z,t)`$ classes in $`\text{NEG}(X)`$ as above, then $`F=F(Z_t^+,td_{Z,t})`$ and $`FF(Z_t^+,td_{Z,t})=d_{Z,t}E_0m_1^{}E_1\mathrm{}m_6^{}E_6`$, so knowing $`\text{NEG}(X)`$ allows us to determine $`d_{Z,t}`$ and the $`m_i^{}`$, and $`Z_t^+`$.
Although Lemma 2.1 gives us a criterion for a class being nef, our method of proof for Theorem 3.1 requires explicit generators for the nef cone; i.e., for the cone $`\text{NEF}(X)`$ of nef divisor classes on a given $`X`$, which by Lemma 2.1 is just the cone of all $`F`$ such that $`FC0`$ for all $`C\text{NEG}(X)`$. Actually, it will turn out that we will need explicit generators only when the anticanonical class, $`K_X=3E_0E_1\mathrm{}E_6`$, is nef. The problem of determining generators of $`\text{NEF}(X)`$ is an example of the general problem of finding generators for the dual of a nonnegative subsemigroup whose generators are given (in this case $`\text{EFF}(X)`$ is the subsemigroup, generated by $`\text{NEG}(X)`$). This is not an easy computation in general, but in case $`K_X`$ is nef the action of the Weyl group, which we now recall, provides a significant simplification.
Let $`r_0=E_0E_1E_2E_3`$ and for $`1i5`$, let $`r_i=E_iE_{i+1}`$. (These are the so-called simple roots of the Lie-theoretic root system of type $`𝐄_6`$.) Each homomorphism $`s_i:\text{Cl}(X)\text{Cl}(X)`$ defined for any $`x\text{Cl}(X)`$ by the so-called reflection $`s_i(x)=x+(xr_i)r_i`$ through $`r_i`$ preserves the intersection product, and moreover $`s_i(K_X)=K_X`$ for all $`i`$. The subgroup of the orthogonal group of $`\text{Cl}(X)`$ generated by the $`s_i`$ is called the Weyl group, denoted $`W_6`$. Since the reflection $`s_i`$ for $`i>0`$ is just the transposition of $`E_i`$ and $`E_{i+1}`$, we see that $`W_6`$ contains the group $`S_6`$ of permutations of $`E_1,\mathrm{},E_6`$. The element $`s_0`$ corresponds to a quadratic transformation.
The group $`W_6`$ is a finite group of order 51,840. The $`W_6`$ orbit of $`E_0`$ is the following list and those obtained from these by permuting the terms involving $`E_i`$ with $`i>0`$:
$`E_0`$ $`4E_02E_12E_22E_3E_4E_5E_6`$
$`2E_0E_1E_2E_3`$ $`5E_02E_12E_22E_32E_42E_52E_6`$
$`3E_02E_1E_2E_3E_4E_5`$
Similarly, the $`W_6`$ orbit of $`E_0E_1`$, up to permutations, is:
$`E_0E_1`$ $`3E_02E_1E_2E_3E_4E_5E_6`$
$`2E_0E_1E_2E_3E_4`$
And the $`W_6`$ orbit of $`2E_0E_1E_2`$, up to permutations, is:
$`2E_0E_1E_2`$ $`4E_03E_1E_2E_3E_4E_5E_6`$
$`3E_02E_1E_2E_3E_4`$ $`5E_03E_12E_22E_32E_4E_5E_6`$
$`4E_02E_12E_22E_3E_4E_5`$ $`6E_03E_13E_22E_32E_42E_52E_6`$
The $`W_6`$ orbit of $`3E_0E_1E_2E_3`$, up to permutations, is:
$`3E_0E_1E_2E_3`$ $`6E_04E_12E_22E_32E_4E_5E_6`$
$`4E_02E_12E_2E_3E_4`$ $`7E_04E_13E_23E_32E_42E_5E_6`$
$`5E_03E_12E_22E_3E_4E_5`$ $`8E_04E_14E_23E_33E_42E_52E_6`$
$`6E_03E_13E_22E_32E_42E_5`$ $`9E_04E_14E_24E_33E_43E_53E_6`$
$`6E_03E_13E_23E_3E_4E_5E_6`$
The $`W_6`$ orbit of $`3E_0E_1E_2E_3E_4`$, up to permutations, is:
$`3E_0E_1E_2E_3E_4`$ $`5E_03E_12E_22E_3E_4E_5E_6`$
$`4E_02E_12E_2E_3E_4E_5`$ $`6E_03E_13E_22E_32E_42E_5E_6`$
$`5E_02E_12E_22E_32E_42E_5`$ $`7E_03E_13E_23E_33E_42E_52E_6`$
The $`W_6`$ orbit of $`3E_0E_1E_2E_3E_4E_5`$, up to permutations, is:
$`3E_0E_1E_2E_3E_4E_5`$ $`5E_02E_12E_22E_32E_42E_5E_6`$
$`4E_02E_12E_2E_3E_4E_5E_6`$
Finally, the $`W_6`$ orbit of $`K_X=3E_0E_1E_2E_3E_4E_5E_6`$ is just itself. The union of these orbits contains 1279 elements. The next lemma says that the nef elements among these 1279 generate the nef cone.
###### Lemma 2.5
Let $`X`$ be a smooth projective rational surface with a birational morphism to P<sup>2</sup> such that $`\text{Cl}(X)`$ has rank 7. If $`K_X`$ is nef, then the set $`\mathrm{\Omega }=\{FW_6G:FC0\text{ for all }C𝒩\}`$ generates $`\text{NEF}(X)`$ as a nonnegative subsemigroup of $`\text{Cl}(X)`$, where $`𝒩`$ is the set of classes of reduced irreducible curves with $`C^2=2`$ (the so-called nodal roots) and $`G`$ is the set consisting of $`E_0`$, $`E_0E_1`$, $`2E_0E_1E_2`$, $`3E_0E_1E_2E_3`$, $`3E_0E_1E_2E_3E_4`$, $`3E_0E_1E_2E_3E_4E_5`$, and $`3E_0E_1E_2E_3E_4E_5E_6`$.
Proof. From the proof of Lemma 2.1, we know the complete list of classes $`C`$ with $`C^2=2`$ and $`CK_X=0`$ and it is not hard to check that they are contained in (and, since $`W_6`$ preserves the intersection form, thus equal to) a single orbit of $`W_6`$; note, for example, $`s_3s_0(E_3E_4)=E_0E_1E_2E_3`$. This orbit is also known as the set of roots of the root system $`𝐄_6`$. It is easy to verify that half of the roots are nonnegative integer linear combinations of the simple roots $`r_0,\mathrm{},r_5`$; the rest are the additive inverses of these. The former are called positive roots; the latter are called negative roots. The class of a reduced irreducible curve $`C`$ with $`C^2=2`$ is necessarily a positive root: it satisfies $`C^2=2`$ and $`CK_X=0`$, so it is a root. Also, since $`E_0`$ is nef, we have $`E_0C0`$. If $`E_0C>0`$, $`C`$ is clearly one of the positive roots. If $`E_0C=0`$, then $`C`$ is a component of one of the exceptional curves $`E_i`$, and thus of the form $`E_iE_j`$ for some $`0<i<j`$, which is a positive root. It is now not hard to check for any two positive roots that $`rr^{}2`$, with $`rr^{}=2`$ if and only if $`r=r^{}`$.
Similarly, we also know the complete list of classes $`C`$ with $`C^2=1`$ and $`CK_X=1`$, and we can again check directly that they form a single orbit $``$ of $`W_6`$; note, for example, $`s_0(E_1)=E_0E_2E_3`$. Since $``$ is preserved under the action of $`W_6`$, so is the nonnegative subsemigroup $`^{}`$ dual to $``$, consisting of all classes $`F`$ such that $`FC0`$ for all $`C`$.
By direct check, $`GC0`$ for all $`C`$, so we have $`G^{}`$, hence $`W_6G^{}`$. Since $`\text{NEG}(X)=𝒩`$, it follows that $`\mathrm{\Omega }\text{NEF}(X)`$. Now we must see that $`\mathrm{\Omega }`$ generates $`\text{NEF}(X)`$. Note that $`\mathrm{\Omega }`$ is $`W_6G^{}𝒩^{}`$, hence it is precisely the set of nef elements in $`W_6G`$.
Since $`W_6`$ is finite, for each $`F^{}`$ there is some $`wW_6`$ such that $`E_0wF`$ is as small as possible. Let $`wF=a_0E_0a_1E_1\mathrm{}a_6E_6`$. Since we can permute the $`a_i`$ with $`i>0`$ by applying $`s_j`$ with $`j>0`$ and this does not affect $`E_0wF`$, we may assume that $`a_1a_2\mathrm{}a_6`$. Since $`E_6wF0`$, we have $`a_60`$. If $`r_0wF<0`$, then we would have $`E_0s_0(wF)<E_0wF`$, so we also have $`r_0wF0`$; i.e., $`a_0a_1+a_2+a_3`$.
This means the $`W_6`$-orbit of every class $`F^{}`$ intersects the subsemigroup $`A`$ of classes $`H=b_0E_0b_1E_1\mathrm{}b_6E_6`$ defined by the conditions $`b_0b_1+b_2+b_3`$ and $`b_1\mathrm{}b_60`$; i.e., by the conditions $`Hr_00`$, $`\mathrm{}`$, $`Hr_50`$. It is not hard to check that the set $`G`$ of classes $`E_0`$, $`E_0E_1`$, $`2E_0E_1E_2`$, $`3E_0E_1E_2E_3`$, $`3E_0E_1E_2E_3E_4`$, $`3E_0E_1E_2E_3E_4E_5`$, $`3E_0E_1E_2E_3E_4E_5E_6`$ generates $`A`$ (which in fact turns out to be a fundamental domain for the action of $`W_6`$ on $`^{}`$). It is easy to check directly that, for every class $`F`$ in $`A`$, $`FC0`$ for every class $`C`$ with $`C^2=1`$ and $`CK_X=1`$. Also, $`Fr_i0`$ holds for all $`i`$ since $`FA`$, hence $`FC0`$ for the class $`C`$ of every reduced irreducible curve with $`C^2=2`$, since each such $`C`$ is a positive root. Thus $`A\text{NEF}(X)`$.
Now let $`F`$ be any nef class. There is a sequence $`r_{i_1},\mathrm{},r_{i_l}`$ of simple roots such that $`F_jC0`$ for all $`C𝒩_j`$, each element of $`𝒩_j`$ is a positive root, and $`F_lA`$, where $`F=F_0`$, $`F_j=s_{i_j}(F_{j1})`$ for $`1jl`$, $`𝒩_0=𝒩`$, and $`𝒩_j=s_{i_j}(𝒩_{j1})`$ for $`1jl`$. For each $`j`$, let $`i_j`$ be the largest $`i`$ such that $`F_{j1}r_i<0`$. If none exist, then $`l=j1`$ and $`F_lA`$, by definition of $`A`$. Otherwise, let $`F_j=s_{i_j}(F_{j1})`$. If $`F=a_0E_0a_1E_1\mathrm{}a_6E_6`$, what the sequence of operations does is to permute $`a_1,\mathrm{},a_6`$ so that they are nondecreasing, and then to decrease $`a_0`$ whenever $`s_0`$ is applied. But the orbit $`W_6F`$ of $`F`$ is contained in $`^{}`$, hence every element $`H`$ of the orbit has $`HE_0=H(E_0E_1E_2)+HE_1+HE_20`$; thus we cannot forever go on reducing the coefficient of $`E_0`$, so eventually we arrive at a class $`F_l`$ for which $`F_lr_i0`$ for all $`i`$, and hence $`F_lA`$. Now, $`F_0C0`$ for all $`C𝒩_0`$ since $`F=F_0`$ is nef. Also, $`wFwC=FC`$ for all $`wW_6`$ since $`W_6`$ preserves the intersection form. It follows that $`F_jC0`$ for all $`C𝒩_j`$ for all $`j`$. Moreover, $`r_{i_j}`$ is never an element of $`𝒩_{j1}`$, since $`F_{j1}r_{i_j}<0`$. It is easy to check directly that reflection by a simple root $`r`$ takes every positive root $`r^{}r`$ to another positive root. Thus each element of $`𝒩_j`$ is a positive root for each $`j`$.
Since $`F_lA`$, $`F_l`$ is a nonnegative integer linear combination of the classes in $`G`$. Moreover, the intersection of each of these classes with every element of $`𝒩_l`$ is nonnegative, since every element of $`𝒩_l`$ is a positive root. Now let $`w=s_{i_l}\mathrm{}s_{i_1}`$; then $`w^1F_l=F`$ and $`w^1H`$ meets every element of $`w^1𝒩_l=𝒩`$ nonnegatively for each $`HG`$. Thus each $`w^1H`$ is nef, hence $`F`$ is an integer linear combination of nef elements in $`W_6G`$, as claimed.
Given a nef divisor $`F`$, we still need a way of verifying that $`\mu _F`$ has maximal rank. Our main tools for doing so involve quantities
$$q(F)=h^0(X,FE_1)\text{ and }l(F)=h^0(X,F(E_0E_1)),$$
and bounds on the dimension of the cokernel of $`\mu _F`$, defined in terms of quantities $`q^{}(F)=h^1(X,FE_1)`$ and $`l^{}(F)=h^1(X,F(E_0E_1))`$, introduced in \[H6\] and \[FHH\]. The following result is Lemma 2.2 of \[FHH\]. (There it is assumed that $`FE_1FE_i`$ for all $`i>1`$, but that is not needed in the proof.)
###### Lemma 2.6
Let $`X`$ be obtained by blowing up distinct points $`p_i\text{P}\text{2}`$, and let $`F`$ be the class of an effective divisor on $`X`$ with $`h^1(X,F)=0`$. Then $`\text{dim ker }\mu _Fq(F)+l(F)`$ and $`\text{dim cok }\mu _Fq^{}(F)+l^{}(F)`$.
###### Remark 2.7
The quantities $`q(F)`$, $`l(F)`$, $`q^{}(F)`$ and $`l^{}(F)`$ are defined in terms of $`E_1`$ and $`E_0E_1`$, but in fact $`E_j`$, $`j>0`$, can be used in place of $`j=1`$, since one can reindex the points.
###### Corollary 2.8
Let $`F`$ and $`G`$ be nef divisors on a surface $`X`$ obtained by blowing up 6 distinct points of P<sup>2</sup>. If $`q(F)>0`$, $`l(F)>0`$ and $`q^{}(F)+l^{}(F)=0`$, then dim cok $`\mu _{F+G}=0`$.
Proof. If more than three points are on a line, then the six points are contained in a conic, and the result follows by Theorem 3.1.2 of \[H4\]. If at most three points lie on a line, then, since there are at most six points and they are distinct, $`K_X`$ is nef. So now we may assume $`K_X`$ is nef.
That $`q(F)>0`$ implies $`q(F+G)>0`$ and $`l(F)>0`$ implies $`l(F+G)>0`$, are clear, since a sum of effective divisors is effective. By Lemma 2.1, $`G+F`$ is effective and $`h^1(X,G+F)=0`$, so by Lemma 2.6 we have $`\text{dim cok }\mu _{F+G}q^{}(F+G)+l^{}(F+G)`$. Thus it’s enough to show $`q^{}(F+G)=0`$ and $`l^{}(F+G)=0`$. By a direct check of the generators listed by Lemma 2.5, $`G`$ is a sum of prime divisors of arithmetic genus at most 1. Hence it is enough by induction to show $`q^{}(F+G)=0`$ and $`l^{}(F+G)=0`$ when $`G`$ is the class of such a curve A. But this follows from $`0𝒪_X(FC)𝒪_X(G+FC)𝒪_A(G+FC)0`$, taking $`C`$ to be $`E_1`$ (for $`q^{}`$) or $`E_0E_1`$ (for $`l^{}`$), since $`h^1(X,FC)=0`$ by hypothesis, and $`h^1(A,G+FC)=0`$. (We have $`A(G+FC)0`$ since $`G`$ is nef, hence $`h^1(A,G+FC)=0`$ if $`A`$ has genus 0, while $`G^2>0`$ holds in each case that $`A`$ has genus 1. Thus $`A(G+FC)>0`$ when the genus is 1, hence again $`h^1(A,G+FC)=0`$.)
Given a nef divisor $`F`$, Corollary 2.8 often applies, in which case $`\mu _{F+G}`$ is surjective for all nef $`G`$. However, not every nef class is an appropriate sum of the form $`F+G`$. In the situations that we will need to deal with, the set of those classes which are not of the appropriate form turns out to be the union of a finite set of exceptions (which we can handle by brute force) with sets of strings of the form $`F+iC`$ (which we can handle by induction on $`i`$).
In order to set up the machinery to carry out the induction, define $`\mathrm{\Gamma }(X)`$ to be the set of all nef classes which are not the sum of two nonzero nef classes. Then $`\mathrm{\Gamma }(X)`$ generates $`\text{NEF}(X)`$ as a subsemigroup (i.e., every element of $`\text{NEF}(X)`$ is a nonnegative integer linear combination of elements of $`\mathrm{\Gamma }(X)`$). For $`i>0`$, let $`\mathrm{\Gamma }_i(X)`$ be the set of all sums with exactly $`i`$ terms, where each term is an element of $`\mathrm{\Gamma }(X)`$. (So, for example, $`\mathrm{\Gamma }_1(X)=\mathrm{\Gamma }(X)`$.) Let $`S(X)`$ be the set of all nef classes $`F`$ such that either $`q(F)=0`$, $`l(F)=0`$ or $`l^{}(F)+q^{}(F)>0`$. Then let $`S_i(X)=S(X)\mathrm{\Gamma }_i(X)`$. By Corollary 2.8 we have $`S_{i+1}(X)S_i(X)+S_1(X)`$.
Thus to show $`\mu _F`$ has maximal rank for every nef class $`F`$, it is enough by Lemma 2.6 to show that $`\mu _F`$ has maximal rank for all $`FS_i(X)`$ for each $`i`$. One checks directly that $`\mu _F`$ has maximal rank for all $`FS_i(X)`$ for small values of $`i`$. (It turns out that it is never necessary to do this for $`i>5`$.) For larger values of $`i`$, one applies Lemma 2.9 (the value of $`k`$ in this lemma never ends up needing to be bigger than 2, although this is not obvious until after the fact) and Lemma 2.10. Also, it turns out that the inclusions $`S_{j+i}(X)\{F+iC_F:FS_j(X)\}`$ in Lemma 2.9 can be chosen to be equalities, but that is more than we will need.
###### Lemma 2.9
Suppose for some $`j`$ there exists a $`k`$ and for each $`FS_j(X)`$ a $`C_FS_1(X)`$ such that $`S_{j+i}(X)\{F+iC_F:FS_j(X)\}`$ for $`0ik`$ and such that whenever $`CS_1(X)`$ but $`CC_F`$, then $`F+kCS_{j+k}(X)`$. Then $`S_{j+i}(X)\{F+iC_F:FS_j(X)\}`$ holds for all $`i0`$.
Proof. By Corollary 2.8, if $`F+kCS_{j+k}(X)`$, then $`F+(k+1)CS_{j+k+1}(X)`$. Thus it is enough by induction to show $`S_{j+k+1}(X)\{F+(k+1)C_F:FS_j(X)\}`$. Say $`G^{}S_{j+k+1}(X)`$. Then $`G^{}=G+C`$, where $`GS_{j+k}(X)`$ and $`CS_1(X)`$. By hypothesis, $`G=F^{}+kC_F^{}`$ for some $`F^{}S_j(X)`$ and $`C_F^{}S_1(X)`$. Since $`G+CS_{j+k+1}(X)`$, it follows by Corollary 2.8 that $`F^{}+CS_{j+1}(X)`$. Let $`H=F^{}+C`$; then $`H=H^{}+C_H^{}`$ for some $`H^{}S_j(X)`$ and $`C_H^{}S_1(X)`$. Now, $`H^{}+kC_F^{}S_{j+k}(X)`$ (since $`H+kC_F^{}=G+CS_{j+k+1}(X)`$), but for $`DS_1(X)`$ we have by hypothesis that $`H^{}+kDS_{j+k}(X)`$ unless $`D=C_H^{}`$. Thus $`C_F^{}=C_H^{}`$, so $`G+C=H^{}+(k+1)C_H^{}\{F+(k+1)C_F:FS_j(X)\}`$, so $`S_{j+k+1}(X)\{F+(k+1)C_F:FS_j(X)\}`$.
###### Lemma 2.10
Let $`X`$ be a blow up of P<sup>2</sup> at 6 distinct points. Let $`F`$ be a nef divisor such that $`\mu _F`$ is surjective, and let $`CX`$ be the class of a smooth rational curve such that $`C^20`$ and $`(F+C)C\text{max}(CE_1,C(E_0E_1))`$. Then $`\mu _{F+C}`$ is surjective.
Proof. Let $`\mathrm{\Lambda }`$ denote $`H^0(X,E_0)`$, and apply the snake lemma to:
$$\begin{array}{ccccccccc}0& & H^0(X,F)\mathrm{\Lambda }& & H^0(X,F+C)\mathrm{\Lambda }& & H^0(C,𝒪_X(F+C)𝒪_C)\mathrm{\Lambda }& & 0\\ & & \mu _1& & \mu _2& & \mu _3& & \\ 0& & H^0(X,F+E_0)& & H^0(X,F+C+E_0)& & H^0(C,𝒪_X(F+C+E_0)𝒪_C)& & 0\end{array}$$
Since $`\mu _F=\mu _1`$ is onto, it is enough to show $`\mu _3`$ is onto also, for which we apply $`(F+C)C\text{max}(CE_1,C(E_0E_1))`$, using the criterion given in \[F2\] (note also \[F3\]).
We will be interested mostly in those $`X`$ such that $`2E_0E_1\mathrm{}E_6`$ is not the class of an effective divisor, since otherwise (i.e., when the points $`p_i`$ lie on a conic, possibly reducible or nonreduced) $`\mu _F`$ is surjective whenever $`F`$ is nef by Theorem 3.1.2 of \[H4\], which in turn depends on Lemma 2.5 of \[H4\]. However, some details were left out of the published proof of this lemma, so we present it here in full. The extra details are indicated by indentation.
###### Lemma 2.11
Let $`X`$ be a smooth projective rational surface, and let $`𝒩`$ be the class of an effective divisor $`N`$ on $`X`$ such that $`h^0(X,𝒩+K_X)=0`$. If $``$ and $`𝒢`$ are the restrictions to $`N`$ of divisor classes $`^{}`$ and $`𝒢^{}`$ on $`X`$ which meet each component of $`N`$ nonnegatively, then $`𝒮(,𝒢)=0`$, where $`𝒮(,𝒢)`$ denotes the cokernel of the natural map $`H^0(N,)H^0(N,𝒢)H^0(N,+𝒢)`$.
Proof. To prove the lemma, induct on the sum $`n`$ of the multiplicities of the components of $`N`$. By Lemma II.9 of \[H3\], $`h^1(N,𝒪_N)=0`$ and every component of $`N`$ is a smooth rational curve. Thus the case $`n=1`$ is trivial (since then $`N=\text{P}\text{1}`$, and the space of polynomials of degree $`f`$ in two variables tensor the space of polynomials of degree $`g`$ in two variables maps onto the space of polynomials of degree $`f+g`$). So say $`n>1`$.
As in the proof of Theorem 1.7 of \[A\], $`N`$ has a component $`C`$ such that $`(NC)C1`$. Let $`L`$ be the effective divisor $`NC`$ and let $``$ be its class. Thus we have an exact sequence $`0𝒪_C()𝒪_N𝒪_L0`$.
> To see this, apply the snake lemma to
>
> $$\begin{array}{ccccccccc}0& & 𝒪_X(N)& & 𝒪_X& & 𝒪_N& & 0\\ & & & & & & & & \\ 0& & 𝒪_X(L)& & 𝒪_X& & 𝒪_L& & 0\end{array}$$
>
> to see that the kernel of $`𝒪_N𝒪_L`$ is just the cokernel of $`𝒪_X(N)𝒪_X(L)`$, which is just $`𝒪_C𝒪_X(L)`$, which we may write as $`𝒪_C(L)`$.
Now, $`LC1`$, and both $`^{}`$ and $`𝒢^{}`$ meet $`C`$ nonnegatively. We may assume $`^{}C𝒢^{}C`$, otherwise reverse the roles of $`^{}`$ and $`𝒢^{}`$. Since $`C=\text{P}\text{1}`$, we see that $`h^1(C,𝒪_C(^{}))`$, $`h^1(C,𝒪_C(𝒢^{}))`$ and $`h^1(C,𝒪_C(^{}+𝒢^{}))`$ all vanish. An argument similar to that used to prove Proposition II.3(a, b) of \[H4\] now shows that we have an exact sequence $`𝒮(𝒪_C(^{}),𝒪_C𝒢^{})𝒮(,𝒢)𝒮(𝒪_L,𝒪_L𝒢)0`$.
> What is actually clear here is that we have $`𝒮(𝒪_C(^{}),𝒢)𝒮(,𝒢)𝒮(𝒪_L,𝒢)0`$. Since $`h^1(C,𝒪_C(𝒢^{}))=0`$, we know $`𝒢𝒪_L𝒢`$ is surjective on global sections, and hence that $`𝒮(𝒪_L,𝒢)`$ is the same as $`𝒮(𝒪_L,𝒪_L𝒢)`$. What needs additional justification here is that $`𝒪_N𝒢^{}𝒪_C𝒢^{}`$ is surjective on global sections, so that we can conclude that $`𝒮(𝒪_C(^{}),𝒢)`$ is the same as $`𝒮(𝒪_C(^{}),𝒪_C𝒢^{})`$.
>
> Now, $`𝒩+K_X`$ is not the class of an effective divisor, and the same will remain true if we replace $`N`$ by any subscheme of $`N`$ obtained by subtracting off irreducible components of $`N`$. Thus any such resulting subscheme $`M`$ of $`N`$ has the property, like $`N`$ itself, that there is a component $`D`$ of $`M`$ such that $`(MD)D1`$. If $`M`$ is just $`N`$ with the reduced induced scheme structure, then by induction on the number of components of $`M`$ it follows (using Lemma II.9 of \[H3\]) that any two components of $`N`$ are smooth rational curves that are either disjoint or meet transversely at a single point, and no sequence $`B_1`$, $`\mathrm{}`$, $`B_i`$ of distinct components exists such that $`B_iB_1>0`$ and $`B_jB_{j+1}>0`$ for $`1j<i`$ (in particular, no three components meet at a single point, and the components of $`M`$ form a disjoint union of trees).
>
> First assume that $`N`$ is reduced; i.e. that $`N=N_{red}`$. Then $`C`$ is not a component of $`NC`$. Choose a section $`\sigma _C`$ of $`𝒪_C𝒢^{}`$, and for each of the other components $`B`$ of $`N`$, choose a section $`\sigma _B`$ of $`𝒪_B𝒢^{}`$ such that $`\sigma _B`$ does not vanish at any of the points where $`B`$ meets another component of $`N`$. (This is possible since $`B`$ is smooth and rational, so $`𝒪_B𝒢^{}`$ is $`𝒪_{\text{P}\text{1}}(d)`$ for some $`d0`$, so a section can always be chosen which does not vanish at any of a given finite set of points of $`B`$.) Since $`N`$ has no cycles and the components meet transversely, it is clear that starting from $`\sigma _C`$ one can patch together appropriate scalar multiples of the sections $`\sigma _B`$ to get a section $`\sigma `$ of $`𝒢`$ which restricts to $`\sigma _C`$. Thus $`𝒪_N𝒢^{}𝒪_C𝒢^{}`$ is surjective on global sections.
>
> Now assume that $`N`$ is not reduced. Let $`M`$ be the union of the components of $`N`$ which have multiplicity greater than 1 (taken with the same multiplicities as they have in $`N`$) together with those multiplicity 1 components of $`N`$ that meet one of these. No multiplicity 1 component $`B`$ of $`M`$ satisfies $`B(MB)1`$, so there must be a component $`B`$ of multiplicity more than 1 that does, and hence we also have $`B(NB)1`$ for some component $`B`$ of $`N`$ of multiplicity more than 1. Now from this and $`0𝒪_B(N+B)𝒢𝒪_N𝒢^{}𝒪_J𝒢^{}0`$, where $`J=NB`$, we see $`h^1(B,𝒪_B(N+B)𝒢)=0`$, so $`𝒪_N𝒢^{}𝒪_J𝒢^{}`$ is surjective on global sections. But $`J`$ still has $`C`$ as a component, because either $`C`$ has multiplicity 1 in $`N`$ (and hence $`CB`$), or $`C`$ has multiplicity more than 1 in $`N`$ (and so even if $`B=C`$, $`C`$ remains a component of $`NB=J`$). By induction on the number of components, we conclude that $`𝒪_N𝒢^{}𝒪_{N_{red}}𝒢^{}`$ is surjective on global sections. But $`C`$ is still a component of $`N_{red}`$, and $`𝒪_{N_{red}}𝒢^{}𝒪_C𝒢^{}`$ is surjective on global sections from above, hence so is $`𝒪_N𝒢^{}𝒪_C𝒢^{}`$.
Since $`𝒮(𝒪_L,𝒪_L𝒢)=0`$ by induction, it suffices to show $`𝒮(𝒪_C(^{}),𝒪_C𝒢^{})=0`$. If $`C(^{})0`$, then the latter is 0 (as in the previous paragraph). Otherwise, we must have $`0=^{}C=𝒢^{}C`$ and $`CL=1`$, so $`𝒪_C(1)=𝒪_C(^{})`$ and $`𝒪_C=𝒪_C𝒢^{}`$, which means $`h^0(𝒪_C,𝒪_C(^{}+𝒢^{}))=0`$ and hence again $`𝒮(𝒪_C(^{}),𝒪_C𝒢^{})=0`$.
## 3 The Main Results
In this section we first determine, up to permuting $`E_1,\mathrm{},E_6`$, which subsets of $`𝒬`$ occur as subsets of the form $`\text{neg}(X)`$, which by Corollary 2.3 is equivalent to determining the configuration types for six distinct points of P<sup>2</sup>. What we find is that the types are precisely those shown in Figure 1, where the classes of the proper transforms of the curves shown in a diagram of Figure 1 give the elements of $`\text{neg}(X)`$ for the corresponding configuration type. We then prove our main result, Theorem 3.1, and finish by explicitly answering, in the case of 6 points, the questions raised in \[GMS\].
To begin, note that the elements $`C`$ of $`\text{neg}(X)`$ satisfy the following three conditions: (i) $`C𝒬`$; (ii) $`C^2<1`$; and (iii) $`CD0`$ whenever $`C,D\text{neg}(X)`$ with $`CD`$.
First, if $`2E_0(E_1+\mathrm{}+E_6)\text{neg}(X)`$, then $`\{2E_0(E_1+\mathrm{}+E_6)\}=\text{neg}(X)`$. (For if $`C\text{neg}(X)`$ but $`C2E_0(E_1+\mathrm{}+E_6)`$, then $`C(2E_0(E_1+\mathrm{}+E_6))0`$ by (iii). But by direct check, every element $`C𝒬`$ with $`C^2<1`$ has $`C(2E_0(E_1+\mathrm{}+E_6))<0`$.) The case that $`\{2E_0(E_1+\mathrm{}+E_6)\}=\text{neg}(X)`$ corresponds to configuration type 11 in Figure 1. It is clear that this possibility actually occurs, since blowing up any six points on a smooth conic results in $`2E_0(E_1+\mathrm{}+E_6)\text{neg}(X)`$, and hence, as we just saw, $`\{2E_0(E_1+\mathrm{}+E_6)\}=\text{neg}(X)`$.
We now classify sets $`M`$ satisfying the conditions: (i) $`M`$ ; (ii) if $`CM`$, then $`C^2<1`$; and (iii) $`CD0`$ whenever $`C,DM`$, $`CD`$. For each such $`M`$, we also will show that there is an $`X`$ with $`M=\text{neg}(X)`$.
In fact, such a subset $`M`$ is just a matroid of rank 3 or less on a six point set, or, in the terminology of \[BCH\], it is a plane 6 point combinatorial geometry. It is not hard to work them all out, but \[BCH\] gives a complete list, saving us the trouble of doing so. The result corresponds precisely with what we show as configuration types 1 through 10 in Figure 1. So now we merely need to see that they all arise.
To show configuration type 1 occurs, we just need to show that one can pick 6 points such that no line passes through any 3 and no conic passes through all 6. Thus we can pick any two distinct points to be $`p_1`$ and $`p_2`$. Then $`p_3`$ can be any point not on the line through $`p_1`$ and $`p_2`$; $`p_4`$ can be any point not on any line through two of the first three points, and $`p_5`$ can be any point not on any line through two of the first four points. Finally, $`p_6`$ can be any point not on any line through two of the first five points nor on the conic through the first five points (of which there is only one). At each step we are allowed to choose any point avoiding a proper closed subset of P<sup>2</sup>. There is no obstruction to doing this, so configuration type 1 occurs.
For configuration type 2, we proceed as before, but the last point must be on exactly one of the lines through two of the previously chosen points. For example, we choose $`p_6`$ to be on the line $`L`$ through $`p_1`$ and $`p_2`$, but not on any other line through two of the previously chosen points. Thus the condition on our choice of $`p_6`$ is that we avoid finitely many points of $`L`$, which clearly we may do.
By similar reasoning, it is easy to check that each of the configurations 1 through 9 occur. With configuration 10, the same reasoning works to choose points $`p_1`$ through $`p_5`$, but the choice of $`p_6`$ is forced, since $`p_1,\mathrm{},p_5`$ uniquely determine $`p_6`$. Since we have no freedom in our choice of $`p_6`$, our previous argument is invalid at the last step. Instead, we take our six points to be the points of intersection of four general lines. Clearly, no four of the points can be collinear. So now we must check that the four lines are the only lines through any three of the points. Suppose there were a fifth line going through three of the points. Given any three of the six points of intersection of four general lines, it is easy to check that one of the four lines passes through two of the three points. So there can be no fifth line through any three of the points. Thus configuration type 10 also occurs.
(The foregoing justifications that the configurations actually occur may at first sight seem unnecessary. To show that they are not, we mention a similar example involving seven points. Suppose we want a configuration of six lines through seven points, arranged such that each line passes through exactly 3 points. Intuitively, we get such a configuration by taking three of the lines to be sides of an equilateral triangle, and the other three to be the angle bisectors. The seven points are the points where any two of the lines meet. This configuration occurs if and only if the ground field does not have characteristic 2. When the characteristic is 2, an additional line through the midpoints on the sides of the triangle is forced.)
We now prove our main result:
###### Theorem 3.1
Let $`X`$ be obtained by blowing up 6 distinct points of P<sup>2</sup>. Let $`E_0`$, $`E_1`$, $`\mathrm{}`$, $`E_6`$ be the corresponding exceptional configuration. Let $`F`$ be a nef divisor on $`X`$. Then $`\mu _F`$ has maximal rank.
Proof. We first consider the two extremes. If no line contains three of the points and no conic contains all 6, then the result follows by \[F1\]. This is the case in which the points are general. If all 6 points are on a conic, the result follows by Theorem 3.1.2 of \[H4\]. Note also that if 4 or more of the points are on a line, then all 6 are on a conic.
So now we are reduced to considering the case that some line contains three points, but no line contains 4 or more of the points and no conic contains all 6. Thus $`\text{neg}(X)`$ consists only of classes of the form $`LE_{i_1}E_{i_2}E_{i_3}`$, hence $`\text{neg}(X)=𝒩`$. If there is more than one line that contains three of the points, then any two such lines must share a point (otherwise all 6 points would lie on the two lines, which is a conic). It follows that the set $`𝒩`$ of nodal roots must, up to indexation of the points $`p_i`$, be one of the following:
$`\{r_0\}`$ — i.e., the first three points are on a line and no other set of three points is on a line (this case corresponds to configuration type 2);
$`\{r_0,E_0E_1E_4E_5\}`$ — i.e., two of the points are on one line, two on another, a fifth point occurs where the two lines meet, and the sixth point is not on any line through any two of the other points (this case corresponds to type 8);
$`\{r_0,E_0E_1E_4E_5,E_0E_3E_5E_6\}`$ — i.e., three lines form a triangle, with three of the points at the vertices, with an additional point on each line, but these last three points are not collinear (this case corresponds to type 9);
$`\{r_0,E_0E_1E_4E_5,E_0E_3E_5E_6,E_0E_2E_4E_6\}`$ — i.e., the 6 points are the points of intersection of four lines, no three of which meet at a single point (this case corresponds to type 10).
We now treat case (iv) in detail. The other cases (and the case that $`𝒩`$ is empty, which thereby recovers the result of \[F1\]) are similar. Using Remark 2.2, from $`𝒩=\text{neg}(X)`$ we determine that $`\text{NEG}(X)`$ consists of the following classes (where we list only the coefficients, so, for example, 1 0 -1 0 -1 0 -1 denotes $`E_0E_2E_4E_6`$):
```
0 1 0 0 0 0 0 1 0 0 -1 -1 0 0 1 -1 -1 -1 0 0 0
0 0 1 0 0 0 0 1 0 -1 0 0 -1 0 1 -1 0 0 -1 -1 0
0 0 0 1 0 0 0 1 -1 0 0 0 0 -1 1 0 0 -1 0 -1 -1
0 0 0 0 1 0 0 1 0 -1 0 -1 0 -1
0 0 0 0 0 1 0
0 0 0 0 0 0 1
```
Next we need to determine generators for $`\text{NEF}(X)`$. By Lemma 2.5, the set of all $`FW_6G`$ such that $`FC0`$ for all $`C𝒩=\{r_0,E_0E_1E_4E_5,E_0E_3E_5E_6,E_0E_2E_4E_6\}`$ generates $`\text{NEF}(X)`$, where $`W_6G`$ is the set of 1279 elements of the $`W_6`$ orbits of the elements of $`G`$ from Lemma 2.5. A tedious but easily coded check results in 212 generators. Many of these 212 are a sum of two other classes among the 212. Removing all classes which occur as such sums, we are left with 39, which therefore generate. Here is a list of these 39:
```
1 0 0 0 0 0 0 2 -1 0 -1 0 -1 0 3 0 0 -1 -2 -1 -1
2 0 -1 -1 -1 0 0 2 -1 0 0 -1 0 -1 3 -1 0 -2 -1 -1 0
2 0 0 -1 -1 -1 0 2 -1 -1 0 0 -1 0 3 0 -1 0 -1 -2 -1
2 0 0 0 -1 -1 -1 2 -1 0 -1 0 0 -1 3 -1 -1 -1 0 0 -2
2 -1 0 -1 -1 0 0 1 -1 0 0 0 0 0 3 -1 -1 -1 -2 0 0
2 0 -1 0 -1 -1 0 1 0 -1 0 0 0 0 3 -1 0 0 -1 -1 -2
2 0 0 -1 -1 0 -1 1 0 0 -1 0 0 0 3 0 -1 -2 -1 0 -1
2 0 -1 0 0 -1 -1 1 0 0 0 -1 0 0 3 -1 -2 0 -1 -1 0
2 0 -1 -1 0 0 -1 1 0 0 0 0 -1 0 3 0 -2 -1 0 -1 -1
2 -1 0 0 0 -1 -1 1 0 0 0 0 0 -1 3 -1 -1 -1 0 -2 0
2 -1 -1 0 0 0 -1 2 0 -1 -1 -1 -1 0 3 -2 0 -1 0 -1 -1
2 -1 -1 0 -1 0 0 2 -1 -1 0 0 -1 -1 3 -2 -1 0 -1 0 -1
2 0 -1 -1 0 -1 0 2 -1 0 -1 -1 0 -1 3 -1 -1 -1 -1 -1 -1
```
For each of these classes $`F`$, we find (after reindexing if need be, as discussed in Remark 2.7, but using the same indexing for $`q`$, $`l`$, $`q^{}`$, and $`l^{}`$) that $`q^{}(F)=0=l^{}(F)`$, hence $`\mu _F`$ is surjective by Lemma 2.6 and Remark 2.7. For example, to see how to compute these quantities, consider $`q^{}(F)`$ for $`F=3E_0E_12E_3E_4E_5`$ from the list above. Then, applying Remark 2.7, we reindex so that $`q(F)=h^0(X,FE_3)`$ and $`l(F)=h^0(X,F(E_0E_3))`$, etc. Since $`r_0\text{NEG}(X)`$ and $`r_0(FE_3)<0`$, we see $`h^0(X,FE_3)=h^0(X,FE_3r_0)`$. But now $`E_2(FE_3r_0)<0`$, so now $`h^0(X,FE_3)=h^0(X,FE_3r_0E_2)`$. Continuing in this way we eventually find that $`h^0(X,F)=\mathrm{}=h^0(X,0)=1`$, hence $`q(F)=1`$. Riemann-Roch now states that $`q(F)q^{}(F)=((FE_3)^2(FE_3)K_X)/2+1=1`$, so $`q^{}(F)=0`$.
Of the 39, all but the following 9 have both $`q`$ and $`l`$ positive, and thus $`S_1(X)`$ is just the set of these 9:
```
1 -1 0 0 0 0 0 1 0 0 0 -1 0 0 2 0 -1 -1 -1 -1 0
1 0 -1 0 0 0 0 1 0 0 0 0 -1 0 2 -1 -1 0 0 -1 -1
1 0 0 -1 0 0 0 1 0 0 0 0 0 -1 2 -1 0 -1 -1 0 -1
```
In each of these cases $`q=0`$. By Corollary 2.8, $`\mu _F`$ is surjective for all nef $`F`$ except possibly those in the subsemigroup generated by these last 9. A direct check shows that the conditions of Lemma 2.9 apply here with $`k=2`$ and $`C_F=F`$, so $`S_i(X)\{iF:FS_1(X)\}`$ for all $`i`$. Surjectivity for $`\mu _{iF}`$ for each $`F`$ and $`i`$ follows by direct check that $`q^{}(iF)=0=l^{}(iF)`$ when $`i=1`$ and 2, and then for all $`i>0`$ by applying Lemma 2.10.
Cases (i), (ii) and (iii) are handled the same way, thereby proving Theorem 3.1. For case (i), $`S_1(X)`$ has 55 elements, $`S_2(X)`$ has 90 elements, and $`S_i(X)`$ has 93 elements for $`i>2`$. Lemma 2.9 applies for $`j=3`$ with $`k=2`$, although this time it is not always true that $`F`$ is a multiple of $`C_F`$. For example, $`F=7E_02E_1\mathrm{}2E_55E_6S_3(X)`$, but $`C_F=3E_01E_1\mathrm{}1E_52E_6`$. For case (ii), $`S_1(X)`$ has 37 elements, $`S_i(X)`$ has 34 elements for $`i>1`$ and Lemma 2.9 applies for $`j=2`$ with $`k=2`$. For case (iii), $`S_1(X)`$ has 22 elements, $`S_i(X)`$ has 12 elements for $`i>1`$ and Lemma 2.9 applies for $`j=2`$ with $`k=2`$. (For the case that $`𝒩`$ is empty, $`S_1(X)`$ has 159 elements, $`S_2(X)`$ has 301 elements, and $`S_i(X)`$ has 316 elements for $`i>2`$. Lemma 2.9 applies for $`j=3`$ with $`k=2`$. Lemma 2.10 then gives the result except for multiples of $`F=5E_02E_1\mathrm{}2E_6`$, since $`\mu _F`$ is injective, and $`l^{}(mF)=1`$ for $`m0`$. Thus one must show ad hoc that $`\mu _{2F}`$ is surjective (see \[F1\]); then Lemma 2.10 applies to show that $`\mu _{mF}`$ is surjective for all $`m>2`$.)
###### Example 3.2
We work out an example to show how to determine the Hilbert function and graded Betti numbers of the ideal of a fat point subscheme. Assume the points are arranged as in case (iv); that is, configuration type 10. Assume the points are indexed so that a line passes through points 1, 2 and 3, and through 1, 4 and 5, and 2, 4 and 6 and 3, 5 and 6. Let $`Z=2p_1+2p_2+6p_3+2p_4+2p_5+2p_6`$. The associated divisor class for degree $`i`$ is $`F(Z,i)=iE_0(2E_1+2E_2+6E_3+2E_4+2E_5+2E_6)`$. Computing $`h_Z(i)=h^0(X,F(Z,i))`$ as in Remark 2.4, we find $`h_Z(5)=0`$, $`h_Z(6)=1`$, $`h_Z(7)=4`$, $`h_Z(8)=11`$, $`h_Z(9)=19`$ and $`h_Z(10)=30`$, so $`\alpha (Z)=6`$. Also, $`h^1(X,F(Z,8))>0`$ but $`h^1(X,F(Z,9))=0`$, hence the regularity $`\sigma (Z)`$ is 10.
Thus $`t_i=0`$ for $`i<\alpha (Z)=6`$ and for $`i>\sigma (Z)=10`$, and since $`h_Z(6)=1`$, we see $`t_6=1`$ and that $`\mu _{F(Z,6)}`$ is injective so $`t_7=h_Z(7)3h_Z(6)=1`$. To find $`t_8`$, note that: $`F(Z,7)C_1<0`$, where $`C_1=E_0E_3E_4`$; $`(F(Z,7)C_1)C_2<0`$ for $`C_2=r_0`$; $`(F(Z,7)C_1C_2)C_3<0`$ for $`C_3=E_0E_3E_4E_5`$; $`(F(Z,7)C_1C_2C_3)C_2<0`$; $`(F(Z,7)C_12C_2C_3)C_3<0`$; and $`F(Z,7)C_12C_22C_3`$ is nef. Thus the divisor class of fixed components of $`F(Z,7)`$ is $`C_1+2C_2+2C_3=5E_02E_12E_25E_3E_42E_52E_6`$, so $`Z_7^{}=2p_1+2p_2+5p_3+p_4+2p_5+2p_6`$, $`d_{Z,7}=5`$, and $`Z_7^+=p_3+p_4`$. Now we have $`t_8=\text{dim cok}(\mu _{F(Z,7)C_12C_22C_3})+(h^0(X,E_0+F(Z,7))h^0(X,E_0+F(Z,7)C_12C_22C_3))`$. But $`F(Z,7)C_12C_22C_3`$ is nef, its $`\mu `$ is onto by Theorem 3.1, and $`h^0(X,E_0+F(Z,7))h^0(X,E_0+F(Z,7)C_12C_22C_3)=h_Z(8)h^0(X,E_0+F(Z,7)C_12C_22C_3)=118=3`$. Similarly, $`t_9=0`$ and $`t_{10}=2`$. From the triple difference $`\mathrm{\Delta }^3h_Z`$, we find $`s_i=0`$ except for $`s_8=1`$, $`s_9=3`$ and $`s_{11}=2`$. Thus the minimal free resolution of $`I_Z`$ is $`0F_1F_0I_Z0`$ where $`F_0=R[6]R[7]R[8]^3R[10]^2`$ and $`F_1=R[8]R[9]^3R[11]^2`$.
It is easy to implement the procedure demonstrated in Example 3.2 as, for example, an awk script. We did so; the resulting script can be run over the web by visiting
http://www.math.unl.edu/$``$bharbour/6ptres/6reswebsite.html .
We used it to determine the Hilbert functions and graded Betti numbers for the ideals defining $`Z=p_1+\mathrm{}+p_6`$ and for $`2Z=2p_1+\mathrm{}+2p_6`$ for each of the 11 configuration types, thereby answering in the case of six points the questions raised in \[GMS\]. We could just as easily run $`mZ`$ for any $`m`$ or for any multiplicities $`m_1p_1+\mathrm{}+m_6p_6`$, if we wished to answer the questions raised by \[GMS\] not only for double points but for points of any given multiplicities. Note that for configuration types 5, 7 and 11, $`Z`$ is a complete intersection, and thus the Hilbert function and graded Betti numbers for $`mZ`$ are already known for all $`m`$ (see, for example, \[BGV1\] and \[BGV2\]). Also, the Hilbert function and graded Betti numbers for $`m_1p_1+\mathrm{}+m_6p_6`$ for any $`m_i`$ are known by \[F1\] for configuration type 1, and by \[H4\] for configurations 3, 4, 5, 6, 7 and 11 (since the points are contained in a conic). Results for configuration types 2, 8, 9 and 10 are new.
For ease of comparison with results of \[GMS\], we give the Hilbert functions $`h_{R/I(Z)}`$ of $`R/I(Z)`$, rather than for $`I(Z)`$. The Hilbert function of $`R/I(Z)`$ in degree 0 is always 1, and then it increases until it achieves the value $`\text{deg}(Z)`$, at which point it becomes constant. In each case we show the value $`h_{R/I(Z)}(t)`$ of the Hilbert function in each degree $`t0`$ until it becomes constant.
Here are the results. There are four different Hilbert functions for $`Z`$, and all together there are six different Hilbert functions for $`2Z`$, two whose support has one Hilbert function, two whose support has another, and one each for the remaining two cases. Note that for each Hilbert function for $`Z`$, there is among the Hilbert functions for $`2Z`$ both a maximum and minimum Hilbert function.
Scheme Type(s) $`h_{R/I(Z)}`$
$`Z`$ 1, 2, 8, 9, 10 1, 3, 6, 6
$`2Z`$ 1, 2, 8, 9 1, 3, 6, 10, 15, 18, 18
$`2Z`$ 10 1, 3, 6, 10, 14, 18, 18
$`Z`$ 3, 6, 7, 11 1, 3, 5, 6, 6
$`2Z`$ 3, 6 1, 3, 6, 10, 14, 16, 17, 18, 18
$`2Z`$ 7, 11 1, 3, 6, 10, 14, 17, 18, 18
$`Z`$ 4 1, 3, 4, 5, 6, 6
$`2Z`$ 4 1, 3, 6, 10, 12, 14, 15, 16, 17, 18, 18
$`Z`$ 5 1, 2, 3, 4, 5, 6, 6
$`2Z`$ 5 1, 3, 5, 7, 9, 11, 13, 14, 15, 16, 17, 18, 18
Scheme Type(s) $`F_1`$ $`F_0`$
$`Z`$ 1, 2, 8, 9, 10 $`R[4]^3`$ $`R[3]^4`$
$`2Z`$ 1, 2 $`R[7]^3`$ $`R[6]R[5]^3`$
$`2Z`$ 8 $`R[7]^3R[6]`$ $`R[6]^2R[5]^3`$
$`2Z`$ 9 $`R[7]^3R[6]^2`$ $`R[6]^3R[5]^3`$
$`2Z`$ 10 $`R[7]^4`$ $`R[6]^4R[4]`$
$`Z`$ 3, 6 $`R[5]R[4]`$ $`R[4]R[3]R[2]`$
$`2Z`$ 3 $`R[9]R[7]R[6]`$ $`R[8]R[5]^2R[4]`$
$`2Z`$ 6 $`R[9]R[7]R[6]^2`$ $`R[8]R[6]R[5]^2R[4]`$
$`Z`$ 7, 11 $`R[5]`$ $`R[3]R[2]`$
$`2Z`$ 7, 11 $`R[8]R[7]`$ $`R[6]R[5]R[4]`$
$`Z`$ 4 $`R[6]R[3]`$ $`R[5]R[2]^2`$
$`2Z`$ 4 $`R[11]R[7]R[5]^2`$ $`R[10]R[6]R[4]^3`$
$`Z`$ 5 $`R[7]`$ $`R[6]R[1]`$
$`2Z`$ 5 $`R[13]R[8]`$ $`R[12]R[7]R[2]`$
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# Homotopical dynamics in symplectic topology.
## 1. Introduction
The main purpose of this paper is to survey a number of Morse-theoretic results which show how to estimate algebraically the high-dimensional moduli spaces of Morse flow lines and to describe some of their recent applications to symplectic topology. We also deduce some new applications.
The paper starts with a brief discussion of the various proofs showing that the differential in the Morse complex is indeed a differential. With this occasion we introduce the main concepts in Morse theory and, in particular, the notion of connecting manifold (or, equivalently, the moduli space of flow lines connecting two critical points) which is the main object of interest in our further constructions. Moreover, an extension of one of these proofs leads naturally to an important result of John Franks which describes the framed cobordism class of connecting manifolds between consecutive critical points as a certain relative attaching map. After describing Franks’ result, we proceed to a stronger result initially proved in which computes a framed bordism class naturally associated to the same connecting manifolds in terms of certain Hopf invariants. While these results only apply to consecutive critical points we then describe a recent method to estimate general connecting manifolds by means of the Serre spectral sequence of the path-loop fibration having as base the ambient manifold . Some interesting topological consequences of these results are briefly mentioned as well as some other methods used in the study of these problems.
The third section discusses a number of symplectic applications. We start with some results which first appeared in . These use the non-vanishing of certain Hopf invariants to deduce the existence of bounded orbits of hamiltonian flows (obviously, inside non-compact manifolds). This is a very “soft” type of result even if difficult to prove. We then continue in §3.2 by describing how to use the Serre spectral sequence result to detect pseudo-holomorphic strips as well as some consequences of the existence of the strips. Most of the results of this part have first appeared in but there are some that are new: we discuss explicitly the detection of pseudoholomorphic strips passing through some submanifold and we present a way to construct in a coherent fashion our theory for lagrangians in general symplectic manifolds as long as we remain under a bubbling threshold. Notice that even the analogue of the classical Floer theory (which is a very particular case of our construction) has not been explicited in the literature in the Lagrangian case even if all the necessary ideas are present in some form - see for the hamiltonian case.
The paper contains a number of open problems and ends with a conjecture which is supported by the results in §3.2 as well as by recent joint results of the second author with François Lalonde.
## 2. Elements of Morse theory
Assume that $`M`$ is a compact, smooth manifold without boundary of dimension $`n`$. Let $`f:M`$ be a smooth Morse function and let $`\gamma :M\times M`$ be a negative gradient Morse-Smale flow associated of $`f`$. In particular, $`f`$ is strictly decreasing along any non-constant flow line of $`\gamma `$ and the stable manifolds
$$W^s(P)=\{xM:\underset{t\mathrm{}}{lim}\gamma _t(x)=P\}$$
and the unstable manifolds
$$W^u(Q)=\{xM:\underset{t\mathrm{}}{lim}\gamma _t(x)=Q\}$$
of any pair of critical points $`P`$ and $`Q`$ of $`f`$ intersect transversally. One of the most useful and simple tools that can be defined in this context is the Morse complex
$$C(\gamma )=(/2<Crit(f)>,d).$$
Here $`/2<S>`$ is the $`/2`$-vector space generated by the set $`S`$, the vector space $`/2<Crit(f)>`$ has a natural grading given by $`|P|=ind_f(P),PCrit(f)`$ and $`d`$ is the differential of the complex which is defined by
$$dx=\underset{|y|=|x|1}{}a_y^xy$$
so that the coefficients $`a_y^x=\mathrm{\#}((W^u(x)W^s(y))/)`$. This definition makes sense because the set $`W^u(P)W^s(Q)`$ which consists of all the points situated on some flow line joining $`P`$ to $`Q`$ is invariant by the $``$-action given by the flow. Moreover, $`W^u(P)`$ and $`W^s(P)`$ are homeomorphic to open disks which implies that the set $`M_Q^P=(W^u(P)W^s(Q))/`$ has the structure of a smooth (in general, non-compact) manifold of dimension $`|P||Q|1`$. We call this space the moduli space of flow lines joining $`P`$ to $`Q`$. It is not difficult to understand the reasons for the non-compactness of $`M_Q^P`$ when $`M`$ is compact as in our setting: this is simply due to the fact that a family of flow lines joining $`P`$ to $`Q`$ may approach a third, intermediate, critical point $`R`$. For this to happen it is necessary (and sufficient - see Smale or Franks ) to have some flow line which joins $`P`$ to $`R`$ and some other joining $`R`$ to $`Q`$. This implies that when $`|P|=|Q|+1`$ the set $`M_Q^P`$ is compact and thus the sum above is finite. For further use let’s define also the unstable sphere of a critical point $`P`$ as $`S_a^u(P)=W^u(P)f^1(a)`$ as well as the stable sphere $`S_a^s(P)=W^s(P)f^1(a)`$ where $`a`$ is a regular value of $`f`$. It should be noted that this names are slightly abusive as these two sets are spheres, in general, only if $`a`$ is sufficiently close to $`f(P)`$. In that case $`S_a^u(P)`$ is homeomorphic to a sphere of dimension $`|P|1`$ and $`S_a^s(P)`$ is homeomorphic to a sphere of dimension $`n|P|1`$. With this notations ths moduli space $`M_Q^P`$ is homeomorphic to $`S_a^u(P)S_a^s(P)`$ for any $`a(f(Q),f(P))`$ which is a regular value of $`f`$.
The main properties of the object defined above are that:
$$d^2=0\mathrm{and}H_{}(C(\gamma ))H_{}(M;/2).$$
We will sometimes denote this complex by $`C(f)`$ and will call it the (classical) Morse complex of $`f`$. The flow $`\gamma `$ may be in fact even a pseudo- (negative) gradient flow of $`f`$. There also exists a version of this complex over $``$ in which the counting of the elements in $`M_y^x`$ takes into account orientations.
There are essentially four methods to prove these properties:
* Deducing the equation $`_ya_y^xa_z^y=0`$ (which is equivalent to $`d^2=0`$) from the properties of the moduli spaces $`M_z^x`$ with $`|x||z|=2`$.
* Comparing $`a_y^x`$ with a certain relative attaching map.
* Expressing $`a_y^x`$ in terms of a connection map in Conley’s index theory.
* A method based on a deformation of the de Rham complex (clearly, in this case the coefficients are required to be in $``$).
For the rest of this paper the two approaches that are of the most interest are (i) and (ii). Therefore we shall first say a few words on the other two methods and will then describe in more detail the first two. Method (iii) consists in regarding two critical points $`x`$, $`y`$ so that $`|x|=|y|+1`$ as an attractor-repellor pair and to apply the general Conley theory of Morse decompositions to this situation . Method (iv) has been introduced by Witten in and is based on a deformation of the differential of the de Rham complex which provides a new differential with respect to which the harmonic forms are in bijection with the critical points of $`f`$.
Method (i) has been probably folklore for a long time but it first appeared explicitly in Witten’s paper. It is based on noticing that the moduli space $`M_Q^P`$ admits a compactification $`\overline{M}_Q^P`$ which is a compact, topological manifold with boundary so that the boundary verifies the formula:
(1)
$$\overline{M}_Q^P=\underset{R}{}\overline{M}_R^P\times \overline{M}_Q^R.$$
There are two main ways to prove this formula. One is analytic and regards a flow line from $`P`$ to $`Q`$ as a solution of a differential equation $`\dot{x}=f(x)`$ and studies the properties of such soultions (this method is presented in the book of Schwarz ). A second approach is more topological/dynamical in nature as is described in detail by Weber . Clearly, from formula (1) we immediately deduce $`_ya_y^xa_z^y=0`$ and hence $`d^2=0`$. Just a little more work is needed to deduce from here the second property.
Method (ii) was the one best known classically and it is essentially implicit in Milnor’s $`h`$-cobordism book . It is based on the observation that $`a_y^x`$ can be viewed as follows. First, to simplify slightly the argument assume that the only critical points in $`f^1([f(y),f(x)])`$ are $`x`$ and $`y`$. It is well known that for $`a(f(y),f(x))`$, there exists a deformation retract
$$r:M(a)=f^1(\mathrm{},a]M(f(y)ϵ)_{\varphi _y}D^{|y|}=M^{}$$
where the attaching map
(2)
$$\varphi _y:S_{f(y)ϵ}^u(y)M(f(y)ϵ)$$
is just the inclusion and $`ϵ`$ is small. This deformation retract follows the flow till reaching of $`U(W^u(y))M(f(y)ϵ)`$ where $`U(W^u(y)`$ is a tubular neighbourhood of $`W^u(y)`$ so that the flow is transverse to its boundary and then collapses this neighbourhood to $`W^u(y)`$ by the canonical projection. Clearly, applying this remark to each critical point of $`f`$ provides a $`CW`$-complex of the same homotopy type as that of $`M`$ and with one cell $`\overline{x}`$ for each element of $`xCrit(f)`$. To this cellular decomposition we may associate a celullar complex $`(C^{}(f),d^{})`$ with the property that $`d\overline{x}=k_{\overline{y}}^{\overline{x}}\overline{y}`$ where $`k_{\overline{y}}^{\overline{x}}`$ is, by definition, the degree of the composition:
(3)
$$\psi _y^x:S_a^u(x)\stackrel{\varphi _x}{}M(a)\stackrel{r}{}M^{}\stackrel{u}{}M^{}/M(f(y)ϵ)$$
with $`\varphi _x:S_a^u(x)M(a)`$ again the inclusion and where the last map, $`u`$, is the projection onto the respective topological quotient space (which is homeomorphic to the sphere $`S^{|y|}`$).
Notice now that $`M_y^xS_a^u(x)`$ is a finite union of points say $`P_1,\mathrm{},P_k`$. Imagine a small disk $`D_iS_a^u(x)`$ around $`P_i`$. The key (but geometrically clear) remark is that the composition of the flow $`\gamma `$ together with the retraction $`r`$ transports $`D_i`$ (if it is chosen sufficiently small) homeomorphically onto a neighbourhood of $`y`$ inside $`W^u(y)`$. Therefore, the degree of $`deg(\psi _y^x)=a_y^x`$ and thus $`d=d^{}`$ which shows that $`d`$ is a differential and that the homology it computes agrees with the homology of $`M`$.
As we shall see further, the points of view reflected in the approaches at (i) and (ii) lead to interesting applications which go much beyond “classical” Morse theory. Method (iv), while striking and inspiring appears for now not to have been exploited efficiently.
### 2.1. Connecting Manifolds
One way to look to the Morse complex is by viewing the coefficients $`a_y^x`$ of the differential as a measure of the $`0`$-dimensional manifold $`M_y^x`$. The question we discuss here is in what way we can measure algebraically the similar higher dimensional moduli spaces. This is clearly a significant issue because, obviously, only a very superficial part of the dynamics of the negative gradient flow of $`f`$ is encoded in the $`0`$-dimensional moduli spaces of connecting flow lines.
As a matter of terminology, the space $`M_Q^P`$ when viewed inside the unstable sphere $`S_a^u(P)`$ (with $`f(P)a`$ positive and very small) is also called the connecting manifold of $`P`$ and $`Q`$.
It was mentioned above that, in general, a connecting manifold $`M_Q^P`$ is not closed. However, if the critical points $`P`$ and $`Q`$ are consecutive in the sense that there does not exist a critical point $`R`$ so that $`M_R^P\times M_Q^R\mathrm{}`$, then $`M_Q^P`$ is closed.
#### 2.1.1. Framed Cobordism Classes
An important remark of John Franks is that connecting manifolds are canonically framed. First recall that a framed manifold $`V`$ is a submanifold $`VS^n`$ which has a trivial normal bundle together with a trivialization of this bundle. Two such trivializations are equivalent (and will generally be identified) if they are restrictions of a trivialization of the normal bundle of $`V\times [0,1]`$ inside $`S^n\times [0,1]`$. We also recall the Thom-Pontryagin construction in this context . Assuming $`VS^n`$ is framed we define a map
$$\varphi _V:S^nS^{codim(V)}$$
as follows: consider a tubular neighbourhood $`U(V)`$ of $`V`$, use the framing to define a homeomorphism $`\psi :U(V)D^k\times V`$ where $`D^k`$ is the closed disk of dimension $`k=codim(V)`$, consider the composition $`\psi ^{}:U(V)\stackrel{\psi }{}D^k\times V\stackrel{p_1}{}D^kD^k/S^{k1}=S^k`$ and define $`\varphi _V`$ by extending $`\psi ^{}`$ outside $`U(V)`$ by sending each $`xS^n\backslash U(V)`$ to the base point in $`D^k/S^{k1}=S^k`$. The homotopy class of this map is the same if two framings are equivalent. It is easy to see that two framed manifolds (of the same dimension) are cobordant iff their associated Thom maps are homotopic.
We return now to Franks’ remark and notice that the manifolds $`M_Q^P`$ are framed. First, we make the convention to view $`M_Q^P`$ as a submanifold of the unstable sphere of $`P`$, $`S_a^u(P)`$ (the other choice would have been to use $`S^s(Q)`$ as ambient manifold). Notice that we have
$$M_Q^P=S_a^u(P)W^s(Q)$$
and this intersection is transversal. Clearly, as $`W^s(Q)`$ is homeomorphic to a disk, its normal bundle in $`M`$ is trivial and any two trivializations of this bundle are equivalent. This implies that the normal bundle of $`M_Q^PS_a^u(P)`$ is also trivial and a trivialization of the normal bundle of $`W^s(Q)`$ provides a trivialization of this bundle which is unique up to equivalence.
Recall that if $`P`$ and $`Q`$ are consecutive critical points, then $`M_Q^P`$ is closed. As we have seen that it is also framed we may associate to it a framed cobordism class
$$\stackrel{~}{M_Q^P}\pi _{|P|1}(S^{|Q|}).$$
Moreover, it is easy to see that the function $`f`$ may be perturbed without modifying the dynamics of the the negative gradient flow so that the cell attachments corresponding to the critical points $`Q`$ and $`P`$ are in succession. Therefore the map
$$\psi _Q^P:S^{|P|1}S^{|Q|}$$
defined as in formula (3) is still defined. The main result of Franks in is:
###### Theorem 2.1.
Assume $`P`$ and $`Q`$ are consecutive critical points of $`f`$. Up to sign $`\stackrel{~}{M_Q^P}`$ coincides with the homotopy class of $`\psi _Q^P`$.
The idea of proof of this result is quite simple. All that is required is to make even more precise the constructions used in the approach (ii) used to show $`d^2=0`$ for the Morse complex. For this we fix for $`W^s(Q)`$ a normal framing $`o`$ which is invariant by translation along the flow $`\gamma `$ and which at $`QW^s(Q)`$ is given by a basis $`e`$ of $`T_QW^u(Q)`$ (this is possible because $`W^u(Q)`$ and $`W^s(Q)`$ intersect transversally at $`Q`$). We also fix the tubular neighbourhood $`U(W^u(Q))`$ so that the projection $`r^{\prime \prime }:U(W^u(Q))W^u(Q)`$ has the property that $`(r^{\prime \prime })^1(Q)=W^s(Q)U(W^u(Q))`$ and, for any point $`y(r^{\prime \prime })^1(Q)`$, we have $`(r^{\prime \prime })_{}(o_y)=e`$. Moreover, we may assume that the normal bundle of $`M_Q^P`$ in $`S_a^u(P)`$ is just the restriction of the normal bundle of $`W^s(Q)`$ (in fact, the two are, in general, only isomorphic and not equal but this is just a minor issue).
Now, follow what happens with the framing of $`M_Q^P`$ along the composition $`ur`$. For this we write $`r=r^{\prime \prime }r^{}`$ where $`r^{}`$ follows the flow till reaching $`U(W^u(Q))`$. Now pick a point in $`xM_Q^P`$ together with its normal frame $`o_x`$ at $`x`$. After applying $`r^{}`$, the pair $`(x,o_x)`$ is taken to a pair $`(x^{},o_x^{})`$ with $`x^{}((r^{\prime \prime })^1(Q))`$. Applying now $`r^{\prime \prime }`$, the image of $`(x^{},o_x^{})`$ is $`(Q,e)`$. Take now $`V`$ a tubular neighbourhood of $`M_Q^PS_a^u(P)`$ together with an identification $`VD^{|Q|}\times M_Q^P`$ which is provided by the framing $`o`$. The argument above implies that if the constant $`ϵ`$ used to construct the map $`u:M(f(Q)ϵ)W^u(Q)=M^{}M^{}/M(f(Q)ϵ)`$ is very small, then the composition $`ur^{\prime \prime }r^{}`$ coincides with the relevant Thom-Pontryagin map.
#### 2.1.2. Framed Bordism Classes and Hopf invariants
It it natural to wonder whether, besides their framing, there are some other properties of the connecting manifolds which can be detected algebraically.
A useful point of view in this respect turns out to be the following: imagine the elements of $`M_Q^P`$ as path or loops on $`M`$. The fact that they are paths is obvious (we parametrize them by the value of $`f`$; the negative sign gives the flow lines the orientation coherent with the negative gradient) but they can be transformed into loops very easily. Indeed, fix a simple path in $`M`$ which joins all the critical points of $`f`$ and contract this to a point thus obtaining a quotient space $`\widehat{M}`$ which has the same homotopy type as $`M`$. Let $`q:M\widehat{M}`$ be the quotient map. We denote by $`\mathrm{\Omega }M`$ the space of based loops on $`M`$ and keep the notation $`q`$ for the induced map $`\mathrm{\Omega }M\mathrm{\Omega }\widehat{M}`$. This discussion shows that there are continuous maps
$$j_Q^P:M_Q^P\mathrm{\Omega }\widehat{M}.$$
These maps have first been defined and used in and they have some interesting properties.
For example, given such a map $`j_Q^P`$ and assuming that $`P`$ and $`Q`$ are consecutive it is natural to ask whether the homology class $`[M_Q^P]H_{|P||Q|1}(\mathrm{\Omega }M;)`$ is computable (here $`[M_Q^P]`$ is the image by $`j_Q^P`$ of the fundamental class of $`M_Q^P`$). We shall see that quite a bit more is indeed possible: the full framed bordism class associated to $`j_Q^P`$ and to the canonical framing on $`M_Q^P`$ can be expressed as a relative Hopf invariant.
To explain this result first recall that if $`VS^n`$ is framed and $`l:VX`$ is a continuous map with $`V`$ a closed manifold we may construct a richer Thom-Pontryagin map as follows. We again consider a tubular neighbourhood $`U(V)`$ of $`V`$ in $`S^n`$ together with an identification $`U(V)D^k\times V`$ where $`k=codim(V)`$ which is provided by the framing. We now define a map
$$\overline{\varphi }_V:U(V)\stackrel{j}{}D^k\times V\stackrel{id\times l}{}D^k\times X(D^k\times X)/(S^{k1}\times X)$$
where $`j:U(V)=D^k\times VV`$ is the projection and the last map is just the quotient map (which identifies $`S^{k1}\times X`$ to the base point). Notice that $`\overline{\varphi }_V(U(V))S^{k1}\times X`$. Therefore, we may extend the definition above to a map
$$\overline{\varphi }_V:S^n(D^k\times X)/(S^{k1}\times X)$$
by sending all the points in the complement of $`U(V)`$ to the base point. It is well-known (and a simple exercise of elementary homotopy theory) that there exists a (canonical) homotopy equivalence $`(D^k\times X)/(S^{k1}\times X)\mathrm{\Sigma }^k(X^+)`$ where $`\mathrm{\Sigma }X`$ is the (reduced) suspension of $`X`$, $`\mathrm{\Sigma }^i`$ is the suspension iterated $`i`$-times and $`X^+`$ is the space $`X`$ with an added disjoint point (notice also that $`\mathrm{\Sigma }^k(X^+)=\mathrm{\Sigma }^kXS^k`$ where $``$ denotes the wedge or the one point union of spaces). This allows us to view the map $`\overline{\varphi }_V`$ as a map with values in $`\mathrm{\Sigma }^k(X^+)`$. The framed bordism class of $`V`$ is simply the homotopy class $`[\overline{\varphi }_V]\pi _n(\mathrm{\Sigma }^k(X^+))`$. This is independent of the various choices made in the construction. Two pairs of data (framings included) $`(V,l)`$ and $`(V^{},l^{})`$ admit an extension to a manifold $`WS^n\times [0,1]`$ with $`W=V\times \{0\}V^{}\times \{1\}`$ iff $`\overline{\varphi }_V\overline{\varphi }_V^{}`$. Notice also that to an element $`\alpha \pi _n(\mathrm{\Sigma }^kX^+)`$ we may associate a homology class $`[\alpha ]H_{nk}(X)`$ obtained by applying the Hurewicz homomorphism, desuspending $`k`$-times and projecting on the $`H_{}(X)`$ term in $`H_{}(X^+)`$.
Returning now to our connecting manifolds $`M_Q^P`$ we again focus on the case when $`P`$ and $`Q`$ are consecutive. The map $`j_Q^P`$ together with the canonical framing provide a homotopy class $`\{M_Q^P\}\pi _{|P|1}(\mathrm{\Sigma }^{|Q|}(\mathrm{\Omega }M^+))`$ (to simplify notation we have replaced $`\widehat{M}`$ with $`M`$ here \- the two are homotopy equivalent).
As indicated above, it turns out that this class can be computed in terms of a relative Hopf invariant. We shall now discuss how this invariant is defined.
Assume that $`S^{q1}\stackrel{\alpha }{}X_0X^{}`$ and $`S^{p1}\stackrel{\beta }{}X^{}X^{\prime \prime }`$ are two successive cell attachments and that $`X^{\prime \prime }`$ is a subspace of some larger space $`X`$. In particular, $`X^{}=X_0_\alpha D^q`$, $`X^{\prime \prime }=X^{}_\beta D^p`$. Let $`SD^q`$ be the $`q1`$-sphere of radius $`1/2`$. There is an important map called the coaction associated to $`\alpha `$ which is defined by the composition
$$:X^{}X^{}/SS^qX^{}$$
where the first map identifies all the points of $`S`$ to a single one and the second is a homeomorphism (in practice it is convenient to also assume that all the maps and spaces involved are pointed and in that case we view $`D^q`$ as the reduced cone over $`S^{q1}`$).
We consider the composition
$$\psi (\beta ,\alpha ):S^{p1}\stackrel{\beta }{}X^{}\stackrel{}{}S^qX^{}\stackrel{idj}{}S^qX^{\prime \prime }S^qX$$
and notice that if $`p_2:S^qXX`$ is the projection on the second factor, then the composition $`p_2\psi (\beta ,\alpha )`$ is null-homotopic. This is due to the fact that this composition is homotopic to $`S^{p1}\stackrel{\beta }{}X^{}X^{\prime \prime }X`$ which is clearly null-homotopic. We now consider the map $`p_2:S^qXX`$. It is well-known in homotopy theory that any map may be transformed into a fibration. In our case this comes down to considering the free path fibration
$$t:\stackrel{~}{P}XX$$
where $`\stackrel{~}{P}X`$ is the set of all continuous path in $`X`$ parametrized by $`[0,1]`$, $`t(\gamma )=\gamma (0)`$. We take the pull back of this fibration over $`p_2`$. The total space $`\stackrel{~}{E}`$ of the resulting fibration has the same homotopy type as $`S^qX`$ and it is endowed with a canonical projection
$$\stackrel{~}{p}:\stackrel{~}{E}X$$
which replaces $`p_2`$, $`\stackrel{~}{p}(z,\gamma )=\gamma (1)`$. It is an exercise in homotopy theory to see that the fibre of the fibration $`\stackrel{~}{p}`$ is homotopic to $`\mathrm{\Sigma }^q((\mathrm{\Omega }X)^+)`$ and that, moreover, the inclusion of this fibre in the total space is injective in homotopy. As the composition $`p_2\psi (\beta ,\alpha )`$ is homotopically trivial, the homotopy exact sequence of the fibration $`\stackrel{~}{E}X`$ implies that $`\psi (\beta ,\alpha )`$ admits a lift to $`\overline{\psi }(\alpha ,\beta ):S^{p1}\mathrm{\Sigma }^q((\mathrm{\Omega }X)^+)`$ whose homotopy class does not depend on the choice of lift. We let $`H(\alpha ,\beta )\pi _{p1}(\mathrm{\Sigma }^q((\mathrm{\Omega }X)^+))`$ be equal to this homotopy class and we call it the relative Hopf invariant associated to the attaching maps $`\alpha `$ and $`\beta `$ (for a discussion of the relations between this Hopf invariants and other variants see Chapters 6 and 7 in ).
To return to Morse theory, recall from (2) that passing through the two consecutive critical points $`Q`$ and $`P`$ leads to two successive attaching maps $`\varphi _Q:S^{|Q|1}M(f(Q)ϵ)`$ and $`\varphi _P:S^{|P|1}M(f(P)ϵ)`$ (we assume again - as we may - that the set $`f^1([f(Q),f(P)])`$ does not contain any other critical points besides $`P`$ and $`Q`$). Moreover, as we know the inclusion $`M^{}=M(f(Q)ϵ)W^u(Q)M(f(P)ϵ)`$ is a homotopy equivalence. Therefore, the construction above can be applied to $`\varphi _Q`$ and $`\varphi _P`$ and it leads to a relative Hopf invariant
$$H(P,Q)\pi _{|P|1}(\mathrm{\Sigma }^{|Q|}(\mathrm{\Omega }M^+)).$$
With these constructions our statement is:
###### Theorem 2.2.
The homotopy class $`H(P,Q)`$ coincides (up to sign) with the bordism class $`\{M_Q^P\}`$. In particular, the homology class $`[M_Q^P]`$ equals (up to sign and desuspension) the Hurewicz image of $`H(P,Q)`$.
The proof of this result can be found in (a variant proved by a slightly different method appears in ). The proof is considerably more complicated than the proof of Theorem 2.1 so we will only present a rough justification here. To simplify notation we let $`M_0=M(f(Q)ϵ)`$. Let $`M_1=M(f(Q)ϵ)U(W^u(Q))`$. Recall, that the inclusions $`M^{}M_1M(f(P)ϵ)`$ are homotopy equivalences.
Let $`𝒫:\mathrm{\Omega }MPMM`$ be the path-loop fibration (of total space the based paths on $`M`$ and of fibre the space of based loops on $`M`$). We denote by $`E_0`$ the total space of the pull-back of the fibration $`𝒫`$ over the inclusion $`M_0M`$. Similarly, we let $`E_1`$ be the total space of the pull-back of $`𝒫`$ over the inclusion $`M_1M`$. The key remark is that the attaching map $`\varphi _P:S^u(P)M_1`$ admits a natural lift to a map $`\stackrel{~}{\varphi _P}:S^u(P)E_1`$. Indeed, we assume here that all the critical points are identified to the base point. The space $`E_1`$ consists of the based paths in $`M`$ that end at points in $`M_1`$. But each element of the image of $`\varphi _P`$ corresponds to precisely such a path which is explicitly given by the corresponding flow line (we need to use here Moore paths and loops which are paths parametrized by arbitrary intervals $`[0,a]`$ and not only the interval $`[0,1]`$). Consider the inclusion $`E_0E_1`$. It is not difficult to see that the quotient topological space $`E_1/E_0`$ admits a canonical homotopy equivalence $`\eta :E_1/E_0\mathrm{\Sigma }^{|Q|}(\mathrm{\Omega }M^+)`$. Therefore, we may consider the composition $`\eta ^{}=\eta \stackrel{~}{\varphi _P}`$. It is possible to show that this map $`\eta ^{}`$ is homotopic to $`H(P,Q)`$. At the same time, we see that the restriction of $`\stackrel{~}{\varphi _P}`$ to $`M_Q^P`$ coincides with $`j_Q^P`$. Moreover, by making explicit $`\eta `$ it is also possible to see that $`\stackrel{~}{\varphi _P}`$ is homotopic to the Thom-Pontryagin map associated to $`M_Q^P`$.
#### 2.1.3. Some topological applications
We now decribe a couple of topological applications of Theorem 2.2. The idea behind both of them is quite simple: the function $`f`$ is also a Morse function and the critical points $`Q`$, $`P`$ are conscutive critical points for $`f`$. Therefore, the connecting manifold $`M_P^Q`$ is well defined as well as its associated bordism homotopy class $`\{M_P^Q\}`$. Clearly, the underlying space for both $`M_Q^P`$ and $`M_P^Q`$ is the same. The map $`j_P^Q`$ is different from $`j_Q^P`$ just by reversing the direction of the loops. The two relevant framings may also be different. The relation between them is somewhat less straightforward but it still may be understood by considering $`M`$ embedded inside a high dimensional euclidean space and taking into account the twisting induced by the stable normal bundle. In all cases, this establishes a relationship between the two Hopf invariants $`H(P,Q)`$ and $`H(Q,P)`$.
A. The first application concerns the construction of examples of non-smoothable, simply-connected, Poincaré duality spaces. The idea is as follows: we construct Poincaré duality spaces which have a simple $`CW`$-decomposition and with the property that for certain two successively attached cells $`e,f`$ the resulting Hopf invariant $`H`$ and the Hopf invariant $`H^{}`$ associated to the dual cells $`f^{},e^{}`$ are not related in the way described above. If the respective Poincaré duality space is smoothable, then the given cell decomposition can be viewed as associated to an appropriate Morse function and this leads to a contradiction. The obstructions to smoothability constructed in this way are obstructions to the lifting of the Spivak normal bundle to $`BO`$. This is an obstruction theory problem but one which can be very difficult to solve directly in the presence of many cells. Thus, this approach is quite efficient to construct examples.
B. The second application concerns the detection of obstructions to the embedding of $`CW`$-complexes in euclidean spaces in low codimension. The argument in this case goes roughly as follows. If the $`CW`$-complex $`X`$ embedds in $`S^n`$, then we may consider a neighbourhood $`U(X)`$ of $`X`$ which is a smooth manifold with boundary. We consider a smooth Morse function $`f:U(X)`$ which is constant, maximal and regular on the boundary of $`U(X)`$. If $`P`$ and $`Q`$ are two consecutive critical points for this funtion we obtain that $`\mathrm{\Sigma }^kH(P,Q)=\mathrm{\Sigma }^k^{}H(Q,P)`$ for certain values of $`k`$ and $`k^{}`$ which can be estimated explicitly - the main reason for this equality is that the Morse function in question is defined on the sphere so all the questions involving the stable normal bundle become irrelevant. If $`X`$ admits some reasonably explicit cell-decomposition it is possible to express $`H(P,Q)`$ as the Hopf invariant $`H`$ of some successive attachment of two cells $`e,f`$ and $`H(Q,P)`$ as $`\mathrm{\Sigma }^{k^{\prime \prime }}H^{}`$ where $`H^{}`$ is another similar Hopf invariant. The obstructions to embedding appear because the low codimension condition translates to the fact that $`k^{}+k^{\prime \prime }>k`$. This can be viewed as an obstruction because it means that after $`k`$ suspensions the homotopy class of $`H`$ has to de-suspend more than $`k`$-times.
###### Remark 2.3.
The applications at A and B are purely of homotopical type. It is natural to expect that the Morse theoretical arguments that were used to establish these statements can be replaced by purely homotopical ones but this has not been done till now.
#### 2.1.4. The Serre spectral sequence
Theorem 2.2 provides considerable information on connecting manifolds for pairs of consecutive critical points. However, it does not shed any light on the case of non-consecutive ones. Clearly, if the critical points are not consecutive the respective connecting manifold is not closed and thus no bordism or cobordism class can be directly associated to it. However, after compactification, the boundary of this connecting manifold has a special structure reflected by equation (1). As we shall see following , this structure is sufficient to construct an algebraic invariant which provides an efficient “measure” of all connecting manifolds.
This construction is based on the fact that the maps
$$j_Q^P:M_Q^P\mathrm{\Omega }M$$
are compatible with compactification and with the formula (1) in the following sense. Recall that here $`\mathrm{\Omega }M`$ are the based Moore loops on $`M`$ (these are loops parametrized by intervals $`[0,a]`$), the critical points of $`f`$ have been identified to a single point and, moreover, in the definition of $`j_Q^P`$ we use the parametrization of the flow lines by the values of $`f`$. Recall that we have a product given by the concatenation of loops
$$\mu :\mathrm{\Omega }M\times \mathrm{\Omega }M\mathrm{\Omega }M.$$
With these notations it is easy to see that we have the following formula:
(4)
$$j_Q^P(u,v)=\mu (j_R^P(u),j_Q^R(v))$$
where $`(u,v)\overline{M}_R^P\times \overline{M}_Q^R\overline{M}_Q^P`$.
We proceed with our consruction. Let $`C_{}(X)`$ be the (reduced) cubical complex of $`X`$ with coefficients in $`/2`$. Notice that there is a natural map
$$C_k(X)C_k^{}(Y)C_{k+k^{}}(X\times Y).$$
A family of cubical chain $`s_y^xC_{|x||y|1|}(\overline{M}_y^x)`$, $`x,yCrit(f)`$ is called a representing chain system for the moduli spaces $`M_y^x`$ if for each pair of critical points $`x,z`$ we have:
* $$ds_z^x=\underset{y}{}s_y^xs_z^y$$
* $`s_z^x`$ represents the fundamental class in $`H_{|x||z|1}(M_z^x,M_z^x)`$.
It is easy to show by induction on the index difference $`|x||z|`$ that such representing chain systems exist. We now fix such a representing chain system $`\{s_y^x\}`$ and we define $`a_y^xC_{|x||y|1}(\mathrm{\Omega }M)`$ by
$$a_y^x=(j_y^x)_{}(s_y^x).$$
Notice that this definition extends the definition of these coefficients in the usual Morse case when $`|x||y|1=0`$. We have a product map
$$:C_k(\mathrm{\Omega }M)C_k^{}(\mathrm{\Omega }M)C_{k+k^{}}(\mathrm{\Omega }M\times \mathrm{\Omega }M)\stackrel{C_{}(\mu )}{}C_{k+k^{}}(\mathrm{\Omega }M)$$
which makes $`C_{}(\mathrm{\Omega }M)`$ into a differential ring. The discussion above shows that inside this ring we have the formula
$$da_z^x=\underset{y}{}a_z^xa_y^z.$$
An elegant way to rephrase this formula is to group these coefficients in a matrix $`A=(a_y^x)`$ and then we have
(5)
$$dA=A^2.$$
We now define a new chain complex $`𝒞(f)`$ associated to $`f`$ by
(6)
$$𝒞(f)=(C_{}(\mathrm{\Omega }M)/2<Crit(f)>,d),dx=\underset{y}{}a_y^xy.$$
We shall call this complex the extended Morse complex of $`f`$. Here, $`C_{}(\mathrm{\Omega }M)/2<Crit(f)>`$ is viewed as a graded $`C_{}(\mathrm{\Omega }M)`$-module and $`d`$ respects this structure in the sense that it verifies $`d(ax)=(da)x+a(dx)`$ (the grading on $`Crit(f)`$ is given, as before, by the Morse index). Choosing orientations on all the stable manifolds of all the critical points induces a co-orientation on all the unstable manifolds, and hence an orientations on the intersections $`W^u(P)W^s(Q)`$ and finally on all the moduli spaces $`M_Q^P`$ : we may then use $``$-coefficients for this complex as well as, of course, for the classical Morse complex. In this case appropriate signs appear in the formulae above. Clearly, $`d^2=0`$ due to (5).
By definition, the coefficients $`a_y^x`$ represent the moduli spaces $`M_y^x`$. However, these coefficients are not invariant with respect to the choices made in their construction. Therefore, it is remarkable that there is a natural construction which extracts from this complex a useful algebraic invariant which is not just the homology of the complex - as it happens, this homology is not too interesting as it coincides with that of a point.
Consider the obvious differential filtration which is defined on this complex by
$$F^k𝒞(f)=C_{}(\mathrm{\Omega }M)/2<xCrit(f):ind_f(x)k>.$$
Denote the associated spectral sequence by
$$E(f)=(E_{p,q}^r(f),d^r).$$
###### Theorem 2.4.
When $`M`$ is simply connected and if $`r2`$ the spectral sequence $`E(f)`$ coincides with the Serre spectral sequence of the path-loop fibration
$$𝒫:\mathrm{\Omega }MPMM.$$
###### Remark 2.5.
a. A similar result can be established even in the absence of the simple-connectivity condition which has been assumed here to avoid some technical complications.
b. The Serre spectral sequence of the path-loop fibration of a space $`X`$ contains considerable information on the homotopy type of the space. In particular, there are spaces with the same cohomology and cup-product but which may be distinguished by their respective Serre spectral sequences.
To outline the proof of the theorem we start by recalling the construction of the Serre spectral sequence in the form which will be of use here. We shall assume here that the Morse function $`f`$ is self-indexed (in the sense that for each critical point $`x`$ we have $`ind_f(x)=kf(x)=k`$) and that it has a single minimum denoted by $`m`$. Let $`M_k=f^1((\mathrm{},k+ϵ])`$. We have
$$M_k=M_{k1}\underset{\varphi _y}{}D_y^k$$
where the union is taken over all the critical points $`yCrit_k(f)`$ and $`\varphi _y:S^u(y)M_{k1}`$ are the respective attaching maps. Denote by $`E_k`$ the total space of the fibration induced by pull-back over the inclusion $`M_kM`$ from the fibration $`𝒫`$. Consider the filtration of $`C_{}(PM)`$ given by $`F^kP=Im(C_{}(E_k)C_{}(PM))`$. The spectral sequence induced by this filtration is invariant after the second page and is precisely the Serre spectral sequence (this spectral sequence may be constructed as above but by using an arbitrary skeletal filtration $`\{X_k\}`$ of a space $`X`$ which has the same homotopy type as that of $`M`$; in our case the particular filtration given by the sets $`M_k`$ is a natural choice).
For further use, we also notice that there is an obvious action of $`\mathrm{\Omega }M`$ on $`PM`$ and this action induces one on each $`E_k`$. Therefore, we may view $`C_{}(E_k)`$ as a $`C_{}(\mathrm{\Omega }M)`$-module.
The first step in proving the theorem is to consider a certain compactification of the unstable manifolds of the critical points of $`f`$. Recall that $`f`$ is self-indexed and that $`m`$ is the unique minimum critical point of $`f`$. Fix $`xCrit(f)`$ and define the following equivalence relation on the set $`\overline{M}_m^x\times [0,f(x)]`$:
$$(a,t)(a^{},t^{})\mathrm{iff}t=t^{}\mathrm{and}a(\tau )=a^{}(\tau )\tau t.$$
Here the elements of $`\overline{M}_m^x`$ are viewed as paths in $`M`$ parametrized by the value of $`f`$ (so that $`f(a(\tau ))=\tau `$).
Denote by $`\widehat{W}(x)`$ the resulting quotient topological space. Notice that, if $`y\overline{W^u(x)}`$, then there exists some $`a\overline{M}_m^x`$ so that $`y`$ is on the (possibly broken) flow line represented by $`a`$. Or, in other words, so that $`a(f(y))=y`$. This path $`a`$ might not be unique. Indeed, inside $`\widehat{W}(x)`$ there is precisely one equivalence class $`[a,f(y)]`$ (with $`a(f(y))=y`$) for each (possibly) broken flow line joining $`x`$ to $`y`$. Clearly, if $`yW^u(x)`$, then there is just one such flow line and so the natural surjection
$$\pi :\widehat{W}(x)\overline{W^u(x)},\pi ([a,t])=a(t)$$
is a homeomorphism when restricted to $`\pi ^1(W^u(x))`$. Thus we may view $`\widehat{W}(x)`$ as a special compactification of $`W^u(x)`$ or as a desingularization of $`\overline{W^u(x)}`$. It is not difficult to believe (but harder to show and will not be proven here see ) that $`\widehat{W}(x)`$ is a topological manifold with boundary and moreover
(7)
$$\widehat{W}(x)=\underset{y}{}M_y^x\times \widehat{W}(y).$$
We continue the proof of the Theorem 2.4 with the remark that there are obvious maps
$$h_x:\widehat{W}(x)PM$$
which associate to $`[a,t]`$ the path in $`M`$ which follows $`a`$ from $`x`$ to $`a(t)`$. These maps and the maps $`j_y^x`$ are compatible with formula (7) in the sense that
(8)
$$h_x(a^{},[a^{\prime \prime },t])=j_y^x(a^{})h_y([a^{\prime \prime },t])$$
where $`(a^{},[a^{\prime \prime },t])M_y^x\times \widehat{W}(y)`$ and $``$ represents the action of $`\mathrm{\Omega }M`$ on $`PM`$.
Now, we may obviously rewrite (7) as:
$$\widehat{W}(x)=\underset{y}{}\overline{M}_y^x\times \widehat{W}(y).$$
Given the representing chain system $`\{s_y^x\}`$ it is easy to construct an associated representing chain system for $`\widehat{W}(x)`$. This is a system of chains $`v(x)C_{|x|}(\widehat{W}(x))`$ so that $`v(x)`$ represents the fundamental class of $`C_{|x|}(\widehat{W}(x),\widehat{W}(x))`$ and we have the formula
$$dv(x)=\underset{y}{}s_y^xv(y).$$
Finally, we define a $`C_{}(\mathrm{\Omega }M)`$-module chain map
$$\alpha :𝒞(f)C_{}(PM)$$
by
$$\alpha (x)=(h_x)_{}(v(x)).$$
The formulas above show that we have
$$d[(h_x)_{}(v(x))]=\underset{y}{}a_y^x(h_y)_{}(v(y))$$
and so $`\alpha `$ is a chain map. It is clear that the map $`\alpha `$ is filtration preserving and it is not difficult to see that it induces an isomorphism at the $`E^2`$ level of the induced spectral sequences and this concludes the proof of Theorem 2.4.
###### Remark 2.6.
a. Another important but immediate property of $`\widehat{W}(x)`$ is that it is a contractible space. Indeed, all the points in $`\overline{M}_m^x\times \{f(x)\}`$ are in the same equivalence class. Moreover, each point $`[a,t]\widehat{W}(x)`$ has the property that it is related by the path $`[a,\tau ],\tau [t,f(x)]`$ to $`=[a,f(x)]`$. The contraction of $`\widehat{W}(x)`$ to $``$ is obtained by deforming $`\widehat{W}(x)`$ along these paths. Given that $`\widehat{W}(x)`$ is a contractible topological manifold with boundary, it is natural to suspect that $`\widehat{W}(x)`$ is homeomorphic to a disk. This is indeed the case as is shown in and is an interesting fact in itself because it implies that the union of the unstable manifolds of a self-indexed Morse-Smale function gives a $`CW`$-decomposition of $`M`$. The attaching map of the cell $`\widehat{W}(x)`$ is simply the restriction of $`\pi `$ to $`\widehat{W}(x)`$.
b. The Serre spectral sequence result above and the bordism result in Theorem 2.2 are obviously related via the central role of the maps $`j_Q^P`$. There is also a more explicit relation. Indeed, (a stable version of) the Hopf invariants appearing in Theorem 2.2 can be interpreted as differentials in the Atyiah-Hirzebruch-Serre spectral sequence of the path-loop fibration with coefficients in the stable homotopy of $`\mathrm{\Omega }M`$. Moreover, the relation (1) can be understood as also keeping track of the framings. This leads to a type of extended Morse complex in which the coefficients of the differential are stable Hopf invariants . All of this strongly suggests that the construction of the complex $`𝒞(f)`$ can be enriched so as to include the framings of the connecting manifolds and, by the same method as above, that the whole Atyiah-Hirzebruch-Serre spectral sequence should be recovered from this construction.
c. Another interesting question, open even for consecutive critical points $`P`$, $`Q`$, is whether there are some additional constraints on the topology of the connecting manifolds $`M_Q^P`$ besides those imposed by Theorem 2.2.
d. Yet another open question is how this machinery can be adapted to the Morse-Bott situation and how it can be extended to general Morse-Smale flows (not only gradient-like ones).
e. It is natural to wonder what is the richest level of information that one can extract out of the moduli spaces of Morse flow lines. At a naive level, the union of all the points situated on the flow lines of $`f`$ is precisely the whole underlying manifold $`M`$ so we expect that there should exist some assembly process producing the manifold $`M`$ out of these moduli spaces. Such a machine has been constructed by Cohen, Jones and Segal . They show that one can form a category out of the moduli spaces of connecting trajectories and that the classifying space of this category is of the homeomorphism type of the underlying manifold. In their construction an essential point is that the glueing of flow lines is associative. This approach is quite different from the techniques above and does not imply the results concerning the extended Morse complex or the Hopf invariants that we have presented. The two points of view are, essentially, complementary.
To end this section it is useful to make explicit a relation between Theorems 2.4 and 2.2 (we assume as above that $`M`$ is simply-connected).
###### Proposition 2.7.
Assume that there are $`q,p`$ so that in the Serre spectral sequence of the path loop fibration of $`M`$ we have $`E_{k,s}^2=0`$ for $`q<k<p`$ and there is an element $`aE_{p,0}^2`$ so that $`d^{pq}a0`$, then any Morse-Smale function on $`M`$ has a pair of consecutive critical points $`P`$, $`Q`$ of indexes at least $`q`$ and at most $`p`$ so that the homology class $`[M_Q^P]H_{|P||Q|1}(\mathrm{\Omega }M)0`$.
Clearly, Theorem 2.4 directly implies that, even without any restriction on $`E^2`$, if we have $`d^ra0`$ with $`aE_{p,0}^r`$, then for any Morse-Smale function $`f`$ there are critical points $`P`$ and $`Q`$ with $`|P|=p`$ and $`|Q|=pr`$ so that $`M_Q^P\mathrm{}`$. Indeed, if this would not be the case, then all the coefficients $`a_y^x`$ in the extended Morse complex of $`f`$ are null whenever $`|x|=p`$, $`|y|=pr`$. By the construction of the associated spectral sequence, this leads to a contradiction. However, the pair $`P`$, $`Q`$ resulting from this argument might have a connecting manifold which is not closed so that its homology class is not even defined and, thus, Proposition 2.7 provides a stronger conclusion. The proof of the Proposition is as follows. Recall that $`E_{s,r}^2H_s(M)H_r(\mathrm{\Omega }M)`$ and so $`H_{}(M)=0`$ for $`q0p`$. If there are some points $`P,QCrit(f)`$ with $`q|Q|,|P|p`$ so that the differential of $`P`$ in the classical Morse complex contains $`Q`$ with a non-trivial coefficient then this pair $`P`$, $`Q`$ may be taken as the one we are looking for. If all such differentials in the classical Morse complex are trivial it follows that the critical points of index $`p`$ and $`q`$ are consecutive. In this case, the geometric arguments used in the proofs of either Theorem 2.2 or 2.4 imply that if for all pairs $`P`$, $`Q`$ with $`|P|=p`$, $`|Q|=q`$ we would have $`[M_Q^P]=0`$, then the differential $`d^{pq}`$ would vanish on $`E_{p,0}^{pq}`$.
###### Remark 2.8.
Notice that the pair of critical points $`P`$ and $`Q`$ constructed in the proposition verify the property that $`|P|`$ and $`|Q|`$ are consecutive inside the set $`\{ind_f(x):xCrit(f)\}`$.
### 2.2. Operations
We discuss here a different and, probably, more familiar approach to understanding connecting manifolds as well as other related Morse theoretic moduli spaces. This point of view has been used extensively by many authors - Fukaya , Betz and Cohen being just a few of them. For this reason we shall review this technique very briefly.
Given two consecutive critical points $`x`$, $`y`$ notice that the set $`T_y^x=W^u(x)W^s(y)`$ is homeomorphic to the un-reduced suspension of $`M_y^x`$. Therefore, we may see this as an obvious inclusion
$$i_y^x:\mathrm{\Sigma }M_y^xM$$
and we may consider the homology class $`[T_y^x]=(i_y^x)_{}(s[M_y^x])`$ where $`s`$ is suspension and $`[M_y^x]`$ is the fundamental class. There exists an obvious evaluation map
$$e:\mathrm{\Sigma }\mathrm{\Omega }MM$$
which is induced by $`\mathrm{\Omega }M\times [0,1]M,(\beta ,t)\beta (t)`$ (the loops here are parametrized by the interval $`[0,1]`$ but this is a minor technical difficulty). It is easy to see, by the definition of this evaluation map, that $`[T_y^x]=e_{}((j_y^x)_{}([M_y^x]))`$. In general the map $`e_{}`$ is not injective in homology. Clearly, the full bordism class $`\{M_y^x\}`$ carries much more information than the homology class $`[T_y^x]`$. Still, there is a direct way to determine $`[T_y^x]`$ without passing through a calculation of $`\{M_y^x\}`$ and we will now describe it.
Consider a second Morse-Smale function $`g:M`$ so that its associated unstable and stable manifolds $`W_g^u()`$, $`W_g^s()`$ intersect transversally the stable and unstable manifolds of $`f`$ and, except if they are of top dimension, they avoid the critical points of $`f`$.
Fix $`x,yCrit(f)`$ and $`sCrit(g)`$ so that $`|x||y|ind_g(s)=0`$. We may define $`k(x,y;s)=\mathrm{\#}(T_y^xW_g^s(s)`$ (where the counting takes into account the relevant orientations if we work over $``$). We now put
$$\overline{k}_y^x=\underset{s}{}k(x,y;s)sC(g).$$
The essentially obvious claim is that:
###### Proposition 2.9.
The chain $`\overline{k}_y^x`$ is a cycle whose homology class is $`[T_y^x]`$.
Indeed we have $`_sk(x,y;s)h_z^s=0`$ where $`ind_g(z)=ind_g(s)1`$ and $`h_z^s`$ are the coefficients in the classical Morse complex of $`g`$. This equality is valid because we may consider the $`1`$-dimensional space $`T_y^xW_g^s(z)`$. This is an open $`1`$-dimensional manifold whose compactification is a $`1`$-manifold whose boundary points are counted precisely by the formula $`_sk(x,y;s)h_z^s=0`$. By basic intersection theory it is immediate to see that the homology class represented by this cycle is $`[T_y^x]`$.
While this construction does not shed a lot of light on the properties of $`M_y^x`$ its role is important once we use it to recover the various homological operations of $`M`$. To see how this is done from our perspective notice that the intersection
$$T_y^xW_g^s(s)=W_f^u(x)W_f^s(y)W_g^s(s)$$
can be viewed as a particular case of the following situation: assume that $`f_1`$, $`f_2`$, $`f_3`$ are three Morse-Smale functions in general position and define
$$T_z^{x,y}=(W_{f_1}^u(x)W_{f_2}^u(y)W_{f_3}^s(z)).$$
If we assume that $`|x|+|y||z|n=0`$ we may again count the points in $`T_z^{x,y}`$ with appropriate signs and we may define coefficients $`t_z^{x,y}=\mathrm{\#}T_z^{x,y}`$.
This leads to an operation
$$C(f_1)C(f_2)C(f_3)$$
given as a linear extension of
$$xy\underset{z}{}t_z^{x,y}z.$$
It is easy to see that this operation descends in homology and that it is in fact the dual of the $``$-product. Moreover, the space $`T_z^{x,y}`$ may be viewed as obtained by considering a graph formed by three oriented edges meeting into a point with the first two entering the point and the other one exiting it and considering all the configurations obtained by mapping this graph into $`M`$ so that to the first edge we associate a flow line of $`f_1`$ which exits $`x`$, to the second edge a flow line of $`f_2`$ which exits $`y`$ and to the third a flow line of $`f_3`$ which enters $`z`$. Clearly, this idea may be pushed further by considering other, more complicated graphs and understanding what are the operations they correspond to as was done by Betz and Cohen .
## 3. Applications to Symplectic topology
We start with some applications that are rather “soft” even if difficult to prove and we shall continue in the main part of the section,§3.2, with some others that go deeper.
### 3.1. Bounded orbits
We fix a symplectic manifold $`(M,\omega )`$ which is not compact. Assume that $`H:M`$ is a smooth hamiltonian whose associated hamiltonian vector field is denoted by $`X_H`$. One of the main questions in hamiltonian dynamics is whether a given regular hypersurface $`A=H^1(a)`$ of $`H`$ has any closed caracteristics, or equivalently, whether the hamiltonian flow of $`H`$ has any periodic orbits in $`A`$. As $`M`$ is not compact, from the point of view of dynamical systems, the first natural question is whether $`X_H`$ has any bounded orbits in $`A`$. Moreover, there is a remarkable result of Pugh and Robinson , the $`C^1`$-closing lemma, which shows that, for a generic choice of $`H`$, the presence of bounded orbits insures the existence of some periodic orbits. Therefore, we shall focus in this subsection on the detection of bounded orbits. It should be noted however that the detection of periodic orbits in this way is not very effective because the periods of the orbits found can not be estimated. Moreover, there is no reasonable test to decide whether a given hamiltonian belongs to the generic family to which the $`C^1`$-closing lemma applies. Finally, it will be clear from the methods of proof described below that these results are also soft in the sense that they are not truly specific to Hamiltonian flows but rather they apply to many other flows.
An example of a bounded orbit result is the following statement .
###### Theorem 3.1.
Assume that $`H`$ is a Morse-Smale function with respect to a riemannian metric $`g`$ on $`M`$ so that $`M`$ is metrically complete and there exists an $`ϵ`$ and a compact set $`KM`$ so that $`_gH(x)ϵ`$ for $`xK`$. Suppose that $`P`$ and $`Q`$ are two critical points of $`H`$ so that $`|P|`$ and $`|Q|`$ are successive in the set $`\{ind_H(x):xCrit(H)\}`$. If the stabilization $`[H(P,Q)]\pi _{|P||Q|1}^S(\mathrm{\Omega }M)`$ of the Hopf invariant $`H(P,Q)`$ is not trivial then there are regular values $`v(H(Q),H(P))`$ so that $`H^1(v)`$ contains bounded orbits of $`X_H`$.
Before describing the proof of this result let’s notice that the theorem is not difficult to apply. Indeed, one simple way to verify that there are pairs $`P`$, $`Q`$ as required is to use Proposition 2.7 together with Remark 2.8 with a minor adaptation required in a non-compact setting. This adaptation consists of replacing the Serre spectral sequence of the path loop fibration with the Serre spectral sequence of a relative fibration $`\mathrm{\Omega }M(E_1,E_0)(N_1,N_0)`$ where $`N_1`$ is an isolating neighbourhood for the gradient flow of $`H`$ which contains $`K`$ and $`N_0`$ is a (regular) exit set for this neighbourhood (to see the precise definition of these Conley index theoretic notions see ). The fibration is induced by pull-back from the path-loop fibration $`\mathrm{\Omega }MPMM`$ over the inclusion $`(N_1,N_0)(M,M)`$. In short, because the gradient of $`H`$ is away from $`0`$ outside of a compact set, pairs $`(N_1,N_0)`$ as above are easy to produce and if the pair $`(N_1,N_0)`$ has some interesting topology it is easy to deduce the existence of non-constant bounded orbits. Here is a concrete example.
###### Corollary 3.2.
Assume that $`M`$ is the contangent bundle of some closed, simply-connected manifold $`N`$ of dimension $`k2`$ and $`\omega `$ is an arbitrary symplectic form. Assume that $`H:M`$ is Morse and that outside of some compact set containing the $`0`$ section, $`H`$ restricts to each fibre of the bundle to a non-degenerate quadratic form. Then, $`X_H`$ has bounded, non constant orbits.
Of course, this result is only interesting when there are no compact level hypersurfaces of $`H`$. This does happen if the quadratic form in question has an index which is neither $`0`$ nor $`k`$. The proof of this result comes down to the fact that as $`N`$ is closed and not a point there exists a lowest dimensional homology class $`uH_t(N)`$ which is transgressive in the Serre spectral sequence ( this means $`d^tu0`$). Using the structure of function quadratic at infinity of $`H`$ it is easy to construct a pair $`(N_1,N_0)`$ where $`N_1`$ is is homotopic to a disk bundle of base $`N`$ and $`N_0`$ is the associated sphere bundle. The spectral sequence associated to this pair can be related by the Thom isomorphism to the Serre spectral sequence of the path-loop fibration over $`N`$ and the element $`\overline{u}H_{}(N_1,N_0)`$ which corresponds to $`u`$ by the Thom isomorphism will have a non-vanishing differential. This means that Proposition 2.7 may be used to show the non-triviality of a homology class $`[M_Q^P]`$ for $`P`$ and $`Q`$ as in Theorem 3.1. By Theorem 2.2, $`[M_Q^P]`$ is the same up to sign as the homology class of the Hopf invariant $`H(P,Q)`$ so Theorem 3.1 is applicable to detect bounded orbits.
We now describe the proof of the theorem. The basic idea of the proof is simpler to present in the particular case when $`H^1(H(Q),H(P))`$ does not contain any critical value. In this case let $`A=H^1(a)`$ where $`a(H(Q),H(P))`$. We intend to show that $`A`$ contains some bounded orbits of $`X_H`$. To do this notice that the two sets $`S_1=W^u(P)A`$, $`S_2=W^s(Q)A`$ are both diffeomorphic to spheres. We now assume that no bounded orbits exist and we consider a compact neighbourhood $`U`$ of $`S_1S_2`$. Assume that we let $`S_2`$ move along the flow $`X_H`$. As this flow has no bounded orbits, each point of $`S_2`$ will leave $`U`$ at some moment. Suppose that we are able to perturb the flow induced by $`X_H`$ to a new deformation $`\eta :M\times M`$ so that for some finite time $`T`$ all the points in $`S_2`$ are taken simultaneously outside $`U`$ (in other words $`\eta _T(S_2)U=\mathrm{}`$) and so that $`\eta `$ leaves $`Q`$ fixed. It is easy to see that this implies that $`S_1S_2`$ is bordant to the empty set which, by Theorem 2.2, is impossible because $`H(P,Q)0`$. This perturbation $`\eta `$ is in fact not hard to construct by using some elements of Conley’s index theory and the fact that the maximal invariant set of $`X_T`$ inside $`U`$ is the empty set (the main step here is to possibly modify also $`U`$ so that it admits a regular exit region $`U_0U`$ and we then construct $`\eta `$ so that it follows the flow lines of $`X_T`$ but stops when reaching $`U_0`$, this eliminates the problem of “bouncing” points which first exit $`U`$ but later re-enter it).
The case when there are critical points in $`H^1(H(Q),H(P))`$ follows the same idea but is considerably more difficult. The main difference comes from the fact that the sets $`S_1`$ and $`S_2`$ might not be closed manifolds. Even their closures $`\overline{S_1}`$ and $`\overline{S_2}`$ are not closed manifolds in general but might be singular sets. To be able to proceed in this case we first replace $`P`$ and $`Q`$ with a pair of critical points of the same index so that for any critical point $`Q^{}H^1(H(Q),H(P))`$ with $`ind_H(Q^{})=ind_H(Q)`$ we have $`[H(P,Q^{})]=0`$. We then take $`a`$ very close to $`H(Q)`$ so that $`S_2`$ at least is diffeomorphic to a sphere. We then first study the stratification of $`\overline{S_1}`$: there is a top stratum of dimension $`|P|1`$ which is $`S_1`$ and a singular stratum $`S^{}`$ of dimension $`|Q|1`$ which is the union of the sets $`W^u(Q^{})A`$ for all $`Q^{}`$ so that $`M_Q^{}^P\mathrm{}`$ and $`|Q^{}|=|Q|`$. Notice that the way to construct the null-bordism of $`S_1S_2`$ is to consider in $`A\times [0,T]`$ the submanifold $`W=(\eta _t(\overline{S_2}),t)`$ and intersect it with $`W^{}=\overline{S_1}\times [0,T]`$ \- we assume here $`\eta _T(S_2)S_1=\mathrm{}`$. Clearly, we need this intersection to be transverse and this can be easily achieved by a perturbation of $`\eta `$. The main technical difficulty is that $`L`$ might intersect the singular part, $`S^{}\times [0,T]`$. Indeed, $`dim(W)=nq`$, $`dim(S^{})=q1`$, $`dim(A)=n1`$ and so generically the intersection $`I`$ between $`W`$ and $`S^{}\times [0,T]`$ is $`0`$-dimensional and not necessarily void. By studying the geometry around each of the points of $`I`$ it can be seen that $`S_1S_2`$ is bordant to the union of the $`M_Q^{}^P`$’s where $`Q^{}H^1(H(Q),H(P))`$ (roughly, this follows by eliminating from the singular bordism $`WW^{}`$ a small closed, cone-like neighbourhood around each singular point and showing that the boundary of this cone-like neighbourhood is homeomorphic to a $`M_Q^{}^P`$). We now use the fact that all the stable bordism classes of the $`M_Q^{}^P`$’s vanish (because $`[H(P,Q^{})]=0`$) and this leads us to a contradiction. Notice also that, at this point, we need to use stable Hopf invariants (or bordism classes) $`\pi ^S(\mathrm{\Omega }M)`$ because, by contrast to the stable case, the unstable Thom-Pontryagin map associated to a disjoint union is not necessarily equal to the sum of the Thom-Pontryagin maps of the terms in the union and hence, unstably, even if we know $`H(P,Q^{})=0,Q^{}`$ we still can not deduce $`H(P,Q)=0`$.
###### Remark 3.3.
It would be interesting to see whether, under some additional assumptions, a condition of the type $`[H(P,Q)]0`$ implies the existence of periodic orbits and not only bounded ones.
### 3.2. Detection of pseudoholomorphic strips and Hofer’s norm
In this subsection we shall again use the Morse theoretic techniques described in §2 and, in particular, Theorem 2.4 to study some symplectic phenomena by showing that Floer’s complex can be enriched in a way similar to the passage from the classical Morse complex to the extended one.
#### 3.2.1. Elements of Floer’s theory.
We start by recalling very briefly some elements from Floer’s construction (for a more complete exposition see, for example, ).
We shall assume from now on that $`(M,\omega )`$ is a symplectic manifold - possibly non-compact but in that case convex at infinity - of dimension $`m=2n`$. We also assume that $`L,L^{}`$ are closed (no boundary, compact) Lagrangian submanifolds of $`M`$ which intersect transversally.
To start the description of our applications it is simplest to assume for now that $`L,L^{}`$ are simply-connected and that $`\omega |_{\pi _2(M)}=c_1|_{\pi _2(M)}=0`$. Cotangent bundles of simply-connected manifolds offer immediate examples of manifolds verifying these conditions.
We fix a path $`\eta 𝒫(L,L^{})=\{\gamma C^{\mathrm{}}([0,1],M):\gamma (0)L`$, $`\gamma (1)L^{}\}`$ and let $`𝒫_\eta (L,L^{})`$ be the path-component of $`𝒫(L,L^{})`$ containing $`\eta `$. This path will be trivial homotopically in most cases, in particular, if $`L`$ is hamiltonian isotopic to $`L^{}`$. We also fix an almost complex structure $`J`$ on $`M`$ compatible with $`\omega `$ in the sense that the bilinear form $`X,Y\omega (X,JY)=\alpha (X,Y)`$ is a Riemannian metric. The set of all the almost complex structures on $`M`$ compatible with $`\omega `$ will be denoted by $`𝒥_\omega `$. Moreover, we also consider a smooth $`1`$-periodic Hamiltonian $`H:[0,1]\times M`$ which is constant outside a compact set and its associated $`1`$-periodic family of hamiltonian vector fields $`X_H`$ determined by the equation
$$\omega (X_H^t,Y)=dH_t(Y),Y.$$
we denote by $`\varphi _t^H`$ the associated Hamiltonian isotopy. We also assume that $`\varphi _1^H(L)`$ intersects transversally $`L^{}`$.
In our setting, the action functional below is well-defined:
(9)
$$𝒜_{L,L^{},H}:𝒫_\eta (L,L^{}),x\overline{x}^{}\omega +_0^1H(t,x(t))𝑑t$$
where $`\overline{x}(s,t):[0,1]\times [0,1]M`$ is such that $`\overline{x}(0,t)=\eta (t)`$, $`\overline{x}(1,t)=x(t)`$, $`t[0,1]`$, $`x([0,1],0)L`$, $`x([0,1],1)L^{}`$. The critical points of $`𝒜`$ are the orbits of $`X_H`$ that start on $`L`$, end on $`L^{}`$ and which belong to $`𝒫_\eta (L,L^{})`$. These orbits are in bijection with a subset of $`\varphi _1^H(L)L^{}`$ so they are finite in number. If $`H`$ is constant these orbits coincide with the intersection points of $`L`$ and $`L^{}`$ which are in the class of $`\eta `$. We denote the set of these orbits by $`I(L,L^{};\eta ,H)`$ or shorted $`I(L,L^{})`$ if $`\eta `$ and $`H`$ are not in doubt.
We now consider the solutions $`u`$ of Floer’s equation:
(10)
$$\frac{u}{s}+J(u)\frac{u}{t}+H(t,u)=0$$
with
$$u(s,t):\times [0,1]M,u(,0)L,u(,1)L^{}.$$
When $`H`$ is constant, these solutions are called *pseudo-holomorphic strips*.
For any strip $`u𝒮(L,L^{})=\{uC^{\mathrm{}}(\times [0,1],M):u(,0)L,u(,1)L^{}\}`$ consider the energy
(11)
$$E_{L,L^{},H}(u)=\frac{1}{2}_{\times [0,1]}\frac{u}{s}^2+\frac{u}{t}X_H^t(u)^2dsdt.$$
For a generic choice of $`J`$, the solutions $`u`$ of (10) which are of finite energy, $`E_{L,L^{},H}(u)<\mathrm{}`$, behave like negative gradient flow lines of $`𝒜`$. In particular, $`𝒜`$ decreases along such solutions. We consider the moduli space
(12)
$$^{}=\{u𝒮(L,L^{}):u\mathrm{verifies}(\text{10}),E_{L,L^{},H}(u)<\mathrm{}\}.$$
The translation $`u(s,t)u(s+k,t)`$ obviously induces an $``$ action on $`^{}`$ and we let $``$ be the quotient space. For each $`u^{}`$ there exist $`x,yI(L,L^{};\eta ,H)`$ such that the (uniform) limits verify
(13)
$$\underset{s\mathrm{}}{lim}u(s,t)=x(t),\underset{s+\mathrm{}}{lim}u(s,t)=y(t).$$
We let $`^{}(x,y)=\{u^{}:u\mathrm{verifies}(\text{13})\}`$ and $`(x,y)=^{}(x,y)/`$ so that $`=_{x,y}(x,y)`$. If needed, to indicate to which pair of Lagrangians, to what Hamiltonian and to what almost complex structure are associated these moduli spaces we shall add $`L`$ and $`L^{}`$, $`H`$, $`J`$ as subscripts (for example, we may write $`_{L,L^{},H,J}(x,y)`$).
For $`x,yI(L,L^{};\eta ,H)`$ we let
$$\begin{array}{cc}𝒮(x,y)=\{uC^{\mathrm{}}([0,1]\times [0,1],M):& u([0,1],0)L,u([0,1],1)L^{},\\ u(0,t)=x(t),u(1,t)=y(t)\}.& \end{array}$$
To each $`uS(x,y)`$ we may associate its Maslov index $`\mu (u)`$ and it can be seen that, in our setting, this number only depends on the points $`x,y`$. Thus, we let $`\mu (x,y)=\mu (u)`$. Moreover, we have the formula
(14)
$$\mu (x,z)=\mu (x,y)+\mu (y,z).$$
According to these relations, the choice of an arbitrary intersection point $`x_0`$ and the normalization $`|x_0|=0`$, defines a grading $`|.|`$ such that :
$$\mu (x,y)=|x||y|.$$
There is a notion of regularity for the pairs of $`(H,J)`$ so that, when regularity is assumed, the spaces $`^{}(x,y)`$ are smooth manifolds (generally non-compact) of dimension $`\mu (x,y)`$ and in this case $`(x,y)`$ is also a smooth manifold of dimension $`\mu (x,y)1`$. Regular pairs $`(H,J)`$ are generic and, in fact, they are so even if $`L`$ and $`L^{}`$ are not transversal (but in that case $`H`$ can not be assumed to be constant), for example, when $`L=L^{}`$.
Floer’s construction is natural in the following sense. Let $`L^{\prime \prime }=(\varphi _1^H)^1(L^{})`$. Consider the map $`b_H:𝒫(L,L^{\prime \prime })𝒫(L,L^{})`$ defined by $`(b_H(x))(t)=\varphi _t^H(x(t))`$. Let $`\eta ^{}𝒫(L,L^{\prime \prime })`$ be such that $`\eta =b_H(\eta ^{})`$. Clearly, $`b_H`$ restricts to a map between $`𝒫_\eta ^{}(L,L^{\prime \prime })`$ and $`𝒫_\eta (L,L^{})`$ and it restricts to a bijection $`I(L,L^{\prime \prime };\eta ^{},0)I(L,L^{};\eta ,H)`$.
It is easy to also check
$$𝒜_{L,L^{},H}(b_H(x))=𝒜_{L,L^{\prime \prime },0}(x)$$
and that the map $`b_H`$ identifies the geometry of the two action functionals. Indeed for $`u:\times [0,1]M`$ with $`u(,0)L`$, $`u(,1)L^{\prime \prime }`$, $`\stackrel{~}{u}(s,t)=\varphi _t(u(s,t))`$, $`\stackrel{~}{J}=\varphi ^{}J`$ we have
$$\varphi _{}(\frac{u}{s}+\stackrel{~}{J}\frac{u}{t})=\frac{\stackrel{~}{u}}{s}+J(\frac{\stackrel{~}{u}}{t}X_H).$$
Therefore, the map $`b_H`$ induces diffeomorphisms:
$$b_H:_{L,L^{\prime \prime },\stackrel{~}{J},0}(x,y)_{L,L^{},J,H}(x,y)$$
where we have identified $`x,yLL^{\prime \prime }`$ with their orbits $`\varphi _t^H(x)`$ and $`\varphi _t^H(y)`$.
Finally, the non-compactness of $`(x,y)`$ for $`x,yI(L,L^{};\eta ,H)`$ is due to the fact that, as in the Morse-Smale case, a sequence of strips $`u_n(x,y)`$ might “converge” (in the sense of Gromov) to a broken strip. There are natural compactifications of the moduli spaces $`(x,y)`$ called Gromov compactifications and denoted by $`\overline{}(x,y)`$ so that each of the spaces $`\overline{}(x,y)`$ is a topological manifold with corners whose boundary verifies:
(15)
$$\overline{}(x,y)=\underset{zI(L,L^{};\eta ,H)}{}\overline{}(x,z)\times \overline{}(z,y).$$
A complete proof of this fact can be found in (when $`dim(_y^x)=1`$ the proof is due to Floer and is now classical).
#### 3.2.2. Pseudoholomorphic strips and the Serre spectral sequence
We will now construct a complex $`𝒞(L,L^{};H,J)`$ by a method that mirrors the construction of $`𝒞(f)`$ in §2.1.4.
This complex, called the extended Floer complex associated to $`L,L^{},H,J`$ has the form:
$$𝒞(L,L^{};H,J)=(C_{}(\mathrm{\Omega }L)/2<I(L,L^{};\eta ,H)>,D)$$
where the cubical chains $`C_{}(\mathrm{\Omega }L)`$ have, as before, $`/2`$-coefficients. If needed, the moduli spaces $`(x,y)`$ can be endowed with orientations which are compatible with formula (15), and so we could as well use $``$-coefficients.
To define the differential we first fix a simple path $`w`$ in $`L`$ which joins all the points $`\gamma (0)`$, $`\gamma I(L,L^{};H)`$ and we identify all these points to a single one by collapsing this path to a single point. We shall continue to denote the resulting space by $`L`$ to simplify notation.
For each moduli space $`(x,y)`$ there is a continuous map
$$l_y^x:(x,y)\mathrm{\Omega }L$$
which is defined by associating to $`u(x,y)`$ the path $`u(,0)`$ parametrized by the (negative) values of the action functional $`𝒜`$. This is a continuous map and it is seen to be compatible with formula (15) in the same sense as in (8).
We pick a representing chain system $`\{k_y^x\}`$ for the moduli spaces $`(x,y)`$ and we let
$$m_y^x=(l_y^x)_{}(k_y^x)C_{}(\mathrm{\Omega }L)$$
and
(16)
$$Dx=\underset{y}{}m_y^xy.$$
As in the case of the extended Morse complex the fact that $`D^2=0`$ is an immediate consequence of formula (15).
###### Remark 3.4.
a. There is an apparent assymetry between the roles of $`L`$ and $`L^{}`$ in the definition of this extended Floer complex. In fact, the coefficients of this complex belong naturally to an even bigger and more symmetric ring than $`C_{}(\mathrm{\Omega }L)`$. Indeed, consider the space $`T(L,L^{})`$ which is the homotopy-pullback of the two inclusions $`LM`$, $`L^{}M`$. This space is homotopy equivalent to the space of all the continuous paths $`\gamma :[0,1]M`$ so that $`\gamma (0)L`$, $`\gamma (1)L^{}`$. By replacing both $`L`$ and $`M`$ by the respective spaces obtained by contracting the path $`w`$ to a point, we see that there are continous maps $`(x,y)\mathrm{\Omega }(T(L,L^{}))`$. We may then use these maps to construct a complex with coefficients in $`C_{}(\mathrm{\Omega }(T(L,L^{}))`$. Clearly, there is an obvious map $`T(L,L^{})L`$ and it is precisely this map which, after looping, changes the coefficients of this complex into those of the extended Floer complex.
b. At this point it is worth mentioning why using representing chain systems is useful in our constructions. Indeed, for the extended Morse complex representing chain systems are not really essential: the moduli spaces $`M_y^x`$ are triangulable in a way compatible with the boundary formula and so, to represent this moduli space inside the loop space $`\mathrm{\Omega }M`$, we could use instead of the chain $`a_y^x`$ a chain given by the sum of the top dimensional simplexes in such a triangulation. This is obviously, a simpler and more natural approach but it has the disadvantage that it does not extend directly to the Floer case. The reason is that it is not known whether the Floer moduli spaces $`(x,y)`$ admit coherent triangulations (even if this is likely to be the case).
The chain complex $`𝒞(L,L^{};H)`$ admits a natural degree filtration which is given by
(17)
$$F^k𝒞(L,L^{};H,J)=C_{}(L)/2<xI(L,L^{};\eta ,H):|x|k>.$$
It is clear that this filtration is differential. Therefore, there is an induced spectral sequence which will be denoted by $`(L,L^{};H,J)=(_{p,q}^r,D^r)`$. We write $`(L,L^{};J)=(L,L^{};0,J)`$. For convenience we have omitted $`\eta `$ from this notation (the relevant components of the paths spaces $`𝒫(L,L^{})`$ will be clear below).
Here is the main result concerning this spectral sequence.
###### Theorem 3.5.
For any two regular pairs $`(H,J),(H^{},J^{})`$, the spectral sequences $`(L,L^{};H,J)`$ and $`(L,L^{};H^{},J^{})`$ are isomorphic up to translation for $`r2`$. Moreover, if $`\varphi `$ is a hamiltonian diffeomorphism , then $`(L,L^{};J)`$ and $`(L,\varphi (L^{});J^{})`$ are also isomorphic up to translation for $`r2`$ (whenever defined). The second term of this spectral sequence is $`^2(L,L^{};H,J)H_{}(\mathrm{\Omega }L)HF_{}(L,L^{})`$ where $`HF_{}(,)`$ is the Floer homology. Finally, if $`L`$ and $`L^{}`$ are hamiltonian isotopic, then $`(L,L^{};J)`$ is isomorphic up to translation to the Serre spectral sequence of the path-loop fibration $`\mathrm{\Omega }LPLL`$.
Isomorphism up to translation of two spectral sequences $`E_{p,q}^r`$, $`F_{p,q}^r`$ means that there exists a $`k`$ and chain isomorphisms $`\varphi ^r:E_{p,q}^rF_{p+k,q}^r`$. This notion appears naturally here because the choice of the element $`x_0I(L,L^{};H)`$ with $`|x_0|=0`$ is arbitrary. A different choice will simply lead to a translated spectral sequence. As follows from the discussion in §3.2.3 B, it is possible to replace this choice of grading with one that only depends on the path $`\eta `$. However, this might make the absolute degrees fractionary and, as the choice of $`\eta `$ is not canonical, the resulting spectral sequence will still be invariant only up to translation.
The outline of the proof of this theorem is as follows (see for details). First, in view of the naturality properties of Floer’s construction, it is easy to see that the second invariance claim in the statement is implied by the first one. Now, we consider a homotopy $`G`$ between $`H`$ and $`H^{}`$ as well as a one-parameter family of almost complex structures $`\overline{J}`$ relating $`J`$ to $`J^{}`$. For $`xI(L,L^{};H,J)`$ and $`yI(L,L^{};H^{},J^{})`$ we define moduli spaces $`𝒩(x,y)`$ which are solutions of an equation similar to (10) but replaces $`H`$ with $`G`$, $`J`$ with $`\overline{J}`$ (and takes into account the additional parameter - this is a standard construction in Floer theory). These moduli spaces have properties similar to the $`(x,y)`$’s. In particular they admit compactifications which are manifolds with boundary so that the following formula is valid
$$\overline{𝒩}(x,y)=\underset{zI(L,L^{};H)}{}\overline{}(x,z)\times \overline{𝒩}(z,y)\underset{z^{}I(L,L^{};H^{})}{}\overline{𝒩}(x,z^{})\times \overline{}(z^{},y).$$
The representing chain idea can again be used in this context and it leads to coefficients $`n_y^xC_{}(\mathrm{\Omega }L)`$. If we group these coefficients in a matrix $``$ and we group the coefficients of the differential of $`𝒞(L,L^{};H,J)`$ in a matrix $`𝒜`$ and the coefficients of $`𝒞(L,L^{};H^{},J^{})`$ in a matrix $`𝒜^{}`$, then the relation above implies that we have:
(18)
$$=𝒜+𝒜^{}.$$
It follows that the module morphism
$$\varphi _{G,\overline{J}}:𝒞(L,L^{};H,J)𝒞(L,L^{};H,J)$$
which is the unique extension of
$$\varphi _{G,\overline{J}}(x)=\underset{y}{}r_y^xy,xI(L,L^{};H)$$
is a chain morphism. Moreover, the chain morphism constructed above preserves filtrations (of course, to for this it is required that the choices for the point $`x_0`$ with $`|x_0|=0`$ for our two sets of data be coherent - this is why the isomorphisms are “up to translation”). After verifying that $`^2H_{}(\mathrm{\Omega }L)FH_{}(L)`$ for both spectral sequences it is not difficult to see that $`\varphi _G`$ induces an isomorphism at the $`^2`$-level of these spectral sequences . Hence it also induces an isomorphism for $`r>2`$.
For the last point of the theorem we use Floer’s reduction of the moduli spaces $`_{J,L,L^{}}(x,y)`$ of pseudoholomorphic strips to moduli spaces of Morse flow lines $`M_y^x(f)`$. In short, this shows that for certain choices of $`J`$, $`f`$ and $`L^{}`$ which is hamiltonian isotopic to $`L`$ we have homeomorphisms $`\psi _{x,y}:(x,y)M_y^x`$ which are compatible with the compactification and with the boundary formulae. This means that with these choices we have an isomorphism $`𝒞(L,L^{})𝒞(f)`$ and it is easy to see that this preserves the filtrations of these two chain complexes. By Theorem 2.4 this completes the outline of proof.
###### Remark 3.6.
It is shown in that the $`^1`$-term of this spectral sequence has also some interesting invariance properties.
#### 3.2.3. Applications.
We will discuss here a number of direct corollaries of Theorem 3.5 most (but not all )of which appear in .
A. Localization and Hofer’s metric. An immediate adaptation of Theorem 2.4 provides a statement which is much more flexible. This is a “change of coefficients” or “localization” phenomenon that we now describe.
Assume that $`f:LX`$ is a smooth map. Then we can consider the induced map $`\mathrm{\Omega }f:\mathrm{\Omega }L\mathrm{\Omega }X`$ and we may use this map to change the coefficients of $`𝒞(L,L^{};H,J)`$ thus getting a new complex
$$𝒞_X(L,L^{};H,J)=(C_{}(\mathrm{\Omega }X)/2<I(L,L^{};H)>,D_X)$$
so that $`D_X(x)=_y(\mathrm{\Omega }f)_{}(m_y^x)y`$ where $`m_y^x`$ are the coefficients in the differential $`D`$ of $`𝒞(L,L^{};H,J)`$ (compare with (16)). This complex behaves very much like the one studied in Theorem 3.5. In particular, this complex admits a similar filtration and the resulting spectral sequence, $`_X(L,L^{})`$, has the same invariance properties as those in the theorem and, moreover, for $`L`$, $`L^{}`$ hamiltonian isoptopic this spectral sequence coincides with the Serre spectral sequence of the fibration $`\mathrm{\Omega }XEL`$ which is obtained from the path-loop fibration $`\mathrm{\Omega }XPXX`$ by pull-back over the map $`f`$.
In particular, the homology of this complex coincides with the singular homology of $`E`$. If $`X`$ is just one point, $``$, it is easy to see that the complex $`𝒞_{}(L,L^{};H,J)`$ coincides with the Floer complex.
The complex $`𝒞_X(L,L^{};H,J)`$ may also be viewed as a sort of localization in the following sense. Assume that we are interested to see what pseudo-holomorphic strips pass through a region $`AL`$. Then we may consider the closure $`C`$ of the complement of this region and the space $`L/C`$ obtained by contracting $`C`$ to a point. There is the obvious projection map $`LL/C`$ which can be used in place of $`f`$ above. Now, if some non-vanishing differentials appear in $`_{L/C}(L,L^{})`$ for $`r2`$, then it means that there are some coefficients $`m_y^x`$ so that $`|(m_y^x)|>0`$ and $`(\mathrm{\Omega }f)_{}(m_y^x)0`$. This means that the map
$$(x,y)\stackrel{l_y^x}{}\mathrm{\Omega }L\stackrel{\mathrm{\Omega }f}{}\mathrm{\Omega }(C/L)$$
carries the representing chain of $`(x,y)`$ to a nonvanishing chain in $`C_{}(\mathrm{\Omega }(L/C))`$. But this means that the intersection $`^{}(x,y)A`$ is of dimension equal to $`\mu (x,y)`$.
The typical choice of region $`A`$ is a tubular neighbourhood of some submanifold $`VL`$. In that case $`L/C`$ is the associated Thom space.
Let
$$(L,L^{})=\underset{H,\varphi _1^H(L)=L^{}}{inf}(\underset{x,t}{\mathrm{max}}H(x,t)\underset{x,t}{\mathrm{min}}H(x,t)).$$
be the Hofer distance between Lagrangians. It has been shown to be non-degenerate by Chekanov for symplectic manifolds with geometry bounded at infinity.
###### Corollary 3.7.
Let $`aH_k(L)`$ be a non-trivial homology class. If a closed submmanifold $`VL`$ represents the class $`a`$, then for any generic $`J`$ and any $`L^{}`$ hamiltonian isotopic to $`L`$, there exists a pseudoholomorphic strip $`u`$ of Maslov index at most $`nk`$ which passes through $`V`$ and which verifies:
$$u^{}\omega (L,L^{}).$$
In view of the discussion above, the proof is simple (we are using here a variant of that used in ).
We start with a simple topological remark. Take $`A`$ to be a tubular neighbourhood of $`V`$. Then $`L/C=TV`$ is the associated Thom space. In the Serre spectral sequence of $`\mathrm{\Omega }(TV)P(TV)TV`$ we have that $`d^{nk}(\tau )0`$ where $`\tau H_{nk}(TV)`$ is the Thom class of $`V`$. By Poincaré duality, there is a class $`bH_{nk}(TV)`$ which is taken to $`\tau `$ by the projection map $`LTV`$ ($`b`$ corresponds to the Poincaré dual $`a^{}`$ of $`a`$ via the isomorphism $`H^{nk}(L)H_k(L)`$). This means that $`D^{nk}(b)`$ is not zero in $`_{TV}(L,L^{};H)`$ (for any Hamiltonian $`H`$).
To proceed with the proof notice that, by the naturality properties of the Floer moduli spaces, it is sufficient to show that for any Hamiltonian $`H`$ (and any generic $`J`$) so that $`\varphi _1^H(L)=L^{}`$ there exists an element $`u_{L,L,J,H}^{}`$ which is of Maslov index $`(nk)`$ and so that $`u(,0)A\mathrm{}`$ and $`E_{L,L,H}(u)H`$ where $`H=\mathrm{max}_{x,t}H(x,t)\mathrm{min}_{x,t}H(x,t)`$. We may assume that $`\mathrm{min}_{x,t}H(x,t)=0`$ and we let $`K=\mathrm{max}_{x,t}H(x,t)`$. We consider a Morse function $`f:L`$ which is very small in $`C^2`$ norm and we extend it to a function (also denoted by $`f`$) which is defined on $`M`$ and remains $`C^2`$ small. In particular, we suppose $`\mathrm{min}_xf(x)=0`$ and $`\mathrm{max}_xf(x)<ϵ`$. We denote $`\underset{¯}{f}=fϵ`$ and $`\overline{f}=f+K`$. It follows that we may construct monotone homotopies $`G:\overline{f}H`$ and $`G^{}:H\underset{¯}{f}`$. Consider the following action filtration of $`𝒞_{TV}(L,L;H)`$
$$F_v𝒞_{TV}(L,L;J,H)=C_{}(\mathrm{\Omega }TV)/2<xI(L,L;H):𝒜_{L,L,H}(x)v>$$
and similarly on the complexes $`𝒞_{TV}(L,L;\overline{f})`$ and $`𝒞_{TV}(L,L;\underset{¯}{f})`$. It is obvious that this is a differential filtration and, if the choice of path $`\eta `$ (used to define the action functional, see (9)) is the same for all the three hamiltonians involved, these monotone homotopies preserve these filtrations.
We now denote
$$𝒞=F_{K+ϵ}𝒞_{TV}(L,L^{};\overline{f})/F_ϵ𝒞_{TV}(L,L^{};\overline{f}),$$
$$𝒞^{}=F_{K+ϵ}𝒞_{TV}(L,L^{};H)/F_ϵ𝒞_{TV}(L,L^{};H),$$
$$𝒞^{\prime \prime }=F_{K+ϵ}𝒞_{TV}(L,L^{};\underset{¯}{f})/F_ϵ𝒞_{TV}(L,L^{};\underset{¯}{f})$$
These three complexes inherit degree filtrations and there are associated spectral sequences $`(𝒞)`$, $`(𝒞^{})`$, $`(𝒞^{\prime \prime })`$. We have induced morphisms $`\varphi _G:𝒞𝒞^{}`$ and $`\varphi _G^{}:𝒞^{}𝒞^{\prime \prime }`$ which also induce morphisms among these spectral sequences. Moreover, as $`𝒞`$ and $`𝒞^{\prime \prime }`$ both coincide with $`𝒞_{TV}(f)`$ (because $`f`$ is very $`C^2`$-small and $`0f(x)<ϵ`$), the composition $`\varphi _G^{}\varphi _G`$ induces an isomorphism of spectral sequences for $`r2`$ (here $`𝒞_{TV}(f)`$ is the extended Morse complex obtained from $`𝒞(f)`$ by changing the coefficients by the map $`LTV`$). But, as the class $`b`$ has the property that its $`D^{nk}`$ differential is not trivial in $`(𝒞)`$, this implies that $`D^{nk}(\varphi _G(b))0`$ which is seen to immediately imply that there is some moduli space $`^{}(x,y)`$ of dimension $`nk`$ with $`𝒜_{L,L^{};H}(x),𝒜_{L,L^{};H}(y)[ϵ,K+ϵ]`$ and $`^{}(x,y)A\mathrm{}`$. Therefore there are $`J`$-strips passing through $`A`$ which have Maslov index $`nk`$ and area less than $`H+2ϵ`$. By letting $`A`$ tend to $`V`$ and $`ϵ0`$, $`H(L,L^{})`$, these strips converge to strips with the properties desired.
We may apply this even to $`1H_0(L)`$ and Corollary 3.7 shows in this case that through each point of $`L`$ passes a strip of Maslov index at most $`n`$ (again, for $`J`$ generic) and of area at most $`(L,L^{})`$. The case $`V=pt`$ was discussed explicitly in .
###### Remark 3.8.
a. It is clear that the strips detected in this corollary actually have a symplectic area which is no larger than $`c(b;H)c(1;H)`$ where $`c(x;H)`$ is the spectral value of the homology class $`x`$ relative to $`H`$,
$$c(x;H)=inf\{v:xIm(H_{}(FC_v)(L,L;H)FH_{}(L,L;H))\}$$
where $`FC_v(L,L;H)`$ is the Floer complex of $`L,L,H`$ generated by all the elements of $`I(L,L;\eta ,H)`$ of action smaller or equal than $`v`$; $`FH_{}(L,L;H)`$ is the Floer homology. Under our assumptions we have a canonical isomorphism (up to translation) between $`HF_{}(L,L,H)`$ and $`H_{}(L)`$ so we may view $`bHF_{}(L,L;H)`$.
b. Clearly, in view of Gromov compactness our result also implies that for any $`J`$ (even non regular) and for any $`L^{}`$ hamiltonian isotopic to $`L`$ and for any $`xL\backslash L^{}`$ there exists a $`J`$-holomorphic strip passing through $`x`$ which has area less than $`(L,L^{})`$. This result (without the area estimate) also follows from independent work of Floer and Hofer . Another method has been mentioned to us by Dietmar Salamon. It is based on starting with disks with their boundary on $`L`$ and which are very close to be constant maps. Therefore, an appropriate evaluation defined on the moduli space of these disks is of degree $`1`$. Each of these disks is made out of two semi-disks which are pseudo-holomorphic and which are joined by a short semi-tube verifying the non-homogenous Floer equation for some given Hamiltonian $`H`$. This middle region is then allowed to expand till, at some point, it will necessarily produce a semi-tube belonging to some $`_H^{}(x,y)`$. It is also possible to use the pair of pants product to produce Floer orbits joining the “top and bottom classes” . Still, having simultaneous area and Maslov index estimates appears to be more difficult by methods different from ours. Of course, detecting strips of lower Maslov index so that they meet some fixed submanifold is harder yet.
Corollary 3.7 has a nice geometric application.
###### Corollary 3.9.
Assume that, as before, $`L`$ and $`L^{}`$ are hamiltonian isotopic. For any symplectic embedding $`e:(B_r,\omega _0)M`$ so that $`e^1(L)=^nB(r)`$ and $`e(B_r)L^{}=\mathrm{}`$ we have $`\pi r^2/2(L,L^{})`$.
This is proven (see ) by using a variant of the standard isoperimetric inequality: a $`J_0`$-pseudoholomorphic surface in the standard ball $`(B_r,\omega _0)`$ of radius $`r`$ whose boundary is on $`B_r^n`$ has area at least $`\pi r^2/2`$.
Clearly, this implies the non-degeneracy result of Chekanov that was mentioned before under the connectivity conditions that we have always assumed till this point.
B. Relaxing the connectivity conditions. We have worked till now under the assumption that
(19)
$$L,L^{}\mathrm{are}\mathrm{simply}\mathrm{connected}\mathrm{and}\omega |_{\pi _2(M)}=c_1|_{\pi _2(M)}=0$$
These requirements were used in a few important places: in the definition of the action functional, the definition of the Maslov index, the boundary product formula (because they forbid bubbling). Of these, only the bubbling isssue is in fact essential: the boundary formula is precisely the reason why $`d^2=0`$ as well as the cause of the invariance of the resulting homology.
We proceed below to extend the corollaries and techniques discussed above to the case when all the connectivity conditions are dropped but we assume that $`L`$ and $`L^{}`$ are hamiltonian isotopic and only work below the minimal energy that could produce some bubbling (this is similar to the last section of but goes beyond the cases treated there).
First, for a time dependent almost complex structure $`J_t`$, $`t[0,1]`$, we define $`\delta _{L,L^{}}(J)`$ as the infimmum of the symplectic areas of the following three types of objects:
* the $`J_t`$-pseudoholomorphic spheres in $`M`$ (for $`t[0,1]`$).
* the $`J_0`$-pseudoholomorphic disks with their boundary on $`L`$.
* the $`J_1`$-pseudoholomorphic disks with their boundary on $`L^{}`$.
By Gromov compactness this number is strictly positive.
We will proceed with the construction in the case when $`L=L^{}`$ and in the presence of a hamiltonian $`H`$. We shall assume that the pair $`(H,J)`$ is regular in the sense that the moduli spaces of strips defined below, $`(x,y)`$, are regular.
We take the fixed reference path $`\eta `$ to be a constant point in $`L`$ (see (9)). Denote $`𝒫_0(L,L)=𝒫_\eta (L,L)`$ and consider in this space the base point given by $`\eta `$. Notice that there is a morphism $`\omega :\pi _1𝒫_0(L,L)`$ obtained by integrating $`\omega `$ over the disk represented by the element $`z\pi _1𝒫_0(L,L)`$ (such an element can be viewed as a disk with boundary in $`L`$). Similarly, let $`\mu :\pi _1𝒫_0(L,L)`$ be the Maslov morphism. Let $`𝒦`$ be the kernel of the morphism
$$\omega \times \mu :\pi _1𝒫_0(L,L)\times .$$
The group
$$\pi =\pi _1(𝒫_0(L,L))/𝒦$$
is an abelian group (as it is a subgroup of $`\times `$) and is of finite rank. Let’s also notice that this group is the quotient of $`\pi _2(M,L)`$ by the equivalence relation $`ab`$ iff $`\omega (a)=\omega (b),\mu (a)=\mu (b)`$ (with this definition this group is also known as the Novikov group). This is a simple homotopical result. First $`𝒫(L,L)`$ is the homotopy pull-back of the map $`LM`$ over the map $`LM`$. But this means that we have a fibre sequence $`F𝒫(L,L)L`$ with $`F`$ the homotopy fibre of $`LM`$ and that this fibre sequence admits a canonical section. This implies that
$$\pi _1𝒫_0(L,L)\pi _1(F)\times \pi _1(L).$$
But $`\pi _1(F)=\pi _2(M,L)`$. As both $`\omega `$ and $`\mu `$ are trivial on $`\pi _1(L)`$ the claim follows. It might not be clear at first sight why $`\mu `$ is null on $`\pi _1(L)`$ here. The reason is that the term $`\pi _1(L)`$ in the product above is the image of the map induced in homotopy by
$$j_L:L𝒫_0(L,L).$$
This map associates to a point in $`L`$ the constant path. Consider a loop $`\gamma (s)`$ in $`L`$. Then $`j_L\gamma `$ is a loop in $`𝒫_0(L,L)`$ which at each moment $`s`$ is a constant path. We now need to view this loop as the image of a disk and $`\mu ([j_L(\gamma )])`$ is the Maslov index of this disk. But this disk is null homotopic so $`\mu ([j_L(\gamma )])=0`$.
Consider the regular covering $`p:𝒫_0^{}(L,L)𝒫_0(L,L)`$ which is associated to the group $`𝒦`$. We fix an element $`\eta _0𝒫_0^{}(L,L)`$ so that $`p(\eta _0)=\eta `$. Clearly, the action functional
$$𝒜_{L,L,H}^{}:𝒫_0^{}(L,L)$$
may be defined by essentially the same formula as in (9):
$$𝒜_{L,L,H}^{}(x)=(pu)^{}\omega +_0^1H(t,(px)(t))𝑑t$$
where $`u:[0,1]𝒫_0^{}(L,L)`$ is such that $`u(0)=\eta _0`$, $`u(1)=x`$ and is now well-defined.
Let $`I^{}(L,L,H)=p^1(I(L,L,H))`$. For $`x,yI^{}(L,L,H)`$ we may define $`\mu (x,y)=\mu (pu)`$ where $`u:[0,1]𝒫_0^{}(L,L)`$ is a path that joins $`x`$ to $`y`$. This is again well defined. For each $`xI^{}(L,L,H)`$ we consider a path $`v_x:[0,1]𝒫_0^{}(L,L)`$ so that $`v_x(0)=\eta _0`$ and $`v_x(1)=x`$. The composition $`pv_x`$ can be viewed as a “semi-disk” whose boundary is resting on the orbit $`p(x)`$ and on $`L`$. Therefore, we may associate to it a Maslov index $`\mu (v_x)`$ and it is easy to see that this only depends on $`x`$. Thus we define $`\mu (x)=\mu (v_x)`$ and we have $`\mu (x,y)=\mu (x)\mu (y)`$ for all $`x,yI^{}(L,L,H)`$.
To summarize what has been done till now: once the choices of $`\eta `$ and $`\eta _0`$ are made, both the action functional $`𝒜^{}:𝒫_0^{}(L,L,H)`$ and the “absolute” Maslov index $`\mu ():I^{}(L,L,H)`$ are well-defined.
Fix an almost complex structure $`J`$. Consider two elements $`x,yI^{}(L,L,H)`$. We may consider the moduli space which consists of all paths $`u:𝒫_0^{}(L,L)`$ which join $`x`$ to $`y`$ and are so that $`pu`$ satisfies Floer’s equation (10) modulo the $``$-action. If regularity is achieved, the dimension of this moduli space is precisely $`\mu (x,y)1`$. The action functional $`𝒜^{}`$ decreases along such a solution $`u`$ and the energy of $`u`$ (which is defined as the energy of $`pu`$) verifies, as in the standard case, $`E(u)=𝒜^{}(x)𝒜^{}(y)`$. Bubbling might of course be present in the compactification of these moduli spaces. As we only intend to work below the minimal bubbling energy $`\delta _{L,L}(J)`$ we artificially put:
$$(x,y)=\mathrm{}\mathrm{if}𝒜^{}(x)𝒜^{}(y)\delta _{L,L}(J)$$
and, of course, for $`𝒜^{}(x)𝒜^{}(y)<\delta _{L,L}(J)`$, $`(x,y)`$ consists of the elements $`u`$ mentioned above. We only require these moduli spaces to be regular.
With this convention, for all $`x,yI^{}(L,L,H)`$ so that $`(x,y)`$ is not void we have the usual boundary formula (15). Notice at the same time that this formula is false for general pairs $`x,y`$ (and so there is no way to define a Floer type complex at this stage).
Now consider a map $`f:LX`$ so that $`X`$ is simply-connected (the only reason to require this is to insure that the Serre spectral sequence does not require local coefficients).
We consider the group:
$$𝒞(L,L,H;X)=C_{}(\mathrm{\Omega }X)/2<I^{}(L,L,H)>.$$
For $`wv`$, we denote $`I_{v,w}^{}=\{xI^{}(L,L,H):w𝒜^{}(x)v\}`$ and we define the subgroup
$$𝒞_{w,v}(L,L,H;X)=C_{}(\mathrm{\Omega }X)/2<I_{v,w}^{}>.$$
Suppose that $`wv\delta _{L,L}(J)ϵ`$. We claim that in this case we may define a differential on $`𝒞_{v,w}(L,L,H;X)`$ by the usual procedure. Consider representing chain systems for all the moduli spaces $`(x,y)`$ and let the image of these chains inside $`C_{}(\mathrm{\Omega }X)`$ be respectively $`\overline{m}_y^x`$. Let $`D`$ be the linear extension of the map given by
$$Dx=\underset{yI_{v,w}^{}}{}\overline{m}_y^xy.$$
###### Proposition 3.10.
The linear map $`D`$ is a differential. A generic monotone homotopy $`G`$ between two hamiltonians $`H`$ and $`H^{}`$
$$\varphi _G^X:𝒞_{v,w}(L,L,H,J;X)𝒞_{v,w}(L,L,H^{},J;X).$$
A monotone homotopy between monotone homotopies $`G`$ and $`G^{}`$ induces a chain homotopy between $`\varphi _G^X`$ and $`\varphi _G^{}^X`$ so that $`H_{}(\varphi _G^X)=H_{}(\varphi _G^{}^X)`$.
Now, $`D^2x=_z(_y\overline{m}_y^x\overline{m}_z^y+\overline{m}_z^x)z.`$ In this formula we have $`𝒜^{}(x)𝒜^{}(z)\delta _{L,L}(J)`$ and, because the usual boundary formula (15) is valid in this range, all the terms vanish so that $`D^2(x)=0`$. The same idea may be applied to a monotone homotopy as well as to a monotone homotopy between monotone homotopies and it implies the claim.
###### Remark 3.11.
a. If we take for the space $`X`$ a single point $``$ we get a chain complex whose differential only takes into account the $`0`$-dimensional moduli spaces and which is a truncated version of Floer homology.
b. The complex $`𝒞_{v,w}(L,L,H,J;X)`$ admits a degree filtration which is perfectly similar to the one given by (17). Let $`𝒞_{v,w}(L,L,H,J;X)`$ be the resulting spectral sequence. Then, under the restrictions in the Proposition 3.10, a monotone homotopy $`G`$ induces a morphism of spectral sequences $`_X(\varphi _G)`$ and two such homotopies $`G`$, $`G^{}`$ which are monotonously homotopic have the property that they induce the same morphism $`_X^r(\varphi _G)=_X^r(\varphi _G^{})`$ for $`r2`$. This last fact follows from Proposition 3.10 by computing $`_X^2(\varphi _G)=id_{H_{}(\mathrm{\Omega }X)}H_{}(\varphi _G^{})=_X^2(\varphi _G^{})`$ where $`\varphi _G^{}:𝒞_{v,w}(L,L,H;)𝒞_{v,w}(L,L,H^{};)`$.
Naturally, the next step is to compare our construction with its Morse theoretical analogue. Consider the map $`j_L:L𝒫_0(L,L)`$ and consider $`p:\stackrel{~}{L}L`$ the regular covering obtained by pull-back from $`𝒫_0^{}(L,L)𝒫_0(L,L)`$. Notice that, because both compositions $`\omega \pi _1(j_L)`$ and $`\mu \pi _1(j_L)`$ are trivial, it follows that the covering $`\stackrel{~}{L}L`$ is trivial. Let $`\overline{f}:\stackrel{~}{L}`$ be defined by $`\overline{f}=fp`$ and consider $`𝒞(\overline{f};X)`$ the extended Morse complex of $`\overline{f}`$ with coefficients changed by the map $`\mathrm{\Omega }\stackrel{~}{L}\mathrm{\Omega }X`$. Notice that, in general, the group $`\pi `$ acts on $`I^{}(L,L,H)`$ and we have the formula:
$$𝒜^{}(gy)=\omega (g)+𝒜^{}(y),\mu (gy)=\mu (g)+\mu (y),yI^{}(L,L,H),g\mathrm{\Pi }.$$
In our particular case, when $`H=f`$, we have $`I^{}(L,L,H)=Crit(\overline{f})`$. For each point $`xCrit(f)`$ let $`\overline{x}Crit(\overline{f})`$ be the element of $`p^1(x)`$ which belongs to the component of $`\stackrel{~}{L}`$ which also contains $`\eta _0`$. We then have $`𝒜^{}(\overline{x})=f(x)`$ and $`\mu (\overline{x})=ind_{\overline{f}}(\overline{x})`$. The extended Morse complex $`𝒞(\overline{f};X)`$ is therefore isomorphic to $`𝒞(f;X)[\pi ]`$ and the action filtration is determined by writting $`𝒜^{}(xg)=f(x)+\omega (g)`$. The degree filtration induces, as usual, a spectral sequence which will be denoted by $`𝒞(f;X)`$.
The remarks above together with Theorem 3.5 show that this spectral sequence consists of copies the Serre spectral sequence of $`\mathrm{\Omega }XEL`$: one copy for each connected component of $`\stackrel{~}{L}`$. We denote by $`𝒞_0(f;X)`$ and $`𝒞_0(f;X)`$ the copies of the extended complex and of the spectral sequence that correspond to the connected component $`L_0`$ of $`\stackrel{~}{L}`$ which contains $`\eta _0`$.
###### Proposition 3.12.
Suppose $`H_0<\delta _J(L,L)`$. There exists a chain morphisms $`\varphi :𝒞_0(f;X)𝒞_{0,H}(L,L,H,J;X)`$ and $`\psi :𝒞_{0,H}(L,L,H,J;X)𝒞_0(f;X)`$ which preserve the respective degree filtrations and so that $`\psi \varphi `$ induces an isomorphism at the $`E^2`$ level of the respective spectral sequences.
To prove this proposition we shall use a different method than the one used in Corollary 3.7. The comparison maps $`\varphi `$, $`\psi `$ will be constructed by the method introduced in and later used in and . Compared to Proposition 3.10 this is particularly efficient because it avoids the need to control the bubbling thershold along deformations of $`J`$.
We fix as before the Morse function $`f`$ as well as the pair $`H,J`$. To simplify the notation we shall assume that $`infH(x,t)=0`$. The construction of $`\varphi `$ is based on defining certain moduli spaces $`𝒲(x,y)`$ with $`xCrit(\overline{f})`$ and $`yI^{}(L,L,H)`$. They consist of pairs $`(u,\gamma )`$ where $`u:𝒫_0^{}(L,L)`$, $`\gamma :(\mathrm{},0]\stackrel{~}{L}`$ and if we put $`u^{}=p(u)`$, $`\gamma ^{}=p(\gamma )`$ ($`p:𝒫_0^{}(L,L)𝒫_0(L,L)`$ is the covering projection) then we have:
$$u^{}(R\times \{0,1\})L,_s(u^{})+J(u^{})_t(u^{})+\beta (s)H(u^{},t)=0,u(+\mathrm{})=y$$
and
$$\frac{d\gamma ^{}}{dt}=_gf(\gamma ^{}),\gamma (\mathrm{})=x,\gamma (0)=u(\mathrm{}).$$
Here $`g`$ is a Riemannian metric so that $`(f,g)`$ is Morse-Smale and $`\beta `$ is a smooth cutt-off function which is increasing and vanishes for $`s1/2`$ and equals $`1`$ for $`s1`$. It is useful to view an element $`(u,\gamma )`$ as before as a semi-tube connecting $`x`$ to $`y`$. Under usual regularity assumptions these moduli spaces are manifolds of dimension $`\mu (x)\mu (y)`$.
The energy of such an element $`(u,\gamma )`$ is defined in the obvious way by $`E(u,\gamma )=_su^{}^2𝑑s𝑑t`$. A simple computation shows that:
$$E(u,\gamma )=I(u)+_{\times [0,1]}(u^{})^{}\omega _0^1H(y(t))𝑑t$$
where $`I(u)=_{\times [0,1]}\beta ^{}(s)H(u^{}(s),t)𝑑s𝑑t`$. If $`xL_0`$, then the energy verifies
$$E(u,\gamma )=I(u)𝒜^{}(y)sup(H)𝒜^{}(y).$$
As before, we only want to work here under the bubbling threshold and we are only interested in the critical points $`xL_0`$ so we put $`𝒲(x,y)=\mathrm{}`$ for all those pairs $`(x,y)`$ with either $`yI^{}(L,L,H)`$ so that $`𝒜^{}(y)[0,H]`$ or with $`xL_0`$. This means that there is no bubbling in our moduli spaces. Thus we may apply the usual procedure: compactification, representing chain systems, representation in the loop space (for this step we need to choose a convenient way to parametrize the paths represented by the elements $`(u,\gamma )`$). Notice that the boundary of $`\overline{𝒲}(x,y)`$ is the union of two types of pieces $`\overline{M}_z^x\times \overline{𝒲}(z,y)`$ and $`\overline{𝒲}(x,z)\times \overline{}(z,y)`$. We then define $`\varphi (x)=w_y^xy`$ where $`w_y^x`$ is a cubical chain representing the moduli space $`𝒲(x,y)`$. This is a chain map as desired.
We now proceed to construct the map $`\psi `$. The construction is prefectly similar: we define moduli spaces $`𝒲^{}(y,x)`$, $`yI^{}(L,L,J)`$, $`xCrit(\overline{f})`$ except that the pairs $`(u,\gamma )`$ considered here, start as semi-tubes and end as flow lines of $`f`$. The equation verified by $`u`$ is similar to the one before but istead of the cut-off function $`\beta `$ we use the cut-off function $`\beta `$. For $`xL_0`$ the energy estimate in this case gives $`E(u,\gamma )𝒜^{}(y)`$. By the same method as above we define $`𝒲^{}(y,x)`$ to be void whenever $`xL_0`$ or $`𝒜^{}(y)>H`$ and we define $`\psi (y)=\overline{w}_x^yx`$ where $`\overline{w}_x^y`$ is a cubical chain representing the moduli space $`W^{}(y,x)`$. Notice that because $`E(u,\gamma )𝒜^{}(y)`$ this map $`\psi `$ does in fact vanish on $`𝒞_0(L,L,H,J;X)`$ and so it induces a chain map (also denoted by $`\psi `$) as desired.
The next step is to notice that the composition $`\psi \varphi `$ induces an isomorphism at $`E^2`$. This is equivalent to showing that $`H_{}(\psi \varphi )`$ is an isomorphism for $`X=`$. In turn, this fact follows by now standard deformation arguments as in .
###### Corollary 3.13.
Assume that $`L`$ and $`L^{}`$ are hamiltonian isotopic and suppose that $`J`$ is generic. If $`(L,L^{})<\delta _{L,L^{}}(J)`$, then the statement of Corollary 3.7 remains true (for $`J`$) without the connectivity assumption (19).
Notice that if $`H`$ is a hamiltonian so that $`\varphi _1^H(L)=L^{}`$ and $`J_H=(\varphi ^H)_{}(J)`$ then, by the naturality described in §3.2, we have:
$$\delta _{L,L}(J_H)=\delta _{L,L^{}}(J).$$
This implies, again by this same naturality argument, that the problem reduces to finding appropriate semi-tubes whose detection comes down to showing the non-vanishing of certain differentials in $`𝒞_{0,w}(L,L,H;TV)`$ for some well chosen $`w<\delta _{L,L}(J_H)`$. But this immediatly follows from Proposition 3.12 by the same topological argument as the one used in the proof of Corollary 3.7
We formulate the geometric consequence which corresponds to Corollary 3.9. For two lagrangians $`L`$ and $`L^{}`$ the following number has been introduced in : $`B(L,L^{})`$ is the supremum of the numbers $`r0`$ so that there exists a symplectic embedding
$$e:(B(r),\omega _0)(M,\omega )$$
so that $`e^1(L)=^nB(r)`$ and $`Im(e)L^{}=\mathrm{}`$.
###### Corollary 3.14.
There exists an almost complex structure $`J`$ so that we have the inequality:
$$(L,L^{})\mathrm{min}\{\delta _{L,L^{}}(J),\frac{\pi }{2}B(L,L^{})^2\}.$$
Clearly, this implies that $`(,)`$ is non-degenerate in full generality (and recovers, in particular, the fact that the usual Hofer norm for Hamiltonians is non-dgenerate). It is useful to also notice as in that this result is a Lagrangian version of the usual capacity - displacement energy inequality . Indeed, this inequality (with the factor $`\frac{1}{2}`$) is implied by the following statement which has been conjectured to hold for any two compact lagrangians in a symplectic manifold :
(20)
$$(L,L^{})\frac{\pi }{2}B(L,L^{})^2.$$
This remains open. An even stronger conjecture is the following:
###### Conjecture 3.15.
For any two hamiltonian isotopic closed lagrangians $`L,L^{}(M,\omega )`$ and for any almost complex structure $`J`$ which compatible with $`\omega `$ and any point $`xL\backslash L^{}`$ there exists a pseudoholomorphic curve $`u`$ which is either a strip resting on $`L`$ and on $`L^{}`$ or a pseudoholomorphic disk with boundary in $`L`$ so that $`xIm(u)`$ and $`u^{}\omega (L,L^{})`$.
By the isoperimetric inequality used earlier in this paper, it follows that this statement implies (20). There is a substantial amount of evidence in favor of this conjecture:
* as explained in this paper, in the absence of any pseudoholomorphic disks (that is, when $`\omega |_{\pi _2(M,L)}=0`$) it was proven in .
* the statement in Corollary 3.14 shows that the area estimate is not unreasonable.
* one striking consequence of Conjecture 3.15 is that if the disjunction energy of the lagrangian $`L`$ is equal to $`E_0<\mathrm{}`$, then, for any $`J`$ as in the statement and any $`xL`$ there is a pseudoholomorpic disk of area at most $`E_0`$ which passes through $`x`$. When $`L`$ is relatively spin, this is indeed true and follows from recent work of the second author joint with François Lalonde. By the same geometric argument as above we deduce a nice consequence. Define the relative (or *real*) Gromov radius of $`L`$, $`Gr(L)`$, to be the supremum of the positive numbers $`r`$ so that there exists a symplectic embedding $`e:(B(r),\omega _0)(M,\omega )`$ with the property that $`e^1(L)=^nB(r)`$, then $`\pi (Gr(L))^2/2E_0`$ (where $`E_0`$, as before, is the disjunction energy of $`L`$). It is also useful to note that if $`L`$ is the zero section of a cotangent bundle, then $`Gr(L)=\mathrm{}`$.
There are numerous other interesting consequences of Conjecture 3.15 besides (20). To conclude, 3.15 appears to be a statement worth investigating.
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# Non-dissipative drag of superflow in a two-component Bose gas
## I Introduction
Macroscopic quantum coherence manifests itself in many specific phenomena. One of them is a non-dissipative drag that takes place in superfluids and superconductors. The non-dissipative drag, also known as the Andreev-Bashkin effect, was considered, for the first time, in Ref. 1 , where a three velocity hydrodynamic model for <sup>3</sup>He-<sup>4</sup>He superfluid mixtures was developed. It was shown that superfluid behavior of such systems can be described under accounting the ”drag” term in the free energy. This term is proportional to the scalar product of the superfluid velocities of two superfluid components. A similar situation may take place in mixtures of superfluids of $`S_z=+1`$ and $`S_z=1`$ pairs in liquid <sup>3</sup>He in the $`A`$-phase 2 . Among other objects, where the non-dissipative drag may be important, are neutron stars, where the mixture of neutron and proton Cooper pair Bose condensates is believed to realize 3 ; 4 . The possibility of realization of the non-dissipative drag in superconductors was considered in 5 . The non-dissipative drag in bilayer Bose systems was treated microscopically in 6 ; 7 for a special case of two equivalent layers of charged bosons. The case of a bilayer system of neutral bosons was studied in my in the limit of small interlayer interaction.
The most promising systems where the non-dissipative drag can be observed experimentally are two-component alkali metal vapors. In such systems the interaction between atoms of different species is of the same order as the interaction between atoms of the same specie and the effect is expected to be larger than in bilayers. In Bose mixtures the components are characterized by different densities, different masses of atoms and different interaction parameters. In this paper we consider such a general case and obtain an analytical expression for the drag current for zero and finite temperatures.
In the system under consideration the drag force influences the dynamics of atoms in the drag component in the same manner as the vector potential of electromagnetic field influences the dynamics of electrons in superconductors. In particular, in neutral superfluids with Josephson links the drag effect may induce the gradient of the phase of the order parameter in the bulk and, as a consequence, control the phase difference between weakly coupled parts of the system. Therefore, on can expect that the effect reveals itself in a modification of Josephson oscillations between weakly coupled Bose gases. In this paper we discuss possible ways for the observation such a modification. We consider the Bose gas confined in a toroidal trap with two Josephson links. In the Fock regime leg the low energy dynamics of the system can be described by the qubit model of general form (the model, where all three components of the pseudomagnetic field can be controlled independently). The parameters of the qubit Hamiltonian depend on the drag factor. The measurement of the state of the qubit under controlled evolution allows to observe the effect caused by the non-dissipative drag and determine the drag factor. In this paper we consider two particular schemes of the measurement. In the first scheme one should determine the time required to transform a reproducible initial state to a given final state. In the second scheme the geometrical (Berry) phase should be detected.
In Sec. II the microscopic theory of the non-dissipative drag in two-component Bose gases is developed. In Sec. III a model of the Bose-Einstein qubit subjected by the drag force is formulated and the schemes of measurement of the drag factor are proposed. Conclusions are given in Sec. IV.
## II Nondissipative drag in a two-component Bose system. Microscopic derivation
Let us consider a uniform two-component atomic Bose gas in a Bose-Einstein condensed state. We will study the most general situation where the densities of atoms in each component are different from one another ($`n_1n_2`$), the atoms of each components have different masses ($`m_1m_2`$) and the interaction between atoms is described by three different scattering lengths ($`a_{11}a_{22}a_{12}`$). The Hamiltonian of the system can be presented in the form
$$H=\underset{i=1,2}{}(E_i\mu _iN_i)+\frac{1}{2}\underset{i,i^{}=1,2}{}E_{ii^{}}^{int},$$
(1)
where
$$E_i=d^3r\frac{\mathrm{}^2}{2m_i}[\widehat{\mathrm{\Psi }}_i^+(𝐫)]\widehat{\mathrm{\Psi }}_i(𝐫)$$
(2)
is the kinetic energy,
$$E_{ii^{}}^{int}=d^3r\widehat{\mathrm{\Psi }}_i^+(𝐫)\widehat{\mathrm{\Psi }}_i^{}^+(𝐫)\gamma _{ii^{}}\widehat{\mathrm{\Psi }}_i^{}(𝐫)\widehat{\mathrm{\Psi }}_i(𝐫)$$
(3)
is the energy of interaction, $`\gamma _{ii}=4\pi \mathrm{}^2a_{ii}/m_i`$ and $`\gamma _{12}=2\pi \mathrm{}^2(m_1+m_2)a_{12}/(m_1m_2)`$ are the interaction parameters, and $`\mu _i`$ are the chemical potentials.
For the further analysis it is convenient to use the density and phase operator approach (see, for instance, 11 ; 12 ). The approach is based on the following representation for the Bose field operators
$$\widehat{\mathrm{\Psi }}_i(𝐫)=\mathrm{exp}\left[i\phi _i(𝐫)+i\widehat{\phi }_i(𝐫)\right]\sqrt{n_i+\widehat{n}_i(𝐫)},$$
(4)
$$\widehat{\mathrm{\Psi }}_i^+(𝐫)=\sqrt{n_i+\widehat{n}_i(𝐫)}\mathrm{exp}\left[i\phi _i(𝐫)i\widehat{\phi }_i(𝐫)\right],$$
(5)
where $`\widehat{n}_i`$ and $`\widehat{\phi }_i`$ are the density and phase fluctuation operators, $`\phi _i(𝐫)`$ are the $`c`$-number terms of the phase operators, which are connected with the superfluid velocities by the relation $`𝐯_i=\mathrm{}\phi _i/m_i`$. In what follows we specify the case of the superfluid velocities independent of $`𝐫`$.
Substituting Eqs. (4), (5) into Eq. (1) and expanding it in series in powers of $`\widehat{n}_i`$ and $`\widehat{\phi }_i`$ we present the Hamiltonian of the system in the following form
$$H=H_0+H_2+\mathrm{}$$
(6)
In (6) the term
$$H_0=V\left(\underset{i=1,2}{}\left[\frac{1}{2}m_in_i𝐯_i^2+\frac{\gamma _{ii}}{2}n_i^2\mu _in_i\right]+\gamma _{12}n_1n_2\right)$$
(7)
does not contain the operator part. Here $`V`$ is the volume of the system. The minimization conditions for the Hamiltonian $`H_0`$ yield the equations
$$\frac{1}{2}m_i𝐯_i^2+\gamma _{ii}n_i+\gamma _{12}n_{3i}\mu _i=0(i=1,2).$$
(8)
Under the conditions (8) the terms, linear in the density fluctuation operators, vanish in the Hamiltonian. Taking into account the $`(n_i\phi _i(𝐫))=0`$, we find that the terms, linear in the phase fluctuation operators, vanish in the Hamiltonian as well.
The part of the Hamiltonian quadratic in $`\widehat{\phi }_i`$ and $`\widehat{n}_i`$ operators reads as
$`H_2={\displaystyle }d𝐫({\displaystyle \underset{i}{}}\{{\displaystyle \frac{\mathrm{}^2}{2m_i}}[{\displaystyle \frac{\left(\widehat{n}_i(𝐫)\right)^2}{4n_i}}+n_i(\widehat{\phi }_i(𝐫))^2]`$ (9)
$`+{\displaystyle \frac{\mathrm{}𝐯_i}{2}}\left(\widehat{n}_i(𝐫)\widehat{\phi }_i(𝐫)+[\widehat{\phi }_i(𝐫)]\widehat{n}_i(𝐫)\right)`$ (10)
$`+{\displaystyle \frac{i\mathrm{}^2}{2m_i}}([\widehat{n}_i(𝐫)]\widehat{\phi }_i(𝐫)[\widehat{\phi }_i(𝐫)]\widehat{n}_i(𝐫))+{\displaystyle \frac{\gamma _{ii}}{2}}\left(\widehat{n}_i(𝐫)\right)^2\}+\gamma _{12}\widehat{n}_1(𝐫)\widehat{n}_2(𝐫)).`$ (11)
The quadratic part of the Hamiltonian determines the spectra of the elementary excitations. Hereafter we will neglect the higher order terms in the Hamiltonian (6). These terms describe the scattering of the quasiparticles and they can be omitted if the temperature is much smaller than the temperature of Bose-Einstein condensation.
Let us rewrite the quadratic part of the Hamiltonian in terms of the operators of creation and annihilation of the elementary excitations. As the first step, we use the substitution
$$\widehat{n}_i(𝐫)=\sqrt{\frac{n_i}{V}}\underset{𝐤}{}e^{i\mathrm{𝐤𝐫}}\sqrt{\frac{ϵ_{ik}}{E_{ik}}}\left[b_i(𝐤)+b_i^+(𝐤)\right],$$
(12)
$$\widehat{\phi }_i(𝐫)=\frac{1}{2i}\sqrt{\frac{1}{n_iV}}\underset{𝐤}{}e^{i\mathrm{𝐤𝐫}}\sqrt{\frac{E_{ik}}{ϵ_{ik}}}\left[b_i(𝐤)b_i^+(𝐤)\right],$$
(13)
where operators $`b_i^+`$, $`b_i`$ satisfy the Bose commutation relations. Here $`ϵ_{ik}=\mathrm{}^2k^2/2m_i`$ is the spectrum of free atoms, and
$$E_{lk}=\sqrt{ϵ_{ik}(ϵ_{ik}+2\gamma _{ii}n_i)}$$
(14)
is the spectrum of the elementary excitations at $`\gamma _{12}=0`$ and $`𝐯_i=0`$. The substitution (12), (13) reduces the Hamiltonian (9) to the form quadratic in $`b_i^+`$ and $`b_i`$ operators:
$`H_2={\displaystyle \underset{i𝐤}{}}\left[_i(𝐤)\left(b_i^+(𝐤)b_i(𝐤)+{\displaystyle \frac{1}{2}}\right){\displaystyle \frac{1}{2}}ϵ_{ik}\right]`$ (15)
$`+{\displaystyle \underset{𝐤}{}}g_k[b_1^+(𝐤)b_2(𝐤)+b_1(𝐤)b_2(𝐤)+h.c.].`$ (16)
Here
$$_i(𝐤)=E_{ik}+\mathrm{}\mathrm{𝐤𝐯}_i$$
(17)
and
$$g_k=\gamma _{12}\sqrt{\frac{ϵ_{1k}ϵ_{2k}n_1n_2}{E_{1k}E_{2k}}}.$$
(18)
The Hamiltonian (15) contains non-diagonal in Bose creation and annihilation operator terms and it can be diagonalized using the standard procedure of u-v transformation 13 . The result is
$$H_2=\underset{𝐤}{}\left[\underset{\lambda =\alpha ,\beta }{}_\lambda (𝐤)\left(\beta _\lambda ^+(𝐤)\beta _\lambda (𝐤)+\frac{1}{2}\right)\frac{1}{2}\underset{i=1,2}{}ϵ_{ik}\right],$$
(19)
where $`\beta _\lambda ^+(k)`$ and $`\beta _\lambda (k)`$ are the operators of creation and annihilation of elementary excitations.
The energies $`_\lambda (𝐤)`$ satisfy the equation
$$det\left(\begin{array}{cc}𝐀𝐈& 𝐁\\ 𝐁& 𝐀+𝐈\end{array}\right)=0,$$
(20)
where
$$𝐀=\left(\begin{array}{cccc}_1(𝐤)& 0& g_k& 0\\ 0& _1(𝐤)& 0& g_k\\ g_k& 0& _2(𝐤)& 0\\ 0& g_k& 0& _2(𝐤)\end{array}\right)$$
(21)
$$𝐁=\left(\begin{array}{cccc}0& 0& 0& g_k\\ 0& 0& g_k& 0\\ 0& g_k& 0& 0\\ g_k& 0& 0& 0\end{array}\right)$$
(22)
and $`𝐈`$ is the identity matrix.
The densities of superfluid currents in two components can be obtained from the relation
$$𝐣_i=\frac{1}{V}\frac{F}{𝐯_i},$$
(23)
where $`F`$ is the free energy of the system. Here the quantity $`𝐣_i`$ is defined as the density of the mass current.
The free energy of the system, described by the Hamiltonian (6), is given by the formula
$$F=H_0+\frac{1}{2}\underset{𝐤}{}\left[\underset{\lambda =\alpha ,\beta }{}_\lambda (𝐤)\underset{i=1,2}{}ϵ_{ik}\right]+T\underset{𝐤}{}\underset{\lambda =\alpha ,\beta }{}\mathrm{ln}\left[1\mathrm{exp}\left(\frac{_\lambda (𝐤)}{T}\right)\right].$$
(24)
The second term in (24) is the energy of the zero-point fluctuations and the third term is the standard temperature dependent part of the free energy for the gas of noninteracting elementary excitations.
We specify the case of small superfluid velocities (much smaller than the critical ones). In this case the currents can be approximated by the expressions linear in $`𝐯_i`$. To obtain these expressions we will find the free energy as series in $`𝐯_i`$, neglecting the terms higher than quadratic.
At $`𝐯_1=𝐯_2=0`$ the equation (20) is easily solved and the spectra are found to be
$$E_{\alpha (\beta )k}=\left(\frac{E_{1k}^2+E_{2k}^2}{2}\pm \sqrt{\frac{(E_{1k}^2E_{2k}^2)^2}{4}+4\gamma _{12}^2n_1n_2ϵ_{1k}ϵ_{2k}}\right)^{1/2}.$$
(25)
As required in the procedure 13 , we take positive valued solutions of Eq. (20). The energies (25) should be real valued quantities. This requirement yields the common condition for the stability of the two-component system: $`\gamma _{12}^2\gamma _{11}\gamma _{22}`$. If this condition were not fulfilled, spatial separation of two components (at positive $`\gamma _{12})`$ or a collapse (at negative $`\gamma _{12}`$) would take place.
At nonzero superfluid velocities we present the solutions of Eq. (20) as series in $`𝐯_i`$:
$`_\alpha (𝐤)=E_{\alpha k}+{\displaystyle \frac{1}{2}}\mathrm{}\mathrm{𝐤𝐯}_1\left(1+{\displaystyle \frac{E_{1k}^2E_{2k}^2}{E_{\alpha k}^2E_{\beta k}^2}}\right)+{\displaystyle \frac{1}{2}}\mathrm{}\mathrm{𝐤𝐯}_2\left(1{\displaystyle \frac{E_{1k}^2E_{2k}^2}{E_{\alpha k}^2E_{\beta k}^2}}\right)`$ (26)
$`+{\displaystyle \frac{2\gamma _{12}^2n_1n_2ϵ_{1k}ϵ_{2k}\left(3E_{\alpha k}^2+E_{\beta k}^2\right)}{E_{\alpha k}\left(E_{\alpha k}^2E_{\beta k}^2\right)^3}}\mathrm{}^2\left(\mathrm{𝐤𝐯}_1\mathrm{𝐤𝐯}_2\right)^2,`$ (27)
$`_\beta (𝐤)=E_{\beta k}+{\displaystyle \frac{1}{2}}\mathrm{}\mathrm{𝐤𝐯}_1\left(1{\displaystyle \frac{E_{1k}^2E_{2k}^2}{E_{\alpha k}^2E_{\beta k}^2}}\right)+{\displaystyle \frac{1}{2}}\mathrm{}\mathrm{𝐤𝐯}_2\left(1+{\displaystyle \frac{E_{1k}^2E_{2k}^2}{E_{\alpha k}^2E_{\beta k}^2}}\right)`$ (28)
$`{\displaystyle \frac{2\gamma _{12}^2n_1n_2ϵ_{1k}ϵ_{2k}\left(E_{\alpha k}^2+3E_{\beta k}^2\right)}{E_{\beta k}\left(E_{\alpha k}^2E_{\beta k}^2\right)^3}}\mathrm{}^2\left(\mathrm{𝐤𝐯}_1\mathrm{𝐤𝐯}_2\right)^2.`$ (29)
Note that at $`𝐯_1=𝐯_2=𝐯`$ the spectra (26), (28) are reduced to common expressions for the energies of quasiparticles in a moving condensate: $`_{\alpha (\beta )}(𝐤)=E_{\alpha (\beta )k}+\mathrm{}\mathrm{𝐤𝐯}`$.
Using Eqs. (24),(26) and (28) we obtain the following expression for the free energy
$$F=F_0+\frac{V}{2}\left[(\rho _1\rho _{n1})𝐯_1^2+(\rho _2\rho _{n2})𝐯_2^2\rho _{\mathrm{dr}}(𝐯_1𝐯_2)^2\right],$$
(30)
where $`F_0`$ does not depend on $`𝐯_i`$. In (30) $`\rho _i=m_in_i`$ are the mass densities, the quantities
$$\rho _{n1}=\frac{m_1}{3V}\underset{𝐤}{}ϵ_{1k}\left[\frac{dN_{\alpha k}}{dE_{\alpha k}}+\frac{dN_{\beta k}}{dE_{\beta k}}+\frac{E_{1k}^2E_{2k}^2}{E_{\alpha k}^2E_{\beta k}^2}\left(\frac{dN_{\alpha k}}{dE_{\alpha k}}\frac{dN_{\beta k}}{dE_{\beta k}}\right)\right],$$
(31)
$$\rho _{n2}=\frac{m_2}{3V}\underset{𝐤}{}ϵ_{2k}\left[\frac{dN_{\alpha k}}{dE_{\alpha k}}+\frac{dN_{\beta k}}{dE_{\beta k}}\frac{E_{1k}^2E_{2k}^2}{E_{\alpha k}^2E_{\beta k}^2}\left(\frac{dN_{\alpha k}}{dE_{\alpha k}}\frac{dN_{\beta k}}{dE_{\beta k}}\right)\right]$$
(32)
describe the thermal reduction of the superfluid densities, and the quantity
$`\rho _{dr}={\displaystyle \frac{4}{3V}}\sqrt{m_1m_2}{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{\gamma _{12}^2n_1n_2\left(ϵ_{1k}ϵ_{2k}\right)^{3/2}}{E_{\alpha k}E_{\beta k}}}[{\displaystyle \frac{1+N_{\alpha k}+N_{\beta k}}{(E_{\alpha k}+E_{\beta k})^3}}{\displaystyle \frac{N_{\alpha k}N_{\beta k}}{(E_{\alpha k}E_{\beta k})^3}}`$ (33)
$`+{\displaystyle \frac{2E_{\alpha k}E_{\beta k}}{\left(E_{\alpha k}^2E_{\beta k}^2\right)^2}}({\displaystyle \frac{dN_{\alpha k}}{dE_{\alpha k}}}+{\displaystyle \frac{dN_{\beta k}}{dE_{\beta k}}})],`$ (34)
which we call the ”drag density,” yields the value of redistribution of the superfluid densities between the components. In Eqs. (31)-(33) $`N_{\alpha (\beta )k}=[\mathrm{exp}(E_{\alpha (\beta )k}/T)1]^1`$ is the Bose distribution function.
Using Eqs. (23), (30) we arrive to the following expressions for the supercurrents
$$𝐣_1=(\rho _1\rho _{n1}\rho _{dr})𝐯_1+\rho _{\mathrm{dr}}𝐯_2,$$
(35)
$$𝐣_2=(\rho _2\rho _{n2}\rho _{dr})𝐯_2+\rho _{\mathrm{dr}}𝐯_1.$$
(36)
One can see that at nonzero $`\rho _{dr}`$ the current of one component contains the term proportional to the superfluid velocity of the other component. It means that there is a transfer of motion between the components. In particular, at $`v_1=0`$ the current in the component 1 ($`𝐣_1=\rho _{dr}𝐯_2`$) is purely the drag current. Since $`\rho _{dr}`$ is the function of $`\gamma _{12}^2`$ (see Eqs. (33) and (25)) the drag current does not depend on the sign of the interaction between the components.
Eq. (33) is the main result of the paper. This equation yields the value of the drag for the general case of two-component Bose system with components of different densities, different masses of atoms, different interaction parameters, and for zero as well as for nonzero temperatures. Moreover, this equation is valid not only for the point interaction between the atoms, but for any central force interaction. In the latter case the interaction parameters $`\gamma _{ik}`$ in Eq. (33) and in the spectra (25), (14) should be replaced with the Fourier components of the corresponding interaction potentials.
To estimate the absolute value of the drag we, for simplicity, specify the case $`m_1=m_2=m`$, that is realized when two components are two hyperfine states of the same atoms.
At $`T=0`$ Eq. (33) is reduced to
$$\rho _{dr}=\frac{4m}{3}_0^{\mathrm{}}𝑑ϵ\frac{\gamma _{12}^2n_1n_2\nu (ϵ)ϵ^{1/2}}{\sqrt{(ϵ+w_1)(ϵ+w_2)}\left(\sqrt{ϵ+w_1}+\sqrt{ϵ+w_2}\right)^3},$$
(37)
where
$$\nu (ϵ)=\frac{m^{3/2}}{\sqrt{2}\pi ^2\mathrm{}^3}\sqrt{ϵ}$$
is the density of states for free atoms, and
$$w_{1(2)}=\gamma _{11}n_1+\gamma _{22}n_2\pm \sqrt{(\gamma _{11}n_1\gamma _{22}n_2)^2+4\gamma _{12}^2n_1n_2}.$$
The integral in (37) can be evaluated analytically. To present the answer in a compact form it is convenient to introduce the dimensionless parameters
$$\eta =\frac{a_{12}^2}{a_{11}a_{22}}\text{ and}\kappa =\sqrt{\frac{n_1a_{11}}{n_2a_{22}}}+\sqrt{\frac{n_2a_{22}}{n_1a_{11}}}.$$
($`0\eta 1`$ and $`\kappa 2`$)
Using these notations we have
$$\rho _{dr}=\sqrt{\rho _1\rho _2}\sqrt[4]{n_1a_{11}^3n_2a_{22}^3}\frac{\eta }{\sqrt{\kappa }}F(\kappa ,\eta ),$$
(38)
where
$`F(\kappa ,\eta )={\displaystyle \frac{256}{45\sqrt{2\pi }}}{\displaystyle \frac{(\kappa +3\sqrt{1\eta })\sqrt{\kappa }}{\left(\sqrt{\kappa +\sqrt{\kappa ^24+4\eta }}+\sqrt{\kappa \sqrt{\kappa ^24+4\eta }}\right)^3}}.`$ (39)
Direct evaluation of Eq. (39) shows that at allowed $`\eta `$ and $`\kappa `$ ($`0\eta 1`$ and $`\kappa 2`$) the factor $`F(\kappa ,\eta )`$ is almost the constant (the range of variation of $`F`$ is $`[0.7÷0.8]`$) and one can neglect the dependence of $`F`$ on the parameter of the system.
At $`a_{11}n_1=a_{22}n_2`$ we obtain from (38) the following approximate relation
$$\rho _{dr}\frac{1}{2}\rho _1\frac{a_{12}^2}{a_{11}a_{22}}\sqrt{n_1a_{11}^3}=\frac{1}{2}\rho _2\frac{a_{12}^2}{a_{11}a_{22}}\sqrt{n_2a_{22}^3}.$$
(40)
If the density of one component is much larger than of the other and $`a_{11}a_{22}`$, the ”drag density” is approximated as
$`\rho _{dr}0.8\rho _1{\displaystyle \frac{a_{12}^2}{a_{22}^2}}\sqrt{n_2a_{22}^3}\text{at}n_1n_2,`$ (41)
$`\rho _{dr}0.8\rho _2{\displaystyle \frac{a_{12}^2}{a_{11}^2}}\sqrt{n_1a_{11}^3}\text{at}n_2n_1.`$ (42)
One can see that the ”drag density” is proportional to the square root of the gas parameter. It means that the drag effect is larger in ”less ideal” Bose gases.
The temperature dependence of the ”drag density” at small $`T`$ can be evaluated analytically from Eq. (33) using the linear approximation for the spectra of the excitations. It yields $`\rho _{dr}(T)=\rho _{dr}(0)(1\alpha _TT^4/T_0^4)`$, where $`T_0=\sqrt{\gamma _{11}n_1\gamma _{22}n_2}`$ and the factor $`\alpha _T`$ is positive. Numerical evaluation of the sum over $`𝐤`$ in Eq. (33) shows that the analytical approximation is valid only at $`TT_0`$. At $`TT_0`$ the ”drag density” decreases much slower under increase of the temperature. As an example, the dependence of $`\rho _{dr}(T)`$ at $`n_1=n_2=n`$, $`\gamma _{11}=\gamma _{22}=\gamma `$ and $`\eta =0.5`$ is shown in Fig. 1.
Now let us discuss how the drag effect can reveal itself in a real physical situation. If one deals with the stationary superflow one implies that it is the circulating superflow, e.g., the tangential superflow in a hole cylinder. In such a case the superfluid velocities satisfy the Onsager-Feymnan quantization condition
$$𝐯_i𝑑𝐥=\frac{2\pi \mathrm{}N_i}{m_i},$$
(43)
where the vorticity parameters $`N_i`$ are integer. Then, the drag effect can be understood as the appearance of the circulating current in the drive component (e.g. specie 1), when the circulation of the superfluid velocity of the drive component (e.g. specie 2) is fixed ($`N_2=const`$). The current of the specie 1 (35) depends on the superfluid velocities of the both species and if the superfluid velocity of the drag component directed anti-parallel to the superfluid velocity of the drive component the current of the drag component might vanish. But since the velocities are quantized it may happen only under certain special conditions (see below). The superfluid velocity of the drag component is determined by that at fixed $`N_2`$ the free energy (30) has a minimum with respect to discrete values of $`𝐯_1=\mathrm{}N_1/(m_1R)`$ (where the $`R`$ is the radius of the contour in (43)). Depending on the value of the parameter $`\alpha =(\rho _{dr}/(\rho _2\rho _{n2}\rho _{dr}))(m_1/m_2)N_2`$ several possibilities can be realized. At $`|\alpha |<1/2`$ the minimum of the energy (30) corresponds to $`N_1=0`$ (and $`𝐯_1=0`$). In this case the current of the drag component is directed along the drive current and it is proportional to the drag density. At $`|\alpha |=p`$ ($`p`$ is natural) the value $`N_1=p`$ minimizes the energy. In this case two terms in Eq. (35) compensate each other and the current in the drag component vanishes. At half-integer $`\alpha `$ the degenerate situation takes place: two state (with co-directed currents, and counter-directed currents) have the same energy. At $`1/2+p<|\alpha |<p`$ the state with counter-directed currents gains the energy and at $`p<|\alpha |<p+1/2`$ the co-directed currents are energetically preferable. In the latter two cases the nonzero vorticity of the drag component ($`N_10`$) is also induced. This behavior is analogous the behavior of a superconducting ring in a magnetic field. We note that since $`\rho _{dr}\rho _2`$, the most realistic case is $`|\alpha |<1/2`$ when the simple picture of the transfer of part of the motion from the drive to the drag component takes place.
In this study we have concentrated on the analytical derivation of the drag effect in the uniform Bose gases. The consideration of the non-uniform case requires the solution of the eigenvalue problem for the elementary excitations in the two-component Bose gas in the external potential. But even for the simplest case of a spherically symmetric trap this problem can be solved analytically only in the long-wavelength limit and the Thomas-Fermi approximation park (the spectrum of elementary excitations in one-component Bose gases was obtained analytically for a number of potentials but also in the same limit 12 ; 14 ; 15 ; 16 ). Since the main contribution to the drag density comes from the excitations with the wave vectors of order of the healing length (see (33)), the rigorous analysis of the drag effect can be done only numerically. Nevertheless, in the Tomas-Fermi situation the drag effect can be evaluated basing on the following arguments. When the linear size of the Bose cloud is much larger than the healing length, the spectrum of the excitations at the wave vectors of order or higher than the inverse healing length is well described by the quasi uniform approximation. Therefore, the drag effect can be described by the same equations, as in the uniform case with the only modification that the quantities $`n_1`$ and $`n_2`$, and, correspondingly, $`\rho _i`$, $`\rho _{ni}`$, $`\rho _{dr}`$ and $`j_i`$ in Eqs. (31)-(36) are understood as functions of coordinates.
At an arbitrary symmetry of the trap potential the superfluid velocity of the drag component cannot be equal to zero in each point. Indeed, in general case of space dependent $`\rho _i`$, $`\rho _{ni}`$, and $`\rho _{dr}`$ the velocity field $`𝐯_2(𝐫)`$ cannot satisfy two independent continuity conditions $`[(\rho _2\rho _{n2}\rho _{dr})𝐯_2]=0`$ and $`(\rho _{dr}𝐯_2)`$. To analyze this case one should find the velocity fields $`𝐯_1(𝐫)`$ and $`𝐯_2(𝐫)`$ that satisfy the continuity conditions and the quantization conditions. To illustrate this point let us consider a simple example of a trap having the shape of a hollow cylinder with the densities that depend only on the polar angle $`\varphi `$. We will seek the velocity fields that do not have radial components. Then, the Eqs. (35), (36), written for the tangential components of the currents and the velocities, can be presented in the matrix form
$$\left(\begin{array}{c}j_1\\ j_2\end{array}\right)=\widehat{R}\left(\begin{array}{c}v_1(r,\varphi )\\ v_2(r,\varphi )\end{array}\right),$$
(44)
where
$$\widehat{R}=\left(\begin{array}{cc}\rho _{s1}(\varphi )\rho _{dr}(\varphi )& \rho _{dr}(\varphi )\\ \rho _{dr}(\varphi )& \rho _{s2}(\varphi )\rho _{dr}(\varphi )\end{array}\right)$$
(45)
with $`\rho _{si}(\varphi )=\rho _i(\varphi )\rho _{ni}(\varphi )`$. Due to the continuity conditions the current $`j_1`$ and $`j_2`$ in (44) do not depend on $`\varphi `$. According to Eq. (44) the velocities $`v_1(r,\varphi )`$ and $`v_2(r,\varphi )`$ are connected with the currents by the equation
$$\left(\begin{array}{c}v_1(r,\varphi )\\ v_2(r,\varphi )\end{array}\right)=\widehat{R}^1\left(\begin{array}{c}j_1(r)\\ j_2(r)\end{array}\right)$$
(46)
Integrating Eq. (46) over $`\varphi `$ and taking into account the quantization conditions (43) we obtain the equation for the currents
$$\widehat{T}\left(\begin{array}{c}j_1(r)\\ j_2(r)\end{array}\right)=\frac{2\pi \mathrm{}}{r}\left(\begin{array}{c}N_1/m_1\\ N_2/m_2\end{array}\right),$$
(47)
where
$$\widehat{T}=\left(\begin{array}{cc}_0^{2\pi }𝑑\varphi \frac{\rho _{s2}\rho _{dr}}{\rho _{s1}\rho _{s2}\rho _{dr}(\rho _{s1}+\rho _{s2})}& _0^{2\pi }𝑑\varphi \frac{\rho _{dr}}{\rho _{s1}\rho _{s2}\rho _{dr}(\rho _{s1}+\rho _{s2})}\\ _0^{2\pi }𝑑\varphi \frac{\rho _{dr}}{\rho _{s1}\rho _{s2}\rho _{dr}(\rho _{s1}+\rho _{s2})}& _0^{2\pi }𝑑\varphi \frac{\rho _{s1}\rho _{dr}}{\rho _{s1}\rho _{s2}\rho _{dr}(\rho _{s1}+\rho _{s2})}\end{array}\right)$$
(48)
If a given vorticity of the drive component $`N_2`$ is not very large the minimum of energy is reached at $`N_1=0`$. In the latter case the solution of Eq. (47) in the leading order in $`\rho _{dr}`$ yields the following expression for the current of the drag component
$$j_1(r)\frac{2\pi \mathrm{}N_2}{m_2r}\frac{_0^{2\pi }𝑑\varphi \frac{\rho _{dr}(\varphi )}{\rho _{s1}(\varphi )\rho _{s2}(\varphi )}}{_0^{2\pi }𝑑\varphi \frac{1}{\rho _{s1}(\varphi )}_0^{2\pi }𝑑\varphi \frac{1}{\rho _{s2}(\varphi )}}$$
(49)
One can see that if at some $`\varphi `$ the density $`\rho _{s1}`$ has a sharp minimum the first factor in denominator in Eq. (49) becomes large. On the other hand, the integral in numerator is not very sensitive to lowering of $`\rho _{s1}`$ (see Eqs. (41)). Thus, in a system with a ”bottle neck” in the drag component the drag current decreases strongly and the main consequence of the drag effect is the emergence of the gradient of the phase of the order parameter of the drag component. Similar situation takes place in a system with a weak link. The latter case is analyzed in the next section. In the uniform case Eq. (49) is reduced to $`j_1=\rho _{dr}v_2`$.
To complete the discussion we emphasize that the crossed term ($`\rho _{dr}𝐯_\mathrm{𝟏}𝐯_\mathrm{𝟐}`$) in the free energy (30) (and, consequently, the drag terms in the currents (35), (36)) comes only from the second and third terms in Eq. (24). Consequently, the drag effect considered in this paper is solely by the excitations. At the mean field level of approximation (which can be also formulated in terms of the Gross-Pitaevsky equation) the effect does not appear, while the coupling between the components is also present at that level of approximation. We would note that at the mean field level the drag effect of another type may emerge. That effect takes place in the case when one of the species is subjected by an asymmetric rotating external potential (see, for instance, loz , where such an effect has been studied with reference to the system of two coupled traps).
## III Model of Bose-Einstein qubit with external drag force
It is known that Bose systems in the Bose-Einstein condensed state may demonstrate Josephson phenomenon leg . It this paper we consider the external Josephson effect that takes place in two-well Bose systems. It was shown in 10 that in such systems one can realize the situation, when two states, that differ in the expectation value of the relative number operator, can be used as qubit states.
To include the drag force into the play we consider the following geometry. Let our two-component system is confined in a toroidal trap and the Bose cloud of the component 1 (the drag component) is situated inside and overlaps with the Bose clouds of the component 2 (the drive component). Such a situation can be realized if $`|\gamma _{12}|<\mathrm{min}(\gamma _{11},\gamma _{22})`$.
Deforming the confining potential one can cut the drag component into two clouds of a half-torus shape (separated by two Josephson links) leaving the Bose cloud of the drive component uncutted ( Fig. 2). In what follows we use the following notations: $`R_t`$ is the large radius of the toroidal trap, $`r_{t1}`$ and $`r_{t2}`$ are the small radiuses of the toroidal Bose clouds of the drag and the drive components, correspondingly.
Rotating this trap one can excite a superflow in the drive component. After the rotation be switched off there will be a circulating superflow in the drive component and no superflow in the drag component (at negligible small Josephson coupling). The superfluid velocity of the drive component is
$$v_2=\frac{N_2\mathrm{}}{m_2R_t}$$
(50)
In (50), we imply that $`R_tr_{t1},r_{t2}`$ and neglect, for simplicity, the effect caused by a dependence of $`r_{t2}`$ on the polar angle.
Since $`j_1=0`$, the phase gradient $`\phi _1`$ should be nonzero to compensate the drag effect. In the polar coordinates the $`\varphi `$ component of the phase gradient is given by the relation
$$(\phi _1)_\varphi =\frac{N_2}{R_t}f_{dr}=f_{dr}(\phi _2)_\varphi ,$$
(51)
where
$$f_{dr}=\frac{m_1}{m_2}\frac{\rho _{dr}}{\rho _{s1}\rho _{dr}}$$
(52)
The quantity $`f_{dr}`$ yields the ratio between the phase gradients in the drag and the drive components in the situation when the drag component is in the open circuit (i.e. the current cannot flow in the circuit). We call this quantity the drag factor.
We imply that $`r_{t1}`$ and $`r_{t2}`$ are much larger than the healing lengths that allows to describe the drag effect in quasi-uniform approximation. For definiteness, we specify the case of $`\rho _1\rho _2`$ and $`\rho _2const`$ in the overlapping region. In this case one can neglect the space dependence the drag factor (see Eqs. (41)).
At nonzero Josephson coupling the current $`j_1`$ can be nonzero, but it cannot exceed the maximum Josephson current $`j_m`$. Relation (51) remains approximately correct at nonzero Josephson coupling, if an inequality $`j_m\mathrm{}\rho _1/(m_1R_t)`$ is satisfied. Here we specify just such a case. It is important to emphasize that we consider the situation, when there is only the external Josephson effect between two half-torus traps, and there is no internal Josephson effect between the drag and the drive species.
The drag force can be considered as an effective vector potential $`𝐀_{dr}=\mathrm{}f_{dr}\phi _2`$ (in units of $`e=c=1`$) that corresponds to an effective magnetic flux $`\mathrm{\Phi }_{dr}=2\pi \mathrm{}f_{dr}N_2`$. Thus, our Bose system is similar to the Cooper pair box system that implements the Josephson charge qubit with the Josephson coupling controlled by an external magnetic flux schon . To extend this analogy we formulate the model of the Bose-Einstein qubit subjected by the drag force. In what follows we use the approach of Ref. 10 .
In the two mode approximation the Bose field operators for the drag component can be presented in the form:
$$\widehat{\mathrm{\Psi }}_1(𝐫,t)=\underset{l=L,R}{}a_l(t)\mathrm{\Psi }_l(𝐫𝐫_l),\widehat{\mathrm{\Psi }}_1^+(𝐫,t)=\underset{l=L,R}{}a_l^+(t)\mathrm{\Psi }_l^{}(𝐫𝐫_l),$$
(53)
where $`a_{L(R)}^+`$ and $`a_{L(R)}`$ are the operators of creation and annihilation of bosons in the condensates confined in the left(right) half-torus, and $`\mathrm{\Psi }_L`$, $`\mathrm{\Psi }_R`$ are two almost orthogonal local mode functions
$$d^3r\mathrm{\Psi }_l^{}(𝐫)\mathrm{\Psi }_l^{}(𝐫)\delta _{ll^{}},l,l^{}=L,R$$
that describe the condensate in the left and right traps rag .
Substituting (53) into Hamiltonian (1), we obtain the following expression for the parts of the Hamiltonian that depends on the operators $`a_l^+`$ and $`a_l`$:
$$H_a=\underset{l=L,R}{}\left(K_la_l^+a_l+\lambda _la_l^+a_l^+a_la_l\right)+(Ja_L^+a_R+J^{}a_R^+a_L).$$
(54)
with
$$K_l=d^3r\mathrm{\Psi }_l^{}\left[\frac{\mathrm{}^2}{2m}^2+V_{tr}+\gamma _{12}\mathrm{\Psi }_2^{}\mathrm{\Psi }_2\right]\mathrm{\Psi }_l,$$
(55)
$$\lambda _l=\frac{\gamma _{11}}{2}d^3r|\mathrm{\Psi }_l|^4,$$
(56)
$$J=d^3r\left[\frac{\mathrm{}^2}{2m}\mathrm{\Psi }_L^{}\mathrm{\Psi }_R+V_{tr}\mathrm{\Psi }_L^{}\mathrm{\Psi }_R\right].$$
(57)
The functions $`\mathrm{\Psi }_L`$ and $`\mathrm{\Psi }_R`$ contain the phase factors $`e^{i\phi _L(𝐫)}`$ and $`e^{i\phi _R(𝐫)}`$, where the phases satisfy Eq. (51). Taking these factors into account, one can choose the following basis for the one mode functions
$$\mathrm{\Psi }_{L(R)}(𝐫)=|\mathrm{\Psi }_{L(R)}(𝐫)|\mathrm{exp}\left[iN_2f_{dr}\varphi _{L(R)}(𝐫)\right],$$
(58)
where $`\varphi _L`$,$`\varphi _R`$ are the polar angles counted from the centers of $`L`$ and $`R`$ half-torus, correspondingly (see Fig. 3). The angles $`\varphi _{L(R)}(𝐫)`$, defined as shown in Fig. 3, satisfy the relation
$$\varphi _R(𝐫_A)\varphi _L(𝐫_A)=\varphi _L(𝐫_B)\varphi _R(𝐫_B)=\pi ,$$
(59)
where $`𝐫_A`$ and $`𝐫_B`$ are the radius-vectors of Josephson links.
Substituting (58) into Eq. (57), using Eq. (59) and taking into account that the functions $`\mathrm{\Psi }_L`$ and $`\mathrm{\Psi }_R`$ overlap in a small vicinity of A and B links, we obtain the following expression for the Josephson coupling parameter:
$$J=(J_A+J_B)\mathrm{cos}\left(\pi \frac{\mathrm{\Phi }_{dr}}{\mathrm{\Phi }_0}\right)+i(J_AJ_B)\mathrm{sin}\left(\pi \frac{\mathrm{\Phi }_{dr}}{\mathrm{\Phi }_0}\right),$$
(60)
where $`\mathrm{\Phi }_0=2\pi \mathrm{}`$ is the ”flux quantum” and
$$J_{A(B)}_{V_{A(B)}}d^3r\left[\frac{\mathrm{}^2}{2m}|\mathrm{\Psi }_L||\mathrm{\Psi }_R|+V_{tr}|\mathrm{\Psi }_L||\mathrm{\Psi }_R|\right].$$
(61)
Here $`V_A`$ and $`V_B`$ are the areas of overlapping of two one mode functions at links $`A`$ and $`B`$, correspondingly.
Considering the Hilbert space in which the total number operator
$$\widehat{N}=a_L^+a_L+a_R^+a_R$$
(62)
is a conservative quantity ($`\widehat{N}=N`$) we present the Hamiltonian (54) in the following form
$$H_a=E_c(\widehat{n}_{RL}n_g)^2+(Ja_L^+a_R+h.c.)+const,$$
(63)
where
$$\widehat{n}_{RL}=\frac{a_R^+a_Ra_L^+a_L}{2}$$
(64)
is the number difference operator,
$$E_c=\lambda _R+\lambda _L$$
(65)
is the interaction energy, and the quantity
$$n_g=\frac{1}{2E_c}\left[K_LK_R+(N1)(\lambda _L\lambda _R)\right]$$
(66)
describes an asymmetry of L and R half-tore.
In what follows we imply that the system is in the Fock regime leg ($`|J|NE_c`$) and use the number representation
$$|n_{RL}|n_R,n_L|\frac{N}{2}+n_{RL},\frac{N}{2}n_{RL}.$$
In this representation the first term in (63) is diagonal. The second term in (63) can be considered as small nondiagonal correction. But if $`n_g`$ is biased near one of degeneracy points
$$n_{deg}=\{\begin{array}{cc}M+\frac{1}{2}\hfill & \text{for even }N\hfill \\ M\hfill & \text{for odd }N\hfill \end{array}$$
(67)
(where $`M`$ is integer and $`|M|<N/2`$), the second term in (63) results in a strong mixing of two lowest states ($`|=|n_{deg}+1/2`$ and $`|=|n_{deg}1/2`$) and the low energy dynamics of the system can be described by a pseudospin Hamiltonian
$$H_{eff}=\frac{\mathrm{\Omega }_x}{2}\widehat{\sigma }_x\frac{\mathrm{\Omega }_y}{2}\widehat{\sigma }_y\frac{\mathrm{\Omega }_z}{2}\widehat{\sigma }_z,$$
(68)
where $`\widehat{\sigma }_i`$ are the Pauli operators, and
$`\mathrm{\Omega }_x=(J_A+J_B)\sqrt{(N+1)^24n_{deg}^2}\mathrm{cos}\left(\pi {\displaystyle \frac{\mathrm{\Phi }_{dr}}{\mathrm{\Phi }_0}}\right),`$ (69)
$`\mathrm{\Omega }_y=(J_AJ_B)\sqrt{(N+1)^24n_{deg}^2}\mathrm{sin}\left(\pi {\displaystyle \frac{\mathrm{\Phi }_{dr}}{\mathrm{\Phi }_0}}\right),`$ (70)
$`\mathrm{\Omega }_z=2E_c(n_gn_{deg})`$ (71)
are the components of the pseudomagnetic field. In experiments one can control the parameters $`n_g`$, $`J_A`$ and $`J_B`$ independently and, consequently, the pseudomagnetic field $`𝛀(t)`$ can be switched arbitrary. It mean that Eq. (68) represents the standard Hamiltonian of the qubit system. The parameters of the qubit (68) depend on the ”drag flux” $`\mathrm{\Phi }_{dr}`$. Therefore, one can determine its value from the measurement of the state of the system after a controlled evolution of a certain reproducible initial state.
Let us consider two possibilities. For definiteness, we specify the case of odd $`N`$ and the degeneracy point $`n_{deg}=0`$.
If the Josephson coupling are switched off and $`n_g`$ is switched on to some positive value (much less than unity) the system is relaxed to the state $`|\psi _{in}=|`$. This state can be used as the reproducible initial state. The quantity should be measured is the expectation value of the number difference operator. In the initial state the expectation value of this operator is $`n_{RL}=1/2`$
When the system is switched suddenly to the degeneracy point $`n_g=0`$ and the Josephson couplings are switched on for some time $`\tau `$ the initial state evolves to another state with another $`n_{RL}`$.
If one sets $`J_A=J_B=J`$ the result of evolution ($`|\psi _f=U|\psi _{in}`$) is described by the unitary operator
$$U_1(\tau )=\left(\begin{array}{cc}\mathrm{cos}(\alpha _1\tau )& i\mathrm{sin}(\alpha _1\tau )\\ i\mathrm{sin}(\alpha _1\tau )& \mathrm{cos}(\alpha _1\tau )\end{array}\right)$$
where $`\alpha _1=(J/\mathrm{})(N+1)\mathrm{cos}(\pi \mathrm{\Phi }_{dr}/\mathrm{\Phi }_0)`$. One can see that at time of evolition $`\tau =\tau _1=\pi /(4|\alpha _1|)`$ the expectation value of the number difference operation will be equal to zero.
For the case $`J_A=J`$ and $`J_B=0`$ the operator of evolution reads as
$$U_2(\tau )=\left(\begin{array}{cc}\mathrm{cos}(\alpha _2\tau )& ie^{i\pi \mathrm{\Phi }_{dr}/\mathrm{\Phi }_0}\mathrm{sin}(\alpha _2\tau )\\ ie^{i\pi \mathrm{\Phi }_{dr}/\mathrm{\Phi }_0}\mathrm{sin}(\alpha _2\tau )& \mathrm{cos}(\alpha _2\tau )\end{array}\right)$$
with $`\alpha _2=(J/2\mathrm{})(N+1)`$. Respectively, the expectation value $`n_{RL}`$ will be equal to zero at $`\tau =\tau _2=\pi /(4\alpha _2)`$
The ratio $`\tau _2/\tau _1=|\mathrm{cos}(\pi \mathrm{\Phi }_{dr}/\mathrm{\Phi }_0)|/2`$ depends only on $`\mathrm{\Phi }_{dr}`$ and the quantity $`\mathrm{\Phi }_{dr}`$ can be extracted from the measurements of $`\tau _1`$ and $`\tau _2`$. It is important to note that to provide this scheme one should control only the ratio of $`J_A`$ and $`J_B`$, but not their absolute values.
Another possibility can be based on detection of the Berry phase ber . Eq. (68) contains all three components of the field $`𝛀`$ and they can be controlled independently. The general scheme of detection of the Berry phase in such a situation was proposed 8 . A concrete realization of this scheme in the Josephson charge qubit was described in 9 . Here we extend the ideas of 8 ; 9 to the case of the ”dragged” Bose-Einstein qubit.
We start from the same initial state and switch to $`J_A=J_B=J`$ and $`n_g=0`$. The initial state $`|`$ can be presented as the superposition of two instantaneous eigenstates $`|e_a=(|+|)/\sqrt{2}`$ and $`|e_b=(||)/\sqrt{2}`$:
$$|\psi _{in}=\frac{1}{\sqrt{2}}(|e_a+|e_b).$$
(72)
An adiabatic cyclic evolution of the parameters of the Hamiltonian (68) results in appearance of the Berry phase in the $`|e_a`$ and $`|e_b`$ eigenstates, if the vector $`𝛀`$ subtends a nonzero solid angle at the origin.
Let us consider the following 4 stage cyclic adiabatic evolution starting from the point $`J_A=J_B=J`$ and $`n_g=0`$: 1 - $`J_B`$ is switched off; 2 - $`J_A`$ is switched off and simultaneously $`n_g`$ is switched to $`n_{g1}>0`$; 3 - $`n_g`$ is returned to the same degeneracy point ($`n_g=0`$) and $`J_B`$ is switched to $`J_B=J`$; 4 - $`J_A`$ is switched to $`J_A=J`$ (all switches should be done slowly: $`\mathrm{}|d𝛀/dt|\mathrm{\Omega }^2`$).
After such an evolution the system arrives at the state
$$|\psi _m=\frac{1}{\sqrt{2}}\left(e^{i\delta _a+i\gamma }|e_a+e^{i\delta _bi\gamma }|e_b\right),$$
(73)
where $`\gamma =\pi \mathrm{\Phi }_{dr}/\mathrm{\Phi }_0`$ is the Berry phase (equals to half of the solid angle subtended by $`𝛀`$) and $`\delta _a`$, $`\delta _b`$ are the dynamical phases.
Elimination of the dynamical phases can performed by swapping the eigenstates ($`\pi `$-transformation) and repeating the same cycle of evolution in a reverse direction (see 8 ).
The $`\pi `$-transformation can be done by fast switching off the Josephson coupling and switching on $`n_g=n_{g2}>0`$ during the time interval $`t_\pi =\mathrm{}\pi /(2E_cn_{g2})`$. After the $`\pi `$-transformation the state becomes
$$|\psi _{m\pi }=\frac{i}{\sqrt{2}}(e^{i\delta _a+i\gamma }|e_b+e^{i\delta _bi\gamma }|e_a).$$
(74)
After the cyclic evolution in the reverse direction we arrive to the state
$$|\psi _f=\frac{i}{\sqrt{2}}e^{i(\delta _a+\delta _b)}(e^{2i\gamma }|e_b+e^{2i\gamma }|e_a).$$
(75)
One can see that the expectation value of the number difference operator in the final state (75) $`n_{RL}=\mathrm{cos}(4\gamma )/2=\mathrm{cos}(4\pi \mathrm{\Phi }_{dr}/\mathrm{\Phi }_0)/2`$ depends only on $`\mathrm{\Phi }_{dr}`$ and the measurement of this difference allows the determine the value of the ”drag flux”.
Thus, the measurements of relative number of atoms in left and right condensates under controlled evolution of the state of the system allows to observe the non-dissipative drag and determine the drag factor (if the vorticity of the drive component is known).
## IV Conclusions
We have investigated the non-dissipative drag effect in three-dimensional weakly interacting two-component superfluid Bose gases. The expression for the drag current is derived microscopically for the general case of two species of different densities, different masses and different interaction parameters. It is shown that the drag current is proportional to the square root of the gas parameter. The drag effect is maximal at zero temperatures and it decreases when the temperature increases, but at temperatures of order of the interaction energy the drag current remains of the same order as at zero temperature.
We have considered the toroidal double-well geometry, where the non-dissipative drag influences significantly on the Josephson coupling between the wells. In the system considered the drag force can be interpreted as an effective vector potential applied to the drag component. The effective vector potential is equal to $`𝐀_{dr}=\mathrm{}f_{dr}\phi _{drv}`$ (in units of $`e=c=1`$), where $`\phi _{drv}`$ is the phase of the drive component, and $`f_{dr}`$ is the drag factor. In the toroidal geometry the effective vector potential can be associated with an effective flux of external field $`\mathrm{\Phi }_{dr}=2\pi \mathrm{}f_{dr}N_v`$, where $`N_v`$ is the vorticity of the drive component. In the Fock regime the system can be considered as a Bose-Einstein counterpart of the Josephson charge qubit in an external magnetic field. The measurement of the state of such a qubit allows to observe the drag effect and determine the drag factor.
## Acknowledgements
This work is supported by the INTAS grant No 01-2344.
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# Galleries, Hall-Littlewood polynomials and structure constants of the spherical Hecke algebra
## 1 Introduction
Let $`\mathrm{\Phi }=(X,\varphi ,X^{},\varphi ^{})`$ be a reduced root datum with Weyl group $`W`$. Let $`V=X`$ and $`V^{}X^{}`$ such that the pairing $`,:V\times V^{}`$ is induced from the perfect pairing between $`X`$ and $`X^{}`$, which will be denoted the same way. Let $`Q^{}`$ be the coroot lattice. Choose simple roots $`\mathrm{\Delta }=\{\alpha _1,\mathrm{},\alpha _l\}`$ and denote by $`\varphi ^+`$ the positive roots with respect to $`\mathrm{\Delta }`$. Let $`X_+^{}=\{xX^{}|\alpha ,x0\text{ for all }\alpha \varphi ^+\}`$ be the dominant cone. There is a natural action of $`W`$ on the group algebra $`[X^{}]`$. The algebra of symmetric polynomials $`\mathrm{\Lambda }`$ is the algebra of invariants under this action. It is closely related to the representation theory of complex algebraic groups.
Let $`G^{}`$ be a complex reductive linear algebraic group with Borel subgroup $`B^{}`$ and maximal torus $`T^{}`$ such that the associated root datum is the dual of $`\mathrm{\Phi }`$. Assigning to a finite dimensional representation of $`G^{}`$ its $`T^{}`$-character yields an isomorphism from the representation ring of $`G^{}`$ to $`\mathrm{\Lambda }`$. For $`\lambda X_+^{}`$ the Schur polynomial $`s_\lambda `$ is the character of the irreducible highest weight module $`V(\lambda )`$ with highest weight $`\lambda `$. The Schur polynomials $`\{s_\lambda \}_{\lambda X_+^{}}`$ are a basis of $`\mathrm{\Lambda }`$. The Kostka number $`k_{\lambda \mu }`$ for $`\lambda ,\mu X_+^{}`$ is the weight multiplicity of $`\mu `$ in $`V(\lambda )`$, i.e. the dimension of the $`\mu `$-weight space $`V(\lambda )_\mu `$. The Littlewood-Richardson coefficient $`c_{\lambda \mu }^\nu `$ for $`\lambda ,\mu ,\nu X_+^{}`$ is the multiplicity of $`V(\nu )`$ in $`V(\lambda )V(\mu )`$.
Combinatorially the $`k_{\lambda \mu }`$ and the $`c_{\lambda \mu }^\nu `$ can be defined as follows: For $`\lambda X_+^{}`$ define the monomial symmetric function $`m_\lambda =_{\mu W\lambda }x^\mu `$ where for $`\mu X^{}`$ we denote the corresponding basis element of $`[X^{}]`$ by $`x^\mu `$. Then $`\{m_\lambda \}_{\lambda X_+^{}}`$ is a basis of $`\mathrm{\Lambda }`$ and the Kostka numbers are the entries of the transition matrix from the $`m_\mu `$ to the $`s_\lambda `$, i.e. we have $`s_\lambda =_{\mu X_+^{}}k_{\lambda \mu }m_\mu `$. The Littlewood-Richardson coefficients are the structure constants of $`\mathrm{\Lambda }`$ with respect to the Schur polynomials, i.e. $`s_\lambda s_\mu =_{\nu X_+^{}}c_{\lambda \mu }^\nu s_\nu `$.
One of the main problems of combinatorial representation theory was to give combinatorial formulas for the $`k_{\lambda \mu }`$ and the $`c_{\lambda \mu }^\nu `$. One of the solutions is Littelmann’s path model in \[Lit94\]. In \[GL05\] Gaussent and Littelmann introduced the gallery model and showed that it is equivalent to the path model. They express the $`k_{\lambda \mu }`$ and the $`c_{\lambda \mu }^\nu `$ as the number of certain galleries.
In this paper we give analogs of these combinatorial formulas for some $`q`$-analog of Schur polynomials, the so called Hall-Littlewood polynomials. We describe the transition matrix from monomial symmetric functions to Hall-Littlewood polynomials and we calculate products of Hall-Littlewood polynomials. Specializing $`q`$ in these formulas we get a new proof for the above mentioned formulas for Schur polynomials in terms of galleries.
Extending the base ring to $`^{}:=[q^1]`$ one gets new interesting bases. The Hall-Littlewood polynomials $`\{P_\lambda (q^1)\}`$ are a basis for $`\mathrm{\Lambda }_q:=^{}[X^{}]^W`$ (as $`^{}`$-module). For $`\lambda X_+^{}`$ they are defined by
$$P_\lambda (q^1)=\frac{1}{W_\lambda (q^1)}\underset{wW}{}w\left(x^\lambda \underset{\alpha \varphi ^+}{}\frac{1q^1x^\alpha ^{}}{1x^\alpha ^{}}\right).$$
Here $`W_\lambda W`$ is the stabilizer of $`\lambda `$ and $`W_\lambda (q^1)=_{wW_\lambda }q^{l(w)}`$ where $`l:W`$ is the usual length function. The Hall-Littlewood polynomials are $`q`$-analogs of the Schur polynomials in the sense that $`P_\lambda (0)=s_\lambda `$. Moreover, we have $`P_\lambda (1)=m_\lambda `$. For these and other properties of the $`P_\lambda (q^1)`$ see the article \[NR03\] of Nelsen and Ram.
Define Laurent polynomials $`L_{\lambda \mu }`$ for $`\lambda ,\mu X_+^{}`$ by
$$P_\lambda (q^1)=\underset{\mu X_+^{}}{}q^{\rho ,\lambda +\mu }L_{\lambda \mu }m_\mu ,$$
where $`\rho :=\frac{1}{2}_{\alpha \varphi ^+}\alpha `$. By definition we have $`q^{\rho ,\lambda +\mu }L_{\lambda \mu }^{}`$. Moreover, since $`P_\lambda (0)=s_\lambda `$, the constant term of $`q^{\rho ,\lambda +\mu }L_{\lambda \mu }`$ is $`k_{\lambda \mu }`$. So a combinatorial description of the $`L_{\lambda \mu }`$ yields a combinatorial description of the $`k_{\lambda \mu }`$. For non-dominant $`\mu X^{}`$ we define $`L_{\lambda \mu }=q^{\rho ,\mu \mu ^+}L_{\lambda \mu ^+}`$, where $`\mu ^+X_+^{}`$ is the unique dominant element in the $`W`$-orbit of $`\mu `$.
The main combinatorial tool for the description of the $`L_{\lambda \mu }`$ are galleries of generalized alcoves of a fixed type. For details on this and other unexplained notation see section 4. Following \[GL05\] we introduce positively folded galleries and associate to each positively folded gallery $`\sigma `$ a combinatorially defined polynomial $`L_\sigma `$ in definition 4.2. In section 5 we prove (see also theorem 4.4 and the paragraph before it)
###### Theorem 1.1.
Let $`\lambda X_+^{}`$ and $`\mu X^{}`$. Denote by $`t^\lambda `$ the type of a minimal gallery from 0 to $`\lambda `$. Then $`L_{\lambda \mu }=q^{l(w_\lambda )}_\sigma q^{l(w_0\iota (\sigma ))}L_\sigma `$. Here $`w_\lambda W_\lambda `$ is the maximal element and the sum is over all positively folded galleries $`\sigma `$ of type $`t^\lambda `$ and weight $`\mu `$ such that the initial direction $`\iota (\sigma )`$ is in the set of minimal representatives $`W^\lambda `$ of $`W/W_\lambda `$.
If $`\mathrm{\Phi }`$ is of type $`A`$ one gets a combinatorial description of the $`L_{\lambda \mu }`$ with Young diagrams as a corollary of the results of Haglund, Haiman and Loehr on the monomial expansion of Macdonald polynomials \[HHL05\].
We get a description of the $`k_{\lambda \mu }`$ in terms of galleries by evaluation at $`q^1=0`$. We introduce LS-galleries in definition 4.7 (roughly speaking these are the galleries which survive the specialization $`q^1=0`$) and get
###### Corollary 1.2.
For $`\lambda ,\mu X_+^{}`$ the Kostka number $`k_{\lambda \mu }`$ is the number of LS-galleries of type $`t^\lambda `$ and weight $`\mu `$.
In definition 4.2 we also introduce a second monic polynomial $`C_\sigma `$ for each gallery $`\sigma `$ which is closely related to $`L_\sigma `$. We prove that with this statistic one can calculate the structure constants of $`\mathrm{\Lambda }_q`$ with respect to the Hall-Littlewood polynomials. More precisely, define $`C_{\lambda \mu }^\nu `$ for $`\lambda ,\mu ,\nu X_+^{}`$ by
$$P_\lambda (q^1)P_\mu (q^1)=\underset{\nu X_+^{}}{}q^{\rho ,\mu \lambda +\nu }C_{\lambda \mu }^\nu P_\nu (q^1).$$
Let $`𝒞:=\{xV^{}|\alpha ,x0\text{ for all }\alpha \varphi ^+\}`$ be the dominant Weyl chamber. We prove in section 6 (see also theorem 4.10 and the paragraph before it)
###### Theorem 1.3.
Let $`\lambda ,\mu ,\nu X_+^{}`$. Then $`C_{\lambda \mu }^\nu =q^{l(w_\mu )}_\sigma q^{l(w_0\iota (\sigma ))}C_\sigma F_{\mu \nu }^{\epsilon (\sigma )}`$. Here the sum is over all positively folded galleries $`\sigma `$ of type $`t^\mu `$ and weight $`\nu `$ starting in $`\lambda `$ such that they are contained in $`𝒞`$ and the final direction $`\epsilon (\sigma )`$ is in $`W_\nu W^{w_0\mu }`$. The correction factor $`F_{\mu \nu }^{\epsilon (\sigma )}`$ is contained in $`^{}`$.
In \[KM04\] a similar formula for $`C_{\lambda \mu }^\nu `$ with some restrictions on $`\lambda ,\mu ,\nu `$ was given by Kapovich and Millson using geodesic triangles in some euclidean building associated to the situation.
For $`q^1=0`$ theorem 1.3 yields a Littlewood-Richardson rule in terms of galleries.
###### Corollary 1.4.
For $`\lambda ,\mu ,\nu `$ in $`X_+^{}`$ the Littlewood-Richardson coefficient $`c_{\lambda \mu }^\nu `$ is the number of LS-galleries $`\sigma `$ of type $`t^\mu `$ and weight $`\nu \lambda `$ such that the translated gallery $`\lambda +\sigma `$ is contained in $`𝒞`$.
The combinatorial descriptions in the corollaries 1.2 and 1.4 are more or less the same as the above mentioned descriptions in \[GL05\]. But our proof is quite different and does not use the combinatorial results stated there. See remark 4.9 for some further details.
For proving theorems 1.1 and 1.3 we use the Satake isomorphism to identify $`[q^{\pm {\scriptscriptstyle \frac{1}{2}}}][X^{}]^W`$ with the spherical Hecke algebra with equal parameters associated to $`\mathrm{\Phi }`$. Under this isomorphism, Hall-Littlewood polynomials correspond (up to some factor) to the Macdonald basis and the monomial symmetric functions correspond to the monomial basis of the spherical Hecke algebra (see remark 3.2). So the above theorems can be proven in the spherical Hecke algebra.
This is done in the slightly more general setting of spherical Hecke algebras with arbitrary parameters. In section 4 we state our combinatorial formulas regarding Satake coefficients, i.e. the entries of the transition matrix from the monomial basis to the Macdonald basis, and the structure constants of the spherical Hecke algebra with respect to the Macdonald basis. The formulas are then proven in sections 5 and 6. We define a basis of the affine Hecke algebra indexed by the generalized alcoves introduced in section 2 and show that right multiplication of this alcove basis by elements of the standard basis can be calculated using galleries.
It is well known that the Satake coefficients form a triangular matrix with respect to the usual partial order on $`X^{}`$, i.e. $`\mu \lambda `$ iff $`\lambda \mu =_{\alpha \mathrm{\Delta }}n_\alpha \alpha ^{}`$ for nonnegative integers $`n_\alpha `$. With our combinatorial description of the Satake coefficients we can show that all remaining entries are in fact nonzero.
Now take the spherical Hecke algebra with coefficients in $``$ such that all parameters are powers of a fixed prime $`p`$. Our description of the $`L_{\lambda \mu }`$ implies that $`L_{\lambda \mu }>0`$ for $`\lambda ,\mu X_+^{}`$ with $`\mu \lambda `$. This gives a combinatorial proof of a positivity result of Rapoport in \[Rap00\].
In section 7 we compute the transition matrix between the alcove basis and the standard basis. This yields a $`q`$-analog of a commutation formula of Pittie and Ram (\[PR99\]) in terms of galleries.
In section 8 we give a geometric interpretation of our combinatorics in the case $`X^{}=Q^{}`$. Fix a prime power $`𝐪`$ and let $`𝔽_𝐪`$ be the finite field with $`𝐪`$ elements. Let $`K`$ be the algebraic closure of $`𝔽_𝐪`$. Let $`G`$ be a semisimple, simply connected algebraic group over $`K`$ with root datum $`\mathrm{\Phi }`$ corresponding to some choice of a Borel $`BG`$ and a maximal torus $`TB`$. Assume that all groups are defined and split over $`𝔽_𝐪`$.
From the definition of the geometric Satake isomorphism (see for example the survey article \[HKP03\] of Haines, Kottwitz and Prassad) it is known, that the evaluation at $`𝐪`$ of the coefficients $`L_{\lambda \mu }`$ gives the number of points of certain intersections in the affine Grassmanian of $`G`$ over $`𝔽_𝐪`$. We show that our combinatorics reflects this interpretation. Using results of Billig and Dyer in \[BD94\] we show that the galleries occurring in theorem 1.1 together with the associated coefficients parameterize decompositions of these intersections. We also give an interpretation of the alcove basis in this context.
*Acknowledgments.* The author would like to thank P. Littelmann and A. Ram for various helpful suggestions and discussions.
## 2 Affine Weyl group and alcoves
In this section we recall some facts on the (extended) affine Weyl group and on alcoves as in \[Bou81\]. Furthermore, we introduce generalized alcoves.
The group $`Q^{}`$ acts on $`V^{}`$ by translations. The affine Weyl group is defined as the semidirect product $`W^𝔞=WQ^{}`$. It acts on $`V^{}`$ by affine transformations. For $`\lambda Q^{}`$ denote by $`\tau _\lambda W^𝔞`$ the associated translation. The affine Weyl group is generated by its affine reflections. Let $`H^𝔞`$ be the union of all reflection hyperplanes of reflections in $`W^𝔞`$. Then $`H^𝔞=_{\alpha \varphi ^+,m}H_{\alpha ,m}`$, where $`H_{\alpha ,m}=\{xV^{}|\alpha ,x=m\}`$. Let $`H_{\alpha ,m}^\pm =\{xV^{}|\alpha ,xm\}`$ be the associated affine half spaces.
The connected components of $`V^{}H^𝔞`$ are called open alcoves. Their closures are the alcoves in $`V^{}`$. Denote by $`𝒜`$ the set of all alcoves. The action of $`W^𝔞`$ on $`𝒜`$ is free and transitive. For $`A𝒜`$ and $`\lambda Q^{}`$ we have $`\tau _\lambda A=\lambda +A=\{\lambda +x|xA\}`$.
The fundamental alcove $`A_f=\{xV^{}|0\alpha ,x1\text{ for all }\alpha \varphi ^+\}𝒜`$ is a fundamental domain for the $`W^𝔞`$-action on $`V^{}`$. We get a bijection $`W^𝔞𝒜,wA_w:=wA_f`$.
A face $`F`$ of an alcove $`A`$ is an intersection $`F=AH`$ such that $`HH^𝔞`$ is a reflection hyperplane and $`F_{\text{aff}}=H`$. Here $`F_{\text{aff}}`$ is the affine subspace spanned by $`F`$. A wall of $`A`$ is some hyperplane $`HH^𝔞`$ such that $`HA`$ is a face of $`A`$. The group $`W^𝔞`$ is generated by the reflections $`S^𝔞`$ at the walls of $`A_f`$. One has $`S^𝔞=S\{s_{01},\mathrm{},s_{0c}\}`$, where $`S=\{s_1,\mathrm{},s_l\}`$ is the set of simple reflections of $`W`$ and $`s_{0k}`$ is the affine reflection at $`H_{\theta _k,1}`$. Here $`\{\theta _k|k=1,\mathrm{},c\}\varphi ^+`$ is the set of maximal roots with respect to the partial order on $`X^{}`$, so $`c`$ is the number of irreducible components of the Dynkin diagram of $`\mathrm{\Phi }`$. Moreover, $`(W^𝔞,S^𝔞)`$ is a Coxeter system.
Let $`F`$ be a face of $`A_f`$. The type of $`F`$ is the reflection at $`F_{\text{aff}}`$. Extend this definition to all faces by demanding that the $`W^𝔞`$-action preserves types.
Right multiplication of $`W^𝔞`$ induces an action of $`W^𝔞`$ on $`𝒜`$ from the right. For $`A𝒜`$ and $`sS^𝔞`$ the alcove $`As`$ is the unique alcove not equal to $`A`$ having a common face of type $`s`$ with $`A`$. Let $`F_sA`$ be the face of type $`s`$ and $`F_s_{\text{aff}}=H_{\alpha ,m}`$ for some $`\alpha \varphi ^+`$ and $`m`$. The hyperplane $`H_{\alpha ,m}`$ is called the separating hyperplane between $`A`$ and $`As`$. Call $`A`$ negative with respect to $`s`$ if $`A`$ is contained in $`H_{\alpha ,m}^{}`$ and denote this by $`AAs`$. Of course $`A`$ is called positive with respect to $`s`$ if $`As`$ is negative with respect to $`s`$. We have $`AAs`$ iff $`\lambda +A\lambda +As`$ for all $`\lambda Q^{}`$.
###### Example 2.1.
* For $`A_w`$ and $`A_{ws}`$ in the dominant chamber $`𝒞`$ we have $`A_wA_{ws}`$ iff $`w<ws`$, where ’$``$’ is the usual Bruhat order on $`W^𝔞`$.
* Let $`wW`$ and $`sS`$. Then $`A_wA_{ws}`$ iff $`w>ws`$.
There is also a natural action of $`X^{}`$ on $`V^{}`$ by translations. So we can extend the above definition and get the extended affine Weyl group $`\stackrel{~}{W}^𝔞:=WX^{}`$. Extending the above notation write $`\tau _\mu `$ for the translation by $`\mu X^{}`$. The action of $`\stackrel{~}{W}^𝔞`$ on $`𝒜`$ is no longer free and type preserving. The stabilizer $`\mathrm{\Omega }`$ of $`A_f`$ is isomorphic to $`X^{}/Q^{}`$. The isomorphism is given by sending $`g\mathrm{\Omega }`$ to the class of $`g(0)`$. So a set of representatives is given by $`X^{}A_f`$. We have $`\stackrel{~}{W}^𝔞\mathrm{\Omega }W^𝔞`$ and every element $`v\stackrel{~}{W}^𝔞`$ can be written as $`v=wg`$ for unique $`wW^𝔞`$ and $`g\mathrm{\Omega }`$. Although $`\stackrel{~}{W}^𝔞`$ is no longer a Coxeter group, we can extend the definition of the length function by setting $`l(v)=l(w)`$. So multiplication by elements of $`\mathrm{\Omega }`$ does not change the length. One also can extend the Bruhat order on $`\stackrel{~}{W}^𝔞`$ as follows: Let $`v=wg`$ and $`v^{}=w^{}g^{}\stackrel{~}{W}^𝔞`$ such that $`w,w^{}W^𝔞`$ and $`g,g^{}\mathrm{\Omega }`$. Then $`vv^{}`$ iff $`g=g^{}`$ and $`ww^{}`$ (in the usual Bruhat order on $`W^𝔞`$).
As mentioned above, the action of $`\stackrel{~}{W}^𝔞`$ on $`𝒜`$ is no longer free. So we introduce generalized alcoves $`\stackrel{~}{𝒜}`$ in order to work with the extended affine Weyl group as follows: Take an alcove $`A𝒜`$. Then some conjugate of $`\mathrm{\Omega }`$ acts transitively on $`AX^{}`$ and this intersection is in natural bijection to $`X^{}/Q^{}`$. Define $`\stackrel{~}{𝒜}:=\{(A,\mu )𝒜\times X^{}|\mu A\}`$. There is a natural embedding $`𝒜\stackrel{~}{𝒜}`$ sending an alcove $`A𝒜`$ to $`(A,\mu )`$ where $`\mu `$ is the unique element in $`AQ^{}`$. We identify $`𝒜`$ with its image in $`\stackrel{~}{𝒜}`$. We have a natural free $`\stackrel{~}{W}^𝔞`$-action on $`\stackrel{~}{𝒜}`$ given by the natural action on the two components. In particular, $`X^{}`$ acts on $`\stackrel{~}{𝒜}`$ by translations in both components. For $`\lambda X^{}`$ and $`A\stackrel{~}{𝒜}`$ we also write $`\lambda +A`$ for $`\tau _\lambda A`$.
We get a bijection $`\stackrel{~}{W}^𝔞\stackrel{~}{𝒜},wA_w:=wA_f`$ extending the bijection $`W^𝔞𝒜`$. In the same way as above we also get a right $`\stackrel{~}{W}^𝔞`$-action on $`\stackrel{~}{𝒜}`$ where $`\mathrm{\Omega }`$ acts only on the second factor. The definitions of face and type of a face carry over to this situation by demanding that $`\stackrel{~}{W}^𝔞`$ acts type preserving.
Every generalized alcove $`A`$ is of the form $`\mu +A_w`$ for unique $`\mu X^{}`$ and $`wW`$. Then $`\mu `$ is called the weight of $`A`$ and $`w`$ its direction. Denote this by $`wt(A):=\mu `$ and $`\delta (A):=w`$.
In various circumstances we will deal with stabilizer subgroups of $`W`$. We use the following notation for some notions related to them.
###### Definition 2.2.
Let $`\mu X^{}`$ and $`W_\mu W`$ its stabilizer. The maximal element of $`W_\mu `$ is denoted by $`w_\mu `$, the minimal representatives of $`W/W_\mu `$ by $`W^\mu `$ and the minimal element in the coset $`\tau _\mu W`$ by $`n^\mu `$.
In particular, $`W=W_0`$ and $`w_0`$ is the longest element in $`W`$. We will frequently use some facts about the length function on $`\stackrel{~}{W}^𝔞`$ summarized in
###### Lemma 2.3.
Let $`\lambda X_+^{}`$.
1. We have $`l(\tau _\lambda )=2\rho ,\lambda `$. In particular, $`l`$ is additive on $`X_+^{}`$.
2. One has $`\tau _\lambda w_\lambda =n^\lambda w_0`$ and $`l(\tau _\lambda )+l(w_\lambda )=l(n^\lambda )+l(w_0)`$. Moreover, $`n^\lambda W\tau _\lambda W`$ is minimal.
## 3 Affine Hecke algebra
Details on the affine Hecke algebra of a root datum with unequal parameters can be found in Lusztig’s article \[Lus89\].
For defining the affine Hecke algebra we first have to fix parameters. Let $`d:S^𝔞`$ be invariant under conjugation by elements of $`\stackrel{~}{W}^𝔞`$. Let $`=\left[q^{\pm {\scriptscriptstyle \frac{1}{2}}}\right]`$ and define $`q_s=q^{d(s)}`$ for $`sS^𝔞`$. For $`vW^𝔞`$ we set $`q_v=_{j=1}^kq_{t_j}`$ where $`v=t_1\mathrm{}t_k`$ with $`t_iS^𝔞`$ is a reduced decomposition of $`v`$. For arbitrary $`v\stackrel{~}{W}^𝔞`$ let $`q_v=q_v^{}`$ where $`v=v^{}g`$ with $`v^{}W^𝔞`$ and $`g\mathrm{\Omega }`$.
For a subset $`HW`$ define $`H(q)=_{wH}q_w`$ and $`H(q^1)=_{wH}q_w^1`$.
The affine Hecke algebra $`\stackrel{~}{}^𝔞`$ associated to the root datum $`\mathrm{\Phi }`$ and the above choice of $`d`$ is a $``$-algebra defined as follows: As a $``$-module it is free with basis $`\{T_w\}_{w\stackrel{~}{W}^𝔞}`$ and multiplication is given by
* $`T_s^2=q_sT_{id}+(q_s1)T_s`$ for all $`sS^𝔞`$.
* $`T_vT_w=T_{vw}`$ for all $`v,w\stackrel{~}{W}^𝔞`$ such that $`l(vw)=l(v)+l(w)`$.
On $`\stackrel{~}{}^𝔞`$ there is a natural $``$-algebra involution $`\overline{}:\stackrel{~}{}^𝔞\stackrel{~}{}^𝔞`$. It is given by $`\overline{T}_w=T_{w^1}^1`$ for $`w\stackrel{~}{W}^𝔞`$ and $`\overline{q^j}=q^j`$.
For $`\lambda X_+^{}`$ define $`q_\lambda =q^{{\scriptscriptstyle \frac{1}{2}}_{j=1}^kd(t_j)}`$ where $`\tau _\lambda =t_1\mathrm{}t_kg`$ is a reduced decomposition with $`g\mathrm{\Omega }`$. So we have $`q_\lambda ^2=q_{\tau _\lambda }`$. For arbitrary $`\mu X^{}`$ define $`q_\mu :=q_\lambda q_\lambda ^{}^1`$ where $`\lambda ,\lambda ^{}X_+^{}`$ such that $`\mu =\lambda \lambda ^{}`$. Then $`q_\mu `$ is independent of the particular choice of $`\lambda ,\lambda ^{}`$ because of the additivity of the length function on $`X_+^{}`$ (see lemma 2.3).
For each $`\mu X^{}`$ define an element $`X_\mu \stackrel{~}{}^𝔞`$ by $`X_\mu :=q_\mu ^1T_{\tau _\lambda }T_{\tau _\lambda ^{}}^1`$ where as above $`\mu =\lambda \lambda ^{}`$ with $`\lambda ,\lambda ^{}X_+^{}`$. By the same reason as above $`X_\mu `$ does not depend on the choice of $`\lambda `$ and $`\lambda ^{}`$. In particular we have $`X_\lambda =q_\lambda ^1T_{\tau _\lambda }`$ for all dominant $`\lambda `$. Thus one gets an inclusion of $``$-algebras
$`[X^{}]`$ $`\stackrel{~}{}^𝔞`$
$`x^\nu `$ $`X_\nu `$
and the image of $`[X^{}]^W`$ is the center of $`\stackrel{~}{}^𝔞`$. We identify $`[X^{}]`$ with its image.
Define $`\mathrm{𝟏}_0=_{wW}T_w\stackrel{~}{}^𝔞`$. We have $`T_w\mathrm{𝟏}_0=q_w\mathrm{𝟏}_0`$ for $`wW`$ and $`\mathrm{𝟏}_0^2=W(q)\mathrm{𝟏}_0`$. The spherical Hecke algebra $`^{sph}`$ is defined by
$$^{sph}=\left\{h\frac{1}{W(q)}\stackrel{~}{}^𝔞\right|T_wh=hT_w=q_wh\text{ for all }wW\}.$$
The Macdonald basis of $`^{sph}`$ is given by $`\{M_\lambda \}_{\lambda X_+^{}}`$ where
$`M_\lambda `$ $`:={\displaystyle \frac{1}{W(q)}}{\displaystyle \underset{wW\tau _\lambda W}{}}T_w={\displaystyle \frac{1}{W(q)W_\lambda (q)}}\mathrm{𝟏}_0T_{n^\lambda }\mathrm{𝟏}_0`$
$`={\displaystyle \frac{q_\lambda q_{w_0}^1}{W(q)W_\lambda (q^1)}}\mathrm{𝟏}_0X_\lambda \mathrm{𝟏}_0.`$
The second equality follows from lemma 2.3 which yields $`X_\lambda =q_\lambda T_{n^\lambda }T_{w_0}\overline{T}_{w_\lambda }`$. Moreover, $`W_\lambda (q^1)=q_{w_\lambda }^1W_\lambda (q)`$. One obtains an isomorphism (Satake)
$`[X^{}]^W`$ $`\stackrel{}{}^{sph}`$
$`x`$ $`{\displaystyle \frac{1}{W(q)}}x\mathrm{𝟏}_0`$
In particular, we have
$`m_\lambda `$ $`Y_\lambda :={\displaystyle \frac{1}{W(q)}}{\displaystyle \underset{\mu W\lambda }{}}X_\mu \mathrm{𝟏}_0`$
So we have two bases for $`^{sph}`$: The natural basis given by the Macdonald basis $`\{M_\lambda \}_{\lambda X_+^{}}`$ and the monomial basis $`\{Y_\lambda \}_{\lambda X_+^{}}`$ given by the images of the monomial symmetric functions under the Satake isomorphism. We are interested in the transition matrix from the monomial basis to the natural basis. (Re)define the family $`\{L_{\lambda \mu }\}_{\lambda ,\mu X_+^{}}`$ as modified entries of this transition matrix. More precisely, we have
$$M_\lambda =\underset{\mu X_+^{}}{}q_\mu ^1L_{\lambda \mu }Y_\mu .$$
For arbitrary $`\mu X^{}`$ and dominant $`\lambda X_+^{}`$ we set $`L_{\lambda \mu }=q_{\mu \mu ^+}L_{\lambda \mu ^+}`$ where $`\mu ^+`$ is the unique dominant element in the $`W`$-orbit of $`\mu `$.
We also calculate the structure constants of the spherical Hecke algebra with respect to the Macdonald basis. For this, (re)define $`\{C_{\lambda \mu }^\nu \}_{\lambda ,\mu ,\nu X_+^{}}`$ as modified structure constants by
$$M_\lambda M_\mu =\underset{\nu X_+^{}}{}q_{\lambda \nu }^2C_{\lambda \mu }^\nu M_\nu .$$
In the next chapter we give a description of the $`L_{\lambda \mu }`$ and the $`C_{\lambda \mu }^\nu `$ using galleries.
###### Remark 3.1.
This is not the most general choice of parameters for which the affine Hecke algebra is defined and where the theorems 4.4 and 4.10 are true. One important example is the following: Replace $``$ by the image of the morphism $``$ evaluating the variable $`q`$ at some fixed prime power. Hecke algebras of reductive groups over local fields are of this form (compare the end of section 8 for the case of equal parameters).
###### Remark 3.2.
Now we want to clarify the relations of this section to symmetric polynomials and their $`q`$-analogs. In particular, we describe the relation between the coefficients defined above and the ones with the same names in section 1.
For this regard the case of equal parameters, i.e. $`d(s)=1`$ for all $`sS^𝔞`$. In this case we have $`q_v=q^{l(v)}`$ for $`v\stackrel{~}{W}^𝔞`$ and $`q_\mu =q^{\rho ,\mu }`$ for $`\mu X^{}`$. It is known (see for example \[NR03, theorem 2.9\]) that the image of $`P_\lambda (q^1)`$ under the Satake isomorphism is $`q_\lambda M_\lambda `$. So comparing the definitions of the $`L_{\lambda \mu }`$ and $`C_{\lambda \mu }^\nu `$ in section 1 with the ones given here shows that the first ones are special cases of the latter ones. So the theorems stated there will follow from theorems 4.4 and 4.10 given in the next section.
## 4 Galleries
In this section we introduce galleries and some polynomials associated to them. We then give a precise meaning to the theorems stated in the introduction in the general setting of the last section. The galleries used here are a slight generalization of the usual galleries in a Coxeter complex since we regard generalized alcoves instead of alcoves.
###### Definition 4.1.
Let $`t=(t_1,\mathrm{},t_k)`$ with $`t_iS^𝔞\mathrm{\Omega }`$. Let $`sS^𝔞`$.
* A gallery $`\sigma `$ of type $`t`$ connecting generalized alcoves $`A`$ and $`B`$ is a sequence $`(A=A_0,\mathrm{},B=A_k)`$ of generalized alcoves such that $`A_{i+1}=A_it_{i+1}`$ if $`t_{i+1}\mathrm{\Omega }`$ and $`A_{i+1}\{A_i,A_it_{i+1}\}`$ if $`t_{i+1}S^𝔞`$. In the case of $`t_{i+1}S^𝔞`$ this means that $`A_i`$ and $`A_{i+1}`$ have a common face of type $`t_{i+1}`$.
* The initial direction $`\iota (\sigma )`$ is defined to be the direction $`\delta (A_0)`$ of the first generalized alcove. The weight $`wt(\sigma )`$ of $`\sigma `$ is $`wt(A_k)`$, the ending $`e(\sigma )`$ is $`A_k`$ and the final direction $`\epsilon (\sigma )`$ is $`\delta (A_k)`$.
* The gallery $`\sigma `$ has a positive $`s`$-direction at $`i`$ if $`t_{i+1}=s`$, $`A_{i+1}=A_is`$ and $`A_i`$ is negative with respect to $`s`$, i.e. $`A_iA_{i+1}`$. The separating hyperplane is the wall of $`A_i`$ corresponding to the face of type $`s`$.
* The gallery $`\sigma `$ is $`s`$-folded at $`i`$ if $`t_{i+1}=s`$ and $`A_{i+1}=A_i`$. The folding hyperplane is the wall of $`A_i`$ corresponding to the face of type $`s`$. The folding is positive if $`A_iA_is`$.
We call $`\sigma `$ positively folded, if all foldings occurring are positive. A gallery is said to be minimal if it is of minimal length among all galleries connecting the same generalized alcoves.
For the precise statement on the $`L_{\lambda \mu }`$ and the $`C_{\lambda \mu }^\nu `$ we need some statistics on galleries.
###### Definition 4.2.
Let $`\sigma `$ be a positively folded gallery of type $`t`$. For $`sS^𝔞`$ define
* $`m_s(\sigma )`$ the number of positive $`s`$-directions.
* $`n_s(\sigma )`$ the number of positive $`s`$-folds.
* $`r_s(\sigma )`$ the number of positive $`s`$-folds such that the folding hyperplane is not a wall of the dominant chamber $`𝒞`$.
* $`p_s(\sigma )`$ the number of positive $`s`$-folds such that the folding hyperplane is a wall of $`𝒞`$.
In particular, $`r_s(\sigma )+p_s(\sigma )=n_s(\sigma )`$. Now we can define
* $`L_\sigma =_{sS^𝔞}q_s^{m_s(\sigma )}(q_s1)^{n_s(\sigma )}`$ and
* $`C_\sigma =_{sS^𝔞}q_s^{m_s(\sigma )+p_s(\sigma )}(q_s1)^{r_s(\sigma )}`$.
For a gallery $`\sigma `$ such that no folding hyperplane is a wall of $`𝒞`$ one has $`L_\sigma =C_\sigma `$. In the case of equal parameters (see remark 3.2) we have $`\mathrm{deg}L_\sigma =\mathrm{deg}C_\sigma `$ for any gallery $`\sigma `$.
Fix some type $`t=(t_1,\mathrm{},t_k)`$. For $`A\stackrel{~}{𝒜}`$ and $`\mu X^{}`$ let $`\mathrm{\Gamma }_t^+(A,\mu )`$ be the set of all positively folded galleries of type $`t`$ starting in $`A`$ with weight $`\mu `$. Further let $`\mathrm{\Gamma }_t^+(\mu )=_{wW}\mathrm{\Gamma }_t^+(A_w,\mu )`$ be the set of all positively folded galleries of weight $`\mu `$ starting in the origin and let $`\mathrm{\Gamma }_t^+`$ be the set of all positively folded galleries starting in the origin. Define
$$L_t(\mu ):=\underset{\sigma \mathrm{\Gamma }_t^+(\mu )}{}q_{w_0\iota (\sigma )}L_\sigma .$$
So there is an additional contribution measuring the distance from $`A_f`$ to the initial alcove.
###### Remark 4.3.
There is an alternative way of defining $`L_t(\mu )`$: For any $`wW`$ choose a minimal gallery $`\sigma _w`$ of type $`t_w`$ which connects $`A_f`$ and $`A_w`$. Then $`\sigma _w`$ is a nonfolded gallery of length $`l(w_0w)=l(w_0)l(w)`$ and it has only positive directions. The positively folded galleries of type $`t_w^{}=(t_w,t)`$ beginning in $`A_f`$ correspond to the positively folded galleries of type $`t`$ starting in $`A_w`$. We get
$$L_t(\mu )=\underset{wW}{}\left(\underset{\sigma \mathrm{\Gamma }_{t_w^{}}^+(A_f,\mu )}{}L_\sigma \right).$$
Now we can give the formula for the Satake coefficients. Let $`\lambda X_+^{}`$ and recall the notation introduced in definition 2.2. Let $`\sigma ^\lambda `$ be a minimal gallery connecting $`A_f`$ and $`A_{n^\lambda }`$ and denote its type by $`t^\lambda `$. Using the last definition we get polynomials $`L_{t^\lambda }(\mu )`$ for all $`\mu X^{}`$. Up to some factor these are the $`L_{\lambda \mu }`$. More precisely we prove in section 5:
###### Theorem 4.4.
For $`\mu X^{}`$ we have
$$L_{\lambda \mu }=\frac{1}{W_\lambda (q)}L_{t^\lambda }(\mu ).$$
Furthermore,
$$L_{\lambda \mu }=q_{w_\lambda }^1\underset{\begin{array}{c}\sigma \mathrm{\Gamma }_{t^\lambda }^+(\mu )\\ \iota (\sigma )W^\lambda \end{array}}{}q_{w_0\iota (\sigma )}L_\sigma .$$
In particular the $`L_{t^\lambda }(\mu )`$ do not depend on the choice of the minimal gallery $`\sigma ^\lambda `$ and $`L_{t^\lambda }(\mu )=q_{\mu w\mu }L_{t^\lambda }(w\mu )`$ for all $`wW`$.
###### Remark 4.5.
One of the surprising implications of the last theorem is the $`W`$-invariance of the $`L_{t^\lambda }(\mu )`$ up to some power of $`q`$. This is surprising because even the cardinality of the sets $`\mathrm{\Gamma }_{t^\lambda }^+(w\mu )`$ depends on $`w`$.
###### Remark 4.6.
Let $`wW^𝔞`$. The choice of a minimal gallery $`\sigma `$ connecting $`A_f`$ and $`A_w`$ is equivalent to the choice of a reduced expression for $`w`$. Let $`t=(t_1,\mathrm{},t_k)`$ be the type of $`\sigma `$. Then we have the reduced expression $`w=t_1\mathrm{}t_k`$.
Let $`v\stackrel{~}{W}^𝔞`$. Then $`v`$ can be written as $`v=wg`$ with $`wW^𝔞`$ and $`g\mathrm{\Omega }`$. A minimal gallery $`\sigma `$ from $`A_f`$ to $`A_v`$ is given by a minimal gallery from $`A_f`$ to $`A_w`$ extended by $`A_v`$. So one can always arrange that at most the last entry of the type of a minimal gallery is in $`\mathrm{\Omega }`$.
Now it is quite natural to ask when $`\mathrm{\Gamma }_{t^\lambda }^+(\mu )\mathrm{}`$. Although the definition of galleries is a combinatorial one, it seems hard to give a combinatorial proof for the existence (or non existence) of a gallery of given type and weight. Let $`\sigma `$ be any gallery of type $`t^\lambda `$ starting in 0, ending in $`A_v`$ of weight $`\mu `$. Since the folding hyperplanes are root hyperplanes we always have $`\lambda \mu Q^{}`$. Moreover, $`v\iota (\sigma )n^\lambda `$ by definition of the Bruhat order on $`\stackrel{~}{W}^𝔞`$. This implies $`\mu ^+\lambda `$. This also follows from the well known fact that the transition matrix from the monomial basis to the Macdonald basis is triangular with respect to the dominance ordering on $`X_+^{}`$.
The question of the existence of a gallery in $`\mathrm{\Gamma }_{t^\lambda }^+(\mu )`$ does not depend on the choice of parameters $`d`$. So we can take $`d=1`$ as in remark 3.2. Since $`P_\lambda (q^1)`$ and $`m_\mu `$ are contained in $`\mathrm{\Lambda }_q`$ we have $`q^{\rho ,\lambda +\mu }L_{\lambda \mu }^{}`$. Moreover, $`q^{l(w_\lambda )}W_\lambda (q)=W_\lambda (q^1)^{}`$ and thus $`q^{\rho ,\lambda +\mu l(w_\lambda )}L_{t^\lambda }(\mu )^{}`$. So we get the upper bound
$$\mathrm{deg}(L_\sigma )+l(w_0\iota (\sigma ))\rho ,\mu +\lambda +l(w_\lambda )$$
(4.1)
for all $`\sigma \mathrm{\Gamma }_{t^\lambda }^+(\mu )`$. The galleries with maximal degree are of special interest. So define
###### Definition 4.7.
A gallery $`\sigma \mathrm{\Gamma }_{t^\lambda }^+`$ is a LS-gallery if we have equality in the above equation, i.e. $`\mathrm{deg}(L_\sigma )+l(w_0\iota (\sigma ))=\rho ,wt(\sigma )+\lambda +l(w_\lambda )`$.
Since $`L_\sigma `$ is monic we get corollary 1.2 by evaluating theorem 4.4 at $`q^1=0`$.
The number of LS-galleries in $`\mathrm{\Gamma }_{t^\lambda }(w\mu )`$ is $`W`$-invariant. This follows from the $`W`$-invariance (up to a power of $`q`$) of $`L_{t^\lambda }(w\mu )`$ in 4.4. Now let $`\mu X^{}`$ such that $`\mu ^+\lambda `$. We know from representation theory that $`k_{\lambda \mu ^+}>0`$. So there exists a LS-gallery in $`\mathrm{\Gamma }_{t^\lambda }^+(\mu ^+)`$ and thus also in $`\mathrm{\Gamma }_{t^\lambda }^+(\mu )`$ by the $`W`$-invariance.
Summarizing all this in the following corollary answers the above question on the existence of galleries with a given weight and sharpens the triangularity.
###### Corollary 4.8.
The number of LS-galleries in $`\mathrm{\Gamma }_{t^\lambda }^+(\mu )`$ is $`k_{\lambda \mu ^+}`$. In particular we have $`\mathrm{\Gamma }_{t^\lambda }^+(\mu )\mathrm{}`$ iff $`\mu `$ occurs as a weight in $`V(\lambda )`$, i.e. $`\mu ^+\lambda `$. Moreover, we have (for arbitrary parameters) $`L_{\lambda \mu }0`$ iff $`\mu \lambda `$.
Specializing $`q`$ at some prime power we get $`L_{\lambda \mu }>0`$ for all $`\mu \lambda `$. This was shown by Rapoport for the case of spherical Hecke algebras of a reductive group over a local field \[Rap00\].
###### Remark 4.9.
For regular $`\lambda `$ the definition of galleries coincides with the one given in \[GL05\]. Instead of using generalized alcoves they regard galleries of alcoves together with an initial and final weight in $`X^{}`$ contained in the first respectively last alcove. This is equivalent to our definition since we can always arrange such that at most the last component of $`t^\lambda `$ is in $`\mathrm{\Omega }`$ (compare remark 4.6). For nonregular $`\lambda `$ they regard degenerate alcoves. This is more or less the same as our choice of the initial direction. See also remark 5.8 for a discussion of this choice.
The proof of corollary 1.2 in \[GL05\] is quite different from here. They define root operators on the set of all galleries of type $`t^\lambda `$ starting in the origin. Then they show that the subset of LS-galleries is closed under these operators and defines the highest weight crystal with highest weight $`\lambda `$.
We now give the formula for the structure constants replacing $`L_\sigma `$ with $`C_\sigma `$. So let $`\lambda X_+^{}`$ and $`t`$ be any type. Define $`\mathrm{\Gamma }_{t,\lambda }^d`$ as the set of all positively folded galleries of type $`t`$ starting in $`\lambda `$ which are contained in the dominant chamber. Here we allow that folding hyperplanes are contained in the walls of $`𝒞`$. For $`\nu X_+^{}`$ let $`\mathrm{\Gamma }_{t,\lambda }^d(\nu )\mathrm{\Gamma }_{t,\lambda }^d`$ be the subset of galleries of weight $`\nu `$. Define
$$C_{\lambda t}(\nu )=\underset{\mathrm{\Gamma }_{t,\lambda }^d(\nu )}{}q_{w_0\iota (\sigma )}C_\sigma .$$
Now let $`\lambda ,\mu X_+^{}`$ and let $`t^\mu `$ be the type of a minimal gallery connecting $`A_f`$ and $`A_{n^\mu }`$ where $`n^\mu \tau _\mu W`$ is the minimal representative in $`\tau _\mu W`$. The above definition yields $`C_{\lambda t^\mu }(\nu )`$ for any $`\nu X_+^{}`$. Define $`F_{\mu \nu }^w:=q_w_{vW^{w_0\mu }W_\nu w}q_v^1`$ for $`\mu ,\nu X_+^{}`$ and $`wW`$. In section 6 we prove:
###### Theorem 4.10.
For $`\lambda ,\mu ,\nu X_+^{}`$ we have
$$C_{\lambda \mu }^\nu =\frac{W_\nu (q^1)}{W_\mu (q)}C_{\lambda t^\mu }(\nu ).$$
Furthermore,
$$C_{\lambda \mu }^\nu =q_{w_\mu }^1\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d(\nu )}{}q_{w_0\iota (\sigma )}C_\sigma F_{\mu \nu }^{\epsilon (\sigma )}.$$
In particular, the $`C_{\lambda t^\mu }(\nu )`$ do not depend on the choice of the minimal gallery.
So in contrast to theorem 4.4 we have a condition on the final direction $`\epsilon (\sigma )`$ since $`F_{\mu \nu }^w=0`$ iff $`wW_\nu W^{w_0\mu }`$.
As above we can give an estimate for the degree of the $`C_{\lambda t^\nu }(\nu )`$ in the case of equal parameters and prove corollary 1.4. From the last theorem we get $`q^{\rho ,\mu \lambda +\nu l(w_\mu )}C_{\lambda t^\mu }(\nu )^{}`$ and thus for any $`\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d(\nu )`$ we have
$$\mathrm{deg}C_\sigma +l(w_0\iota (\sigma ))\rho ,\mu \lambda +\nu +l(w_\mu ).$$
Since $`\mathrm{deg}L_\sigma =\mathrm{deg}C_\sigma `$ and translating a gallery by an element of $`X^{}`$ does not change $`L_\sigma `$ and the initial direction, corollary 1.4 is proven and we get
###### Corollary 4.11.
For $`\lambda ,\mu ,\nu X_+^{}`$ we have $`C_{\lambda \mu }^\nu 0`$ if $`c_{\lambda \mu }^\nu 0`$.
Specializing $`q`$ to a prime power one gets that $`C_{\lambda \mu }^\nu >0`$ if $`c_{\lambda \mu }^\nu >0`$. For equal parameters this is proven in \[KM04\] and also by Haines in \[Hai03\] by geometric arguments using the affine Grassmanian of the Langlands dual $`G`$ of $`G^{}`$ to calculate the degree and the leading coefficients of $`C_{\lambda \mu }^\nu `$.
Another interpretation of the structure constants was given by Parkinson \[Par06\] using regular affine buildings.
## 5 Satake coefficients
In this section we introduce the alcove basis of the extended affine Hecke algebra and show that right multiplication of this alcove basis by elements of the standard basis can be calculated using positively folded galleries. From this theorem 4.4 follows. We also show that one can replace positively folded galleries by negatively folded galleries.
###### Definition 5.1.
Let $`A\stackrel{~}{𝒜}`$. Define $`X_A=q_{wt(A)}q_{\delta (A)}X_{wt(A)}\overline{T}_{\delta (A)}`$.
The set $`\{X_A\}_{A\stackrel{~}{𝒜}}`$ is a basis of $`\stackrel{~}{}^𝔞`$. Before we proceed, we need some properties of this basis. First let $`\lambda X^{}`$ and $`A\stackrel{~}{𝒜}`$. One calculates
$$X_\lambda X_A=q_\lambda X_{\lambda +A}.$$
(5.1)
Now assume $`A=A_v`$ to be dominant such that $`\lambda :=wt(A)`$ is regular. Then $`v=\tau _\lambda \delta (A)`$. Moreover, $`\tau _\lambda `$ is of maximal length in $`\tau _\lambda W`$ by lemma 2.3 and $`l(v)=l(\tau _\lambda )l(\delta (A))`$. So we get $`T_{\tau _\lambda }\overline{T}_{\delta (A)}=T_{\tau _\lambda \delta (A)}=T_v`$ and thus
$$X_A=q_\lambda q_{\delta (A)}X_\lambda \overline{T}_{\delta (A)}=q_{\tau _\lambda }^1q_{\delta (A)}T_{\tau _\lambda }\overline{T}_{\delta (A)}=q_v^1T_v.$$
(5.2)
Multiplying the elements of the alcove basis with $`T_s`$ from the right can be expressed in terms of the alcove order. It is a $`q`$-analog of the $`\stackrel{~}{W}^𝔞`$-action on $`\stackrel{~}{𝒜}`$.
###### Lemma 5.2.
Let $`A\stackrel{~}{𝒜}`$. In $`\stackrel{~}{}^𝔞`$ we have
$$X_AT_s=\{\begin{array}{cc}q_sX_{As}\hfill & \text{ if }AAs\hfill \\ X_{As}+(q_s1)X_A\hfill & \text{ if }AAs.\hfill \end{array}$$
###### Proof.
By (5.1) the assertion is invariant under translation, i.e. under left multiplication with some $`X_\mu `$. So it is enough to show the assertion for alcoves $`A=A_v`$ such that $`wt(A)\alpha ^{}`$ is dominant and regular for all $`\alpha \varphi `$. By (5.2) we have $`X_A=q_v^1T_v`$ and the multiplication law in $`\stackrel{~}{}^𝔞`$ yields
$$T_vT_s=\{\begin{array}{cc}T_{vs}\hfill & \text{ if }l(v)<l(vs)\hfill \\ q_sT_{vs}+(q_s1)T_v\hfill & \text{ if }l(v)>l(vs).\hfill \end{array}$$
But for generalized alcoves in the dominant chamber increasing in the alcove order is equivalent to increasing the length of the corresponding elements of $`\stackrel{~}{W}^𝔞`$ (see example 2.1). Moreover, by the choice of $`A`$ we get $`X_{As}=q_{vs}^1T_{vs}`$ as elements in $`\stackrel{~}{}^𝔞`$ again by (5.2) and the assertion follows. ∎
Using the same arguments and the fact that multiplying by $`T_g`$ for $`g\mathrm{\Omega }`$ does not change the length we get
###### Lemma 5.3.
For $`A\stackrel{~}{𝒜}`$ we have $`X_AT_g=X_{Ag}`$ as elements in $`\stackrel{~}{}^𝔞`$.
Now we can connect the multiplication in $`\stackrel{~}{}^𝔞`$ to the $`L`$-polynomials. For generalized alcoves $`A`$ and $`B`$ and any type $`t`$ define $`\mathrm{\Gamma }_t^+(A,B)`$ to be the set of all positively folded galleries of type $`t`$ connecting $`A`$ and $`B`$ and set $`L_t(A,B)=_{\sigma \mathrm{\Gamma }_t^+(A,B)}L_\sigma `$.
###### Lemma 5.4.
Let $`t=(t_1,\mathrm{},t_k),sS^𝔞`$, $`t^{}=(t_1,\mathrm{},t_k,s)`$, and fix generalized alcoves $`A`$ and $`B`$. We have
$$L_t^{}(A,Bs)=\{\begin{array}{cc}L_t(A,B)\hfill & \text{ if }BBs\hfill \\ q_sL_t(A,B)+(q_s1)L_t(A,Bs)\hfill & \text{ if }BBs.\hfill \end{array}$$
###### Proof.
Let $`\sigma ^{}=(A,\mathrm{},C,Bs)\mathrm{\Gamma }_t^{}^+(A,Bs)`$. Then $`C\{B,Bs\}`$. We have $`C=Bs`$ iff $`\sigma ^{}`$ is $`s`$-folded at $`k+1`$. Let $`\sigma =(A,\mathrm{},C)`$ and distinguish two cases:
$`BBs`$: We then have $`C=B`$ and $`\sigma ^{}`$ is negative at $`k+1`$. So $`\sigma \mathrm{\Gamma }_t^+(A,B)`$ and $`L_\sigma ^{}=L_\sigma `$. Moreover, all galleries in $`\mathrm{\Gamma }_t^+(A,B)`$ are obtained this way.
$`BBs`$: If $`C=B`$ we have $`\sigma \mathrm{\Gamma }_t^+(A,B)`$ and $`\sigma ^{}`$ is positive at $`k+1`$. So $`L_\sigma ^{}=q_sL_\sigma `$ and one gets all galleries in $`\mathrm{\Gamma }_t^+(A,B)`$ this way. If $`C=Bs`$ we have $`\sigma \mathrm{\Gamma }_t^+(A,Bs)`$, $`L_\sigma ^{}=(q_s1)L_\sigma `$ and one obtains all galleries in $`\mathrm{\Gamma }_t^+(A,Bs)`$ this way.
The lemma follows. ∎
Let $`v\stackrel{~}{W}^𝔞`$ and $`\sigma `$ be a minimal gallery of type $`t`$ connecting $`A_f`$ and $`A_v`$.
###### Theorem 5.5.
Given $`A\stackrel{~}{𝒜}`$ one has $`X_AT_v=_{B\stackrel{~}{𝒜}}L_t(A,B)X_B`$.
###### Proof.
Because of lemma 5.3 and since the $`L`$-polynomials are not affected by elements of $`\mathrm{\Omega }`$ in the type it is enough to show the theorem for $`vW^𝔞`$. The proof is done by induction on $`l(v)`$.
Let first $`v=sS^𝔞`$: Distinguish two cases.
$`AAs`$: In this case $`L_{(s)}(A,A)=0`$, $`L_{(s)}(A,As)=q_s`$ and $`L_{(s)}(A,B)=0`$ for all other $`B`$ and $`X_AT_s=q_sX_{As}`$.
$`AAs`$: In this case $`L_{(s)}(A,A)=q_s1`$, $`L_{(s)}(A,As)=1`$ and $`L_{(s)}(A,B)=0`$ for all other $`B`$ and $`X_AT_s=X_{As}+(q_s1)X_A`$.
Now let $`vW^𝔞`$, $`sS^𝔞`$ such that $`l(v)<l(vs)`$ and $`\sigma ^{}=(A_0,\mathrm{},A_v,A_{vs})`$ is a minimal gallery of type $`t^{}`$. Using the last lemma we get
$`X_AT_{vs}`$ $`=X_AT_vT_s=\left({\displaystyle \underset{BW^𝔞}{}}L_t(A,B)X_B\right)T_s`$
$`={\displaystyle \underset{BBs}{}}q_sL_t(A,B)X_{Bs}+{\displaystyle \underset{BBs}{}}L_t(A,B)X_{Bs}+{\displaystyle \underset{BBs}{}}(q_s1)L_t(A,B)X_B`$
$`={\displaystyle \underset{BBs}{}}\left(q_sL_t(A,B)+(q_s1)L_t(A,Bs)\right)X_{Bs}+{\displaystyle \underset{BBs}{}}L_t(A,B)X_{Bs}`$
$`={\displaystyle \underset{B\stackrel{~}{𝒜}}{}}L_t^{}(A,Bs)X_{Bs}={\displaystyle \underset{B\stackrel{~}{𝒜}}{}}L_t^{}(A,B)X_B`$
In particular we get that $`L_t(A,B)`$ does not depend on $`\sigma `$ and $`t`$ but only on $`v`$. Thus the following definition is well defined.
###### Definition 5.6.
For $`v\stackrel{~}{W}^𝔞`$ define $`L_v(A,B):=L_t(A,B)`$ where $`t`$ is the type of a minimal gallery from $`A_f`$ to $`A_v`$.
With these results we now can prove proposition 4.4.
###### Lemma 5.7.
For $`\lambda X_+^{}`$ we have
$$\mathrm{𝟏}_0T_{n^\lambda }\mathrm{𝟏}_0=\underset{\mu X^{}}{}q_\mu L_{t^\lambda }(\mu )X_\mu \mathrm{𝟏}_0.$$
###### Proof.
We use the last theorem and the facts that $`\overline{\mathrm{𝟏}_0}=q_{w_0}^1\mathrm{𝟏}_0`$ and $`\overline{T}_w\mathrm{𝟏}_0=q_w^1\mathrm{𝟏}_0`$ for all $`wW`$. So one calculates
$`\mathrm{𝟏}_0T_{n^\lambda }\mathrm{𝟏}_0`$ $`=q_{w_0}{\displaystyle \underset{wW}{}}\overline{T}_wT_{n^\lambda }\mathrm{𝟏}_0=q_{w_0}{\displaystyle \underset{wW}{}}q_w^1X_{A_w}T_{n^\lambda }\mathrm{𝟏}_0`$
$`=q_{w_0}{\displaystyle \underset{wW}{}}q_w^1{\displaystyle \underset{\sigma \mathrm{\Gamma }_{t^\lambda }^+,\iota (\sigma )=w}{}}q_{wt(\sigma )}q_{\epsilon (\sigma )}L_\sigma X_{wt(\sigma )}\overline{T}_{\epsilon (\sigma )}\mathrm{𝟏}_0`$
$`={\displaystyle \underset{wW}{}}q_{w_0w}{\displaystyle \underset{\sigma \mathrm{\Gamma }_{t^\lambda }^+,\iota (\sigma )=w}{}}q_{wt(\sigma )}L_\sigma X_{wt(\sigma )}\mathrm{𝟏}_0`$
$`={\displaystyle \underset{\sigma \mathrm{\Gamma }_{t^\lambda }^+}{}}q_{w_0\iota (\sigma )}q_{wt(\sigma )}L_\sigma X_{wt(\sigma )}\mathrm{𝟏}_0={\displaystyle \underset{\mu X^{}}{}}q_\mu L_{t^\lambda }(\mu )X_\mu \mathrm{𝟏}_0`$
where the last equality holds by the definition of $`L_{t^\lambda }(\mu )`$ in section 4. ∎
From this we get
$$M_\lambda =\frac{1}{W(q)W_\lambda (q)}\underset{\mu X^{}}{}q_\mu L_{t^\lambda }(\mu )X_\mu \mathrm{𝟏}_0.$$
But on the other hand $`q_\mu L_{\lambda \mu }`$ for dominant $`\mu `$ is the coefficient of $`M_\lambda `$ with respect to $`Y_\mu `$. Moreover, for arbitrary $`\nu X^{}`$ we defined $`L_{\lambda \nu }=q_{\nu \nu ^+}L_{\lambda \nu ^+}`$. So we get
$`M_\lambda `$ $`={\displaystyle \underset{\mu X_+^{}}{}}q_\mu L_{\lambda \mu }Y_\mu ={\displaystyle \frac{1}{W(q)}}{\displaystyle \underset{\mu X_+^{}}{}}\left({\displaystyle \underset{\nu W\mu }{}}q_\nu L_{\lambda \nu }X_\nu \mathrm{𝟏}_0\right)`$
$`={\displaystyle \frac{1}{W(q)}}{\displaystyle \underset{\mu X^{}}{}}q_\mu L_{\lambda \mu }X_\mu \mathrm{𝟏}_0.`$
Comparing coefficients of these two expansions we get
$$L_{\lambda \mu }=\frac{1}{W_\lambda (q)}L_{t^\lambda }(\mu )$$
which proves the first statement in 4.4. The second statement can be obtained as follows: Every $`wW`$ can be written as $`w=w_1w_2`$ for unique $`w_1W^\lambda `$ and $`w_2W_\lambda `$ such that $`l(w)=l(w_1)+l(w_2)`$ (using the notation introduced in definition 2.2). Define $`\mathrm{𝟏}_\lambda =_{wW_\lambda }T_w`$. Since $`\overline{T}_v\overline{T}_w=\overline{T}_{vw}`$ for $`v,wW`$ with $`l(v)+l(w)=l(vw)`$ and $`\overline{\mathrm{𝟏}_\lambda }=q_{w_\lambda }^1\mathrm{𝟏}_\lambda `$ we get
$$\mathrm{𝟏}_0=q_{w_0}\underset{wW^\lambda }{}\overline{T}_w\overline{\mathrm{𝟏}_\lambda }=q_{w_0w_\lambda }\underset{wW^\lambda }{}\overline{T}_w\mathrm{𝟏}_\lambda .$$
If $`vW_\lambda `$ we have $`l(v)+l(n^\lambda )=l(vn^\lambda )`$. Moreover, $`vn^\lambda =v\tau _\lambda w_\lambda w_0=\tau _\lambda vw_\lambda w_0=n^\lambda v^{}`$ with $`v^{}=w_0w_\lambda vw_\lambda w_0`$ by lemma 2.3. Then $`l(v^{})=l(v)`$ and $`q_v=q_v^{}`$. Thus $`T_vT_{n^\lambda }\mathrm{𝟏}_0=T_{n^\lambda }T_v^{}\mathrm{𝟏}_0=q_vT_{n^\lambda }\mathrm{𝟏}_0`$ and we get
$$\mathrm{𝟏}_0T_{n^\lambda }\mathrm{𝟏}_0=q_{w_0w_\lambda }W_\lambda (q)\underset{wW^\lambda }{}\overline{T}_wT_{n^\lambda }\mathrm{𝟏}_0.$$
Now the second statement follows the same way as in the proof of lemma 5.7 using $`q_{w_0w_\lambda }_{wW^\lambda }\overline{T}_w=q_{w_\lambda }^1_{wW^\lambda }q_{w_0w}X_{A_w}`$.
###### Remark 5.8.
In the above considerations there are various other choices for the condition on the initial alcove. Let $`vW_\lambda `$. We have
$$\mathrm{𝟏}_0=q_{w_0w_\lambda v}\underset{wW^\lambda }{}\overline{T}_{wv}\mathrm{𝟏}_\lambda $$
since $`\overline{T}_v\mathrm{𝟏}_\lambda =q_v^1\mathrm{𝟏}_\lambda `$ and the last equation becomes
$$\mathrm{𝟏}_0T_{n^\lambda }\mathrm{𝟏}_0=q_{w_0w_\lambda v}W_\lambda (q)\underset{wW^\lambda v}{}\overline{T}_wT_{n^\lambda }\mathrm{𝟏}_0.$$
Thus one gets
$$L_{\lambda \mu }=q_{w_\lambda v}^1\underset{\begin{array}{c}\sigma \mathrm{\Gamma }_{t^\lambda }^+(\mu )\\ \iota (\sigma )W^\lambda v\end{array}}{}q_{w_0\iota (\sigma )}L_\sigma .$$
The case considered above was $`v=id`$. In the case of equal parameters we get for any gallery $`\sigma \mathrm{\Gamma }_{t^\lambda }^+`$ such that $`\iota (\sigma )W^\lambda v`$ the upper bound
$$\mathrm{deg}L_\sigma +l(w_0\iota (\sigma ))\rho ,\lambda +wt(\sigma )+l(w_\lambda v).$$
One could define LS-galleries to be the ones such that $`\iota (\sigma )W^\lambda v`$ and where there is equality in the last equation. But only with the choice $`v=id`$ it is enough to impose this equality. The condition on the initial direction follows from this. In particular, for a LS-gallery $`\sigma `$ we have $`\iota (\sigma )W^\lambda `$.
For the definition of the $`L_{t^\lambda }(\mu )`$ we started with the minimal representative $`n^\lambda `$ and we showed that $`L_{t^\lambda }(\mu )`$ is independent of the initially chosen minimal gallery. One can allow even more freedom in this initial choice. Let $`vW\tau _\lambda W`$ and let $`w,w^{}W_\lambda `$ such that $`v=wn^\lambda w^{^{}}`$ and $`l(w)+l(n^\lambda )+l(w^{})=l(v)`$. If instead of $`t^\lambda `$ we use the type $`t`$ of a minimal gallery from $`A_f`$ to $`A_v`$ we get from the proof of 5.7 that $`L_t(\mu )=q_wq_w^{}L_{t^\lambda }(\mu )`$ for any $`\mu X^{}`$. It is clear that the number of LS-galleries in $`\mathrm{\Gamma }_t^+(\mu )`$ (with the appropriate changes of the degree condition in the definition) is the same as in $`\mathrm{\Gamma }_{t^\lambda }^+(\mu )`$ since they always encode $`s_\lambda `$. One also has a canonical bijection between these different sets of LS-galleries. But the total number of galleries in $`\mathrm{\Gamma }_t^+(\mu )`$ really depends on the choice of $`v`$ and this number is minimal if we choose $`n^\lambda `$. There is another fact that singles out $`n^\lambda `$: All the nonfolded galleries are LS-galleries.
###### Remark 5.9.
In definition 4.1 one can replace positive (respectively positively folded) by negative (respectively negatively folded), i.e. one gets $`m_s^{}(\sigma )`$ and $`n_s^{}(\sigma )`$ for each negatively folded gallery $`\sigma `$. With the obvious changes this yields polynomials $`L_\sigma ^{}`$ nonzero only for negatively folded galleries. Going further, one gets $`\mathrm{\Gamma }_t^{}(A,B)`$, $`L_t^{}(A,B)`$ and recursions (using the same notations as in 5.4)
$$L_t^{}^{}(A,Bs)=\{\begin{array}{cc}L_t^{}(A,B)\hfill & \text{ if }BBs\hfill \\ q_sL_t^{}(A,B)+(q_s1)L_t^{}(A,Bs)\hfill & \text{ if }BBs.\hfill \end{array}$$
Since $`\overline{T}_s=q_s^1(T_s+(1q_s)T_{id})`$ for $`sS^𝔞`$ we get from lemma 5.2 that
$$X_A\overline{T}_s=\{\begin{array}{cc}X_{As}+(q_s^11)\hfill & \text{ if }AAs\hfill \\ q_s^1X_{As}\hfill & \text{ if }AAs\hfill \end{array}$$
for any $`A\stackrel{~}{𝒜}`$ and $`sS^𝔞`$. Under the hypotheses of theorem 5.5 we get
$$X_A\overline{T}_v=\underset{B𝒜}{}\overline{L_t^{}(A,B)}X_B.$$
If one defines
$$L_t^{}(\mu )=\underset{\sigma \mathrm{\Gamma }_t^{}(\mu )}{}q_{\iota (\sigma )}L_\sigma ^{}$$
we also can express the $`L_{\lambda \mu }`$ with negatively folded galleries. For this note that left multiplication by $`w_0`$ on $`\stackrel{~}{𝒜}`$ induces a type preserving bijection $`\varphi :\mathrm{\Gamma }_t^+\mathrm{\Gamma }_t^{}`$ for any type $`t`$. We have $`L_{\varphi (\sigma )}^{}=L_\sigma `$ and $`\iota (\varphi (\sigma ))=w_0\iota (\sigma )`$. In particular, we get the equality $`L_t(\mu )=L_t^{}(w_0\mu )`$. Combining this with the semi-invariance of the $`L_{\lambda \mu }`$ with respect to $`\mu `$ we get
$$L_{\lambda \mu }=q_{\mu w_0\mu }L_{\lambda ,w_0\mu }=\frac{q_\mu ^2}{W_\lambda (q)}L_{t^\lambda }(w_0\mu )=\frac{q_\mu ^2}{W_\lambda (q)}L_{t^\lambda }^{}(\mu )$$
which gives an expression of $`L_{\lambda \mu }`$ in terms of negatively folded galleries by the definition of $`L_{t^\lambda }^{}(\mu )`$.
## 6 Structure constants
In this section we calculate the structure constants of the spherical Hecke algebra with respect to the Macdonald basis and prove theorem 4.10 and thus theorem 1.3 and its corollary.
###### Lemma 6.1.
Let $`A=\mu +A_w`$ be a dominant generalized alcove such that $`As`$ is no longer dominant. Let $`H_{\alpha _i,0}`$ be the hyperplane separating $`A`$ and $`As`$. Then we have $`X_AT_s=T_{s_i}X_A`$.
###### Proof.
We have $`s_iA=As`$ and $`AAs`$. So $`s_i`$ and $`s`$ are conjugate in $`\stackrel{~}{W}^𝔞`$ and thus $`q_{s_i}=q_s`$. Distinguish two cases:
If $`s=s_{\theta ,1}`$ with $`\theta \mathrm{\Theta }`$ we have $`\alpha _i,\mu =1`$ and thus $`s_i(\mu )=\mu \alpha _i^{}`$ and $`s_iA=s_i(\mu )+A_{s_iw}`$. But on the other hand we have $`As=(\mu +w\theta _k^{})+A_{ws_\theta }`$ and so $`w\theta ^{}=\alpha _i^{}`$. In particular, $`s_iw<w`$. From \[Lus89, lemma 2.7(d) and proposition 3.6\] we know that $`q_{\alpha _i^{}}=q_s`$ in this case and
$$T_{s_i}X_\mu =X_{\mu \alpha _i^{}}T_{s_i}+(q_{s_i}1)X_\mu .$$
Together with $`s_iw<w`$ this yields
$$T_{s_i}X_\mu \overline{T}_w=X_{\mu \alpha _i^{}}\overline{T}_{s_iw}+(q_{s_i}1)X_\mu \overline{T}_w$$
and thus
$$T_{s_i}X_A=X_{A_s}+(q_{s_i}1)X_A=X_AT_s$$
where the last equality follows from $`AAs`$.
If $`s=s_jS`$ we have $`s_i(\mu )=\mu `$ and $`w^1(\alpha _i)=\alpha _j`$. So here $`s_iw>w`$. Using $`T_{s_i}X_\mu =X_\mu T_{s_i}`$ one obtains the desired equality as above. ∎
We keep the notation of the last lemma and get $`\mathrm{𝟏}_0X_AT_s=\mathrm{𝟏}_0T_{s_i}X_A=q_s\mathrm{𝟏}_0X_A`$ (recall that $`q_{s_i}=q_s`$). For a generalized alcove $`A`$ and a type $`t`$ define $`\mathrm{\Gamma }_{t,A}^+`$ to be the set of all positively folded galleries of type $`t`$ with initial alcove $`A`$.
Let $`t=(t_1,\mathrm{},t_k)`$ be a type and define $`T_t=T_{t_1}\mathrm{}T_{t_k}`$. From theorem 5.5 we get
$$X_AT_t=\underset{\sigma \mathrm{\Gamma }_{t,A}^+}{}L_\sigma X_{e(\sigma )}$$
where $`e(\sigma )`$ is the ending of $`\sigma `$ as introduced in definition 4.1. This yields
$$\mathrm{𝟏}_0X_AT_t=\underset{\sigma \mathrm{\Gamma }_{t,A}^+}{}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}.$$
Setting $`t^{}=(s,t)`$ we obtain by the same arguments
$$\mathrm{𝟏}_0X_AT_sT_t=\underset{\sigma \mathrm{\Gamma }_{t^{},A}^+}{}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}.$$
Since $`\mathrm{𝟏}_0X_AT_sT_t=q_s\mathrm{𝟏}_0X_AT_t`$ we get the following
###### Lemma 6.2.
Let $`t`$ be any type and let $`A`$ be a dominant generalized alcove such that $`As`$ is no longer dominant. Setting $`t^{}=(s,t)`$ we have
$$q_s\underset{\sigma \mathrm{\Gamma }_{t,A}^+}{}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}=\underset{\sigma \mathrm{\Gamma }_{t^{},A}^+}{}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}.$$
Now let $`\lambda X_+^{}`$. Then the generalized alcove $`A:=\lambda +A_w`$ is dominant iff $`w^1W^\lambda `$. Let $`w^1W^\lambda `$ and $`vW_\lambda `$. Since $`\overline{T}_vX_\lambda =X_\lambda \overline{T}_v`$ we get
$$\overline{T}_vX_A=q_v^1X_{\lambda +A_{vw}}=q_v^1X_{vA}.$$
Since $`vW_\lambda `$ we get the equality (using the notation introduced before the last lemma)
$$\mathrm{𝟏}_0X_{vA}T_t=q_v\mathrm{𝟏}_0\overline{T}_vX_AT_t=\mathrm{𝟏}_0X_AT_t.$$
For later use observe that $`vA=\lambda +A_{vw}`$ and thus $`vA`$ is no longer dominant. We get
###### Lemma 6.3.
Let $`\lambda X_+^{}`$, $`w^1W^\lambda `$ and $`vW_\lambda `$. Let $`A=\lambda +A_w`$. For any type $`t`$ we have
$$\underset{\sigma \mathrm{\Gamma }_{t,A}^+}{}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}=\underset{\sigma \mathrm{\Gamma }_{t,vA}^+}{}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}.$$
Now let $`\lambda ,\mu X_+^{}`$. Let $`w_\mu W_\mu `$ and $`n^\mu \tau _\mu W`$ as in definition 2.2. Let $`t^\mu `$ denote the type of a minimal gallery from $`A_f`$ to $`A_{n^\mu }`$. As in the proof of lemma 5.7 we get
$`\mathrm{𝟏}_0X_\lambda \mathbf{\hspace{0.33em}1}_0T_{n^\mu }`$ $`=\mathrm{𝟏}_0X_\lambda {\displaystyle \underset{\sigma \mathrm{\Gamma }_{t^\mu }^+}{}}q_{w_0\iota (\sigma )}L_\sigma X_{e(\sigma )}`$ (6.1)
$`=q_\lambda {\displaystyle \underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^+}{}}q_{w_0\iota (\sigma )}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}.`$ (6.2)
Here $`\mathrm{\Gamma }_{t^\mu ,\lambda }^+`$ is the set of all galleries of type $`t^\mu `$ starting in $`\lambda `$ and the last equality holds since translating a gallery $`\sigma `$ by $`\lambda `$ does not change $`L_\sigma `$. So we have an expansion for the product in terms of $`X_A`$ for $`A\stackrel{~}{𝒜}`$. But we need the expansion in terms of $`X_A`$ for dominant $`A`$ to compute the structure constants.
###### Theorem 6.4.
For $`\lambda ,\mu X_+^{}`$ we have
$$\mathrm{𝟏}_0X_\lambda \mathrm{𝟏}_0T_{n^\mu }=q_\lambda W_\lambda (q^1)\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d}{}q_{w_0\iota (\sigma )}C_\sigma \mathrm{𝟏}_0X_{e(\sigma )}.$$
###### Proof.
For the proof of this theorem we use lemmas 6.2 and 6.3 to show that the contribution of the galleries with non-dominant weights in the formula (6.2) is exactly the contribution of the $`p_s`$.
First assume $`\lambda `$ is regular. Then the first generalized alcove of every gallery starting in $`\lambda `$ is dominant. Let $`\eta \mathrm{\Gamma }_{t^\mu ,\lambda }^+`$ be a gallery leaving the dominant chamber. Let $`\gamma `$ be the maximal initial subgallery of $`\eta `$ contained in $`𝒞`$ and let $`A`$ be $`e(\gamma )`$. Then $`\eta `$ is not folded after $`A`$ and the next generalized alcove in $`\eta `$ is of the form $`As`$ for some $`sS^𝔞`$. Denote by $`\mathrm{\Gamma }_\gamma ^+\mathrm{\Gamma }_{t^\mu ,\lambda }^+`$ the set of galleries starting with $`\gamma `$. By lemma 6.2 we have that
$$\frac{q_s}{q_s1}\underset{\sigma \mathrm{\Gamma }_\gamma ^+,\sigma \text{ folded at }A}{}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}=\underset{\sigma \mathrm{\Gamma }_\gamma ^+}{}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}.$$
So the contribution of all galleries starting with $`\gamma `$ is the same as the contribution of the galleries starting with $`\gamma `$ and staying in $`𝒞`$ at $`A`$, if the contribution of the folding at $`A`$ is $`q_s`$ instead of $`q_s1`$. Iteration of this procedure eventually yields
$$\underset{\sigma \mathrm{\Gamma }_\gamma ^+}{}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}=\underset{\sigma \mathrm{\Gamma }_\gamma ^+,\sigma 𝒞}{}C_\sigma \mathrm{𝟏}_0X_{e(\sigma )}$$
which proves the theorem for regular $`\lambda `$.
If $`\lambda `$ is non-regular we have to apply lemma 6.3 to obtain the theorem because in this case the first alcove of a gallery starting in $`\lambda `$ can be non-dominant. In this case its contribution has a part coming from the initial direction, which we did not need to consider in the regular case. But lemma 6.3 tells us that the contribution arising from these alcoves is the same as the contribution from the dominant ones. More precisely we have for $`w^1W^\lambda `$ and $`vW_\lambda `$
$$\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^+,\iota (\sigma )=w}{}q_{w_0w}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}=q_v\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^+,\iota (\sigma )=vw}{}q_{w_0vw}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}$$
and thus
$$W_\lambda (q^1)\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^+,\iota (\sigma )=w}{}q_{w_0w}L_\sigma \mathrm{𝟏}_0X_{e(\sigma )}=\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^+,\iota (\sigma )W_\lambda w}{}q_{w_0\iota (\sigma )}\mathrm{𝟏}_0X_{e(\sigma )}.$$
Since the sum over all $`w^1W^\lambda `$ of the left hand side of the last equation is exactly the contribution of the galleries starting in $`𝒞`$, the theorem follows. ∎
###### Remark 6.5.
The proofs for multiplying Schur polynomials using paths are of a similar type as above (see for example \[Lit94, section 6\]). First one gets a formula involving also Schur polynomials associated to paths leaving the dominant chamber. Then one shows that the contributions of the leaving paths cancel each other. This is done by combinatorial arguments, i.e. one can see which paths cancel each other. In contrast to this we do not have any concrete information about this cancellation process.
Now we can prove the first part of theorem 4.10 respectively theorem 1.3. We multiply the equation of the last theorem from the right by $`\mathrm{𝟏}_0`$ and get by the definition of the Macdonald basis
$`M_\lambda M_\mu `$ $`={\displaystyle \frac{q_\lambda q_{w_0}^1}{W(q)W_\lambda (q^1)}}{\displaystyle \frac{1}{W(q)W_\mu (q)}}\mathrm{𝟏}_0X_\lambda \mathrm{𝟏}_0\mathbf{\hspace{0.17em}1}_0T_{n^\mu }\mathrm{𝟏}_0`$
$`={\displaystyle \frac{q_\lambda ^2q_{w_0}^1}{W(q)W_\mu (q)}}{\displaystyle \underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d}{}}q_{w_0\iota (\sigma )}C_\sigma \mathrm{𝟏}_0X_{e(\sigma )}\mathrm{𝟏}_0`$
$`={\displaystyle \frac{q_\lambda ^2q_{w_0}^1}{W(q)W_\mu (q)}}{\displaystyle \underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d}{}}q_{wt(\sigma )}q_{w_0\iota (\sigma )}C_\sigma \mathrm{𝟏}_0X_{wt(\sigma )}\mathrm{𝟏}_0`$
$`={\displaystyle \frac{q_\lambda ^2}{W_\mu (q)}}{\displaystyle \underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d}{}}q_{wt(\sigma )}^2q_{w_0\iota (\sigma )}C_\sigma W_{wt(\sigma )}(q^1)M_{wt(\sigma )}`$
$`={\displaystyle \frac{q_\lambda ^2}{W_\mu (q)}}{\displaystyle \underset{\nu X_+^{}}{}}q_\nu ^2W_\nu (q^1)C_{t^\mu ,\lambda }(\nu )M_\nu .`$
To prove the second part of theorem 4.10 and thus theorem 1.3 we need one more step. It is not possible to impose conditions on the initial direction as in theorem 4.4. Instead we impose conditions on the final direction to get rid of the fraction $`\frac{1}{W_\mu (q)}`$. For doing this we need some preparation. The situation is more difficult than the case of Satake coefficients since now two stabilizers instead of one are involved. So we first need some information on the interplay between them.
We use the notation for stabilizer subgroups introduced in definition 2.2. Moreover, for any $`\nu X^{}`$ let $`\mathrm{𝟏}_\nu =_{wW_\nu }T_w`$ be the corresponding symmetrizer. Note that $`W_{w_0\mu }=w_0W_\mu w_0`$ and thus $`q_{w_\mu }=q_{w_{w_0\mu }}`$ and $`W_\mu (q)=W_{w_0\mu }(q)`$.
Let $`Y=_{wW}R_w\overline{T}_w\stackrel{~}{}^𝔞`$ with $`R_w`$. Assume $`Y\stackrel{~}{}^𝔞\mathrm{𝟏}_{w_0\mu }`$. Then $`R_w=R_{wv}`$ for any $`wW`$ and $`vW_{w_0\mu }`$ and thus
$$Y=q_{w_\mu }^1\underset{wW^{w_0\mu }}{}R_w\overline{T}_w\mathrm{𝟏}_{w_0\mu }$$
(6.3)
since for $`wW^{w_0\mu }`$ we have $`\overline{T}_w\overline{\mathrm{𝟏}_{w_0\mu }}=_{vW_{w_0\mu }}\overline{T}_{wv}`$ and $`\overline{\mathrm{𝟏}_{w_0\mu }}=q_{w_\mu }^1\mathrm{𝟏}_{w_0\mu }`$.
Let $`\nu X_+^{}`$ and take $`Y`$ of a special form, namely $`Y=_{w^1W^\nu }R_w\mathrm{𝟏}_\nu \overline{T}_w`$. For $`wW`$ denote by $`w^\nu `$ the minimal element of the coset $`W_\nu w`$. In particular $`(w^\nu )^1W^\nu `$. Expanding $`Y`$ in terms of the $`\overline{T}_w`$ yields
$$Y=q_{w_\nu }\underset{wW}{}R_{w^\nu }\overline{T}_w.$$
So if in addition $`Y\stackrel{~}{}^𝔞\mathrm{𝟏}_{w_0\mu }`$ we get $`Y=q_{w_\nu }q_{w_\mu }^1_{wW^{w_0\mu }}R_{w^\nu }\overline{T}_w\mathrm{𝟏}_{w_0\mu }`$ by the considerations above.
We calculate $`Y\mathrm{𝟏}_0`$ and get $`Y\mathrm{𝟏}_0=q_{w_\nu }q_{w_\mu }^1W_\mu (q)_{wW^{w_0\mu }}q_w^1R_{w^\nu }\mathrm{𝟏}_0`$. Thus
$`Y\mathrm{𝟏}_0=q_{w_\nu }W_\mu (q^1){\displaystyle \underset{w^1W^\nu }{}}q_w^1F_{\mu \nu }^wR_w\mathrm{𝟏}_0`$ (6.4)
where $`F_{\mu \nu }^w:=q_w_{vW^{w_0\mu }W_\nu w}q_v^1`$. Observe that $`W^{w_0\mu }W_\nu w\mathrm{}`$ iff $`wW_\nu W^{w_0\mu }`$. In particular, we get for regular $`\nu `$ that $`F_{\mu \nu }^w=1`$ if $`wW^{w_0\mu }`$ and 0 else.
Now we relate this to our problem. We have
$$W_\mu (q)\mathrm{𝟏}_0T_{n^\mu }=\mathrm{𝟏}_0\mathrm{𝟏}_\mu T_{\tau _\mu }T_{w_\mu }\overline{T}_{w_0}=\mathrm{𝟏}_0T_{\tau _\mu }T_{w_\mu }\mathrm{𝟏}_\mu \overline{T}_{w_0}=\mathrm{𝟏}_0T_{n^\mu }T_{w_0}\mathrm{𝟏}_\mu \overline{T}_{w_0}.$$
But $`T_{w_0}T_w\overline{T}_{w_0}=T_{w_0ww_0}`$ for all $`wW`$ and thus $`T_{w_0}\mathrm{𝟏}_\mu \overline{T}_{w_0}=\mathrm{𝟏}_{w_0\mu }`$. So $`\mathrm{𝟏}_0T_{n^\mu }\stackrel{~}{}^𝔞\mathrm{𝟏}_{w_0\mu }`$.
We have $`W_\nu (q)\mathrm{𝟏}_0X_\nu =\mathrm{𝟏}_0X_\nu \mathrm{𝟏}_\nu `$ since $`T_wX_\nu =X_\nu T_w`$ for any $`wW_\nu `$. So the contribution of $`\mathrm{𝟏}_0X_\nu `$ in theorem 6.4 is given by
$$\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d(\nu )}{}q_{w_0\iota (\sigma )}C_\sigma \mathrm{𝟏}_0X_{e(\sigma )}=\frac{q_\nu }{W_\nu (q)}\mathrm{𝟏}_0X_\nu \underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d(\nu )}{}q_{w_0\iota (\sigma )}q_{\epsilon (\sigma )}C_\sigma \mathrm{𝟏}_\nu \overline{T}_{\epsilon (\sigma )}.$$
As already observed before, $`\nu +A_v𝒞`$ with $`vW`$ iff $`v^1W^\nu `$. So the final directions of the galleries $`\sigma `$ occurring in the last equation satisfy $`(\epsilon (\sigma ))^1W^\nu `$. If we define $`Y:=_{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d(\nu )}q_{w_0\iota (\sigma )}q_{\epsilon (\sigma )}C_\sigma \mathrm{𝟏}_\nu \overline{T}_{\epsilon (\sigma )}`$ then $`Y`$ is of the kind considered above and $`Y\stackrel{~}{}^𝔞\mathrm{𝟏}_{w_0\mu }`$. So we can apply (6.4) and get
$$\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d(\nu )}{}q_{w_0\iota (\sigma )}q_{\epsilon (\sigma )}C_\sigma \mathrm{𝟏}_\nu \overline{T}_{\epsilon (\sigma )}\mathrm{𝟏}_0=q_{w_\nu }W_\mu (q^1)\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d(\nu )}{}q_{w_0\iota (\sigma )}C_\sigma F_{\mu \nu }^{\epsilon (\sigma )}\mathrm{𝟏}_0.$$
Bringing all this together we can multiply the assertion of theorem 6.4 from the right by $`\mathrm{𝟏}_0`$ and get
$$\mathrm{𝟏}_0X_\lambda \mathrm{𝟏}_0T_{n^\mu }\mathrm{𝟏}_0=q_\lambda W_\lambda (q^1)W_\mu (q^1)\underset{\nu X_+^{}}{}\frac{q_\nu }{W_\nu (q^1)}\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d(\nu )}{}q_{w_0\iota (\sigma )}C_\sigma F_{\mu \nu }^{\epsilon (\sigma )}\mathrm{𝟏}_0X_\nu \mathrm{𝟏}_0.$$
Now we can calculate the coefficient of $`M_\nu `$ in the product $`M_\lambda M_\mu `$ as above. It is equal to
$$q_{\lambda \nu }^2q_{w_\mu }^1\underset{\sigma \mathrm{\Gamma }_{t^\mu ,\lambda }^d(\nu )}{}q_{w_0\iota (\sigma )}C_\sigma F_{\mu \nu }^{\epsilon (\sigma )}$$
which proves the second part of theorem 4.10 and thus 1.3.
###### Remark 6.6.
Consider the case of equal parameters and let $`w^1W^\nu `$. Then we have $`F_{\mu \nu }^w=q^{l(w)}_{vW^{w_0\mu }W_\nu w}q^{l(v)}`$. By definition of $`W^\nu `$ we have $`l(v)l(w)`$ for all $`vW_\nu w`$ and thus $`F_{\mu \nu }^w^{}`$. Moreover, the constant term of $`F_{\mu \nu }^w`$ is 1 iff $`wW^{w_0\mu }`$.
###### Remark 6.7.
One can proceed the same way to obtain a formula for the Satake coefficients as in the second part of theorem 4.4 with a condition on the final direction. For stating the results we consider again the situation of section 5. So $`\lambda X_+^{}`$ and $`t^\lambda `$ is the type of a minimal gallery from $`A_f`$ to $`A_{n^\lambda }`$. Applying the above considerations (for $`\lambda `$ instead of $`\mu `$) yields $`\mathrm{𝟏}_0T_{n^\lambda }\stackrel{~}{}^𝔞\mathrm{𝟏}_{w_0\lambda }`$. A formula for $`\mathrm{𝟏}_0T_{n^\lambda }`$ is given by (see the proof of lemma 5.7) $`_{\sigma \mathrm{\Gamma }_{t^\lambda }^+}q_{wt(\sigma )}q_{\epsilon (\sigma )}q_{w_0\iota (\sigma )}L_\sigma X_{wt(\sigma )}\overline{T}_{\epsilon (\sigma )}`$. So we get by (6.3)
$$\mathrm{𝟏}_0T_{n^\lambda }=q_{w_\lambda }^1\underset{\begin{array}{c}\sigma \mathrm{\Gamma }_{t^\lambda }^+\\ \epsilon (\sigma )W^{w_0\lambda }\end{array}}{}q_{wt(\sigma )}q_{\epsilon (\sigma )}q_{w_0\iota (\sigma )}L_\sigma X_{wt(\sigma )}\overline{T}_{\epsilon (\sigma )}\mathrm{𝟏}_{w_0\lambda }.$$
Multiplying by $`\mathrm{𝟏}_0`$ from the right then yields
$$M_\lambda =\frac{q_{w_\lambda }^1}{W(q)}\underset{\mu X^{}}{}q_\mu \underset{\begin{array}{c}\sigma \mathrm{\Gamma }_{t^\lambda }^+(\mu )\\ \epsilon (\sigma )W^{w_0\lambda }\end{array}}{}q_{w_0\iota (\sigma )}L_\sigma X_\mu \mathrm{𝟏}_0$$
and thus $`L_{\lambda \mu }=q_{w_\lambda }^1_{\begin{array}{c}\sigma \mathrm{\Gamma }_{t^\lambda }^+(\mu )\\ \epsilon (\sigma )W^{w_0\lambda }\end{array}}q_{w_0\iota (\sigma )}L_\sigma `$. Moreover, we see that for a LS-gallery $`\sigma `$ we have $`\epsilon (\sigma )W^{w_0\lambda }`$.
## 7 Commutation formula
Here we give two other applications of theorem 5.5. The first is a nice combinatorial description of the base change matrix between the standard basis $`\{T_w\}_{w\stackrel{~}{W}^𝔞}`$ and the basis $`\{X_\lambda \overline{T}_w\}_{\mu X^{},wW}`$.
###### Corollary 7.1.
Let $`v\stackrel{~}{W}^𝔞`$ and fix some minimal gallery of type $`t`$ connecting $`A_f`$ and $`A_v`$. Then
$$T_v=\underset{\sigma \mathrm{\Gamma }_t^+,\iota (\sigma )=id}{}L_\sigma X_{e(\sigma )}=\underset{\sigma \mathrm{\Gamma }_t^+,\iota (\sigma )=id}{}q_{wt(\sigma )}^1q_{\epsilon (\sigma )}L_\sigma X_{wt(\sigma )}\overline{T}_{\epsilon (\sigma )}.$$
Now let $`wW^\lambda `$ and recall the notation from definition 2.2. Then $`l(ww_\lambda )=l(w)+l(w_\lambda )`$ and thus $`T_{ww_\lambda }=T_wT_{w_\lambda }`$. We also have $`T_{\tau _\lambda }=T_{n^\lambda }T_{w_0}\overline{T}_{w_\lambda }`$ by lemma 2.3. Since $`T_{w_\lambda }X_\lambda =X_\lambda T_{w_\lambda }`$ we get
$$T_wT_{w_\lambda }X_\lambda =T_wX_\lambda T_{w_\lambda }=q_\lambda ^1T_wT_{n^\lambda }T_{w_0}\overline{T}_{w_\lambda }T_{w_\lambda }=q_\lambda ^1T_{wn^\lambda }T_{w_0}.$$
Applying corollary 7.1 to $`v=wn^\lambda `$ and using $`\overline{T}_yT_{w_0}=T_{yw_0}`$ for $`yW`$ we get a $`q`$-analog of a commutation formula of Pittie and Ram in the nil-affine Hecke algebra \[PR99\].
###### Corollary 7.2.
Let $`\lambda X_+^{}`$ and $`wW^\lambda `$. Let $`t`$ be the type of a minimal gallery connecting $`A_f`$ and $`A_{wn^\lambda }`$. Then
$$T_{ww_\lambda }X_\lambda =q_\lambda ^1T_{wn^\lambda }T_{w_0}=\underset{\sigma }{}q_{wt(\sigma )+\lambda }^1q_{\epsilon (\sigma )}L_\sigma X_{wt(\sigma )}T_{\epsilon (\sigma )w_0}$$
where the sum is over all $`\sigma \mathrm{\Gamma }_t^+`$ starting in $`A_f`$.
## 8 Geometric interpretation
In this section we reformulate results of \[BD94\] using galleries to show some relations of our combinatorics to geometry. We show that this geometrical interpretation is compatible with the geometrical interpretation of the $`L_{\lambda \mu }`$ arising from the geometric definition of the Satake isomorphism (see \[HKP03\]). Since there are no new results in this section we are rather sketchy.
Assume that $`X^{}=Q^{}`$. So $`\stackrel{~}{W}^𝔞=W^𝔞`$ and generalized alcoves are alcoves. We use negatively folded galleries as in remark 5.9. We regard the affine Hecke algebra $`\stackrel{~}{}^𝔞`$ specialized at some prime power $`𝐪`$ and all polynomials evaluated at $`𝐪`$.
Details of the following constructions and their relation to affine Kac-Moody algebras can be found in Kumar’s book \[Kum02\]. We use the notation from the introduction, i.e. $`G`$ is a semisimple, simply connected algebraic group over $`K`$ with root datum $`\mathrm{\Phi }`$ associated to a Borel $`BG`$ and a maximal torus $`TB`$ and all groups are defined and split over $`𝔽_𝐪`$. Let $`𝒦=K((t))`$ be the field of Laurent series and denote by $`𝒪=K[[t]]𝒦`$ the ring of formal power series. The evaluation at $`0`$ induces a map $`ev:𝒪K,t0`$. This induces a morphism of groups $`ev:G(𝒪)G`$. Define $`=ev^1(B)`$. Further we set $`𝒢=G(𝒦)`$ and let $`𝒩𝒢`$ be the normalizer of $`T`$ in $`𝒢`$. Then $`(𝒢,,𝒩,T(𝒪))`$ is a Tits system with Weyl group $`W^𝔞`$.
From the theory of Tits system one has a Bruhat decomposition $`𝒢=_{wW^𝔞}𝒰_ww`$ with $`𝒰_w`$ isomorphic to $`𝔸_K^{l(w)}`$. On the other hand there is the Iwasawa decomposition $`𝒢=_{wW^𝔞}U(𝒦)w`$ where $`UB`$ is the unipotent radical. These decompositions are compared by Billig and Dyer in \[BD94\].
###### Theorem 8.1 (\[BD94\]).
For $`wW^𝔞`$ and $`s=s_iS^𝔞`$ one has
$$U(𝒦)ws=\{\begin{array}{cc}U(𝒦)ws\hfill & \text{if }A_wA_{ws}\hfill \\ U(𝒦)wU(𝒦)ws\hfill & \text{if }A_wA_{ws}.\hfill \end{array}$$
Let $`wW^𝔞`$ and let $`\sigma `$ be a minimal gallery of type $`t=(t_1,\mathrm{},t_k)`$ which connects $`A_f`$ and $`A_w`$. Define a map $`\eta :𝒰_w\mathrm{\Gamma }_t^{}`$ by
$$\eta (u)=(A_f,A_{w_1},\mathrm{},A_{w_k})\text{ iff }ut_1\mathrm{}t_jU(𝒦)w_j\text{ for }j\{1,\mathrm{},k\}.$$
It follows from the last theorem that $`\eta `$ is well defined.
The affine flag variety is the quotient $`𝒢/`$. It is the set of closed points of an ind-variety defined over $`𝔽_𝐪`$ such that the Bruhat cells $`w`$ for $`wW^𝔞`$ are isomorphic to affine spaces of dimension $`l(w)`$. For $`Y𝒢`$ denote by $`Y`$ its image in $`𝒢/`$. For a negatively folded gallery $`\sigma `$ denote by $`m^{}(\sigma )`$ the total number of negative directions and by $`n^{}(\sigma )`$ the total number of negative folds. Then one has
###### Theorem 8.2 (\[BD94\]).
1. If $`\sigma \mathrm{\Gamma }_t^{}`$ then $`\eta ^1(\sigma )K^{m^{}(\sigma )}\times (K^{})^{n^{}(\sigma )}`$.
2. If $`vW^𝔞`$ then $`wU(𝒦)v=_{\sigma \mathrm{\Gamma }_t^{}(A_f,A_v)}\eta ^1(\sigma )w`$.
For $`Y𝒢/`$ denote by $`|Y|`$ the number of $`𝔽_𝐪`$-rational points. Let $`L_w^{}(A_f,A_v)`$ be the analog of definition 5.6 for negatively folded galleries. From the last theorem we get
###### Corollary 8.3.
If $`v,wW^𝔞`$ then $`|wU(𝒦)v|=L_w^{}(A_f,A_v)`$.
###### Remark 8.4.
Looking at positively folded galleries starting in $`A_f`$ one can calculate intersections $`^{}wU^{}(𝒦)v`$ in the same way as in corollary 8.3. Here $`^{}`$ is obtained from the opposite Borel $`B^{}`$ in the same way as $``$ from $`B`$ and $`U^{}`$ is the unipotent radical of $`B`$.
As already mentioned in the introduction the definition of the *geometric* Satake isomorphism yields a geometric interpretation of the $`L_{\lambda \mu }`$. For stating this we have to consider intersections in the affine Grassmanian $`𝒢/G(𝒪)`$. For $`\lambda ,\mu X^{}`$ let $`X_{\lambda \mu }^{}:=G(𝒪)\tau _\lambda G(𝒪)U^{}(𝒦)\tau _\mu G(𝒪)`$. Then $`L_{\lambda \mu }=|X_{\lambda \mu }^{}|`$. This interpretation can also be recovered using the last corollary as follows.
The group $`G(𝒪)`$ is the parabolic subgroup of $`𝒢`$ associated to the classical Weyl group $`WW^𝔞`$, i.e. $`G(𝒪)=_{wW}w`$. Let $`\pi :𝒢/𝒢/G(𝒪)`$ be the canonical projection. From general theory of Tits systems one knows that
$$\pi _{|w}:wwG(𝒪)$$
is an isomorphism iff $`w`$ is of minimal length in $`wW`$. So we get an isomorphism
$$\pi _{|_{wW^\lambda }wn^\lambda }:\underset{wW^\lambda }{}wn^\lambda G(𝒪)\tau _\lambda G(𝒪).$$
If $`\mu Q^{}`$ then
$$\pi ^1(U(𝒦)\tau _\mu G(𝒪))=\underset{wW}{}U(𝒦)\tau _\mu w.$$
Bringing this together and setting $`X_{\lambda \mu }=G(𝒪)\tau _\lambda G(𝒪)U(𝒦)\tau _\mu G(𝒪)`$ we get
$`|X_{\lambda \mu }|`$ $`={\displaystyle \underset{\begin{array}{c}wW^\lambda \\ vW\end{array}}{}}|wn^\lambda U(𝒦)\tau _\mu v|`$
$`={\displaystyle \underset{\begin{array}{c}wW^\lambda \\ vW\end{array}}{}}L_{wn^\lambda }^{}(A_f,\tau _\mu v)={\displaystyle \underset{wW^\lambda }{}}L_{wn^\lambda }^{}(A_f,\mu )`$
$`={\displaystyle \frac{1}{W_\lambda (𝐪)}}L_{t^\lambda }^{}(\mu )=𝐪^{2\rho ,\mu }L_{\lambda \mu }.`$
The last equalities follow from the analog of the second statement in lemma 5.7 for negatively folded galleries. Using the fact that $`X_{\lambda \mu }^{}=w_0X_{\lambda ,w_0\mu }`$ and $`L_{\lambda ,\mu }=𝐪^{\rho ,\mu w_0\mu }L_{\lambda ,w_0\mu }=𝐪^{2\rho ,\mu }L_{\lambda ,w_0\mu }`$ we get as in \[GL05\] that
$$|X_{\lambda \mu }^{}|=\underset{\sigma \mathrm{\Gamma }_{t^\lambda }^+(\mu ),\iota (\sigma )W^\lambda }{}𝐪^{l(w_0\iota (\sigma ))}L_\sigma .$$
It is well known that the affine Hecke algebra $`\stackrel{~}{}^𝔞`$ specialized at $`𝐪`$ can be interpreted as the convolution algebra of $`_𝐪`$-invariant functions with finite support on the affine flag variety $`𝒢_𝐪/_𝐪`$ (see for example \[HKP03\]). In this setting the generator $`T_w`$ corresponds to the characteristic function on $`_𝐪w_𝐪`$. Using theorem 8.1 one can give a similar interpretation for the alcove basis.
Let $`t_wF`$ be the characteristic function on $`U^{}(𝒦_𝐪)w_𝐪`$ (which in general does not have finite support) and let $`M`$ be the subspace spanned by all $`t_w`$. Using theorem 8.1 one can show that $`M`$ is closed under right convolution by $`\stackrel{~}{}^𝔞`$. Moreover, for $`wW^𝔞`$ and $`sS^𝔞`$ one gets
$$t_wT_s=\{\begin{array}{cc}t_{ws}\hfill & \text{if }A_wA_{ws}\hfill \\ 𝐪t_{ws}+(𝐪1)t_w\hfill & \text{if }A_wA_{ws}.\hfill \end{array}$$
So by lemma 5.2 the map
$`M`$ $`\stackrel{~}{}^𝔞`$
$`t_v`$ $`𝐪^{\rho ,\mu (A_v)}X_{wt(A_v)}\overline{T}_{\delta (A_v)}=𝐪^{2\rho ,\mu (A_v)l(\delta (A_v))}X_{A_v}`$
is an isomorphism of right $`\stackrel{~}{}^𝔞`$-modules.
###### Remark 8.5.
In \[BD94\] the cited results are shown for any Kac-Moody group and any generalized system of positive roots. Theorem 8.1 is then formulated with distinguished subexpressions instead of folded galleries. All facts also follow quite immediately using the methods from \[GL05\] in this case.
The explicit formulas for the action of $`\stackrel{~}{}^𝔞`$ on $`M`$ are known (at least to experts). It is essentially the action of $`\stackrel{~}{}^𝔞`$ on the space of Iwahori fixed vectors of the universal unramified principal series. Corollary 8.3 follows from these explicit formulas.
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# Mutually unbiased phase states, phase uncertainties, and Gauss sums
## I INTRODUCTION
In quantum mechanics, orthogonal bases of a Hilbert space $`_q`$ of finite dimension $`q`$ are mutually unbiased if inner products between all possible pairs of vectors of distinct bases are all equal to $`1/\sqrt{q}`$. Eigenvectors of ordinary Pauli spin matrices (i.e. in dimension $`q=2`$) provide the best known example. It has been shown that in dimension $`q=p^m`$ which is the power of a prime $`p`$, the complete sets of mutually unbiased bases (MUBs) result from Fourier analysis over a Galois field $`F_q`$ (in odd characteristic $`p`$)Woott89 or of Galois ring $`R_{4^m}`$(in even characteristic $`2`$)Klapp03 . InWoott04b ; Planat04 , one can find an exhaustive literature on MUBs . Complete sets of MUBs have an intrinsic geometrical interpretation, and were related to discrete phase spacesWoott04b ; Paz ; Pittinger05 , finite projective planesSaniga ; Saniga2 , convex polytopes Bengtson , and complex projective $`2`$-designsBarnum02 ; Klap2 . There are hints on the relation to symmetric informationally complete positive operator measures (SIC-POVMs)Wootters04 ; Grassl04 ; Appleby04 ; Klap05 , and to Latin squaresWocjan04 .
There are strong motivations to embark on detailed studies of MUBs. First, they enter rigorous treatments of Bohr’s principle of complementarity that distinguishes between quantum and classical systems at the practical level of measurements. This fundamental quantum principle introduces the idea of complementary pairs of observables in the sense that precise measurement of one of them implies that possible outcomes of the other (when measured) are equally probable. In the nondegenerate case, if an observable $`O`$ represented by a $`q`$ times $`q`$ hermitian matrix is measured in a quantum system prepared in the eigenbase of its complementary counterpart $`O_c`$, then the probability to find the system in one of the eigenstates of $`O`$ is just $`1/q`$ as corresponding to mutually unbiased inner products. Another domain of applications where MUBs have been found to play an important role is the field of secure quantum key exchange (quantum cryptography). In the area of quantum state tomography, one should use MUBs for a complete reconstruction of an unknown quantum state quipro .
In this paper we approach the MUBs theory from the point of view of the theory of additive and multiplicative characters in Galois number field theory. The multiplicative characters $`\psi _k(n)=\mathrm{exp}(\frac{2i\pi nk}{q1})`$, $`k=0\mathrm{}q2`$, are well known since they constitute the basis for the ordinary discrete Fourier transform. But in order to construct MUBs, the additive characters introduced below are the ones which are useful. This construction is implicit in some previous papersWoott89 ; Klapp03 ; Planat04 , and is now being fully recognizedDurt04 ; Klimov04 .
An interesting consequence is the following: the discrete Fourier transform in $`𝒵_q`$ has been used by Pegg & Burnett Pegg89 as a definition of phase states $`|\theta _k`$, $`k=0\mathrm{}q1`$, in $`_q`$. The phase states $`|\theta _k`$ could be considered as eigenvectors of a properly defined Hermitian phase operator $`\mathrm{\Theta }_{PB}`$. Phase properties and phase fluctuations attached to particular field states were extensively described. In particular the classical phase variance $`\pi ^2/3`$ could be recovered.
We construct here a phase operator $`\mathrm{\Theta }_{\mathrm{Gal}}`$ having phase MUBs as eigenvectors. In contrast to the case of $`\mathrm{\Theta }_{PB}`$, we find that the phase fluctuations of $`\mathrm{\Theta }_{\mathrm{Gal}}`$ can be expressed in terms of Gauss sums over the finite number field $`F_q`$, and could be in principle smaller than those due to $`\mathrm{\Theta }_{PB}`$. This points to the fact that the phase MUBs may be of interest for quantum signal processing. Character sums and Gauss sums which are useful for optimal bases of $`m`$-qudits ($`p`$ odd) are also generalized to optimal bases of $`m`$-qubits ($`p=2`$).
## II Phase MUBs in odd prime characteristic
### II.1 Mathematical preliminaries
The key relation between Galois fields $`F_q`$ and MUBs is the theory of characters. This has not been recognized before and here we use the standpoint of characters as the most general way of considering previous results and also as a better criterium for elaborating on future results.
A Galois field is a finite set structure endowed with two group operations, the addition “$`+`$” and the multiplication “$``$”. The field $`F_q`$ can be represented as classes of polynomials obtained by computing modulo an irreducible polynomial over the ground field $`F_p=𝒵_p`$, the integers modulo $`p`$Lidl83 . A Galois field exists if and only if $`q=p^m`$. We also recall that $`F_q[x]`$ is the standard notation for the set of polynomials in $`x`$ with coefficients in $`F_q`$.
A character $`\kappa (g)`$ over an abelian group $`G`$ is a (continuous) map from $`G`$ to the field of complex numbers $`𝒞`$ of unit modulus, i.e. such that $`|\kappa (g)|=1`$, $`gG`$.
We start with a map from the extended field $`F_q`$ to the ground field $`F_p`$ which is called the trace function
$$tr(x)=x+x^p+\mathrm{}+x^{p^{m1}}F_p,xF_q.$$
(1)
Using (1), an additive character over $`F_q`$ is defined as
$$\kappa (x)=\omega _p^{tr(x)},\omega _p=\mathrm{exp}(\frac{2i\pi }{p}),xF_q.$$
(2)
The main property is that it satisfies $`\kappa (x+y)=\kappa (x)\kappa (y),x,yF_q`$.
On the other hand, the multiplicative characters are of the form
$$\psi _k(n)=\omega _{q1}^{nk},k=0\mathrm{}q2,n=0\mathrm{}q2.$$
(3)
In the present research, the construction of Galois phase MUBs will be related to character sums with polynomial arguments $`f(x)`$ also called Weil sumsKlapp03
$$W_f=\underset{xF_q}{}\kappa (f(x)).$$
(4)
In particular, ( theorem 5.38 in Lidl83 ), for a polynomial $`f_d(x)F_q[x]`$ of degree $`d1`$, with $`gcd(d,q)=1`$, one gets $`W_{f_d}(d1)q^{1/2}`$.
The quantum fluctuations arising from the phase MUBs will be found to be related to Gauss sums of the form
$$G(\psi ,\kappa )=\underset{xF_q^{}}{}\psi (x)\kappa (x),$$
(5)
where $`F_q^{}=F_q\{0\}`$. Using the notation $`\psi _0`$ for a trivial multiplicative character $`\psi =1`$, and $`\kappa _0`$ for a trivial additive character $`\kappa =1`$ the Gaussian sums (5) satisfy $`G(\psi _0,\kappa _0)=q1`$; $`G(\psi _0,\kappa )=1`$; $`G(\psi ,\kappa _0)=0`$ and $`|G(\psi ,\kappa )|=q^{1/2}`$ for nontrivial characters $`\kappa `$ and $`\psi `$.
### II.2 Galois quantum phase states
We now introduce a class of quantum phase states as a “Galois” discrete quantum Fourier transform of the Galois number kets
$$|\theta ^{(y)}=\frac{1}{\sqrt{q}}\underset{nF_q}{}\psi _k(n)\kappa (yn)|n,yF_q$$
(6)
in which the coefficient in the computational base $`\{|0,|1,\mathrm{},|q1\}`$ represents the product of an arbitrary multiplicative character $`\psi _k(n)`$ by an arbitrary additive character $`\kappa (yn)`$.
It is easy to show that previous basic results in this area can be obtained as particular cases of (6). Indeed:
Pegg & Barnett (1989): For $`\kappa =\kappa _0`$ and $`\psi \psi _k(n)`$, one recovers the ordinary quantum Fourier transform over $`𝒵_q`$. It has been shownPegg89 that the corresponding states
$$|\theta _k=\frac{1}{\sqrt{q}}\underset{n𝒵_q}{}\psi _k(n)|n,$$
(7)
are eigenstates of the Hermitian phase operator
$$\mathrm{\Theta }_{PB}=\underset{k𝒵_q}{}\theta _k|\theta _k\theta _k|,$$
(8)
with eigenvalues $`\theta _k=\theta _o+\frac{2\pi k}{q}`$, $`\theta _0`$ an arbitrary initial phase.
Wootters & Fields (1989): We recover the result of Wootters and Fields in a more general form by employing the Euclidean division theorem (see theorem 11.19 in Lidl98 ) for the field $`F_q`$, which says that given any two polynomials $`y`$ and $`n`$ in $`F_q`$, there exists a uniquely determined pair $`(a,b)F_q\times F_q`$, such that $`y=an+b`$, $`deg(b)<deg(a)`$. Using the decomposition of the exponent in (6), we obtain
$$|\theta _b^a=\frac{1}{\sqrt{q}}\underset{nF_q}{}\psi _k(n)\kappa (an^2+bn)|n,a,bF_q.$$
(9)
(The result of Wootters & Fields corresponds to the trivial multiplicative character $`\psi _0=1`$). Eq. (9) defines a set of $`q`$ bases (with index $`a`$) of $`q`$ vectors (with index $`b`$). Using Weil sums (4) it is easily shown that, for $`q`$ odd, so that $`gcd(2,q)=1`$, the bases are orthogonal and mutually unbiased to each other and to the computational base
$$|\theta _b^a|\theta _d^c|=|\frac{1}{q}\underset{nF_q}{}\omega _p^{tr((ca)n^2+(db)n}|=\{\begin{array}{cc}& \delta _{bd}\text{if}c=a(\text{orthogonality})\hfill \\ & \frac{1}{\sqrt{q}}\text{if}ca(\text{unbiasedness}).\hfill \end{array}$$
(10)
## III Quantum fluctuations of phase MUBs in odd prime characteristic
Following Pegg and Barnett, a good procedure to examine the phase properties of a quantized electromagnetic field state is by introducing a phase operator and this was one of the reasons that led them to introduce their famous Hermitian phase operator $`\mathrm{\Theta }_{PB}`$. In Section 6 of their seminal paper they showed “for future reference” how their phase operator could be employed to achieve this goal. In this section we proceed along the same lines using the phase form of the Wootters-Field MUBs.
### III.1 The Galois phase operator
On the other hand, the phase MUBs as given in (9) are eigenstates of a “Galois” quantum phase operator
$$\mathrm{\Theta }_{\mathrm{Gal}}=\underset{bF_q}{}\theta _b|\theta _b^a\theta _b^a|,a,bF_q.$$
(11)
with eigenvalues $`\theta _b=\frac{2\pi b}{q}`$. We use this fact to perform several calculations of quantum phase expectation values and phase variances for these MUBs.
Using (9) in (11) and the properties of the field theoretical trace the Galois quantum phase operator reads
$$\mathrm{\Theta }_{\mathrm{Gal}}=\frac{2\pi }{q^2}\underset{m,nF_q}{}\psi _k(nm)\omega _p^{tr[a(n^2m^2)]}S(n,m)|nm|,$$
(12)
where $`S(n,m)=_{bF_q}b\omega _p^{tr[b(nm)]}`$. In the diagonal matrix elements, we have the partial sums
$$S(n,n)=\frac{q(q1)}{2},$$
(13)
so that $`n|\mathrm{\Theta }_{\mathrm{Gal}}|n=\frac{\pi (q1)}{q}`$. In the non-diagonal matrix elements, the partial sums can be calculated from
$$\underset{bF_q}{}bx^b=x(1+2x+3x^2+\mathrm{}+qx^{q1})=x\left[\frac{1x^q}{(1x)^2}\frac{qx^q}{1x}\right]=\frac{xq}{x1},$$
(14)
where we introduced $`x=\omega _p^{tr(nm)}`$ and we made use of the relation $`x^q=1`$. Finally, we get
$$S(m,n)=\frac{q}{1\omega _p^{tr(mn)}}.$$
(15)
### III.2 The Galois phase-number commutator
Using (12) and the Galois number operator
$$N=\underset{lF_q}{}l|ll|,$$
(16)
the matrix elements of the phase-number commutator $`[\mathrm{\Theta }_{\mathrm{Gal}},N]`$ are calculated as
$$u_{\mathrm{Gal}}(n,m)=\frac{2\pi }{q^2}(nm)\psi _k(nm)\omega _p^{tr[a(n^2m^2)]}S(n,m).$$
(17)
The diagonal elements vanish, the corresponding matrix is anti-Hermitian since $`u_{\mathrm{Gal}}(n,m)=u_{\mathrm{Gal}}^{}(m,n)`$, and the states are pseudo-classical since $`lim_q\mathrm{}u_{\mathrm{Gal}}(n,m)=0`$. These properties are similar to those of the Pegg & Barnett commutator.
### III.3 Galois phase properties of a pure quantum electromagnetic state
For the evaluation of the phase properties of a general pure state of an electromagnetic field mode in the Galois number field we proceed similarly to Pegg & Barnett. Thus, we consider the pure state of the form
$$|f=\underset{nF_q}{}u_n|n,\mathrm{with}u_n=\frac{1}{\sqrt{q}}\mathrm{exp}(in\beta ),$$
(18)
where $`\beta `$ is a real parameter, and we sketch the computation of the phase probability distribution $`|<\theta _b|f>|^2`$, the phase expectation value $`<\mathrm{\Theta }_{\mathrm{Gal}}>=_{bF_q}\theta _b|<\theta _b|f>|^2`$ and the phase variance $`<\mathrm{\Delta }\mathrm{\Theta }_{\mathrm{Gal}}^2>=_{bF_q}(\theta _b<\mathrm{\Theta }_{\mathrm{Gal}}>)^2|<\theta _b|f>|^2`$, respectively (the upper index $`a`$ for the base is implicit and we discard it for simplicity).
The two factors in the expression for the probability distribution
$$\frac{1}{q^2}[\underset{nF_q}{}\psi _k(n)\mathrm{exp}(in\beta )\kappa (an^2bn)][\underset{mF_q}{}\psi _k(m)\mathrm{exp}(im\beta )\kappa (am^2+bm)],$$
(19)
have absolute values bounded by the absolute value of generalized Gauss sums $`G(\psi ,\kappa )=_{xF_q}\psi (g(x))\kappa (f(x))`$, with $`f,gF_q[x]`$. Weil W48 showed that for $`f(x)`$ of degree $`d`$ with $`gcd(d,q)=1`$ as in (4), under the constraint that for the multiplicative character $`\psi `$ of order $`s`$, the polynomial $`g(x)`$ should not be a $`s`$th power in $`F_q[x]`$ and with $`\nu `$ distinct roots in the algebraic closure of $`F_q`$, the order of magnitude of the sums is $`(d+\nu 1)\sqrt{q}`$. For a trivial multiplicative character $`\psi _0`$, and $`\beta =0`$, the overall bound is $`|<\theta _b|f>|^2\frac{1}{q}`$ and it follows that the absolute value of the Galois phase expectation value is bounded from above as expected for a common phase operator
$$|<\mathrm{\Theta }_{\mathrm{Gal}}>|\frac{2\pi }{q^2}\underset{bF_q}{}b\pi .$$
(20)
The exact formula for the phase expectation value reads
$$<\mathrm{\Theta }_{\mathrm{Gal}}>=\frac{2\pi }{q^3}\underset{m,nF_q}{}e^\beta (m,n)S(m,n),$$
(21)
where $`e^\beta (m,n)=\psi _k(mn)\mathrm{exp}[i(nm)\beta ]\chi [a(m^2n^2)]`$ and the sums $`S(m,n)`$ were defined in (13) and (15). The set of all the $`q`$ diagonal terms $`m=n`$ in $`<\mathrm{\Theta }_{\mathrm{Gal}}>`$ contributes an order of magnitude $`\frac{2\pi }{q^3}qS(n,n)\pi `$. The contribution from off-diagonal terms in (21) are not easy to evaluate analytically; we were able to show that for them one has $`|S(m,n)|=\frac{q}{2}|\mathrm{sin}[\frac{\pi }{p}tr(nm)]|^1`$.
The phase variance can be written as
$$<\mathrm{\Delta }\mathrm{\Theta }_{\mathrm{Gal}}^2>=\underset{bF_q}{}(\theta _b^22\theta _b<\mathrm{\Theta }_{\mathrm{Gal}}>)|<\theta _b|f>|^2.$$
(22)
The term $`<\mathrm{\Theta }_{\mathrm{Gal}}>^2_{bF_q}|<\theta _b|f>|^2`$ does not contribute since it is proportional to the Weil sum $`_{bF_q}\omega _p^{tr(b(nm)}=0`$. As a result a cancellation of the quantum phase fluctuations may occur in (22) from the two extra terms of opposite sign. But the calculation are again not easy to perform analytically. For the first term one gets $`2(2\pi /q^2)^2_{m,nF_q}e^\beta (m,n)|S(m,n)|^2.`$ The second term in (22) is $`2_{bF_q}\theta _b<\mathrm{\Theta }_{\mathrm{Gal}}>|<\theta _b|f>|^2=2<\mathrm{\Theta }_{\mathrm{Gal}}>^2`$. Partial cancellation occurs in the diagonal terms of (22) leading to the contribution $`\frac{2\pi ^2}{3}`$ which is still twice (in absolute value) the amount of phase fluctuations in the classical regime. A closed form for the estimate of the non-diagonal terms is still an open problem.
## IV Phase MUBs for $`\mathrm{m}`$-qubits
### IV.1 Mathematical preliminaries
The Weil sums (4) which have been proved useful in the construction of MUBs in odd characteristic $`p`$ (and odd dimension $`q=p^m`$), are not useful in characteristic $`p=2`$, since in this case the degree $`2`$ of the polynomial $`f_d(x)`$ is such that $`gcd(2,q)`$=2.
An elegant method for constructing complete sets of MUBs of $`m`$-qubits was found by Klappenecker and RöttelerKlapp03 . It makes use of objects belonging to the context of quaternary codes Wan97 , the so-called Galois rings $`R_{4^m}`$; we refer the interested reader to their paper for more mathematical details. We present a brief sketch in the following.
Any element $`yR_{4^m}`$ can be uniquely determined in the form $`y=a+2b`$, where $`a`$ and $`b`$ belong to the so-called Teichmüller set $`𝒯_m=(0,1,\xi ,\mathrm{},\xi ^{2^m2})`$, where $`\xi `$ is a nonzero element of the ring which is a root of the so-called basic primitive polynomial $`h(x)`$ Klapp03 . Moreover, one finds that $`a=y^{2^m}`$. We can also define the trace to the base ring $`𝒵_4`$ by the map
$$\stackrel{~}{tr}(y)=\underset{k=0}{\overset{m1}{}}\sigma ^k(y),$$
(23)
where the summation runs over $`R_{4^m}`$ and the Frobenius automorphism $`\sigma `$ reads
$$\sigma (a+2b)=a^2+2b^2.$$
(24)
In the Galois ring of characteristic $`4`$ the additive characters are
$$\stackrel{~}{\kappa }(x)=\omega _4^{\stackrel{~}{\mathrm{tr}}(x)}=i^{\stackrel{~}{\mathrm{tr}}(x)}.$$
(25)
The Weil sums (4) are replaced by the exponential sums Klapp03
$$\mathrm{\Gamma }(y)=\underset{u𝒯_m}{}\stackrel{~}{\kappa }(yu),yR_{4^m}$$
(26)
which satisfy
$$|\mathrm{\Gamma }(y)|=\{\begin{array}{cc}& 0\text{if}y2𝒯_m,y0\hfill \\ & 2^m\text{if}y=0\hfill \\ & \sqrt{2^m}\text{otherwise}.\hfill \end{array}$$
(27)
Gauss sums for Galois rings were constructed Oh01
$$G_y(\stackrel{~}{\psi },\stackrel{~}{\kappa })=\underset{xR_{4^m}}{}\stackrel{~}{\psi }(x)\stackrel{~}{\kappa }(yx),yR_{4^m},$$
(28)
where the multiplicative character $`\overline{\psi }(x)`$ can be made explicit Oh01 .
Using the notation $`\overline{\psi _0}`$ for a trivial multiplicative character and $`\stackrel{~}{\kappa _0}`$ for a trivial additive character, the Gaussian sums (28) satisfy $`G_y(\stackrel{~}{\psi _0},\stackrel{~}{\kappa _0})=4^m`$; $`G_y(\stackrel{~}{\psi },\stackrel{~}{\kappa _0})=0`$ and $`|G_y(\stackrel{~}{\psi },\stackrel{~}{\kappa })|2^m`$.
### IV.2 Phase states for $`m`$-qubits
The quantum phase states for $`m`$-qubits can be found as the “Galois ring” Fourier transform
$$|\theta ^{(y)}=\frac{1}{\sqrt{2^m}}\underset{n𝒯_m}{}\stackrel{~}{\psi }_k(n)\stackrel{~}{\kappa }(yn)|n,yR_{4^m}.$$
(29)
Using the Teichmüller decomposition in the character function $`\stackrel{~}{\kappa }`$ one obtains
$$|\theta _b^a=\frac{1}{\sqrt{2^m}}\underset{n𝒯_m}{}\stackrel{~}{\psi }_k(n)\stackrel{~}{\kappa }[(a+2b)n]|n,a,b𝒯_m.$$
(30)
This defines a set of $`2^m`$ bases (with index $`a`$) of $`2^m`$ vectors (with index $`b`$). Using the exponential sums (26), it is easy to show that the bases are orthogonal and mutually unbiased to each other and to the computational base. The case $`\overline{\psi }\overline{\psi _0}`$ was obtained before Klapp03 .
### IV.3 Phase MUBs for m-qubits: $`m=1`$, $`2`$ and $`3`$
For the special case of qubits, one uses $`\stackrel{~}{\mathrm{tr}}(x)=x`$ in (30) so that the three pairs of MUBs are given as
$$[|0,|1];\frac{1}{\sqrt{2}}[|0+|1,|0|1];\frac{1}{\sqrt{2}}[|0+i|1,|0i|1].$$
For $`2`$-qubits one gets a complete set of $`5`$ bases as follows
$`(|0,|1,|2,|3);`$
$`{\displaystyle \frac{1}{2}}[|0+|1+|2+|3,|0+|1|2|3,|0|1|2+|3,|0|1+|2|3]`$
$`{\displaystyle \frac{1}{2}}[|0|1i|2i|3,|0|1+i|2+i|3,|0+|1+i|2i|3,|0+|1i|2+i|3]`$
$`{\displaystyle \frac{1}{2}}[|0i|1i|2|3,|0i|1+i|2+|3,|0+i|1+i|2|3,|0+i|1i|2+|3]`$
$`{\displaystyle \frac{1}{2}}[|0i|1|2i|3,|0i|1+|2+i|3,|0+i|1+|2i|3,|0+i|1|2+i|3],`$
(31)
and for $`3`$-qubits a complete set of $`9`$ bases
$`(|0,|1,|2,|3,|4,|5,|6,|7);`$
$`{\displaystyle \frac{1}{4}}[|0+|1+|2+|3+|4+|5+|6+|7,|0+|1|2+|3|4|5|6+|7,`$
$`|0|1+|2|3|4|5+|6|7,|0+|1|2|3|4+|5+|6|7,`$
$`|0|1|2|3+|4+|5|6+|7,|0|1|2+|3+|4|5+|6|7,`$
$`|0|1+|2+|3|4+|5|6|7,|0+|1+|2|3+|4|5|6|7],`$
$`\mathrm{}`$ (32)
where only the first two bases have been written down for brevity reasons.
Quantum phase states of $`m`$-qubits (30) are eigenstates of a “Galois ring” quantum phase operator as in (11), and calculations of the same type as to those performed in Sect. (III) can be done, since the $`\stackrel{~}{\mathrm{tr}}`$ operator (23) fulfills rules similar to the $`\mathrm{tr}`$ operator (1). By analogy to the case of qudits in dimension $`p^m`$, $`p`$ an odd prime, phase properties for sets of $`m`$-qubits heavily rely on the Gauss sums (28). The calculations are tedious once again but can in principle be achieved in specific cases.
## V Mutual unbiasedness and maximal entanglement
Roughly speaking, entangled states in $`_q`$ cannot be factored into tensorial products of states in Hilbert spaces of lower dimension. We show now that there is an intrinsic relation between MUBs and maximal entanglement (see below).
We start with the familiar Bell states
$`(|_{0,0},|_{0,1})={\displaystyle \frac{1}{\sqrt{2}}}(|00+|11,|00|11),`$ $`(|_{1,0},|_{1,1})={\displaystyle \frac{1}{\sqrt{2}}}(|01+|10,|01|10),`$
where the compact notation $`|00=|0|0`$, $`|01=|0|1`$,…, is employed for the tensorial products.
These states are both orthonormal and maximally entangled, i.e., such that $`trace_2|_{h,k}_{h,k}|=\frac{1}{2}I_2`$, where $`trace_2`$ means the partial trace over the second qubit Nielsen00 .
One can define more general Bell states using the multiplicative Fourier transform (7) applied to the tensorial products of two qudits Cerf01 Durt04 ,
$$|_{h,k}=\frac{1}{\sqrt{q}}\underset{n=0}{\overset{q1}{}}\omega _q^{kn}|n,n+h,$$
(33)
These states are both orthonormal, $`_{h,k}|_{h^{},k^{}}=\delta _{hh^{}}\delta _{kk^{}}`$, and maximally entangled, $`trace_2|_{h,k}_{h,k}|=\frac{1}{q}I_q`$.
We define here an even more general class of maximally entangled states using the Fourier transform (9) over $`F_q`$ as follows
$$|_{h,b}^a=\frac{1}{\sqrt{q}}\underset{n=0}{\overset{q1}{}}\omega _p^{tr[(an+b)n]}|n,n+h.$$
(34)
The $`h`$ we use here has nothing to do with the polynomial $`h(x)`$ of Section (II). A list of the generalized Bell states of qutrits for the base $`a=0`$ can be found in Fujii01 which is a work that relies on a coherent state formulation of entanglement. In general, for $`q`$ a power of a prime, starting from (34) one obtains $`q^2`$ bases of $`q`$ maximally entangled states. Each set of the $`q`$ bases (with $`h`$ fixed) has the property of mutual unbiasedness.
Similarly, for sets of maximally entangled m-qubits one uses the Fourier transform over Galois rings (30) so that
$$|_{h,b}^a=\frac{1}{\sqrt{2^m}}\underset{n=0}{\overset{2^m1}{}}i^{tr[(a+2b)n]}|n,n+h.$$
(35)
For qubits ($`m=1`$) one gets the following bases of maximally entangled states (in matrix form, up to the proportionality factor)
$$\left[\begin{array}{cc}(|00+|11,|00|11)& (|01+|10,|01|10)\\ (|00+i|11,|00i|11)& (|01+i|10,|01)i|10)\end{array}\right].$$
(36)
Two bases in one column are mutually unbiased, while vectors in two bases on the same line are orthogonal to each other.
For two-particle sets of quartits, using Eqs. (31) and (35), one gets $`4`$ sets of $`(|_{h,b}^a`$, $`h=0,\mathrm{},3)`$, see them below, each entailing $`4`$ MUBs $`(a=0,\mathrm{},3)`$:
$`\{(|00+|11+|22+|33|++|++|++);`$
$`(|00|11i|22i|33|+(+i)(+i)|++(+i)(i)|++(i)(+i));\mathrm{}\}`$
$`\{(|01+|12+|23+|30|++|++|++);`$
$`(|01|12i|23i|30|+(+i)(+i)|++(+i)(i)|++(i)(+i));\mathrm{}\}`$
$$\{(|02+|13+|20+|31|++|++|++);\mathrm{}\}$$
$$\{(|03+|10+|21+|32|++|++|++);\mathrm{}\}$$
where, for the sake of brevity, we omitted the normalization factor ($`1/2`$) and the bases in the sets have been labeled by their coefficients unless for the first base. Thus, in the first set $`|00+|11+|22+|33|++++`$. Within each set, the four bases are mutually unbiased, as in (31), while the vectors of the bases from different sets are orthogonal.
As a conclusion, the two related concepts of mutual unbiasedness and maximal entanglement derive from the study of lifts of the base field $`𝒵_p`$ to Galois fields of prime characteristic $`p>2`$ (in odd dimension), or of lifts of the base ring $`𝒵_4`$ to Galois rings of characteristic $`4`$ (in even dimension). One wonders if lifts to more general algebraic structures would play a role in the study of non maximal entanglement. We have first in mind the nearfields that are used for deriving efficient classical codes and which have a strong underlying geometrygpiltz .
## VI Conclusion
In this research, we approached the MUBs fundamental topic from the point of view of the additive and multiplicative characters over finite fields in number theory. We consider that this framework is the most general including previous results in the literature as particular cases. Since MUBs are essentially generalized discrete Fourier transforms over finite number field kets, we formulated a quantum phase interpretation and illustrated several calculations of the phase properties of pure quantum states of the electromagnetic field in this finite number field mathematical context. Various types of Gauss sums get involved in this type of calculations of the MUBs phase properties of a pure quantum state and the generalization to the mixed states, although straightforward through the usage of the density matrix formalism, could lead to even more complicated calculations involving such sums. We hope to evaluate them in future works. We also mentioned in the last section a possible application to phase MUBs states of Bell type. This could lead to finite number field measures of the degree of entanglement.
Note: The authors acknowledge Dr. Igor Shparlinski for suggesting important corrections on this paper. On June 20/2005 he sent us an e-message pointing to some mathematical inconsistencies in subsection III.C (i.e., subsection 3.3 in the epjd version). Parameter $`\beta `$ therein should be itself an element in $`F_q`$ to perform the calculations.
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# Modular deformations of analytic polyhedra11footnote 1This research was supported by a grant from the GIF, German-Israeli Foundation of Scientific Research and Development.
## 1 Introduction
The notion of modular deformation helps to define a canonical analytic structure in the category of analytic objects. In this paper, we discuss modular deformations of complex analytic spaces. The category of germs of complex spaces is not appropriate for this purpose since the fiber of a deformation of a germ is no more a germ. Instead, we consider the category of analytic polyhedrons (see Sec.4). Let $`f:𝒳S`$ be a deformation of an analytic space or of a polyhedron. Take a point $``$ in the base $`S`$ and consider the pair $`(S,)`$ as a germ of complex space. A modular stratum is an analytic subspace $`(M,)(S,)`$ such that the uniqueness property holds: any germ morphisms $`g:R(S,)`$ and $`h:R(M,)`$ induce isomorphic deformations $`f\times _Sgf\times _Sh`$ only if $`g=h.`$ The modular deformation is the restriction of $`f`$ to $`M,`$ that is the deformation $`f_Mf\times _SM.`$ Suppose that $`f:𝒳(S,)`$ is a versal deformation of the fiber $`X_{}=f^1()`$ and $`M`$ is the maximal modular stratum. It has the special feature: for any other versal deformation $`f^{}:𝒳^{}S^{}`$ of $`X`$ and an isomorphism $`i:SS^{}`$ such that $`i\times f^{}f`$, the restriction of $`i`$ to $`M`$ is uniquely defined and $`M^{}i\left(M\right)`$ is the maximal modular stratum in $`S^{}`$. The isomorphism $`\left(i\times f^{}\right)_Mf_M`$ is also canonically defined. Moreover, the maximal modular stratum $`M`$ (if it exists) possesses the semi-local property: for any point $`sM`$ close to the marked point $``$ the germ $`(M,s)`$ is also maximal modular for the fiber $`X_sf^1\left(s\right)`$. Therefore two versal deformations $`f`$ and $`f^{}`$ can be glued together (amalgamated) along its modular strata, say $`M`$ and $`M^{},`$ provided there is given an isomorphism $`\varphi :f^1\left(s\right)\stackrel{~}{}f^1\left(s^{}\right)`$ for $`sM,s^{}M^{}.`$ This gives rise to a flat morphism $`F:XMM^{}`$ where $`MM^{}`$ denotes the amalgam over some open neighborhood of the points $`s`$ and $`s^{}.`$ Repeating this construction, we can get a flat morphism of complex spaces $`F:𝒳`$ which is a maximal modular deformation of each its fiber. This conception is close to that of fine moduli space in the sense of Mumford, . For each point $`s,`$ the automorphism group of the space $`X_s`$ acts on $``$ which may possible further amalgamation and the global amalgam may have complicated structure.
We formulate here some general results on existence and properties of maximal modular deformations in the category of complex analytic polyhedra and discuss several examples of modular deformation of isolated singularities. V.Arnold’s classification table of hypersurface singularities is a rich source of such examples. For many examples of complete intersection singularities modular deformations are contained in the paper of A. Aleksandrov , further study see in .
Several complicated examples, which have been beyond the reach of ‘by hand’ calculations, were computed and studied by B.Martin and by T. Hirsch and B. Martin . These authors applied a specialized computer algebra program based on SINGULAR. We shall see how these new local modular families glue together and global modular families emerge. In several cases the base of a modular family can be made compact by gluing together sufficiently many local modular deformations.
I thank Bernd Martin for many helpful discussions.
## 2 Rudiments of the deformation theory
Remind some basic definitions. A germ of complex analytic spaces (which will be called simply ‘germ’) is a pair $`(Z,𝒪(Z))`$, where $`Z`$ is the germ of complex analytic set in $`(^n,)`$ for some $`n`$ and $`𝒪(Z)`$ is the coherent sheaf of analytic $``$-algebras on $`(^n,)`$ such that $`\mathrm{supp}𝒪\left(Z\right)=Z`$. We use the notation $``$ for the marked point in $`^n`$ and in $`Z`$. Any morphism of germs $`(W,𝒪(W))(Z,𝒪(Z))`$ is a pair $`(f,\varphi )`$ where $`f:(W,)(Z,)`$ is a mapping of germs of analytic sets and $`\varphi :f^{}\left(𝒪(Z)\right)𝒪(W)`$ is a morphism of sheaves of analytic $``$-algebras over $`W`$. Restricting $`\varphi `$ to the marked point yields the morphism of analytic $``$-algebras $`\varphi _{}:𝒪_{}(Z)𝒪_{}(W)`$. The morphism $`(f,\varphi )`$ is called embedding, if $`\varphi _{}`$ is a surjection. Vice versa, for any analytic $``$-algebra $`A`$ there exists and is uniquely defined the germ $`(Z,𝒪\left(Z\right))`$ such that $`A𝒪_{}\left(Z\right)`$. In particular, the algebra corresponding to the simple point $``$ is the field $``$. For an arbitrary germ $`Z`$ there is the canonical embedding $`Z`$; the corresponding algebra homomorphism is the canonical morphism $`𝒪_{}(Z)`$. The kernel of this morphism is the maximal ideal in the algebra $`A=𝒪_{}(Z);`$ it is denoted $`𝔪(Z)`$. The germ $`Z=`$ is called simple point; a germ $`Z`$ is called fat point, if the corresponding algebra has a finite dimension over $``$.
The fiber product operation is well defined in the category of analytic spaces; it is denoted $`f\times _Xg`$ or $`Y\times _XZ`$ for germ morphism $`f:YX,g:ZX.`$ If $`g`$ is an embedding, we call the product $`f\times _Xg`$ restriction $`f`$ to $`Z`$ and denote it $`f|Z`$. Let $`V`$ be a $``$-vector space of finite dimension; we denote by $`𝔻\left[V\right]`$ the $``$-algebra that is isomorphic to $`V`$ as vector space with the multiplication rule $`(a,u)(b,v)=(ab,av+bu)`$. This is an analytic algebra of a fat point denoted $`P\left(V\right)`$. All germs of this kind form a category denoted D, where for any objects $`P\left(U\right),P\left(V\right)`$ the set $`\mathrm{Hom}(P\left(U\right),P\left(V\right))`$ is bijective to $`\mathrm{Hom}_{}(U,V).`$ This induces a natural structure of vector space over $``$ in the set $`\mathrm{Hom}(P\left(U\right),P\left(V\right))`$; the composition of morphisms is a bilinear operation. For any germ $`Z`$ there is defined the germ $`T(Z)\text{D}`$ and the embedding $`t(Z):T(Z)Z`$ such that $`𝒪(T(Z))=𝒪(Z)/𝔪^2(Z).`$ The germ $`T(Z)`$ is called the tangent space of the germ $`Z`$.
Definition 1. Let $`X`$ be a complex analytic space. A deformation of $`X`$ over a germ $`(Z,)`$ is a pair $`(f,i)`$, where $`f:𝒳(Z,)`$ is a flat morphism and $`i:Xf^1()=f\times `$ is an isomorphism of complex analytic spaces. For any morphism of germs $`h:WZ`$ and a deformation $`(f,i)`$ of $`X`$ over $`Z`$ the fiber product $`f\times _Zh:𝒳\times _ZWW`$ is also flat. The pair $`(f\times _Zh,ji)`$ is a deformation of $`X`$ with base $`W,`$ where $`j:f^1()\left(f\times _Zh\right)^1()`$ is the canonical bijection. In the same way, one can treat deformations of $`_2`$-graded, $``$-graded analytic spaces, deformations of germs, of fiber bundles, of coherent sheaves and the like.
A versal deformation of the space $`X`$ is a deformation $`(f,i),`$ $`f:𝒳S`$ such that for any deformation $`(g,j),g:𝒴R`$ there exists a morphism of germs $`h:RS`$ and a isomorphism $`\gamma :𝒳\times _SR𝒴`$ over $`R`$ such that the diagram commutes
| $`𝒳\times _S`$ | | $`\stackrel{\gamma \times _R}{}`$ | | $`𝒴\times _R`$ |
| --- | --- | --- | --- | --- |
| | $`\stackrel{𝑖}{}`$ | | $`\stackrel{𝑗}{}`$ | |
| | | $`X`$ | | |
A pair $`(f,i)`$ is called universal deformation, if the morphism $`h`$ is unique. The above definition can be applied for the category D. Suppose that a space $`X`$ possesses a universal deformation $`\delta :𝒳T_X`$ in the category D. Take an arbitrary germ $`S`$ and a deformation $`(f,i)`$ of $`X`$ over $`S`$; consider the morphism $`f_T=`$ $`f\times _ST_{}\left(S\right)`$. It is a deformation of $`X`$ with the base $`T_{}\left(S\right)\text{D}`$. Due to the universality property of $`\delta `$ a morphism $`\mathrm{D}_{}f:T_{}(S)T_X`$is defined in the category D such that $`\mathrm{D}_{}f\times \delta f_T.`$ It is called Kodaira-Spencer mapping; this mapping is linear since it belongs to D. Let $`(f,i)`$ be a versal deformation of $`X`$ with a base $`S.`$ There exists a morphism $`t:T_XS`$ such that $`f\times _St\delta .`$ Then $`\mathrm{D}_{}ft=\mathrm{id}`$ since $`\delta `$ is universal, consequently $`\mathrm{D}_{}f`$ is surjective. The pair $`(f,i)`$ is called miniversal (or minimal versal), if $`\mathrm{D}_{}f`$ is a injection, hence a bijection.
Definition 2. , Let $`f`$ be a deformation of $`X`$ as above. A subgerm (stratum) $`MS`$ is called modular, if for any morphisms of germs $`h:RM,g:RS`$ the equation $`f\times _Sh=f\times _Sg`$ implies that $`h=g`$. A modular stratum $`MS`$ is called maximal, if it contains any other modular stratum. If $`f:𝒳S`$ is a miniversal deformation and $`M`$ is the maximal modular stratum, the restriction $`f_M`$ is called maximal modular deformation of $`X.`$ Any universal deformation is, of course, maximal modular.
A similar conception (prorepresenting stratum) was introduced in for the formal deformation theory of affine schemes. See for further results. The notion of modular deformation was also treated in in a more general setting.
###### Proposition 2.1
Let $`f:𝒳S`$ and $`f^{}:𝒳^{}S^{}`$ be miniversal deformations of a space $`X`$ and $`M,M^{}`$ are respective maximal modular strata. There exists a isomorphism $`h:M^{}M`$ such that $`f^{}|M^{}f|M\times h.`$ The morphism $`h`$ is unique.
$``$ There exist morphisms $`g:S^{}S`$ and $`g^{}:SS^{}`$such that $`f^{}\times g^{}f,f\times gf^{}.`$ We have $`\mathrm{D}_{}f\mathrm{d}g=\mathrm{D}_{}f^{}`$ and $`\mathrm{D}_{}f^{}\mathrm{d}g^{}=\mathrm{D}_{}f.`$ The differentials $`\mathrm{d}g,`$ $`\mathrm{d}g^{}`$ are inverse one to another, since $`\mathrm{D}_{}f,\mathrm{D}_{}f^{}`$ are bijections. Therefore $`g`$ is a isomorphism of germs and the mapping $`g:M^{}S`$ is a modular germ. It is factorized through a uniquely defined morphism $`h:M^{}M,`$ since $`M`$ is maximal modular. $``$
## 3 Tangent cohomology and criterion of a modular stratum
Let $`g:X`$ be a complex analytic space; the graded tangent sheaf $`_{q=0}𝒯^q(X)`$ is defined on $`X.`$ It is a $`_+`$-graded sheaf algebra Lie, which means that there is defined a bracket operation that satisfies the graded commutation and Jacobi identities. Moreover, this has a natural structure of $`𝒪\left(X\right)`$-module which agrees with the Lie algebra structure in a natural way. For any $`q0`$ the sheaf $`𝒯^q`$ is a coherent $`𝒪\left(X\right)`$-sheaf. In particular, the term $`𝒯^0\left(X\right)`$ is the sheaf of tangent fields on $`X`$.
The global tangent cohomology $`T^{}(X)=_0^{\mathrm{}}T^n(X)`$ is a $`_+`$-graded algebra and there is a spectral sequence $`E^{}`$ that converges to the tangent cohomology $`T^{}(X)`$ with the second term $`E_2^{pq}=H^p(X,𝒯^q(X))`$, see . The term $`T^0\left(X\right)`$ is the Lie algebra of tangent fields on $`X.`$ If the term $`T^1\left(X\right)`$ is of finite dimension, it represents the base $`T_X`$ of the universal deformation $`\delta `$ of $`X`$ in the category D as in Sec.2. If $`g:XY`$ is a morphism of complex spaces, a vertical tangent field $`t`$ on $`g`$ is a tangent field on $`X`$ with the property $`t\left(g^{}\left(a\right)\right)=0`$ for an arbitrary $`a𝒪\left(Y\right).`$ The notation $`T^0\left(X/Y\right)`$ means the space of vertical tangent fields. For any point $`yY`$ the restriction mapping $`T^0\left(X/Y\right)T^0\left(X_y\right)`$ is canonically defined where $`X_y=g^1\left(y\right)`$ is the complex subspace of $`X.`$ There is a general criterion of modularity:
###### Theorem 3.1
Let $`f:𝒳(S,)`$ be a deformation of a complex analytic space $`X`$. Then
(i) the simple point $`S`$ is modular, if and only if the Kodaira-Spencer mapping $`\mathrm{D}_{}f`$ is injective;
(ii) if $`\mathrm{D}_{}f`$ is injective, then a subgerm $`MS`$ is modular, if for any fat point $`ZM`$ the restriction mapping $`T^0(𝒳\times _MZ/Z)T^0(X)`$ is surjective; if $`\mathrm{D}_{}f`$ is bijective, this condition is necessary as well;
(iii) $`T_{}(M)`$ coincides with the space of tangent vectors $`tT_{}\left(S\right)`$ that satisfy the equation $`[\mathrm{D}_{}f(t),v]=0`$ for any $`vT^0(X)`$.
A proof is similar to that of , Proposition 5.1.
## 4 Analytic polyhedrons
The existence of maximal modular deformation and some its properties were stated in for an arbitrary compact complex analytic space. Here we formulate similar results for analytic polyhedra. First we recall some definitions of .
Definition 3. For an arbitrary integer $`n`$ we fix a coordinate space $`^n`$, i.e. a complex vector space with a marked system of linear coordinates $`w_1,\mathrm{},w_n`$. Denote by $`D^n`$ the closed unit polydisk in $`^n`$. A complex analytic $`n`$-polyhedron is a pair $`(X,\phi )`$, where $`X`$ is a complex analytic subspace of a complex space $`\stackrel{~}{X}`$ and $`\phi :\stackrel{~}{X}^n`$ is a holomorphic mapping (called a barrier map ) such that the set $`\phi ^1\left(\overline{D}^n\right)`$ is compact and $`X=\phi ^1\left(D^n\right).`$ The neighborhood $`\stackrel{~}{X}`$ of $`X`$ can be contracted, i.e. $`X`$ is thought as the germ of a complex analytic space $`\stackrel{~}{X}`$ on the compact set $`\overline{X}=\phi ^1\left(\overline{D}^n\right).`$ The set $`X\phi ^1\left(D^n\right)`$ is the boundary of the polyhedron $`X.`$ A morphism of polyhedra $`(X,\phi )`$ $`(Y,\psi )`$ is a pair $`(f,p),`$ where $`f:\overline{X}\overline{Y}`$ of germs of complex spaces and $`p:^n^m`$ is a coordinate projection that make the commutative diagram: $`p\phi =\psi f`$.
Let $`S`$ be a complex analytic spaces; we call a relative analytic $`n`$-polyhedron (r.p.) over $`S`$ any pair $`(X,\phi )`$, where $`X`$ is an open subspace a complex space $`\stackrel{~}{X}`$ and $`\phi :\stackrel{~}{X}^n\times S`$ is a holomorphic mapping such that the set $`\phi ^1\left(\overline{D}^n\times S\right)`$ is proper over $`S`$ and $`\phi ^1\left(D^n\times S\right)=X`$. Let $`(X,\phi )`$ be a r.p. over $`S`$ and $`h:RS`$ be a morphism of complex spaces. The fiber product $`(X_R,\phi _R)(X,\phi )\times _SR`$ is a r.p. over $`R.`$ In particular, for any point $`sS`$ the fiber of $`(X,\phi )`$ is defined as the product $`(X_s,\phi _s)=(X,\phi )\times _Ss`$, which is an (absolute) analytic polyhedron, called the fiber of $`(X,\phi ).`$
A morphism $`(X,\phi )(Y,\psi )`$ of relative polyhedra over $`S`$ is any pair $`(f,p)`$ of holomorphic mappings of germs $`f:\overline{X}\overline{Y},p:\overline{D}^n\times S\overline{D}^m\times S`$ such that $`p\phi =\psi f`$ and $`p`$ commutes with the projections to $`S.`$ If $`S=`$ this is a definition of morphism of (absolute) polyhedra.
Definition 4. Let $`S`$ be a germ of complex space with the marked point $``$, and $`(X,\phi )`$ be a polyhedron; we call a deformation of this polyhedron with the base $`S`$ any triple $`(𝒳,\mathrm{\Phi };\theta )`$, where
(i) $`(𝒳,\mathrm{\Phi })`$ is a relative polyhedron over $`S`$,
(ii) $`(𝒳,\mathrm{\Phi })`$ is flat over $`S,`$ i.e. the composition $`\pi \mathrm{\Phi }:𝒳^n\times SS`$ is a flat morphism of complex analytic spaces, where $`\pi `$ is the projection.
(iii) $`\theta :(X,\phi )(𝒳,\mathrm{\Phi })\times _S\{\}`$ is an isomorphism of polyhedra. A deformation $`(𝒳,\mathrm{\Phi };\theta )`$ of a polyhedron $`(X,\phi )`$ is called versal, if for any deformation $`(𝒴,\mathrm{\Psi };\eta )`$ of the same polyhedron with a base $`R`$ there exist a morphism of germs $`h:RS`$ and a isomorphism $`\alpha :(𝒴,\mathrm{\Psi })(𝒳,\mathrm{\Phi })\times _SR`$ such that $`(\alpha \times )\theta =\eta `$. Definitions of modular stratum, maximal modular deformation of an analytic polyhedron are given in the same lines as in Definition 2.
## 5 Versal and modular deformations
Let $`^n`$ be again a coordinate space and $`D^n`$ be the closed unit polydisk. For an arbitrary coordinate projection $`p:^n^m`$ we have $`p(D^n)=D^m`$ and the sheaf $`𝒪(D^n)`$ is endowed with a structure of module over the sheaf algebra $`𝒪(D^m)`$.
###### Theorem 5.1
Let $`(X,\phi )`$ be an analytic $`n`$-polyhedron such that
(i) for any coordinate projection $`p:^n^m`$ the equations hold:
$$Tor_k^{𝒪(D^m}(𝒯_\phi ^{k+1},)=0,k=0,1,\mathrm{},m$$
(1)
in each point of the set $`(D)^m\times D^{nm}`$, where $`𝒯^{}𝒯^{}(X)`$ and
(ii) $`\varphi `$ is finite over $`D^n.`$ The polyhedron $`(X,\phi )`$ has a miniversal deformation $`F:𝒳(S,0)`$ where $`Sg^1(0)`$ and $`g:(T^1\left(X\right),0)T^2(X)`$ is a holomorphic mapping of germs such that $`g(0)=0,\mathrm{d}g(0)=0`$.
Remark. The conditions (1) are fulfilled for the trivial barrier $`\phi =0`$, since $`𝒪(D^0)=.`$ This implies the existence of miniversal deformation for any compact complex space. This case was studied by H.Grauert and other authors, see the survey . Theorem 5.1 was proved in . Note that the space $`X`$ that satisfies (1) need not to be smooth at the boundary of the polyhedron. The condition (1) for $`k=0`$ means that the sheaf $`𝒯_\phi ^1`$ vanishes on the boundary of $`D^n`$, hence the support of this sheaf is a compact subset of the interior of $`D^n`$. The spectral sequence $`E_2^{pq}=H^p(X,𝒯^q)T^{}\left(X\right)`$ yields
$$0H^1(X,𝒯^0)T^1(X)\mathrm{\Gamma }(X,𝒯^1)H^2(X,𝒯^0)\mathrm{}$$
The term $`\mathrm{\Gamma }(X,𝒯^1)`$ is of finite dimension since $`𝒯^1`$ is supported by a compact set. From the Leray sequence we find the exact sequence
$$H^1(X,𝒯^0)H^0(D^n,^1\phi \left(𝒯^0\right))\stackrel{d_2}{}H^2(D^n,^0\phi \left(𝒯^0\right))H^1(D^n,^0\phi \left(𝒯^0\right))$$
(2)
The sheaves $`^q\phi \left(𝒯^0\right),q=0,1`$ are coherent by Grauert’s Theorem, hence the terms $`H^1`$ and $`H^2`$ vanish in the right-hand side. Therefore we have
$$H^1(X,𝒯^0)H^0(D^n,^0\phi \left(𝒯^0\right)).$$
The sheaf $`^1\phi \left(𝒯^0\right)`$ vanishes at the boundary of the polydisk since of (i). This yields that $`H^1(X,𝒯^0)`$ has finite dimension. Finally $`\tau dim_{}T^1(X)<\mathrm{}`$.
The second condition of Theorem 5.1 can be weaken as follows
(ii)’$`^1\phi \left(𝒯^0\right)|D^n=0.`$
###### Theorem 5.2
Let $`(X,\phi )`$ be a polyhedron as in the previous theorem and $`f:𝒳(S,)`$ be its miniversal deformation. Then there exists a neighborhood $`S^{}`$ of the marked point $``$ in $`S`$ and a closed subspace $`MS^{}`$ such that:
(i) for any $`sM`$ the germ $`(M,s)`$ is maximal modular for the deformation $`f`$ of the fiber $`X_s`$.
(ii) the restriction mapping $`T^0(𝒳\times _SM/M)T^0(X)`$ is surjective and $`M`$ is maximal with this property;
(iii) the Zariski tangent space $`T_{}(M)`$ is the space of all vectors $`tT_{}\left(S\right)`$ satisfying the equation $`[\mathrm{D}_{}f\left(t\right),v]=0`$ for any $`vT^0(X)`$;
(iv) the support of $`M`$ is the set of points $`sS^{}`$, where the mapping $`\mathrm{D}_sf:T_s\left(S\right)T^1\left(X_s\right)`$ is injective.
The statement (i) is essential; it means that any point of $`M^{}`$ can be taken as marked one, hence it is a globally defined complex analytic space. A proof can be done on the lines of and . The parts (ii) and (iv) then follow from Theorem 3.1.(i). The part (iii) is a corollary of Theorem 3.1.(iii), since the bracket $`[\mathrm{D}f\left(t\right),v]`$ is for any $`tT_{}\left(S\right)`$ the first obstruction for extension of $`v`$ to a tangent field on $`𝒳\times _SM/M.`$ The statement (iv) implies that
(v) if $`S`$ is regular, then $`\mathrm{supp}M`$ is given by the equation $`dimT^1\left(X_s\right)=dimS,sS^{}`$.
Definition 6. Take a versal deformation $`(𝒳,S)`$ as in Theorem 5.1 and the maximal modular stratum $`MS`$ as in Theorem 5.2. We call $`f|M:𝒳\times _SMM`$ maximal modular deformation. Let $``$ be a complex analytic space and $`Y`$ be an analytic polyhedron over $``$; we call a modular family the mapping $`𝒴`$, if it is maximal modular deformation of each fiber $`Y_s,s`$.
## 6 Shrinking a polyhedron
Let $`(g,q):(Y,\psi )(X,\phi )`$ be a morphism of analytic polyhedra. We call it imbedding, if $`g`$ is an imbedding, $`q`$is a coordinate projection and there exists a proper holomorphic mapping $`\rho :X^n`$ that makes the commutative diagram
| $`Y`$ | $`\stackrel{𝑔}{}`$ | $`X`$ |
| --- | --- | --- |
| $`\psi `$ | $`\rho `$ | $`\phi `$ |
| $`^n`$ | $`\stackrel{𝑞}{}`$ | $`^m`$ |
In geometric terms, this means that $`\phi _i=q_i\rho `$ for $`i=1,\mathrm{},m`$ and
$$Y=\left\{xX,\left|\psi _{m+1}\left(x\right)\right|<1,\mathrm{},\left|\psi _n\left(x\right)\right|<1\right\}$$
for an appropriate numeration of coordinates in $`^n.`$
###### Proposition 6.1
Let $`(g,q,\rho )`$ be an imbedding of polyhedra such that
(i) $`\mathrm{supp}𝒯^1\left(X\right)g\left(Y\right)`$ and
(ii) the mapping
$$\rho :X\backslash g\left(Y\right)D^m\times ^{nm}D^{nm}$$
is finite. Then we have $`T^1\left(Y\right)T^1\left(X\right).`$
Proof. Compare cohomology of tangent sheaves in $`X`$ and $`Y.`$ Consider the diagram
| $`H^1(Y,𝒯^0\left(Y\right))`$ | $``$ | $`H^0(D^n,\mathrm{R}^1\psi \left(𝒯^0\left(Y\right)\right))`$ | $`\stackrel{𝑖}{}`$ | $`H^0(D^m,\mathrm{R}^0q\mathrm{R}^1\psi \left(𝒯^0\left(Y\right)\right))`$ |
| --- | --- | --- | --- | --- |
| $`h_1`$ | | | | $`\gamma `$ |
| $`H^1(X,𝒯^0\left(X\right))`$ | $``$ | $`H^0(D^m,\mathrm{R}^1\phi \left(𝒯^0\left(X\right)\right))`$ | $`\stackrel{𝑗}{}`$ | $`H^0(D^m,\mathrm{R}^0q\mathrm{R}^1\rho \left(𝒯^0\left(X\right)\right))`$ |
The left-side isomorphisms follow from (2). The Leray spectral sequence
$$H^r(D^m,\mathrm{R}^sq_{}\left(\right))H^{r+s}(D^n,),\mathrm{R}^1\psi \left(𝒯^0\left(Y\right)\right)$$
yields the isomorphism $`i,`$ since the sheaf $``$ is coherent. We have $`\mathrm{R}^rq\mathrm{R}^0\rho \left(𝒯^0\left(X\right)\right)=0`$ for $`r=1,2`$ since the sheaf $`\mathrm{R}^0\rho \left(𝒯^0\left(X\right)\right)`$ is coherent in $`D^m\times ^{nm}`$ and $`q`$ is a coordinate projection. By applying the Leray sequence to the composition $`\phi =q\rho `$ we obtain the isomorphism in $`D^n`$
$$\mathrm{R}^1\phi \left(𝒯^0\left(X\right)\right)=\mathrm{R}^0q\mathrm{R}^1\rho \left(𝒯^0\left(X\right)\right)$$
This yields the bijection $`j.`$ Further we have
$$\mathrm{R}^1\psi \left(𝒯^0\left(Y\right)\right)=\mathrm{R}^1\rho g\left(𝒯^0\left(Y\right)\right)=\mathrm{R}^1\rho \mathrm{R}^0g\left(𝒯^0\left(Y\right)\right)=\mathrm{R}^1\rho \left(𝒯^0\left(X\right)|g\left(Y\right)\right)$$
since $`\mathrm{R}^1g=0.`$ On the other hand, the sheaf $`\mathrm{R}^1\rho \left(𝒯^0\left(X\right)\right)`$ vanishes in $`X\backslash g\left(Y\right)`$ since $`\rho `$ is finite there. Therefore the right-hand side is isomorphic to the sheaf $`\mathrm{R}^1\rho \left(𝒯^0\left(X\right)\right).`$ This yields the bijection $`\gamma `$ and also the bijection $`h_1.`$ Consider the commutative diagram
| $`\mathrm{Ker}d_2\left(Y\right)`$ | $``$ | $`H^0(Y,𝒯^1\left(Y\right))`$ | $`\stackrel{d_2\left(Y\right)}{}`$ | $`H^2(Y,𝒯^0\left(Y\right))`$ |
| --- | --- | --- | --- | --- |
| $`k_2`$ | | $`h_0`$ | | $`h_2`$ |
| $`\mathrm{Ker}d_2\left(X\right)`$ | $``$ | $`H^0(X,𝒯^1\left(X\right))`$ | $`\stackrel{d_2\left(X\right)}{}`$ | $`H^2(X,𝒯^0\left(X\right))`$ |
The isomorphism $`h_2`$ is defined similarly to $`h_1`$ and the bijective mapping $`h_0`$ follows from (i). This diagram yields the isomorphism $`k_2.`$ This together with $`h_1`$ gives
$$T^1\left(Y\right)H^1(Y,𝒯^0\left(Y\right))\mathrm{Ker}d_2\left(Y\right)H^1(X,𝒯^0\left(X\right))\mathrm{Ker}d_2\left(X\right)T^1\left(X\right).$$
###### Theorem 6.2
Let $`(Y,\psi ),(X,\phi )`$ be analytic polyhedra that fulfil the conditions of Theorem 5.1 and $`(𝒳,\mathrm{\Phi },S_X)`$ and $`(𝒴,\mathrm{\Psi },S_Y)`$ be their versal deformations. Let $`(g,q,\rho ):(Y,\psi )(X,\phi )`$ be a imbedding of analytic polyhedra that fulfil the conditions (i), (ii) of Proposition 6.1. Then there exists an imbedding of relative polyhedra $`(G,Q,R):(𝒴,\mathrm{\Psi },S_Y^{})(𝒳,\mathrm{\Phi },S_X^{})`$ such that $`G\times =g,Q\times =q,`$ $`R\times =\rho `$ and the diagram commutes
| $`Y`$ | $``$ | $`𝒴`$ | $`\stackrel{𝐺}{}`$ | $`𝒳`$ | $``$ | $`X`$ |
| --- | --- | --- | --- | --- | --- | --- |
| $`\psi `$ | | $`\mathrm{\Psi }`$ | $`R`$ | $`\mathrm{\Phi }`$ | | $`\phi `$ |
| $`^n\times `$ | $``$ | $`^n\times S_Y^{}`$ | $`\stackrel{𝑄}{}`$ | $`^m\times S_X^{}`$ | $``$ | $`^m\times `$ |
| | $``$ | $``$ | | $``$ | $``$ | |
| | | $`^n`$ | $`\stackrel{𝑞}{}`$ | $`^m`$ | | |
(3)
where $`S_X^{}`$ and $`S_Y^{}`$ are some neighborhoods of marked points in $`S_X,`$ respectively in $`S_Y.`$
Sketch of the proof. We can write $`\rho =\phi \times \xi `$ where $`\xi :X^{nm}`$ is a bounded holomorphic mapping, say $`\left|\xi \right|<b`$ in $`\overline{X}`$ for some positive $`b.`$ Consider the polyhedron $`(X,\stackrel{~}{\rho })`$ where $`\stackrel{~}{\rho }=\phi \times b^1\xi .`$ It has the same boundary as $`X`$ and fulfils the conditions of Theorem 5.1. Take a versal deformation $`(𝒳,R,S)`$ of this polyhedron, where $`R:𝒳^n\times S`$ is the barrier function. We can write $`R=R_m\times R_{nm}\times F,`$ where $`R_m:𝒳^m,R_{nm}:𝒳^{nm},F:𝒳S,`$ and define the mapping $`\mathrm{\Phi }=R_m\times F:𝒳^m\times S.`$ We have $`R_m\times =\phi `$ and the triple $`(𝒳,R_m,S)`$ is a deformation of the polyhedron $`(X,\phi ).`$ Set $`\mathrm{\Psi }=R_m\times bR_{nm}\times F:𝒳^n\times S;`$ the polyhedron $`𝒴\left\{xX;\left|bR_{nm}\right|<1\right\}`$ is embedded in $`𝒳`$ and $`𝒴\times =Y,\mathrm{\Psi }\times =\psi ,`$ that $`(𝒴,\mathrm{\Psi },S)`$ is a deformation of $`(Y,\psi ).`$ The diagram (3) commutes, where $`S_X^{}=S_Y^{}=S`$ and $`Q=q\times \mathrm{id}_S.`$ Analyzing the construction of $`(𝒳,R,S)`$ (see ) we can see that $`(𝒳,R_m,S)`$ is isomorphic to a versal deformation $`(𝒳,\mathrm{\Phi },S_X)`$ of the polyhedron $`(X,\phi )`$. Next we prove that $`(𝒴,\mathrm{\Psi },S)`$ is a versal deformation of $`(Y,\psi )`$ and Theorem follows. $``$
Definition. The class of imbeddings in $`(X,\phi )`$ that satisfy the conditions of Proposition 6.1 is called the germ of $`(X,\phi );`$ each polyhedron $`(Y,\psi )`$ is a representative of this germ. By the above theorem the versal deformations of all representative are naturally isomorphic. So are the maximal modular strata $`M`$ and maximal modular deformations.
## 7 Amalgams and automorphisms
Amalgam of modular deformations. Let $`f:𝒳M`$ and $`g:𝒴N`$ be modular families. If there are points $`sM,tN`$ such that the fibers $`X_s`$ and $`Y_t`$ are isomorphic, then there exist neighborhoods $`U`$ of $`s`$ and $`V`$ of $`t`$ and a uniquely defined commutative diagram
| $`f^1\left(U\right)`$ | $`\stackrel{𝐴}{}`$ | $`g^1\left(V\right)`$ |
| --- | --- | --- |
| $``$ | | $``$ |
| $`U`$ | $`\stackrel{𝑎}{}`$ | $`V`$ |
This statement follows from Theorem 6.2. The above families can be patched together along $`a`$ giving rise to the modular family $`𝒳_a𝒴`$ with the base $`M_AN,`$ which is the amalgam (coproduct) of $`M`$ and $`N`$:
$$\begin{array}{ccccc}f^1\left(U\right)\hfill & & \stackrel{𝐴}{}\hfill & & g^1\left(V\right)\hfill \\ \hfill & & & & \hfill \\ 𝒳\hfill & \hfill & 𝒳_a𝒴\hfill & \hfill & 𝒴\hfill \\ \hfill & & \hfill & & \hfill \\ M\hfill & \stackrel{𝜇}{}\hfill & M_AN\hfill & \stackrel{𝜈}{}\hfill & N\hfill \end{array}.$$
where $`\mu `$ and $`\nu `$ are the natural morphisms. For given $`f`$ and $`g`$ several local isomorphisms $`(a,A)`$ may occur which can make the amalgam a complicated occasionally non-Hausdorff space. This is just the case, if the automorphism group of $`f`$ or of $`g`$ is non-trivial.
Automorphisms. Let $`(f,i)`$ be a miniversal deformation of an analytic polyhedron $`X`$ as above and $`M`$ be the maximal modular stratum in its base $`S`$. Let $`X_K`$ be the germ of $`X`$ on the compact set $`Ksupp𝒯^1\left(X\right)X`$. We show that an arbitrary automorphism $`a`$ of $`X_K`$ generates an automorphism of the stratum $`M`$. Choose a polyhedron $`X^{}X`$ that contains the germ $`X_K`$ such that $`a`$ defines a morphism $`a:X^{}X`$. Take a miniversal deformation $`(f^{},i^{})`$ of $`X^{}`$; let $`S^{}`$ be the base of $`f^{}`$. By Theorem 6.2 there exists an imbedding $`j:S^{}S`$ that induces an isomorphism $`f\times jf^{}`$ and maps the modular stratum $`M^{}`$ to the modular stratum $`M`$. Consider the deformation $`(f,ia)`$. This is a miniversal deformation of $`X^{},`$ hence it is induced from $`(f^{},i^{})`$ by a germ endomorphism $`\alpha :SS^{}`$. It is an automorphism, since $`\mathrm{D}\alpha `$ is bijective at the marked point. The restriction $`\alpha |M`$ is uniquely defined and the composition $`\alpha _Mj^1\alpha |M:MM`$ is an automorphism of the germ $`M`$.
###### Corollary 7.1
The mapping $`a\alpha _M`$ defines a homomorphism of the group $`\mathrm{Aut}\left(X_K\right)`$ to the group $`\mathrm{Aut}\left(M\right)`$ of automorphisms of the maximal modular germ $`M.`$
The kernel $`\mathrm{Aut}_0\left(X_K\right)`$ of this homomorphism is a normal subgroup and the quotient $`\mathrm{G}\left(X_K\right)=\mathrm{Aut}\left(X_K\right)/\mathrm{Aut}_0\left(X_K\right)`$ acts faithfully; we call it active automorphism group. Any automorphism $`a`$ of $`X_K`$ generated by a tangent field $`vT^0\left(X_K\right)`$ acts trivially on $`M.`$ The field $`v`$ generates a local holomorphic subgroup $`a\left(\zeta \right)`$ of automorphisms. The infinitesimal action of this group in $`T_{}(S)`$ is given by the commutator $`t[\mathrm{D}_{}f\left(t\right),v]`$, which is trivial in $`M`$ due to Theorem 5.2.(iii). Therefore $`a_M=\mathrm{id}`$ that is $`a`$ $`\mathrm{Aut}_0\left(X_K\right)`$. The active group G$`\left(X_K\right)`$ is discontinuous in several examples (see below).
The fibers $`X_s`$ are, of course, isomorphic for points $`s`$ in any coset of the automorphism group. The inverse statement is by no means obvious.
Problem 1 Let $`f:𝒳(M,)`$ be a maximal modular deformation of a space $`XX_{}`$. When does the existence of a isomorphism $`X_sX_t`$ for some points $`s,tM`$ imply $`\alpha _M\left(s\right)=t`$ for an element $`a\mathrm{G}\left(X_s\right)`$? In particular, when an isomorphism $`X_sX,sM^{}`$ implies $`s=`$?
Is it true, if the group $`\mathrm{G}(X_s)`$ is finite?
Let $`𝒳M`$ be a modular family with a irreducible base $`M`$ and $`sM`$. By Corollary 7.1 any element $`\alpha \mathrm{G}\left(X_s\right)`$ of the fiber $`X_s=f^1\left(s\right)`$ generates an automorphism $`\alpha _M`$ of the germ $`(M,s)`$. By analytic continuation, we can extend $`\alpha _M`$ to a uniquely defined automorphism $`\alpha `$ of the base $`M`$. Let $`\mathrm{G}`$ be automorphism group of $`M`$ generated by all the elements $`\alpha _M`$ for $`\alpha \mathrm{G}\left(X_s\right),sM`$. The quotient $`M/\mathrm{G}`$ can be taken as a candidate for ‘true’ moduli space.
## 8 Modular deformation of compact spaces
We start with very classical families that appear to be modular.
Example 1. The family $`T=\left\{T\left(\lambda \right)\right\},\lambda _+`$ of $`1`$-tori parameterized by the upper half-plane $`_+`$; $`T\left(\lambda \right)`$ is the quotient of the plane $``$ by the lattice generated by $`1`$ and $`\lambda .`$ This family is maximal modular for each point $`\lambda _+`$. The automorphism group of the family is the group $`\mathrm{Sl}(2,)`$ which acts in $`_+`$ by fractional linear transformations. The group $`\mathrm{Aut}_0\left(T\left(\lambda \right)\right)`$ is equal to the semi-direct product of the group $`T\left(\lambda \right)`$ generated by tangent fields on $`T\left(\lambda \right).`$ The group $`_2`$ generated by the mapping $`zz`$ in $``$ acts trivially in the base $`_+`$. Take the standard fundamental domain $`D=\{\lambda :\left|\mathrm{Re}\lambda \right|<1/2,\left|\lambda \right|>1\}`$ and consider the restriction of the family to $`D.`$ There are two points $`\lambda _2=ı,\lambda _3=\sqrt[3]{1}`$ on the boundary of $`D`$ such that the active automorphism group $`\mathrm{G}(T\left(\lambda \right))`$ has elements with non-trivial action: these are elements $`a_j\mathrm{G}\left(T\left(\lambda _j\right)\right)`$ of order $`2`$, respectively $`3.`$ These elements are generated by rotation of the unit square by $`\pi /2`$ and by rotation of the rhombus by $`\pi /3,`$ respectively. According to previous section, these groups generate transformation group $`\mathrm{G}`$ of $`D.`$ The quotient space $`D/\mathrm{G}`$ is isomorphic to the complex plane $``$ and the family $`T`$ gives rise to a family $`\stackrel{~}{T}`$ of tori on $`.`$ It can be compacted by means of amalgam with the deformation $`f:𝒴U`$ of the singular curve $`Y_0,`$ which is the projective line with one point of transversal self-intersection. Here $`U`$ is the unit disk and the deformation is given in an affine chart by
$$w^2z(zs)(z1)=0,\left|s\right|<1/2$$
The fibers $`Y_s`$ are non-isomorphic tori for $`s0`$, and $`dimT^1\left(Y_0\right)=1,`$ hence the family $`𝒴`$ is maximal modular. It can be patched to $`\stackrel{~}{T}`$ giving rise to a modular family with base $`\{s=0\}=^1`$. The point $`s=0`$ corresponds to infinity, since a ratio of periods of the surface $`Y_s`$ tends to infinity as $`s0`$.
Example 2. Generalize the above construction for curves of an arbitrary genus $`g>1`$. Consider the family of hyperelliptic Riemann surfaces $`X_a,a^m`$, where $`X_a`$ is given in the affine chart by
$$w^2p(z)=0,p(z)=z^m+a_1z^{m1}+\mathrm{}+a_m,m=2g+2.$$
(4)
Suppose that $`g4`$, take the singular surface $`X_0`$ for which $`a_1=\mathrm{}=a_m=0`$ and check that
$$dimT^1(X_0)=3g3$$
(5)
We have
$$dimT^1(X)=dimH^0(X,𝒯^1)+dimH^1(X,𝒯^0),$$
since $`H^2(X,𝒯^1)=0`$. Further,
$$dimH^0(X,𝒯^1)=dim𝒯_s^1=m1=2g+1,$$
since the only singular point $`s=(0,0)`$ is of multiplicity $`m1`$. The surface $`X`$ is union of two spheres that are tangent one to another at the origin. The tangent sheaf $`𝒯^0`$ is generated at the origin by two fields
$$t_1=mz\frac{}{z}+2w\frac{}{w},t_2=2w\frac{}{z}mz^{m1}\frac{}{w}$$
over the algebra over $`𝒪(X)_s`$. It helps to check that $`dimH^1(X,𝒯^0)=g4`$, which implies (5).
Let $`F_0:𝒴S`$ be a miniversal deformation of $`X_0`$. Prove that it is universal. The base $`S`$ is a piece of $`^{3g3}`$ and any non-singular fiber $`Y_s`$ is a Riemann surface of genus $`g`$. Therefore, we have again $`dimT^1(Y_s)=dimH^1(Y_s,𝒯^0)=3g3`$. This implies that the dimension of $`T^1(Y_s)`$ is constant in the deformation $`F,`$ since it is a upper semi-continuous function on the base. By Theorem 5.2.(iv) the deformation $`F`$ is maximal modular and therefore universal. The family of surfaces $`X_a`$ is obviously a deformation of $`X_0`$ which yields that $`\mathrm{dim}T^1(X_a)=3g3`$. Therefore any miniversal deformation $`F_a`$ of $`X_a`$ is universal too. They can be amalgamated together in a maximal modular family $`F`$ with a base $`S`$ which is an open subset in $`^{3g3}`$. On the other hand, $`S`$ contains the base $`^m`$ of the family (4). The set $`DS`$ of critical values of this deformation coincides with the variety $`\mathrm{\Delta }(p)=0`$ in $`^m`$, where $`\mathrm{\Delta }(p)`$ is the discriminant of the polynomial $`p`$.
On the other hand, the Teichm̈uller space $`T(g)`$ is the base of a universal family $`R(g)`$, whose fibers are all non-singular Riemann surfaces of genus $`g`$ with an additional structure. The automorphism group $`\mathrm{\Gamma }(g)`$ of the family $`R(g)`$, called modular group, acts discontinuously in $`T(g)`$. The deformation $`F`$ is amalgamated with the family $`R(g)`$ and the base $`S\backslash D`$ is patched to $`T(g)/\mathrm{\Gamma }(g)`$. This gives a compacting of the space $`T(g)/\mathrm{\Gamma }(g)`$, whereas the discriminant set $`D`$ covers the boundary.
In the cases $`g=2,3`$ the deformation $`F`$ as above is no more modular for the surface $`X_0`$, since $`dimT^1(X_0)>3g3`$. We take the surface $`X(q)`$ for $`q(z)=(z^21)z^{2g}`$ instead. The miniversal deformation of $`X(q)`$ is universal, like in the case $`g=1`$.
This compacting of the space $`T(g)`$ is, apparently, different from that of W.Baily .
Example 3. Fix an integer $`n>1`$ and consider a proper holomorphic mapping of manifolds $`F:TS`$ whose fibers are complex analytic $`n`$-tori, see . The dimension of $`S`$ is equal to $`n^2`$ and the mapping is maximal modular deformation of each fiber as in the case $`n=1.`$ On the other hand, the active group $`\mathrm{G}`$ of $`F`$ is not discontinuous. Moreover, any non-empty open set $`U`$ contains a point $`s`$ such that the coset $`\mathrm{G}sU`$ is infinite. The moduli space $`M/\mathrm{G}`$ is not separable.
Example 4. Take the family $`F:𝒱P(m)`$ of algebraic hypersurfaces of degree $`m`$ in $`^n,n3`$, where $`P(m)`$ is the space of all homogeneous polynomials of degree $`m`$ and $`V(f)F^1(f)`$ is the hypersurface defined by the equation $`f=0`$ in $`^n`$. Take a polynomial $`fP(m)`$ such that the variety $`V(f)`$ is non-singular and choose an affine subspace $`SP(m)`$ that contains the point $`f`$ and is transversal to the linear span of polynomials $`z_if/z_j`$, $`i,j=0,\mathrm{},n.`$ Consider the restriction $`F|S`$ of this family; the Kodaira-Spencer mapping $`DF|S:T_f(S)T^1(V(f))`$ is injective. This mapping is surjective and the family $`F|S`$ is versal, except in the case $`n=3,m=4`$, see . Moreover, this family is maximal modular, since $`dimS=(m+n)!/n!m!(n+1)^2`$ does not depend on $`f`$. The automorphism group of any fiber $`X`$ is finite for the same cases.
In the exceptional case $`n=3,m=4`$ the miniversal deformation $`F`$ of any surface contains non-algebraic fibers $`X_s`$ that are also $`\mathrm{K3}`$-surfaces (the canonical bundle is trivial). The base $`S`$ has dimension $`20`$ and the algebraic fibers only appear for points $`s`$ in the dense subset $`S_{alg}`$ which is a countable union of $`19`$-dimensional subspaces in the base, see more details in . Therefore $`dimT^1(X_s)=20`$ on the germ $`S`$ and the formation $`F`$ is again maximal modular. On the other hand, there is no non-trivial tangent fields on a $`\mathrm{K3}`$-surface, but its automorphism group can be infinite .
## 9 Modular deformations of singular points
The modular deformations of polyhedra with isolated singularities look similar to the above examples. Moreover, there is a relation between deformations in these two categories.
Let $`D^n`$ be the open unit polydisk in a coordinate space $`^n`$ centered at the origin. Take holomorphic functions $`f_1,\mathrm{},f_k`$ in $`\overline{D}^n`$ and consider the analytic polyhedron $`X=\{zD^n,f_1(z)=\mathrm{}=f_k(z)=0\}`$ endowed with the sheaf $`𝒪(X)=𝒪(D_\epsilon ^n)/`$, where $``$ is the sheaf-ideal generated by these functions. The polyhedron $`X`$ satisfies the condition 1, if it contains only finite number of singular points.
Suppose that $`k=1`$, the generating function $`f`$ is weighted homogeneous for certain coordinate system $`z_1,\mathrm{},z_n`$ in $`^n`$ and there is only one singular point $`z=0`$ in $`X`$ (otherwise the dimension of the singular set of $`X`$ is positive). The base $`M`$ of the maximal modular deformation is isomorphic to the subspace of $`T^1\left(X\right)`$ of elements whose weight is equal to that of $`f`$. In particular, $`M`$ is a simple point for singularities of types $`A,D`$ and $`E`$ that are characterized by the inequality $`weight\left(\tau \right)<weight\left(f\right)`$ for all $`\tau T^1\left(X\right)`$.
Assume now that the weights of the coordinates are equal to $`1`$ so that $`f`$ is a homogeneous polynomial of degree $`m>0`$. It defines a non-singular hypersurface $`V(f)`$ in $`^{n1}`$ and any deformation of $`V(f)`$ as in Example 4 generates a deformation of the analytic polyhedron $`X(f)\{zD^n,f(z)=0\}`$. We have
$$T^1(X(f))[z_1,\mathrm{},z_n]/(f/z_1,\mathrm{},f/z_n),$$
(6)
The weight grading generates a grading in the space (6).
###### Proposition 9.1
Suppose that $`n4`$ and the case $`n=m=4`$ is excluded. Then there exists a linear injective mapping
$$j:T^1(V(f))T^1(X(f)),$$
(7)
whose image coincides with the subspace of weight $`m`$.
$``$ By diagram (9.3) of , p.109 we have $`T^1(V(f))=P(m)/J`$, where $`P(m)`$ is the space of polynomials in $`^n`$ of degree $`m`$ and $`J`$ is the linear envelope of the polynomials $`z_if/z_j,i,j=1,\mathrm{},n`$. Note that the term $`H^1(V(f),\theta )`$ vanishes in the diagram due to Proposition 9.2. The isomorphism (6) makes the mapping (7) obvious. $``$
###### Corollary 9.2
For any family of homogeneous polynomials $`\{f_s,sS\}`$ of $`n`$ variables with the only critical point $`z=0`$ the numbers dim $`T^1\left(V(f_s)\right)`$ and $`\mathrm{dim}T^1\left(X(f_s)\right)`$ stay constant. If one of the families is modular, another is also modular.
$``$ For the first statement we note that $`\mathrm{dim}T^1(X(f_s))=(m1)^n.`$ By the previous Proposition $`dimT^1(V(f_s))`$ is equal to the number of monomials of degree $`m`$ minus $`n,`$ that is the dimensions depend only on $`n`$ and $`m.`$ We have the equation for the Kodaira-Spencer mappings $`\mathrm{D}\mathrm{\Phi }=j\mathrm{D}F`$ where $`F`$ and $`\mathrm{\Phi }`$ are the families of polyhedrons and of projective hypersurfaces, respectively. The second statement follows from injectivity of $`j.`$
The modular stratum of a non weighted-homogeneous hypersurface may be singular and have imbedded primary components.
Example 5. Consider the polyhedron in $`D^3`$ defined by the equation
$$f_{p,q,r}(\lambda ;x,y,z)x^p+y^q+z^r+\lambda xyz=0$$
(8)
in the complex coordinates $`x,y,z`$. We denote the singular germ at the origin by $`T_{p,q,r}(\lambda )`$, if the equation holds $`1/p+1q+1/r=1`$. The generating function is then weighted homogeneous and the parameter $`\lambda `$ is the coordinate in a maximal modular deformation, where $`\lambda `$ runs over the complex plane with few gaps, where the polyhedron contains a singular curve. This follows from Theorem 5.2.(iv), since the space of elements $`tT^1\left(X_\lambda \right)`$ such that $`[t,v]=0`$ for all $`vT^0\left(X_\lambda \right)`$ is one-dimensional. There are just three modular families of this type: $`T_{3,3,3}(\lambda ),T_{4,4,2}(\lambda )`$ and $`T_{6,3,2}(\lambda )`$ whose Tyurina numbers $`\tau \left(X\right)dimT^1\left(X\right)`$ are $`8,9,10,`$ respectively.
In the case $`1/p+1q+1/r<1`$ the surfaces (8) are all isomorphic for $`\lambda 0`$. We set $`\lambda =1`$ and use the notation $`T_{p,q,r}`$.
Example 6. Start with the family $`T_{3,3,3}\left(\lambda \right)`$ defined for $`\lambda `$. We have $`\tau \left(X\right)=8`$ and the family is maximal modular for $`\lambda ^327`$. The active group $`\mathrm{G}=_3`$ acts in the family by $`\lambda \epsilon \lambda `$, where $`\epsilon ^3=1`$. Take the deformation $`YN`$ of the germ $`T_{4,3,3}`$ defined in $`D^3\times N`$ by the function
$$F(r,s,t;x,y,z)=x^4+y^3+z^3+\mu xyz+rx^3+s_1y+s_2y^2+t_1z+t_2z^2,$$
where the number $`\mu 0`$ is fixed. It is maximal modular for the base $`NM\times ^3`$ given by the ideal $`I=I_1I_0`$ where the ideal $`I_1=(s_1,s_2,t_1,t_2)`$ defines the line parameterized by the coordinate $`r`$ and we have an isomorphism $`Y_{r,0,0}T_{3,3,3}(\lambda )`$ for $`\lambda =r^{1/3}`$. The ideal $`I_0`$ determines a fat point at the origin in $`^5`$, whose embedding dimension is equal to $`5`$. Therefore the modular deformations $`T_{3,3,3}(\lambda )`$ and $`Y`$ are patched together to a maximal modular deformation with a compact base $`M`$.
Example 7. Consider the family $`𝒳`$ of curves $`X_\lambda D_1^2`$ given by the equation
$$(x^3xy^2)(axby)=0,\lambda =b/a^1.$$
The family is maximal modular for all $`\lambda `$, except in three points $`\lambda =1,0,1`$. These gaps can be patched by means of another modular deformations. B.Martin has considered the family $`Y^7`$ of polyhedra in $`D^2\times ^7`$, given by the equation
$$x^4x^2y^2+s_1x+s_2y+s_3xy+s_4y^2+s_5y^3+s_6xy^2+ty^4+y^5=0.$$
He has shown that this family is maximal modular over the base $`M,`$ whose ideal is $`I(M)J_1J_0`$, where $`J_1=(s_1,\mathrm{},s_6)`$ defines the line and the ideal $`J_0`$ has $`9`$ generators and defines a fat point at the origin. For any point $`(0,t)^7`$, $`t0`$ the germ $`Y_{0,t}`$ is isomorphic to $`X_\lambda `$ for $`\lambda =2\sqrt{t}`$. Therefore the family $`𝒴`$ can be amalgamated with the family $`\left\{X_\lambda \right\}`$ patching the gap $`\lambda =0`$. The gaps $`\lambda =\pm 1`$ are patched in a similar way. The resulting amalgam is a maximal modular family $`\stackrel{~}{T}_{3,3,3}`$ with the compact base $`^1`$.
Note that the Milnor number of a fiber of this family is not constant; it is equal to $`9`$, except in the points $`\lambda =1,0,1,`$ where it equals $`10.`$
Example 8. Consider the family $`𝒵`$ of complete intersection curves $`Z_sD^3\times `$ defined by two polynomials:
$$x^4+y^4+2z^2=szxy=0,s.$$
It is modular with $`\tau \left(Z_s\right)=9`$, . The curve $`Z_s`$ is isomorphic for $`s0`$ to the plane curve $`x^4+y^4+2s^2x^2y^2=0`$. The germ at the origin of this curve is isomorphic to the fiber $`X_\lambda `$ of the family in Example 7 with $`\lambda =s^2`$. Now the family $`𝒵`$ is amalgamated with $`𝒳`$ mending the gap at the point $`\lambda =0`$. This provides a compacting of the family $`𝒳`$ at this point, which is different form that of Example 7.
Example 9. () The maximal modular deformations of polyhedra $`Y_{7,3,2},Y_{6,4,2}`$ and $`Y_{6,3,3}`$ can be amalgamated with the modular family $`T_{6,3,2}(\lambda )`$. The modular deformation of $`Y_{6,4,2}`$ is given by the function
$$x^6+y^4+z^2+xyz+rx^4+sx^5+ty^3+vz=0$$
in $`D^3\times ^4`$. It is maximal modular over the base germ $`M`$ defined by the ideal $`I(M)=J_1J_2J_3,`$ where
$$J_1=(r,s,v),J_2=(r,s^2,t^2,20vst),J_3=(t,4rs^2,v).$$
The zero set of $`J_1`$ is the line with the coordinate $`t`$; the fiber $`Y(0,0,t,0)`$, $`t0`$ is a surface whose germ is isomorphic to $`T_{6,3,2}(\lambda )`$ for $`\lambda =t^{1/3}.`$ The ideal $`J_2`$ defines the fat point at the origin and $`J_3`$ defines a smooth curve parameterized by $`s_5`$; the fiber of deformation over this curve is the union of the germ $`T_{4,4,2}(\lambda )`$ for some $`\lambda `$ and of a singular germ of multiplicity $`1`$ at the point $`(s/2,0,0)`$.
The maximal modular deformation $`Y_{6,3,3}`$ has the base, which is the union of three components. Two of them are straight lines that are symmetric under the permutation $`xy`$; the fibers are of type $`T_{6,3,2}(\lambda )`$. The third component is a curve; the general fiber of the deformation over this curve has the singular germ of type $`Y_{4,3,3}`$ at the origin and another singular point of multiplicity $`2`$.
The modular deformation $`Y_{7,3,2}`$ is similar to the previous two cases.
Problem 2. These examples give rise to the general questions: let $`f^{}`$ be a modular deformation over the punctured disk $`D^{}`$. When there exists a modular deformation $`f`$ over $`D`$ such that $`f|D^{}f^{}`$? Note that the deformation $`f`$ need not to be unique as we saw in Examples 8 and 9.
## 10 Concluding remarks
Let $`X`$ be an arbitrary hypersurface germ with isolated singularity which is not weighted homogeneous. Then the tangent space $`T\left(M\right)`$ to the maximal modular stratum has always positive embedding dimension. Indeed, the action of the Lie algebra $`T^0\left(X\right)`$ in $`T^1\left(X\right)`$ is nilpotent and by Engel’s theorem, there is an element $`tT^1\left(X\right)`$, $`t0`$ that vanishes under action of the Lie algebra. By Theorem 5.2 this yields $`tT\left(M\right)`$. This is the case for any singular germ of type $`T_{p,q,r}`$.
More examples can be extracted from the classification from Arnold’s list , where several families of singularities with constant Milnor number are listed. Such a family has also constant Tyurina number, except for a proper subvariety of the family base, where this number jumps up. Restricting the family to a certain affine subvariety $`M`$, yields a modular deformation. The ‘gaps’ in $`M`$ could be filled by means of amalgams with appropriate modular deformations. The completed variety $`\stackrel{~}{M}`$ is expected to be a projective algebraic. Besides, any ‘solitary’ singularity in Arnold’s list which is not included in families may have the maximal modular deformation with the base $`M`$ which is either a fat point or a germ of positive dimension with splitting singularity like that of $`T_{6,4,2}`$ and $`T_{6,3,3}`$ in Example 9.
A.Aleksandrov has found a series of modular families of complete intersection curves in $`^3`$. Examples of families with constant Tyurina number $`34`$ were calculated in the book . B.Martin . T.Hirsch and B.Martin are found sophisticated examples of modular deformations.
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# Intra-cellular transport of single-headed molecular motors KIF1A
## Abstract
Motivated by experiments on single-headed kinesin KIF1A, we develop a model of intra-cellular transport by interacting molecular motors. It captures explicitly not only the effects of ATP hydrolysis, but also the ratchet mechanism which drives individual motors. Our model accounts for the experimentally observed single molecule properties in the low density limit and also predicts a phase diagram that shows the influence of hydrolysis and Langmuir kinetics on the collective spatio-temporal organization of the motors. Finally, we provide experimental evidence for the existence of domain walls in our in-vitro experiment with fluorescently labeled KIF1A.
Intra-cellular transport of a wide variety of cargo in eucaryotic cells is made possible by motor proteins, like kinesin and dynein, which move on filamentary tracks called microtubules (MT) howard ; schliwa . However, often a single MT is used simultaneously by many motors and, in such circumstances, the inter-motor interactions cannot be ignored. Fundamental understanding of these collective physical phenomena may also expose the causes of motor-related diseases (e.g., Alzheimer’s disease) hirotaked thereby helping, possibly, also in their control and cure. Some of the most recent theoretical models of interacting molecular motors frey ; santen ; popkov ; lipo utilize the similarities between molecular motor traffic on MT and vehicular traffic on highways css both of which can be modelled by appropriate extensions of driven diffusive lattice gases sz ; schuetz . In those models the motor is represented by a self-driven particle and the dynamics of the model is essentially an extension of that of the asymmetric simple exclusion processes (ASEP) sz ; schuetz that includes Langmuir-like kinetics of adsorption and desorption of the motors. In reality, a motor protein is an enzyme whose mechanical movement is loosely coupled with its biochemical cycle. In this letter we consider specifically the single-headed kinesin motor, KIF1A okada1 ; okada3 ; unpub ; Nitta ; the movement of a single KIF1A motor has been modelled recently with a Brownian ratchet mechanism julicher ; reimann . In contrast to the earlier models frey ; santen ; popkov ; lipo of molecular motor traffic, which take into account only the mutual interactions of the motors, our model explicitly incorporates also the Brownian ratchet mechanism of individual KIF1A motors, including its biochemical cycle that involves adenosine triphosphate(ATP) hydrolysis.
The ASEP-like models successfully explain the occurrence of shocks. But since most of the bio-chemistry is captured in these models through a single effective hopping rate, it is difficult to make direct quantitative comparison with experimental data which depend on such chemical processes. In contrast, the model we propose incorporates the essential steps in the biochemical processes of KIF1A as well as their mutual interactions and involves parameters that have one-to-one correspondence with experimentally controllable quantities.
The biochemical processes of kinesin-type molecular motors can be described by the four states model shown in Fig. 1 okada1 ; Nitta : bare kinesin (K), kinesin bound with ATP (KT), kinesin bound with the products of hydrolysis, i.e., adenosine diphosphate(ADP) and phosphate (KDP), and, finally, kinesin bound with ADP (KD) after releasing phosphate. Recent experiments okada1 ; Nitta revealed that both K and KT bind to the MT in a stereotypic manner (historically called “strongly bound state”, and here we refer to this mechanical state as “state 1”). KDP has a very short lifetime and the release of phosphate transiently detaches kinesin from MT Nitta . Then, KD re-binds to the MT and executes Brownian motion along the track (historically called “weakly bound state”, and here referred to as “state 2”). Finally, KD releases ADP when it steps forward to the next binding site on the MT utilizing a Brownian ratchet mechanism, and thereby returns to the state K.
Model definition. — A single protofilament of MT is modelled by a one-dimensional lattice of $`L`$ sites each of which corresponds to one KIF1A-binding site on the MT; the lattice spacing is equivalent to $`8`$ nm which is the separation between the successive binding sites on a MT howard . Each kinesin is represented by a particle with two possible internal states labelled by the indices $`1`$ and $`2`$. Attachment of a motor to the MT occurs stochastically whenever a binding site on the latter is empty. Attachment and detachment at the two ends of the lattice need careful treatment and will be specified below. Thus, each of the lattice sites can be in one of three possible allowed states (Fig. 2): empty (denoted by $`0`$), occupied by a kinesin in state $`1`$, or occupied by a kinesin in state $`2`$.
For the dynamical evolution of the system, one of the $`L`$ sites is picked up randomly and updated according to the rules given below together with the corresponding probabilities (Fig. 2):
$`\mathrm{Attachment}:\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}1\mathrm{with}\omega _adt`$ (1)
$`\mathrm{Detachment}:\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}0\mathrm{with}\omega _ddt`$ (2)
$`\mathrm{Hydrolysis}:\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}1}2\mathrm{with}\omega _hdt`$ (3)
$`\mathrm{Ratchet}:\{\begin{array}{c}21\mathrm{with}\omega _sdt\\ 2001\mathrm{with}\omega _fdt\end{array}`$ (6)
$`\mathrm{Brownian}\mathrm{motion}:\{\begin{array}{c}2002\mathrm{with}\omega _bdt\\ 0220\mathrm{with}\omega _bdt\end{array}`$ (9)
The probabilities of detachment and attachment at the two ends of the MT may be different from those at any bulk site. We choose $`\alpha `$ and $`\delta `$, instead of $`\omega _a`$, as the probabilities of attachment at the left and right ends, respectively. Similarly, we take $`\gamma _1`$ and $`\beta _1`$, instead of $`\omega _d`$, as probabilities of detachments at the two ends (Fig. 2). Finally, $`\gamma _2`$ and $`\beta _2`$, instead of $`\omega _b`$, are the probabilities of exit of the motors through the two ends by random Brownian movements.
Let us relate the rate constants $`\omega _f`$, $`\omega _s`$ and $`\omega _b`$ with the corresponding physical processes in the Brownian ratchet mechanism of a single KIF1A motor. Suppose, just like models of flashing ratchets julicher ; reimann , the motor “sees” a time-dependent effective potential which, over each biochemical cycle, switches back and forth between (i) a periodic but asymmetric sawtooth like form and (ii) a constant. The rate constant $`\omega _h`$ in our model corresponds to the rate of the transition of the potential from the form (i) to the form (ii). The transition from (i) to (ii) happens soon after ATP hydrolysis, while the transition from (ii) to (i) happens when ATP attaches to a bare kinesinokada1 . The rate constant $`\omega _b`$ of the motor in state $`2`$ captures the Brownian motion of the free particle subjected to the flat potential (ii). The rate constants $`\omega _s`$ and $`\omega _f`$ are proportional to the overlaps of the Gaussian probability distribution of the free Brownian particle with, respectively, the original well and the well immediately in front of the original well of the sawtooth potential.
Let us denote the probabilities of finding a KIF1A molecule in the states $`1`$ and $`2`$ at the lattice site $`i`$ at time $`t`$ by the symbols $`r_i`$ and $`h_i`$, respectively. In mean-field approximation the master equations for the dynamics of motors in the bulk of the system are given by
$`{\displaystyle \frac{dr_i}{dt}}`$ $`=`$ $`\omega _a(1r_ih_i)\omega _hr_i\omega _dr_i+\omega _sh_i`$ (10)
$`+\omega _fh_{i1}(1r_ih_i),`$
$`{\displaystyle \frac{dh_i}{dt}}`$ $`=`$ $`\omega _sh_i+\omega _hr_i\omega _fh_i(1r_{i+1}h_{i+1})`$ (11)
$`\omega _bh_i(2r_{i+1}h_{i+1}r_{i1}h_{i1})`$
$`+\omega _b(h_{i1}+h_{i+1})(1r_ih_i).`$
The corresponding equations for the boundaries, which depend on the rate constants $`\alpha `$, $`\delta `$, $`\gamma _i`$ and $`\beta _i`$ for entry and exit (Fig. 2), are similar and will be presented elsewhere unpub .
From experimental data okada1 ; okada3 , good estimates for the parameters of the suggested model can be obtained. Assuming that one timestep corresponds to 1 ms, each simulation run had a duration of 1 minute in real time. The length of MT is fixed as $`L=600`$. The detachment rate $`\omega _d0.0001`$ ms<sup>-1</sup> is found to be independent of the kinesin population. On the other hand, $`\omega _a=10^7`$ $`C`$/M$``$s depends on the concentration $`C`$ (in M) of the kinesin motors. In typical eucaryotic cells in-vivo the kinesin concentration can vary between 10 and 1000 nM. Therefore, the allowed range of $`\omega _a`$ is $`0.0001`$ ms$`{}_{}{}^{1}\omega _a0.01`$ ms<sup>-1</sup>. The rate $`\omega _b^1`$ must be such that the Brownian diffusion coefficient $`D`$ in state 2 is of the order of $`40000`$ nm<sup>2</sup>/s; using the the relation $`\omega _bD/(8\text{nm})^2`$, we get $`\omega _b0.6`$ ms<sup>-1</sup>. Moreover, from the experimental observations that $`\omega _f/\omega _s3/8`$ and $`\omega _s+\omega _f0.2`$ ms<sup>-1</sup>, we get the individual estimates $`\omega _s0.145`$ ms<sup>-1</sup> and $`\omega _f0.055`$ ms<sup>-1</sup>. The experimental data on the Michaelis-Menten type kinetics of hydrolysis howard suggest that
$$\omega _h^1\left[4+9\left(\frac{0.1\text{mM}}{\text{ATP concentration (in mM)}}\right)\right]\text{ms}$$
(12)
so that the allowed biologically relevant range of $`\omega _h`$ is $`0\omega _h0.25`$ ms<sup>-1</sup>.
Single-molecule properties. — An important test for the model is provided by a quantitative comparision of the low density properties with empirical results. Single molecule experiments okada1 on KIF1A have established that
(i) $`v`$, the mean speed of the kinesins, is about $`0.2`$ nm/ms if the supply of ATP is sufficient, and that $`v`$ decreases with the lowering of ATP concentration following a Michaelis-Menten type relation like (12);
(ii) $`D/v190`$ nm, irrespective of the ATP concentration, where $`D`$ is the diffusion constant;
(iii) $`\tau `$, the mean duration of the movement of a kinesin on the MT, is more than $`5`$ s, irrespective of the ATP concentration.
The corresponding predictions of our model (see Table 1) for $`\omega _a=\alpha =1.0\times 10^6`$ ms<sup>-1</sup>, which allows realization of the condition of low density of kinesins, are in excellent agreement with the experimental results.
Collective properties. — Assuming periodic boundary conditions, the solutions $`(r_i,h_i)=(r,h)`$ of the mean-field equations (11) in the steady-state are found to be
$`r`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_h\mathrm{\Omega }_s(\mathrm{\Omega }_s1)K+\sqrt{D}}{2K(1+K)}},`$ (13)
$`h`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }_h+\mathrm{\Omega }_s+(\mathrm{\Omega }_s+1)K\sqrt{D}}{2K}}`$ (14)
where $`K=\omega _d/\omega _a`$, $`\mathrm{\Omega }_h=\omega _h/\omega _f`$, $`\mathrm{\Omega }_s=\omega _s/\omega _f`$, and
$$D=4\mathrm{\Omega }_sK(1+K)+\left(\mathrm{\Omega }_h+\mathrm{\Omega }_s+(\mathrm{\Omega }_s1)K\right)^2.$$
(15)
The probability of finding an empty binding site on a MT is $`Kr`$ as the stationary solution satisfies the equation $`r+h+Kr=1`$. The steady-state flux of the motors along their MT tracks is then given by $`J=\omega _fh(1rh).`$ It is interesting to note that in the low ATP concentration limit ($`\omega _h\omega _s\omega _f`$) of our model, the flux of the motors is well approximated by $`J_{\mathrm{low}}=q_{\text{eff}}\rho (1\rho )`$, which formally looks like the corresponding expression for the totally asymmetric exclusion process, where $`\rho `$ is close to the Langmuir limit $`1/(1+K)`$ and,
$$q_{\text{eff}}=\frac{\omega _h(1+K)}{\mathrm{\Omega }_s(1+K)+K}$$
(16)
as the effective hopping probabilityunpub .
Although the system with periodic boundary conditions is fictitious, the results provide good estimates of the density and flux in the corresponding system with open boundary conditions, particularly, in the high $`\omega _a`$ regime (Fig. 3) which corresponds to jammed traffic of kinesin on MT (see Fig. 4). We also see that, for a given $`\omega _a`$, the bulk density of motors in state 2 exceeds that of those in state 1 as $`\omega _h`$ increases beyond a certain value.
Phase diagram. — In contrast to the phase diagrams in the $`\alpha \beta `$-plane reported by earlier investigators frey ; santen ; lipo , we have drawn the phase diagram of our model (Fig. 4) in the $`\omega _a\omega _h`$ plane by carrying out extensive computer simulations for realistic parameter values of the model with open boundary conditions.
The phase diagram shows the strong influence of hydrolysis on the spatial distribution of the motors along the MT. For very low $`\omega _h`$ no kinesins can exist in state 2; the kinesins, all of which are in state 1, are distributed rather homogeneously over the entire system. In this case the only dynamics present is due to the Langmuir kinetics.
Even a small, but finite, rate $`\omega _h`$ is sufficient to change this scenario. In this case both the density profiles $`\rho _j^1`$ and $`\rho _j^2`$ of kinesins in the states 1 and 2 exhibit a shock. As in the case of the ASEP-like models with Langmuir kinetics frey ; santen , these shocks are localized. In computer simulations we have observed that the shocks in density profiles of kinesins in the states 1 and 2 always appear at the same position. Note that if the individual density profiles $`\rho _j^1`$ and $`\rho _j^2`$ exhibited shocks at two different locations, two shocks would appear in the total density profile $`\rho _j=\rho _j^1+\rho _j^2`$ violating the usual arguments shockform that ASEP-type models exhibit exactly one shock. Moreover, we have found that the position of the immobile shock depends on the concentration of the motors as well as that of ATP; the shock moves towards the minus end of the MT with the increase of the concentration of kinesin or ATP or both (Fig. 4).
Finally, we present direct experimental evidence that support of the formation of the shock. The “comet-like structure”, shown in the middle of Fig. 5, is the collective pattern formed by the red fluorescent labelled kinesins where a domain wall separates the low-density region from the high-density region. The position of the domain wall depends on both ATP and KIF1A concentrations. Moreover, as we increase the concentration of KIF1A, the transition from the regime of free flow of kinesins to the formation of the shock is observed(top and middle in Fig. 5). Furthermore, we observe jammed traffic of kinesins at sufficiently high concentration (bottom in Fig. 5). The position of the shock in our simulation agrees well with the location of the domain wall in the comet-like structure observed in experimentsunpub .
In this letter we have developed a stochastic model for the collective intra-cellular transport by KIF1A motors, by taking into account the biochemical cycle of individual motors involving ATP hydrolysis and their mutual steric interactions. We have been able to identify the biologically relevant ranges of values of all the model parameters from the empirical data. In contrast to some earlier oversimplified models, the predictions of our model are in good quantitative agreement with the corresponding experimental data. Moreover, we have mapped the phase diagram of the model in a plane spanned by the concentrations of ATP and KIF1A, both of which are experimentally controllable quantities. Finally, we have reported the experimental observation of a comet-like collective pattern formed by the kinesin motors KIF1A and identified the domain wall in the pattern with the shock predicted by our model.
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# Precession of the Isolated Neutron Star PSR B1828–11
## 1 Introduction
Pulse arrival times of neutron stars can be found very accurately, which allows for the determination of the spin period and period derivative to very high precision. Normally, the time of arrival *residuals* which are calculated by subtracting the period and the period derivative (and in some cases the period second derivative) are mostly white noise. However, residuals from a small number of rotating neutron stars are found to exhibit long term cyclical, but non-oscillatory, variations with characteristic timescales of order months to years (Cordes 1993). The variability may be temporary (e.g. the Vela pulsar during its Christmas glitch \[McCulloch et al. 1990\]) or persistent (e.g. the accreting neutron star Her X-1 \[Tannanbaum et al. 1972\], the Crab pulsar \[Lyne, Pritchard and Smith 1988\], and the pulsars PSR 1642-03 \[Blaskiewicz 1992\], PSR B0959-54 \[D’Alessandro and McCulloch 1997\] and PSR B1828–11 \[Stairs, Lyne and Shemar 2000\]). The long timescales that characterize the observed variations would arise naturally from precession<sup>1</sup><sup>1</sup>1Throughout this paper, we call this phenomenon precession, as has become common in the literature, although purists might prefer the term nutation., when the principal axes of a body (defined through the moments of inertia, which we will take as $`I_1I_2I_3`$) revolve periodically around the angular momentum, as viewed in an inertial frame. An ellipticity $`ϵ=(I_3I_1)/I_11`$ would be expected to produce variations in the timing residuals of an axisymmetric body with a period $`P_p=P_{}/ϵ3.2P_{}(\mathrm{sec})(10^8ϵ)^1`$ years, where $`P_{}`$ is the rotation period. The arrival time variations characteristic of precession would be strictly periodic, but not sinusoidal for a triaxial rotator.
There are two physical causes for time of arrival residuals ($`\mathrm{\Delta }t`$) in a precessing neutron star (Cordes 1993). One is directly geometrical: as the rotating star precesses, the symmetry axis of its radiation beam crosses the plane defined by the angular momentum of the star and the direction to the observer at times that vary periodically over the precession cycle. The magnitude of the variability is set by the amplitude of the precession, which is roughly the *wobble angle*, $`\theta `$ (defined as the angle between the angular momentum and the principal axis corresponding to the largest moment of inertia, Fig. 1). Typically $`\theta 1`$, and the amplitude of the time of arrival residuals is $`\mathrm{\Delta }t_{\mathrm{geo}}\theta P_{}`$. In addition, the dependence of the spindown torque acting on the pulsar on the angle between its spin and magnetic axes produces a timing residual that can be comparable to and even exceed $`\mathrm{\Delta }t_{\mathrm{geo}}`$. If we assume that the spindown torque is proportional to $`\widehat{𝛀}a(\widehat{𝛀}\widehat{𝐛})\widehat{𝐛}`$, where $`𝛀=\mathrm{\Omega }\widehat{𝛀}`$ is the angular velocity, $`\widehat{𝐛}`$ is the magnetic axis, and the dimensionless parameter $`a`$ is a measure of the angular dependence ($`a1`$ for a spinning magnetic dipole radiating into vacuum) then the spindown rate varies over the precession cycle as well, producing a timing residual $`\mathrm{\Delta }t_{\mathrm{sd}}a\theta P_p^2/t_{\mathrm{sd}}(aP_p^2/P_{}t_{\mathrm{sd}})\mathrm{\Delta }t_{\mathrm{geo}}`$, where $`t_{\mathrm{sd}}`$ is the spindown timescale for the pulsar; the dimensionless parameter $`\mathrm{\Gamma }_{\mathrm{sd}}P_p^2/P_{}t_{\mathrm{sd}}3.2P_p^2(\mathrm{years})/P_{}(\mathrm{sec})(t_{\mathrm{sd}}/10^7\mathrm{years})`$ may be large. Associated with these arrival time residuals are period residuals $`(\mathrm{\Delta }P/P_{})_{\mathrm{geo}}\theta P_{}/P_p3.2\times 10^8\theta P_{}(\mathrm{sec})/P_p(\mathrm{years})`$ and $`(\mathrm{\Delta }P/P_{})_{\mathrm{sd}}a\mathrm{\Gamma }_{\mathrm{sd}}(\mathrm{\Delta }P/P_{})_{\mathrm{geo}}`$.
The best candidate to date for truly periodic long term variations in arrival times is PSR B1828–11 (Stairs, Lyne & Shemar 2000; Stairs et. al. 2003). Fourier analysis of these variations reveals harmonically related periodicities at approximately $`1000`$, $`500`$ and $`250`$ days (Stairs, Lyne & Shemar 2000; Fig. 2), with the latter two somewhat more pronounced than the first. The length of the timescale of these variations implies that they are probably not of magnetospheric origin, since the natural timescale in the magnetosphere is of the order of the spin period, which in this case is $`P_{}=0.405`$ sec. Even the $`𝑬`$$`\mathbf{\times }`$$`𝑩`$ drift of subpulses (e.g. Ruderman & Sutherland 1975) does not exceed $`10P_{}`$. (However, Ruderman has suggested the possibility of drifts with periods of the order of a year.) As of the time of writing, there are no quantitative models for the data based solely on magnetospheric effects, but there are successful models based on precession (e.g. Jones & Andersson 2001, Link & Epstein 2001, and Wasserman 2003). Link & Epstein (2001) previously modelled the timing residuals from this pulsar in terms of precession of an axisymmetric, oblate rotating rigid body slowing down according to the vacuum magnetic dipole radiation formula. They found that the observations could be accounted for in this model provided that the underlying pulsar is nearly an orthogonal rotator (magnetic obliquity $`\chi 89^{}`$ to the body’s symmetry axis) and nearly aligned angular momentum (wobble angle $`\theta 3^{}`$ between angular velocity and symmetry axis). These are in accordance with the conclusions reached by Jones & Andersson (2001). Although the precession amplitude is small, it may suffice to unpin superfluid vortex lines (Link & Cutler 2002), thus avoiding a potential impediment to precession: pinning was shown to shorten the precession period to about 100 spin periods, and precession itself is dissipated over a timescale of 100-10,000 precession periods (Shaham 1977, 1986; Sedrakian, Wasserman & Cordes 1999). Wasserman (2003) argued that the data could also be accounted for if the underlying neutron star has either a type II superconductor or a strong toroidal magnetic field in its core. In these models, the angle between the angular velocity and symmetry axis of the star could be larger than found by Link & Epstein (2001), and the star does not have to be a nearly orthogonal rotator; spindown variations were found to dominate the timing residuals in this case as well. Link (2003) showed that the standard picture of the core of type II superconducting protons coexisting with superfluid neutrons is inconsistent with long-period precession; pinning of the neutron vortices to the proton flux tubes makes the precession frequency comparable to the rotation frequency of the star, a factor of $`10^8`$ too fast. Possible implications include a normal core (both neutrons and protons), superfluid neutrons and normal protons, normal neutrons and superfluid protons (type I or type II), or superfluid neutrons and type I protons. Sedrakian (2005) studied the last possibility. He calculated the drag on neutron vortices moving in a type I superconducting core, and found that the drag is sufficiently small (that is, the vortices are sufficiently mobile with respect to the protons) that long-period precession is indeed possible in this scenario. Irrespective of the details, magnetic stresses in excess of the relatively weak ones that would arise from the pulsar’s apparent dipole field strength, together with crustal stresses, would render the neutron star effectively triaxial in shape (Cutler 2002; Wasserman 2003; Cutler, Ushomirsky & Link 2003).
Link & Epstein (2001) and Wasserman (2003) gave two alternative models that interpret the timing of PSR B1828–11 as precession. These models fit the data well, thus providing strong evidence that the observed timing variations do indeed represent precession. These two models, however, are special cases. The purpose of this paper is to do a thorough search of the parameter space to see what we can learn about the properties of the spindown torque and the stellar figure. To this end, we analyze the period residuals from this pulsar in terms of a simple model in which the rotating neutron star is assumed to be a triaxial rigid body. Obviously, precise axisymmetry is a special case, and we do not expect it to hold generally, particularly if the crust of the star is not in a relaxed state (e.g. Cutler, Ushomirsky & Link 2003), or has substantial internal magnetic stresses that may not be axisymmetric to begin with. Thus, one of our goals is to see what the data from PSR B1828–11 reveal about the shape of the neutron star crust.
In this paper, we model a precessing neutron star as a single (rigid) body, rotating uniformly. Realistic neutron star modelling should take account of at least two different components – its solid crust, and (super)fluid core. Bondi & Gold (1955) considered the precession of a body consisting of a solid crust coupled frictionally to a fluid core. Their work showed that the long-term precession of the composite system depends on the timescale $`t_{cc}`$ on which the crust and core couple to one another. If $`\mathrm{\Omega }t_{cc}1`$ then the crust and core are very tightly coupled to one another on timescales smaller than a rotation period, and the moment of inertia tensor relevant to precession is that of the entire system, crust plus core. In this case, precession damps out slowly, on a timescale $`(\mathrm{\Omega }t_{cc})^1`$ precession periods. If $`\mathrm{\Omega }t_{cc}1`$ the crust and core only couple on timescales long compared to a rotation period, and the moment of inertia tensor relevant to precession is that of the crust alone. In this case, precession also damps slowly, with a characteristic decay timescale $`\mathrm{\Omega }t_{cc}`$ precession periods. Estimates of the crust-core coupling timescale vary but the consensus is that the coupling is weak, with $`\mathrm{\Omega }t_{cc}10^210^41`$ (e.g. Alpar, Langer & Sauls 1984, Alpar & Sauls 1988, Sedrakian & Sedrakian 1995), so the precession dynamics are governed by the moment of inertia tensor of the crust alone. However, it also turns out that as long as the crust and core couple on a timescale short compared with a precession period, but still long compared to the spin period, the relevant moment of inertia for all spindown effects, including those that vary periodically over a precession period, is the total stellar moment of inertia (Akgun, Link & Wasserman 2005, in preparation); this is the appropriate regime as long as $`\mathrm{\Omega }t_{cc}10^8`$, which appears to be the case. Thus, our one component model for precession is justified, apart from slow decay of the precession, which we neglect.
Another issue is that the neutron star crust is not perfectly rigid, but has a finite shear modulus. For a biaxial precessing star, the crust must be strained in order for the star to precess with a period of order a year (Cutler, Ushomirsky & Link 2003). In addition, the strain field will vary with time as the star precesses, making the star’s moment of inertia tensor time-dependent in the body frame. For simplicity, we neglect these effects, and assume that the rotation of an imperfectly rigid, triaxial star is well-described by the Euler equations for a rigid body, but with a moment of inertia tensor that is rescaled to account for the finite shear modulus.
Since $`P_{}=0.405`$ seconds and $`t_{\mathrm{sd}}10^5`$ years for PSR B1828–11, $`\mathrm{\Gamma }_{sd}10^3`$, and the spindown contribution to the precession-induced timing residuals is particularly important. As a result, we may also hope to use these data to probe the value of $`a`$, that is, to probe the angular dependence of the spindown torque. While it is common to assume that $`a1`$ for rough analysis, on theoretical grounds we should not expect this to be true, for even an aligned rotator surrounded by a magnetosphere radiates energy, a process whose source is ultimately the rotational energy of the star, resulting in spindown at a rate presumably not much different from its luminosity divided by its rotational frequency. One of our chief findings is evidence that the external torque that spins down a pulsar does indeed possess at least some angular dependence.
Our analysis uses the same segment of the data<sup>2</sup><sup>2</sup>2We thank I. H. Stairs, A. G. Lyne and S. L. Shemar for generously sharing their timing residual data with us. that was the basis of Link & Epstein (2001); this facilitates direct comparisons between the results of the two studies. We focus on the period residuals because we can derive an analytic formula for them in terms of elliptic functions. (An analogous formula for an oblate axisymmetric rotator has been derived previously by Bisnovatyi-Kogan, Mersov & Sheffer 1990, and Bisnovatyi-Kogan & Kahabka 1993, and has been applied to the 35-day cycle of Her X-1.) Because the underlying triaxial model involves numerous parameters, using an analytic formula speeds up the computation considerably, which is a distinct advantage. By contrast, direct analysis of the timing residuals would require numerical integration of the model equations, a distinct disadvantage. Thus, for computational convenience, we analyze the period residuals rather than the arrival time residuals. However, a straightforward analysis of the period residuals using their tabulated uncertainties yielded very large values of $`\chi ^2`$ ($`\mathrm{a}\mathrm{few}\times 10^3`$) under the assumption that the residuals are due solely to precession. This indicated to us that there is extra noise in the period residuals, either because their estimated uncertainties are too small (which we regard as unlikely to account for all the noise) or because there is a physical source of period noise that smears out the smooth contribution from precession systematically. In order to account for this extra noise simply, we multiplied the tabulated uncertainties by a (single) factor $`F`$, and then marginalized over $`F`$ to obtain posterior distributions of the (more interesting) parameters of the precession model. This method – Student’s t-test – represents a computational realization of “chi-by-eye” for data whose uncertainties may only be known incompletely. Details are given in Appendix D.
Section 2 contains basic features of our model; further details may be found in Appendices A, B and C. Section 3 contains results and implications of our analysis; statistical details, including the “chi by eye” method mentioned above, are found in Appendix D. Section 4 is a short digression on the pulse shape of PSR B1828–11, which is seen to vary systematically with precession phase (Stairs et al. 2000). Although we do not use this information in our statistical analysis, precession samples different regions in a pulsar’s radio-emitting region, and offers the possibility of mapping out its shape (as has been done for PSR 1913+16, which exhibits geodetic precession, by Weisberg, Romani & Taylor , and previously by Link & Epstein for PSR B1828–11).
## 2 Models
Here we review the models that we use briefly, highlighting some of the more important parameters. The derivations for the period residuals are lengthy and are left for the Appendices; we present the geometric model in Appendix A, and the spindown model is derived in Appendix C. In what is next, we will follow the notation used in these Appendices. The parameters of our models are listed in Table 1.
### 2.1 Geometric Model
What we refer to as the *geometric* model is the effect of triaxiality alone (i.e. torque-free precession). In this case, Euler’s equation for the angular momentum can be solved analytically in terms of Jacobian elliptic functions (Landau & Lifshitz, 1976). As we show in Appendix A, the period residuals are then found to be of the form $`\mathrm{\Delta }P_{ge}/P_{}\varpi _pf_n`$ where $`P_{}`$ is the rotation period at a fiducial epoch, $`\varpi _p`$ is a dimensionless quantity of order $`P_{}/P_pϵ`$, and $`f_n`$ is a complicated combination of the elliptic functions. Because of the inherent form of $`f_n`$ the amplitude of the residuals is not trivial to predict in general, and they can exhibit very rich behavior.
We denote the principal moments of inertia by $`I_i`$; the associated axes serve as the basis for the rotating (body) frame. We then define the following parameters: $`ϵ=(I_3I_1)/I_1`$ which measures the deviation from sphericity; $`e^2=[I_3(I_2I_1)]/[I_1(I_3I_2)]`$ which measures the degree of triaxiality; $`k^2`$ which is the parameter of the Jacobian elliptic functions, and depends on the angular momentum and the moments of inertia; and $`\lambda `$ which determines the components of the angular momentum (and would be simply $`\lambda =L_1/L_3`$ for an axisymmetric star, but is slightly different in the more general case, see Eq. (12)). The last three are not independent: $`k=e\lambda `$. Note that $`k^2`$ does not depend on $`ϵ`$, but only on $`e^2`$ and $`\lambda `$.
Note that for $`e^2=0`$ the body is *oblate* and axisymmetric, and $`\lambda =0`$ when there is no precession. In both cases $`k=0`$, and the Jacobian elliptic functions reduce to the regular trigonometric functions. At $`k=1`$ they become hyperbolic functions, and the angular momentum exponentially aligns with the principal axis corresponding to the intermediate moment of inertia, $`I_2`$. Between these two extremes, the precession of the angular momentum takes place along the intersection of the sphere defined by the conservation of angular momentum ($`L^2=L_iL_i`$), and the ellipsoid defined by the conservation of energy ($`E=L_i^2/2I_i`$). The resulting shape of the trajectories is the Binet ellipsoid (Landau & Lifshitz, 1976). In the limit $`e^2\mathrm{}`$ the body becomes *prolate* axisymmetric.
### 2.2 Spindown Model
The rotation of an isolated neutron star is not torque-free, but slows down, resulting in a gradual increase in the rotation period. It is thought that because of the rotating magnetic field, angular momentum is lost to radiation. In the simple model of a rotating dipole in a vacuum, the pulses are emitted at the poles of the magnetic field and the torque has the form $`𝐍N_o[\widehat{𝛀}(\widehat{𝛀}\widehat{𝐛})\widehat{𝐛}]`$, where $`\widehat{𝛀}`$ is the instantaneous rotation axis, and $`\widehat{𝐛}`$ is the pulse (and magnetic) axis.
This is a very crude model. The magnetic field may have non-dipolar components of considerable amplitude, which we will not consider here. The pulsar is also not in a perfect vacuum, but is surrounded by a plasma-filled magnetosphere (Goldreich & Julian, 1969). The vacuum torque vanishes when the angular velocity and the magnetic axis are aligned, while the presence of a magnetosphere would require a loss of angular momentum, no matter what the orientation is. Therefore, the vacuum dipole torque should give only an incomplete description at best. We adopt a general spindown torque of the form $`𝐍_{sd}=N_o[\widehat{𝛀}a(\widehat{𝛀}\widehat{𝐛})\widehat{𝐛}]`$, where we introduce an additional parameter $`a`$. Loss of angular momentum mandates that $`a1`$. It should also be positive, or we would have an angle dependent torque that is opposite in sign to the dipole contribution. The vacuum case is retrieved by setting $`a=1`$. The case of $`a=0`$ corresponds to an external torque with no angular dependence, which would not produce periodic time of arrival residuals.
The torque-modified Euler’s equation can be solved approximately for small $`ϵ`$, and a second contribution to the period residuals arises due to the torque, $`\mathrm{\Delta }P_{sd}/P_{}\stackrel{~}{\mathrm{\Gamma }}_{sd}\mathrm{\Delta }\stackrel{~}{\mathrm{}}`$. We will refer to this as the *spindown* model; and the sum of both geometric and spindown contributions will be referred to as the *full* model. Here, $`\stackrel{~}{\mathrm{\Gamma }}_{sd}=I_3N_o/\varpi _pL^2P_p/t_{sd}10^5`$ and is determined by the spindown properties of the neutron star (in particular, the period derivative, $`\dot{P}`$ or the characteristic age, $`t_{sd}`$); $`\mathrm{\Delta }\stackrel{~}{\mathrm{}}`$ is another complicated function of Jacobian elliptic functions and Legendre integrals. For an axisymmetric star, $`\mathrm{\Delta }\stackrel{~}{\mathrm{}}=a_1\mathrm{sin}\omega _pt+a_2\mathrm{sin}2\omega _pt`$, where $`\omega _p`$ is now the precession frequency, and $`a_i`$ are some coefficients; the axisymmetric model thus has two harmonically related components. We take the two peaks in the spectrum of the period residuals of PSR B1828–11 with periods of $`500`$ and $`250`$ days to be the most significant. From the form of the axisymmetric model, we thus conclude that the precession period must be $`500`$ days. It can also be shown that for an axisymmetric star, in the region of interest, the geometric contribution is quite negligible compared to the spindown contribution for $`a1`$ (see Appendix C).
If the 1000–day period represented the precession period, we would expect to see variations of the pulsar beam width at the same period. While such changes were reported by Stairs, Lyne & Shemar (2000), subsequent more detailed analysis has not shown a 1000–day period in the beam width data (Parry et al. 2005). We thus assume that the precession period is $`500`$ days, and attribute the 1000–day component in the timing data to something unrelated to precession, such as timing noise. We note that our model cannot provide satisfactory fits to the data if the precession period is $`1000`$ days.
### 2.3 Constraints and Statistical Analysis
The orientation of the angular momentum, $`\widehat{𝐋}`$ is fixed in the inertial frame. This is still true even in the presence of the spindown torque, if $`ϵ`$ is sufficiently small (see Appendix C). Then the requirement that the pulse beam, which we assume to be centered along the magnetic axis, $`\widehat{𝐛}`$ never precesses entirely out of our line of sight means that the angle between $`\widehat{𝐋}`$ and $`\widehat{𝐛}`$ should not vary by more than the angular width of the pulse itself. We will refer to the angle between $`\widehat{𝐋}`$ and $`\widehat{𝐛}`$ as the *beam swing angle* and denote it by $`\vartheta `$. If the angular radius of the pulse is $`\rho `$, then the above constraint can be expressed as $`\mathrm{\Delta }\vartheta =\vartheta _{max}\vartheta _{min}2\rho `$; in general, we will require the beam swing variation to be less than some value $`\mathrm{\Delta }\vartheta _{max}`$. We also will define the *wobble angle*, $`\theta `$ as the angle between the angular momentum and the body $`z`$ axis (see Fig. 1 and Appendix B).
The duty cycle allows us to estimate the angular extent of the pulse, and for PSR B1828–11 this varies between $`5^{}`$ to $`7^{}`$. For a circular pulse, this implies that the beam swing angle cannot be varying by more than a few degrees. Larger variations would require a more elongated pulse shape. Yet, at this time, there is not enough evidence to elaborate more on this. In particular, polarization data might be quite useful to determine the extent of the pulse. Stairs et al. (2000) also report periodic variations in the average pulse shape. We offer a possible explanation in a following section.
Another restriction may be that PSR B1828–11 does not have an interpulse. That can further restrict the relative orientation of the angular momentum and the magnetic axis. However, due to uncertainties in the structure of the magnetic field, it is not clear that an interpulse will necessarily appear. Therefore, we do not impose this restriction. The observer’s location is an additional parameter, and can be independently fixed.
We apply the two models - geometric $`(a=0)`$ and full $`(a0)`$ \- under the given constraints to PSR B1828–11, using a Bayesian approach to obtain probability distribution functions (pdfs) for individual parameters. We assume specific priors in the full multi-dimensional parameter space (to be discussed next), but the effective priors exhibited in the projected (marginalized) 1-D posterior pdfs shown in the figures below are integrals over the constraints. Because there appears to be systematic noise in the period residuals larger than their tabulated uncertainties, we scale the latter and then marginalize over the scaling factor as detailed in Appendix D. Once the likelihood is determined, the individual pdfs are obtained through integration over the remaining parameters and normalization.
Let $`\{p_k\}`$ denote the set of the $`n`$ parameters. Then the likelihood, $`(\{p_k\})`$ and the volume of integration, $`𝒱(\{p_k\})`$ are functions of this set. The latter also depends on the beam swing angle constraint, which itself is a function of a subset of the parameters. The priors, $`g_i(p_i)`$ are functions only of the single parameter they refer to. Then the projected 1-D posterior pdf for the $`i`$-th parameter can be expressed as an integral of the likelihood over the remaining parameters, over the volume defined by the constraints,
$`f_i(\{p_k\})={\displaystyle _{𝒱(\{p_k\})}}g_i(p_i)(\{p_k\}){\displaystyle \underset{ji}{\overset{n}{}}}g_j(p_j)dp_j.`$ (1)
Similarly, the projected 1-D prior can be expressed as,
$`h_i(\{p_k\})={\displaystyle _{𝒱(\{p_k\})}}g_i(p_i){\displaystyle \underset{ji}{\overset{n}{}}}g_j(p_j)dp_j.`$ (2)
It is these two quantities ($`f_i`$ and $`h_i`$) that are plotted in Figs. 4-9. Note that, if the volume of integration had not depended on the constraints, then we would simply have $`h_i=g_i`$.
Within the context of our precession model we can use the pdf to compare how well different sets of model parameters fit the data. However, we cannot assess the extent to which the data demand explanation in terms of precession, as opposed to some other, completely different physical model. Any model for the data will lead to a pdf with local maxima at certain values of the parameters of the model, and we can assess the relative significance of these peaks to quantify the extent to which the model parameters are determined by the data. Whether or not another model that is just as well-motivated physically as our precession model can fit the data better is outside the scope of our analysis. Given a competitor model - of which we are unaware - Bayesian methods could be used for making model comparisons.
## 3 Results and Discussion
The physical parameters that determine the form of the residuals are the two angles that specify the orientation of the magnetic axis in the body frame (the polar angle, $`\chi `$ and the azimuthal angle, $`\phi `$; see Fig. 14); any two of $`e^2`$, $`\lambda `$ and $`k^2`$; and $`a`$. There is also a $`\tau _o`$ (measured in precession cycles) that determines the initial phase. Thus, the total number of parameters is six. $`\chi `$ varies between 0 and $`\pi /2`$, and its prior is taken to be flat over $`\mathrm{cos}\chi `$; and $`\varphi `$ varies between 0 and $`2\pi `$, and has a flat prior. In other words, we assume that orientations of the magnetic axis are equally likely over all solid angles. Priors for $`\tau _o`$ and $`a`$ are flat between 0 and 1. On the other hand, $`e^2`$ and $`\lambda `$ can have any positive values, as long as the beam swing angle is constrained and $`k^2<1`$; therefore, we have to introduce cutoffs in their priors. $`\lambda `$ is related to the wobble angle, and due to the beam swing angle constraint it cannot be too large; we take $`\lambda 0.2`$ with a flat prior.
The situation is slightly more complicated for $`e^2`$. The crust of a neutron star (which in our model is the only component since we do not consider the liquid interior) can relax only through shearing motions as it spins down and so must be triaxial (Link, Franco & Epstein 1998; Franco, Link & Epstein 2000). Adding the magnetic stresses, which result from the multi-polar field near the surface would produce a very complicated figure. It is, therefore, quite unlikely that the star is oblate axisymmetric ($`e^2=0`$) or prolate axisymmetric ($`e^2\mathrm{}`$) to a very high precision. On theoretical grounds one might expect $`e^2`$ to be close to unity. Thus, we first take $`e^22`$ with a flat prior (Figs. 4-6). However, we find that within this range the pdf for $`e^2`$ is not confined. Because of that, we also consider a second case where we allow for larger values of $`e^2`$ (Figs. 7-9). The spindown model we use is derived under the assumption that $`e^2`$ is not exceedingly large (see Appendix C). Therefore, we take $`e^24000`$, which seems to encompass the regions of interest, without violating our assumptions. Since the volume of integration is considerably larger at large values of $`e^2`$ than at small values, taking a flat pdf over $`e^2`$ in this case would greatly suppress the importance of small $`e^2`$. Therefore, we need to incorporate a prior that is fair for both regimes: we take a flat prior over $`\mathrm{ln}(1+e^2)`$, i.e. the prior for $`e^2`$ is $`1/(1+e^2)`$. For both $`e^22`$ and $`e^24000`$, we calculate pdfs for three different values of the maximum beam swing angle constraint: $`\mathrm{\Delta }\vartheta _{max}=1^{}`$, $`3^{}`$ and $`5^{}`$.
In Fig. 3 we show the data that we use, together with some sample models. The parameters for these models are listed in Table 1. The best fit that we find is a purely geometrical model, which has a very large beam swing angle that is, in fact, outside our prior range (which was relaxed for determining an unconstrained, global “favorite” model). The axisymmetric model given here is similar to that of Link & Epstein (2001), except that the beam swing angle is constrained to be below $`5^{}`$; in fact, as we show in Appendix C, any axisymmetric model is assured to yield quite similar results even when we introduce the additional torque parameter $`a`$.
Because of the large number of parameters, numerical integration for the pdfs is quite time consuming. The figures presented here typically have a resolution of about 100 points per parameter or less. This means that fine structure in the pdfs may have been missed. Nevertheless, the most notable structures in the pdfs are expected to remain.
In Figs. 4-9 we show the projected 1-D pdfs for each parameter, computed by integrating the multidimensional pdf over all other parameters. The prior is also a function of the entire set of parameters and is not separable for all except $`a`$. We define the 1-D prior for a given parameter by integrating over the rest. A comparison with the full 1-D pdf illustrates the importance of the period residuals in determining the pdf. Keep in mind that both the prior and the posterior pdfs also include the beam swing angle constraint. In these figures, the dotted lines are for the prior pdfs, the dashed lines are for the geometric model alone (Eq. (35)), and the solid lines are for the geometric model and the spindown model (Eq. (96)) both combined.
We now discuss some of the main characteristics and implications of our analysis.
*The torque parameter $`a`$:* For $`e^22`$ Figs. 4-6 show considerable probability over the whole range of acceptable values, with a peak at low $`a`$ that becomes more prominent as $`\mathrm{\Delta }\vartheta `$ increases. Another lower and wider peak appears for $`\mathrm{\Delta }\vartheta _{max}=5^{}`$ at larger values of $`a`$ (Fig. 6). Nevertheless, neither of the peaks is highly significant, because they do not contain most of the probability. Therefore, we conclude that the data do not constrain $`a`$ strongly, and it can be quite different from the vacuum spindown value $`a1`$. As larger values of $`e^2`$ are permitted, Figs. 7-9 show a peak at $`a1`$, but with a large tail extending over most of the parameter space. With increasing values of $`\mathrm{\Delta }\vartheta `$, lower $`a`$ values become likelier, but most of the probability ($`>90\%`$) still lies at $`a0.25`$. Thus, the value of $`a`$ is not well-determined, but there is evidence for an angle-dependent torque. The parameter $`a`$ is truly a measure of the angle dependence of the spin-down torque; it does not depend on the geometric effect at all.
*The magnetic inclination $`\chi `$:* The axisymmetric model discussed by Link & Epstein (2001) requires $`\chi `$ to be extremely close to $`90^{}`$. As discussed in Appendix C, this is true even when we allow $`a1`$. For a triaxial model, we find the range of acceptable $`\chi `$ values to be much larger. The geometric model has a peak at small $`\chi `$, which moves on to higher values of $`\chi `$ and broadens with increasing beam swing angle. This trend is still seen at $`e^22`$ when $`a0`$ is turned on. The spindown produces a narrow but strong peak in the vicinity of $`90^{}`$, which becomes more pronounced as $`\mathrm{\Delta }\vartheta `$ is allowed to be bigger. This peak corresponds to the axisymmetric case, and implies that it requires larger $`\mathrm{\Delta }\vartheta `$ values; in fact, the model discussed by Link & Epstein has $`\mathrm{\Delta }\vartheta 6.4^{}`$. For $`e^24000`$ (Figs. 7-9), the peak at large $`\chi `$ remains apparent, though now it is quite broad. The inclusion of points beyond $`e^22`$ seems to favor a more important spindown contribution, and the geometric effect is further suppressed. There is also a small cusp that appears in the pdf for $`\mathrm{\Delta }\vartheta <5^{}`$, at $`\chi `$ very near $`90^{}`$, corresponding to the axisymmetric case. Yet, this peak is quite narrow, and the vast majority of the probability lies outside of it.
*The triaxiality parameter $`e^2`$:* For $`e^22`$, the pdf looks quite similar to the prior, implying that the data do not differentiate among values of $`e^2`$ (Figs. 4-6). There is a narrow sharp peak at $`e^2=0`$ for the full model, corresponding to the *oblate* axisymmetric case, but the probability enclosed within the peak is very small. Note that, at the same time, the pdf for the geometric effect alone almost vanishes, i.e. for the axisymmetric case, spindown is essential. When we allow $`e^2>2`$, another peak appears in the pdf (Figs. 7-9). At these values of $`e^2`$ the star is once again almost axisymmetric, except that now $`I_1<I_2I_3`$, i.e. the star is *prolate*. Qualitatively $`e^2=1`$ separates oblate and prolate shapes. Then, comparing the probabilities enclosed in the two regions, $`e^2<1`$ and $`e^2>1`$ , we find that the prolate case contains, by far, most of the probability, even though we have adopted a prior which somewhat disfavors large $`e^2`$ values. Note that as the beam swing angle constraint is relaxed, the probability becomes quite evenly distributed over a wide region in $`e^2`$. This leads to the conclusion that there are many triaxial models with a wide range of $`e^2`$ that can fit the data. In other words, the data do not discriminate among values of $`e^2`$, especially when $`\mathrm{\Delta }\vartheta `$ is relatively large.
*The wobble parameter $`\lambda `$:* For both $`e^22`$ and $`e^24000`$, the pdf is contained in a region which seems to be mostly confined by the beam swing angle. The constraint results in a dramatic cutoff at the high end of $`\lambda `$. Beyond that point, we cannot find a value of $`e^2`$ for which $`\mathrm{\Delta }\vartheta `$ will be smaller than the constraint. Within that region, there is considerable probability distributed over the whole range of $`\lambda `$; both models follow the same trend. We can conclude that the oblate case favors somewhat larger values of $`\lambda `$, which produces a peak that is partially visible for $`\mathrm{\Delta }\vartheta _{max}=5^{}`$ (Fig. 6). The prolate case, on the other hand, favors smaller values of $`\lambda `$, resulting in a second peak, which appears when we allow $`e^2`$ to be large (Fig. 9), but is absent when $`e^2`$ is confined to low values (Fig. 6).
Fig. 10 shows multivariate pdfs for the full model, plotted as surfaces at constant $`e^2`$, as functions of $`a`$ and $`\chi `$. The remaining parameters ($`\varphi `$, $`\lambda `$ and $`\tau _o`$) are integrated out, and the pdfs are normalized over all surfaces. (The special case $`e^2=0`$ can be done with considerably higher resolution, because of the freedom of choice of $`\varphi `$.) The ripples that are present are artifacts of the integration and subside as the resolution is improved. The amplitude of the ripples serves as an implicit way of evaluating the significance of the peaks in the pdf. Relatively large amplitude implies that there are no significant peaks in the pdf, meaning that no points or regions in the allowed parameter space are favored strongly.
As $`\mathrm{\Delta }\vartheta _{max}`$ is increased, two regions acquire prominence: a very narrow ridge at large $`\chi `$, which extends over a wide region in $`a`$ and includes the axisymmetric case, and a smaller peak at small $`\chi `$ and small $`a`$. Keep in mind that this second peak is further discriminated against by the prior, which is flat over $`\mathrm{cos}\chi `$, i.e. there is a factor of $`\mathrm{sin}\chi `$ that also enters the pdf. Consequently, when the one-dimensional pdfs are calculated, the second peak is considerably suppressed.
## 4 Shape of the Pulse
The pulse profile of PSR B1828–11 alternates between two different modes, one narrow and the other broad, and has Fourier power at both 250 and 500 days (Stairs, Lyne & Shemar 2000, Stairs et al. 2003). Stairs et al. (2003) describe how the pulse profile is determined. During a particular observing session, 16 pulse averages may be either broad or narrow, with a shape parameter $`S`$ defined to be the fraction of the mean pulse shape for that session attributed to the narrow component. The shape parameter $`S`$ varies systematically between $`0`$ (all wide) and $`1`$ (all narrow) over the precession cycle, with a strong Fourier component at the “first harmonic” $`1/250\mathrm{days}`$ of the “fundamental” precession frequency $`1/500\mathrm{days}`$ (Stairs, Lyne & Shemar 2000, Stairs et al. 2003). Link & Epstein (2001) suggested that the emission beam of the pulsar must have an hourglass shape in order for $`S`$ to exhibit substantial variability on the 250 day timescale. (A similar elongated shape was inferred from studies of the geodetic precession of PSR B1913+16 by Weisberg & Taylor 2002.) However, they did not address the issue of mode switching during individual observing sessions at all. Here we present an alternative viewpoint centered around modeling profile mode switching within individual observing sessions, and argue that it may be able to produce some aspects of the required harmonic structure shown in Fig. 2.
The basic geometrical picture is shown in Fig. 11. We attribute the narrow component of the pulse to core emission centered around the beam axis. The broader profile is a superposition of core and conal emission. Thus, in the parlance of Rankin (1990, 1993), the pulsar alternates between presenting a core single ($`\mathrm{S}_\mathrm{t}`$) and triple (T) pulse profile. The relatively young spindown age of PSR B1828–11 (about 0.11 Myr) and its large value of $`B_{12}/P^2`$ are consistent with this categorization. However, the apparent pulse width in the narrow state appears to be anomalous: Rankin (1990) finds a FWHM pulse width $`W_{\mathrm{core}}=2.45^{}P^{1/2}/\mathrm{sin}\chi `$, where $`\chi `$ is the angle between a pulsar’s spin and magnetic axes. For PSR B1828–11, the FWHM of the narrow state is about $`2.3^{}`$, as opposed to $`W_{\mathrm{core}}3.85^{}/\mathrm{sin}\chi `$ from Rankin’s (1990) formula. We note that the bounding relationship $`W_{\mathrm{core}}2.45^{}/P^{1/2}`$ was derived from a set of “interpulsars” thought to be nearly orthogonal rotators, so we would have expected Rankin’s (1990) formula to work especially well if $`\chi `$ is near $`90^{}`$. as was suggested by Link & Epstein (2001); on the other hand, the discrepancy is also smallest for $`\chi 90^{}`$, which may be circumstantial evidence that PSR B1828–11 really is nearly an orthogonal rotator. (Moreover, the frequency dependence of the core width is relatively weak at high frequencies, so that the fact that the Rankin formula is for 1 GHz emission, whereas the Stairs et al. observations were at 1.6 or 1.7 GHz is not responsible for the discrepancy.) We note that there are other exceptions to Rankin’s (1990) bound, but not many (see e.g. Fig. 23b in Graham-Smith , adapted from Gould ); given uncertainties in $`\chi `$ there may be other pulsars with $`W_{\mathrm{core}}>2.45^{}/\sqrt{P}`$ but also $`W_{\mathrm{core}}<2.45^{}/\sqrt{P}\mathrm{sin}\chi `$.
If pulsar core emission intensity were Gaussian, we would expect an observed intensity of the form,
$$I_{\mathrm{core}}=I_{\mathrm{core}}(0)\mathrm{exp}\left[\frac{\beta ^2}{2\rho _1^2}\frac{(\varphi \mathrm{sin}\alpha )^2}{2\rho _2^2}\right],$$
(3)
where $`\beta `$ is the impact parameter of the observer’s line of sight relative to the beam axis, $`\varphi `$ is pulse phase (centered on epoch of closest passage relative to the axis), and $`\alpha =\chi +\beta `$ is the angle between the line of sight to the observer and the stellar spin axis. Here $`\rho _1`$ and $`\rho _2`$ define the extent and the shape of the beam, which would be elliptical when they are not equal. This formula assumes that $`\beta \chi `$, and that emission is strongly beamed along magnetic field lines, but does not presume that the emission pattern is circularly symmetric with respect to the beam axis. The two directions 1 and 2 are relative to a coordinate system in which $`\widehat{𝐞}_3=\widehat{𝐛}`$ coincides with the magnetic moment of the star, which is assumed to be the beam axis, $`\widehat{𝐞}_2=\widehat{𝐋}\mathbf{\times }\widehat{𝐛}/|\widehat{𝐋}\mathbf{\times }\widehat{𝐛}|`$, and $`\widehat{𝐞}_1=(\widehat{𝐛}\mathbf{}\widehat{𝐋}\widehat{𝐛}\widehat{𝐋})/|\widehat{𝐋}\mathbf{\times }\widehat{𝐛}|`$. For a Gaussian beam, Eq. (3) shows that the core component width is independent of the impact angle $`\beta `$, although the peak intensity is not (e.g. Rankin 1990). However, this is a unique property of a Gaussian profile. We can well imagine that the emission cuts off sharply (even discontinuously) for sufficiently large $`\beta `$, in which case the core width could be narrower than normal. Because the peak core intensity would also be lower in such cases, it would be harder to detect, which may account for the rarity of exceptions to Rankin’s (1990) bound.
A sharper cutoff to the core emission beam would not only allow narrower core pulse profiles, but would also introduce $`\beta `$ dependence into the width. As a simple example, suppose that the beam profile is,
$`I_{\mathrm{core}}`$ $`=`$ $`I_{\mathrm{core}}(0)\mathrm{exp}\left({\displaystyle \frac{u}{2}}{\displaystyle \frac{\kappa u^2}{4}}\right)`$
$`u`$ $`=`$ $`{\displaystyle \frac{\beta ^2}{\rho _1^2}}+{\displaystyle \frac{(\varphi \mathrm{sin}\alpha )^2}{\rho _2^2}},`$ (4)
i.e. still a self-similar function but with a sharper cutoff than a Gaussian profile. The peak intensity is at $`\varphi =0`$, where $`u=u_{\mathrm{min}}=\beta ^2/\rho _1^2`$; the FWHM is at phases $`\pm \varphi _{1/2}`$, where
$`{\displaystyle \frac{\varphi _{1/2}\mathrm{sin}\alpha }{\rho _2}}`$ $`=`$ $`\left[\sqrt{\left(u_{\mathrm{min}}+{\displaystyle \frac{1}{\kappa }}\right)^2+{\displaystyle \frac{4\mathrm{ln}2}{\kappa }}}\left(u_{\mathrm{min}}+{\displaystyle \frac{1}{\kappa }}\right)\right]^{1/2}`$ (5)
$``$ $`\sqrt{{\displaystyle \frac{2\mathrm{ln}2}{1+\kappa u_{\mathrm{min}}}}}=\sqrt{{\displaystyle \frac{2\mathrm{ln}2}{1+\kappa \beta ^2/\rho _1^2}}}`$
where the approximation is valid for small values of $`\kappa `$, irrespective of $`\kappa u_{\mathrm{min}}`$. The cutoff becomes important once $`\kappa u_{\mathrm{min}}1`$ i.e. for $`\beta \rho _1/\sqrt{\kappa }`$. The core width decreases with increasing $`\beta `$, as does the peak intensity observed from the core.
As we mentioned above, we ascribe the broader pulse profile state to a superposition of core and conal components. In keeping with the schematic Fig. 11, we assume that the conal emitting pattern is patchy and, as we discuss further below, probably only stationary in the mean. Consider an individual conal emitting region (hereafter “blob”) $`i`$; we assume that it is centered at $`(x_i^1,x_i^2)=\rho _i(\mathrm{cos}\sigma _i,\mathrm{sin}\sigma _i)`$. The emission pattern of blob $`i`$ may be anisotropic in a complicated fashion, with possible preferred directions not only along the $`\widehat{𝐞}_{1,2}`$ axes, but also along and perpendicular to $`\widehat{𝐱}_i=(\mathrm{cos}\sigma _i,\mathrm{sin}\sigma _i)`$. The observer sees an intensity that is a function of the two variables $`\beta \rho _i\mathrm{cos}\sigma _i`$ and $`(\varphi +\varphi _d)\mathrm{sin}\alpha \rho _i\mathrm{sin}\sigma _i`$, where the phase delay is $`\varphi _d=h_i/cP_{}1^{}(h_i/340\mathrm{km})`$, where $`h_i`$ is the height of blob $`i`$ above the core emitting region. The peak value of the intensity of radiation seen from any given blob is only a function of $`|\beta \rho _i\mathrm{cos}\sigma _i|`$, though, and we assume that blob $`i`$ is detectable provided that
$$|\beta \rho _i\mathrm{cos}\sigma _i|\delta _i,$$
(6)
where $`\delta _i`$ may depend on $`\sigma _i`$. Thus, the detectability of an individual blob varies through the precession cycle, and the probability of seeing any blobs at all also varies, thus affecting the observed beam width.
The problem of modeling the detectability of a given blob can be quite complex. To illustrate, suppose that each blob is at $`\rho _i=\rho `$, and has the same anisotropic shape. Simplify even more by assuming that the emission profiles of the blobs have a characteristic length $`\delta _r`$ along the (radial) direction from the beam axis to its center, and a different length $`\delta _t`$ in the direction tangential to it. Then we can detect the blob if,
$`(\beta \rho \mathrm{cos}\sigma _i)^2\delta _t^2\mathrm{sin}^2\sigma _i+\delta _r^2\mathrm{cos}^2\sigma _i`$ (7)
For convenience, we define $`\delta ^2=\frac{1}{2}(\delta _t^2+\delta _r^2)`$ and $`q\delta ^2=\frac{1}{2}(\delta _t^2\delta _r^2)`$. Note that $`q`$ may be positive or negative, and $`|q|1`$; $`q=0`$ for a circular beam. Since we expect that in general $`\delta \rho `$, and except for special values, $`\delta \beta `$, we only expect blobs within a small range $`\mathrm{\Delta }\sigma (\beta )`$ to be visible.
Fig. 12 shows the result of solving Eq. (7) for the range of observable $`\mathrm{\Delta }\sigma /2\pi `$ as a function of impact parameter $`\beta `$. (The solutions were not extended beyond $`\beta /\rho =1+\delta \sqrt{1q}`$, where $`\mathrm{\Delta }\sigma 0`$.) The results exhibit complicated behavior even in this simple model. Given $`N`$ conal blobs, the probability of seeing the broader pulse profile is large when $`2\times N\mathrm{\Delta }\sigma /2\pi 1`$, and is small when $`N\mathrm{\Delta }\sigma /\pi 1`$. (The factor of two is because the observer’s line of sight crosses the cone twice.) Thus, we may expect the shape parameter $`S`$ to be small for impact parameters where $`\mathrm{\Delta }\sigma /2\pi `$ is large, and vice-versa; during a precession cycle, both regimes may be sampled.
Fig. 13 shows an example of how the probability of seeing only the narrow pulse would vary with precession phase in this model; this example captures the main features of the observed beam width variations shown in Fig. 2. For constructing the figure we adopted $`q=0.8`$, $`\delta =0.1\rho `$, and assumed a sinusoidal variation of the impact parameter with precession phase:
$$\beta (\varphi _p)=\beta _0+\beta _1\mathrm{sin}\varphi _p$$
(8)
with $`\beta _0=0.95\rho `$ and $`\beta _1=0.15\rho `$ assumed for graphical purposes. The “shape parameter” is taken to be $`S=(1\mathrm{\Delta }\sigma /\pi )^N`$ i.e. the probability that no blobs are detected; Fig. 13 assumes $`N=6`$. Clearly, $`S`$ varies periodically but not sinusoidally in this model, and also varies substantially in half a precession cycle. This distinctive “doubly periodic” variation is only seen if the observer’s line of sight crosses near $`\beta =\rho `$. This is consistent with our earlier discussion of core widths if the core emission is still visible but starting to cut off at such impact angles. Presumably, the peak intensity of the core emission must also far exceed that of any conal blob for this model to be viable; there are some indications that conal emission becomes more prominent as pulsars age (e.g. Rankin 1990). Although the range of variation of $`S`$ in this example is smaller than in PSR B1828–11, extensions of the model, such as different assumptions about the conal emission (e.g. an hourglass-shaped cone as in Link & Epstein , or a more complicated version of blob anisotropy) may possibly yield a better account of the data.
If the conal blobs were stationary in the rotating frame, then the observer would see pulse profile variations as a function of precession phase, but would not see any variations at a given precession phase. However, it is likely that the conal blobs are not at fixed positions but rather circulate around the cone in a rotating carousel (Deshpande & Rankin 1999, 2001; see Rankin & Wright 2003 for a review). In this picture, which has empirical support (Deshpande & Rankin 1999, 2001), conal emission is from beams that circulate with a frequency $`\mathrm{\Omega }_d=f_d\mathrm{\Omega }_{}`$ relative to a reference frame rotaing with the star; the circulation is probably the consequence of $`𝑬`$$`\mathbf{\times }`$$`𝑩`$ drift (e.g. Ruderman & Sutherland 1975, Gil, Melikidze & Geppert 2003, Wright 2003), and $`f_d0.1`$ is a reasonable value. By contrast, the core emission is stationary, and from a much lower altitude than the conal emission (possibly from near the polar cap). In this picture, during a given observing session the core component is always visible, but the conal component fluctuates as emitting blobs pop in and out of the observer’s line of sight periodically. The probability that the observer sees conal emission at all varies systematically during the precession cycle, and is fixed during any observing session lasting a day or so (i.e. far less than the precession period).
If this model is correct at least in a broad outline, then the total beam swing during a precession cycle is $`13^{}`$, given expected core radii (Rankin 1993). Larger beam swings could be accommodated by a more complex model for the pulse shape (eg. the hourglass shape of Link & Epstein 2001).
## 5 Conclusion
We find a wide range of triaxial models that may explain the period residuals of PSR B1828–11 in terms of precession, even under the stringent constraints we have imposed. We find many fits that are as good as, or better than the axisymmetric model considered before (Link & Epstein 2001). In general, fits improve with larger beam swing angle variations ($`\mathrm{\Delta }\vartheta `$), but if we assume that the pulse is confined to a region a few degrees in size, we have to rule them out. Prolate and oblate axisymmetric models seem to be favored, especially for small $`\mathrm{\Delta }\vartheta `$, but that is not sufficient to rule out other triaxial models. Both the geometric and spindown effects contribute to the fits. Oddly, if we relax our beam swing constraint completely, the data prefer a best fit that has $`a=0`$ (no spindown contribution), but $`\mathrm{\Delta }\vartheta `$ for that model is unreasonably large (Fig. 3 and Table 1), so it is merely an unphysical curiosity.
In the oblate axisymmetric model, spindown is the dominant effect, but it requires parameters (in particular, a large $`\chi `$ value) that could be expected to produce an interpulse, which is not seen in PSR B1828–11. If we were to impose the absence of an interpulse as a constraint, some of the models we have permitted in our analysis would be excluded, particularly those at large $`\chi `$. Conceivably, the magnetic field and core beam structure of PSR B1828–11 are sufficiently complex that an interpulse would be absent even at $`\chi 90^{}`$. We note that our model for shape variations suggests that we are only viewing the outskirts of the core emission in the component we detect, which may enhance the probability that emission from the opposite pole is undetectable. Thus, we do not impose the absence of an interpulse as a constraint on our analysis.
Our models do not require $`a=1`$, so substantial deviations from the vacuum spindown formula are allowed. In fact, rather small values of $`a`$ are permitted for an oblate star ($`e^2<1`$). However, for a prolate star ($`e^2>1`$) we find that larger values are favored ($`a>0.25`$), thus providing evidence for an angle-dependent torque. To our knowledge, our analysis provides the first evidence that pulsars are spun down by a torque that depends on the angle between the magnetic moment and the instantaneous angular velocity.
The magnetic obliquity, $`\chi `$ is no longer required to be extremely close to $`90^{}`$, and we find $`\lambda `$ (which is related to the wobble angle) to be restricted mainly by the beam swing angle. Two peaks in $`e^2`$ are prominent, corresponding to the oblate axisymmetric ($`e^2=0`$) and prolate nearly axisymmetric (large $`e^2`$) cases. These are especially evident for small beam swing angle variations. For larger beam swing variations, the data do not discriminate among values of $`e^2`$ very much.
In summary, our precession model fits the data equally well for a broad range of parameters. We cannot constrain the shape of the star, but we do find evidence for angle-dependent spin-down torque. Overall, the ability of our physically-motivated model to account for the principal features of PSR B1828–11’s timing data without special choices of the parameters reinforces the idea that the pulsar is precessing. In particular, we have shown that the data can be fit without resorting to a nearly orthogonal rotator with a vacuum-like dipole torque as Link & Epstein (2001) did in their preliminary work. Though we cannot strongly constrain the angular dependence of the spin-down torque with the present data, the potential remains for learning more about this important aspect of the neutron star magnetosphere from future observations. The parameters of our model are not very tightly constrained. Two possible reasons are that we have only three cycles of data, and that the data have a large degree of intrinsic scatter that our simple model cannot account for, creating wide pdfs. It will be interesting to see if the parameters can be more tightly constrained as more data become available over the next decade.
Our analysis does not employ the data on the shape parameter variations, because constructing a comprehensive mathematical description would require a reliable model for the pulsar beam. Although we do not possess such an accurate model, we offer an explanation using a compound pulse structure, with core and cone components.
Precession has interesting implications for pulsar observations, which so far have not been widely discussed. One immediate, and very obvious effect would be *disappearing pulsars*, i.e. pulsars that due to precession, would at some point leave the line of sight of the observer, but excluding other effects, would eventually come back. The timescale of such changes could be months to years.
## Acknowledgments
We thank Ingrid Stairs for providing us with the data and for valuable discussion. This research is supported in part by NSF AST-0307273 (Cornell University), NSF AST-0098728 (Montana State University) and IGPP 1222R from LANL.
## Appendix A Period Residuals for a Triaxial Rigid Star in the Absence of External Torques: The Geometric Contribution
Euler’s equation for a freely precessing rigid body is,
$`{\displaystyle \frac{d𝐋}{dt}}+𝛀\times 𝐋=0`$ (9)
and can be solved analytically in terms of Jacobian elliptic functions (see Landau & Lifshitz, 1976). Using the principal axes ($`I_1I_2I_3`$) as the basis for the body (rotating) frame, we can express the components of the angular momentum unit vector $`\widehat{𝐋}`$ as,
$`L_1=\mathrm{\Lambda }_1\text{cn}(\tau ,k^2)\text{ , }\mathrm{\Lambda }_1=\sqrt{{\displaystyle \frac{I_1(2EI_3L^2)}{L^2(I_3I_1)}}}`$
$`L_2=\mathrm{\Lambda }_2\text{sn}(\tau ,k^2)\text{ , }\mathrm{\Lambda }_2=\sqrt{{\displaystyle \frac{I_2(2EI_3L^2)}{L^2(I_3I_2)}}}=\mathrm{\Lambda }_1\sqrt{1+e^2}`$ (10)
$`L_3=\mathrm{\Lambda }_3\text{dn}(\tau ,k^2)\text{ , }\mathrm{\Lambda }_3=\sqrt{{\displaystyle \frac{I_3(L^22EI_1)}{L^2(I_3I_1)}}}=\sqrt{1\mathrm{\Lambda }_1^2}`$
Here, the argument of the elliptic functions is,
$`\tau =t\omega _p\text{where}\omega _p=\sqrt{{\displaystyle \frac{(I_3I_2)(L^22EI_1)}{I_1I_2I_3}}}={\displaystyle \frac{ϵL\mathrm{\Lambda }_3}{I_3\sqrt{1+e^2}}}`$ (11)
and, the parameter of the elliptic functions is,
$`k^2={\displaystyle \frac{(I_2I_1)(2EI_3L^2)}{(I_3I_2)(L^22EI_1)}}={\displaystyle \frac{e^2\mathrm{\Lambda }_1^2}{\mathrm{\Lambda }_3^2}}=e^2\lambda ^2`$ (12)
where, we make use of the following auxiliary definitions,
$`ϵ=(I_3I_1)/I_1\text{ and }e^2={\displaystyle \frac{I_3(I_2I_1)}{I_1(I_3I_2)}}`$ (13)
The minus signs that we have explicitly included in our definitions are due to our choice of the initial orientation of axes.
Note that $`\omega _p`$ is not the precession frequency, since the elliptic functions do not have a period of $`2\pi `$, or more precisely, $`\omega _p`$ is not the time derivative of the angular displacement. Instead, the precession frequency is given through,
$`\mathrm{\Omega }_p={\displaystyle \frac{2\pi }{P_p}}={\displaystyle \frac{\omega _p\pi }{\stackrel{~}{\pi }}}`$ (14)
where 2$`\stackrel{~}{\pi }`$ is the period of the elliptic functions and can be calculated using the Legendre elliptic integral of the first kind, $`\tau =F(\varphi ,k)`$ where $`\mathrm{sin}\varphi =\text{sn}\tau `$,
$`\stackrel{~}{\pi }/2=F(\pi /2,k)`$ (15)
The values of the parameter $`k^2`$ are unrestricted, though different regimes require careful handling. $`k^2<1`$ corresponds to precession around the body $`z`$ axis; $`k^2=1`$ corresponds to the unstable trajectories of the angular momentum, which decay exponentially towards the intermediate axis, $`y`$; and $`k^2>1`$ is precession around the $`x`$ axis (see Binet ellipsoid). For now we will confine ourselves to the first case, and the other two will be left for a later section.
We are interested in an isolated neutron star, and we want to determine the time of arrival (TOA) of pulses (produced along the magnetic axis) for an inertial observer. Let the inertial $`z`$ axis be along the angular momentum vector, which remains constant; and let the inertial $`x`$ axis be defined by the orientation of the observer, whom we choose to locate in the first quadrant of the inertial $`xz`$ plane. Let a unit vector $`\widehat{𝐛}`$ denote the orientation of the magnetic axis, and $`b_i`$ be the rotating frame components. Then, whenever the inertial $`y`$ component (which we choose to denote by $`b_y`$) vanishes, while the inertial $`x`$ component ($`b_x`$) is positive, we get a pulse. The two frames are related through a rotation matrix constructed from the Euler angles $`\theta ,\psi `$ and $`\varphi `$ (see Goldstein 1980), whence the two conditions can be expressed as,
$`b_y=b_1(\mathrm{cos}\psi \mathrm{sin}\varphi +\mathrm{cos}\theta \mathrm{cos}\varphi \mathrm{sin}\psi )+b_2(\mathrm{sin}\psi \mathrm{sin}\varphi +\mathrm{cos}\theta \mathrm{cos}\varphi \mathrm{cos}\psi )b_3\mathrm{sin}\theta \mathrm{cos}\varphi =0`$ (16)
and,
$`b_x=b_1(\mathrm{cos}\psi \mathrm{cos}\varphi \mathrm{cos}\theta \mathrm{sin}\varphi \mathrm{sin}\psi )b_2(\mathrm{sin}\psi \mathrm{cos}\varphi +\mathrm{cos}\theta \mathrm{sin}\varphi \mathrm{cos}\psi )+b_3\mathrm{sin}\theta \mathrm{sin}\varphi >0`$ (17)
Using the solution for the angular momentum (Eq. (10)), the Euler angles are given through,
$`\mathrm{cos}\theta =\mathrm{\Lambda }_3\text{dn}\tau \text{ , }\mathrm{sin}\theta =\mathrm{\Lambda }_1\sqrt{1+e^2\text{sn}^2\tau }\text{ , }\mathrm{cos}\psi ={\displaystyle \frac{\sqrt{1+e^2}\text{sn}\tau }{\sqrt{1+e^2\text{sn}^2\tau }}}\text{ , }\mathrm{sin}\psi ={\displaystyle \frac{\text{cn}\tau }{\sqrt{1+e^2\text{sn}^2\tau }}}\text{ and}`$ (18)
$`{\displaystyle \frac{d\varphi }{dt}}={\displaystyle \frac{L}{I_3}}\left(1+{\displaystyle \frac{ϵ}{1+e^2\text{sn}^2\tau }}\right)`$
The last equation can be written as,
$`\varphi (t)=\varphi _o+{\displaystyle \frac{L}{I_3}}t+{\displaystyle \frac{\sqrt{1+e^2}}{\mathrm{\Lambda }_3}}{\displaystyle _0^\tau }{\displaystyle \frac{d\tau }{1+e^2\text{sn}^2\tau }}`$ (19)
Eq. (16) also implies,
$`\mathrm{tan}\varphi ={\displaystyle \frac{N}{D}}`$
$`N=b_3\mathrm{\Lambda }_1\left(1+e^2\text{sn}^2\tau \right)+b_2\mathrm{\Lambda }_3\text{sn}\tau \text{dn}\tau \sqrt{1+e^2}+b_1\mathrm{\Lambda }_3\text{cn}\tau \text{dn}\tau `$ (20)
$`D=b_2\text{cn}\tau b_1\text{sn}\tau \sqrt{1+e^2}`$
Pulses are seen when Eqs. (19) and (20) are both satisfied. In the absence of precession ($`\mathrm{\Lambda }_1=0`$ and $`k^2=0`$) the solution of Eq. (20) for $`\varphi `$ is simply given through,
$`\varphi =2\pi n+\eta +\mathrm{tan}^1(\sqrt{1+e^2}\mathrm{tan}\tau )`$ (21)
where $`\eta =\pi /2\phi `$ (Fig. 14); and we can further restrict it to lie anywhere between 0 and $`2\pi `$. Note that we have implicitly included the second requirement, Eq. (17), by skipping every other possible solution for $`\varphi `$. This is a non-trivial assumption, and would break down if the magnetic axis $`𝐛`$ happens to lie between $`𝛀`$ and $`𝐋`$. However, for a pulsar these two vectors are very nearly aligned since $`ϵ`$ is extremely small. Therefore, we do not need to worry about such a case.
We express the general solution for $`\varphi `$ as,
$`\varphi =2\pi n+\eta +\zeta `$ (22)
where $`\zeta `$ is confined to lie within a period of tangent (i.e. $`\pi `$). Using $`\mathrm{tan}\eta =b_1/b_2`$ one can show that,
$`\mathrm{tan}\zeta ={\displaystyle \frac{Nb_2Db_1}{Nb_1+Db_2}}`$ (23)
We now have two equations for $`\varphi `$ (Eqs. (19) and (22)) which we can combine to get the times of arrival of pulses,
$`{\displaystyle \frac{L}{I_3}}t_n=2\pi n+\zeta _n\zeta _o{\displaystyle \frac{\sqrt{1+e^2}}{\mathrm{\Lambda }_3}}{\displaystyle _0^{\tau _n}}{\displaystyle \frac{d\tau }{1+e^2\text{sn}^2\tau }}`$ (24)
$`\zeta _o`$ (for $`\tau =0`$) appears as a consequence of the fact that $`\varphi _o=\eta +\zeta _o`$. We have,
$`\mathrm{tan}\zeta _o={\displaystyle \frac{b_2b_3\mathrm{\Lambda }_1b_1b_2(1\mathrm{\Lambda }_3)}{b_1b_3\mathrm{\Lambda }_1+b_2^2+b_1^2\mathrm{\Lambda }_3}}`$ (25)
Note that in the absence of precession (i.e. when $`\mathrm{\Lambda }_1=0`$) the times of arrival reduce to the form,
$`t_n={\displaystyle \frac{2\pi I_3n}{L}}`$ (26)
which is the solution for pure rotation.
The period (between two pulses) is given as,
$`P_n=t_nt_{n1}`$ (27)
whence,
$`{\displaystyle \frac{L}{I_3}}P_n2\pi ={\displaystyle \frac{L}{I_3}}\mathrm{\Delta }P_n=\zeta _n\zeta _{n1}{\displaystyle \frac{\sqrt{1+e^2}}{\mathrm{\Lambda }_3}}{\displaystyle _{\tau _{n1}}^{\tau _n}}{\displaystyle \frac{d\tau }{1+e^2\text{sn}^2\tau }}`$ (28)
If the precession period is much longer than the pulse period (as is the case for a neutron star) we can approximate the differences by derivatives,
$`{\displaystyle \frac{L}{I_3}}\mathrm{\Delta }P_n={\displaystyle \frac{d\zeta _n}{dn}}{\displaystyle \frac{\sqrt{1+e^2}/\mathrm{\Lambda }_3}{1+e^2\text{sn}^2\tau _n}}{\displaystyle \frac{d\tau _n}{dn}}=\left[{\displaystyle \frac{d\zeta _n}{d\tau _n}}{\displaystyle \frac{\sqrt{1+e^2}/\mathrm{\Lambda }_3}{1+e^2\text{sn}^2\tau _n}}\right]{\displaystyle \frac{d\tau _n}{dn}}`$ (29)
We will find it convenient to define the expression inside the parentheses as a new function,
$`f_n={\displaystyle \frac{d\zeta _n}{d\tau _n}}{\displaystyle \frac{\sqrt{1+e^2}/\mathrm{\Lambda }_3}{1+e^2\text{sn}^2\tau _n}}`$ (30)
The derivative $`d\zeta /d\tau `$ is given through (from Eq. (23)),
$`{\displaystyle \frac{d\zeta }{d\tau }}={\displaystyle \frac{DN^{}ND^{}}{N^2+D^2}}`$ (31)
where, from Eq. (20),
$`{\displaystyle \frac{dN}{d\tau }}=2b_3\mathrm{\Lambda }_1e^2\text{sn}\tau \text{cn}\tau \text{dn}\tau +b_2\mathrm{\Lambda }_3\text{cn}\tau (\text{dn}^2\tau k^2\text{sn}^2\tau )\sqrt{1+e^2}b_1\mathrm{\Lambda }_3\text{sn}\tau (\text{dn}^2\tau +k^2\text{cn}^2\tau )`$
$`{\displaystyle \frac{dD}{d\tau }}=b_2\text{sn}\tau \text{dn}\tau b_1\text{cn}\tau \text{dn}\tau \sqrt{1+e^2}`$
To evaluate $`d\tau _n/dn`$ we will make use of the time of arrival equation, Eq. (24). Taking the derivative of both sides with respect to $`n`$, we get,
$`{\displaystyle \frac{L}{I_3}}{\displaystyle \frac{dt_n}{dn}}={\displaystyle \frac{L}{I_3\omega _p}}{\displaystyle \frac{d\tau _n}{dn}}=2\pi +f_n{\displaystyle \frac{d\tau _n}{dn}}\text{ so that }{\displaystyle \frac{d\tau _n}{dn}}={\displaystyle \frac{2\pi \varpi _p}{1\varpi _pf_n}}`$ (32)
where we have defined a new dimensionless quantity,
$`\varpi _p={\displaystyle \frac{I_3\omega _p}{L}}={\displaystyle \frac{ϵ\mathrm{\Lambda }_3}{\sqrt{1+e^2}}}`$ (33)
The pulse period will be given through,
$`P_n=t_nt_{n1}{\displaystyle \frac{dt_n}{dn}}={\displaystyle \frac{1}{\omega _p}}{\displaystyle \frac{d\tau _n}{dn}}={\displaystyle \frac{P_{}}{1\varpi _pf_n}}`$ (34)
whence the period residuals can be found to be,
$`{\displaystyle \frac{\mathrm{\Delta }P_n}{P_{}}}{\displaystyle \frac{P_n}{P_{}}}1={\displaystyle \frac{\varpi _pf_n}{1\varpi _pf_n}}\varpi _pf_n`$ (35)
where $`P_{}=2\pi I_3/L`$ is the rotation period of the star, and the last approximation results from our anticipation that $`ϵ`$ will be sufficiently small. Indeed, from the above definitions we get, for small $`k^2`$,
$`{\displaystyle \frac{ϵ}{\sqrt{1+e^2}}}{\displaystyle \frac{P_{}}{P_p}}3.2\times 10^8{\displaystyle \frac{P_{}\text{ (sec)}}{P_p\text{ (yrs)}}}`$ (36)
The coefficient of the function $`f_n`$ in Eq. (35) is,
$`B=P_{}\varpi _p\text{where}\varpi _p={\displaystyle \frac{I_3\omega _p}{L}}={\displaystyle \frac{2\stackrel{~}{\pi }I_3}{LP_p}}={\displaystyle \frac{\stackrel{~}{\pi }P_{}}{\pi P_p}}`$ (37)
in other words,
$`B={\displaystyle \frac{\stackrel{~}{\pi }P_{}}{\pi P_p}}^2={\displaystyle \frac{\stackrel{~}{\pi }B_o}{\pi }}`$ (38)
For PSR B1828–11 the rotation period is 405.04 ms, and the precession period is about 511 days, so that $`B_o3.8`$ ns.
### A.1 The Axisymmetric Body
For an axisymmetric body we have $`e^2=k^2=0`$. We can also set $`b_2=0`$ which is equivalent to introducing some initial phase shift in the definition of $`\tau `$. We thus get, after some rearrangement,
$`f_n={\displaystyle \frac{d\zeta }{d\tau }}{\displaystyle \frac{1}{\mathrm{\Lambda }_3}}={\displaystyle \frac{\mathrm{\Lambda }_1}{\mathrm{\Lambda }_3}}{\displaystyle \frac{b_1b_3\mathrm{\Lambda }_3\mathrm{cos}\tau +b_1^2\mathrm{\Lambda }_1\mathrm{sin}^2\tau +b_3^2\mathrm{\Lambda }_1}{(b_3\mathrm{\Lambda }_1+b_1\mathrm{\Lambda }_3\mathrm{cos}\tau )^2+b_1^2\mathrm{sin}^2\tau }}`$ (39)
Let’s now assume that the angle $`\theta `$ between the symmetry axis and the angular momentum is small, i.e. $`\mathrm{\Lambda }_1\theta `$ and $`1\mathrm{\Lambda }_3\theta ^2/2`$, and working to second order compute the period residuals,
$`f_n\theta \left({\displaystyle \frac{b_3}{b_1}}\right)\mathrm{cos}\tau {\displaystyle \frac{\theta }{2}}^2+\theta ^2\left[{\displaystyle \frac{1}{2}}+\left({\displaystyle \frac{b_3}{b_1}}\right)^2\right]\mathrm{cos}2\tau `$ (40)
Also let $`b_3=\mathrm{cos}\chi `$. Then,
$`\mathrm{\Delta }P_n/P_{}\varpi _pf_n\theta ϵ\mathrm{cot}\chi \mathrm{cos}\tau {\displaystyle \frac{\theta ^2ϵ}{2}}+\theta ^2ϵ\left[{\displaystyle \frac{1}{2}}+\mathrm{cot}^2\chi \right]\mathrm{cos}2\tau `$ (41)
(Here $`P_{}=2\pi I_3/L`$ and $`\varpi _p=I_3\omega _p/L=ϵ\mathrm{\Lambda }_3`$.) Note that a harmonic arises from geometrical effects.
### A.2 Precession Around the Principal Axis Corresponding to the Smallest Moment of Inertia ($`k^2>1`$)
We will now look at this case in more detail. The solution given through Eq. (10) is still valid. However, it will be mathematically and computationally convenient to carry out a transformation of the Jacobian elliptic functions with $`k^2>1`$ into functions with parameter $`1/k^2<1`$. (That, in the limit $`\mathrm{\Lambda }_1=1`$ and $`\mathrm{\Lambda }_3=0`$, we get $`\omega _p=0`$ whence $`\tau =0`$ and $`k^2=\mathrm{}`$, provides further motivation.) Then, we can write the general solution as,
$`L_1=\mathrm{\Lambda }_1\text{dn}\widehat{\tau }`$
$`L_2=\mathrm{\Lambda }_3\sqrt{1+\widehat{e}^2}\text{sn}\widehat{\tau }`$ (42)
$`L_3=\mathrm{\Lambda }_3\text{cn}\widehat{\tau }`$
where, $`\widehat{\tau }=\tau k`$, $`\widehat{e}^2=1/e^2`$ and the parameter of the elliptic functions is now $`\widehat{k}^2=1/k^2`$. We can also define a new frequency from Eq. (11),
$`\widehat{\omega }_p=\omega _pk=\sqrt{{\displaystyle \frac{(I_2I_1)(2EI_3L^2)}{I_1I_2I_3}}}`$ (43)
Through an appropriate redefinition of axes, the solution can be expressed in a form identical to the $`k^2<1`$ case, except that now $`\widehat{I}_1>\widehat{I}_2>\widehat{I}_3`$. Define a new right-handed coordinate basis for the rotating frame,
$`\widehat{𝐞}_1=𝐞_3`$
$`\widehat{𝐞}_2=𝐞_2`$ (44)
$`\widehat{𝐞}_3=𝐞_1`$
Let $`\widehat{\mathrm{\Lambda }}_1=\mathrm{\Lambda }_3`$ and $`\widehat{\mathrm{\Lambda }}_3=\mathrm{\Lambda }_1`$, then the components of the angular momentum can be expressed as,
$`\widehat{L}_1=L_3=\widehat{\mathrm{\Lambda }}_1\text{cn}\widehat{\tau }`$
$`\widehat{L}_2=L_2=\widehat{\mathrm{\Lambda }}_1\sqrt{1+\widehat{e}^2}\text{sn}\widehat{\tau }`$ (45)
$`\widehat{L}_3=L_1=\widehat{\mathrm{\Lambda }}_3\text{dn}\widehat{\tau }`$
The precession is now clockwise, as can also be verified from Euler’s equation. $`\widehat{\mathrm{\Lambda }}_i`$ have exactly the same form as before, in terms of the new moments of inertia,
$`\widehat{\mathrm{\Lambda }}_1=\sqrt{{\displaystyle \frac{\widehat{I}_1(L^22E\widehat{I}_3)}{L^2(\widehat{I}_1\widehat{I}_3)}}}`$
$`\widehat{\mathrm{\Lambda }}_2=\sqrt{{\displaystyle \frac{\widehat{I}_2(L^22E\widehat{I}_3)}{L^2(\widehat{I}_2\widehat{I}_3)}}}`$ (46)
$`\widehat{\mathrm{\Lambda }}_3=\sqrt{{\displaystyle \frac{\widehat{I}_3(2E\widehat{I}_1L^2)}{L^2(\widehat{I}_1\widehat{I}_3)}}}`$
So do $`\widehat{\omega }_p`$, $`\widehat{e}^2`$ and $`\widehat{k}^2`$, as can be verified from the equations above. We have thus transformed the problem from a “$`k^2>1`$ case for an $`I_3>I_1`$ body” into a “$`\widehat{k}^2<1`$ case for an $`\widehat{I}_1>\widehat{I}_3`$ body”, which should not be surprising.
The equations for the Euler angles (Eqs. (18)) remain of the same form, with the exception of $`\mathrm{cos}\widehat{\psi }`$. This is effectively a sign change, $`\tau \widehat{\tau }`$, in the argument,
$`\widehat{\zeta }(\widehat{\tau })=\zeta (\widehat{\tau })\text{ whence }{\displaystyle \frac{d\widehat{\zeta }}{d\widehat{\tau }}}={\displaystyle \frac{d\zeta (\widehat{\tau })}{d\widehat{\tau }}}`$ (47)
One must be careful with Eq. (19) as well, where there is also a sign change due to the fact that now $`\widehat{I}_1>\widehat{I}_3`$,
$`{\displaystyle \frac{L(\widehat{I}_3\widehat{I}_1)}{\widehat{\omega }_p\widehat{I}_1\widehat{I}_3}}={\displaystyle \frac{\sqrt{1+\widehat{e}^2}}{\widehat{\mathrm{\Lambda }}_3}}={\displaystyle \frac{\widehat{\mathrm{\Lambda }}_2}{\widehat{\mathrm{\Lambda }}_1\widehat{\mathrm{\Lambda }}_3}}`$ (48)
These two effects add up to modify the function $`f_n`$ defined through Eq. (30),
$`\widehat{f}_n(\widehat{\tau })=f_n(\widehat{\tau })`$ (49)
## Appendix B The Wobble and Beam Swing Angles
We will define the wobble ($`\theta `$) and beam swing ($`\vartheta `$) angles according to (see Fig. 1),
$`\mathrm{cos}\theta =\widehat{𝐋}\widehat{𝐳}=\mathrm{\Lambda }_3\text{dn}\tau `$ (50)
$`\mathrm{cos}\vartheta =\widehat{𝐋}\widehat{𝐛}=b_1\mathrm{\Lambda }_1\text{cn}\tau b_2\mathrm{\Lambda }_2\text{sn}\tau +b_3\mathrm{\Lambda }_3\text{dn}\tau `$ (51)
By definition, the wobble angle is equivalent to Euler’s angle $`\theta `$ and is constant for an axisymmetric star. The beam swing angle is related to the angle between the beam and the observer. It could exceed $`90^{}`$, but since the pulse will have a limited angular size, there is a restriction on how much it can vary throughout a precession period. Otherwise, the beam will leave the observer’s line of sight. Therefore, the span of the beam swing angle serves as a constraint. The angle can be further restricted by imposing the conditions for an interpulse.
We define the widest span as $`\mathrm{\Delta }\vartheta =\vartheta _{max}\vartheta _{min}`$. Then the constraint is that this be smaller than some value $`\mathrm{\Delta }\vartheta _{max}`$, which is estimated based on information about the pulse width and shape. Note that the beam angle depends on four parameters: the two angles determining the orientation of the pulse, and any two of $`k^2`$, $`e^2`$ and $`\lambda =\mathrm{\Lambda }_1/\mathrm{\Lambda }_3`$.
## Appendix C Period Residuals for the Spindown Torque
When external torques are present Euler’s equation becomes,
$`{\displaystyle \frac{d𝐋}{dt}}+𝛀\times 𝐋=𝐍`$ (52)
Taking the dot product with the angular momentum, we get the equation governing the evolution of its magnitude,
$`{\displaystyle \frac{dL}{dt}}=\widehat{𝐋}𝐍`$ (53)
If we now substitute $`𝐋=L\widehat{𝐋}`$ in Euler’s equation, we get, after some rearrangement,
$`L{\displaystyle \frac{d\widehat{𝐋}}{dt}}+L^2(𝐈^1\widehat{𝐋})\times \widehat{𝐋}=𝐍(\widehat{𝐋}𝐍)\widehat{𝐋}`$ (54)
which governs the evolution of the orientation of the angular momentum. These two are the basic equations that need to be solved. Of course, only three (of the total of four components) are independent equations.
In the classical rotating magnetic-dipole model of pulsars, the angular momentum is lost to radiation. The electromagnetic torque for a spherical rigid star in vacuum is (Davis and Goldstein, 1970),
$`𝐍_{vac}={\displaystyle \frac{2\mu ^2\mathrm{\Omega }^3}{3c^3}}\widehat{𝐛}\times (\widehat{𝛀}\times \widehat{𝐛})={\displaystyle \frac{2\mu ^2\mathrm{\Omega }^3}{3c^3}}\left[\widehat{𝛀}(\widehat{𝛀}\widehat{𝐛})\widehat{𝐛}\right]`$ (55)
Note that the torque vanishes when $`\widehat{𝛀}`$ and $`\widehat{𝐛}`$ are aligned. However, the pulsar is not in a perfect vacuum, and is surrounded by a magnetosphere. Therefore, there should be loss of angular momentum even when these two vectors are aligned. We will therefore adopt a general spindown torque of the form,
$`𝐍_{sd}=N_o\left[\widehat{𝛀}a(\widehat{𝛀}\widehat{𝐛})\widehat{𝐛}\right]`$ (56)
where $`a`$ is a dimensionless parameter that measures the relative importance of the two components. The amplitude of the spindown torque can be estimated from observed values of the spindown time, and is small. We will be interested in a particular example (PSR B1828–11) where the spindown time is,
$`t_{sd}{\displaystyle \frac{L}{N_o}}10^5\text{ yrs}`$ (57)
Compare this with the observed precession period for the same pulsar,
$`P_p={\displaystyle \frac{2\pi }{\mathrm{\Omega }_p}}{\displaystyle \frac{I_3}{ϵL}}1\text{ yr}`$ (58)
The ratio of the two gives,
$`{\displaystyle \frac{t_{sd}}{P_p}}{\displaystyle \frac{ϵL^2}{I_3N_o}}10^5`$ (59)
The second term in Eq. (54) has a magnitude of $`ϵL^2/I`$, therefore the RHS of that equation is quite negligible for the case of interest (as will be discussed below).
The above form of the torque is true for spherical stars. This is nevertheless a good approximation, given how small $`N_o`$ and $`ϵ`$ are. In fact, we will neglect all combinations of $`N_o`$ with $`ϵ`$. This is equivalent to taking $`\widehat{𝛀}\widehat{𝐋}`$ within all torque terms, since the angle between these two vectors is of the order of $`ϵ`$. We therefore have,
$`𝐍=N_o\left[\widehat{𝐋}a(\widehat{𝐋}\widehat{𝐛})\widehat{𝐛}\right]`$ (60)
$`{\displaystyle \frac{dL}{dt}}=\widehat{𝐋}𝐍=N_o\left[1a(\widehat{𝐋}\widehat{𝐛})^2\right]`$ (61)
$`L{\displaystyle \frac{d\widehat{𝐋}}{dt}}+L^2(𝐈^1\widehat{𝐋})\times \widehat{𝐋}=𝐍(\widehat{𝐋}𝐍)\widehat{𝐋}=N_oa(\widehat{𝐋}\widehat{𝐛})\left[\widehat{𝐛}(\widehat{𝐋}\widehat{𝐛})\widehat{𝐋}\right]`$ (62)
Loss of energy (or angular momentum, given through the second equation) now clearly requires that $`a1`$. Finally, we will also neglect any time dependence within $`N_o`$ itself.
The third equation demands careful thought. In component form,
$`L{\displaystyle \frac{d}{dt}}\left(\begin{array}{c}L_1\\ L_2\\ L_3\end{array}\right)+{\displaystyle \frac{ϵL^2}{I_1}}\left(\begin{array}{c}L_2L_3(1s)\\ L_1L_3\\ L_1L_2s\end{array}\right)=N_oa\mathrm{cos}\vartheta \left(\begin{array}{c}b_1L_1\mathrm{cos}\vartheta \\ b_2L_2\mathrm{cos}\vartheta \\ b_3L_3\mathrm{cos}\vartheta \end{array}\right)`$ (72)
where $`L_i`$ are the components of $`\widehat{𝐋}`$, $`s=(I_2I_1)/(I_3I_1)`$, $`\mathrm{cos}\vartheta =\widehat{𝐋}\widehat{𝐛}`$, and we have already neglected second order terms in $`ϵ`$. As long as the star is sufficiently non-spherical and the angular momentum is sufficiently misaligned with the body $`z`$ axis, we can neglect the RHS, as it causes changes in the orientation of the angular momentum smaller (by many orders of magnitude) than the second term. In other words, $`ϵ`$ and $`\mathrm{\Lambda }_1`$ are small but not zero. (Keep in mind that there is no precession if either one is zero.) Also, $`s`$ cannot be too close to unity, i.e. $`e^2`$ cannot be exceptionally large.
The same cannot be done in the RHS of the equation for the magnitude of the angular momentum, as it is the only term we have. Incidentally, setting $`N_o=0`$ would take us back to the torque-free precession case.
In order to write the equations in a dimensionless form, let’s divide all sides by a frequency $`\omega _p`$, defined in accordance with Eq. (11),
$`\omega _p(t)={\displaystyle \frac{ϵL(t)\mathrm{\Lambda }_3}{I_3\sqrt{1+e^2}}}`$ (73)
but where the magnitude of the angular momentum is no longer constant. Also define,
$`d\tau =\omega _p(t)dt`$ (74)
which, for a constant $`\omega _p`$, reduces to the familiar form of the torque-free case. Now, the differential equations become, after some rearrangement,
$`{\displaystyle \frac{dL}{d\tau }}=\widehat{𝐋}𝐍/\omega _p=(N_o/\omega _p)\left[1a(\widehat{𝐋}\widehat{𝐛})^2\right]`$ (75)
$`{\displaystyle \frac{d\widehat{𝐋}}{d\tau }}+(L/\omega _p)(𝐈^1\widehat{𝐋})\times \widehat{𝐋}=0`$ (76)
Since $`L/\omega _p`$ is time-independent, the second equation has exactly the same solution as before, except that $`\tau `$ is now different, and given through a differential equation on its own. In other words, $`L_i`$’s remain of the same form. Thus, we only need to solve Eqs. (74) and (75).
Define a new dimensionless constant,
$`\varpi _p=I_3\omega _p/L=ϵ\mathrm{\Lambda }_3/\sqrt{1+e^2}ϵ`$ (77)
and let’s write,
$`L=L_o[1\mathrm{}(\tau )]`$
$`𝐍=N_o𝐧(\tau )`$ (78)
$`t=[\tau +\delta (\tau )]/\omega _{po}`$
where $`\omega _{po}=\varpi _pL_o/I_3`$. The differential equations now become,
$`{\displaystyle \frac{d\mathrm{}}{d\tau }}=\left({\displaystyle \frac{I_3N_o}{\varpi _pL_o^2}}\right){\displaystyle \frac{\widehat{𝐋}𝐧}{1\mathrm{}}}\text{ and }{\displaystyle \frac{d\delta }{d\tau }}={\displaystyle \frac{\mathrm{}}{1\mathrm{}}}`$ (79)
It’s worth noting that we make no assumptions in these substitutions.
If we finally define one more dimensionless constant,
$`\stackrel{~}{\mathrm{\Gamma }}_{sd}={\displaystyle \frac{I_3N_o}{\varpi _pL_o^2}}`$ (80)
and let $`\mathrm{}=\stackrel{~}{\mathrm{\Gamma }}_{sd}\stackrel{~}{\mathrm{}}`$ and $`\delta =\stackrel{~}{\mathrm{\Gamma }}_{sd}\stackrel{~}{\delta }`$, then the two equations can be written as,
$`{\displaystyle \frac{d\stackrel{~}{\mathrm{}}}{d\tau }}={\displaystyle \frac{\widehat{𝐋}𝐧}{1\stackrel{~}{\mathrm{\Gamma }}_{sd}\stackrel{~}{\mathrm{}}}}\text{ and }{\displaystyle \frac{d\stackrel{~}{\delta }}{d\tau }}={\displaystyle \frac{\stackrel{~}{\mathrm{}}}{1\stackrel{~}{\mathrm{\Gamma }}_{sd}\stackrel{~}{\mathrm{}}}}`$ (81)
For the pulsar that we discuss here, we have,
$`\stackrel{~}{\mathrm{\Gamma }}_{sd}{\displaystyle \frac{I_3N_o}{ϵL_o^2}}{\displaystyle \frac{P_p}{t_{sd}}}10^5`$ (82)
This means that one may safely ignore the denominators of the two equations, thus further simplifying the results,
$`{\displaystyle \frac{d\stackrel{~}{\mathrm{}}}{d\tau }}=\widehat{𝐋}𝐧\text{ and }{\displaystyle \frac{d\stackrel{~}{\delta }}{d\tau }}=\stackrel{~}{\mathrm{}}`$ (83)
Note that $`\widehat{𝐋}𝐧=1a(\widehat{𝐋}\widehat{𝐛})^20`$, thus assuring that $`L`$ is monotonically decreasing (as required by loss of angular momentum).
### C.1 Time of Arrival Residuals
Since $`L_i`$ remain of the same form, with the only difference being that $`\tau `$ is now determined through a differential equation (Eq. (74)), the Euler angles remain the same (Eqs. (18)). However, due to the time dependence of $`L`$, it is more convenient to express $`\varphi `$ as a function of $`\tau `$, and we need to replace Eq. (19) by,
$`\varphi (\tau )=\varphi _o+{\displaystyle _0^\tau }{\displaystyle \frac{d\tau }{\varpi _p}}\left(1+{\displaystyle \frac{ϵ}{1+e^2\text{sn}^2\tau }}\right)=\varphi _o+{\displaystyle \frac{\tau }{\varpi _p}}+{\displaystyle \frac{ϵ}{\varpi _p}}{\displaystyle _0^\tau }{\displaystyle \frac{d\tau }{1+e^2\text{sn}^2\tau }}`$ (84)
Thus, all we have to do is to replace $`Lt_n/I_3`$ by $`\tau _n/\varpi _p`$ on the LHS of Eq. (24). The period (given through Eq. (34)) remains the same as well, as does the calculation of $`d\tau _n/dn`$. In fact, we run into trouble only with the period residuals, since the magnitude of the angular momentum is now changing. Define,
$`\mathrm{\Delta }P_{res}=P_n{\displaystyle \frac{2\pi I_3}{L_o}}\text{ and }\mathrm{\Delta }P_n=P_n{\displaystyle \frac{2\pi I_3}{L}}`$ (85)
$`\mathrm{\Delta }P_n`$ is formally the same as the torque-free case. However, observations give us only information about $`\mathrm{\Delta }P_{res}`$. In practice, one first determines the period $`(P_{})`$ at some epoch $`(t_o)`$, and then finds the period derivative ($`\dot{P}_{}`$) which is the secular term attributed to spindown, and subtracts both contributions, so that the residuals are then given through,
$`\mathrm{\Delta }P_{res}=P(t)P_{}\dot{P}_{}(tt_o)`$ (86)
Consider the difference between the two definitions above,
$`\mathrm{\Delta }P_{res}\mathrm{\Delta }P_n=2\pi I_3\left({\displaystyle \frac{1}{L}}{\displaystyle \frac{1}{L_o}}\right)={\displaystyle \frac{2\pi I_3\stackrel{~}{\mathrm{\Gamma }}_{sd}}{L_o}}{\displaystyle \frac{d\stackrel{~}{\delta }}{d\tau }}`$ (87)
where we make use of Eqs. (79) and (81). We thus get,
$`{\displaystyle \frac{\mathrm{\Delta }P_{res}}{P_{}}}={\displaystyle \frac{\varpi _pf_n}{1\varpi _pf_n}}+{\displaystyle \frac{\stackrel{~}{\mathrm{\Gamma }}_{sd}\stackrel{~}{\mathrm{}}}{1\stackrel{~}{\mathrm{\Gamma }}_{sd}\stackrel{~}{\mathrm{}}}}\varpi _pf_n+\stackrel{~}{\mathrm{\Gamma }}_{sd}\stackrel{~}{\mathrm{}}`$ (88)
where $`P_{}=2\pi /\mathrm{\Omega }_{}=2\pi I_3/L_o`$. The first term is the geometric effect ($`\mathrm{\Delta }P_{ge}/P_{}\varpi _pf_n`$) and the second term is the spindown term ($`\mathrm{\Delta }P_{sd}/P_{}\stackrel{~}{\mathrm{\Gamma }}_{sd}\stackrel{~}{\mathrm{}}`$). There are still secular terms present in the spindown term that need to be subtracted. This will be taken care of below. The relative amplitude of these two terms cannot be simply determined, and it is possible that either one is dominant, or that they are comparable.
We now turn our attention to the calculation of $`\stackrel{~}{\mathrm{}}`$. From Eq. (83) we have,
$`\stackrel{~}{\mathrm{}}={\displaystyle _0^\tau }\widehat{𝐋}𝐧𝑑\tau \text{ where, }\widehat{𝐋}𝐧=1a(\widehat{𝐋}\widehat{𝐛})^2`$ (89)
and,
$`\widehat{𝐋}\widehat{𝐛}=b_1\mathrm{\Lambda }_1\text{cn}\tau b_2\mathrm{\Lambda }_2\text{sn}\tau +b_3\mathrm{\Lambda }_3\text{dn}\tau =\mathrm{cos}\vartheta `$ (90)
Carrying out the integrals of the Jacobian elliptic functions, and substituting the values of $`k^2`$ and $`\mathrm{\Lambda }_2`$, we get, after some rearrangement,
$`{\displaystyle _0^\tau }(\widehat{𝐋}\widehat{𝐛})^2d\tau ={\displaystyle \frac{\mathrm{\Lambda }_3^2}{e^2}}[b_2^2(1+e^2)b_1^2k_1^2]\tau +{\displaystyle \frac{\mathrm{\Lambda }_3^2}{e^2}}[b_1^2b_2^2(1+e^2)+b_3^2e^2]E(\text{am }\tau ,k)`$ (91)
$`+{\displaystyle \frac{2b_1b_2\mathrm{\Lambda }_3^2\sqrt{1+e^2}}{e^2}}(1\text{dn}\tau )2\mathrm{\Lambda }_1\mathrm{\Lambda }_3\left[b_2b_3\sqrt{1+e^2}(1\text{cn}\tau )+b_1b_3\text{sn}\tau \right]`$
where we have introduced the complementary parameter $`k_1^2=1k^2`$. $`E(\text{am}\tau ,k)`$ is the Legendre elliptic integral of the second kind, and am$`\tau `$ is the Jacobi amplitude, am$`\tau =\mathrm{sin}^1\text{sn}\tau `$. Note that we have implicitly assumed that at the zero of time, the angular momentum is in the $`xz`$ plane of the body frame. This is a non-trivial assumption, and in general one does not have the freedom of randomly setting the initial orientation of the angular momentum. Therefore, in general we would have $`\widehat{𝐋}=\widehat{𝐋}(\tau \tau _o)`$, where $`\tau _o`$ is the phase offset and is an additional parameter, and the overall result would be to replace $`\stackrel{~}{\mathrm{}}(\tau )`$ above by $`\stackrel{~}{\mathrm{}}(\tau \tau _o)\stackrel{~}{\mathrm{}}(\tau _o)`$, where $`\stackrel{~}{\mathrm{}}(\tau _o)`$ is just a constant. For simplicity, we will continue to assume $`\tau _o=0`$ in the rest of our derivations, but the general case should be kept in mind.
The oscillatory part of $`\stackrel{~}{\mathrm{}}`$, after secular terms have been removed, is given through,
$`\mathrm{\Delta }\stackrel{~}{\mathrm{}}=\stackrel{~}{\mathrm{}}\widehat{𝐋}𝐧\tau `$ (92)
The average is carried over a precession period, $`2\stackrel{~}{\pi }=4F(\pi /2,k)`$,
$`\widehat{𝐋}𝐧=1a(\widehat{𝐋}\widehat{𝐛})^2`$ (93)
where, using Eq. (91) we get,
$`(\widehat{𝐋}\widehat{𝐛})^2={\displaystyle \frac{1}{2\stackrel{~}{\pi }}}{\displaystyle _0^{2\stackrel{~}{\pi }}}(\widehat{𝐋}\widehat{𝐛})^2d\tau ={\displaystyle \frac{\mathrm{\Lambda }_3^2}{e^2}}[b_2^2(1+e^2)b_1^2k_1^2]+{\displaystyle \frac{\mathrm{\Lambda }_3^2}{e^2}}[b_1^2b_2^2(1+e^2)+b_3^2e^2]{\displaystyle \frac{E(\pi /2,k)}{F(\pi /2,k)}}`$ (94)
and we have made use of the relations $`\text{am}(2\stackrel{~}{\pi })=2\pi `$ and $`E(2\pi ,k)=4E(\pi /2,k)`$. We thus get,
$`\mathrm{\Delta }\stackrel{~}{\mathrm{}}=a(\widehat{𝐋}\widehat{𝐛})^2\tau a{\displaystyle _0^\tau }(\widehat{𝐋}\widehat{𝐛})^2𝑑\tau `$ (95)
which can now be used in Eq. (88) to calculate the time of arrival residuals,
$`{\displaystyle \frac{\mathrm{\Delta }P_{sd}}{P_{}}}\stackrel{~}{\mathrm{\Gamma }}_{sd}\mathrm{\Delta }\stackrel{~}{\mathrm{}}`$ (96)
We will find it convenient to express this equation in the following form,
$`\mathrm{\Delta }\stackrel{~}{\mathrm{}}/a=c_1(1\text{cn}\tau )+c_2\text{sn}\tau +{\displaystyle \frac{c_3}{k^2}}(1\text{dn}\tau )+{\displaystyle \frac{c_4}{k^2}}\left[{\displaystyle \frac{E(\pi /2,k)}{F(\pi /2,k)}}\tau E(\text{am }\tau ,k)\right]`$ (97)
where,
$`c_1=2\mathrm{\Lambda }_2\mathrm{\Lambda }_3b_2b_3\text{ , }c_2=2\mathrm{\Lambda }_1\mathrm{\Lambda }_3b_1b_3\text{ , }c_3=2\mathrm{\Lambda }_1\mathrm{\Lambda }_2b_1b_2\text{ and }c_4=\mathrm{\Lambda }_1^2[b_1^2b_2^2(1+e^2)+b_3^2e^2]`$
These coefficients are related to each other through,
$`c_4={\displaystyle \frac{c_2c_3}{2c_1}}+{\displaystyle \frac{c_1c_3}{2c_2}}{\displaystyle \frac{c_1c_2}{2c_3}}k^2`$ (98)
It is also interesting to note that $`\mathrm{\Delta }\stackrel{~}{\mathrm{}}`$ has a non-zero average over a precession period. The residuals may have non-zero average depending on when and how the period and its derivatives are calculated. This becomes particularly important when calculating the time of arrival residuals, which can be obtained by integrating the period residuals, and if the period residuals have a constant term, then the time of arrival residuals will have a linear term. Therefore, in calculating the time of arrival residuals one will have to subtract any constant terms from the period residuals.
### C.2 Amplitude of the Residuals
Consider the period derivative, which is given through the secular terms in $`\stackrel{~}{\mathrm{}}`$,
$`\dot{P}_{}=\eta P_{}\stackrel{~}{\mathrm{\Gamma }}_{sd}\omega _p`$ (99)
where $`\eta =\widehat{𝐋}𝐧=1ac_o`$ and $`c_o=\mathrm{cos}^2\vartheta `$. The amplitude of the period residuals thus becomes, from Eq. (96),
$`A=aP_{}\stackrel{~}{\mathrm{\Gamma }}_{sd}={\displaystyle \frac{a\dot{P}_{}}{\omega _p(1ac_o)}}={\displaystyle \frac{aA_o}{1ac_o}}\text{where}A_o={\displaystyle \frac{P_{}}{2\omega _pt_{sd}}}`$ (100)
and $`t_{sd}=P_{}/2\dot{P}_{}`$ is the spindown time. Recall that $`a`$ measures the strength of the oscillating part of the spindown torque, and must be $`1`$.
For PSR B1828–11 the period is 405.04 ms, the precession period is about 511 days, and the spindown time is 0.11 Myr, so that we get $`A_o409.95`$ ns.
### C.3 The Axisymmetric Body
For an axisymmetric star $`e^2=k^2=0`$, but $`\lambda =k/e0`$. Due to the symmetry we can set $`b_2=0`$ by shifting the zero of time through some phase $`\tau _o`$. (Note that the same cannot be done in a triaxial body, where we chose to fix the axes according to the principal moments of inertia.) In this case Eq. (97) reduces to the form,
$`\mathrm{\Delta }\stackrel{~}{\mathrm{}}/a={\displaystyle \frac{\lambda }{1+\lambda ^2}}\left[\mathrm{sin}2\chi \mathrm{sin}(\tau \tau _o){\displaystyle \frac{\lambda }{4}}\mathrm{sin}^2\chi \mathrm{sin}2(\tau \tau _o)\right]`$ (101)
Note that the non-linearity of the dipole contribution of the torque naturally brings in a harmonic. The period residuals are then,
$`{\displaystyle \frac{\mathrm{\Delta }P_{sd}}{P_o}}=\stackrel{~}{\mathrm{\Gamma }}_{sd}\mathrm{\Delta }\stackrel{~}{\mathrm{}}\text{ where }\stackrel{~}{\mathrm{\Gamma }}_{sd}={\displaystyle \frac{N_o}{\omega _pL_o}}={\displaystyle \frac{1}{2\omega _p\tau _c}}`$ (102)
Here $`\tau _c`$ is the characteristic time,
$`\tau _c={\displaystyle \frac{3c^3I_3}{4\mu ^2\mathrm{\Omega }_o^2}}`$ (103)
It turns out that, for the axisymmetric case, the geometric term is quite negligible compared to the spindown term, for the range of physical parameters of interest (Jones & Andersson, 2001; Link & Epstein, 2001).
To convert our result for period residuals ($`\mathrm{\Delta }P/P_o`$) into residuals of the derivative of the angular velocity ($`\mathrm{\Delta }\dot{\mathrm{\Omega }}/\mathrm{\Omega }_o`$) given by Link & Epstein, we make use of,
$`\mathrm{\Delta }\dot{P}={\displaystyle \frac{d\mathrm{\Delta }P}{dt}}=\omega _p{\displaystyle \frac{d\mathrm{\Delta }P}{d\tau }}\text{ and }{\displaystyle \frac{\mathrm{\Delta }\dot{\mathrm{\Omega }}}{\mathrm{\Omega }_o}}={\displaystyle \frac{\mathrm{\Delta }\dot{P}}{P_o}}`$ (104)
which indeed gives the correct results, together with the initial phase difference of $`\pi `$ between the two definitions,
$`{\displaystyle \frac{\mathrm{\Delta }\dot{\mathrm{\Omega }}}{\mathrm{\Omega }_o}}={\displaystyle \frac{a\lambda }{2\tau _c(1+\lambda ^2)}}\left[\mathrm{sin}2\chi \mathrm{cos}(\tau \tau _o)+{\displaystyle \frac{\lambda }{2}}\mathrm{sin}^2\chi \mathrm{cos}2(\tau \tau _o)\right]`$ (105)
There is one difference between the two derivations and that is the presence of the coefficient $`a`$ which measures the strength of the spindown torque. With the addition of this new element, the number of unknowns increases to three ($`a`$, $`\lambda `$ and $`\chi `$), while a fit to data will yield only two coefficients ($`a_1`$ and $`a_2`$; $`\tau _o`$ does not contain any further information). This implies that there is a certain level of freedom in the choice of the physical coefficients. Let’s denote the fitting function by $`f`$,
$`f=a_1\mathrm{sin}(\tau \tau _o)a_2\mathrm{sin}2(\tau \tau _o)`$ (106)
Then, the relations between the coefficients of this function and the physical parameters that we actually seek would be,
$`a_1={\displaystyle \frac{a\lambda \mathrm{sin}2\chi }{1+\lambda ^2}}\text{ and }a_2={\displaystyle \frac{a\lambda ^2\mathrm{sin}^2\chi }{4(1+\lambda ^2)}}`$ (107)
Define $`\lambda =\mathrm{tan}\theta `$, and the ratio of the two coefficients gives,
$`\mathrm{tan}\chi \mathrm{tan}\theta ={\displaystyle \frac{8a_2}{a_1}}`$ (108)
It is also possible to express $`\chi `$ and $`\theta `$ as functions of $`a`$. However, as it turns out, the range of the physical parameters is severely restricted by the beam swing angle constraint, which in the axisymmetric case is given through,
$`\mathrm{\Delta }\vartheta =2\mathrm{min}(\chi ,\theta )<\mathrm{\Delta }\vartheta _{max}`$ (109)
This forces one of the two angles to be small (which will have $`\mathrm{tan}_1<0.09`$ even if we let $`\mathrm{\Delta }\vartheta _{max}=10^{}`$); while Eq. (108) ensures that the other remains very close to $`90^{}`$. (The ratio of the coefficients is found to be $`a_2/a_10.4`$ for the data used by Link & Epstein (2001). This yields the condition $`\mathrm{tan}_2>36`$, i.e. the second angle has to be larger than $`88^{}`$, in accordance with previous findings.)
## Appendix D Statistical Inference
Denote the set of parameters by $`\stackrel{}{x}`$. Then the pdf for the parameters can be calculated as, by Bayes’s theorem,
$`P(\stackrel{}{x}|D,M)={\displaystyle \frac{P(\stackrel{}{x}|M)P(D|\stackrel{}{x},M)}{P(D|M)}}`$ (110)
where $`D`$ stands for data, $`M`$ stands for the model and also takes into account any other information that is available on the problem apart from data (in this case the beam swing angle, which we impose as a restriction on the parameter space). $`P(\stackrel{}{x}|M)`$ is the prior probability for the parameters; $`P(D|\stackrel{}{x},M)`$ is the likelihood; and $`P(D|M)`$ is effectively a normalization constant.
To find the pdf for a certain parameter, or a subset of parameters, we integrate Eq. (110) over the remaining parameters. For that we need to know the likelihood. For each data point $`y_i`$ at time $`t_i`$, we have a theoretical prediction $`f_i=f(t_i|\stackrel{}{x},M)`$. For well-known uncertainties with Gaussian distribution, we would then have,
$`P(D|\stackrel{}{x},M)={\displaystyle \underset{i}{}}(\sigma _i\sqrt{2\pi })^1\mathrm{exp}\left[{\displaystyle \frac{(y_if_i)^2}{2\sigma _i^2}}\right]`$ (111)
If we assume that the error bars are not well-determined and rescale them through some number $`F`$, the above equation becomes,
$`P(D|\stackrel{}{x},M)={\displaystyle \underset{i}{}}(F\sigma _i\sqrt{2\pi })^1\mathrm{exp}\left[{\displaystyle \frac{(y_if_i)^2}{2F^2\sigma _i^2}}\right]`$ (112)
and we regard $`F`$ as an additional parameter. We have to introduce a prior for $`F`$. Since we do not want it to depend much on the endpoints, we take it to be flat over $`d\mathrm{ln}F`$, i.e. proportional to $`dF/F`$, and integrate over all values of $`F`$. Define,
$`{\displaystyle \underset{i}{}}{\displaystyle \frac{(y_if_i)^2}{2\sigma _i^2}}\chi _o^2(\stackrel{}{x})`$ (113)
whence we get, for $`d`$ data points,
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dFP(D|\stackrel{}{x},M)}{F}}=\left({\displaystyle \underset{i}{}}{\displaystyle \frac{1}{\sigma _i\sqrt{2\pi }}}\right){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{e^{\chi _o^2/F^2}dF}{F^{d+1}}}\left[\chi _o^2(\stackrel{}{x})\right]^{d/2}`$ (114)
where we have dropped anything that does not depend on the remaining parameters $`\stackrel{}{x}`$, including integrals that give constants, products of the original sigmas, and factors of $`\sqrt{2\pi }`$. Thus, our final result is,
$`P(\stackrel{}{x}|D,M)P(\stackrel{}{x}|M)\left[\chi _o^2(\stackrel{}{x})\right]^{d/2}`$ (115)
where the first term is the prior probability, and the constant of proportionality can be computed from the condition that the final pdf is normalized to one.
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# Tracing Galaxy Formation with Stellar Halos I: Methods
## 1. Introduction
There has been a long tradition of searching in the stellar halo of our Galaxy for signatures of its formation. Stars in the halo provide an important avenue for testing theories of galaxy formation because they have long orbital time periods, have likely suffered little from dissipation effects, and tend to inhabit the outer regions of the Galaxy where the potential is relatively smooth and slowly evolving. The currently favored Dark Energy + Cold Dark Matter ($`\mathrm{\Lambda }`$CDM) model of structure formation makes the specific prediction that galaxies like the Milky Way form hierarchically, from a series of accretion events involving lower-mass systems. This leads naturally to the expectation that the stellar halo should be formed primarily from disrupted, accreted systems. In this work, we develop an explicit, cosmologically-motivated model for stellar halo formation using a hybrid N-body plus semi-analytic approach. Set within the context of $`\mathrm{\Lambda }`$CDM, we use this model to test the general consistency of the hierarchical formation scenario for the stellar halo and to provide predictions for upcoming surveys aimed at probing the accretion history of the Milky Way and nearby galaxies.
In a classic study, Eggen, Lynden-Bell, & Sandage (1962) used proper motions and radial velocities of 221 dwarfs to show that those with lower metallicity (i.e. halo stars) tended to move on more highly eccentric orbits. They interpreted this trend as a signature of formation of the lower metallicity stars during a rapid radial collapse. In contrast, Searle & Zinn (1978) suggested that the wide range of metallicites found in a sample of 19 globular clusters at a variety of Galactocentric radii instead indicated that the Galaxy formed from the gradual agglomeration of many sub-galactic sized pieces. A recent analysis of 1203 metal-poor Solar neighborhood stars, selected without kinematic bias (Chiba & Beers, 2000), points to the truth being some combination of these two pictures: this sample contained a small concentration of very low metallicity stars on highly eccentric orbits (reminiscent of Eggen, Lynden-Bell & Sandage’s, 1962 work) but otherwise showed no correlation of increasing orbital eccentricity with decreasing metallicity.
In the last decade, much more direct evidence for the lumpy build-up of the Galaxy has emerged in the form of clumps of stars in phase-space (and, in some cases, metallicity) both relatively nearby (Majewski, Munn, & Hawley, 1996; Helmi, White, de Zeeuw, & Zhao, 1999) and at much larger distances. The most striking example in the latter category is the discovery of the Sagittarius dwarf galaxy (Ibata, Gilmore, & Irwin, 1994; Ibata, Gilmore & Irwin, 1995) — hereafter Sgr — and its associated trails of debris (see Majewski et al., 2003, for an overview of the many detections) which have now been traced entirely around the Galaxy (Ibata et al., 2001; Majewski et al., 2003). Large scale surveys of the stellar halo are now underway (Majewski, Ostheimer, Kunkel, & Patterson, 2000; Morrison et al., 2000; Yanny et al., 2000; Ivezić et al., 2000; Newberg et al., 2002), and have uncovered additional structures, not associated with Sgr (Newberg et al., 2003; Martin et al., 2004; Rocha-Pinto et al., 2004). Moreover, recent advances in instrumentation are now permitting searches for and discoveries of analogous structures around other galaxies in the form of overdensities in integrated light (Shang et al., 1998; Zheng et al., 1999; Forbes et al., 2003; Pohlen et al., 2003) or, in the case of M31, star counts (Ibata et al., 2001; Ferguson et al., 2002; Zucker et al., 2004). Given this plethora of discoveries, there can be little doubt that the accretion of satellites has been an important contributor to the formation of our and other stellar halos. In addition, both theoretical (Abadi et al., 2003; Brook et al., 2005a; Robertson et al., 2005) and observational (Gilmore, Wyse, & Norris, 2002; Yanny et al., 2003; Ibata et al., 2003; Crane et al., 2003; Rocha-Pinto, Majewski, Skrutskie, & Crane, 2003; Frinchaboy et al., 2004; Helmi et al., 2005) work is beginning to suggest that some significant fraction of the Galactic disk could also have been formed this way.
All of the above discoveries are in qualitative agreement with the expectations of hierarchical structure formation (Peebles, 1965; Press & Schechter, 1974; Blumenthal et al., 1984). As the prevailing variant of this picture, $`\mathrm{\Lambda }`$CDM is remarkably successful at reproducing a wide range of observations, especially on large scales (e.g. Eisenstein et al., 2005; Maller et al., 2005; Tegmark et al., 2004; Spergel et al., 2003; Percival et al., 2002). On sub-galactic scales, however, the agreement between theory and observation is not as obvious (e.g. Simon et al., 2005; Kazantzidis et al., 2004; D’Onghia & Burkert, 2004). Indeed, the problems explaining galaxy rotation curve data, dwarf galaxy counts, and galaxy disk sizes have lead some to suggest modifications to the standard paradigm, including an allowance for warm dark matter (e.g. Sommer-Larsen et al., 2004), early-decaying dark matter (Kaplinghat, 2005), or non-standard inflation (Zentner & Bullock, 2002, 2003). These modifications generally suppress fluctuation amplitudes on small scales, driving sub-galactic structure formation towards a more monolithic, non-hierarchical collapse. These issues bring to sharper focus a fundamental question in cosmology today: is structure formation truly hierarchical on small scales? Stellar halo surveys offer powerful data sets for directly answering this question.
Numerical simulations of individual satellites disrupting about parent galaxies can in many cases provide convincing similarities to the observed phase-space lumps. These models allow the observations to be interpreted in terms of the mass and orbit of the progenitor satellite (e.g., Velazquez & White, 1995; Johnston, Spergel, & Hernquist, 1995; Johnston, Sigurdsson, & Hernquist, 1999; Johnston et al., 1999; Helmi & White, 2001; Helmi et al., 2003; Law, Johnston & Majewski, 2004), and even the potential of the galaxy in which it is orbiting (Johnston, Zhao, Spergel, & Hernquist, 1999; Murali & Dubinski, 1999; Ibata et al., 2001; Ibata et al., 2004; Johnston, Law & Majewski, 2005). Nevertheless, a true test of hierarchical galaxy formation will require robust predictions for the frequency and character of the expected phase space structure of the halo.
Going beyond qualitative statements to model the full stellar halo (including substructure) within a cosmological context is non-trivial. The largest contributor of substructure to our own halo is Sgr, estimated to have a currently-bound mass of order $`3\times 10^8M_{}`$ (Law, Johnston & Majewski, 2004). Even the highest resolution cosmological N-body simulations would not resolve such an object with more than a few hundred particles, which would permit only a poor representation of the phase-space structure of its debris (see Helmi, White, & Springel, 2003, for an example of what can currently be done in this field). Such simulations are computationally intensive, so the cost of examining more than a handful of halos is prohibitive and it is difficult to make statements about the variance of properties of halos that might be seen in a large sample of galaxies. Moreover, such simulations in general only follow the dark matter component of each galaxy not the stellar component. In their studies of thick disk and inner halo formation, Brook and collaborators (Brook et al., 2003, 2004a, 2004b, 2005a, 2005b) have modeled the stellar components directly by simulating the evolution of individual galaxies as isolated spheres of dark matter and gas with small-scale density fluctuations superimposed to account for the large-scale cosmology. However, their sample size remains small and, though they are able to make general statements about the properties of their stellar halos, their resolution would prohibit a detailed phase-space analysis.
An alternative is to take an analytic or semi-analytic approach to halo building (e.g. Bullock, Kravtsov, & Weinberg, 2001; Johnston, Sackett, & Bullock, 2001; Taylor, 2004). This allows the production of many halos, and the potential of including prescriptions to follow the stars separately from the dark matter. However, such techniques use only approximate descriptions of the dynamics and are unable to follow the fine details of the phase-space structure accurately.
In this study we develop a hybrid scheme, which draws on the strengths of each of the former techniques to build high resolution, full phase-space models of a statistical sample of stellar halos. Our approach is to vastly decrease the computational cost of a full cosmological simulation by modeling only those accretion events that contribute directly to the stellar halo in detail with N-body simulations, and to represent the rest of the galaxy with smoothly-evolving analytic functions. The baryonic component of each contributing event is followed using semi-analytic prescriptions.
The purpose of this paper is to describe our method, its strengths and limitations (§2), present the results of tests of the consistency of our models with general properties of galaxies and their satellite systems (§3) and outline some implications (§4). We summarize the conclusions in §5. In further work we will go on to compare the full phase-space structure of our halos in detail to observations and to examine the evolution of dark and light matter in satellite galaxies after their accretion.
## 2. Methods
Our methods can be broadly separated into: (I) a simulation phase, which follows the phase-space evolution of the dark matter; and (II) a prescription phase, which embeds a stellar mass with each dark matter particle. Specifically:
Phase I: Simulations
1. We generate merger trees for our parent galaxies using the method outlined in Somerville & Kolatt (1999) based on the Extended-Press-Schechter (EPS) formalism (Lacey & Cole, 1993, — see §2.1)).
2. For each event in step IA, we run a high-resolution N-body simulation that tracks the evolution of the dark matter component of a satellite disrupting within an analytic, time-dependent, parent galaxy + host halo potential (see §2.2).
Phase II: Prescriptions
1. We follow the gas accretion history of each satellite prior to falling into the parent and track its star-formation rate using cosmologically-motivated, semi-analytic prescriptions (see §2.3).
2. We embed the stellar components generated in step IIA within each dark matter satellite by assigning a variable mass-to-light ratio to every particle that is tracked in the (Phase I) N-body simulations (see §2.4).
We consider the two-phase approach a necessary and acceptable simplification since it allows us to separate well-understood and justified approximations in Phase I from prescriptions that can be adjusted and refined during Phase II. In addition, this separation allows us to save computational time and use just one set of dark matter simulations to explore the effect of varying the details of how baryons are assigned to each satellite. A more complete discussion of the strengths and limitations of our scheme is given in §2.5.
### 2.1. Cosmological Framework
Throughout this work we assume a $`\mathrm{\Lambda }`$CDM cosmology with $`\mathrm{\Omega }_\mathrm{m}=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`\mathrm{\Omega }_\mathrm{b}h^2=0.024`$, $`h=0.7`$, and $`\sigma _8=0.9`$. The implied baryon fraction is $`\mathrm{\Omega }_b/\mathrm{\Omega }_\mathrm{m}=0.16`$.
We focus on the formation of stellar halos for “Milky-Way” type galaxies. In all cases our $`z=0`$ host dark matter halos have virial masses $`M_{\mathrm{vir},0}=1.4\times 10^{12}M_{}`$, corresponding virial radii $`R_{\mathrm{vir},0}=282`$ kpc, and virial velocities $`V_{\mathrm{vir}}=144\mathrm{k}\mathrm{m}\mathrm{s}^1`$. The quantities $`M_{\mathrm{vir}}`$ and $`R_{\mathrm{vir}}`$ are related by
$$M_{\mathrm{vir}}=\frac{4\pi }{3}\rho _\mathrm{M}(z)\mathrm{\Delta }_{\mathrm{vir}}(z)R_{\mathrm{vir}}^3,$$
(1)
where $`\rho _\mathrm{M}`$ is the average matter density of the Universe and $`\mathrm{\Delta }_{\mathrm{vir}}`$ is the “virial overdensity”. In the cosmology considered here, $`\mathrm{\Delta }_{\mathrm{vir}}(z=0)337`$, and $`\mathrm{\Delta }_{\mathrm{vir}}178`$ at $`z\stackrel{>}{}1`$ (Bryan & Norman, 1998). The virial velocity is defined as $`V_{\mathrm{vir}}\sqrt{GM_{\mathrm{vir}}/R_{\mathrm{vir}}}`$.
We generate a total of eleven random realizations of stellar halos. General properties of all eleven are summarized in Table 1. Any variations in our results for stellar halos among these are determined by differences in their accretion histories. In all subsequent figures we present results for four stellar halos (1,2,6, and 9) chosen to span the range of properties seen in our full sample.
#### 2.1.1 Semi-analytic accretion histories
We track the mass accretion and satellite acquisition of each parent galaxy by constructing merger trees using the statistical Monte Carlo method of Somerville & Kolatt (1999) based on the EPS formalism (Lacey & Cole, 1993). This method gives us a record of the the masses and accretion times of all satellite halos and hence allows us to follow the mass accretion history of each parent as a function of lookback time. We explicitly note all satellites more massive than $`M_{\mathrm{min}}=5\times 10^6M_{}`$ and treat all smaller accretion events as diffuse mass accretion. Column 2 of Table 1 lists the total number of such events for each simulated halo. For further details see Lacey & Cole (1993); Somerville & Kolatt (1999); Zentner & Bullock (2003). Four examples of the cumulative mass accretion histories of parent galaxies generated in this manner are shown by the (jagged) solid lines in Figure 1.
#### 2.1.2 Satellite orbits
Upon accretion onto the host, each satellite is assigned an initial orbital energy based on the range of binding energies observed in cosmological simulations (Klypin et al., 1999). This is done by placing each satellite on an initial orbit of energy equal to the energy of a circular orbit of radius $`R_{\mathrm{circ}}=\eta R_{\mathrm{vir}}`$, with $`\eta `$ drawn randomly from a uniform distribution on the interval $`[0.4,0.8]`$. Here $`R_{\mathrm{vir}}`$ is the virial radius of the host halo at the time of accretion. We assign each subhalo an initial specific angular momentum $`J=ϵJ_{\mathrm{circ}}`$, where $`J_{\mathrm{circ}}`$ is the specific angular momentum of the aforementioned circular orbit and $`ϵ`$ is the orbital circularity, which takes a value between $`0`$ and $`1`$. We choose $`ϵ`$ from the binned distribution shown in Figure 2 of Zentner & Bullock (2003), which was designed to match the cosmological N-body simulation results of Ghigna et al (1998), and is similar to the circularity distributions found in more recent N-body analyses (Zentner et al., 2004; Benson, 2005). Finally, the plane of the orbit is drawn from a uniform distribution covering the halo sphere.
#### 2.1.3 Dark matter density distributions
We model all satellite and parent halos with the spherically averaged density profile of Navarro, Frenk & White (1996) (NFW):
$$\rho _{{}_{\mathrm{NFW}}{}^{}(r)}=\rho _\mathrm{s}\left(\frac{r}{r_{\mathrm{halo}}}\right)^1\left(1+\frac{r}{r_{\mathrm{halo}}}\right)^2,$$
(2)
where $`r_{\mathrm{halo}}`$ ($`r_s`$ in NFW) is the characteristic inner scale radius of the halo. The normalization, $`\rho _\mathrm{s}`$, is set by the requirement that the mass interior to $`R_{\mathrm{vir}}`$ be equal to $`M_{\mathrm{vir}}`$. The value of $`r_{\mathrm{halo}}`$ is usually characterized in terms of of the halo “concentration” parameter: $`cR_{\mathrm{vir}}/r_{\mathrm{halo}}`$. The implied maximum circular velocity for this profile occurs at a radius $`r_{\mathrm{max}}2.15r_{\mathrm{halo}}`$ and takes the value $`V_{\mathrm{max}}0.466V_{\mathrm{vir}}F(c)`$, where $`F(c)=\sqrt{c/[\mathrm{ln}(1+c)c/(1+c)]}`$.
For satellites, we set the value of $`c`$ using the simulation results of Bullock et al. (2001) and the corresponding relationship between halo mass, redshift, and concentration summarized by their analytic model. The median $`c`$ relation for halos of mass $`M_{\mathrm{vir}}`$ at redshift $`z`$ is given approximately by
$$c9.6\left(\frac{M_{\mathrm{vir}}}{10^{13}M_{}}\right)^{0.13}(1+z)^1,$$
(3)
although in practice we use the full analytic model discussed in Bullock et al. (2001).
For parent halos, we allow their concentrations to evolve self-consistently as their virial masses increase, as has been seen in the N-body simulations of Wechsler et al. (2002). Rather than represent the halo growth as a series of discrete accretion events, we smooth over the Monte Carlo EPS merger tree by fitting the following functional form to the Monte Carlo mass accretion history for each halo:
$`M_{\mathrm{vir}}(a)`$ $`=`$ $`M_{\mathrm{vir}}(a_0)\mathrm{exp}\left[2a_c\left({\displaystyle \frac{a_0}{a}}1\right)\right].`$ (4)
Here $`a(1+z)^1`$ is the expansion factor, and $`a_c`$ is the fitting parameter, corresponding to the value of the expansion factor at a characteristic “epoch of collapse”. Wechsler et al. (2002) demonstrated that the value of $`a_c`$ connects in a one-to-one fashion with the halo concentration parameter (Wechsler et al., 2002):
$$c(a)=5.1\frac{a}{a_c}.$$
(5)
Halos that form earlier (smaller $`a_c`$’s) are more concentrated.
Example fits to four of our halo mass accretion histories are shown by the smooth solid lines in Figure 1. The $`a_c`$ values for each of the halos in this analysis are listed in the third column of Table 1. Typical host halos in our sample have $`c14`$ at $`z=0`$, scale radii $`r_{\mathrm{halo}}20`$kpc, and maximum circular velocities $`V_{\mathrm{max}}190\mathrm{k}\mathrm{m}\mathrm{s}^1`$.
### 2.2. N-body simulations of dark matter evolution
Having determined the mass, accretion time and orbit of each satellite (§2.1.1 and §2.1.2), and the evolution the potential into which it is falling (§2.1.3), we next run individual N-body simulations to track the dynamical evolution of each satellite halo separately. We follow only those that contain a significant stellar component (see §2.3 below). In practice, this restricts our analysis to satellite halos more massive than $`M_{\mathrm{vir}}\stackrel{>}{}10^8M_{}`$ — the number of such satellites infalling into each parent is listed in column 5 of Table 1. Based on our star-formation prescription discussed in §2.3, systems smaller than this never contain an appreciable number of stars and thus don’t contribute significantly to the stellar halo.
#### 2.2.1 The parent galaxy potential
The parent galaxy is represented by a three-component bulge/disk/dark halo potential which we allow to evolve with time as the halo accretes mass. The (spherically-symmetric) dark halo potential at each epoch $`a`$ is given by the NFW potential generated by the dark matter distribution in equation (2)
$`\mathrm{\Phi }_{\mathrm{halo}}(r)={\displaystyle \frac{GM_{\mathrm{halo}}}{r_{\mathrm{halo}}}}{\displaystyle \frac{1}{(r/r_{\mathrm{halo}})}}\mathrm{ln}\left({\displaystyle \frac{r}{r_{\mathrm{halo}}}}+1\right),`$ (6)
equation where $`M_{\mathrm{halo}}=M_{\mathrm{halo}}(a)`$ and $`r_{\mathrm{halo}}=r_{\mathrm{halo}}(a)`$ are the instantaneous mass and length scales of the halo respectively. The halo mass scale is related to the virial mass via
$`M_{\mathrm{halo}}`$ $`=`$ $`{\displaystyle \frac{M_{\mathrm{vir}}}{\mathrm{ln}(c+1)c/(c+1)}}.`$ (7)
The disk and bulge are assumed to grow in mass and scale with the halo virial mass and radius:
$$\mathrm{\Phi }_{\mathrm{disk}}(R,Z)=\frac{GM_{\mathrm{disk}}}{\sqrt{R^2+\left(R_{\mathrm{disk}}+\sqrt{Z^2+Z_{\mathrm{disk}}^2}\right)^2}},$$
(8)
$$\mathrm{\Phi }_{\mathrm{sphere}}(r)=\frac{GM_{\mathrm{sphere}}}{r+r_{\mathrm{sphere}}},$$
(9)
where $`M_{\mathrm{disk}}(a)=1.0\times 10^{11}(M_{\mathrm{vir}}/M_{\mathrm{vir},0})M_{}`$, $`M_{\mathrm{sphere}}(a)=3.4\times 10^{10}(M_{\mathrm{vir}}/M_{\mathrm{vir},0})M_{}`$, $`R_{\mathrm{disk}}=6.5(r_{\mathrm{vir}}/r_{\mathrm{vir},0})`$ kpc, $`Z_{\mathrm{disk}}=0.26(r_{\mathrm{vir}}/r_{\mathrm{vir},0})`$ kpc and $`r_{\mathrm{sphere}}=0.7(r_{\mathrm{vir}}/r_{\mathrm{vir},0})`$ kpc.
#### 2.2.2 Satellite initial conditions
We use $`10^5`$ particles to represent the dark matter in each accreted satellite. Particles are initially distributed as an isotropic NFW model, with mass and scale chosen as described in §2.1.2. The phase-space distribution function is derived by integrating over the density and potential distributions
$$f(ϵ)=\frac{1}{8\pi ^2}\left[_0^ϵ\frac{d^2\rho }{d\mathrm{\Psi }^2}\frac{d\mathrm{\Psi }}{\sqrt{ϵ\mathrm{\Psi }}}+\frac{1}{\sqrt{ϵ}}\left(\frac{d\rho }{d\mathrm{\Psi }}\right)_{\mathrm{\Psi }=0}\right].$$
(10)
with $`\rho =\rho _{_{\mathrm{NFW}}}`$ and where $`\mathrm{\Psi }=\mathrm{\Phi }_{_{\mathrm{NFW}}}+\mathrm{\Phi }_0`$ is the relative potential (such that $`\mathrm{\Psi }0`$ as $`r\mathrm{}`$) and $`ϵ=\mathrm{\Psi }v^2/2`$ is the relative energy (see Binney & Tremaine, 1987, for discussion). This distribution function is used (in tabulated form) to generate a random realization. This ensures a stable satellite configuration — initial conditions generated by instead assuming a local Maxwellian velocity distribution have been shown to evolve (Kazantzidis, Magorrian, & Moore, 2004). Given $`f(ϵ)`$, the differential energy distribution follows in a straightforward manner from the density of states, $`g(ϵ)`$,
$$\frac{dM}{dϵ}=f(ϵ)g(ϵ),g(ϵ)16\pi ^2_0^{r_ϵ}\sqrt{2(\mathrm{\Psi }ϵ)}r^2𝑑r,$$
(11)
where $`r_ϵ`$ is the largest energy that can be reached by a star of relative energy $`ϵ`$. The differential energy distribution for our initial halo is shown by the solid histogram in Figure 2. We see that the majority of the (dark matter) material in an infalling satellite is quite loosely bound.
Rather than generating a unique $`f(ϵ)`$ and particle distribution for each satellite in each accretion history, a single initial conditions file with unit mass and scale, and outer radius $`R_{\mathrm{out}}=35r_{\mathrm{halo}}`$ ($`=35`$ in our units) is used for all simulations with masses and scales appropriately rescaled for each run. Since all of our accreted satellites have concentrations $`c<35`$, our set up effectively allows each accreted satellite’s mass profile to extend beyond its virial radius for several scale lengths. We do not expect this simplification to significantly affect our results because the the light matter is always embedded at the very central regions of the halo ($`r_{}\stackrel{<}{}r_{\mathrm{halo}}`$) and the outer material is always quickly stripped away from the outer parts of the halos upon accretion.
In §2.4 we discuss our method of “embedding” star particles within the cores the accreted satellite dark halos.
#### 2.2.3 Satellite evolution
The mutual interactions of the satellite particles are calculated using a basis function expansion code (Hernquist & Ostriker, 1992). The initial conditions file for the satellite is allowed to relax in isolation for ten dynamical times using this code to confirm stability. For each accretion event a single simulation is run, following the evolution of the relaxed satellite under the influence of its own and the parent galaxy’s potential, for the time since it was accreted (as generated by methods in §2.1.1) along the orbit chosen at random from the distribution discussed in §2.1.2. (Note that simulations of satellite accretions in static NFW potentials using this code produced results identical to those reported in Hayashi et al., 2003)
Using this approach, the satellites are not influenced by each other, other than through the smooth growth of the parent galaxy potential. Nor does the parent galaxy react to the satellite directly. In order to mimic the expected decay of the satellite orbits due to dynamical friction (i.e. the interaction with the parent), we include a drag term on all particles within two tidal radii $`r_{\mathrm{tide}}`$ of the satellite’s center, of the form proposed by Hashimoto, Funato, & Makino (2003) and modified for NFW hosts by Zentner & Bullock (2003). This approach includes a slight modification to the standard Chandrasekhar dynamical friction formula (e.g. Binney & Tremaine, 1987). The tidal radius $`r_{\mathrm{tide}}`$ is calculated from the instantaneous bound mass of the satellite $`m_{\mathrm{sat}}`$, the distance $`r`$ of the satellite to the center of the parent galaxy and the mass of the parent galaxy within that radius $`M_r`$ as $`r_{\mathrm{tide}}=r(m_{\mathrm{sat}}/M_r)^{1/3}`$.
#### 2.2.4 Increasing phase-space resolution with test particles
In this study, we are most interested in following the phase-space evolution of the stellar material associated with each satellite. This is assumed to be embedded deep within each dark matter halo (see §2.4) — typically only of order $`10^4`$ of the N-body particles in each satellite have any light associated with them at all. In order to increase the statistical accuracy our analysis we sample the inner 12% of the energy distribution with an additional $`1.2\times 10^5`$ test particles. This does not increase the dynamic range our simulation, but does allow us to more finely resolve the low surface brightness features we are interested in with only a modest increase in computational cost: we gain a factor of 10 in particle resolution with an increase of $``$25% in computing time. In this paper, we have used test particles only in generating the images shown in Figures 13 \- 16.
### 2.3. Following the satellites’ baryonic component
We follow each satellite’s baryonic component using the expected mass accretion history of each satellite halo (prior to falling into the parent galaxy) in order to track the inflow of gas. The gas mass is then used to determine the instantaneous star formation rate and to track the buildup of stars within each halo. The physics of galaxy formation is poorly understood, and any attempt to model star formation and gas inflow into galaxies (whether semi-analytic or hydrodynamic) necessarily require free parameters. Our own prescription requires three ”free” parameters: $`z_{\mathrm{re}}`$, the redshift of reionization (see §2.3.1); $`f_{\mathrm{gas}}`$, the fraction of baryonic material in the form of cold gas (i.e. capable of forming stars) that remains bound to each satellite at accretion (see §2.3.2); and $`t_{}`$, the globally-averaged star formation timescale (see §2.3.3).
In the following subsections we describe how these parameters enter into our prescriptions, and choose a value of $`f_{\mathrm{gas}}`$ consistent with observations. In §3 we go on to demonstrate that the observed characteristics of the stellar halo (e.g. its mass, and radial profile) and the Milky Way’s satellite system (e.g. their number and distribution in structural parameters) provide strong constraints on the remaining free parameters and hence the efficiency of star formation in low-mass dark matter halos in general.
#### 2.3.1 Reionization
Any attempt to model stellar halo buildup within the context of $`\mathrm{\Lambda }`$CDM must first confront the so-called “missing satellite problem” — the apparent over-prediction of low-mass halos compared to the abundance of satellite galaxies around the Milky Way and M31. For example, there are eleven know satellites of the Milky Way — nine classified as dwarf spheroidal and two as dwarf irregulars — yet numerical work predicts several hundred dark matter satellite halos in a similar mass range (Klypin et al., 1999; Moore et al., 1999). It is quite likely that our inventory of stellar satellites is not complete given the luminosity and surface brightness limits of prior searches (as the recent discovery of the Ursa Minor dwarf spheroidal demonstrates, see Willman et al., 2005), but incompleteness is not seen as a viable solution for a problem of this scale (see Willman et al., 2004, for a discussion).
The simplest solution to this problem is to postulate that only a small fraction of the satellite halos orbiting the Milky Way host an observable galaxy. In this work, we solve the missing satellite problem using the suggestion of Bullock, Kravtsov, & Weinberg (2000), which maintains that only the $`10\%`$ of low-mass galaxies ($`V_{\mathrm{max}}<30\mathrm{k}\mathrm{m}\mathrm{s}^1`$) that had accreted a substantial fraction of their gas before the epoch of reionization host observable galaxies (see also Chiu et al., 2001; Somerville, 2002; Benson et al., 2002; Kravtsov et al., 2004). The key assumption is that after the redshift of hydrogen reionization, $`z_{\mathrm{re}}`$, gas accretion is suppressed in halos with $`V_{\mathrm{max}}<50\mathrm{k}\mathrm{m}\mathrm{s}^1`$, and completely stopped in halos with $`V_{\mathrm{max}}<30\mathrm{k}\mathrm{m}\mathrm{s}^1`$. These thresholds follow from the results of Thoul & Weinberg (1996) and Gnedin (2000) who used hydrodynamic simulations to show that gas accretion in low-mass halos is indeed suppressed in the presence of an ionizing background.
We also impose a low-mass cutoff for tracking galaxy formation in satellite halos with $`V_{\mathrm{max}}<15\mathrm{k}\mathrm{m}\mathrm{s}^1`$. Two processes and one practical consideration motivate us to ignore galaxy formation in these tiny halos: first, photo-evaporation acts to eliminate any gas that was accreted before reionization in halos with $`V_{\mathrm{max}}\stackrel{<}{}15\mathrm{k}\mathrm{m}\mathrm{s}^1`$ (Barkana & Loeb, 1999; Shaviv & Dekel, 2003); second, the cooling barrier below virial temperatures of $`10^4`$K (corresponding to $`V_{\mathrm{max}}16\mathrm{k}\mathrm{m}\mathrm{s}^1`$) prevents any gas that could remain bound to these halos from cooling and forming stars (Kepner et al., 1997; Dekel & Woo, 2003); finally, even if we were to allow star formation in these systems, their contribution to the stellar halo mass would be negligible. Once we are more confident of our inventory of the lowest luminosity and lowest surface brightness satellites (Willman et al., 2004)of the Milky Way we should be able to confirm these physical arguments with observational constraints.
The epoch of reionization $`z_{\mathrm{re}}`$ determines the numbers of galaxies that have collapsed in each of the above $`V_{\mathrm{max}}`$ limits, and hence the number of luminous satellites that will be accreted, whether they disrupt to form the stellar halo or survive to form the Galaxy’s satellite system. We discuss limits on this parameter in §3.1.1.
#### 2.3.2 Gas accretion following reionization
The virial mass of each satellite, $`M_{\mathrm{vir}}^{\mathrm{sat}}`$, at the time of its accretion, $`a_{\mathrm{ac}}`$, is set by our merger tree initial conditions (§2.1.1). We assume that each satellite halo has had a mass accumulation history set by Equation 4 up to the time of its merger into the “Milky Way” host, with $`a_0=a_{\mathrm{ac}}`$. After accretion, all mass accumulation onto the satellite is truncated (see §2.3.3). For massive satellites, $`V_{\mathrm{max}}>50\mathrm{k}\mathrm{m}\mathrm{s}^1`$, we set $`a_c`$ in Equation 4 using the satellite’s mass-defined concentration parameter via Equation 5 (see Bullock et al., 2001; Wechsler et al., 2002). This provides a “typical” formation history for each satellite. For low-mass satellites, we are necessarily interested in where $`a_c`$ falls in the distribution of halo formation epochs because this determines the fraction of mass in place at reionization. Therefore, if $`V_{\mathrm{max}}<50\mathrm{k}\mathrm{m}\mathrm{s}^1`$, we use the methods of Lacey & Cole (1993) in order to derive the fraction of the satellite’s mass that was in place at the epoch of reionization, $`z_{\mathrm{re}}`$, and use this to set the value of $`a_c`$. Given $`a_c`$ for each satellite, we determine the instantaneous accretion rate of dark matter $`h(t)`$ in to this system as a function of cosmic time via
$$h(t)=\frac{\mathrm{d}M_{\mathrm{vir}}^{\mathrm{sat}}}{\mathrm{d}t}.$$
(12)
In the absence of radiative feedback effects, cooling is extremely efficient in pre-merged satellites of the size we consider (see, e.g. Maller & Bullock, 2004). Therefore we expect the cold gas inflow rate to track the dark matter accretion rate, $`h(t)`$ — at least in the absence of the effects of reionization — and take it to be
$$Cf_{\mathrm{gas}}h(tt_{\mathrm{in}}).$$
(13)
The time lag within $`h(t)`$ accounts for the finite time it takes for gas to settle into the center of the satellite after being accreted. We assume this occurs in roughly a halo orbital time at the virial radius: $`t_{\mathrm{in}}=\pi R_\mathrm{h}/V_{\mathrm{vir}}6\mathrm{Gyr}(1+z)^{3/2}`$. We have introduced the constant $`C`$ in order to account for the suppression of gas accretion in low-mass halos (as alluded to in §2.3.1). Before the epoch of reionization, we set $`C=1`$ for all galaxies. For systems with $`V_{\mathrm{max}}>50\mathrm{k}\mathrm{m}\mathrm{s}^1`$, $`C=1`$ at all times. After reionization, $`C=0`$ in systems with $`V_{\mathrm{max}}<30\mathrm{k}\mathrm{m}\mathrm{s}^1`$, and $`C`$ varies linearly in $`V_{\mathrm{max}}`$ between $`0`$ and $`1`$ if $`V_{\mathrm{max}}`$ falls between $`30`$ and $`50\mathrm{k}\mathrm{m}\mathrm{s}^1`$ (see Thoul & Weinberg, 1996).
The fraction of mass in each satellite in the form of cold, accreting baryons, $`f_{\mathrm{gas}}`$, determines the total stellar mass plus cold gas mass associated with each dark matter halo. In what follows, we adopt $`f_{\mathrm{gas}}=0.02`$, which is an upper limit on the range of cold baryonic mass fraction in observed galaxies (Bell et al. 2003).
#### 2.3.3 Star Formation
If we assume that cold gas forms stars over a timescale $`t_{}`$, then the evolution of stellar mass $`M_{}`$ and cold gas mass $`M_{\mathrm{gas}}`$ follows a simple set of equations:
$`{\displaystyle \frac{\mathrm{d}M_{}}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{M_{\mathrm{gas}}}{t_{}}},`$ (14)
$`{\displaystyle \frac{\mathrm{d}M_{\mathrm{gas}}}{\mathrm{d}t}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}M_{}}{\mathrm{d}t}}+Cf_{\mathrm{gas}}h(tt_{\mathrm{in}}).`$ (15)
For simplicity, the star formation is truncated soon after each satellite halo is accreted onto the Milky Way host. Physically, this could result from gas loss via ram-pressure stripping from the background hot gas halo (Lin & Faber, 1983; Moore & Davis, 1994; Blitz & Robishaw, 2000; Maller & Bullock, 2004; Mayer et al., 2005). This model is broadly consistent with observations that demonstrate that the gas fraction in satellites of the Milky Way and Andromeda is typically far less than that in field dwarfs in the Local Group, as illustrated by the separation of the open (satellites) and filled (field dwarfs) symbols in Figure 3 (plotting data taken from Grebel, Gallagher, & Harbeck, 2003). Of course, this assumption is over-simplified, but it allows us to capture in general both the expectations of the hierarchical picture and the observational constraints. We note that this is likely a bad approximation for massive satellites, whose deep potential wells will tend to resist the effects of ram-pressure stripping. However, we expect that this will have little impact on our stellar halo predictions, since most of the stellar halo is formed from satellites that are accreted early and destroyed soon after.
The star formation timescale, $`t_{}`$ determines the star to cold gas fraction in each satellite upon accretion and, for a given value of $`f_{\mathrm{gas}}`$, total stellar luminosity associated with each surviving satellite and the stellar halo. We discuss limits on this parameter in §3.1.2.
### 2.4. Embedding baryons within the dark matter satellites
We model the evolution of a two-component population of stellar matter and dark matter in each satellite by associating stellar matter with the more tightly bound material in the halo. As discussed in §2.1.3, the mass profile of the satellite is assumed to take the NFW form. Mass-to-light ratios for each particle are picked based on the particle energy in order to produce a realistic stellar profile for a dwarf galaxy.
A phenomenologically-motivated approximation for the stellar distribution in dwarf galaxies is the spherically symmetric King profile (King 1962):
$$\rho _{}(r)=\frac{K}{x^2}\left(\frac{\mathrm{cos}^1(x)}{x}\sqrt{1x^2}\right),x\frac{1+(r/r_\mathrm{c})^2}{1+(r_\mathrm{t}/r_\mathrm{c})^2}.$$
(16)
The core radius is $`r_\mathrm{c}`$ and $`r_\mathrm{t}`$ is the tidal radius, where $`\rho _{}(r>r_\mathrm{t})=0`$. The normalization, $`K`$, is set by the average density of the satellite, determined by its mass (2.3.3) and size scales (discussed below).
For each satellite, we assume a stellar mass to light ratio of $`M_{}/L_\mathrm{V}=2`$, and use the stellar mass calculated in §2.3.3 in order to assign a median King core radius
$$r_c=160\mathrm{p}\mathrm{c}\left(\frac{L_{}}{10^6L_{}}\right)^{0.19},$$
(17)
where throughout $`L_{}`$ is assumed to be the V-band stellar luminosity. We allow scatter about the relation using a uniform logarithmic deviate between $`0.3\mathrm{\Delta }\mathrm{log}_{10}L0.3`$. This slope and normalization was determined by least-square fit to the luminosity and core size correlation for the dwarf satellite data presented in Mateo (1998), and the scatter was determined by a “by-eye” comparison to the scatter in the data about the relation. Our adopted relation between $`r_c`$ and $`L_{}`$ is also consistent with the relevant projection of the fundamental plane for dwarf galaxies (e.g. Kormendy 1985). For all satellites we adopt $`r_\mathrm{t}/r_\mathrm{c}=10.`$
Assuming isotropic orbits for the stars and that the gravitational potential is completely dominated by the dark matter, the stellar energy distribution function corresponding to the King profile $`f_{}(ϵ)`$ is determined by setting $`\rho =\rho _{}`$ and $`\mathrm{\Psi }=\mathrm{\Phi }_{_{\mathrm{NFW}}}+\mathrm{\Phi }_0`$ in equation (10). The mass-to-light ratio of a particle of energy $`ϵ`$ is then simply $`f_{}(ϵ)/f(ϵ)=(dM/dϵ)_{}/(dM/dϵ)`$. Three examples are given Figure 2.
### 2.5. Limitations of our method
The main limitation of our method is that it only follows the smooth growth of the parent potential analytically — the satellite/satellite interactions and reaction of the parent to the satellite are not modeled self-consistently. Hence we do not anticipate following the evolution of the field or satellite particles during a major or even minor merger event with great accuracy. Given this limitation, we only simulate the accretion histories of halos generated from the Monte Carlo merger tree code that have not suffered a significant merger ($`>`$10 % of the parent halo mass) in the recent past ($`<`$7 Gyr) — 11 of the 20 accretion histories generated met this criterion. In addition, we consider results from simulations of accretion events that have occurred prior to the last significant merger to be less reliable. We label the halos used in this work 1-11. The five left-hand columns of Table 1 summarize the properties of the simulations run for each halo.
Even with these restrictions, we consider our study to be a useful approach for exploring substructure in galaxy halos because: (i) the highest surface brightness features in halos are likely to have come from recent events, whose debris has had a shorter time to phase-mix and/or be dispersed by oscillations in the parent galaxy potential; and (ii) substructure should be more readily detectable around spiral (rather than elliptical) galaxies because their stellar distributions are less extended — the existence of disks in spirals suggests that these are the ones the more quiescent accretion histories.
## 3. Results I: Tests of the Model
As outlined in §2.5, our method most accurately follows the phase-space evolution of debris from accretion events that occur during relatively quiescent times in a galaxy’s history (which we define as being after the last $`>`$10% merger event). In future work we concentrate on those events. In this paper, we analyze the results from simulations of the full accretion histories of our halos. While not accurate in following the phase-space properties of debris material from events occurring before the epoch of major merging, the fact that these systems are disrupted is predicted robustly, and we are able to record the time of disruption and the cumulative mass in those disrupted events as well.
In what follows, we first constrain the remaining free parameters $`z_{\mathrm{re}}`$ (§3.1.1) and $`t_{}`$ (§3.1.2) (with $`f_{\mathrm{gas}}=0.02`$), by requiring that the general properties of our surviving satellite populations are consistent with those of the Milky Way’s own satellites. We then go on to demonstrate that these parameter choices naturally produce the observed distributions in and correlations of the structural parameters of surviving satellites (see §3.2.1 and 3.2.2), as well as stellar halos with total luminosity and radial profiles consistent with the Milky Way (§3.2.3).
### 3.1. Primary constraints on parameters
#### 3.1.1 Satellite number counts
As described in §2.3.1, we have chosen to solve the missing satellite problem by suppressing gas accretion in small halos after the epoch of reionization, $`z_{\mathrm{re}}`$, and suppressing gas accumulation all together in satellites smaller than $`15\mathrm{k}\mathrm{m}\mathrm{s}^1`$. The number of satellites that host stars is then set by choosing $`z_{\mathrm{re}}`$. In the work presented in this paper, we assume that reionization occurred at a redshift $`z_{\mathrm{re}}=10`$ or at a lookback time of $`13`$Gyr. The fifth column of Table 1 gives the number of luminous satellites accreted over the lifetime of each halo and the sixth column gives the number of luminous satellites that survive disruption in each. (The numbers in brackets are for those events since the last $`>`$10% merger.) We see that our reionization prescription leads to agreement within a factor of $`2`$ with the number of satellites observed orbiting the Milky Way. Our results are roughly insensitive to this choice as long as $`8\stackrel{<}{}z_{\mathrm{re}}\stackrel{<}{}15`$.
#### 3.1.2 Infalling satellite gas content
When reviewing the properties of Local Group dwarf galaxies it is striking that — with the notable exceptions of the Large and Small Magellanic Clouds (hereafter LMC and SMC) — satellites of the Milky Way and Andromeda galaxies are exceedingly gas-poor compared to their field counterparts (Mateo, 1998; Grebel, Gallagher, & Harbeck, 2003). Figure 3 emphasizes this point by plotting the V-band luminosity vs gas fraction from the compilation by Grebel, Gallagher, & Harbeck (2003) for satellites (open squares) and field dwarfs (filled squares). We see that field dwarfs tend to have $`M_{HI}/L_V0.33`$, whereas satellite dwarfs have gas fractions $`0.0010.1`$.
While our star formation model assumes that most of the gas in accreted dwarfs is lost shortly after a dwarfs becomes a satellite galaxy, consistency with the field dwarf population requires that the most recent events in our simulations have gas-to-star ratios of order unity immediately prior to their accretion. This requirement forces us to choose a long star formation timescale, $`t_{}=15`$ Gyr, comparable to the Hubble time. Figure 4 shows the ratio $`M_{\mathrm{gas}}/L_\mathrm{V}`$ each satellite at the time it was accreted for our four example halos. The clear trend with accretion time follows because early accreted systems have not had time to turn their gas into stars. Solid points indicate satellites that survive until the present day. We see that the most recently accreted systems ($`t_{\mathrm{accr}}12`$ Gyr, those that should correspond most closely with true “field” dwarfs today) have $`M_{\mathrm{gas}}/L_{}12`$, which is in reasonable agreement with the gas content of field dwarfs. The points along the lower edge of the trend have lower gas fractions at a fixed accretion time because they stopped accreting gas at reionization (see §2.3.1).
Our choice of $`t_{}=15`$Gyr is much longer than is typical for semi-analytic prescriptions of galaxy formation set within the CDM context (e.g., Somerville & Primack, 1999), but these usually focus on much larger galaxies than the dwarfs we focus on here, where star formation is likely to have proceeded more efficiently. Observations suggest that the dwarf spheroidal satellites of the Milky Way have rather bursty, sporadic star formation histories, with recent star formation in some cases (Grebel, 2000; Smecker-Hane & Mc William, 1999; Gallart et al., 1999). This effectively demands that the star formation timescales must be long in these systems: our model can be viewed as smoothing over these histories with an average low-level of star formation.
Note that we do not explicitly include supernova feedback in our star formation histories, but it is implicitly included by requiring a very low level of efficiency in our model (i.e. a large value of $`t_{}`$). In two companion papers we do include the effects of feedback (accounting for both gas gained due to mass loss from stars during normal stellar evolutionary phases and gas lost via winds driven by supernovae) in order to accurately model chemical enrichment in our accreted satellites (Robertson et al., 2005; Font et al., 2005). With feedback included, a choice of $`t_{}=6.75`$Gyr provides nearly identical distributions of gas and stellar mass in satellite galaxies as does our non-feedback choice of $`t_{}=15`$Gyr.
### 3.2. Verification of Model’s Validity
#### 3.2.1 Distributions in satellite structural parameters
Figure 5 shows histograms of the fractional number of satellites as a function of central surface brightness $`\mu _0`$, total luminosity $`L_{}`$ and central line-of-sight velocity dispersion $`\sigma _{}`$ for the Milky Way dwarf spheroidal satellites in solid lines. The dashed lines represent our simulated distribution of surviving satellite properties, derived by combining the structural properties of the 156 surviving satellites from all eleven halos. The histograms are visually similar. (Note that the LMC and SMC are not included in the observational data set since they are rotationally supported and our models are restricted to hot systems. They would be equivalent to the most luminous, highest velocity dispersion systems in our model data set that appear to be missing from the Milky Way distribution.)
To quantify the level of similarity of the simulated and observed data sets we use the 3-dimensional KS-statistic (Fasano & Franceschini, 1987)
$$Z_{n,3\mathrm{D}}=d_{\mathrm{max}}\sqrt{n},$$
(18)
where $`n`$ is the number in the sample tested against our model parent distribution of all 156 surviving satellites. In this method $`d_{\mathrm{max}}`$ is defined as the maximum difference between the observed and predicted normalized integral distributions, cumulated within the eight volumes of the three-dimensional space defined for each data point $`(X_i,Y_i,Z_i)=(\mu _{0,i},L,i,\sigma _{,i})`$ by
$$(x<X_i,y<Y_i,z<Z_i),\mathrm{},(x>X_i,y>Y_i,z>Z_i).$$
(19)
Fasano & Franceschini (1987) present assessments of the significance level of values obtained for $`Z_{n,3\mathrm{D}}`$ as a function of $`n`$ and of the degree of correlation of the data. Since we already have eleven similarly-sized samples drawn from the same parent distribution, we instead quantify the significance level of $`Z_{n,3\mathrm{D}}`$ found for the Milky Way satellites by comparing it to the distribution of $`Z_{n,3\mathrm{D}}`$ for our simulated samples. Figure 6 shows a histogram of the results for our simulated halos, with the dotted line indicating where the Milky Way satellite distribution falls. According to this test only one of the eleven simulated populations is more similar to the simulated parent population than the observed satellites. (Note that $``$80% of our simulated samples have $`Z_{n,3\mathrm{D}}<1.2`$. This significance level is similar to those derived by Fasano & Franceschini (1987) for 3-dimensional samples with $`n=10`$ and a moderate degree of correlation in the distribution — see their Figure 7 — as might be expected given the expected relation between $`\sigma _0`$ and $`L_{\mathrm{tot}}`$, see §3.2.2.)
#### 3.2.2 Correlations in satellite structural parameters
Figure 7 shows the relationship between the central ($`<r_c`$), 1-D light-weighted velocity dispersion and satellite stellar mass, $`M_{}`$, for model galaxies and observed galaxies in the Local Group. Crosses show surviving model satellites for all halos and open circles show the relationship for the same set of satellites before they were accreted into the host dark matter halo. Solid triangles show the relationship for Local Group satellites as compiled by Dekel & Woo (2003). The two nearly identical solid lines show the best-fit regressions for the initial and final model populations. The dashed line shows the best-fit line for the data. Our model galaxies reproduce a trend quite similar to that seen in the data. The relative agreement is significant for two reasons. First, the stellar velocity dispersion of our initialized satellites is set by the underlying potential well of their dark matter halos convolved with their associated King profile parameters. While in §2.4 we set King profile parameters using a phenomenological relation based on the stellar luminosity ($`L_{}r_c`$), there was no guarantee that the dark matter potential associated with a given luminosity would provide a consistent stellar velocity dispersion ($`r_c+\rho _{_{DM}}\sigma _{}`$). In this sense, the general agreement between model satellites and the data is a success of our star formation prescription, which varies based on the mass accretion histories of halos of a given size (and therefore density structure).
A second interesting feature shown in Figure 7 is that final surviving satellites obey the same relation as the initial satellites. Most of these systems have experienced significant dark matter mass loss, but since the star particles are more tightly bound, their velocity dispersion does not significantly evolve. This point is emphasized in Figure 8, where we plot the central ($`<r_c`$), 1-D velocity dispersion for the dark matter in halos, again as a function of the satellite galaxy’s stellar mass, $`M_{}`$. As in Figure 7, open circles show the relationship for the final, surviving satellites, and crosses show the relationship for those same satellites before they were accreted. Unlike in the case of light-weighted dispersions, the dark matter dispersion velocities in the surviving systems is systematically lower than in the initial halos owing to the loss of the most energetic particles. They also exhibit a broader scatter at fixed stellar mass, reflecting variations in their mass-loss histories. Comparing again to Figure 7 we see that most of the particles associated with light in these systems remains bound to the satellites and their velocity dispersions do not evolve significantly. This result may have important implications for interpreting the nature of the dark matter halos of dwarf galaxies in the Local Group and for understanding the regularity in observed dwarf properties irrespective of their environments. In future work we will return to a more detailed structural and evolutionary analysis of the light matter and the dark matter halos in which the stars are embedded.
The results presented in this and the previous sub-section clearly indicate that our star formation scenario coupled with setting the King parameters of our infalling dwarfs to match Local Group observations lead to surviving satellite populations consistent both in number and structural properties with the Milky Way’s.
#### 3.2.3 The stellar halo’s mass and density profile
Estimates for the size, shape and extent of the Milky Way’s stellar halo come either from star count surveys (Morrison et al., 2000; Chiba & Beers, 2000; Yanny et al., 2000; Siegel, Majewski, Reid, & Thompson, 2002) or studies where distances could be estimated using RR Lyraes (Wetterer & McGraw, 1996; Ivezić et al., 2000) . These studies agree on a total luminosity of order $`L_\mathrm{V}10^9L_{}`$ (or mass $`2\times 10^9M_{}`$), which is in good agreement with the unbound stellar luminosity for all eleven of our model stellar halos, listed in Column 6 of Table 1 (numbers in brackets again refer to stars from accretion events since the last $`>`$10% merger). The match between predicted and observed total halo mass is non-trivial and depends sensitively on the mass accretion history of the dark matter halo along with the value of the star-formation timescale, $`t_{}`$. Specifically, we show in §4.1.1 that the majority of dwarf galaxies that make up the stellar halo were accreted early, more than $`8`$Gyr ago. The total stellar halo mass ($`10^9M_{}`$) is relatively small compared to the total cold baryonic mass in accreted satellites ($`10^{10}M_{}`$), because the star formation timescale is long compared to the age of the Universe at typical accretion times, and the stellar mass fractions are correspondingly low (see Figure 4). If we would have chosen a star formation timescale short compared the time of typical accretion for a destroyed system (e.g. $`5`$Gyr) this would have resulted in a stellar halo of stripped stellar much more massive than that observed for the Milky Way. This is in agreement with the results of Brook et al. (2004b), who found that a strong feedback model (effectively slowing the star formation rate in dwarf galaxies) in their smoothed particle hydrodynamical simulations of galaxy formation was necessary in order to build relatively small halo components in their models.
The observational studies find density profiles falling more steeply than the dark matter halo (a power law index in the range -2.5 to -3.5, compared to $`2`$ for the dark matter at relevant scales). Some of the variance between results from different groups can be attributed to substructure in the halo since these studies have commonly been limited in sky-coverage with surveys covering significant portions of the sky only now becoming feasible. Figure 9 plots the density profiles generated (arbitrarily normalized) from our four representative stellar halo models (light solid curves), which transition between slopes of -2 within $`10`$kpc to -4 around 50kpc and fall off even more steeply beyond this. To illustrate the general agreement with observations, the dotted line is a power law with exponent of -3. Note that there is some variation in the total luminosity (about a factor of 2) and slopes of our model halos, as might be expected given their different accretion histories. There is also a clear roll-over below the power law in the outer parts of the stellar halo, sometimes at radii as small as 30kpc.
To contrast to the light, the density profiles of the dark matter in our models are plotted in bold lines Figure 9 (also with arbitrary normalization). The dark matter profiles are all close to an NFW profile with $`m_{\mathrm{halo}}=1.4\times 10^{12}M_{}`$ and $`r_{\mathrm{halo}}=10`$kpc. Within $`30`$kpc of the Galactic center it appears that our stellar halos roughly track the dark matter, but beyond this they tend to fall more steeply. The difference in profile shapes — and the steep roll-over in the light matter at moderate to large radii — is a natural consequence of embedding the light matter deep within the dark matter satellites: The satellites’ orbits can decay significantly before any of the more tightly bound material is lost. Hence we anticipate that more/less extended stellar satellites would result in a more/less extended stellar halo. Studies of the distant Milky Way halo are still sufficiently limited that it is not possible to say whether the location of the roll-over in our model stellar halos is in agreement with observations and this could be an interesting test of our models in the near future (see, e.g. Ivezić et al., 2003).
## 4. Results II: Model Predictions
We have now fixed our free parameters to be $`z_{\mathrm{re}}=10`$, $`t_{}=15`$Gyr and $`f_{\mathrm{gas}}`$=0.02. By limiting our description of the evolution of the baryons associated with each dark matter satellite to depend on only these parameters, we find we have little freedom in how we choose them. For example, if we were to choose a shorter star formation timescale $`t_{}`$, we would over-produce the mass of the stellar halo, form dwarf galaxies that were over-luminous at fixed velocity dispersions, and form dwarfs with low gas fractions compared to isolated dwarfs observed in the Local Group. The first two problems could be adjusted by adapting $`f_{\mathrm{gas}}`$, but the last problem is independent of this.
Despite its simplicity, our model reproduces observations of the Milky Way in some detail. In particular, we recover the full distribution of satellites in structural properties. This suggests both that have we assigned the right fraction of dark matter halos to be luminous and that our luminous satellites are sitting inside the right mass dark matter halos.
We can now go on with some confidence to discuss the implications of our model for the mass accretion history of the halo and satellite systems (§4.1) and the level of substructure in the stellar halo (§4.2).
### 4.1. Building up the stellar halo and satellite systems
#### 4.1.1 Accretion times and mass contributions of infalling satellites
The stellar halo in our model is formed from stars originally born in accreted satellites. Once accreted, satellites lose mass with time until the satellite is destroyed. Once a particle becomes unbound from a satellite, we associate its stellar mass with the stellar halo. Figure 10 shows the cumulative luminosity fraction of the stellar halo (solid lines) coming from accreted satellites as a function of the accretion time of the satellite for halos 1,2, 6, and 9. Clearly most of the mass in the stellar halos originates in satellites that were accreted more than $`8`$ Gyr ago. The dotted lines show the contribution to the stellar halo from satellite halos more massive than $`M_{\mathrm{vir}}>2\times 10^{10}M_{}`$ at the time of their accretion. While only $`1020`$ of the $`150`$ accreted satellites meet this mass requirement, we see that $`7590\%`$ of the mass associated with each stellar halo originated within massive satellites of this type.
Compare these to the dashed lines, which show the cumulative number fraction of surviving satellite galaxies as a function of the time they were accreted for the entire population (long-dashed lines) and restricted to satellite halos that were more massive than $`M_{\mathrm{vir}}\stackrel{>}{}5\times 10^9M_{}`$ at the time of their accretion (short dashed lines). We see that surviving satellites are accreted much later ($`35`$ Gyr lookback) than their destroyed counterparts and that the most massive satellites that survive tend to be accreted even later because the destructive effects of dynamical friction are more important for massive satellites.
#### 4.1.2 Spatial growth
Studies of dark matter halos in N-body simulations show that they are built from the inside out (e.g. Helmi, White, & Springel, 2003). The top panel of Figure 11 confirms that this idea holds for our model stellar halos: it shows the average over all our halos of the fraction of material in each spherical shell from all accretion events (solid line) and from those that have occurred since the last major ($`>10`$%) merger (dotted line — the time when this occurred is given in column 4 of Table 1). Although the recent events represent only a fraction of the total halo luminosity ($``$5-50%, see Table 1), they become the dominant contributor at radii of 30-60kpc and beyond. There is some suggestion of this being the case for the Milky Way’s halo globular cluster population, which can can fairly clearly be separated into an ’old’ , inner population (which exhibits some rotation, is slightly flattened and has a metallicity gradient) and a ’young’ outer one (which is more extended and has a higher velocity dispersion — see Zinn, 1993).
One implication of the inside-out growth of stellar halos, combined with the late accretion time of surviving satellites is that the two should each follow different radial distributions. The dashed lines in the bottom panel of Figure 11 shows the number fraction of all surviving satellites in our models as a function of radius — the distribution is much flatter than the one shown for the halo in the upper panel. In fact, all satellites of our own Milky Way (except Sgr) lie at or beyond 50kpc from its center, with most 50-150kpc away, — as shown by the solid line in the lower panel. Hence, the radial distribution of the observed satellites is consistent with our models and suggests that they do indeed represent recent accretion events.
#### 4.1.3 Implications for the abundance distributions of the stellar halo and satellites
Studies which contrast abundance patterns and stellar populations in the stellar halo with those in dwarf galaxies seem to be at odds with models (such as ours) that build the stellar halo from satellite accretion (Unavane, Wyse, & Gilmore, 1996). For example, both field and satellite populations have similar metallicity ranges, but the former typically have higher alpha-element abundances than the latter (Tolstoy et al., 2003; Shetrone et al., 2003; Venn et al., 2004). Clearly, it is not possible to build the halo from present-day satellites.
We have already shown (in Figure 10 ) that we would expect a random sample of halo stars to have been accreted 8-10Gyears ago from satellites with masses $`M_{\mathrm{vir}}\stackrel{>}{}10^{10}M_{}`$, while surviving satellites are accreted much more recently. (Note that Figure 10 deliberately compares the cumulative luminosity fraction of the stellar halo to the number fraction of satellites. This is the most relevant comparison to make when interpreting observations because any sample of halo stars will be weighted by the luminosity of the contributing satellites, while samples of satellite stars are often composed of a few stars from each satellite.) Figure 12 explores the number and luminosity contribution of different luminosity satellites to each population in more detail. It shows the number fraction of satellites in different luminosity ranges contributing to the stellar halo (dotted lines) and satellites system (dashed lines): the peak of the dotted/dashed lines at lower/higher luminosities is a reflection of the much later accretion time — and hence longer time available for growth of the individual contributors — of the satellite system relative to the stellar halo. However, as discussed above, it is more meaningful to compare the number fraction of surviving satellites to the luminosity fraction of the halo (solid lines) contributed by satellites of a given luminosity range. The solid line emphasizes (as noted above) that most of the stellar halo comes from the few most massive (and hence most luminous) satellites, with luminosities in the range $`10^710^9L_{}`$. In contrast, Galactic satellite stellar samples would likely be dominated by stars born in $`10^510^7L_{}`$ systems.
Overall our results provide a simple explanation of the difference between halo and satellite stellar populations and abundance patterns. The bulk of the stellar halo comes from massive satellites that were accreted early, and hence had star formation histories that must be short (because of their early disruption) and intense (in order to build a significant luminosity in the time before disruption). In contrast, surviving satellites are lower mass and accreted much later, and hence have more extended, lower level star formation histories. Stars formed in these latter environments represent a negligible fraction of the stellar halo in all our models. This is confirmed by the last column of Table 1, which lists the percentage contribution of surviving satellites to the total halo (less than 10% in every case). Note that the contributions of surviving satellites to the local halo (i.e. within 10-20kpc of the Sun), which is the only region of the halo where detailed abundance studies have been performed, are even lower (less than 1% in every case).
A more quantitative investigation of the consequences of the difference between the ”accretion age” of stars and satellites in the halo is underway (Robertson et al., 2005; Font et al., 2005).
### 4.2. Substructure
Abundant substructure is one of the most basic expectations for a hierarchically-formed stellar halo. Here we give a short description of the substructure we see in our simulations, and reserve more detailed and quantitative explorations for future work. Recall, our study (by design) follows the more recent accretion events in our halo more accurately than the earlier ones — we showed in §4.1.2 that these are the dominant contributors to the halo at radii of 30-60kpc and beyond. Hence we can expect our study make fairly accurate predictions of expectations of the level of substructure in the outer parts of galactic halos — precisely the region where substructure should be more dominant and easier to detect.
Figures 13 and 14 show external galactic views for halo realizations 1, 2, 6, and 9. The color code reflects surface brightness per pixel: white, 24 magnitudes per square arcsecond, to light blue at 30 magnitudes per square arcsecond, to black which is (fainter than) 38 magnitudes per square arcsecond. The darkest blue features are of course too faint to be seen (except by star counts). We have simply set the scale in order to reveal all the spatial features that are there in principle. We mention that our test particles (§2.2.4) were used in making these images.
In addition to spatial structure, we also expect significant structure in phase space. A two-dimensional slice of the full six dimensional phase space is illustrated in Figure 15, where we plot radial velocity $`V_r`$ versus radius $`r`$ for all of halo 1 (left) and halo 9 (right). Each point represents 1000 solar luminosities. <sup>1</sup><sup>1</sup>1 In most cases, we subsample our luminous particles in order to plot a single point for every $`1000L_{}`$. However, some of the particles in our simulations have luminosity weights greater than 1000 $`L_{}`$. In these cases we plot a number of points ($`=L_{\mathrm{particle}}/1000`$) with the same $`V_r`$ and $`r`$ as the relevant particle (using small random offsets to give the effect of a “bigger” point on the plot). The color code reflects the time the particle became unbound from its original satellite: dark blue for particles that became unbound more than 12 Gyr ago and white for particles that either remain bound or became unbound less than 1.5 Gyr ago. The radial gradient in color reflects the “inside out” formation of the stellar halo discussed in previous sections.
Note that significant coherent structure is visible in Figure 15 even without any spatial slicing of the halos. Except for particles belonging to bound satellites (white streaks), the structure strongly resembles a nested series of orbit diagrams. This is not surprising since the halo was formed by particles brought in on satellite orbits. A direct test of this prediction should be possible with SDSS data and other similar surveys. Indeed, if the phase space structure of stellar halo stars does reveal this kind of orbit-type structure it will be a direct indication that the stellar halo formed hierarchically. Figure 16 shows the same diagram for halo 1, now subdivided into two distinct quarters of the sky.
## 5. Summary and Conclusions
We have presented a cosmologically self-consistent model for the formation of the stellar halo in Milky-Way type galaxies. Our approach is hybrid. We use a semi-analytic formalism to calculate a statistical ensemble of accretion histories for Milky Way size halos and to model star formation in each accreted system. We use a self-consistent N-body approach to follow the dynamical evolution of the accreted satellite galaxies. A crucial ingredient in our model is the explicit distinction between the evolution of light and dark matter in accreted galaxies. Stellar material is much more tightly bound than the majority of the dark matter in accreted halos and this plays an important role in the final density distribution of stripped stellar material as well as the evolution in the observable quantities of satellite galaxies.
A primary goal of this, our first paper on stellar halo formation, was to normalize our model to, and demonstrate consistency with, the gross properties of the Milky Way stellar halo and its satellite galaxy population. We constrained our two main star formation parameters, the redshift of reionization $`z_{\mathrm{re}}`$, and the star formation timescale of cold gas $`t_{}`$, using the observed number counts of Milky Way satellites and the gas mass fraction of isolated dwarf galaxies. With these parameters fixed, the model reproduces many of the observed structural properties of (surviving) Milky Way satellites: the luminosity function, the luminosity-velocity dispersion relation, and the surface brightness distribution. The satellite galaxies that are accreted and destroyed in our model produce stellar halos of material with a total luminosity in line with estimates for the stellar halo of the Milky Way ($`10^9L_{}`$).
The success of our model lends support to the hierarchical stellar halo formation scenario, where the stellar halos of large galaxies form mainly via the accretion subsequent disruption of smaller galaxies. More specifically, it allows us to make more confident predictions concerning the precise nature of stellar halos and their associated satellite systems in Milky Way type galaxies. These include:
* The density profile of the stellar halo should follow a varying power-law distribution, changing in radial slope from $`2`$ within $`20`$kpc to $`4`$ beyond 50kpc. The distribution is expected to be much more centrally concentrated than the dark matter, owing to the fact that the stars that build the stellar halo were much more tightly bound to their host systems than the dark material responsible for building up the dark matter halo.
* Stellar halos (like dark matter halos) are expected to form from the inside out, with the majority of mass being deposited from the $`15`$ most massive accretion events, typically dwarf-irregular size halos with mass $`10^{10}M_{}`$ and luminosities of order $`10^710^9L_{}`$.
* Destroyed satellites contributing mass to the stellar halo tend to be accreted earlier than satellites that survive as present-day dwarf satellites ($`9`$Gyr compared to $`5`$ Gyr in the past).
* Substructure, visible both spatially and in phase space diagrams, should be abundant in the outer parts of galaxies. Proper counts of this structure, both in our galaxy and external systems, should provide important constraints on the late-time accretion histories of galaxies and a test of hierarchical structure formation.
Together, the second and third points imply that most of the stars in the inner halo are associated with massive satellites that were accreted $`\stackrel{>}{}9`$Gyr ago. Dwarf satellites, on the other hand, tend to be lower mass and are associated with later time accretion events. This suggests that classic “stellar halo” stars should be quite distinct chemically from stars in surviving dwarf satellites. We explore this point further in two companion papers (Robertson et al. 05; Font et al. 05).
KVJ’s contribution was supported through NASA grant NAG5-9064 and NSF CAREER award AST-0133617.
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# A direct approach to Bergman kernel asymptotics for positive line bundles
## 1. Introduction
Let $`L`$ be a holomorphic line bundle with a positively curved hermitian metric $`\varphi `$, over a complex manifold $`X.`$ Then $`i/2`$ times the curvature form $`\overline{}\varphi `$ of $`L`$ defines a Kähler metric on $`X,`$ that induces a scalar product on the space of global sections with values in $`L.`$ The orthogonal projection $`P`$ from $`L^2(X,L)`$ onto $`H^0(X,L),`$ the subspace of holomorphic sections, is the Bergman projection. Its kernel with respect to the scalar product is the Bergman kernel $`K`$ of $`H^0(X,L)`$; it is a section of $`\overline{L}L`$ over $`X\times X.`$ It can also be characterized as a reproducing kernel for the Hilbert space $`H^0(X,L),`$ i.e
(1.1)
$$\alpha (x)=(\alpha ,K_x)$$
for any element $`\alpha `$ of $`H^0(X,L),`$ where $`K_x=K(,x)`$ is identified with a holomorphic section of $`L\overline{L_x},`$ where $`L_x`$ denotes the fiber of $`L`$ over $`x`$. The restriction of $`K`$ to the diagonal is a section of $`\overline{L}L`$ and we let $`B(x)=\left|K(x,x)\right|`$ be its pointwise norm.
Even though the Bergman kernel is of course impossible to compute in general, quite precise asymptotic formulas, when we replace $`L`$ by $`L^k`$ and $`\varphi `$ by $`k\varphi `$, are known, see Zelditch and Catlin . Namely:
(1.2)
$$K_x(y)e^{k\psi (y,x)}=\frac{k^n}{\pi ^n}(1+\frac{b_1(x,y)}{k}+\frac{b_2(x,y)}{k^2}+\mathrm{})$$
as $`k`$ goes to infinity.
Here $`\psi `$ is an (almost) holomorphic extension of $`\varphi `$ and the $`b_j`$s are certain hermitian functions. In particular, on the diagonal $`y=x`$ we have an asymptotic series expansion for
$$K_x(x)e^{k\varphi (x)}=\frac{k^n}{\pi ^n}(1+\frac{b_1(x,x)}{k}+\frac{b_2(x,x)}{k^2}+\mathrm{}).$$
The functions $`b_j(x,x)`$ contain interesting geometric information of $`X`$ with the Kahler metric $`i\overline{}\varphi /2`$, see Lu .
In and the existence of an asymptotic expansion is proved using a formula, due to Boutet de Monvel and Sjöstrand, for the boundary behaviour of the Bergman-Szegö kernel for a strictly pseudoconvex domain, , extending an earlier result of C Fefferman, , to include also the off-diagonal behaviour. The purpose of this paper is to give a direct proof of the existence of this expansion, the main point being that it is actually simpler to construct an asymptotic formula directly. Our method also gives an effective way of computing the terms $`b_j`$ in the asymptotic expanison. Even though the inspiration for the construction comes from the calculus of Fourier integral operators with complex phase, the arguments in this paper are elementary.
The method of proof uses localization near an arbitrary point of $`X`$. Local holomorphic sections of $`L^k`$ in a small coordinate neighbourhood , $`U`$, are just holomorphic functions on $`U`$, and the local norm is a weighted $`L^2`$-norm over $`U`$ with weight function $`e^{k\varphi }`$ where $`\varphi `$ is a strictly plurisubharmonic function. Using the ideas from we then compute local asymptotic Bergman kernels on $`U`$. These are holomorphic kernel functions, and the scalar product with such a kernel function reproduces the values of holomorphic functions on $`U`$ up to an error that is small as $`k`$ tends to infinity. Assuming the bundle is globally positive it is then quite easy to see that the global Bergman kernel must be asymptotically equal to the local kernels.
Many essential ideas of our approach were already contained in the book written by the third authour. Here we use them in order to find a short derivation of the Bergman kernel asymptotics. For the closely related problem of finding the Bergman kernel for exponentially weighted spaces of holomorphic functions, this was done by A. Melin and the third author , but in the present paper we replace a square root procedure used in that paper by a more direct procedure, which we think is more convenient for the actual computations of the coefficients in the asymptotic expansions. There are also close relations to the subject of weighted integral formulas in complex analysis . We have tried to make the presentation almost self-contained, hoping that it may serve as an elementary introduction to certain micro-local techniques with applications to complex analysis and differential geometry.
## 2. The local asymptotic Bergman kernel
The local situation is as follows. Fix a point in $`X.`$ We may choose local holomorhic coordinates $`x`$ centered at the point and a holomorphic trivialization $`s`$ of $`L`$ such that
(2.1)
$$\left|s\right|^2=e^{\varphi (x)},$$
where $`\varphi `$ is a smooth real valued function. $`L`$ is positive if and only if all local functions $`\varphi `$ arising this way are strictly plurisubharmonic. We will call $`\varphi _0(x)=\left|x\right|^2`$ the *model fiber metric*, since it may be identified with the fiber metric of a line bundle of constant curvature on $`^n.`$ The Kähler form, $`\omega `$, of the metric on our base manifold $`X`$ is given by $`i/2`$ times the curvature form of $`L`$,
$$\omega =i\overline{}\varphi /2.$$
The induced volume form on $`X`$ is equal to $`\omega _n:=\omega ^n/n!`$. Now any local holomorphic section $`u`$ of $`L^k`$ may be written as $`us^k`$ where $`u`$ is a holomorphic function. The local expression of the norm of a section of $`L^k`$ over $`U`$ is then given by
$$u_{k\varphi }^2:=_U\left|u\right|^2e^{k\varphi }\omega _n$$
where $`u`$ is a holomorphic function un $`U.`$ The class of all such functions $`u`$ with finite norm is denoted by $`H_{k\varphi }(U).`$
We now turn to the construction of local asymptotic Bergman kernels. The main idea is that since a posteriori Bergman kernels will be quite concentrated near the diagonal, we require a local Bergman kernel to satisfy the reproducing formula (1.1) locally, up to an error which is exponentially small in $`k`$. In the sequel we fix our coordinate neighbourhood to be the unit ball of $`^n`$. Let $`\chi `$ be a smooth function supported in the unit ball $`B`$ and equal to one on the ball of radius $`1/2`$. We will say that $`K_k`$ is a *reproducing kernel* mod $`O(e^{\delta k})`$ for $`H_{k\varphi }`$ if for any fixed $`x`$ in some neighbourhood of the origin we have that for any local holomorphic function $`u_k,`$
(2.2)
$$u_k(x)=(\chi u_k,K_{k,x})_{k\varphi }+O(e^{k(\varphi (x)/2\delta )})u_{k\varphi },$$
uniformly in some neighbourhood of the origin. Furthermore, if $`K_{k,x}(y)`$ is holomorphic in $`y`$ we say that $`K_{k,x}`$ is a *Bergman kernel* mod $`O(e^{\delta k})`$.
Given a positive integer $`N,`$ Bergman and reproducing kernels mod $`O(k^N)`$ are similarly defined.
### 2.1. Local reproducing kernels mod $`e^{\delta k}`$
Let $`\varphi `$ be a strictly plurisubharmonic function in the unit ball. and let $`u`$ be a holomorphic function in the ball such that
$$u^2:=_B|u|^2e^{k\varphi }<\mathrm{}.$$
We shall first show that (cf ) integrals of the form
(2.3)
$$c_n(k/2\pi )^n_\mathrm{\Lambda }e^{k\theta (xy)}u(y)𝑑\theta dy,$$
define reproducing kernels mod $`e^{\delta k}`$ for suitably choosen contours
$$\mathrm{\Lambda }=\{(y,\theta );\theta =\theta (x,y)\}.$$
Here we think of $`x`$ as being fixed (close to the origin) and let $`y`$ range over the unit ball, so that $`\mathrm{\Lambda }`$ is a $`2n`$-dimensional submanifold of $`B_y\times _\theta ^n`$, and $`c_n=i^n(1)^{n(n+1)/2}=i^{n^2}`$ is a constant of modulus 1 chosen so that $`c_nd\overline{y}dy`$ is a positive form. Let us say that such a contour is good if uniformly on $`\mathrm{\Lambda }`$ for $`x`$ in some neighbourhood of the origin and $`|y|1`$
$$2\mathrm{R}\mathrm{e}\theta (xy)\delta |xy|^2\varphi (y)+\varphi (x).$$
Note that, by Taylor’s formula
$$\varphi (x)=\varphi (y)+2\mathrm{R}\mathrm{e}Q_i(x,y)(x_iy_i)+\varphi _{i\overline{j}}(x_iy_i)\overline{(x_jy_j)}+o(xy)^2,$$
where $`Q_i(x,y)(x_iy_i)`$ is the part of the second order Taylor expansion which is holomorphic in $`x`$. Hence, if $`\varphi `$ is strictly plurisubharmonic,
$$\theta (x,y)=Q(x,y)$$
is a good contour, depending on $`x`$ in a holomorphic way. In particular, $`\theta =\overline{y}`$ defines a good contour for $`\varphi (x)=|x|^2`$.
###### Proposition 2.1.
For any good contour,
$$u(x)=(k/2\pi )^nc_n_\mathrm{\Lambda }e^{k\theta (xy)}u(y)\chi (y)𝑑\theta dy+O(e^{k(\varphi (x)/2\delta )})u_{k\varphi },$$
for $`x`$ in some fixed neighbourhood of 0 if $`u`$ is an element of $`H_{k\varphi }(B)`$.
###### Proof.
For $`s`$ a real variable between 0 and $`\mathrm{}`$, we let
$$\mathrm{\Lambda }_s=\{(x,y,\theta );\theta +s(\overline{x}\overline{y})\mathrm{\Lambda }\},$$
and denote by $`\eta =\eta _k`$ the differential form
$$\eta =c_n(k/2\pi )^ne^{k\theta (xy)}u(y)\chi (y)d\theta dy.$$
Our presumptive reproducing formula is the integral, $`I_0`$, of $`\eta `$ over $`\mathrm{\Lambda }_0`$ and it is easy to see that the limit of
$$I_s:=_{\mathrm{\Lambda }_s}\eta $$
as $`s`$ goes to infinity equals $`u(x)`$. (This is because $`c_n(s/2\pi )^ne^{s|xy|^2}d\overline{y}dy`$ tends to a Dirac measure at $`x`$ as $`s`$ tends to infinity.) The difference between $`I_0`$ and $`I_s`$ is by Stokes formula
$$I_0I_s=_{B\times [0,s]}𝑑h^{}(\eta )$$
where $`h`$ is the homotopy map
$$h(y,\lambda )=(y,\theta (x,y)\lambda (\overline{x}\overline{y})).$$
Now,
$$d\eta =c_n(k/2\pi )^ne^{k\theta (xy)}ud\chi d\theta dy.$$
This equals 0 if $`|y|<1/2`$, and since $`\theta `$ is good we have the estimate
$$|dh^{}(\eta )|Ck^ne^{k((\delta /2+\lambda )|xy|^2\varphi (y)/2+\varphi (x)/2)}(1+\lambda )^n|u(y)|.$$
If $`|x|`$ is, say smaller than 1/4, $`|xy|1/4`$ when $`d\eta `$ is different from 0, so we get
$$|𝑑h^{}(\eta )|Ck^ne^{k(\varphi (x)/2\delta )}_{|y|>1/2}|u(y)|e^{k\varphi (y)/2}_0^s(1+\lambda )^ne^{k\lambda }𝑑\lambda $$
with a smaller $`\delta `$. By the Cauchy inequality the first integral in the right hand side is dominated by
$$u_{k\varphi }.$$
Since the last integral is bounded by a constant independent of $`k`$ we get the desired estimate. ∎
Thus we have a family of reproducing kernels mod $`e^{\delta k}`$. When $`\varphi =|y|^2`$ and $`\theta =\overline{y}`$, the kernel in the representation \- $`e^{\overline{x}y}`$ \- is also holomorphic in $`y`$ so we even have an asymptotic Bergman kernel mod $`e^{\delta k}`$. To achieve the same thing for general weights we need to introduce a bit more flexibility in the construction by allowing a more general class of amplitudes in the integral. We will replace the function $`\chi `$ in the integral by $`\chi (1+a)`$ for a suitably choosen function $`a`$, where $`a`$ has to be chosen to give an exponentially small contribution to the integral.
For this we consider differential forms
$$A=A(x,y,\theta ,k)=A_j(x,y,\theta ,k)\widehat{d\theta _j}$$
of bidegree $`(n1,0)`$. By $`\widehat{d\theta _j}`$ we mean the wedge product of all the differentials $`d\theta _i`$ except $`d\theta _j`$, with a sign chosen so that $`d\theta _j\widehat{d\theta _j}=d\theta `$. We assume that $`A`$ has an asymptotic expansion of order 0
$$AA_0+k^1A_1+\mathrm{}$$
By this we mean that for any $`N0`$
$$A\underset{0}{\overset{N}{}}A_mk^m=O(k^{N1})$$
uniformly as $`k`$ goes to infinity.
We assume also that the coefficients are holomorphic (in the smooth case almost holomorphic) for $`x`$, $`y`$ and $`\theta `$ of norm smaller than 2. Let
$$ad\theta =e^{k\theta (xy)}d_\theta e^{k\theta (xy)}A,$$
so that
(2.4)
$$a=D_\theta A+k(xy)A=:A,$$
where $`D_\theta =/\theta `$. We will say that a function $`a`$ arising in this way is a negligible amplitude. In the applications we will also need to consider finite order approximations to amplitude functions. Let
$$A^{(N)}=\underset{0}{\overset{N}{}}A_m/k^m$$
and similarily
$$a^{(N)}=\underset{0}{\overset{N}{}}a_m/k^m.$$
Then
$$a^{(N)}=A^{(N+1)}D_\theta A_{N+1}/k^{N+1},$$
so $`a^{(N)}`$ is a negligible amplitude modulo an error term which is $`O(1/k^{N+1})`$.
###### Proposition 2.2.
For any good contour $`\mathrm{\Lambda }`$ and any negligible amplitude $`a`$,
$$u(x)=c_n(k/2\pi )^n_\mathrm{\Lambda }e^{k\theta (xy)}u(y)\chi (y)(1+a)𝑑\theta dy+O(e^{k(\varphi (x)/2\delta )})u_{k\varphi },$$
for all $`x`$ in a sufficiently small neighbourhood of the origin if $`u`$ is an element of $`H_{k\varphi }(B)`$. Moreover
$$u(x)=c_n(k/2\pi )^n_\mathrm{\Lambda }e^{k\theta (xy)}u(y)\chi (y)(1+a^{(N)})𝑑\theta dy+O(e^{k\varphi (x)/2}/k^{N+1n})u_{k\varphi },$$
###### Proof.
For the first statement we need to verify that the contribution from $`a`$ is exponentially small as $`k`$ tends to infinity. But
$$_\mathrm{\Lambda }e^{k\theta (xy)}u(y)\chi (y)a𝑑\theta dy=_\mathrm{\Lambda }u(y)\chi (y)d_\theta (e^{k\theta (xy)}A)dy=$$
$$=_\mathrm{\Lambda }\chi d\left(u(y)e^{k\theta (xy)}Ady\right)=_\mathrm{\Lambda }𝑑\chi u(y)e^{k\theta (xy)}Ady.$$
Again, the last integrand vanishes for $`|y|<1/2`$ and is, since $`\mathrm{\Lambda }`$ is good, dominated by a constant times
$$|u(y)|e^{k(\delta |xy|^2\varphi (y)/2+\varphi (x)/2)}$$
The last integral is therefore smaller than
$$uO(e^{k(\varphi (x)/2\delta )})$$
so the first formula is proved. The second formula follows since by the remark immediately preceeding the proposition, $`a^{(N)}`$ is a good amplitude modulo an error of order $`1/k^{N+1}`$. ∎
The condition that an amplitude function $`a`$ can be written in the form (2.6) can be given in an equivalent very useful way. For this we will use the infinite order differential operator
$$Sa=\underset{0}{\overset{\mathrm{}}{}}\frac{1}{(k)^m(m!)}(D_\theta D_y)^m.$$
This is basically the classical operator that appears in the theory of pseudodifferential operators when we want to replace an amplitude $`a(x,y,\theta )`$ by an amplitude $`b(x,\theta )`$ independent of $`y`$, see . We let $`S`$ act on $`(n1)`$-forms $`A`$ as above componentwise. We say that $`Sa=b`$ for $`a`$ and $`b`$ admitting asymptotic expansions if all the coefficients of the powers $`(1/k)^m`$ in the expansion obtained by applying $`S`$ to $`a`$ formally equal the corresponding coefficients in the expansion of $`b`$. No convergence of any kind is implied. That $`Sa`$ equals $`b`$ to order $`N`$ means that the same thing holds for $`mN`$. Note also that since formally
$$S=e^{D_\theta D_y/k},$$
we have that
$$S^1=e^{D_\theta D_y/k}=\underset{0}{\overset{\mathrm{}}{}}\frac{1}{(k)^m(m!)}(D_\theta D_y)^m.$$
###### Lemma 2.3.
Let
$$a(1/k)^ma_m(x,y,\theta )$$
be given. Then there exists an $`A`$ satisfying (2.6) asymptotically if and only if
$$Sa|_{x=y}=0.$$
Moreover the last equation holds to order $`N`$ if and only if $`a^{(N)}`$ can be written
(2.5)
$$a^{(N)}=A^{(N+1)}+O(1/k^{N+1}).$$
###### Proof.
Note first that $`S`$ commutes with $`D_\theta `$ and that
$$S\left((xy)A\right)=(xy)SA(1/k)D_\theta SA.$$
Moreover
$$A=D_\theta A+k(xy)A,$$
so it follows that
(2.6)
$$SA=SD_\theta A+kS(xy)A=$$
$$D_\theta SA+k(xy)SAD_\theta SA=k(xy)SA.$$
Thus, if $`a`$ admits a representation $`a=A`$, then $`Sa`$ must vanish for $`x=y`$. Similarily, if
$$a^{(N)}=A^{(N+1)}+O(1/k^{N+1}).$$
it follows that $`Sa^{(N)}|_{y=x}=0`$ to order $`N`$.
Conversely, assume $`Sa|_{y=x}=0`$ . Then $`Sa=(xy)B`$ for some form $`B`$. But (2.6) implies that
$$S^1=kS^1(xy)$$
so
$$a=S^1\left((xy)B\right)=(1/k)S^1B$$
and (2.4) holds with $`A=1/kS^1B`$. If the equation $`Sa|_{y=x}=0`$ only holds to order $`N`$, then
$$Sa^{(N)}=(xy)B^{(N)}$$
to order $`N`$. Hence
$$a^{(N)}=S^1((xy)B^{(N)})=(1/k)S^1B^{(N)}=(1/k)(S^1B)^{(N)}$$
to order $`N`$, so (2.7) holds with $`A^{(N+1)}=1/k(S^1B)^{(N)}`$. ∎
### 2.2. The phase
Let us now see how to choose the contour $`\mathrm{\Lambda }`$ to get the phase function $`\psi (x,\overline{y})`$ appear in the expression
$$e^{\theta (xy)}.$$
In this section we still assume that the plurisubharmonic function $`\varphi `$ is real analytic and let $`\psi (x,y)`$ be the unique holomorphic function of $`2n`$ variables such that
$$\psi (x,\overline{x})=\varphi (x).$$
By looking at the Taylor expansions of $`\psi `$ and $`\varphi `$ one can verify that
(2.7)
$$2\mathrm{R}\mathrm{e}\psi (x,\overline{y})\varphi (x)\varphi (y)\delta |xy|^2$$
for $`x`$ and $`y`$ sufficiently small. Following an idea of Kuranishi, see , , we now find a holomorphic function of $`3n`$ variables, $`\theta (x,y,z)`$ that solves the division problem
(2.8)
$$\theta (xy)=\psi (x,z)\psi (y,z),$$
This can be done in many ways, but any choice of $`\theta `$ satisfies
$$\theta (x,x,z)=\psi _x(x,z).$$
To fix ideas, we take
$$\theta (x,y,z)=_0^1\psi (tx+(1t)y,z)dt$$
with $``$ denoting the differential of $`\psi `$ with respect to the first $`n`$ variables.
Since $`\theta (x,x,z)=\psi _x(x,z)`$ it follows that
$$\theta _z(0,0,0)=\psi _{xz}(0,0)=\varphi _{x\overline{x}}(0,0)$$
is a nonsingular matrix. Therefore
$$(x,y,z)(x,y,\theta )$$
defines a biholomorphic change of coordinates near the origin. After rescaling we may assume that $`\psi `$ is defined and satisfies (2.7) and that the above change of coordinates is well defined when $`|x|`$, $`|y|`$ and $`|z|`$ are all smaller than 2. We now define $`\mathrm{\Lambda }`$ by
$$\mathrm{\Lambda }=\{(y,\theta );z=\overline{y}\}.$$
Thus, on $`\mathrm{\Lambda }`$, $`\theta `$ is a holomorphic function of $`x`$, $`y`$ and $`\overline{y}`$. The point of this choice is that by (2.8), on $`\mathrm{\Lambda }`$,
$$\theta (xy)=\psi (x,\overline{y})\psi (y,\overline{y}).$$
Therefore we get the right phase function in our kernel and by (2.7)
$$2\mathrm{R}\mathrm{e}\theta (xy)=2\mathrm{R}\mathrm{e}\psi (x,\overline{y})2\varphi (y)\varphi (x)\varphi (y)\delta |xy|^2,$$
which means that $`\mathrm{\Lambda }`$ is a good contour in the sense of the previous section. By Proposition 2.2 we therefore get the following proposition, where we use the notation $`\beta `$ for the standard Kähler form in $`^n`$,
$$\beta =i/2dy_jd\overline{y}_j.$$
###### Proposition 2.4.
Suppose that $`u`$ is in $`H_{k\varphi }.`$ If $`a(x,y,\theta ,1/k)`$ is a negligible amplitude, we have
(2.9)
$$u(x)=$$
$$=(k/\pi )^n\chi _xe^{k(\psi (x,\overline{y})\psi (y,\overline{y}))}(det\theta _{\overline{y}})u(y)(1+a)\beta _n+$$
$$+O(e^{k(\varphi (x)/2+\delta )})u_{k\varphi },$$
with $`a=a(x,y,\theta (x,y,\overline{y})`$. Moreover
$$u(x)=$$
$$=(k/\pi )^n\chi _xe^{k(\psi (x,\overline{y})\psi (y,\overline{y}))}(det\theta _{\overline{y}})u(y)(1+a^{(N)})\beta _n+$$
$$+O(e^{k\varphi (x)/2}k^{nN1})u_{k\varphi }$$
### 2.3. The amplitude
In order to get an asymptotic Bergman kernel from (2.10) we need to choose the amplitude $`a`$ so that
$$det\theta _{\overline{y}}(1+a)=B(x,\overline{y})det\psi _{y\overline{y}},$$
with $`B`$ analytic. Polarizing in the $`y`$-variable, i e replacing $`\overline{y}`$ by $`z`$, this means that
$$(1+a(x,y,\theta (x,y,z,1/k))=B(x,z,1/k)det\psi _{yz}(y,z)/det\theta _z(x,y,z),$$
where $`B`$ is a analytic and independent of $`y`$. Consider this as an equation between functions of the variables $`x`$, $`y`$ and $`\theta `$. Let
$$\mathrm{\Delta }_0(x,y,\theta )=det\psi _{yz}(y,z)/det\theta _z(x,y,z)=det_\theta \psi _y.$$
Since $`\psi _y=\theta `$ when $`y=x`$ we have that $`\mathrm{\Delta }_0=1`$ for $`y=x`$. We need $`a`$ to be representable in the form (2.6) which by the previous lemma means that $`Sa=0`$ for $`y=x`$. Equivalently, $`S(1+a)=1`$ for $`y=x`$, so we must solve
(2.10)
$$S\left(B(x,z(x,y,\theta ),1/k)\mathrm{\Delta }_0(x,y,\theta )\right)=1$$
for $`y=x`$. This equation should hold in the sense of formal power series which means that the coefficient of $`1/k^0`$ must equal 1, whereas the coefficient of each power $`1/k^m`$ must vanish for $`m>0`$. In the computations $`x`$ is held fixed and $`z=z(y,\theta )`$. The first equation is
(2.11)
$$b_0(x,z(x,y,\theta ))\mathrm{\Delta }_0(x,x,\theta )=1.$$
This means that $`b_0(x,z(x,\theta ))=1`$, for all $`\theta `$ which implies that $`b_0`$ is identically equal to 1.
The second condition is
(2.12)
$$(D_\theta D_y)\left(b_0\mathrm{\Delta }_0\right)+b_1\mathrm{\Delta }_0=0$$
for $`y=x`$. Since we already know that $`b_0=1`$ this means that
$$b_1(z(x,\theta ))=(D_\theta D_y)\left(\mathrm{\Delta }_0\right)|_{y=x},$$
which again determines $`b_1`$ uniquely. Continuing in this way, using the recursive formula
(2.13)
$$\underset{0}{\overset{m}{}}\frac{(D_\theta D_y)^l}{l!}\left(b_{ml}\mathrm{\Delta }_0\right)|_{x=y}=0$$
for $`m>0`$ we can determine all the coefficients $`b_m`$, and hence $`a`$. Then $`Sa|_{y=x}=0`$ so $`Sa^{(N)}|_{y=x}=0`$ to order $`N`$, and the next proposition follows from Propositions 2.4 and 2.3.
###### Proposition 2.5.
Suppose that $`\varphi `$ is analytic. Then there are analytic functions $`b_m(x,z)`$ defined in a fixed neighbourhood of $`x`$ so that for each $`N`$
(2.14)
$$(k/\pi )^n(1+b_1(x,\overline{y})k^1+\mathrm{}+b_N(x,\overline{y})k^N)e^{k\psi (x,\overline{y})},$$
is an asymptotic Bergman kernel mod $`O(k^{(N+1)})`$.
### 2.4. Computing $`b_1`$
Let us first recall how to express some Riemannian curvature notions in Hermitian geometry. The Hermitian metric two-form $`\omega :=\frac{i}{2}H_{ij}dy^i\overline{dy^j}`$ determines a connection $`\eta `$ on the complex tangent bundle $`TX`$ with connection matrix (with respect to a holomorphic frame)
(2.15)
$$\eta =H^1H=:\eta _jdy_j.$$
The curvature is the matrix valued two-form $`\overline{}\eta `$ and the scalar curvature $`s`$ is $`\mathrm{\Lambda }\text{Tr}\overline{}\eta `$ where $`\mathrm{\Lambda }`$ is contraction with the metric form $`\omega .`$ Hence, in coordinates centered at $`x`$ where $`H(0)=I`$ the scalar curvature $`s`$ at $`0`$ is given by
(2.16)
$$s(0)=\text{Tr}(\frac{}{\overline{y_j}}\eta _j),$$
considering $`\eta `$ as matrix. We now turn to the computation of the coefficient $`b_1`$ in the expansion 2.14. By the definition of $`\theta `$ we have that
(2.17)
$$\theta _i(x,y,z)=\psi _{y_i}(y,z)+\frac{1}{2}\underset{k}{}(\frac{}{y_k}\psi _{y_i})(y,z)(x^ky^k)+\mathrm{}$$
Differentiating with respect to $`z`$ gives
$$\theta _z=H+\frac{1}{2}_yH(xy)+\mathrm{}.,$$
where $`H=H(y,z).`$ Multiplying both sides by $`H^1`$ and inverting the relation we get
(2.18)
$$\theta _z^1H=I\frac{1}{2}(H^1H)(xy)+\mathrm{}.,$$
Taking the determinant of both sides in formula 2.18 gives
(2.19)
$$\mathrm{\Delta }_0=1\frac{1}{2}\text{Tr}\eta (xy)+\mathrm{}.,$$
Hence, equation (2.14) now gives, since $`\frac{}{y}(xy)=1,`$ that
$$b_1(0,0)=(\frac{}{\theta }(\frac{1}{2}\text{Tr}\eta )_{x=y}=\frac{1}{2}\frac{}{\overline{y}}\text{Tr}\eta $$
showing that $`b_1(x,\overline{x})=\frac{s}{2},`$ according to 2.16.
### 2.5. Twisting with a vector bundle $`E`$
We here indicate how to extend the previous calculation to the case of sections with values in $`L^kE,`$ where $`E`$ is a holomorphic vector bundle with a hermitian metric $`G`$ (see also ). First observe that $`u(x)`$ is now, locally, a holomorphic vector and the Bergman kernel may be identified with a matrix $`K(x,\overline{y})`$ such that
$$u(x)=K(x,\overline{y})G(y,\overline{y})u(y)\psi _{y\overline{y}}e^{k\psi (y,\overline{y}))}𝑑\overline{y}dy.$$
To determine $`K`$ one now uses the ansatz
$$K(x,\overline{y})=c_n(k/2\pi )^ne^{k(\psi (x,\overline{y})}B(x,\overline{y},k^1)G(x,\overline{y})^1.$$
Then the condition on the amplitude function becomes
(2.20)
$$(1+a(x,y,\theta (x,y,z),1/k)det(\frac{\theta }{z}(x,y,z))=$$
$$=B(x,z,1/k)G(x,z)^1G(y,z)det(\psi _{yz}).$$
where $`a`$ now is a matrix valued form , i.e. $`\mathrm{\Delta }_0`$ in section 2.3 is replaced by the matrix $`\mathrm{\Delta }_G:=\mathrm{\Delta }_0G(x,z)^1G(y,z).`$ Note that
$$G(x,z)^1G(y,z)=IG^1(y,z)\frac{}{y}G(y,z)(xy)+\mathrm{}=:I\eta _E(y,z)(xy)+\mathrm{},$$
where $`\eta _E:=G^1\frac{}{y}G`$ is the connection matrix of $`E.`$ Hence, the equation 2.19 is replaced by
$$\mathrm{\Delta }_G=1\text{-(Tr}\eta /2I+\eta _E)(xy)+\mathrm{}$$
The same calculation as before then shows that the matrix $`b_1(0,0)`$ is given by
$$b_1(0,0)=\frac{1}{2}\frac{}{\overline{y}}Tr\eta I\frac{}{\overline{y}}\eta _E=\frac{s}{2}I+\mathrm{\Lambda }\mathrm{\Theta }_E,$$
where $`\mathrm{\Theta }_E:=\overline{}\eta _E`$ is the curvature matrix of $`E`$ and $`\mathrm{\Lambda }`$ denotes contraction with the metric two-form $`\omega .`$
Remark: Let $`K_k`$ be the Bergman kernel of $`H^0(X,L^k),`$ defined with respect a general volume form $`\mu _n.`$ Then the function $`G:=\mu _n/\omega _n`$ defines a hermitian metric on the trivial line bundle $`E`$ and the asymptotics of $`K_k`$ can then be obtained as above.
### 2.6. Smooth metrics.
Denote by $`\psi (y,z)`$ any almost holomorphic extension of $`\varphi `$ from $`\overline{\mathrm{\Delta }}=\{z=\overline{y}\}`$, i.e. an extension such that the anti-holomorphic derivatives vanish to infinite order on $`\overline{\mathrm{\Delta }}.`$ We may also assume that $`\overline{\psi (y,\overline{z})}=\psi (z,\overline{y}).`$ That $`\psi `$ is almost holomorphic means that for any multi index $`\alpha `$
(2.21)
$$D^\alpha (\overline{}\psi )=0,$$
(where $`D^\alpha `$ is the local real derivative of order $`\alpha )`$ when evaluated at a point in $`\overline{\mathrm{\Delta }},`$ i.e. when $`z=\overline{y}.`$ Let now
(2.22)
$$\theta =_0^1(_y\psi )(tx+(1t)y,z)𝑑t,\theta ^{}=_0^1(\overline{}_y\psi )(tx+(1t)y,z)𝑑t,$$
(where $`_y\psi `$ denotes the vector of partial holomorphic derivatives w.r.t the first argument of $`\psi )`$ so that
(2.23)
$$(xy)\theta +\overline{(xy)}\theta ^{}=\psi (x,z)\psi (y,z)$$
Then the smooth map corresponding to $`(x,y,z)(x,y,\theta )`$ is locally smoothly invertible for the same reason as in the analytic case, since $`\overline{}_z\theta =0`$ when $`x=y=\overline{z}.`$ Define the algebra $`𝒜`$ of all functions almost holomorphic when $`x=y=\overline{z}`$ as the set of smooth functions, $`f`$, of $`x`$ $`y`$ and $`z`$, such that
$$D^\alpha \overline{}f=0$$
for all multi indices $`\alpha ,`$ when $`x=y=\overline{z}`$. For a vector valued function we will say that it is in $`𝒜`$, if its components are in $`𝒜.`$ We also define the vanishing ideal $`^{\mathrm{}}`$ as the set of smooth functions $`f`$ such that
$$D^\alpha f=0,$$
for all multi indices $`\alpha ,`$ when $`x=y=\overline{z}`$. Hence, if $`f`$ belongs to $`𝒜`$ then (the coefficients of) $`\overline{}f`$ will belong to $`^{\mathrm{}}`$.
Note that $`\psi (tx+(1t)y,z)`$ is in $`𝒜`$ for each fixed $`t.`$ Hence $`\theta `$ is in $`𝒜`$ and $`\theta ^{}`$ is in $`^{\mathrm{}},`$ so that 2.23 gives
(2.24)
$$(xy)\theta =\psi (x,\overline{y})\psi (y,\overline{y})+O(\left|xy\right|^{\mathrm{}})$$
when $`\overline{z}=y.`$ The next simple lemma is used to show that the contribution of elements in $`^{\mathrm{}}`$ to the phase function and the amplitude is negligable.
###### Lemma 2.6.
Let $`f_i`$ be elements of the vanishing ideal $`^{\mathrm{}}`$ and let $`b(x,y)`$ be a local smooth function. Then
$$\chi _x(y)e^{k(\psi (x,\overline{y})\psi (y,\overline{y})+f_1(x,y,\overline{y}))}(b(x,y)+f_2(x,y,\overline{y}))u(y)\beta _n(y)=$$
$$=\chi _x(y)e^{k(\psi (x,\overline{y})\psi (y,\overline{y})}b(x,y)u(y)\beta _n(y)+O(k^{\mathrm{}})u_{k\varphi }$$
###### Proof.
First observe that if $`f_i`$ is an element of $`^{\mathrm{}},`$ then $`f_i(x,y,\overline{y})=O(|xy|^{\mathrm{}}.`$ Moreover, we have
(2.25)
$$\left|\right|xy|^{2N}e^{k(\psi (x,\overline{y})\varphi (x)/2\varphi (y)/2)}|Ck^{2N}(k|xy|)^2)^Ne^{k\delta \left|xy\right|^2}=O(k^{2N}),$$
where we have used (2.7). Combining this bound with the Cauchy-Schwartz inequality proves the lemma when $`f_1=0.`$ Now write
$$e^{k(\psi (x,\overline{y})\psi (y,\overline{y})+f_1(x,y,\overline{y}))}=e^{k(\psi (x,\overline{y})\psi (y,\overline{y})}+_0^1_t(e^{(\psi (x,\overline{y})\psi (y,\overline{y})+tf_1(x,y,\overline{y}))})dt.$$
By 2.25 the second term gives a contribution which is of the order $`O(k^{\mathrm{}}).`$ Hence the general case follows. ∎
.
###### Proposition 2.7.
Suppose that $`L`$ is smooth. Then there exists an asymptotic reproducing kernel *$`K_k^{(N)}`$ mod $`O(k^{nN1})`$* for $`H_{k\varphi }`$ , such that
(2.26)
$$K_k^{(N)}(x,\overline{y})=e^{k\psi (x,\overline{y})}(b_0+b_1k^1+\mathrm{}+b_Nk^N)$$
where $`b_i`$ is a polynomial in the derivatives $`_x^\alpha \overline{}_y^\beta \psi (x,\overline{y})`$ of the almost holomorphic extension $`\psi `$ of $`\varphi .`$ In particular,
(2.27)
$$e^{k(\varphi (x)/2+\varphi (y)/2)}(D_{x,y}^\alpha (\overline{}_x,_y))K(x,y)=O(k^{\mathrm{}})$$
uniformly in $`x`$ and $`y`$ for any given $`\alpha .`$
###### Proof.
We go through the steps in the proof of the analytic case and indicate the necessary modifications.
First we determine the coefficients $`b_m(x,z)`$ in the same way as in the analytic case, i e by fixing $`x`$ and solving
$$S(B(z)\mathrm{\Delta }_0)|_{y=x}=1$$
Here $`S`$ has the same meaning as before and in particular contains only derivatives with respect to $`\theta `$ and no derivatives with respect to $`\overline{\theta }`$. The difference is that $`\mathrm{\Delta }_0`$ is no longer analytic so $`B`$ will not be holomorphic, but it will still belong to $`𝒜`$ since $`\mathrm{\Delta }_0`$ does.
We next need to consider lemma 2.3 with $`a𝒜.`$ Then we get that
$$a𝒜,(Sa)_{y=x}=O(k^{N1})A𝒜:a=A+O(k^{N1})\text{mod}^{\mathrm{}}$$
Indeed, this follows from the argument in the analytic case and the fact that if $`c(=c(x,y,z))𝒜,`$ then
$$c(x,x,z)=0d𝒜:c=(xy)d\text{mod}^{\mathrm{}}$$
as can be seen by defining $`d`$ by
$$d=_0^1(_yc)(x,x+(1t)y,z)𝑑t.$$
Here $``$ also has the same meaning as before and contains only a derivative with respect to $`\theta `$ and no derivative with respect to $`\overline{\theta }`$. Then Proposition 2.2 holds as before except that there will be one extra contribution in the application of Stokes theorem coming from $`\overline{}_\theta A`$ (when $`z=\overline{y}).`$ Since $`\overline{}_\theta A`$ vanishes to infinite order when $`x=y=\overline{z}`$, it gives a contribution to the integral which is $`O(1/k^N)`$ for any $`N`$ by lemma 2.6.
We therefore get from Proposition 2.2 a reproducing kernel of the form claimed in (2.26) except that the phase function equals
$$k\theta (xy)=k(\psi (x,\overline{y})\varphi (y))+f)$$
with $`f`$ in $`𝒜`$. Again by lemma 2.6 we may remove $`f`$ at the expense of adding a contribution to the to the integral which is negligable, i.e which is $`O(1/k^N)`$ for any $`N`$. ∎
## 3. The global Bergman kernel
In this section we will show that, if the curvature of $`L`$ is positive everywhere on $`X,`$ then the global Bergman kernel $`K_k`$ of $`H^0(X,L^k)`$ is asymptotically equal to the local Bergman kernel $`K_k^{(N)}`$ of $`H_{k\varphi }`$ (constructed in section 2).
Recall (section 1) that the Bergman kernel $`K`$ associated to $`L`$ is a section of $`\overline{L}L`$ over $`X\times X.`$ By restriction $`K_x`$ is identified with a holomorphic section of $`\overline{L_x}L,`$ where $`L_x`$ is the fiber of $`L`$ over $`x.`$ Given any two vector spaces $`E`$ and $`F,`$ the scalar product on $`L`$ extends uniquely to a pairing
(3.1)
$$(,):LE\times LFEF,$$
linear over $`E`$ and anti-linear over $`F.`$ In terms of this pairing $`K_y`$ has the global reproducing property
(3.2)
$$\alpha (y)=(\alpha ,K_y)$$
for any element $`\alpha `$ of $`H^0(X,L).`$ By taking $`\alpha =K_x`$ (so that $`E=L_x`$ and $`F=\overline{L_y}`$ in 3.1) one gets
(3.3)
$$K(y,x):=K_x(y)=(K_x,K_y).$$
This also implies that $`\overline{K(x,y)}=K(y,x)`$ and that
(3.4)
$$K(x,x)=(K_x,K_x)=K_x^2.$$
$`K(x,x)`$ is a section of $`\overline{L}L`$. Its norm as a section to this bundle is the Bergman function, which in a local frame with respect to which the metric on $`L`$ is given by $`e^\varphi `$ equals
$$B(x)=K(x,x)e^{\varphi (x)}.$$
Notice also that by the Cauchy inequality we have an extremal characterization of the Bergman function:
$$B(x)=sup|s(x)|^2$$
where the supremum is taken over all holomorphic sections to $`L`$ of norm not greater than 1.
We now denote by $`K_k`$ the Bergman kernel associated to $`L^k`$, and write $`B_k`$ for the associated Bergman function. It follows from the extremal characterization of the Bergman function and the submeanvalue inequality for a holomorphic section $`s`$ over a small ball with radius roughly $`1/k^{1/2}`$ that
$$B_kCk^n,$$
uniformly on $`X`$ (see e g ).
Let now $`K_x^{(N)}(y)`$ be the local Bergman kernel of propositions 2.5-2.6, where the coefficients $`b_m`$ are given by (2.15),
(3.5)
$$\overline{K_x^{(N)}(y)}=(k/\pi )^n(1+b_1(x,\overline{y})k^1+\mathrm{}+b_N(x,\overline{y})k^N)e^{k\psi (x,\overline{y})}.$$
By construction, the coefficients $`b_m(x,z)`$ are holomorphic if the metric on $`L`$ \- locally represented by $`\varphi `$ \- is real analytic. In case $`\varphi `$ is only smooth the $`b_m`$s are almost holomorphic, meaning that
$$\overline{}_{xz}b_m$$
vanishes to infinite order when $`z=\overline{x}`$.
Replacing $`K_y`$ in the relation 3.3 with the local Bergman kernel $`K_k^{(N)}`$ will now show that $`K_k=K_k^{(N)}`$ up to a small error term.
###### Theorem 3.1.
Assume that the smooth line bundle $`L`$ is globally positive. Let $`K_k^{(N)}`$ be defined by (3.5), where the coefficients $`b_m`$ are determined by the recursion (2.15).
If the distance $`d(x,y)`$ is sufficiently small, then
(3.6)
$$K_k(x,y)=K_k^{(N)}(x,y)+O(k^{nN1})e^{k(\varphi (x)/2+\varphi (y)/2)}\text{,}$$
Moreover,
$$D^\alpha (K_k(x,y)K_k^{(N)}(x,y))=O(k^{m+nN1})e^{k(\varphi (x)/2+\varphi (y)/2)}$$
if $`D^\alpha `$ is any differential operator with respect to $`x`$ and $`y`$ of order at most $`m`$.
###### Proof.
Let us first show that
(3.7)
$$K_k(y,x)=(\chi K_{k,x},K_{k,y}^{(N)})+O(k^{nN1})e^{k(\varphi (x)/2+\varphi (y)/2)}$$
where $`\chi `$ is a cut-off function equal to 1 in a neighbourhood of $`x`$ which is large enough to contain $`y`$. Fixing $`x`$ and applying Proposition 2.5 to $`u_k=K_{k,x}`$ gives 3.7 with the error term
$$e^{\varphi (y)/2}O(k^{N1})K_x.$$
Now, by 3.4 and the estimate for $`B_k`$
$$K_{k,x}^2=B_k(x)e^{k\varphi (x)}Ck^ne^{k\varphi (x)},$$
This proves 3.7 with uniform convergence.
Next we estimate the difference
$$u_{k,y}(x):=K_{k,y}^{(N)}(x)(\chi K_y^{(N)},K_{k,x}).$$
Since the scalar product in this expression is the Bergman projection,
$$P_k(\chi K_{k,y}^{(N)})(x),$$
$`u_{k,y}`$ is the $`L^2`$-minimal solution to the $`\overline{}`$-equation
$$\overline{}u_{k,y}=\overline{}(\chi K_{k,y}^{(N)}).$$
The right hand side equals
$$(\overline{}\chi )K_{k,y}^{(N)}+\chi \overline{}K_{k,y}^{(N)}.$$
Since $`\chi `$ equals 1 near $`y`$ it follows from (2.9) and the explicit form of $`K_{k,y}^{(N)}`$ that the first term is dominated by
$$e^{\delta k}e^{k(\varphi ()/2+\varphi (y)/2)}.$$
The second term vanishes identically in the analytic case. In the smooth case $`\overline{}K_k^{(N)}`$ can by proposition 2.7 be estimated by
(3.8)
$$O(1/k^{\mathrm{}})e^{k(\varphi ()/2+\varphi (y)/2)}.$$
Altogether $`\overline{}u_{k,y}`$ is therefore bounded by (3.7), so by the Hörmander $`L^2`$-estimate we get that
$$u_{k,y}^2O(1/k^{\mathrm{}})e^{k\varphi (y)/2}.$$
But, since the estimate on $`\overline{}u_{k,y}`$ is even uniform, we get by a standard argument involving the Cauchy integral formula in a ball around $`x`$ of radius roughly $`1/k^{1/2}`$ that $`u_{k,y}`$ satisfies a pointwise estimate
$$|u_{k,y}(x)|^2O(1/k^{\mathrm{}})e^{k(\varphi (y)/2+\varphi (x)/2)}.$$
Combining this estimate for $`u_{k,y}(x)`$ with (3.6) we finally get
(3.9)
$$|K_{k,y}^{(N)}(x)\overline{K_k(y,x)}|e^{k\varphi (x)/2k\varphi (y)/2}O(1/k^{N+1}).$$
Since $`K_k`$ is hermitian ( i e $`K_k(x,y)=\overline{K_k(y,x)}`$) this proves the proposition except for the statement on convergence of derivatives.
In the analytic case the convergence of derivatives is, by the Cauchy estimates, an automatic consequence of the uniform convergence, since the kernels are holomorphic in $`x`$ and $`\overline{y}`$. In the smooth case, we have that
$$\overline{}K_k^{(N)}(x,\overline{z})=O(1/k^{\mathrm{}})e^{k(\varphi ()/2+\varphi (y)/2)}.$$
This implies that the Cauchy estimates still hold for the difference between $`K_k`$ and $`K_k^{(N)}`$, up to an error which is $`O(1/k^{\mathrm{}})`$, and so we get the convergence of derivatives even in the smooth case. ∎
###### Remark 3.2.
The proof above actually shows that the asymptotic expansion for the global Bergman kernel $`K_k(x,y)`$ holds close to any point $`x`$ which is in $`X(0)`$ (the open subset of $`X`$ where the curvature form of $`\varphi `$ is positive) and is such that $`x`$ satisfies the following global condition: for any given $`\overline{}`$closed $`(0,1)`$form $`g_k`$ with values in $`L^k`$ supported in some fixed neighbourhood of $`x`$ we may find sections $`u_k`$ with values in $`L^k`$ such that
(3.10)
$$\overline{}u_k=g_k,$$
on $`X`$ and
(3.11)
$$u_k_{k\varphi }Cg_k_{k\varphi }$$
After this paper was written this observation was used in to obtain an asymptotitc expansion of $`K_k(x,y)`$ on a certain subset of $`X(0)`$ for any Hermitian line bundle $`(L,\varphi )`$ over a projective manifold $`X.`$
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# Photo-disintegration cross section measurements on 186W, 187Re and 188Os: implications to the Re-Os cosmochronology
## I Introduction
Owing to the long half-life of <sup>187</sup>Re, the <sup>187</sup>Re-<sup>187</sup>Os pair may serve as a cosmochronometer to measure the duration of stellar nucleosynthesis that precedes the solidification of the solar system Clay64 . By adding the age of the solar system ($``$4.6 Gyr), it provides the age of the Galaxy. The facts that both <sup>186</sup>Os and <sup>187</sup>Os are produced only by the s-process nucleosynthesis apart from the cosmoradiogenic yield of <sup>187</sup>Os and that the isotopic solar abundance ratio of <sup>186</sup>Os and <sup>187</sup>Os AnGr89 is available make this chronometer potentially reliable in the sense that it is independent of r-process models. The quantitative interpretation is, however, complicated by the possible enhancement of <sup>187</sup>Re-<sup>187</sup>Os transmutation rates in stellar condition, the stellar production and destruction of <sup>187</sup>Re and <sup>187</sup>Os during the chemical evolution of the Galaxy, the possible existence of s-process branchings at <sup>185</sup>W and <sup>186</sup>Re, and the neutron capture by the 9.75 keV first excited state in <sup>187</sup>Os Yoko83 ; Arno84 ; Taka03 ; Woos79 .
The last issue on the effect of neutron capture on the 9.75 keV state in <sup>187</sup>Os was raised in Ref. Woos79 . In the local approximation, the ratio of the s-process yields of <sup>186</sup>Os and <sup>187</sup>Os, $`N_s(^{186}\mathrm{Os})/N_s(^{187}\mathrm{Os})`$, can be expressed as
$$\frac{N_s(^{187}\mathrm{Os})}{N_s(^{186}\mathrm{Os})}F_\sigma \frac{\sigma ({}_{}{}^{186}\mathrm{Os})}{\sigma ({}_{}{}^{187}\mathrm{Os})}.$$
(1)
Here, $`\sigma ({}_{}{}^{186}\mathrm{Os})`$ and $`\sigma ({}_{}{}^{187}\mathrm{Os})`$ are the Maxwellian-averaged neutron capture cross sections on <sup>186</sup>Os and <sup>187</sup>Os in the ground states, respectively. The $`F_\sigma `$ value accounts for the correction to the cross section due to the neutron capture on the 9.75 keV state in <sup>187</sup>Os which is substantially populated at typical s-process temperatures $`T`$1-3$`\times 10^8`$ K. It is defined by
$$F_\sigma =\frac{\sigma ({}_{}{}^{187}\mathrm{Os})}{\sigma ^{}({}_{}{}^{187}\mathrm{Os})}$$
(2)
with $`\sigma ^{}(^{187}\mathrm{Os})`$ being the Maxwellian-averaged neutron capture cross section on <sup>187</sup>Os at a given stellar temperature. Here, the first excited state in <sup>186</sup>Os and the second excited state in <sup>187</sup>Os lie at higher excitation energies of 137 keV and 74 keV, respectively, and therefore their contributions to neutron capture in the stellar condition may safely be ignored. In 1970s, it was of critical concern whether or not the $`F_\sigma `$ value exceeds unity, because it has a great impact on the age of the Galaxy; the larger the $`F_\sigma `$, the smaller the age. However, there was a large spread in the early estimate (0.80-1.10 Woos79 , 0.8 Holm76 , and $``$1.5 Fowl73 ).
This concern, combined with the fact that a direct measurement of neutron capture on the 9.75 keV state is virtually impossible, led to measurements of neutron capture on <sup>186,187,188</sup>Os Wint80 ; Wint82 and neutron inelastic scattering to the 9.75 keV state in <sup>187</sup>Os Hers83 ; Mack83 . The measured neutron capture cross sections were in good agreement with those of the earlier measurements BrLS76 ; Brow81 . But, the statistical analysis of the capture data gave a lower bound of 0.30 $`b`$ to the inelastic scattering cross section $`\sigma _{nn^{}}`$ at a neutron energy of 30 keV, and by combining with an upper limit of $`\sigma _{nn^{}}=0.5`$ $`b`$ Wint74 , gave $`F_\sigma 1`$. In contrast, the two neutron inelastic scattering data, $`\sigma _{nn^{}}=1.13\pm 0.2`$ $`b`$ at 60 keV Hers83 and $`1.5\pm 0.2`$ $`b`$ at 34 keV Mack83 , are consistent with $`F_\sigma =`$0.80-0.83 and 0.80, respectively, within the statistical models.
More recent efforts have been made toward a unified statistical model analysis of all available data including measurements of ($`n,\gamma `$) cross sections WiMH87 and elastic/inelastic scattering cross sections McEl89 on a neighboring nucleus <sup>189</sup>Os in which the ground state with $`J^\pi `$ = 3/2<sup>-</sup> and the first excited state with 1/2<sup>-</sup> at 36 keV appear in the reverse order of the corresponding states in <sup>187</sup>Os. These analyses showed that $`F_\sigma `$=0.79-0.83.
In the Hauser-Feshbach statistical model calculations, however, large uncertainties may arise from the $`\gamma `$-ray transmission coefficients rather than the neutron optical potential and the level density. Information on the $`\gamma `$-ray transmission coefficients for neutron captured states in the low-energy tail of the giant dipole resonance in <sup>188</sup>Os can be obtained in the inverse photo-disintegration of <sup>188</sup>Os. But, the conventional Lorentzian model based on the previous photo-disintegration data on <sup>188</sup>Os Ber69 may not be satisfactory for two reasons. One reason is that the data taken with $`\gamma `$-ray beams from the positron annihilation in flight exhibit non-vanishing cross sections even below the neutron threshold (see sec. III.2). The non-vanishing cross sections may be attributed to contributions from the positron bremsstrahlung. The other is that microscopic models can predict the $`E1`$ $`\gamma `$ strength function more reliably than the Lorentzian model.
Besides $`F_\sigma `$, the effect of the s-process branchings at <sup>185</sup>W and/or <sup>186</sup>Re was parameterized as $`F_b`$ and investigated within the framework of the schematic s-process models Arno84 . More recently, neutron capture cross sections were measured for neighboring <sup>185</sup>Re and <sup>187</sup>Re nuclei to derive the statistical model parameters from a consistent systematics Kaep91 . With the improved parameters for s-process analysis, a stellar model calculation for low-mass AGB stars showed that the local approximation was disturbed by the branchings at <sup>185</sup>W and <sup>186</sup>Re. However, the precise physical conditions of the AGB model need to be scrutinized before any definite conclusion on $`F_b`$ is drawn. This is clearly beyond the scope of the present study. Instead, the relevant statistical parameters, particularly the $`E1`$ $`\gamma `$ strength function, can further be improved by the photo-disintegration measurement on <sup>186</sup>W.
In the present study, we have measured photoneutron cross sections of <sup>186</sup>W, <sup>187</sup>Re and <sup>188</sup>Os using tunable quasi-monochromatic $`\gamma `$-ray beams from laser Compton scattering (LCS). The photo-disintegration data allow us to constrain $`F_\sigma `$ values within the Hauser-Feshbach statistical model and discuss their implications to the Re-Os cosmochronology. Part of the present data has already been published Mohr04 .
## II Experimental procedure
Photoneutron cross section measurements on <sup>186</sup>W, <sup>187</sup>Re and <sup>188</sup>Os were performed at the National Institute of Advanced Industrial Science and Technology (AIST). Tunable quasi-monochromatic photon beams were generated by Compton scattering of laser photons, with relativistic electrons circulating in the storage ring TERAS Ohg91 . A Nd:YLF Q-switch laser at a wavelength of 527 nm in second harmonics was operated at a frequency of 2 kHz. The electron energy was varied in the range from 450 to 588 MeV to produce LCS photons with the average energy from 7.3 and 10.9 MeV. A 20 cm lead collimator with a small hole of 2 mm in diameter was placed at approximately 6 m downstream from the interaction area which defines a scattering cone of the LCS photons. The typical energy resolution was 10 % in FWHM. Further details on the experimental setup can be found in Ref. Utsu03 .
Figure 1 shows an energy spectrum of the LCS photons measured by an HP-Ge detector with a relative efficiency of 120 %. The energy was calibrated at 1460.8 and 2614.5 keV with <sup>40</sup>K and <sup>208</sup>Tl radioactive isotopes of natural origin. A Monte Carlo simulation was performed with the EGS4 code Nel85 to analyze the response of the Ge detector. The energy distribution of incident LCS photons (dotted line in Fig. 1) was determined so as to reproduce the observed response function (solid line in Fig. 1). Photons with energies higher than the neutron threshold ($`S_n`$=7.19 MeV for <sup>186</sup>W, $`S_n`$=7.36 MeV for <sup>187</sup>Re and $`S_n`$=7.99 MeV for <sup>188</sup>Os) are responsible for the ($`\gamma `$,$`n`$) reactions. The fraction of these photons in the total photon flux and the average photon energy were obtained from the original LCS photon spectrum.
The beam current of the electron storage ring decreases exponentially in a normal condition with a lifetime $``$ 6 hours. In the EGS4 Monte Carlo simulation, it was found that the electron beam size in the region of the interaction with laser photons varied with time: for example, from 2.2 mm in diameter at 166 mA to 1.2 mm at 72 mA in the <sup>187</sup>Re measurement. A space-charge effect is considered to be a main cause for the decrease in the beam size. This beam size effect, which was evident in long runs near the neutron thresholds, introduced uncertainties in the fraction of the LCS photon beam above the threshold. The resultant uncertainty was estimated to be 1 - 6 % in the present experiment. On the other hand, the average photon energy was determined well within 40 keV.
The number of LCS photons was monitored during the experiment using a large volume (8”$`\times `$12”) NaI(Tl) scintillation detector placed behind targets. A typical pile-up spectrum is shown in Fig. 2. The pulse height of the spectrum is proportional to the number of LCS photons per beam pulse. The photon flux was determined with 3 % uncertainty based on a statistical analysis on the pile-up spectrum Toyo00 .
Metallic powders of 1246 mg <sup>186</sup>W, 996 mg <sup>187</sup>Re, and 693 mg <sup>188</sup>Os enriched to 99.79, 99.52 % and 94.99 %, respectively, were pressed to self-supporting tablets with a diameter of 8 mm. The <sup>188</sup>Os powder included major contaminants of <sup>189</sup>Os (2.55 %), <sup>190</sup>Os (1.27 %) and <sup>192</sup>Os (0.97 %). These tablets were mounted inside thin containers made of aluminum, and were irradiated with the LCS photon beams. The threshold energy of the $`(\gamma ,n)`$ reaction on <sup>27</sup>Al is 13.06 MeV, which is higher than those for <sup>186</sup>W, <sup>187</sup>Re and <sup>188</sup>Os. The present photoneutron cross section measurements were performed at energies below the threshold energy and therefore undisturbed by the <sup>27</sup>Al$`(\gamma ,n)`$ reaction. Further, measurements with an empty aluminum container (blank target) showed that no background neutrons were produced from photo-disintegration of possible impurities in the aluminum.
Emitted neutrons were detected by sixteen <sup>3</sup>He proportional counter (EURISYS MESURES 96NH45) embedded into a polyethylene moderator. Two sets of eight counters were placed in double concentric (inner and outer) rings at 7 and 10 cm from the beam axis. Time correlations (Fig. 3) between the neutron signal and the laser pulse were measured to estimate the number of background neutrons that arrived randomly at the <sup>3</sup>He detectors. These background neutrons were most likely produced by bremsstrahlung arising from collisions of electrons with residual gaseous molecules in the storage ring. In the moderation time distribution, constant events above 400 $`\mu `$s and at small correlation times were taken to be background neutrons. The constant background was further confirmed by using a 1 kHz laser and a wider (1 ms) time range (see, for example, Fig. 3 of Ref. Hara03 ). The background subtraction is included in the statistical uncertainties through the error propagation.
The neutron detection efficiency was measured at the average neutron energy 2.14 MeV with a standard <sup>252</sup>Cf source. The dependence of the efficiency on neutron energy was determined by a Monte Carlo MCNP simulation with the statistical accuracy less than 0.5 % Hara03 . The so-called ring ratio between the count rates of inner and outer rings was used to determine the average energy of emitted neutrons Utsu03 . A polynomial fit to the energy dependence of the ring ratios was made as in Utsu03 , where an asymptotic value of the ring ratio (8.0) at 10 keV simulated for the present neutron detector Hara03 was used as a constraint. The ring ratio varied between 7.4 and 2.9, indicating the neutron energies of several tens to a few hundreds keV for measurements close above the thresholds and up to 0.82 MeV at higher photon energies. The uncertainty in the neutron energy thus determined was estimated to be 10-15 keV, resulting in the uncertainty less than 0.8 % in the total neutron detection efficiency. Note that the total efficiency is nearly constant (44.4-44.3 %) over the neutron energy from 1 MeV to 400 keV and that it slowly decreases to 39.5 % at 50 keV. The overall systematic uncertainty for cross sections was estimated to be 5.9-8.3 %, which was determined by the neutron emission rate of the calibration source (5 %), the number of the incident LCS photons (3 %), and the beam size effect.
## III Results
### III.1 Data reduction
The photoneutron cross section measured with a monochromatic photon beam is given by
$$\sigma =\frac{n_n}{N_\mathrm{t}N_\gamma fϵ_n(E_n)}$$
(3)
where $`n_n`$ is the number of neutrons detected with the <sup>3</sup>He counters, $`N_\mathrm{t}`$ is the number of target nuclei per unit area, $`N_\gamma `$ is the number of incident photons, $`f`$ is the correction factor for a thick-target measurement, and $`ϵ(E_n)`$ is the neutron detection efficiency. The correction factor is given by $`f=(1e^{\mu d})/(\mu d)`$ with the linear attenuation coefficient of photons, $`\mu `$, and the target thickness, $`d`$. The attenuation coefficient was taken by interpolation from Jaeger68 for the average energy of the LCS photons. The energy spread of the LCS photon beam in the full-width at half maximum makes negligible contributions ($`0.2`$ %) to the determination of the correction factor $`f`$, which deviates from unity by no more than 6 % in the present measurements.
Recently, a methodology was developed to determine cross sections for reactions induced by a quasi-monochromatic photon beam Mohr04 . When the photon beam has an energy distribution of $`n_\gamma (E)`$, $`N_\gamma \sigma `$ in Eq. (3) has to be replaced by the integral $`n_\gamma (E)\sigma (E)𝑑E`$. By expanding the cross section $`\sigma (E)`$ in the Taylor series at the average photon energy $`E_0`$, the first term in the Taylor series $`\sigma (E_0)`$ (the cross section at the average energy) was numerically evaluated along with the higher-order terms. The new methodology determines $`\sigma (E_0)`$ in the energy region of astrophysical importance near threshold within 6 % corrections from the monochromatic approximation (Eq. (3)). Photoneutron cross sections presented in this paper are based on this methodology.
### III.2 Photoneutron cross sections
Photoneutron cross sections measured for <sup>186</sup>W, <sup>187</sup>Re and <sup>188</sup>Os are shown in Fig. 4 and Table 1. The contributions from the reactions with the main contaminants (<sup>189</sup>Os, <sup>190</sup>Os and <sup>192</sup>Os) of the <sup>188</sup>Os target were estimated from the previous data from Ref. Ber79 . The error bars in Fig. 4 include both the statistical and systematic uncertainties. Previously, cross section data were taken for <sup>186</sup>W Ber69 and <sup>188</sup>Os Ber79 with quasi-monochromatic photons from positron annihilation in flight, as shown in Fig. 4 for comparison. In addition, data for <sup>186</sup>W and <sup>187</sup>Re were taken with bremsstrahlung Sonn03 ; Mull04 ; gory73 . In Ref. gory73 , yield curves obtained in small increments of the electron beam energy with 1-MeV spacing were converted to cross sections through an unfolding procedure based on the Penfold-Leiss method Penf59 . As mentioned in Ref. gory73 , it is a known fact Tikh63 that cross sections are not obtained correctly because of swings of the solution, though uncertainties resulting from the swings may be reduced significantly.
The present <sup>188</sup>Os data are in reasonable agreement with the previous data Ber79 except at the low energy where the previous data exhibit non-vanishing cross sections even below the neutron threshold of 7.99 MeV. The non-vanishing cross sections may be attributed to leftover in the subtraction of contributions of the positron bremsstrahlung admixed with the positron annihilation photons.
The ($`\gamma ,n`$) cross section exhibits the threshold behavior Wign48 ; Brei58 of
$$\sigma (E)=\sigma _0\left(\frac{E_\gamma S_n}{S_n}\right)^p.$$
(4)
Here, $`p`$ is related to the neutron orbital angular momentum $`\mathrm{}`$ through $`p=\mathrm{}+1/2`$. The values of $`p=0.5`$ and 1.5 are expected for the s- and p-wave neutron decays, respectively. The best fits to the present experimental data gave $`p=0.47`$ and $`\sigma _0=78`$ mb for the <sup>186</sup>W($`\gamma ,n`$) reaction, $`p=0.67`$ and $`\sigma _0=117`$ mb for the <sup>187</sup>Re($`\gamma ,n`$) reaction, and $`p=0.53`$ and $`\sigma _0=143`$ mb for the <sup>188</sup>Os($`\gamma ,n`$) reaction. The results show a rather pure s-wave character for the <sup>186</sup>W and <sup>188</sup>Os($`\gamma ,n`$) reactions, and suggest an admixture of p- and/or d-wave neutron emissions following the $`E1`$ excitation of <sup>187</sup>Re. The cross-section parametrization (Eq. (4)) was also made in the bremsstrahlung measurements Sonn03 ; Mull04 ; the results are in rough agreement with our new experimental data.
## IV Comparison with theory
### IV.1 Theoretical framework
The cross sections measured in the present work are now compared with the predictions of the Hauser-Feshbach (HF) compound nucleus theory hau52 ; Holm76 . The uncertainties involved in HF cross section calculation are known not to be related to the theory of compound nucleus emission itself, but rather to the uncertainties associated with the evaluation of the nuclear properties entering the calculation of the transmission coefficients. It is therefore of prime importance to compare the effects of different nuclear inputs to estimate the reliability and accuracy of the predictions, especially when considering the reverse rates, i.e, in the present case, the radiative neutron capture rate which might be sensitive to different input parameters than the rate measured. In the present work, the nuclear level densities are derived from two models, either the widely used back-shifted Fermi gas (BSFG) model based on the global parametrization of Gori02 or the microscopic calculations taking into account the discrete structure of the single-particle spectra associated with Hartree-Fock+BCS (HFBCS) potentials dem01 . This model has the advantage of treating shell, pairing and deformation effects consistently, and for practical applications, has been renormalized on existing experimental information (low-lying levels and s-wave neutron resonance spacings whenever available as in the cases considered here). The transmission coefficients for particle emission is calculated either with the so-called JLMB semi-microscopic potential of bdg01 derived from the Brückner–Hartree–Fock approximation based on a Reid’s hard core nucleon–nucleon interaction, or with the global phenomenological mass- and energy-dependent potential of Woods-Saxon type developed by koning03 .
The photon transmission function of particular interest in photoemission data is calculated assuming the dominance of dipole transitions in the photon channel. The electric- and magnetic-dipole (GDR) transition strength functions are usually described by a Lorentz-type function where the energies and widths are determined by experimental data, whenever they exist, or by appropriate parametrizations. However, the calculation of the radiative capture or photoabsorption at low energies (and particularly in stellar conditions where the excited states of the target nucleus are thermally populated) is particularly sensitive to the low-energy tail of the GDR of the compound system. The shape of the GDR is expressed most frequently by a generalized energy-dependent-width Lorentzian function adjusted on low-energy data ko90 . To test such models, we consider here the Hybrid model go98 which couples the GDR Lorentzian description at high energies with an analytical approximation to the theory of finite Fermi systems at energies below the neutron separation energy ka83 .
In addition to the Hybrid model go98 , the Quasi-Particle Random Phase Approximation (QRPA) model of kh01 is also considered here for estimating the photon transmission coefficients. These QRPA calculations are self-consistently built on a ground state derived with the HFBCS approximation. The final $`E1`$-strength functions is obtained by folding the QRPA strength with a Lorentzian function to account for the damping of collective motions and the deformation effects. This global calculation based on the SLy4 Skyrme interaction has been shown to reproduce relatively well photoabsorption and average resonance capture data at low energies Utsu03 ; kh01 . Note that both the Hybrid and QRPA models differ not only in the predictions of the position and width of the GDR, but also in the energy dependence of its tail, which is an important quantity to derive the reaction rate. Since we are here mainly concerned with the GDR tail at low energies and the prediction of the reverse neutron capture cross section, both $`E1`$-strength functions are renormalized on the available experimental information on the position of the GDR peak and the corresponding maximum absorption cross section. In the case of the HFBCS+QRPA model, this adjustment is achieved within the folding procedure introduced to account for damping and deformation effects.
### IV.2 Comparison between experimental and theoretical rates
The final nuclear inputs considered in the present analysis are summarized in Table 2. Four different sets are used to estimate the photoemission cross section as well as the reverse radiative neutron capture cross sections. The comparison between these 4 sets allow us to estimate the sensitivity of the cross sections to the various input quantities, but also the uncertainties affecting the final prediction of the neutron capture rate of astrophysical interest. In Fig. 5, our new experimental data are compared with the theoretical photoemission cross sections. Also shown are previous measurements obtained in the vicinity of the GDR peak energy Ber79 ; gory73 . Most of the calculations agree relatively well with experimental data, down to energies close to the neutron threshold. The neutron optical potential and nuclear level densities influence the photoneutron cross section only in a small energy range of no more than 1 MeV above the neutron threshold, so that the global behavior of cross section is almost entirely dictated by the $`E1`$-strength. Some specific comments can be made for each reaction:
* In the <sup>186</sup>W($`\gamma ,n`$)<sup>185</sup>W case, the HFBCS+QRPA model predicts some extra strength at energies around 7.5-10 MeV with respect to the Hybrid model. This extra strength is clearly seen experimentally below 8 MeV but not above. However, all models overestimate the 7mb cross section at $`E=7.26`$ MeV.
* In the case of the <sup>187</sup>Re($`\gamma ,n`$)<sup>186</sup>Re reaction, the low-energy data can only be reproduced when adopting the BSFG model of nuclear level densities. The microscopic HFBCS-based model fails to describe the fast rise of the cross section at the neutron threshold. Around 11 MeV, the Hybrid and HFBCS+QRPA strength predict a relatively different cross section, the former one being compatible with the Goryachev et al. gory73 and the later with our more accurate measurements.
* As far as <sup>188</sup>Os($`\gamma ,n`$)<sup>187</sup>Os is concerned, all HF calculations reproduce relatively well the data, though the Hybrid model gives a lower cross section in the 8.5-10 MeV energy range.
### IV.3 Determination of the neutron capture cross sections
We now estimate the stellar Maxwellian-averaged neutron capture cross section $`\sigma ^{}`$ of astrophysics interest on the basis of the calculations presented above, i.e constrained by the reverse photo-disintegration rate compatible with the new measurements. It should be recalled that the neutron capture cross section at energies of a few tens of keV is mainly sensitive to photon transmission coefficient at an energy close to and even below the neutron separation energy. For this reason, only the calculations that reproduce the photoemission rate at the neutron separation energy relatively well are retained. This leads us to reject the calculation ‘INP-1’ for the <sup>187</sup>Re($`\gamma ,n`$)<sup>186</sup>Re reaction.
In the particular case of the stable <sup>187</sup>Os target, direct experimental data are available for the <sup>187</sup>Os($`n,\gamma `$)<sup>188</sup>Os reaction Wint80 ; Brow81 and used as additional constraints on the nuclear ingredients, namely the combination of the nuclear level density and optical potential. The resulting cross sections obtained with the 4 sets of nuclear inputs given in Table 2 are shown in Fig. 6 and seen to agree relatively well with experimental data, except in the specific ‘INP-3’ calculation where the use of the global optical potential koning03 give rise to a cross section with an energy dependence relatively different from the one measured.
The stellar Maxwellian-averaged neutron capture cross sections for the three reactions studied here are shown in Fig. 7. It is evident that although the photoneutron cross section is relatively insensitive to the nuclear level density and neutron-nucleus optical potential, particularly in the <sup>186</sup>W($`\gamma ,n`$)<sup>185</sup>W case, these quantities can affect the reverse rate significantly. If we characterize the remaining uncertainty affecting the prediction of the neutron capture rate by the ratio between the upper and lower limits obtained in Fig. 7, we find at an energy of 25 keV a factor of 1.9 for the <sup>185</sup>W($`n,\gamma `$)<sup>186</sup>W reaction, of 2.1 for <sup>186</sup>Re($`n,\gamma `$)<sup>187</sup>Re and only 1.1 for <sup>187</sup>Os($`n,\gamma `$)<sup>188</sup>Os. In the last case, the small error bars arise from the additional constraints made available through the experimental <sup>187</sup>Os($`n,\gamma `$)<sup>188</sup>Os cross section.
## V Implications to the Re-Os chronometry
At stellar temperatures relevant to the s-process nucleosynthesis ($`T`$1-3$`\times 10^8`$ K), the <sup>187</sup>Os first excited state at 9.75keV is strongly populated and can significantly affect the estimate of the stellar neutron capture rate on <sup>187</sup>Os. The correction to the cross sections due to the neutron capture on the 9.75 keV state is introduced by $`F_\sigma `$ in Eq. (2).
On the basis of the present photoneutron data and the calculations, the temperature dependence of the $`F_\sigma `$ factor has been re-estimated (Fig. 8 and Table 3). At the s-process temperature of $`3\times 10^8`$ K, all calculations converge to the $`F_\sigma `$ value of about 0.87. However, at lower temperatures, the model INP-3 based on the global phenomenological optical potential predicts significantly larger $`F_\sigma `$ values. This difference mainly originates from deviations seen in the laboratory cross section (cf. Fig. 6) and should therefore be given a lower credibility. The present calculation agrees relatively well with the value of $`0.80F_\sigma 0.83`$ at $`kT`$=30 keV from Ref. Hers83 , but not with the value of $`1F_\sigma 1.15`$ from Ref. Wint82 .
We now consider implications of the $`F_\sigma `$ values constrained by the present study to the Re-Os cosmochronology. The most tantalizing aspect of the Re-Os chronology is that it requires a rather detailed model of the chemical evolution of the Galaxy. Since constructing a reasonable model of the Galactic chemical evolution is beyond the scope of the present study, we rather focus on the uncertainty in the estimate of the age of the Galaxy within a schematic model.
We recall here that abundances of elements in the relevant mass region can be symbolically expressed as follows:
$`{}_{}{}^{\mathit{186}}Os^{}={}_{}{}^{\mathit{186}}Os^s+{}_{}{}^{\mathit{186}}Os^p`$
$`{}_{}{}^{\mathit{187}}Os^{}={}_{}{}^{\mathit{187}}Os^s+{}_{}{}^{\mathit{187}}Os^c`$
$`{}_{}{}^{\mathit{187}}Re^{}={}_{}{}^{\mathit{187}}Re^r+{}_{}{}^{\mathit{187}}Re^s{}_{}{}^{\mathit{187}}Re^c`$
$`{}_{}{}^{\mathit{187}}Os^c={}_{}{}^{\mathit{187}}Re^c`$
Here, $``$, $`s`$, $`p`$, $`r`$, and $`c`$ represent the solar, s-process, p-process, r-process, and cosmoradiogenic origins, respectively. We introduce the following approximations:
$`{}_{}{}^{\mathit{186}}Os^p=p\times {}_{}{}^{\mathit{186}}Os^{}`$
$`{}_{}{}^{\mathit{187}}Re^s=0`$
$`{}_{}{}^{\mathit{187}}Os^s=[F_\sigma \sigma (^{186}`$Os)/$`\sigma (^{187}`$Os)\] $`{}_{}{}^{\mathit{186}}Os^s`$ (local approximation: Eq. (1))
Following the recent p-process calculations Arno03 , we estimated $`p`$ as lying within the 0.01-0.04 range, being in agreement with $`p`$=0.02 Kaep91 . As noted earlier, we ignore the possible s-process contribution to <sup>187</sup>Re through the branchings at <sup>185</sup>W and <sup>186</sup>Re.
Thus, the cosmoradiogenic component of <sup>187</sup>Os can be obtained as
$${}_{}{}^{\mathit{187}}Os^c={}_{}{}^{\mathit{187}}Os^{}(1p)[F_\sigma \sigma (^{186}\mathrm{Os})/\sigma (^{187}\mathrm{Os})]{}_{}{}^{\mathit{186}}Os^{}.$$
(5)
A consideration of the simplest model of a closed system would lead us to the assumption that the evolution of <sup>187</sup>Re can be effectively described by
$$\frac{d{}_{}{}^{\mathit{187}}Re(t)}{dt}=\lambda {}_{\beta }{}^{eff}{}_{}{}^{\mathit{187}}Re(t)+Y(t)$$
(6)
where $`\lambda _\beta ^{eff}`$ is the effective $`\beta `$-decay rate of <sup>187</sup>Re in consideration of some enhancement by astration from the laboratory decay rate of $`\lambda _\beta `$ = ln2/(41.6 Gyr), whereas $`Y(t)`$ term represents the net r-process yield. As in Ref. Kaep91 , we adopt a simple form of $`Y(t)=y\mathrm{exp}(\lambda t)`$ where $`\lambda `$ is a free parameter. Thus, we have
$${}_{}{}^{\mathit{187}}Re(t)=\frac{y[e^{\lambda _\beta ^{eff}t}e^{\lambda t}]}{\lambda \lambda _\beta ^{eff}}.$$
(7)
Using $`d^{\mathit{187}}Os^c(t)/dt`$ = + $`\lambda _\beta ^{eff}{}_{}{}^{\mathit{187}}Re(t)`$, the abundance ratio between $`{}_{}{}^{\mathit{187}}Os^c(t)`$ and $`{}_{}{}^{\mathit{187}}Re(t)`$ can be obtained as
$$\frac{{}_{}{}^{\mathit{187}}Os^c(t)}{{}_{}{}^{\mathit{187}}Re(t)}=\frac{B}{A}$$
(8)
where
$$A=e^{\lambda _\beta ^{eff}t}e^{\lambda t}$$
(9)
$$B=[1e^{\lambda _\beta ^{eff}t}][1e^{\lambda t}]\lambda _\beta ^{eff}/\lambda .$$
(10)
The abundance ratio in Eq. (8) at 4.55 Gyr ago is matched with that from Eq. (5).
For the meteoritic abundance of cosmoradiogenic $`{}_{}{}^{\mathit{187}}Os^c`$ relative to $`{}_{}{}^{\mathit{187}}Re^{}`$, we used the following solar abundances Faes98 :
$`{}_{}{}^{\mathit{187}}Re^{}`$/$`{}_{}{}^{\mathit{186}}Os^{}`$ = 3.51 $`\pm `$ 0.09
$`{}_{}{}^{\mathit{187}}Os^{}`$/$`{}_{}{}^{\mathit{186}}Os^{}`$ = 0.793 $`\pm `$ 0.001.
The ($`n,\gamma `$) cross section ratios at s-process temperatures, $`\sigma (^{186}\mathrm{Os})/\sigma (^{187}\mathrm{Os})`$, were taken from Ref. Kaep91 . Thus, the meteoritic quantity is determined with the present $`F_\sigma `$ value being a unique parameter. As for $`\lambda _\beta ^{eff}`$/$`\lambda _\beta `$ is used to calculate the abundance ratio in Eq. (8) at $`t`$ = $`T_G`$ (the age of the Galaxy) $`4.55`$ Gyr, Clayton Clay88 derived 1.4 from the work of Yokoi et al. Yoko83 , whereas more recent analyses suggest considerably lower values Taka04 . We adopt here 1.2 as our standard value for the net enhancement by astration of the <sup>187</sup>Re $`\beta `$-decay rate.
Matching conditions in $`{}_{}{}^{\mathit{187}}Os^c`$/$`{}_{}{}^{\mathit{187}}Re^{}`$ were investigated in the $`T_G`$ range of 11 - 15 Gyr, showing good agreement with $`0\lambda 2`$ Gyr<sup>-1</sup>. We summarize some conclusions obtained under the simplest assumption of r-process nucleosynthesis yields varying exponentially in time.
* $`F_\sigma `$ values at typical s-process temperatures are in the range of 0.86-0.94 (Table 3).
* The probable range of the differential coefficient $`dT_G/dF_\sigma `$ is $``$(5.0-12.8) Gyr.
* Consequently, the remaining uncertainty of $`T_G`$ that stems from that of $`F_\sigma `$ values is less than 1 Gyr. When the temperature dependence of $`\sigma (^{186}\mathrm{Os})/\sigma (^{187}\mathrm{Os})`$ is considered along with that of $`F_\sigma `$, the uncertainty in $`T_G`$ is approximately halved.
We note here that the model of Yokoi et al. Yoko83 cannot be reconciled with the present $`F_\sigma `$ data. The model of chemical evolution developed there favored $`F_\sigma `$ values much higher than unity. If we mimic the results in terms of Eq. (6), the corresponding values of $`\lambda `$ become negative. This also explains the much larger $`|dT_G/dF_\sigma |`$ values of up to 100 Gyr (as seen in the slope of Fig. 9 in Ref. Yoko83 ). Finally, it is noted that if we would use the age of the universe as derived by the WMAP ($`T_U`$ = 13.7 $`\pm `$ 0.2 Gyr) WMAP , and assume $`T_U`$ \- $`T_G`$ $``$ 1 Gyr, $`\lambda `$ values that are consistent with the present $`F_\sigma `$ value are narrowed in the ranges of 0.16-0.46 Gyr<sup>-1</sup>, inclusively, and 0.22-0.34 Gyr<sup>-1</sup>, exclusively. The exclusive values are commonly allowed in calculations with all possible combinations of $`kT`$, $`F_\sigma `$, and $`T_G`$.
## VI Conclusion
Photoneutron cross sections were measured with accuracy for <sup>186</sup>W, <sup>187</sup>Re and <sup>188</sup>Os using quasi-monochromatic photon beams from laser Compton scattering (LCS) at energies near the neutron thresholds. The cross sections were used to constrain the model parameters in the framework of the Hauser-Feshbach model. Four different sets of nuclear ingredients were adopted to estimate the photoneutron cross sections and the reverse radiative neutron capture cross sections. When no experimental data on the direct ($`n,\gamma `$) cross section is available, an accuracy of about a factor of 2 was achieved in the predictions. The influence of the neutron capture by the 9.75 keV first excited state in <sup>187</sup>Os which is substantially populated in stellar plasmas at typical s-process temperatures has been estimated in connection with the <sup>187</sup>Re-<sup>187</sup>Os cosmochronology and shown to lead to an increase of the neutron capture rate by a factor of about 1.15 at a temperature of $`3\times 10^8`$ K. Uncertainties by about 10% associated with the neutron-nucleus optical potential still affect the stellar rate at temperatures between 1 and $`2\times 10^8`$ K.
The correction factor $`F_\sigma `$ to be used in the local approximation (Eqs. (1) and (2)) was constrained well in the present study (Table 3). Based on the simplest assumption of r-process nucleosynthesis yields varying exponentially in time, the cosmochronological uncertainty in the age of the Galaxy arising from the $`F_\sigma `$ values is estimated to be less than 1 Gyr; when the temperature dependences of both $`\sigma (^{186}`$Os)/$`\sigma (^{187}`$Os) and $`F_\sigma `$ are considered, the uncertainty is less than 0.5 Gyr.
###### Acknowledgements.
H.U. is grateful to Kohji Takahashi for helpful suggestions. This work was done within the Konan-ULB joint project and supported by the Japan Atomic Energy Research Institute (the REIMEI Research Resource), the Japan Private School Promotion Foundation, and the Japan Society for the Promotion of Science. S.G. is FNRS Research Associate.
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# Introduction to cirquent calculus and abstract resource semantics
## 1 Introduction
This paper introduces a refinement of the sequent calculus approach called cirquent calculus. Roughly speaking, the difference between the two is that, while in Gentzen-style proof trees sibling (or cousin, etc.) sequents are disjoint and independent sequences of formulas, in cirquent calculus they are permitted to share elements. Explicitly allowing or disallowing shared resources and thus taking to a more subtle level the resource-awareness intuitions underlying substructural logics, cirquent calculus offers much greater flexibility and power than sequent calculus does.
A need for substantially new deductive tools came with the recent (2003) birth of computability logic (CL), characterized in as “a formal theory of computability in the same sense as classical logic is a formal theory of truth”. Indeed, formulas in CL are seen as computational problems rather than propositions or predicates, and their “truth” seen as algorithmic solvability. In turn, computational problems, understood in their most general — interactive — sense, are defined as games played by an interactive Turing machine against its environment, with “algorithmic solvability” meaning existence of a machine that wins the game against any possible (behavior of the) environment. A core collection of the most basic and natural operations on computational problems forms the logical vocabulary of the theory, with some of those operations, as logical operators, resembling those of linear logic. With this semantics, CL provides a systematic answer to the fundamental question “what can be computed? ”, just as classical logic is a systematic tool for telling what is true. Furthermore, as it turns out, in positive cases “what can be computed” always allows itself to be replaced by “how can be computed”, which makes CL of interest in not only theoretical computer science, but some more applied areas as well, including interactive knowledge base systems and resource oriented systems for planning and action. On the logical side, CL can serve as a basis for constructive applied theories. This is a very brief summary. See , or for elaborated expositions of the philosophy, motivations and techniques of computability logic.<sup>1</sup><sup>1</sup>1A comprehensive online source on CL can be found at http://www.cis.upenn.edu/$``$giorgi/cl.html
The above-mentioned fact of resemblance between computability-logic and linear-logic operators is no accident. Both logics claim to be “logics of resources”, with their philosophies and ambitions thus having a significant overlap. The ways this common philosophy is materialized, however, are rather different. Computability logic directly captures resource intuitions through its semantics. Resources, understood in the specific sense of computational resources, are dual/symmetric to computational problems: what is a problem for the machine, is a resource for the environment (=user), and vice versa. So, as a logic of computational problems, CL also automatically is a logic of computational resources. The scheme that CL follows can be characterized as “from semantics to syntax”: it starts with a clear concept of resources (=computational problems) and resource-semantical validity (=universal algorithmic solvability), and only after that, as a natural second step, asks what the corresponding syntax is, i.e. how the set of valid formulas can be axiomatized. On the other hand, it would be accurate to say that linear logic, as a logic of resources (rather than that of phases or coherence spaces), has started directly from the second step, essentially by taking classical sequent calculus and deleting the structural rules unsound from a naive, purely intuitive resource point of view. For simplicity, in this discussion we narrow linear logic down to its multiplicative fragment; furthermore, taking some terminological liberty, by “linear logic” we mean the version of it more commonly known as affine logic, which is classical sequent calculus without the contraction rule (Girard’s canonical system for linear logic further deletes the rule of weakening as well). Even the most naive and vague resource intuitions are sufficient to see that the deleted rule of contraction, responsible for the principle $`PPP`$, was indeed wrong: having $1 does not imply having $1 and $1, i.e. $2. Such intuitions can also be safely relied upon in deeming all the other rules of classical sequent calculus “right”. To summarize, linear logic is undoubtedly sound as a logic of resources. But even more so is … the empty logic. Completeness is thus a crucial issue. This is where the need for a mathematically strict and intuitively convincing resource semantics becomes critical, without which a question on completeness cannot even be meaningfully asked. Despite intensive efforts, however, such a semantics has never really been found for linear logic. And apparently the reason for this failure is as straightforward as it could possibly be: linear logic, as a resource logic, is simply incomplete. At least, this is what CL believes, for it has been shown () that the semantics of the latter, with its well-justified claims to be a semantics of resources, validates a strictly bigger class of formulas than linear (=affine) logic does.
Taking pride in the meaningfulness of its semantics, computability logic, at the same time, has been suffering from one apparent disadvantage: the absence of a good syntax, i.e. proof-theoretically “reasonable” and nice deductive systems, as opposed to the beauty and harmony of the Gentzen-style axiomatizations for linear logic and its variations, let alone proof nets. True, certain sound and complete systems, named CL1, CL2, CL3 and CL4, have been constructed for incrementally expressive (and rather expressive) fragments of CL in , and probably more results in the same style are still to come. Yet, hardly many would perceive those systems as “logical calculi”, and perhaps not everyone would even call them “deductive systems”. Rather, those somewhat bizarre constructions — one of which (CL2) will be reproduced later in Section 9 — might be seen as just ad hoc syntactic characterizations, offering sound and complete decision or enumeration procedures for the corresponding sets of valid formulas of CL, but otherwise providing no real proof-theoretic insights into this new logic. Repeated attempts to find Gentzen- or Hilbert-style equivalents of those systems have hopelessly failed even at the most basic, $`\neg ,,,`$ (“multiplicative”) level. And probably this failure, just like the failure to find a good resource semantics for linear logic, is no accident. The traditional deductive methods have been originally developed with traditional logics in mind. There are no reasons to expect for those methods to be general and flexible enough to just as successfully accommodate the needs of finer-level semantic approaches, such as the computational semantics of CL, or resource semantics in general. Switching to a novel vision in semantics may require doing the same in syntax.
This is where cirquent calculus as a nontraditional syntax comes in, breaking the stubborn resistance of CL to axiomatization attempts. While the full collection of its rules just offers an alternative axiomatization for the kind old classical logic, removing (the cirquent-calculus version of) contraction from that collection — we call the resulting system CL5 — yields a sound and complete system for the $`(\neg ,,,)`$-fragment of CL, the very core of the logic previously appearing to be “most unaxiomatizable”. Being complete, CL5 is thus strictly stronger than the incomplete affine logic. The latter, by merely deleting the offending rule of contraction without otherwise trying to first appropriately re(de)fine ordinary sequent calculus, has thrown out the baby with the bath water. Among the innocent victims expelled together with contraction is Blass’s principle
$$(PQ)(RS)(PR)(QS),$$
provable in CL5 but not in affine logic, which, in fact, even fails to prove the less general formula
$$(PP)(PP)(PP)(PP).$$
To strengthen the implied claim of computability logic that it is CL5 rather than affine logic that adequately materializes the resource philosophy traditionally associated with the latter, the present paper further introduces an abstract resource semantics and shows that CL5 is sound and complete with respect to that semantics as well. Unlike the semantics of computability logic, which understands resources in the special — computational — sense, abstract resource semantics can be seen as a direct formalization of the more general intuitions in the style “having $1 does not imply having $1 and $1” or “one cannot get both a candy and an apple for a dollar even if one dollar can buy either”. As noted earlier, the inherent incompleteness of linear logic, resulting from the fundamental limitations of the underlying sequent-calculus approach, is the reason why such intuitions and examples, while so heavily relied on in the popular linear-logic literature, have never really found a good explication in the form of a mathematically well-defined semantics.
The set of theorems of CL5 admits an alternative, simple yet non-deductive characterization, according to which this is the set of all binary tautologies and their substitutional instances. Here binary tautologies mean tautologies of classical propositional logic in which no propositional letter occurs more than twice. The class of such formulas has naturally emerged in the past in several unrelated contexts. The earliest relevant piece of literature of which the author is aware is , dating back to 1963, where Jaśkowski studied binary tautologies as the solution to the problem of characterizing the provable formulas of a certain deductive system. Andreas Blass came across the same class of formulas twice. In he introduced a game semantics for linear-logic connectives and found that the multiplicative fragment of the corresponding logic was exactly the class of the substitutional instances of binary tautologies. In the same paper he argued that this class was inherently unaxiomatizable — using his words, “entirely foreign to proof theory”. Such an assessment was both right and wrong, depending on whether proof theory is understood in the strictly traditional (sequent calculus) or a more generous (cirquent calculus) sense. 11 years later, in , using Herbrand’s Theorem, Blass introduced the concept of universal simple Herbrand validity, a natural sort of resource consciousness that makes sense in classical logic. Blass found in that this (non-game) semantics validates exactly the same class of propositional formulas as his unrelated game semantics for the multiplicative fragment of the language of linear logic does. While independently experimenting with various semantical approaches prior to the invention of computability logic, the author of the present paper, too, had found game-semantical soundness and completeness of the class of binary tautologies and their substitutional instances. Once this happened in and then, again, in . The underlying semantics in those two cases were rather different from each other, as well as different from that of CL or Blass’s game semantics. The fact that the set of the theorems of CL5 arises in different approaches by different authors with various motivations and traditions, serves as additional empirical evidence for the naturalness of CL5. This is somewhat in the same sense as the existence of various models of computation that eventually yield the same class of computable functions speaks in favor of the Church-Turing Thesis.
The version of cirquent calculus presented in this paper captures the most basic yet only a modest fraction of the otherwise very expressive language of computability logic. For instance, the formalism of the earlier-mentioned system CL4, in addition to $`\neg ,,,`$ called parallel connectives, has the connectives $`,`$ (choice connectives, resembling the additives of linear logic), and the two groups $`,`$ (choice) and $`,`$ (blind) of quantifiers. Among the other operators officially introduced within the framework of CL so far are the parallel (“multiplicative”) quantifiers $`,`$, and the two groups $`\text{ },\text{ }`$ (branching) and $`\text{ },\text{ }`$ (parallel) of recurrence (“exponential”) operators. Extending the cirquent-calculus approach so as to accommodate incrementally expressive fragments of CL is a task for the future. The results of the present paper could be seen just as first steps on that long road. What is important is that the syntactic ice cover of computability logic, previously having seemed to be unbreakable, is now cracked.
It should be noted that, even though computability logic has been at the center of discussion almost throughout this introductory section, essentially its relevance to the present paper is limited to being the primary source of motivation and inspiration. Cirquent calculus (the very idea of it), abstract resource semantics and all related technical results presented in this paper are new and, as the author wishes to hope, valuable in their own rights.<sup>2</sup><sup>2</sup>2For the exception of the soundness and completeness theorem for CL5 with respect to the CL semantics, of course. There is no overlap with any prior work on CL, and familiarity with the latter, while desirable, is not at all necessary for understanding the present material.
## 2 Cirquents
Throughout the rest of this paper, unless otherwise specified, by a formula we mean one of the language of classical propositional logic. We consider the version of this language that has infinitely many non-logical atoms (also called propositional letters), for which we use the metavariables $`P,Q,R,S`$, and no logical atoms such as $``$ or $``$. The propositional connectives are limited to the unary $`\neg `$ and the binary $`,`$. If we write $`FG`$, it is to be understood as an abbreviation of $`\neg FG`$. Furthermore, we officially allow $`\neg `$ to be applied only to atoms. $`\neg \neg F`$ is to be understood as $`F`$, $`\neg (FG)`$ as $`\neg F\neg G`$, and $`\neg (FG)`$ as $`\neg F\neg G`$.
Where $`k`$ is a natural number, by a $`k`$-ary pool we mean a sequence $`F_1,\mathrm{},F_k`$ of formulas. Such a sequence may have repetitions, and we refer to a particular occurrence of a formula in a pool as an oformula, with the prefix “o” derived from “occurrence”. This prefix will as well be used with a similar meaning in other words and contexts where same objects — such as, say, atoms or subformulas — may have several occurrences. Thus, the pool $`F,G,F`$ has two formulas but three oformulas; and the formula $`P(P\neg P)`$ has one atom but three oatoms. For readability, we usually refer to oformulas (or oatoms, etc.) by the name of the corresponding formula (atom, etc.), as in the phrase “the oformula $`F`$”, assuming that it is clear from the context which of the possibly many occurrences of $`F`$ we have in mind.
A $`k`$-ary cirquentstructure, or simply structure, is a finite sequence $`\mathrm{𝐒𝐭}=\mathrm{\Gamma }_1,`$ $`\mathrm{},\mathrm{\Gamma }_m`$ ($`m0`$), where each $`\mathrm{\Gamma }_i`$, called a group of $`\mathrm{𝐒𝐭}`$, is a subset of $`\{1,\mathrm{},k\}`$. As in pools, here we may have $`\mathrm{\Gamma }_i=\mathrm{\Gamma }_j`$ for some $`ij`$. Again, to differentiate between a particular occurrence of a group in the structure from the group as such, we use the term ogroup. The structure $`\{1,2\},\mathrm{},\{1,3\},\{1,2\}`$ thus has three groups but four ogroups. Yet, as in the case of oformulas or oatoms, we may just say “the ogroup $`\{1,2\}`$” if it is clear which of the two occurrences of the group $`\{1,2\}`$ is meant.
###### Definition 2.1
A $`k`$-ary ($`k0`$) cirquent is a pair $`C=(\text{St}^C,\text{Pl}^C)`$, where $`\text{St}^C`$, called the structure of $`C`$, is a $`k`$-ary cirquentstructure, and $`\text{Pl}^C`$, called the pool of $`C`$, is a $`k`$-ary pool.
An (o)group of such a $`C`$ will mean an (o)group of $`\text{St}^C`$, and an (o)formula of $`C`$ mean an (o)formula of $`\text{Pl}^C`$. Also, we often prefer to think of the groups of $`C`$ as sets of its oformulas rather than sets of the corresponding ordinal numbers. For example, if $`\text{Pl}^C=F,G,H`$ and $`\mathrm{\Gamma }=\{1,3\}`$, we can understand the same $`\mathrm{\Gamma }`$ as the set $`\{F,H\}`$ of oformulas. In this case we say that $`\mathrm{\Gamma }`$ contains $`F`$ and $`H`$. When $`\mathrm{\Gamma }`$ is seen as an ogroup rather than a group, we may also refer to such a set $`\{F,H\}`$ as the content of $`\mathrm{\Gamma }`$.
An $`1`$-ary cirquent whose only ogroup is $`\{1\}`$ is said to be a singleton.
We represent cirquents using diagrams, such as the one shown below:
This diagram represents the cirquent whose pool is $`F,G,H,F`$ and whose structure is $`\{1\},\{2,3\},\{3,4\}`$. We typically do not terminologically distinguish between cirquents and diagrams: for us, a diagram is (rather than represents) a cirquent, and a cirquent is a diagram. The top level of a diagram thus lists the oformulas of the cirquent, and the bottom level lists its ogroups, with each ogroup represented by (and identified with) a $``$, where the arcs (lines connecting the $``$ with oformulas) are pointing to the oformulas that a given ogroup contains. The horizontal line at the top of the diagram is just to indicate that this is one cirquent rather than, say, two cirquents (one $`1`$-ary and one $`3`$-ary) put together. Our convention is that such a line should be present even if there is no potential ambiguity. It is required to be long enough — and OK if longer than necessary — to cover all of the oformulas and ogroups of the cirquent.
The term “cirquent” is a hybrid of “circuit” and “sequent”. So is, in a sense, its meaning. Cirquents can be seen to generalize sequents by imposing circuit-style structures on their oformulas. In a preliminary attempt to see some familiar meaning in cirquents, it might be helpful to think of them as Boolean circuits of depth 2, with oformulas serving as inputs, all first-level gates — representing ogroups — being $``$-gates, and the only second-level gate, connected to each first-level gate, being an $``$-gate. This is illustrated in Figure 1:
In traditional logic, circuits are interesting only in the context of representational complexity, and otherwise they do not offer any additional expressive power, for duplicating or merging identical nodes creates no difference when Boolean functions are all one sees in circuits. So, from the classical perspective, the circuit of Figure 1 is equivalent to either circuit of Figure 2, with the tree-like circuit on the right being a direct reading of the formula $`F(GH)(HF)`$ expressing the Boolean function of the circuit of Figure 1, and the circuit on the left being a most economical representation of the same Boolean function:
Linear logic, understanding the nodes of the circuit as representing resources rather than just Boolean values, would not agree with such an equivalence though: the first and the fourth upper-level nodes of the circuit of Figure 1, even though having the same type, would be seen as two different individual resources. What linear logic generally fails to account for, however, is the possibility of resource sharing. $`H`$ is a resource shared by two different compound resources — the resources represented by #2 and #3 $``$-gates of Figure 1. Allowing shared resources in cirquent calculus refines the otherwise crude approach of linear logic. And by no means does it mean departing from the idea that resources should be accurately book-kept. Indeed, a shared resource does not mean a duplicated resource. Imagine Victor has $10,000 on his bank account. One day he decides to give his wife access to the account. From now on the $10,000 is shared. Two persons can use it, either at once, or portion by portion. Yet, this does not turn the $10,000 into $20,000, as the aggregate possible usage remains limited to $10,000.
## 3 Core cirquent calculus rules
Different cirquent calculus systems will differ in what logical operators and atoms their underlying formal languages include, and what rules of inference they allow. The underlying language is fixed in this paper (it will only be slightly extended in the last paragraph of Section 9). And all of the rules will come from the ones introduced in the present section. We explain those rules in a relaxed fashion, in terms of inserting arcs, swapping oformulas, etc. Such explanations are rather clear, and translating them into rigorous formulations in the style and terms of Definition 2.1, while possible, is hardly necessary.
We need to agree on some additional terminology first. Adjacent oformulas of a given cirquent are two oformulas $`F`$ and $`G`$ with $`G`$ appearing next to (to the right of) $`F`$ in the pool of the cirquent. We say that $`F`$ immediately precedes $`G`$, and that $`G`$ immediately follows $`F`$. Similarly for adjacent ogroups.
By merging two adjacent ogroups $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ in a given cirquent $`C`$ we mean replacing in $`C`$ the two ogroups $`\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$ by the one ogroup $`\mathrm{\Gamma }\mathrm{\Delta }`$, leaving the rest of the cirquent unchanged. The resulting cirquent will thus only differ from $`C`$ in that it will have one $``$ where $`C`$ had the two adjacent $``$s, with the arcs of this new $``$ pointing exactly to the oformulas to which the arcs of one or both of the old $``$s were pointing. For example, the right cirquent of the following figure is the result of merging ogroups #2 and #3 in the left cirquent:
Merging two adjacent oformulas $`F`$ and $`G`$ into $`H`$ means replacing those two oformulas by the one oformula $`H`$, and redirecting to it all arcs that were pointing to $`F`$ or $`G`$. For example, the right cirquent of the following figure is the result of merging, in the left cirquent, (the first) $`F`$ and $`G`$ into $`H`$:
Now we are ready to look at the rules.
### 3.1 Axioms (A)
Axioms are “rules” with no premises. There are two axioms, called the empty cirquent axiom and the identity axiom. The first one introduces the empty cirquent $`(,)`$ (both the pool and the structure are empty); the second one which — just like the rest of the rules — is, in fact, a scheme of rules because $`F`$ can be an arbitrary formula, introduces the cirquent $`(\{1,2\},\neg F,F)`$.
The letter “A” next to the horizontal line stands for the name of the rule by which the conclusion is obtained. We will follow the same notational practice for the other rules.
### 3.2 Mix (M)
This rule takes two premises. The conclusion is obtained by simply putting one premise next to the other, thus creating one cirquent out of the two, as illustrated below:
### 3.3 Exchange (E)
This and all of the remaining rules take a single premise. The exchange rule comes in two flavors: oformula exchange and ogroup exchange. The conclusion of oformula (resp. ogroup) exchange is the result of swapping in the premise two adjacent oformulas (resp. ogroups) and correspondingly redirecting all arcs. The following oformula exchange example swaps $`F`$ with $`G`$; and the ogroup exchange example swaps ogroup #2 with ogroup #3:
The presence of oformula exchange essentially allows us to treat the pool of a cirquent as a multiset rather than a sequence of formulas. Similarly, the presence of ogroup exchange makes it possible to see the structure of a cirquent as a multiset rather than a sequence of groups.
### 3.4 Weakening (W)
This rule, too, comes in two flavors: ogroup weakening and pool weakening. In the first case the conclusion is the result of inserting a new arc between an existing ogroup and an existing oformula of the premise. In the second case, the conclusion is the result of inserting a new oformula anywhere in the pool of the premise.
### 3.5 Duplication (D)
This rule comes in two versions as well: downward duplication and upward duplication. The conclusion (resp. premise) of downward (resp. upward) duplication is the result of replacing in the premise (resp. conclusion) some ogroup $`\mathrm{\Gamma }`$ by two adjacent ogroups that, as groups, are identical with $`\mathrm{\Gamma }`$.
Note that the presence of duplication together with ogroup exchange further allows us to think of the structure of a cirquent as a set rather than a sequence or multiset of groups.
### 3.6 Contraction (C)
The premise of this rule is a cirquent with two adjacent oformulas $`F,F`$ that are identical as formulas. The conclusion is obtained from the premise by merging those two oformulas into $`F`$. The following two examples illustrate applications of contraction.
### 3.7 $``$-introduction ($``$)
The conclusion of this rule is obtained by merging in the premise some two adjacent oformulas $`F`$ and $`G`$ into $`FG`$. We say that this application of the rule introduces $`FG`$. Below are three illustrations:
In what we call conservative $``$-introduction (the rightmost example), which is a special case of $``$-introduction, the situation is that whenever an ogroup of the conclusion contains the introduced $`FG`$, the corresponding ogroup of the premise contains both $`F`$ and $`G`$. In a general case (the first two examples), this is not necessary. What is always necessary, however, is that if an ogroup of the conclusion contains the introduced $`FG`$, then the corresponding ogroup of the premise should contain at least one of the oformulas $`F,G`$.
We have just used and will continue to use the jargon “the corresponding ogroup”, whose meaning should be clear: the present rule does not change the number or order of ogroups, and it only modifies the contents of some of those ogroups. So, to ogroup #$`i`$ of the conclusion corresponds ogroup #$`i`$ of the premise, and vice versa. The same applies to the rules of oformula exchange, weakening and contraction. In an application of ogroup exchange that swaps ogroups #$`i`$ and #$`i+1`$, to ogroup #$`i`$ of the premise corresponds ogroup #$`i+1`$ of the conclusion, and vice versa; to ogroup #$`i+1`$ of the premise corresponds ogroup #$`i`$ of the conclusion, and vice versa; and to any other ogroup #$`j`$ of the premise corresponds ogroup #$`j`$ of the conclusion and vice versa. Finally, in an application of mix, to ogroup #$`i`$ of the first premise corresponds ogroup #$`i`$ of the conclusion, and vice versa; and, where $`n`$ is the number of the ogroups of the first premise, to ogroup #$`i`$ of the second premise corresponds ogroup #$`n+i`$ of the conclusion, and vice versa.
### 3.8 $``$-introduction ($``$)
The premise of this rule is a cirquent with adjacent oformulas $`F`$ and $`G`$, such that the following two conditions are satisfied:
* No ogroup contains both $`F`$ and $`G`$.
* Every ogroup containing $`F`$ is immediately followed by an ogroup containing $`G`$, and every ogroup containing $`G`$ is immediately preceeded by an ogroup containing $`F`$.
The conclusion is obtained from the premise by merging each ogroup containing $`F`$ with the immediately following ogroup (containing $`G`$) and then, in the resulting cirquent, merging $`F`$ and $`G`$ into $`FG`$. In this case we say that the the rule introduces $`FG`$.
Below are three examples for the simple case when there is only one ogroup in the conclusion that contains the introduced $`FG`$:
Perhaps this rule is easier to comprehend in the bottom-up (from conclusion to premise) view. To obtain a premise from the conclusion (where $`FG`$ is the introduced conjunction), we “split” every ogroup $`\mathrm{\Gamma }`$ containing $`FG`$ into two adjacent ogroups $`\mathrm{\Gamma }^F`$ and $`\mathrm{\Gamma }^G`$, where $`\mathrm{\Gamma }^F`$ contains $`F`$ (but not $`G`$), and $`\mathrm{\Gamma }^G`$ contains $`G`$ (but not $`F`$); all other ($`FG`$) oformulas of $`\mathrm{\Gamma }`$ — and only such oformulas — should be included in either $`\mathrm{\Gamma }^F`$, or $`\mathrm{\Gamma }^G`$, or both. In what we call conservative $``$-introduction, all of the non-$`FG`$ oformulas of $`\mathrm{\Gamma }`$ should be included in both $`\mathrm{\Gamma }^F`$ and $`\mathrm{\Gamma }^G`$.
The following is an example of an application of the $``$-introduction rule in a little bit more complex case where the conclusion has two ogroups containing the introduced conjunction. It is not a conservative one. To make this application conservative, we should add two more arcs to the premise: one connecting ogroup #3 with $`J`$, and one connecting ogroup #4 with $`E`$.
## 4 Cirquent calculus systems
By a cirquent calculus system in the present context we mean any subset of the set of the eight rules of the previous section. The one that has the full collection of all eight rules we denote by CCC (“Classical Cirquent Calculus”), and the one that has all rules but contraction we denote by CL5. Any other system we denote by placing the abbreviated names of the corresponding rules between parentheses. For instance, (AME) stands for the system that has the axioms, mix and exchange.
Let $`S`$ be a cirquent calculus system, and $`C,A_1,\mathrm{},A_n`$ (possibly $`n=0`$) any cirquents. A derivation of $`C`$ from $`A_1,\mathrm{},A_n`$ in $`S`$ is a tree of cirquents with $`C`$ at its root, where each node is a cirquent that either follows from its children by one of the rules of $`S`$, or is among $`A_1,\mathrm{},A_n`$ (and has no children). A derivation of $`C`$ in $`S`$ from the empty set of cirquents is said to be a proof of $`C`$ in $`S`$. Of course, if $`S`$ does not contain axioms, then there will be no proofs in it.
Throughout this paper we identify each formula $`F`$ with the singleton cirquent $`(\{1\},F)`$, i.e. the cirquent
Correspondingly, a proof or derivation of a given formula $`F`$ in a given system $`S`$ is a proof or derivation of $`(\{1\},F)`$. The following is an example of a proof of $`\neg F(FF)`$ in (AMEC$``$):
It is our convention that if a proof is a proof of a formula $`F`$, then the last cirquent we simply represent as “$`F`$” rather than through a diagram. Just to save space. In a similar space-saving spirit, we will often combine several obvious steps together, labeling the combined application of a “rule” by the name of the system which contains all of the rules that have been combined. For instance, the above derivation of $`\neg F(FF)`$ we might want to rewrite in a more compact yet clear way as follows:
Below is an (AME$``$)-proof of Blass’s principle mentioned in Section 1:
## 5 Classical and affine logics
In sequent calculus (where a sequent means a nonempty sequence of formulas), classical logic can be axiomatized by the following six rules, where $`F,G`$ stand for any formulas and $`\mathrm{\Gamma },\mathrm{\Delta }`$ stand for any — possibly empty — sequences of formulas:
Affine logic is obtained from classical logic by deleting contraction. As noted earlier, the term “affine logic” in this paper refers to what is called the multiplicative fragment of this otherwise more expressive logic. A sequent calculus system, in general, is any subset of the above six rules. The definition of provability of a sequent $`\mathrm{\Gamma }`$ in a sequent calculus system $`S`$ is standard: this means existence of a tree of sequents — called a proof tree for $`\mathrm{\Gamma }`$ — with $`\mathrm{\Gamma }`$ at its root, in which every node of the tree follows from its children (where the set of children may be empty in the case of axiom) by one of the rules of $`S`$. A formula $`F`$ is considered provable in a sequent calculus system iff $`F`$, viewed as a one-element sequent, is provable.
At the end of Section 4 we saw that cirquent calculus needs neither weakening nor contraction (nor duplication) to prove Blass’s principle. Replacing all atoms by $`P`$ in our proof tree for Blass’s principle also yields an (AME$``$)-proof of
$$\left((\neg P\neg P)(\neg P\neg P)\right)\left((PP)(PP)\right).$$
(1)
The following Fact 5.1 establishes that, in contrast, sequent calculus needs both weakening and contraction to prove (1), let alone the more general Blass’s principle.<sup>3</sup><sup>3</sup>3That affine logic does not prove (1) was shown by Blass in .
###### Fact 5.1
Any proof of (1) in sequent calculus would have to use both weakening and contraction.
Proof. First, let us attempt to construct, in a bottom-up fashion, a proof of (1) in affine logic to see that such a proof does not exist. The only rule that can yield (1) is $``$-introduction, so the premise should be the sequent
$$(\neg P\neg P)(\neg P\neg P),(PP)(PP).$$
Weakening is not applicable to the above sequent, for both of its formulas are non-valid in the classical sense and hence, in view of the known fact that all of the sequent calculus rules preserve classical validity, those formulas, in isolation, are not provable. $``$-introduction is not applicable, either, for there is no disjunction on the surface of the sequent. And exchange, of course, would not take us closer to our goal of finding a proof. This leaves us with $``$-introduction. The sequent is symmetric, so we may assume that the introduced conjunction is, say, the first one. The non-active formula $`(PP)(PP)`$ of the conclusion can then only be inherited by one of the premises, meaning that the other premise will be just $`\neg P\neg P`$. Now we are stuck with that premise, as it is a non-tautological formula which cannot be proven.
Next, for a contradiction, assume that there is a weakening-free (but not necessarily contraction-free) sequent calculus proof of (1). Consider any branch $`\mathrm{\Gamma }_1,\mathrm{},\mathrm{\Gamma }_n`$ of the proof tree, where $`\mathrm{\Gamma }_1`$ should be $`\neg F,F`$ for some formula $`F`$, and $`\mathrm{\Gamma }_n`$ be the sequent consisting just of (1). Notice that once a given sequent $`\mathrm{\Gamma }_i`$ of the above sequence contains a formula $`G`$, $`G`$ will be inherited by each of the subsequent sequents $`\mathrm{\Gamma }_{i+1},\mathrm{},\mathrm{\Gamma }_n`$ — either as a formula of the sequent, or as a subformula of such. So, both $`F`$ and $`\neg F`$ should be subformulas of (1). This leaves us only with the possibility $`\{F,\neg F\}=\{P,\neg P\}`$, because (1) does not contain any other subformula $`F`$ together with $`\neg F`$. Let $`i`$ be the greatest number among $`1,\mathrm{},n1`$ such that $`\mathrm{\Gamma }_{i+1}`$ is neither $`\neg P,P`$ nor $`P,\neg P`$. $`\mathrm{\Gamma }_{i+1}`$ cannot be derived from $`\mathrm{\Gamma }_i`$ by exchange because then $`\mathrm{\Gamma }_{i+1}`$ would again be $`\neg P,P`$ or $`P,\neg P`$. Nor can it be derived by contraction which is simply not applicable to $`\mathrm{\Gamma }_i`$. Nor can $`\mathrm{\Gamma }_{i+1}`$ be derived by $``$-introduction, because then $`\mathrm{\Gamma }_{i+1}`$ would be $`\neg PP`$ or $`P\neg P`$, which is not a subformula of (1). Finally, $`\mathrm{\Gamma }_{i+1}`$ cannot be derived from $`\mathrm{\Gamma }_i`$ (and an arbitrary other premise) by $``$-introduction either. This is so because an application of this rule would introduce a conjunction where $`\neg P`$ or $`P`$ is a conjunct; but, again, (1) does not have such a subformula. $`\mathrm{}`$
As we just saw, cirquent calculus indeed offers a substantially more flexible machinery for constructing (substructural) deductive systems than sequent calculus does. Sequent calculus can be seen as a simple special case of cirquent calculus that we call “primitive”. Specifically, we say that a cirquent is primitive iff all of its ogroups are (pairwise) disjoint. The groups of such a cirquent can be thought of as — and identified with — sequents: in this section we will not terminologically distinguish between a group $`\mathrm{\Gamma }`$ of a primitive cirquent and the sequent $`\mathrm{\Gamma }`$ consisting exactly of the oformulas that $`\mathrm{\Gamma }`$ contains, arranged in the same order as they appear in the pool of the cirquent.
For any given cirquent calculus system $`S`$, we let $`S^{}`$ denote the version of $`S`$ where the definition of a proof or a derivation has the additional condition that every cirquent in the proof or derivation should be primitive. So, $`S^{}`$ can be called the “primitive version” of $`S`$. Of course, $`S`$ proves or derives everything that $`S^{}`$ does.
Strictly speaking, a sequent or cirquent calculus system is the particular collection of its rules, so that even if two systems prove exactly the same formulas or sequents or cirquents, they should count as different systems. Yet, often we identify a sequent or cirquent calculus system with the set of formulas (or sequents, or cirquents) provable in it, as done in the following Theorem 5.2. The equalities in the left column of that theorem, as can be seen from our subsequent proof of it, extend to all other natural pairs of systems obtained by allowing/disallowing various rules, such as affine logic without weakening (i.e. linear logic in the proper sense) vs. the primitive version of CL5 without weakening, or classical logic without weakening vs. the primitive version of CCC without weakening. It is such equalities that allow us to say that sequent calculus is nothing but the primitive version of cirquent calculus. That primitiveness makes cirquent calculus degenerate to sequent calculus is no surprise. The former owes its special power to the ability to express resource sharing, and it is exactly resource sharing that primitive cirquents forbid.
###### Theorem 5.2
With the following sequent calculus and cirquent calculus systems identified with the sets of formulas that they prove, we have:
$$\begin{array}{ccccccc}\text{1. Affine logic}\hfill & =& 𝐂L5^{}\hfill & & 𝐂L5\hfill & & \text{Affine logic.}\hfill \\ \text{2. Classical logic}\hfill & =& 𝐂CC^{}\hfill & & 𝐂CC\hfill & =& \text{Classical logic.}\hfill \end{array}$$
Proof. As noted earlier, the inclusions of the type $`S^{}S`$ are trivial. The inequality CL5 $``$ Affine logic immediately follow from Fact 5.1 together with the earlier-established provability of (1) in (AME$``$). The equality CCC $`=`$ Classical logic follows from Theorem 6.3, which will be proven in the next section. The latter implies that a formula is provable in CCC iff it is a tautology in the classical sense, and it just remains to remember that the same is known to be true for Classical logic.
Our task now is to verify the equalities Affine logic = CL5 and Classical logic = CCC.
The inclusions Affine logic $``$ CL5 and Classical logic $``$ CCC can be proven by showing that whenever either sequent calculus system proves a sequent $`\mathrm{\Gamma }`$, the corresponding primitive cirquent calculus system proves the cirquent whose only group — as well as pool — is $`\mathrm{\Gamma }`$. This can be easily done by induction on the heights of proof trees. The steps of such induction are rather straightforward, for every application of a sequent calculus rule — except weakening and $``$-introduction — directly translates into an application of the same-name rule of cirquent calculus as shown below:
As for weakening and $``$-introduction, their sequent-calculus to cirquent-calculus translations take two steps:
The inclusions CL5 $``$ Affine logic and CCC $``$ Classical logic can be verified in a rather similar way. Specifically, this can be done by showing that, whenever the primitive version of either cirquent calculus system proves a given cirquent $`C`$, the corresponding sequent calculus system proves each of the groups of $`C`$ understood as sequents. Induction on the heights of proof trees is again the way to proceed. The basis of induction is straightforward, taking into account that the translation shown earlier for identity axiom works in either direction, and that the case of the empty cirquent axiom is vacuous as there are no groups in its “conclusion”. Duplication does not need to be considered, for either the premise or the conclusion of it has to be non-primitive, which automatically bans this rule in primitive cirquent calculus systems. The inductive steps dealing with mix or pool weakening are trivial, because these rules do not create new groups or affect the contents of the existing groups. The same is true for ogroup exchange, as well as oformula exchange if it is external, i.e. swaps oformulas that are in different groups, as opposed to internal exchange that swaps oformulas that are in the same group. What internal oformula exchange, ogroup weakening, contraction, $``$-introduction and $``$-introduction do in primitive cirquents is that they modify one or two of the groups of the premise without affecting any other groups (if there are such). This local behavior allows us to pretend for our present purposes that simply there are no other groups in the cirquent under question. So, in inductive steps dealing with internal oformula exchange, contraction and $``$-introduction, we can rely on the fact that the above-illustrated translations between the sequent- and cirquent-calculus versions of these rules work in either direction. As for ogroup weakening and $``$-introduction, their cirquent-calculus to sequent-calculus translations work as follows:
$`\mathrm{}`$
## 6 Tautologies
By a classical model, or simply model, we mean a function $`M`$ that assigns a truth value — true ($`1`$) or false ($`0`$) — to each atom, and extends to compound formulas in the standard classical way. The traditional concepts of truth and tautologicity naturally extend from formulas to groups and cirquents. Let $`M`$ be a model, and $`C`$ a cirquent. We say that a group $`\mathrm{\Gamma }`$ of $`C`$ is true in $`M`$ iff at least one of its oformulas is so. And $`C`$ is true in $`M`$ if every group of $`C`$ is so. “False”, as always, means “not true”. Finally, $`C`$ or a group $`\mathrm{\Gamma }`$ of it is a tautology iff it is true in every model. Identifying each formula $`F`$ with the singleton cirquent $`(\{1\},F)`$, our concepts of truth and tautologicity of cirquents preserve the standard meaning of these concepts for formulas. Let us mark the evident fact that a cirquent is tautological if and only if all of its groups are so. Note also that a cirquent containing the empty group is always false, while a cirquent with no groups, such as the empty cirquent $`(,)`$, is always true.
###### Lemma 6.1
All of the rules of Section 3 preserve truth in the top-down direction — that is, whenever the premise(s) of an application of any given rule is (are) true in a given model, so is the conclusion. Taking no premises, (the conclusions of) axioms are thus tautologies.
Proof. A routine examination of those rules and our definition of truth for cirquents. $`\mathrm{}`$
###### Lemma 6.2
The rules of mix, exchange, duplication, contraction, conservative $``$-introduction and conservative $``$-introduction preserve truth in the bottom-up direction as well — that is, whenever the conclusion of an application of such a rule is true in a given model, so is (are) the premise(s).
Proof. The above statement for mix, exchange, duplication and contraction is rather obvious. Let us only examine it for the conservative versions of $``$\- and $``$-introduction.
Conservative $``$-introduction: Assume the disjunction that the rule introduced is $`FG`$. Notice that the only difference between the conclusion and the premise is that wherever the conclusion has an ogroup $`\mathrm{\Gamma }`$ containing the oformula $`FG`$, the premise has the ogroup $`\mathrm{\Gamma }^{}=(\mathrm{\Gamma }\{FG\})\{F,G\}`$ instead. Since truth-semantically a group is nothing but the disjunction of its oformulas, the truth values of $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }^{}`$ (in whatever model $`M`$) are the same. Hence so are those of the conclusion and the premise.
Conservative $``$-introduction: Assume the conjunction that the rule introduced is $`FG`$. The only difference between the conclusion and the premise is that wherever the conclusion has an ogroup $`\mathrm{\Gamma }`$ containing $`FG`$, the premise has the two ogroups $`\mathrm{\Gamma }_F=(\mathrm{\Gamma }\{FG\})\{F\}`$ and $`\mathrm{\Gamma }_G=(\mathrm{\Gamma }\{FG\})\{G\}`$ instead. Obviously this implies that if $`\mathrm{\Gamma }`$ is true in a given model, then so are both $`\mathrm{\Gamma }_F`$ and $`\mathrm{\Gamma }_G`$. The above, in turn, implies that if the conclusion is true and hence all of its groups are true, then so are all of the groups of the premise, and hence the premise itself. $`\mathrm{}`$
We say that an oformula $`F`$ of (the pool of) a cirquent $`C`$ is homeless iff no group of $`C`$ contains $`F`$.
A literal means $`P`$ (positive literal of type $`P`$) or $`\neg P`$ (negative literal of type $`P`$) for some atom $`P`$. The term oliteral has the expected meaning: this is a particular occurrence of a literal in a formula or in (the pool or an oformula of) a cirquent. A literal cirquent is a cirquent whose pool contains only literals. An essentially literal cirquent is one every oformula of whose pool either is an oliteral or is homeless.
###### Theorem 6.3
A cirquent is provable in CCC iff it is a tautology.
Proof. The soundness part of this theorem is an immediate corollary of Lemma 6.1. For the completeness part, consider any tautological cirquent $`A`$. In the bottom-up sense, keep applying to it<sup>4</sup><sup>4</sup>4As this can be understood, here and later in similar contexts, “it” means $`A`$ only in the beginning, then “it” becomes the premise of $`A`$, then the premise of the premise of $`A`$, and so on. conservative $``$-introduction and conservative $``$-introduction — in whatever order you like — until you hit an essentially literal cirquent $`B`$, such as the one shown in the following example:
The above procedure will indeed always hit an essentially literal cirquent because conservative $``$-introduction is obviously always applicable when a given cirquent (conclusion) has a non-homeless oformula $`FG`$, and so is conservative $``$-introduction whenever such a cirquent has a non-homeless oformula $`FG`$. $`A`$ thus follows from $`B`$ in ($``$). In view of Lemma 6.2, $`B`$ is a tautology, and since all of its non-homeless oformulas are literals, the tautologicity of $`B`$ obviously means that every group of it contains at least one pair of $`P,\neg P`$ of same-type positive and negative oliterals. Fix one such pair for each group, and then apply (in the bottom-up sense) to $`B`$ a series of weakenings to first delete, in each group, all arcs but the two arcs pointing to the two chosen oliterals, and next delete all homeless oformulas if any such oformulas are present. This is illustrated below:
Every group of the resulting cirquent $`C`$ will thus have exactly two oformulas: some atom and its negation. $`B`$ follows from $`C`$ in (W), so that $`A`$ follows from $`C`$ in (W$``$). Now apply (again in the bottom-up fashion) a series of contractions to $`C`$ to separate all shared oliterals, as illustrated in the example below, with the resulting cirquent called $`D`$:
So, our original cirquent $`A`$ is derivable from $`D`$ in (WC$``$). Every ogroup of $`D`$ is disjoint from every other ogroup and, as in $`C`$, each such ogroup contains exactly two oformulas: $`P`$ and $`\neg P`$ for some atom $`P`$. Therefore $`D`$ is provable in (AME) as illustrated below:
(In the pathological case of $`D`$ having no groups at all, it is simply the empty cirquent and hence an axiom.)
We conclude that $`A`$ is provable in (AMEWC$``$) and hence in CCC. $`\mathrm{}`$
###### Remark 6.4
Note that our completeness proof of Theorem 6.3 does not appeal to duplication, which means that CCC=(AMEWC$``$). Duplication is thus syntactic sugar for CCC. Furthermore, from the completeness proof given in the next section it can be seen that downward (though not upward) duplication does not really add anything to the deductive power of CL5, either. How sweet is the sugar of duplication? It can certainly improve proof sizes, but the possible magnitude of that improvement is unknown at this point. One claim that we are making without a proof, however, is that CL5 minus duplication has polynomial-size proofs while still remaining strictly stronger than affine logic (remember that Blass’s principle is provable in such a system but not in affine logic).
## 7 Binary tautologies and their instances
Let $`C`$ be a cirquent. An oatom $`P`$ of $`C`$, i.e. an occurrence of an atom $`P`$ in $`C`$, is an occurrence of $`P`$ in an oformula of (the diagram of) $`C`$. Such an oatom is negative if it comes with a $`\neg `$; otherwise it is positive. When $`C`$ is a cirquent or formula, by an “atom of $`C`$” we mean an atom that has at least one occurrence in $`C`$.
A substitution is a function $`\sigma `$ that sends every atom $`P`$ to some formula $`\sigma (P)`$; if such a $`\sigma (P)`$ always (for every $`P`$) is an atom, then $`\sigma `$ is said to be an atomic-level substitution. Function $`\sigma `$ extends from atoms to all formulas in the expected way: $`\sigma (\neg P)=\neg \sigma (P)`$; $`\sigma (FG)=\sigma (F)\sigma (G)`$; $`\sigma (FG)=\sigma (F)\sigma (G)`$. $`\sigma `$ also extends to cirquents $`C`$ by stipulating that $`\sigma (C)`$ is the result of replacing in $`C`$ every oformula $`F`$ by $`\sigma (F)`$.
Let $`A`$ and $`B`$ be cirquents. We say that $`B`$ is a (substitutional) instance of $`A`$ iff $`B=\sigma (A)`$ for some substitution $`\sigma `$; and $`B`$ is an atomic-level instance of $`A`$ iff $`B=\sigma (A)`$ for some atomic-level substitution $`\sigma `$. Example: the second cirquent of Figure 3 is an instance — though not an atomic-level one — of the first cirquent; the (relevant part of the) substitution $`\sigma `$ used here is defined by $`\sigma (P)=QP`$, $`\sigma (Q)=Q`$ and $`\sigma (R)=P`$.
###### Lemma 7.1
If a given cirquent calculus system proves a cirquent $`C`$, then it also proves every instance of $`C`$.
Proof. Consider a proof tree $`T`$ of an arbitrary cirquent $`C`$, and an arbitrary instance $`C^{}`$ of $`C`$. Let $`\sigma `$ be a substitution with $`\sigma (C)=C^{}`$. Replace every oformula $`F`$ of every cirquent of $`T`$ by $`\sigma (F)`$. It is not hard to see that the resulting tree $`T^{}`$, which uses exactly the same rules as $`T`$ does, is a proof of $`C^{}`$. $`\mathrm{}`$
A cirquent is said to be binary iff no atom has more than two occurrences in it. A binary cirquent is said to be normal iff, whenever it has two occurrences of an atom, one occurrence is negative and the other is positive. A binary tautology (resp. normal binary tautology) is a binary (resp. normal binary) cirquent that is a tautology in the sense of the previous section. This terminology also extends to formulas understood as cirquents. The left cirquent of Figure 3 is an example of a normal binary tautology.
###### Lemma 7.2
A cirquent is an instance of some binary tautology iff it is an atomic-level instance of some normal binary tautology.
Proof. The “if” part is trivial. For the “only if” part, assume $`A`$ is an instance of a binary tautology $`B`$. Let $`P_1,\mathrm{},P_n`$ be all of the atoms of $`B`$ that have two positive or two negative occurrences in $`B`$. Let $`Q_1,\mathrm{},Q_n`$ be any pairwise distinct atoms not occurring in $`B`$. Let $`C`$ be the result of replacing in $`B`$ one of the two occurrences of $`P_i`$ by $`Q_i`$, for each $`i=1,\mathrm{},n`$. Then obviously $`C`$ is a normal binary cirquent, and $`B`$ an instance of it. By transitivity, $`A`$ (as an instance of $`B`$) is also an instance of $`C`$.
We want to see that $`C`$ is a tautology. Deny this. Then there is a classical model $`M`$ in which $`C`$ is false. Let $`M^{}`$ be the model such that:
* $`M^{}`$ agrees with $`M`$ on all atoms that are not among $`P_1,\mathrm{},P_n,Q_1,\mathrm{},Q_n`$;
* for each $`i\{1,\mathrm{},n\}`$, $`M^{}(P_i)=M^{}(Q_i)=\{\begin{array}{c}\text{false}\text{ if }P_i\text{ and }Q_i\text{ are positive in }C\hfill \\ \text{true}\text{ if }P_i\text{ and }Q_i\text{ are negative in }C.\hfill \end{array}`$
By induction on complexity, it can be easily seen that, for every subformula $`F`$ of a formula of $`C`$, whenever $`F`$ is false in $`M`$, so is it in $`M^{}`$. This extends from (sub)formulas to groups of $`C`$ and hence $`C`$ itself. Thus $`C`$ is false in $`M^{}`$ because it is false in $`M`$. But $`M^{}`$ does not distinguish between $`P_i`$ and $`Q_i`$ (any $`1in`$). This clearly implies that $`C`$ and $`B`$ have the same truth value in $`M^{}`$. That is, $`B`$ is false in $`M^{}`$, which is however impossible because $`B`$ is a tautology. From this contradiction we conclude that $`C`$ is a (normal binary) tautology.
Let $`\sigma `$ be a substitution such that $`A=\sigma (C)`$. Let $`\sigma ^{}`$ be a substitution such that, for each atom $`P`$ of $`C`$, $`\sigma ^{}(P)`$ is the result of replacing in $`\sigma (P)`$ each occurrence of each atom by a new atom in such a way that: (1) no atom occurs more than once in $`\sigma ^{}(P)`$, and (2) whenever $`PQ`$, no atom occurs in both $`\sigma ^{}(P)`$ and $`\sigma ^{}(Q)`$. As an instance of the tautological $`C`$, $`\sigma ^{}(C)`$ remains a tautology (this follows from Lemma 7.1 and Theorem 6.3). $`\sigma ^{}(C)`$ can also be easily seen to be a normal binary cirquent, because $`C`$ is so. Finally, with a little thought, $`A`$ can be seen to be an atomic-level instance of $`\sigma ^{}(C)`$. $`\mathrm{}`$
###### Lemma 7.3
The rules of mix, exchange, duplication, $``$-introduction and $``$-introduction preserve binarity and normal binarity in both top-down and bottom-up directions.
Proof. This is so because the above five rules in no way affect what atoms occur in a cirquent and how many times they occur. $`\mathrm{}`$
###### Lemma 7.4
Weakening preserves binarity and normal binarity in the bottom-up direction.
Proof. This is so because, in the bottom-up view, weakening can (delete but) never create any new occurrences of atoms. $`\mathrm{}`$
###### Theorem 7.5
A cirquent is provable in CL5 iff it is an instance of a binary tautology.
Proof. ($``$:) Consider an arbitrary cirquent $`A`$ provable in CL5. By induction on the height of its proof tree, we want to show that $`A`$ is an instance of a binary tautology.
The above is obvious when $`A`$ is an axiom.
Suppose now $`A`$ is derived by exchange from $`B`$. Let us just consider oformula exchange, with ogroup exchange being similar. By the induction hypothesis, $`B`$ is an instance of a binary tautology $`B^{}`$. Let $`A^{}`$ be the result of applying exchange to $`B^{}`$ “at the same place” as it was applied to $`B`$ when deriving $`A`$ from it, as illustrated in the following example:
Obviously $`A`$ will be an instance of $`A^{}`$. It remains to note that, by Lemmas 6.1 and 7.3, $`A^{}`$ is a binary tautology.
The rules of duplication, $``$-introduction and $``$-introduction can be handled in a similar way.
Next, suppose $`A`$ is derived from $`B`$ and $`C`$ by mix. By the induction hypothesis, $`B`$ and $`C`$ are instances of some binary tautologies $`B^{}`$ and $`C^{}`$, respectively. We may assume that no atom $`P`$ occurs in both $`B^{}`$ and $`C^{}`$, for otherwise, in one of the cirquents, rename $`P`$ into something different from everything else. Let $`A^{}`$ be the result of applying mix to $`B^{}`$ and $`C^{}`$. By Lemmas 6.1 and 7.3, $`A^{}`$ is a binary tautology. And, as in the cases of the other rules, it is evident that $`A`$ is an instance of $`A^{}`$.
Finally, suppose $`A`$ is derived from $`B`$ by weakening. If this is ogroup weakening, the conclusion is an instance of a binary tautology for the same reasons as in the case of exchange, duplication, $``$-introduction or $``$-introduction. Assume now we are dealing with pool weakening, so that $`A`$ is the result of inserting a new oformula $`F`$ into $`B`$. By the induction hypothesis, $`B`$ is an instance of a binary tautology $`B^{}`$. Let $`P`$ be an atom not occurring in $`B^{}`$. And let $`A^{}`$ be the result of applying weakening to $`B^{}`$ that inserts $`P`$ in the same place into $`B^{}`$ as the above application of weakening inserted $`F`$ into $`B`$ when deriving $`A`$. Obviously $`A^{}`$ inherits binarity from $`B^{}`$; by Lemma 6.1, it inherits from $`B^{}`$ tautologicity as well. And, for the same reasons as in all previous cases, $`A`$ is an instance of $`A^{}`$.
($``$:) Consider an arbitrary cirquent $`A`$ that is an instance of a binary tautology $`A^{}`$. In view of Lemma 7.1, it would suffice to show that CL5 proves $`A^{}`$. We construct a proof of $`A^{}`$, in the bottom-up fashion, as follows. Starting from $`A^{}`$, we keep applying conservative $``$-introduction and conservative $``$-introduction until we hit an essentially literal cirquent $`B`$. As in the proof of Theorem 6.3, such a cirquent $`B`$ is guaranteed to be a tautology, and $`A^{}`$ follows from it in ($``$). Furthermore, in view of Lemma 7.3, $`B`$ is in fact a binary tautology. Continuing as in the proof of Theorem 6.3, we apply to $`B`$ a series of weakenings and hit a tautological cirquent $`C`$ with no homeless oformulas, where every group only has two oformulas: $`P`$ and $`\neg P`$ for some atom $`P`$. By Lemma 7.4, $`C`$ remains binary. Our target cirquent $`A^{}`$ is thus derivable from $`C`$ in (W$``$). In the proof of Theorem 6.3 we next applied a series of contractions to separate shared oformulas. In the present case it suffices to use downward duplication (preceeded with ogroup exchange if necessary) instead of contraction: as it is easy to see, the binarity of $`C`$ implies that there are no shared oformulas in it except the cases when oformulas are shared by identical-content ogroups. Applying to $`C`$ a series of ogroup exchanges and downward duplications, as illustrated below, yields a cirquent $`D`$ that no longer has identical-content ogroups and hence no longer has any shared oformulas.
$`A^{}`$ is thus derivable from $`D`$ in (EWD$``$). In turn, as in the proof of Theorem 6.3, $`D`$ is provable in (AME). So, CL5 proves $`A^{}`$. $`\mathrm{}`$
## 8 Abstract resource semantics
### 8.1 Elaborating the abstract resource intuitions
Among the basic notions used in our presentation of abstract resource semantics is that of atomic resource. This is an undefined concept, and we can only point at some examples of what might be intuitively considered atomic resources. These can be a specified amount of money; electric power of a specified voltage and amperage; a specified task performed by a computer, such as providing Internet browsing capabilities; a specified number of bits of memory; the standard collection of tasks/duties of an employee in a given enterprise; the choice between a candy and an apple that a vending machine offers to whoever inserts a $1 bill into it; etc.
From atomic resources we will be building compound resources. Of course, whether a resource is considered atomic or thought of as a combination of some more basic resources depends on the degree of abstraction or encapsulation we choose in a given treatment. For instance, $2 can be treated as an atomic resource, but it can as well be understood as a combination of $1 and $1 — specifically, the combination $`\$1\$1`$, with $`\alpha \beta `$ generally being the resource having which intuitively means having both $`\alpha `$ and $`\beta `$. Similarly, a multiple-piece software package can be encapsulated and treated as an atomic resource, but in more subtle considerations it can be seen as a combination of the programs, data, etc. of which the package consists. And, had we extended our present approach to choice (“additive”) connectives, the earlier listed atomic resource providing a choice between a candy and an apple could be deatomized and understood as the choice conjunction $``$ of Candy and Apple.
When talking about resources, we always have two parties in mind: the resource provider and the resource user. Correspondingly, every entity that we call a resource comes in two flavors, depending on who is “responsible” for providing the resource. Suppose Victor received a salary of $3,000 in the morning, and paid a $3,000 mortgage bill in the afternoon. We are talking about the same resource of $3,000, but in one case it came to Victor as an income (input), while in the other case it was an expense (output). In the morning Victor was the user (and his employer the provider), while in the afternoon he was the provider (and the mortgage company the user) of the resource $3,000.
Or, imagine a car dealer selling a Toyota to his customer for $20,000. Two atomic resources can be seen involved in this transaction: the Toyota and the $20,000. The $20,000 is an income/input from the dealer’s perspective while an expense/output from the customer’s perspective; on the other hand, the Toyota is an expense/output for the dealer while an income/input for the customer.
Or, compare the two devices: a mini power generator that produces 100 watts of electric power, and a lamp with a 100-watt light bulb. Power is an output for the generator while an input for the lamp. In turn, the generator does not produce power for free: it takes/consumes certain input such as fuel in a specified quantity. Similarly, the lamp outputs light in exchange for power input.
The power generator and the lamp are also resources, in our treatment being compound ones unlike the atomic Power, Fuel or Light. Analyzing our intuitive concept of resources, one can notice that we are willing to call a resource anything that can be used — perhaps in combination with some other resources — to achieve certain goals. And “achieving a goal”, in turn, can be understood as nothing but obtaining/generating certain resources. $20,000 is a resource because it can be used — in combination with the resource Car dealer — to obtain the resource Toyota. Similarly, a generator and a lamp can help us — in combination with the resource Fuel — obtain Light. The component atomic resources of Generator are Fuel and Power, the former being an input as already noted, and the latter being an output. We will be using the term port as a common name for inputs and outputs. To indicate that a given port is an input, the name of the corresponding atomic resource will be prefixed with a “$``$”; the absence of such a prefix will mean that the port is an output. It should be noted that “$``$” is merely an indication of the input/output status of a port, and nothing more; it should not be mistaken for an operation on resources as, say, the later-defined operation $`\neg `$ is.
The sequence of all ports of a given compound resource we call its interface. Thus, the interface of the resource Generator is $`\text{Fuel, Power}`$, and the interface of Lamp is $`\text{Power, Light}`$. Generally, a compound resource may take any number of inputs and outputs. For example, Victor possessing both a generator and a lamp can be seen as possessing just one compound resource Generator$``$Lamp. The interface of this compound resource is then the concatenation of the interfaces of its two components, i.e. $``$$``$Fuel, Power, $``$Power, Light$``$. The first and the third ports of this interface stand for the resources that are expected to be provided by the user Victor, so they are inputs and hence come with a “$``$”; the second and the fourth ports, on the other hand, stand for the resources that Victor expects to receive, so they are outputs and hence come without a “$``$”. As for Victor (as opposed to the provider of his resource Generator$``$Lamp), he sees the same ports, but in negative colors: for him, the first and the third ports are outputs while the second and the fourth ports are inputs. To visualize Generator$``$Lamp as a one resource, it may be helpful for us to imagine a generator and a lamp mounted together on one common board, with the $``$Fuel port in the form of a pipe, the Power port in the form of a socket, the $``$Power port in the form of a plug, and the Light port in the form of a light bulb. Multiple inputs and/or outputs are very common. A TV set takes two inputs: power and cable. And a look at the back panel of a personal computer will show a whole array of what the engineers indeed call “ports”. This explains the choice of some of our terminology.
The list of all inputs and outputs is only a half of a full description of a compound resource. The other half is what we call the resource’s truth function. Formally, the latter is a function that returns a truth value — 0 or 1 — for each assignment of truth values to the ports of the compound resource. Intuitively, the value $`1`$ for a given resource — whether it be atomic or compound — means that the resource is “functioning”, or “doing its job”, or “keeping its promise”. In such cases we will simply say that the resource is true. And the value $`0`$, as expected, means false, i.e. not true. For instance, the value $`1`$ for the resource Power means that power is indeed generated/supplied, and the value $`0`$ means that there is no power supply. If Victor plugs the plug of his lamp into a (functioning, i.e. true) outlet, then the input port $``$Power of the resource Lamp becomes true; otherwise the port will probably remain false. We call assignments of truth values to the ports of a given compound resource situations. In these terms, the truth function $`F`$ of the compound resource tells us in which situations $`𝐬`$ the resource is considered true ($`F(𝐬)=1`$) and in which situations $`𝐬`$ it is false ($`F(𝐬)=0`$). Intuitively, such an $`F`$ can be seen as a description of the job that the resource is “supposed” (or “promises”) to perform. Specifically, the job/promise of the resource is to be true, i.e. ensure that no situations $`𝐬`$ with $`F(𝐬)=0`$ will arise. Jobs are not always done and promises not always kept however. So a resource, whether elementary or compound, may or may not be true.
Going back to the resource Generator, its job is to ensure that whenever there is fuel input, there is also power output. That is, whenever the (input) $``$Fuel port is true, so is the (output) Power port. In more intuitive but less precise terms, this job can be characterized as “turning fuel into power”. The following two tables contain full descriptions of the two resources Generator and Lamp, each table showing both the interface and the truth function (the rightmost column) of the corresponding compound resource:
As we see from the above figure, the only situation in which Generator is false, i.e. considered to have failed to do its job, is $`10`$ — the situation in which fuel was supplied to the generator but the latter did not produce power. While the situation $`01`$ is unlikely to occur with a real generator, Generator is considered true in it. For this is a situation in which the generator not only did not break its promise, but in fact did even more than promised. A customer who ended up receiving a $20,000-priced Toyota for something less than $20,000 (or for free) would hardly be upset and call the generous dealer a deal-breaker. The philosophy “it never hurts to do more than necessary” is inherent to our present approach, which is formalized in the requirement that the truth function of a resource should always be monotone, in the sense that changing $`0`$ to $`1`$ in an output or $`1`$ to $`0`$ in an input can never turn a true resource into a false one. Since Generator is true in situation $`00`$ (no input, no output), so should it be in $`01`$, for the generator “produced even more than expected”; an equally good reason why Generator is true in situation $`01`$ is that it is true in $`11`$, so that, in $`01`$, the generator “consumed even less than expected”.
We remember that the interface of the combination $`\alpha \beta `$ of resources is the concatenation of those of $`\alpha `$ and $`\beta `$. As for the truth function of $`\alpha \beta `$, it should account for the intuition that $`\alpha \beta `$ is considered to be doing its job iff both $`\alpha `$ and $`\beta `$ are doing their jobs. This can be seen from the following table for the resource Generator$``$Lamp:
The resource Generator$``$Lamp is false in situations $`1000`$, $`1001`$ and $`1011`$ because its Generator component failed to do its job: there was fuel input but no power output. The reason why Generator$``$Lamp is false in situations $`0010`$, $`0110`$ and $`1110`$ is that Lamp malfunctioned: there was power input for it but no light output. And in situation $`1010`$ Generator$``$Lamp is false because neither Generator nor Lamp kept the promise: there were both fuel input (for the generator) and power input (for the lamp), yet neither power output nor light output was generated.
The value of Generator$``$Lamp in a situation $`xyzt`$ only depends on the value $`u`$ of Generator in situation $`xy`$ and the value $`v`$ of Lamp in situation $`zt`$. So, such a $`xyzt`$ can be simply seen as $`uv`$, and the $``$$``$Fuel, Power$``$ and $``$$``$Power, Light$``$ parts of the interface seen as simply Generator and Lamp, respectively. That is, the two compound conjuncts of Generator$``$Lamp can be encapsulated and treated as atomic resources, which yields the following, simpler table:
Ignoring the minor — at least seemingly so — technical detail that columns are required to contain an additional bit of information indicating input/output status, our tables in the above style bear resemblance with those used in classical logic. Yet there is one crucial difference, which does not show itself in Figure 6 but catches the eye in Figure 5. In our tables, the same atom — such as Power in Figure 5 — may occur more than once, and the rows that assign different truth values to different occurrences of the same atom are as meaningful as any other rows. This is so because the expressionsPower ”, “Fuel ”, etc., as such, stand just for resource types, while particular occurrences of such expressions in a table (or a formula, cirquent etc.) stand for individual resources of those types. That is, we (again) deal with the necessity to differentiate between ports and oports, inputs and oinputs, etc. It is possible that the generator is producing power but the lamp is not receiving any: maybe Victor was not smart enough to insert the lamp’s input plug into the generator’s output socket, or it was a hot and sunny noontime, and he decided to use those 100 watts to feed his fan instead of the lamp. In the particular case of Figure 5, the two Power-labeled columns happen to be of different genders: one an output and the other an input. This, however, would not always be the case, and generally a table can contain any number of columns with identical labels of either gender. For example, the non-simplified table for Generator$``$Generator would have two (input) columns labeled $``$Fuel, and two (output) columns labeled Power. Such a 16-row table would not be at all the same as the 4-row one for the resource Generator seen in Figure 4. On the other hand, in classical logic, the truth table for any formula $`F`$ would be no different from that for $`FF`$, because for classical logic, which sees formulas as propositions or Boolean functions, $`F`$ and $`FF`$ are indistinguishable. Generator and Generator$``$Generator, on the other hand, are certainly not the same as resources. Provided that Victor has enough fuel to feed two generators, with Generator$``$Generator he can produce 200 watts of electricity while with just Generator only 100 watts. It may, however, happen that Victor decides to provide input for the first generator but not for the second one, so that rows with a $`1`$ in the first $``$Fuel column and a $`0`$ in the second $``$Fuel column cannot be dismissed as meaningless or impossible. And even the fully classical-looking table of Figure 6 would no longer look classical with Generator$``$Generator instead of Generator$``$Lamp. Such a table would still have 4 rows, while the classical table for $`PP`$ (atomic $`P`$) would only have 2 rows.
The meaning of the disjunction $`\alpha \beta `$ of resources must be easy to guess. The interface of $`\alpha \beta `$ is the same as that of $`\alpha \beta `$: just the interfaces of $`\alpha `$ and $`\beta `$ put together. As for the truth function, it corresponds to the intuition that the resource $`\alpha \beta `$ is considered to have failed its job if and only if so did both of its components $`\alpha `$ and $`\beta `$. Hence, say, the table for Generator$``$Lamp would differ from the one of Figure 5 in that the last column would have a $`0`$ only in one — the $`1010`$ — row. The job of Generator$``$Lamp is thus to generate either power output or light output (or both) whenever both fuel and power inputs are received.
The implicative combination $`\alpha \beta `$ of resources, in rough intuitive terms, can be characterized as the resource that “consumes $`\alpha `$ and produces $`\beta `$”. To see this, it would suffice to point out that our old friends Generator and Lamp (Figure 4) are nothing but Fuel$``$Power and Power$``$Light, respectively. Another intuitive characterization of $``$ is to say that this is a resource reduction operation. For example, the resource Lamp = Power$``$Light reduces (the task of generating) Light to (the task of generating) Power. Generally, the interface of $`\alpha \beta `$ is the concatenation of those of $`\alpha `$ and $`\beta `$, only, in the $`\alpha `$ part, the input/output status of each port is reversed. That is, in the antecedent the roles of the provider and the user are interchanged: the resource provider of $`\alpha \beta `$ acts as a provider in the $`\beta `$ part but as a user in the $`\alpha `$ part. Indeed, Generator provides Power, but uses Fuel. This is why Fuel — the atomic resource which, in isolation, is its own “output” — is an input rather than an output of Generator. The promise that the resource $`\alpha \beta `$ carries is to make $`\beta `$ true as long as $`\alpha `$ is true. In other words, to guarantee that either $`\beta `$ true, or $`\alpha `$ is false (or both). Imagine a car dealer who promised to his customer to sell to him a Toyota for a “to be negotiated” price. One way to keep this promise is, of course, to actually sell a Toyota. But what if the dealer has run out of cars by the time the customer arrives? A way out for the dealer is to request an unreasonable price that he believes the customer would never be able or willing to pay.
The above discussion makes it clear that $`\alpha \beta `$ is, in fact, a disjunction. Specifically, it is $`\neg \alpha \beta `$, where $`\neg \alpha `$, intuitively, is “the opposite of $`\alpha `$”: the interface of $`\neg \alpha `$ is that of $`\alpha `$ with all inputs turned into outputs and all outputs turned into inputs; and $`\neg \alpha `$ is true in exactly the situations in which $`\alpha `$ is false. Alternatively and equivalently, we can define $`\neg \alpha `$ as $`\alpha \mathrm{𝟎}`$; here $`\mathrm{𝟎}`$ is the empty-interface resource that is just (“always”) false, intuitively meaning a resource that no one can ever provide.
We close this subsection with a look at an intuitive example where $``$ takes a compound antecedent. Let us imagine that Victor has Fuel and Lamp. Could he generate light if these two resources are all of his possessions? Not really. What he needs is a generator. That is, Victor cannot (successfully) provide the resource Light, but with his Fuel, Lamp and some thought, he can provide the weaker resource Generator$``$Light, i.e. $`(`$Fuel$``$Power$`)`$ $``$ Light. Or can he? “Some thought” was not originally listed among the resources of Victor, and he might have a hard time exercising it should the example be more complex than it is. This is where CL5 can come to help. As it turns out, CL5 is exactly the logic that provides a systematic, sound and complete answer to the question on what and how Victor can generally do. Back to our present example, with perfect knowledge of CL5, Victor has a guarantee of success because CL5 proves
$$\text{Fuel}(\text{Power}\text{Light})\left((\text{Fuel}\text{Power})\text{Light}\right).$$
### 8.2 Resources and resource operations defined formally
Before we move any further, let us summarize, as formal definitions, the explanations given in the previous subsection. First of all, we agree that what we have been calling “atomic resources” are nothing but propositional letters, i.e. atoms of the language underlying cirquent calculus. This, of course, is some abuse of concepts because, strictly speaking, the atoms of the language are variables ranging over atomic resources rather than atomic resources as such. Similar terminological liberty extends to the concepts formally defined below as “ports”, (compound) “resources”, etc.
A port is $`P`$ or $`P`$, where $`P`$ is an atom called the type of the port. A port which is just an atom is said to be an output, and a port which is a “$``$”-prefixed atom is said to be an input. The input/output status of a port is said to be the gender of that port. The two genders input, output are said to be opposite.
An interface is a finite sequence of ports. A particular occurrence of a port (input, output) in an interface will be referred to as an oport (oinput, ooutput). As we did with oformulas in the context of cirquents, we usually refer to an oport by the name of the corresponding port (as in the phrase “the oport $`P`$”), even though different oports may be identical as ports.
Let $`I=X_1,\mathrm{},X_n`$ be an interface. A situation for $`I`$ is a function $`𝐬`$ of the type $`\{1,\mathrm{},n\}\{0,1\}`$. We identify such a function $`𝐬`$ with the bit string $`a_1\mathrm{}a_n`$, where $`a_1=𝐬(1),\mathrm{},a_n=𝐬(n)`$; we can also write $`𝐬(X_i)`$ instead of $`𝐬(i)`$, thinking of $`𝐬`$ as a function assigning truth values to oports. When $`𝐬(X_i)=1`$, we say that $`X_i`$ is true in $`𝐬`$, and when $`𝐬(X_i)=0`$, we say that $`X_i`$ is false in $`𝐬`$. We define the relation $`_I`$ on situations for $`I`$ by stipulating that $`𝐬_I𝐬^{}`$ iff, for each $`i\{1,\mathrm{},n\}`$, we have:
* if $`X_i`$ is an output, then $`𝐬(X_i)𝐬^{}(X_i)`$;
* if $`X_i`$ is an input, then $`𝐬^{}(X_i)𝐬(X_i)`$.
The relations $`_I,<_I`$ and $`>_I`$ have the expected meanings: $`𝐬_I𝐬^{}`$ iff $`𝐬^{}_I𝐬`$; $`𝐬<_I𝐬^{}`$ iff $`𝐬_I𝐬^{}`$ and $`𝐬𝐬^{}`$; and $`𝐬>_I𝐬^{}`$ iff $`𝐬^{}<_I𝐬`$.
###### Definition 8.1
An abstract resource — henceforth simply “resource” — is a pair $`\alpha =(𝐈nt^\alpha ,\text{Tfn}^\alpha )`$, where:
1. $`𝐈nt^\alpha `$, called the interface of $`\alpha `$, is an interface.
2. $`\text{Tfn}^\alpha `$, called the truth function of $`\alpha `$, is a function that sends every situation $`𝐬`$ for $`𝐈nt^\alpha `$ to $`0`$ or $`1`$, such that the following monotonicity condition is satisfied:
* Whenever $`𝐬_{𝐈nt^\alpha }𝐬^{}`$, we have $`\text{Tfn}^\alpha (𝐬)\text{Tfn}^\alpha (𝐬^{})`$.
When $`\text{Tfn}^\alpha (𝐬)=1`$ (resp. $`=0`$), we say that $`\alpha `$ is true (resp. false) in situation $`𝐬`$. “Output of $`\alpha `$”, “oport of $`\alpha `$”, “situation for $`\alpha `$”, etc. mean “output of $`𝐈nt^\alpha `$”, “oport of $`𝐈nt^\alpha `$”, “situation for $`𝐈nt^\alpha `$”, etc.
We need to agree on some notation used in the following definition of the basic resource operations. Let $`I,I_1,I_2`$ be interfaces. We write $`I`$ to mean the interface that is the same as $`I`$ only with the genders of all oports reversed; that is, $`I`$ is obtained from $`I`$ by deleting the prefix “$``$” wherever it was present and, simultaneously, adding such a prefix wherever it was absent. Next, $`I_1I_2`$ will stand for the concatenation of $`I_1`$ and $`I_2`$, i.e. the result of appending the oports of $`I_2`$ to those of $`I_1`$. Note that every situation (understood as a bit string) for $`I_1I_2`$ will have the form $`𝐬_1𝐬_2`$ (the concatenation of strings $`𝐬_1`$ and $`𝐬_2`$), where $`𝐬_1`$ is a situation for $`I_1`$ and $`𝐬_2`$ is a situation for $`I_2`$. The empty interface will be denoted by $``$. There is only one possible situation for it, and we denote that situation by $`ϵ`$. With situations understood as bit strings, $`ϵ`$ is thus the empty bit string.
###### Definition 8.2
Let $`\alpha `$, $`\alpha _1`$, $`\alpha _2`$ be resources. The operations $`\neg `$, $``$, $``$, $``$, $`\mathrm{𝟎}`$, $`\mathrm{𝟏}`$ are defined as follows:
1. Negation $`\neg \alpha `$:
* $`𝐈nt^{\neg \alpha }=𝐈nt^\alpha `$;
* For any situation $`𝐬`$ for $`\neg \alpha `$, $`\text{Tfn}^{\neg \alpha }(𝐬)=1`$ iff $`\text{Tfn}^\alpha (𝐬)=0`$.
2. Conjunction $`\alpha _1\alpha _2`$:
* $`𝐈nt^{\alpha _1\alpha _2}=𝐈nt^{\alpha _1}𝐈nt^{\alpha _2}`$;
* For any situation $`𝐬=𝐬_1𝐬_2`$ for $`\alpha _1\alpha _2`$, where $`𝐬_1`$ is a situation for $`\alpha _1`$ and $`𝐬_2`$ is a situation for $`\alpha _2`$, we have $`\text{Tfn}^{\alpha _1\alpha _2}(𝐬)=1`$ iff $`\text{Tfn}^{\alpha _1}(𝐬_1)=1`$ and $`\text{Tfn}^{\alpha _2}(𝐬_2)=1`$.
3. Disjunction $`\alpha _1\alpha _2`$:
* $`𝐈nt^{\alpha _1\alpha _2}=𝐈nt^{\alpha _1}𝐈nt^{\alpha _2}`$;
* For any situation $`𝐬=𝐬_1𝐬_2`$ for $`\alpha _1\alpha _2`$, where $`𝐬_1`$ is a situation for $`\alpha _1`$ and $`𝐬_2`$ is a situation for $`\alpha _2`$, we have $`\text{Tfn}^{\alpha _1\alpha _2}(𝐬)=1`$ iff $`\text{Tfn}^{\alpha _1}(𝐬_1)=1`$ or $`\text{Tfn}^{\alpha _2}(𝐬_2)=1`$.
4. Implication $`\alpha _1\alpha _2`$:
* $`𝐈nt^{\alpha _1\alpha _2}=(𝐈nt^{\alpha _1})𝐈nt^{\alpha _2}`$;
* For any situation $`𝐬=𝐬_1𝐬_2`$ for $`\alpha _1\alpha _2`$, where $`𝐬_1`$ is a situation for $`\alpha _1`$ and $`𝐬_2`$ is a situation for $`\alpha _2`$, we have $`\text{Tfn}^{\alpha _1\alpha _2}(𝐬)=1`$ iff $`\text{Tfn}^{\alpha _1}(𝐬_1)=0`$ or $`\text{Tfn}^{\alpha _2}(𝐬_2)=1`$.
5. Empty-interface (constant) resources $`\mathrm{𝟎}`$ and $`\mathrm{𝟏}`$:
* $`𝐈nt^\mathrm{𝟎}=𝐈nt^\mathrm{𝟏}=`$;
* $`\text{Tfn}^\mathrm{𝟎}(ϵ)=0`$; $`\text{Tfn}^\mathrm{𝟏}(ϵ)=1`$.
For safety, we need to verify that the above operations (with $`\mathrm{𝟎}`$ and $`\mathrm{𝟏}`$ seen as $`0`$-ary operations) are indeed operations on resources — that is, that they do not violate the monotonicity condition of Definition 8.1. The truth function of $`\mathrm{𝟎}`$ (and similarly for $`\mathrm{𝟏}`$) is trivially monotone, because there is only one situation for the empty interface. Of the other operations, it would be sufficient to only consider $``$ for, in view of Theorem 8.3, just as this is the case in classical logic, $`\neg `$, $``$ and $``$ are definable in terms of $``$ and $`\mathrm{𝟎}`$.
Consider any two situations $`𝐬,𝐬^{}`$ for $`\alpha _1\alpha _2`$ such that $`𝐬_{𝐈nt^{\alpha _1\alpha _2}}𝐬^{}`$. Our goal is to show that if $`\alpha _1\alpha _2`$ is false in $`𝐬^{}`$, then so is it in $`𝐬`$. Let $`𝐬_1,𝐬_1^{}`$ be the situations for $`\alpha _1`$, and $`𝐬_2,𝐬_2^{}`$ the situations for $`\alpha _2`$, such that $`𝐬=𝐬_1𝐬_2`$ and $`𝐬^{}=𝐬_1^{}𝐬_2^{}`$. Then the condition $`𝐬_{𝐈nt^{\alpha _1\alpha _2}}𝐬^{}`$ implies that $`𝐬_1_{𝐈nt^{\alpha _1}}𝐬_1^{}`$ and $`𝐬_2_{𝐈nt^{\alpha _2}}𝐬_2^{}`$. Assume $`\alpha _1\alpha _2`$ is false in $`𝐬^{}`$. This means that $`\text{Tfn}^{\alpha _1}(𝐬_1^{})=1`$ and $`\text{Tfn}^{\alpha _2}(𝐬_2^{})=0`$. From $`\text{Tfn}^{\alpha _1}(𝐬_1^{})=1`$ and $`𝐬_1_{𝐈nt^{\alpha _1}}𝐬_1^{}`$ the monotonicity of $`\text{Tfn}^{\alpha _1}`$ implies $`\text{Tfn}^{\alpha _1}(𝐬_1)=1`$. Similarly, from $`\text{Tfn}^{\alpha _2}(𝐬_2^{})=0`$ and $`𝐬_2_{𝐈nt^{\alpha _2}}𝐬_2^{}`$ the monotonicity of $`\text{Tfn}^{\alpha _2}`$ implies $`\text{Tfn}^{\alpha _2}(𝐬_2)=0`$. This means that $`\alpha _1\alpha _2`$ is false in $`𝐬=𝐬_1𝐬_2`$. Done.
The following theorem can be verified by a routine examination of the relevant definitions, left as an exercise for the reader:
###### Theorem 8.3
For any resources $`\alpha ,\alpha _1,\alpha _2`$, the following equalities hold:
$$\begin{array}{ccc}\hfill \neg \alpha & =& \alpha \mathrm{𝟎};\hfill \\ \hfill \neg \neg \alpha & =& \alpha ;\hfill \\ \hfill \neg (\alpha _1\alpha _2)& =& \neg \alpha _1\neg \alpha _2;\hfill \\ \hfill \neg (\alpha _1\alpha _2)& =& \neg \alpha _1\neg \alpha _2;\hfill \\ \hfill \alpha _1\alpha _2& =& \neg \alpha _1\alpha _2.\hfill \end{array}$$
### 8.3 Formulas as resources
By a situation for a formula $`F`$ we mean an assignment $`𝐬`$ of truth values ($`0`$ or $`1`$) to its oatoms. Such an $`𝐬`$ will as well be understood as a situation for any osubformula $`G`$ of $`F`$ by mechanically restricting its domain to the oatoms of $`G`$. Note the difference between a situation and a classical model. The latter is an assignment of truth values to atoms rather than oatoms. So, say, (the relevant part of) a classical model for $`\neg P(PP)`$ would only have to assign a value to $`P`$, while a situation for this formula would have to list three — not necessarily identical — values for the three oatoms of the formula. Classical models can be viewed as special cases of situations that assign identical truth values to oatoms that are identical as atoms. Disregarding this difference, the truth status of a formula $`F`$ in a situation $`𝐬`$ for $`F`$ is determined in the “standard” way by recursion on (the occurrences of) its subformulas. That is:
* $`\neg P`$ is true in $`𝐬`$ iff $`P`$ is false in $`𝐬`$; here and later “false”, as always, means “not true”.
* $`GH`$ is true in $`𝐬`$ iff so are both $`G`$ and $`H`$;
* $`GH`$ is true in $`𝐬`$ iff so is $`G`$ or $`H`$ (or both).
While the above definition looks like exactly the classical definition of truth, once again we emphasize the implicit and important difference: in the present definition, $`\neg P`$, $`GH`$, $`GH`$ and their subformulas are particular occurrences of (sub)formulas rather than formulas as such. This makes it possible for, say, $`P\neg P`$ to be true, which will be the case when the first oatom $`P`$ is true and the second oatom $`P`$ is false.
Notice that, graphically, literals and ports are almost the same, with the difference that ports take the prefix “$``$” where literals take the prefix “$`\neg `$”. Since we often convert between literals and ports, here we introduce two functions Port and Literal, with Port for literal-to-port conversion and Literal for port-to-literal conversion. Specifically, for any atom $`P`$, we have $`\text{Port}(P)=\text{Literal}(P)=P`$, $`\text{Port}(\neg P)=P`$ and $`\text{Literal}(P)=\neg P`$.
As in the previous subsection, a situation $`𝐬`$ for a formula $`F`$ can and will be understood as a bit string — specifically, a bit string of length $`n`$, where $`n`$ is the number of oatoms of $`F`$. Then the same $`𝐬`$ can also be considered a situation for any resource whose interface has $`n`$ ports. The second clause of the following definition implicitly relies on this seeing no difference between situations for formulas and situations for resources.
###### Definition 8.4
With each formula $`F`$ we associate the resource $`F^{\mathrm{}}`$, called the resource represented by $`F`$, defined as follows:
* $`𝐈nt^F^{\mathrm{}}=\text{Port}(L_1),\mathrm{},\text{Port}(L_n)`$, where $`L_1,\mathrm{},L_n`$ is the sequence of the oliterals of $`F`$ listed in the order of their appearance in $`F`$.
* For every situation $`𝐬`$ for $`F^{\mathrm{}}`$, $`\text{Tfn}^F^{\mathrm{}}(𝐬)=1`$ iff $`F`$ is true in $`𝐬`$.
According to the following theorem, respects the meanings of $`\neg ,,`$ (and, as we may guess, also of $``$ if the latter was officially allowed in formulas) as operations on resources. Hence, with propositional letters understood as representing atomic resources, $`F^{\mathrm{}}`$ is indeed “the resource represented by $`F`$”. The theorem can be verified by a straightforward analysis of the relevant definitions, so we state it without a proof:
###### Theorem 8.5
For any formulas $`F,G`$, we have $`(\neg F)^{\mathrm{}}=\neg (F^{\mathrm{}})`$, $`(FG)^{\mathrm{}}=F^{\mathrm{}}G^{\mathrm{}}`$ and $`(FG)^{\mathrm{}}=F^{\mathrm{}}G^{\mathrm{}}`$.
While every formula corresponds to a resource, vice versa does not hold. For example, it is obvious that no formula in our present sense does represent the resource $`\mathrm{𝟎}`$ or $`\mathrm{𝟏}`$, for there are no atomless formulas. This is not a serious issue of course, for we could painlessly add $`\mathrm{𝟎}`$ and $`\mathrm{𝟏}`$ (perhaps using the symbols $``$ and $``$ instead) to our formal language, laziness being the only reason for not having done so.
The resource whose interface is $`P,Q,R`$ and which is true when so is $`Q`$ and at least one of $`P,R`$ cannot be expressed by a formula, either. For example, $`Q(PR)`$ would not fit the bill because of the wrong order of its atoms. Again, this is not a “serious” problem, for we may be willing to not distinguish between permutationally equivalent resources — resources that, informally speaking, only differ in the order in which their interfaces list the ports.
There are, however, really serious reasons that make it impossible to capture all resources with formulas, reasons that essentially call for switching to non-traditional means of expression such as cirquents. The closure of atomic resources under $`\neg `$, $``$, $``$, $``$ and any other operations in a similar style does not yield the class of all resources, and this is so because, vaguely speaking, such operations do not allow us to account for the possibility of resource sharing. An analysis of the following example can make this point clear.
###### Example 8.6
Let $`\beta `$ be the resource defined by:
* $`𝐈nt^\beta =P,Q,R`$;
* $`\text{Tfn}^\beta (𝐬)=1`$ iff at least two of the three bits of $`𝐬`$ are $`1`$s.
The “two out of three” Boolean function in classical logic would be expressed by the formula $`F=(PQ)(PR)(QR)`$. $`F^{\mathrm{}}`$, however, is not $`\beta `$, for the former has six oports rather than three. We need yet do not have means to indicate that, say, the two occurrences of $`P`$ in $`F`$ stand for the same individual resource rather than two different resources of the same type. In other words, we need yet do not have means to indicate that $`P`$ is shared between the two subresources $`PQ`$ and $`PR`$. It should not be hard for the reader to convince himself or herself that generally for no $`\neg ,,`$-combination $`F`$ of $`P,Q,R`$ do we have $`F^{\mathrm{}}=\beta `$.
Thus, while the collection $`(\neg ,,)`$ enjoys what is called functional completeness in classical logic, in no reasonable sense is the expressive power of this collection complete when it comes to resource semantics. Nor is there any easy remedy such as adding some extra connectives to the language. What we need is to go substantially beyond the traditional formalisms of logic. As we are going to see in the next subsection (Theorem 8.8), the formalism of cirquent calculus turns out to be sufficient.
### 8.4 Cirquents as resources
Let $`C`$ be a cirquent. By a situation for $`C`$ we mean an assignment $`𝐬`$ of truth values ($`0`$ or $`1`$) to the oatoms of $`C`$. Understood as a bit string, such an $`𝐬`$ is the concatenation of situations for the oformulas of $`C`$. Then an oformula of the cirquent is considered true or false in $`𝐬`$ if it is so in the corresponding substring of $`𝐬`$. For instance, if the pool of $`C`$ is $`G,H`$ where $`G`$ has 3 oatoms and $`H`$ has 2 oatoms, then $`G`$ is true in situation $`10111`$ iff it is true (in the sense of the previous subsection) in $`101`$, and $`H`$ is true in $`10111`$ iff it is true in $`11`$. Next, we consider a group $`\mathrm{\Gamma }`$ of $`C`$ true in a given situation $`𝐬`$ for $`C`$ iff at least one of the oformulas of $`\mathrm{\Gamma }`$ is true. Notice that, unlike oformulas, the truth values of two different ogroups that are identical as groups would always be the same. So, truth can be considered a property of groups rather than ogroups. Finally, $`C`$ is true in $`𝐬`$ iff all of its groups are so.
###### Definition 8.7
With each cirquent $`C`$ we associate the resource $`C^{\mathrm{}}`$, called the resource represented by $`C`$, defined as follows:
1. $`𝐈nt^C^{\mathrm{}}=\text{Port}(L_1),\mathrm{},\text{Port}(L_n)`$, where $`L_1,\mathrm{},L_n`$ is the sequence of the oliterals of $`C`$ listed in the order of their appearance in $`C`$.
2. For every situation $`𝐬`$ for $`C^{\mathrm{}}`$, $`\text{Tfn}^C^{\mathrm{}}(𝐬)=1`$ iff $`C`$ is true in $`𝐬`$.
###### Theorem 8.8
For every resource $`\alpha `$ there is a cirquent $`C`$ — in fact a literal one — which represents $`\alpha `$, i.e. such that $`C^{\mathrm{}}=\alpha `$.
Proof. Consider an arbitrary resource $`\alpha `$ with $`𝐈nt^\alpha =X_1,\mathrm{},X_n`$. We construct a corresponding cirquent $`C`$ as follows. The pool of $`C`$ is $`L_1,\mathrm{},L_n=\text{Literal}(X_1),`$ $`\mathrm{},\text{Literal}(X_n)`$. In view of condition 1 of Definition 8.7, it is already clear that $`𝐈nt^C^{\mathrm{}}=𝐈nt^\alpha `$. We now need to define the structure of $`C`$ and show that we also have $`\text{Tfn}^C^{\mathrm{}}=\text{Tfn}^\alpha `$.
Let us call a situation $`𝐬`$ for $`\alpha `$ critical iff $`\alpha `$ is false in $`𝐬`$ and, for every situation $`𝐬^{}`$ for $`\alpha `$ with $`𝐬<_{𝐈nt^\alpha }𝐬^{}`$, $`\alpha `$ is true in $`𝐬^{}`$. Let $`𝐬_1,\mathrm{},𝐬_k`$ be a list of all critical situations for $`\alpha `$. Seeing them as bit strings, these situations, of course, are also situations for $`C`$ (no matter what the structure of $`C`$ is). For each $`i\{1,\mathrm{},k\}`$, we define the group $`\mathrm{\Gamma }_i`$ by
$$\mathrm{\Gamma }_i=\{L_j|1jn,L_j\text{ is false in }𝐬_i\}.$$
Now, we define the structure of $`C`$ to be $`\mathrm{\Gamma }_1,\mathrm{},\mathrm{\Gamma }_k`$. To see that $`\text{Tfn}^C^{\mathrm{}}=\text{Tfn}^\alpha `$, consider any situation $`𝐬`$ for $`C`$ (and hence for $`C^{\mathrm{}}`$ and $`\alpha `$).
Suppose $`\text{Tfn}^C^{\mathrm{}}(𝐬)=0`$, by condition 2 of Definition 8.7 meaning that $`C`$ is false in $`𝐬`$. Then there is a group $`\mathrm{\Gamma }_i`$ ($`1ik`$) which is false in $`𝐬`$, and this, in turn, means that every $`L_j`$ which is false in $`𝐬_i`$ (i.e. every oformula of $`\mathrm{\Gamma }_i`$) is false in $`𝐬`$. In other words, $`𝐬_{𝐈nt^\alpha }𝐬_i`$. But $`\alpha `$ is false in $`𝐬_i`$ because $`𝐬_i`$ is critical. Hence, by the monotonicity of $`\text{Tfn}^\alpha `$, we have $`\text{Tfn}^\alpha (𝐬)=0`$.
Now suppose $`\text{Tfn}^\alpha (𝐬)=0`$. We may assume that $`𝐬`$ is critical, for otherwise replace it with a critical situation $`𝐬^{}`$ such that $`𝐬_{𝐈nt^\alpha }𝐬^{}`$ (by monotonicity, such an $`𝐬^{}`$ is guaranteed to exist). Since $`𝐬`$ is critical, $`𝐬=𝐬_i`$ for some $`1ik`$. Remembering now how the group $`\mathrm{\Gamma }_i`$ was chosen, it is clear that $`\mathrm{\Gamma }_i`$ is false in $`𝐬`$, whence $`C`$ is false in $`𝐬`$, i.e. $`\text{Tfn}^C^{\mathrm{}}(𝐬)=0`$. $`\mathrm{}`$
###### Example 8.9
Applying the construction from our proof of Theorem 8.8 to the resource $`\beta `$ of Example 8.6, $`\beta `$ is represented by the following cirquent:
### 8.5 Resource-semantical validity
Every logical semantics has a concept of validity, and so does our abstract resource semantics. In classical semantics valid formulas are called tautologies, and in abstract resource semantics they will be called trivialities. Trivialities and tautologies are similar in many respects. Asserting a tautology means “asserting nothing”. Likewise, possessing (or providing) a triviality means “possessing (or providing) nothing”. Of course, the word “nothing” has a negative flavor. But there is a positive side as well. If possessing $`\alpha `$ amounts to possessing nothing, this means that, in fact, everyone possesses $`\alpha `$. Tautologicity is a guarantee of truth, and can be eventually used in finding true (and non-tautological) statements. Similarly, triviality is a guarantee of success in providing resources, and can be eventually used in finding what (nontrivial) resources can be generated. The formula given at the end of Subsection 8.1 is a triviality. And it is exactly this fact that allowed us to be confident that Victor, who possesses the nontrivial resources represented by the antecedent of the formula, can successfully generate the nontrivial resource represented by the consequent.
The simplest examples of trivial resources are those of the form $`\alpha \alpha `$. What makes such a resource trivial is that it merely returns back what it takes. Even the poorest person in the world would be able to provide the resource $3,000$``$$3,000, i.e. pay $3,000 if he or she receives $3,000. And it would not take a power plant to support the resource Power$``$Power : assuming here that the input port comes to the resource provider in the form of a socket and the output port in the form of a plug, enough to just insert the plug into the socket. A similar trick works with Power$``$Power$``$Power, providing which not only does not require any spending, but in fact can be even done with a benefit. While $`\alpha \alpha \alpha `$ is a triviality, $`\alpha \alpha \alpha `$ is not. Only someone having $3,000 of his or her own would be able to pay two $3,000 bills while only receiving an income of $3,000. And a duplex power adapter that seemingly turns Power into Power$``$Power, does so only seemingly: while the voltage being accurately reproduced, the amperage in the two output (o)ports of the adapter would be inevitably lower than in its input port, as 100w cannot be converted into 200w “for free”.
What makes $`PP`$, i.e. $`\neg PP`$ valid in classical logic is that the truth table for this formula has a $`1`$ in both of its two rows. The resource-semantical table for $`\neg PP`$, however, has a $`0`$ in one — $`10`$ — of its four rows. Generally, no formula in our present sense would have only $`1`$s in its resource-semantical table. Yet Victor, acting as a provider of $`\neg PP`$, has a way to make sure that the falsifying situation $`10`$ never occurs. This way is to allocate the input oport/resource $`P`$ to the output oport/resource $`P`$. If here $`P`$ is Power with oport $`P`$ given in the form of a socket and oport $`P`$ in the form of a plug, the physical meaning of allocation, as noted above, could be inserting the plug into the socket, which guarantees that whenever the socket has power supply, so does the plug. If $`P`$ is computer memory, then allocation may literally mean allocation. If $`P`$ is $3,000, then allocation probably means redirection or mutual cancellation. If $`P`$ is Light, allocating probably means just letting the user utilize the light generated by himself or herself. For yet more diversity, assume $`P`$ is the resource Chess whose promise is to play the game of chess white and win. Then providing the resource $`\neg \text{Chess}\text{Chess}`$ means to play chess on two boards and win on at least one of them; specifically, to play white on the right (“output”) board, and play black on the left board, which is considered an “input” and hence the roles — colors — of the two players are interchanged. In this case, a way to “allocate” $`\text{Chess}`$ and Chess to each other would be to mimic on one board the moves made by the adversary on the other board, and vice versa. Obviously this strategy, amounting to having the user/adversary play against himself, guarantees that whenever the game is won by white (i.e. the promise of Chess kept) on the left board, so is it on the right board.
The concept of allocation thus allows various, rather different particular interpretations. Attempting to formalize those meanings (if so, which one?) is beyond the scope and dignity of our resource semantics which is meant to be an abstract, general-purpose formal framework. Essentially we treat allocation as a basic, undefined concept, just as we treat the concepts of atomic resources, inputs, outputs, ports, or truth (for atomic resources). Our semantics only focuses on the conditions that all intended interpretations of its basic concepts would satisfy. There are three such conditions pertaining to allocation:
Only same-type and opposite-gender oports can be allocated to each other.
Allocations should be “monogamous”, in the sense that no oport can be involved in more than one allocation.
Each allocation guarantees that whenever its input oport is true, so is the output oport.
Notice the egalitarian view of the genders implied by our phrases such as “allocating to each other”: if oinput $`X`$ is allocated to ooutput $`Y`$, we also say that ooutput $`Y`$ is allocated to oinput $`X`$. This is in concordance with the fact that the intuitive distinction between input and output is often blurred and, when modeling real-life situations, initial decisions regarding whether a certain port should be listed as an input or an output can be arbitrary. For instance, when representing a computer as a compound resource, the (resource provided by the) monitor would most likely become an output and the keyboard an input; the fates of the modem or the floppy disk drive, however, are not just as clear. And, in our earlier example with $`\neg \text{Chess}\text{Chess}`$, there were hardly any reasons in favor of treating the left rather than the right board as an “input”. In each case, however, we still deal with two opposite — provider’s and user’s — perspectives of the same resource, clearly dictating whether any given pair of same-type oports should have the same or opposite genders, even if there is flexibility in choosing what particular genders they have. So, under any choice, the resource modeling what was originally modeled as $`\neg \text{Chess}\text{Chess}`$ would have two same-type and opposite-gender oports. Note, however, that if genders are reversed, so should be the meanings of the truth values and all aspects of the provider/user perspectives of the corresponding atomic resources. For example, in the case of Chess, such a reversal would mean considering Chess (or whatever new name we use for it) true when black rather than white wins the game, with the role of the provider of this resource now being that of the black rather than the white player. This interchangeability and symmetry between input and output explains why many-inputs-to-one-output allocations would generally be just as inadmissible or impossible as many-outputs-to-one-input allocations. The fact that to-be-allocated resources more often than not would involve elements of both “income” and “expense” (encapsulated Generator being a simplest example), with some thought, can also be seen to speak against the admissibility of any dual standards for inputs and outputs when it comes to allocations.
Back to the topic of validity, triviality of a resource is understood as a possibility to establish allocations that rule out all situations in which the resource is false. Assuming that making allocations only requires intellect and no other, “external” resources, and that Victor does possess intellect, the triviality of a resource indeed means a guarantee that Victor can always successfully provide the resource. The following formal definition summarizes these intuitions:
###### Definition 8.10
Let $`\alpha `$ be a resource.
1. An allocation for $`\alpha `$ is a pair $`(X,Y)`$, where $`X`$ is an oinput of $`\alpha `$, and $`Y`$ is an ooutput of $`\alpha `$ of the same type as $`X`$. We say that the allocation $`(X,Y)`$ utilizes $`X`$ and $`Y`$.
2. An arrangement for $`\alpha `$ is a set of allocations for $`\alpha `$.
3. An arrangement $`𝒜`$ for $`\alpha `$ is said to be monogamous iff no oport of $`\alpha `$ is utilized by more than one allocation of $`𝒜`$.
4. Let $`𝐬`$ be a situation for $`\alpha `$, and $`𝒜`$ an arrangement for $`\alpha `$. We say that $`𝐬`$ is consistent with $`𝒜`$ iff, whenever $`(X,Y)𝒜`$, $`𝐬(X)𝐬(Y)`$.
5. We say that an arrangement $`𝒜`$ for $`\alpha `$ is trivializing (for $`\alpha `$) iff $`\alpha `$ is true in every situation for $`\alpha `$ consistent with $`𝒜`$.
6. We say that resource $`\alpha `$ is trivial, or is a triviality, iff there is a monogamous trivializing arrangement for it.
7. We say that a cirquent (or formula) $`C`$ is trivial, or is a triviality, iff the resource represented by such a cirquent (formula) is trivial. That is, $`C`$ is trivial iff there is a monogamous trivializing arrangement for $`C^{\mathrm{}}`$.
###### Exercise 8.11
Show that dropping the monogamicity requirement in the above definition of triviality yields the ordinary concept of tautologicity for cirquents and formulas. That is, a cirquent or formula $`C`$ is a tautology in the sense of Section 6 iff there is a (not necessarily monogamous) trivializing arrangement for $`C^{\mathrm{}}`$.
Hint: Let us say that an arrangement $`𝒜`$ for a given resource is greedy iff $`𝒜`$ contains every possible allocation for that resource. Exploit the fact that, in a sense, classical models for $`C`$ are nothing but situations consistent with the greedy arrangement for $`C^{\mathrm{}}`$.
###### Lemma 8.12
A cirquent (or formula) is trivial iff it is an instance of a binary tautology.
Proof. Formulas are special cases of cirquents, so let us limit our attention to cirquents in general. Consider an arbitrary cirquent $`C`$. For readability, here we will be identifying the literals of $`C`$ with the corresponding ports of $`C^{\mathrm{}}`$, so that any arrangement for $`C^{\mathrm{}}`$ can be seen as a set of pairs $`(\neg P,P)`$ of oliterals of $`C`$.
($``$): Assume $`C`$ is trivial. Let $`𝒜`$ be a monogamous trivializing arrangement for $`C^{\mathrm{}}`$. Whenever two oatoms of $`C`$ are in literals utilized by the same allocation of $`𝒜`$, we call such oatoms coupled. We replace oatoms of $`C`$ by new atoms in such a way that coupled oatoms are replaced by the same atom, and any pair of non-coupled atoms end up being replaced by different atoms. Call the resulting cirquent $`D`$. Of course, $`D`$ is a normal binary cirquent and $`C`$ is an (atomic-level) instance of it. What remains to show is that $`D`$ is a tautology. Suppose it is not. Then there must be a classical model $`M`$ in which $`D`$ is false. Let $`𝐬`$ be the situation for $`C`$ that sends every oatom $`P`$ of $`C`$ to the same truth value as $`M`$ sends the atom that replaced $`P`$ when obtaining $`D`$ from $`C`$. Clearly $`𝐬`$ is consistent with $`𝒜`$. It is also obvious that every osubformula of every oformula of $`C`$ has the same truth value in $`𝐬`$ as the corresponding osubformula of $`D`$ in $`M`$. We conclude that $`C`$ — and hence $`C^{\mathrm{}}`$ — is false in $`𝐬`$. This contradicts our assumption that $`𝒜`$ is trivializing for $`C^{\mathrm{}}`$.
($``$): Assume $`C`$ is an instance of a binary tautology $`D`$. In view of Lemma 7.2, we may assume that $`D`$ is normal and $`C`$ is an atomic-level instance of it. Let $`𝒜`$ be the arrangement for $`C^{\mathrm{}}`$ that contains an allocation $`(\neg P,P)`$ if and only if the oatoms of the corresponding two oliterals of $`D`$ are identical as atoms (“corresponding oliterals of $`D`$” means the ones that were replaced by $`\neg P`$ and $`P`$ when obtaining $`C`$ from $`D`$). Clearly $`𝒜`$ is monogamous. We want to verify that $`𝒜`$ is trivializing for $`C^{\mathrm{}}`$. Suppose it is not. Then there is a situation $`𝐬`$ for $`C^{\mathrm{}}`$ consistent with $`𝒜`$ such that $`C^{\mathrm{}}`$ is false in $`𝐬`$. In view of the monotonicity of $`\text{Tfn}^C^{\mathrm{}}`$, we may assume that whenever two oatoms of $`C`$ are coupled (“coupled” in the same sense as in the previous paragraph), $`𝐬`$ assigns identical truth values to them. Let then $`M`$ be a classical model such that, for any atom $`P`$ of $`D`$, the truth value of $`P`$ in $`M`$ is the same as the truth value of the corresponding oatom(s) of $`C`$ in $`𝐬`$ (again, “corresponding oatom(s)” mean(s) the one(s) that replaced $`P`$ when obtaining $`C`$ from $`D`$). It is not hard to see that then $`D`$ has the same truth value — particularly, the value false — in $`M`$ as $`C`$ has in $`𝐬`$. So, $`D`$ is not a tautology, which is a contradiction. $`\mathrm{}`$
The following theorem, establishing the soundness and completeness of CL5 with respect to our abstract resource semantics, is an immediate corollary of Theorem 7.5 and Lemma 8.12:
###### Theorem 8.13
A cirquent (or formula) is provable in CL5 iff it is trivial.
###### Remark 8.14
At the end of Subsection 8.1 we promised that CL5 would provide an ultimate answer to the question what and how Victor can achieve purely by means of smart resource management. Theorem 8.13 directly pertains only to the what part of this question: Victor can successfully provide the resource represented by a formula or cirquent $`C`$ (i.e. $`C`$ is trivial) iff CL5$`C`$. Specifically, he can achieve his success by just setting up a monogamous trivializing arrangement for $`C^{\mathrm{}}`$ (making allocations is exactly what “resource management” means). So, the how part of the above question reduces to finding such an arrangement. From our proof of the $`(`$) part of Lemma 8.12 one can see that finding a monogamous trivializing arrangement for $`C^{\mathrm{}}`$, in turn, essentially means nothing but finding a normal binary cirquent $`D`$ such that $`C`$ is an atomic-level instance of $`D`$. Our proof of the soundness part of Theorem 7.5 implicitly provides an easy way to turn a CL5-proof of $`C`$ into a CL5-proof of a binary (and, in fact, normal) tautological cirquent $`C^{}`$ such that $`C`$ is an instance of $`C^{}`$. $`C`$ is not necessarily an atomic-level instance of such a $`C^{}`$ though. No problem: our constructive proof of Lemma 7.2 shows how to turn $`C^{}`$ into a(nother) normal binary tautology $`D`$ such that $`C`$ is an atomic-level instance of $`D`$. Putting the above steps and observations together yields a polynomial-time algorithm that turns a CL5-proof of an arbitrary cirquent $`C`$ into a monogamous trivializing arrangement for $`C^{\mathrm{}}`$.
## 9 Cirquent calculus and computability logic
The claim that the semantics of computability logic is a semantics of resources has been explicitly or implicitly present in every work on CL. The abstract resource semantics introduced in this paper makes the same claim. While the two semantics are far from being the same, there is no contradiction or competition here. This is so not only due to the forthcoming Theorem 9.2 which, in conjunction with Theorem 8.13, implies that the two semantics validate the same principles. Abstract resource semantics, as noted, is a general-purpose framework, with its basic notions such as (atomic) resources, their truth values, or allocations being open to various specific interpretations. Computability logic offers one of such interpretations, and its semantics can be seen as a materialization of abstract resource semantics. A “resource” in CL has a very specific meaning as briefly explained in Section 1. It is a derived concept defined in terms of some other, more basic and subtler-level entities (specifically, games). The same can be said about the other basic semantical concepts: in CL, to what we call “successfully providing the resource” (truth value $`1`$) corresponds winning the game; to our “allocations” correspond copy-cat subroutines in game-playing strategies; etc. While being more special, however, the semantics of computability logic still turns out to be general enough to invalidate anything that abstract resource semantics does.
The earlier-promised soundness and completeness of CL5 with respect to the semantics of computability logic can be proven purely syntactically, based on the fact of the soundness and completeness of system CL2 known from . The propositional language of the latter is considerably more expressive than the one in which the formulas of CL5 are written. One difference is that the language of CL2 includes the choice operators $``$ and $``$. The syntactic behavior of these operators, as noted in Section 1, is somewhat reminiscent of that of the additive operators of linear logic, just as $`,,`$, called parallel operators in CL, are relatives of the multiplicative operators of linear logic. Another significant difference is that the language of CL2 has two sorts of atoms: general and elementary.<sup>5</sup><sup>5</sup>5In computability logic, general atoms stand for computational problems (games) of arbitrary degrees of interactivity, while the interpretations of elementary atoms are limited to a special, “moveless” sort of games, i.e. computational problems with no interaction at all. Such problems are in fact nothing but propositions or predicates in the standard sense. General atoms are the same as in the formulas of CL5, and we continue using the uppercase $`P,Q,R,S`$ as metavariables for them. As for elementary atoms, they are foreign to CL5. Elementary atoms are divided into logical and non-logical. There are two logical elementary atoms $``$ and $``$, and infinitely many non-logical elementary atoms, for which we will be using the lowercase $`p,q,r,s`$ as metavariables. There is also a minor difference: in the official version of CL2 given in , formulas are allowed to contain $``$, and also the scope of $`\neg `$ is not limited to atomic formulas. We may safely pretend that this difference does not exist: whether it be classical, linear or computability logic, $`FG`$ is virtually the same as $`\neg FG`$ and hence can be understood just as an abbreviation. For similar reasons, it is irrelevant whether $`\neg `$ is allowed to be applied only to atoms or compound formulas as well. Another insignificant difference is that the language of CL2 officially treats $``$ and $``$ (as well as $``$ and $``$) as variable-arity operators, as opposed to the strictly binary treatment chosen in the present paper. In view of the associativity of these operators in classical, linear and computability logics, this difference, too, can be safely ignored. So, in this presentation we assume that the official language of CL2 does not include $``$, that $`\neg `$ is only allowed to be applied to atoms, and that $`,,,`$ are strictly binary.
What we have been referring to as “formulas” in the previous sections we will now call CL5-formulas, and formulas of the language of CL2 we will call CL2-formulas. CL5-formulas are thus nothing but CL2-formulas that do not contain $`,`$ and elementary atoms. A CL2-formula is said to be elementary iff it contains neither general atoms nor $``$,$``$. Ignoring any differences (in their status) between the general and elementary sorts of atoms, in this section $`,`$-free CL2-formulas — including elementary formulas — will be seen as formulas of classical propositional logic, with $``$ and $``$ having their standard meanings (“truth” and “falsity”, respectively). The terms “positive occurrence” and “negative occurrence” will also be used with their standard meanings as we did before. A surface occurrence of a subformula of a CL2-formula is an occurrence which is not in the scope of $`,`$. The elementarization of a CL2-formula $`F`$ is the result of replacing in $`F`$ every positive surface occurrence of each general atom by $``$, every negative surface occurrence of each general atom by $``$, every surface occurrence of each $``$-subformula by $``$, and every surface occurrence of each $``$-subformula by $``$. A CL2-formula is said to be stable iff its elementarization is a tautology of classical logic; otherwise the formula is said to be instable. In these terms, and with $`𝒫C`$ meaning “from premise(s) $`𝒫`$ conclude $`C`$”, system CL2 is given in by the following three rules of inference:
$`\stackrel{}{H}F`$, where $`F`$ is stable and $`\stackrel{}{H}`$ is the smallest set of formulas such that, whenever $`F`$ has a surface occurrence of a subformula $`G_1G_2`$, for both $`i\{1,2\}`$, $`\stackrel{}{H}`$ contains the result of replacing that occurrence in $`F`$ by $`G_i`$.
$`HF`$, where $`H`$ is the result of replacing in $`F`$ a surface occurrence of a subformula $`G_1G_2`$ by $`G_1`$ or $`G_2`$.
$`HF`$, where $`H`$ is the result of replacing in $`F`$ two — one positive and one negative — surface occurrences of some general atom by a nonlogical elementary atom that does not occur in $`F`$.
Axioms are not explicitly stated, but notice that the set of the premises of Rule (a) can be empty, in which case the conclusion acts as an axiom. Let us look at a couple of examples before moving any further. CL2 proves $`PPP`$, i.e. $`(\neg P\neg P)P`$. This formula follows from $`(\neg P\neg p)p`$ by Rule (c). In turn, the stable formula $`(\neg P\neg p)p`$ is derivable from the empty set of premises by Rule (a). On the other hand, CL2 does not prove $`PPP`$, i.e. $`\neg P(PP)`$. Indeed, this formula is instable and does not contain choice operators, so the only rule by which it could be derived is (c). The premise should be $`\neg p(pP)`$ or $`\neg p(Pp)`$ for some nonlogical elementary atom $`p`$. In either case we deal with an instable formula without choice operators and with only one occurrence of a general atom. Hence it cannot be the conclusion of any of the three rules of CL2.
###### Lemma 9.1
A CL5-formula is provable in CL2 iff it is an instance of a binary tautology.
Proof. In this proof, as agreed a while ago, we treat both general and elementary atoms as atoms of classical propositional logic. With this minor and irrelevant difference, the terms “instance”, “binary tautology” etc. have the same meanings as before.
($``$:) Consider an arbitrary CL5-formula $`F`$ provable in CL2. Fix a CL2-proof of $`F`$ in the form of a sequence (rather than tree) $`F_n,F_{n1},\mathrm{},F_1`$ of formulas, with $`F_1=F`$. We may assume that this sequence has no repetitions or other redundancies. We claim that, for each $`i`$ with $`1in`$, the following conditions are satisfied:
$`F`$ is an atomic-level instance of $`F_i`$, and (hence) $`F_i`$ does not contain $`,`$.
Whenever $`F_i`$ contains an elementary atom, that atom is non-logical, and has exactly two — one positive and one negative — occurrences in $`F_i`$.
If $`i<n`$, then $`F_i`$ is derived from $`F_{i+1}`$ by Rule (c).
$`F_n`$ is derived (from the empty set of premises) by Rule (a).
Condition 4 is obvious, because it is only Rule (a) that may take no premises. That Conditions 1-3 are also satisfied can be verified by induction on $`i`$. For the basis case of $`i=1`$, Conditions 1 and 2 are immediate because $`F`$ ($`=F_1`$) is its own atomic-level instance and, as a CL5-formula, it contains neither $`,`$ nor elementary atoms. $`F_1`$ cannot be derived by Rule (b) because, by Condition 1, $`F_1`$ does not contain any $``$. Nor can it be derived by Rule (a) unless $`n=1`$, for otherwise either $`F_1`$ would have to contain a $``$ (which is not the case according to Condition 1), or the proof of $`F`$ would have redundancies as $`F_1`$ would not really need any premises. Thus, if $`1<n`$, the only possibility for $`F_1`$ is to be derived from $`F_2`$ by Rule (c). For the induction step, assume $`i<n`$ and the above conditions are satisfied for $`F_i`$. According to Condition 3, $`F_i`$ is derived by Rule (c) from $`F_{i+1}`$. This obviously implies that $`F_{i+1}`$ inherits Conditions 1 and 2 from $`F_i`$. And that Condition 3 also holds for $`F_{i+1}`$ can be shown in the same way as we did for $`F_1`$.
Condition 1 implies that $`F`$ is an instance of $`F_n`$. Therefore, in order to complete our proof of the $`()`$ part of the lemma, it would now suffice to show that $`F_n`$ is an instance of a binary tautology. As the conclusion of Rule (a) (Condition 4), $`F_n`$ is stable. Let $`G`$ be the elementarization of $`F_n`$. The stability of $`F_n`$ means that $`G`$ is a tautology. Let $`H`$ be the result of replacing in $`G`$ every occurrence of $``$ and $``$ by an elementary nonlogical atom not occurring in $`G`$, in such a way that different occurrences of $`,`$ are replaced by different atoms. In view of Condition 2 (applied to $`F_n`$), it is obvious that $`H`$ is binary. By the same condition, $`F_n`$ did not contain $`,`$. This means that every occurrence of $``$ in $`G`$ comes from replacing a negative occurrence of a general atom in $`F_n`$, and every occurrence of $``$ in $`G`$ comes from replacing a positive occurrence of a general atom in $`F_n`$. It is not hard to see that, for this reason, $`F_n`$ is an instance of $`H`$. What remains to show is that $`H`$ is a tautology. But this is indeed so because $`H`$ results from the tautological $`G`$ by replacing positive occurrences of $``$ and negative occurrences of $``$. It is known from classical logic that such replacements do not destroy truth and hence tautologicity of formulas.
($``$:) Assume $`F`$ is a CL5-formula which is an instance of a binary tautology $`T`$. In view of Lemma 7.2, we may assume that $`T`$ is normal and $`F`$ is an atomic-level instance of it. For simplicity, we may also assume that all atoms of $`T`$ are elementary (elementary and non-logical, that is). Let us call the atoms that only have one occurrence in $`T`$ single, and the atoms that have two occurrences married. Let $`\sigma `$ be the substitution with $`\sigma (T)=F`$. Let $`G`$ be the formula resulting from $`T`$ by substituting each single atom $`q`$ by $`\sigma (q)`$. It is clear that then $`F`$ can be derived from $`G`$ by a series of applications of Rule (c), with each such application replacing two — a positive and a negative — occurrences of some married atom $`p`$ by $`\sigma (p)`$. So, in order to show that CL2 proves $`F`$, it would suffice to verify that $`G`$ is stable and hence it can be derived from the empty set of premises by Rule (a). But $`G`$ is indeed stable. To see this, consider the elementarization $`G^{}`$ of $`G`$. It results from $`T`$ by replacing its single atoms by $``$ (if negative) or $``$ (if positive). It is known from classical logic that replacing a single atom by whatever formula does not destroy tautologicity of formulas. Hence, as $`T`$ is a tautology, so is $`G^{}`$, meaning that $`G`$ is stable. $`\mathrm{}`$
###### Theorem 9.2
A formula is provable in CL5 iff it is valid in computability logic.
Proof. This theorem is an immediate corollary of Lemma 9.1, Theorem 7.5 and the known fact (proven in ) that CL2 is sound and complete with respect to the semantics of computability logic. $`\mathrm{}`$
Theorems 8.13 and 9.2 are similar in that both establish a soundness and completeness of CL5. A difference that may catch the eye is that Theorem 8.13 talks about all cirquents while the statement of Theorem 9.2 is restricted only to formulas. This is so because cirquents first emerged in the present paper, and hence the semantics of CL elaborated earlier did not extend to them. Here we briefly outline how to fill this gap. This explanation is very informal, and is only meant for readers already familiar with the semantics of CL. Every $`n`$-ary cirquentstructure (see Section 2) can be seen as an $`n`$-ary operation on games, and every cirquent $`C`$ with that structure then seen as the result of applying that operation to the games $`A_1,\mathrm{},A_n`$ represented by the oformulas of the pool of the cirquent. Specifically, $`C`$ is a game playing which means playing the $`n`$ games in parallel, i.e. it is a game on $`n`$ boards, where $`A_1`$ is played on board #1, $`A_2`$ is played on board #2, etc. The machine (player $``$) is considered the winner iff, for each group $`\mathrm{\Gamma }`$ of the cirquent, it wins at least one of the games (represented by the oformulas) of the group. This is the whole story. As we see, cirquents can be understood as parallel combinations of games in the same style as $`,`$-combinations. Not every combination represented by cirquents can be expressed using the ordinary (such as $`,`$) operators of CL. Among such combinations is the one expressed by the cirquent of Example 8.9. We claim without a proof that, with the just-outlined extended semantics of CL, Theorem 9.2 generalizes to all cirquents.
We further claim that Theorem 9.2 can be strengthened by replacing CL5 with the more expressive cirquent calculus system CL6. What makes the language of CL6 stronger than that of CL5 is that, along with general atoms, CL6-formulas may also contain elementary atoms including the logical atoms $``$ and $``$, with $`\neg `$ now always written as $``$ and $`\neg `$ as $``$. In accordance with the earlier-established meaning of the term “elementary”, we call a CL6-formula elementary iff it does not contain general atoms. With “formulas” now meaning any CL6-formulas, the set of the rules of CL6 is obtained from that of CL5 by adding to it $``$ (understood as a singleton cirquent) as an additional axiom, plus the rule of contraction limited only to elementary formulas (that is, the contracted formula $`F`$ is required to be elementary). One could show that Theorem 9.2 remains true with CL6 instead of CL5. So does it with “formula” further replaced by “cirquent”, where computability-logic validity of cirquents is understood in the sense of the previous paragraph.
## 10 What is next?
The results of the previous section can be seen as first steps within the program of developing cirquent-calculus deductive systems for ever more expressive fragments of computability logic. Probably the same applies to our present version of abstract resource semantics. Linear-logic literature abounds with examples illustrating the intended resource intuitions for additive and exponential operators. Yet, just as in the case of multiplicatives, such intuitions have never found a good formal explication, the reason for which, again, being the inherent incompleteness of linear logic. If and when abstract resource semantics is successfully and naturally extended to additive- and exponential-style operators, one should probably expect it to validate the same class of formulas as computability logic does, thus yielding a logic properly stronger than linear (affine) logic in its full language.
It would also be interesting to see if abstract resource semantics can be modified so as to make a meaningful semantics for relevance logic, specifically, for the obviously relevance-logic-style system which results from CCC by deleting weakening (but not contraction) and perhaps duplication as well.
Back to proof theory, besides adopting more expressive underlying languages for formulas, a promising direction for cirquent calculus to grow might be relaxing its concept of a cirquentstructure. As illustrated in Section 2, cirquentstructures in the sense of the present paper are special sorts of circuits — ones of depth 2, with all level-1 gates being $``$-gates and the single level-2 gate being an $``$-gate. Allowing circuits of arbitrary depths and forms (including circuits with non-traditional types of gates) as underlying cirquentstructures would yield more general concepts of cirquents. It is appropriate to still use the same term “cirquent” in these cases. After all, the gist of the cirquent calculus approach is its ability to syntactically capture resource-sharing, and resource-sharing is exactly what circuits are all about.
As long as obtaining soundness and completeness is the only goal, the semantics of computability logic — of its $`(\neg ,,)`$-fragment for sure — does not call for generalizations in the above style, just as classical logic needs nothing beyond ordinary sequent calculus. But even if, in some good sense, the present form of cirquent calculus is natural and sufficiently powerful as a “universal syntax” (not that the author really wants to make such a claim), studying generalizations of it could be still worthwhile. Switching to the new syntactic vision known as the calculus of structures () introduced by Alessio Guglielmi has proved very beneficial, among the benefits being exponential speedups of proofs, and the possibility of richer combinatorial analysis of proofs for diverse logical systems, including the old and well-axiomatized classical, linear or modal logics. Achieving or enhancing similar effects could be some of the possible motivations for tackling general forms of cirquent calculus. A central idea in the calculus of structures is what is called deep inference. It means deducing inside a formula at any depth, as opposed to the shallow inference of sequent calculus where rules only see the roots of formula trees. Notice that all three rules of CL2 modify subformulas at any ($`,`$)-depths, thus essentially being deep inference rules. Some time ago this observation gave rise to the hope that the calculus of structures could be a well-suited deductive framework for CL.<sup>6</sup><sup>6</sup>6See http://alessio.guglielmi.name/res/cos/crt.html#ComLog However, efforts to axiomatize even the simplest fragments of CL in it have not been successful so far. While CL obviously does need deep inference, what is making the calculus of structures apparently insufficient is the absence of resource-sharing capabilities. Just like the rules of CL2, our cirquent calculus rules from Section 3 can be seen as limited sorts of deep inference, specifically, inferences that always take place at depth 2 (when they modify oformulas) or 1 (when they modify ogroups). With arbitrary-depth underlying cirquentstructures, generalized cirquent calculus would naturally invite inferences at any structural levels, thus enjoying the full generality, flexibility and power offered by both deep inference and resource-sharing.
Last but maybe not least, it would be interesting to understand the impact of sharing on cut elimination, the idol worshiped by all proof-theoreticians. Of course, before attempting to even ask this question, one needs to figure out what a cut rule should exactly mean in the context of cirquent calculus.
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# Pólya Theory for Orbiquotient Sets
## 1 Introduction
Assume that a finite group $`G`$ acts on a finite set $`X.`$ The quotient $`X/G`$ of the action of $`G`$ on $`X`$ is a rich and subtle concept, traditionally $`X/G=\{\overline{x}xX\}`$ where $`\overline{x}=\{gxxX\}.`$ In recent years it has proven convenient to modify this notion in various contexts. For example one may think of $`X/G`$ as the groupoid whose set of objects is $`X`$ and with morphisms given by $`X/G(a,b)=\{gGga=b\}`$ for $`a,bG.`$ Following Connes the groupoid $`X/G`$ is studied with the methods of non-commutative geometry, i.e. looking at the convolution (incidence) algebra of $`X/G.`$ Another approach to quotient sets became rather popular after Vafa and Witten introduced in the so called stringy Euler numbers. In a nutshell they considered the Euler numbers of orbiquotient sets
$$X/^{orb}G=\underset{\overline{g}C(G)}{}X^g/Z(g),$$
where $`C(G)`$ is the set of conjugacy classes of $`G`$, $`X^gX`$ is the set of points fixed by $`gG`$, and $`Z(g)`$ is the centralizer of $`g`$ in $`G.`$ We will assume that a representative $`gG`$ has been chosen for each conjugacy class $`\overline{g}`$ of $`G`$. Orbiquotient sets first appeared, see , in the context of equivariant K-theory in the works of Atiyah and Segal. The goal of this paper is to bring the notion of orbiquotient sets into combinatorial waters. Let us provide a combinatorial motivation for the study of orbiquotient sets inspired by an analogue topological construction given by Hirzebruch and H$`\ddot{\text{o}}`$fer in . Consider the set of $`n`$-cycles in $`X/G,`$ i.e., the set of maps $`f:X/G`$ such that $`f(k)=f(k+n)`$ for $`k.`$ Suppose we want to lift $`f`$ to a map $`l:X`$ such that $`\pi l=f`$, where $`\pi :XX/G`$ is the canonical projection. The lift $`l`$ will not be unique, indeed if $`l`$ is a lift then $`gl`$ is another a lift; also we have that $`l(k)=gl(k+n)`$ for all $`k`$ and some $`gG`$. Thus the set of lifts of $`n`$-periodic maps $`X/G`$ may be identified with
$$\{l:Xl(k)=gl(k+n)\text{ for }k\text{ and some }gG\}/G.$$
Inside the later set sits $`I(G,X)/G`$ the set of constant maps, where
$$I(G,X)=\{(g,x)G\times X|gx=x\}$$
is the so called inertial set . The group $`G`$ acts on $`I(G,X)`$ as $`k(g,x)=(kgk^1,kx),`$ and it is not hard to see that
$$I(G,X)/G=\underset{\overline{g}C(G)}{}X^g/Z(g)=X/^{orb}G.$$
In this paper we develop orbianalogues for two main results in elementary combinatorics, the orbit counting lemma and the Polya-Redfield theorem, see . We fix a commutative ring $`𝔸`$ and consider the category of $`𝔸`$-weighted sets whose objects are pairs $`(X,f)`$ where $`X`$ is a finite set and $`f:X𝔸`$ is an arbitrary map called the weight of $`X`$. Morphisms between $`𝔸`$-weighted sets are weight preserving bijections. The cardinality $`|X|_f`$ of a weighted set $`(X,f)`$ is given by
$$|X|_f=\underset{xX}{}f(x).$$
A finite group $`G`$ acts on $`(X,f)`$ if $`G`$ acts on $`X`$ and $`f(gx)=f(x)`$ for all $`gG,`$ $`xX.`$ The Cauchy-Frobenius-Burnside orbit counting lemma gives us a way to compute $`|X/G|_f`$ as follows:
$$|X/G|_f=\frac{1}{|G|}\underset{gG}{}|X^g|_f.$$
Suppose now that $`G`$ is a group of permutations $`GS_m`$. The cardinality of $`X^m/G`$ is determined by the Pólya-Redfield theorem:
$$|X^m/G|_f=P_G(|X|_f,|X|_{f^2},\mathrm{},|X|_{f^m}),$$
where $`P_G`$ is the cycle index polynomial of $`G`$ given by
$$P_G(x_1,x_2,\mathrm{},x_m)=\frac{1}{|G|}\underset{gG}{}x_1^{c_1(g)}\mathrm{}x_m^{c_m(g)},$$
and $`c_i(g)`$ is the number of $`g`$-cycles of length $`i.`$ If $`X=[n]`$ and $`f(i)=x_i`$ for $`iX`$, then directly from the definition of quotient sets we get that
$$|[n]^m/G|_f=\underset{(i_1,\mathrm{},i_n)^n}{}c_G(i_1,\mathrm{},i_n)x_1^{i_1}\mathrm{}x_n^{i_n},$$
where $`c_G(i_1,\mathrm{},i_n)`$ counts the colorations of $`[m]`$ with $`i_k`$ elements of color $`k[n],`$ and two colorations are identified if they are linked by the action of $`G`$. The Pólya-Redfield theorem allows us to compute the coefficients $`c_G(i_1,\mathrm{},i_n)`$ in a different way, namely we have that
$$|[n]^m/G|_f=P_G(\underset{j=1}{\overset{n}{}}x_j,\underset{j=1}{\overset{n}{}}x_j^2,\mathrm{},\underset{j=1}{\overset{n}{}}x_j^m).$$
The rest of this work is organized as follows. In Section 2 we provide an orbi-analogue of the orbit counting lemma. In Section 3 we provide an orbi-analogue of the Pólya-Redfield theorem in full generality, we shall see that lattice of partitions plays a fundamental role in our presentation. In the remaining sections we explicitly compute the orbicycle index polynomial for various groups in increasing order of difficulty. In Section 4 we consider the case of cyclic groups, and apply it to the study of orbicycles in orbiquotient sets. In Section 5 we consider the full symmetric group. In a rather dull fashion we may regard combinatorics as geometry in dimension zero. It is thus rather interesting when one can show that the zero dimensional combinatorial case determines the higher dimensional situation. A theorem of this sort is proved at the end of Section 5 which provides a strong motivation for the study of orbiquotient sets. In Section 6 we compute the orbicycle index polynomial for the dihedral groups.
## 2 Orbi-analogue of the orbit counting lemma
If $`SG`$ and $`G`$ acts on $`X`$, then we set $`X^S=\{xXgx=x\text{for }gS\}.`$ Also we let $`g_1,\mathrm{},g_n`$ be the subgroup of $`G`$ generated by $`\{g_1,\mathrm{},g_n\}G.`$
###### Definition 1.
The orbiquotient of $`X`$ by the action of $`G`$ is the set given by
$$X/^{orb}G=\underset{\overline{g}C(G)}{}X^g/Z(g).$$
The orbiquotient $`X/^{orb}G`$ is well defined up to canonical bijections. Indeed if $`h=kgk^1`$ then the map $`\psi :X^gX^h`$ given by $`\psi (x)=kx`$ induces a bijection
$$\psi :X^g/Z(g)X^h/Z(h).$$
If $`G`$ acts on a weighted set $`(X,f)`$ then $`X/^{orb}G`$ is also weighted: $`X^g`$ is weighted by $`f_{X^g}`$ and $`X^g/Z(g)`$ is weighted by $`f(\overline{x})=f(x)`$ for $`\overline{x}X^g/Z(g).`$ Our next result is the orbi-analogue of the orbit counting lemma, let us first introduce a notation that will be used repeatedly
$$P(G)=\{(\overline{g},h)|\overline{g}C(G)\text{ and }hZ(g\}.$$
###### Theorem 2.
If $`G`$ acts on $`(X,f),`$ then the cardinality of $`X/^{orb}G`$ is given by
$$\left|X/^{orb}G\right|_f=\frac{1}{|G|}\underset{(\overline{g},h)P(G)}{}|\overline{g}|\left|X^{g,h}\right|_f.$$
The proof of this result is quite simple:
$`\left|X/^{orb}G\right|_f`$ $`=`$ $`{\displaystyle \underset{\overline{g}C(G)}{}}\left|X^g/Z(g)\right|_f`$
$`=`$ $`{\displaystyle \underset{\overline{g}C(G)}{}}{\displaystyle \frac{|\overline{g}|}{|G|}}{\displaystyle \underset{hZ(g)}{}}|X^gX^h|_f`$
$`=`$ $`{\displaystyle \frac{1}{|G|}}{\displaystyle \underset{(\overline{g},h)P(G)}{}}|\overline{g}|\left|X^{g,h}\right|_f.`$
## 3 Orbi-analogue of Pólya-Redfield theorem
Let $`\mathrm{Par}(X)`$ be lattice of partitions of $`X.`$ The minimal and maximal elements of $`\mathrm{Par}(X)`$ are $`\{\{x\}xX\}`$ and $`\{X\},`$ respectively. The $`\mathrm{joint}`$ $`\pi \rho `$ of partitions $`\pi `$ and $`\rho `$ is defined by demanding that $`i,jX`$ belong to a block of $`\pi \rho `$ if there exists a sequence $`i=a_0,a_1,\mathrm{},a_n=j,`$ such that for $`0in1`$ either $`a_i`$ and $`a_{i+1}`$ belong to a block in $`\pi ,`$ or $`a_i`$ and $`a_{i+1}`$ belong a block in $`\rho .`$ The $`\mathrm{meet}`$ of partitions $`\pi `$ and $`\rho `$ is $`\pi \rho =\{BCB\pi ,C\rho ,BC0\}.`$ Let the group $`G`$ act on a set $`X`$ with $`n`$-elements. Each $`gG`$ induces a partition $`C(g)`$ on $`X`$ such that $`C(g)=_{i=1}^nC_i(g),`$ where $`C_i(g)=\{g\text{-}\mathrm{cycles}\mathrm{on}\mathrm{X}\mathrm{of}\mathrm{length}i\}`$ for $`1in.`$ We use the notation $`c(g)=|C(g)|`$ and $`c_i(g)=|C_i(g)|`$ for $`1in.`$ If $`\pi `$ is a partition of $`X`$ we let $`b_k(\pi )`$ be the number of blocks of $`\pi `$ of cardinality $`k.`$
###### Definition 3.
The orbicycle index polynomial $`P_G^{orb}(x_1,x_2,\mathrm{})[x_1,x_2,\mathrm{}]`$ is given by
$$P_G^{orb}(x_1,x_2,\mathrm{})=\frac{1}{|G|}\underset{(\overline{g},h)P(G)}{}|\overline{g}|x^{C(g)C(h)},$$
where
$$x^{C(g)C(h)}=\underset{k1}{}x_k^{b_k(C(g)C(h))}.$$
If $`GS_m`$ then $`G`$ acts on $`X^m`$. Suppose that $`g,hG`$ commute, then $`i,j[m]`$ belong to the same block of $`C(g)C(h)`$ if and only if there exist $`a,b`$ such that $`j=(g^ah^b)(i).`$ It is easy to check that $`fX^m`$ is fixed by $`g`$ and $`h`$ if and only if $`f`$ is constant on each block of $`C(g)C(h).`$
###### Theorem 4.
Let $`(X,f)`$ be an $`𝔸`$-weighted set and $`GS_m.`$ The cardinality of $`X^m/^{orb}G`$ is given by
$$|X^m/^{orb}G|_f=P_G^{orb}(|X|_f,|X|_{f^2},\mathrm{}).$$
###### Proof.
$`|X^m/^{orb}G|_f`$ $`=`$ $`{\displaystyle \frac{1}{|G|}}{\displaystyle \underset{(\overline{g},h)P(G)}{}}|\overline{g}|\left|\left(X^{[m]}\right)^g\left(X^{[m]}\right)^h\right|_{\stackrel{~}{f}}`$
$`=`$ $`{\displaystyle \frac{1}{|G|}}{\displaystyle \underset{(\overline{g},h)P(G)}{}}|\overline{g}|{\displaystyle \underset{\alpha :[m]X,\alpha g=\alpha h=\alpha }{}}{\displaystyle \underset{x[m]}{}}f(\alpha (x))`$
$`=`$ $`{\displaystyle \frac{1}{|G|}}{\displaystyle \underset{(\overline{g},h)P(G)}{}}|\overline{g}|{\displaystyle \underset{\alpha :C(g)C(h)X}{}}{\displaystyle \underset{BC(g)C(h)}{}}{\displaystyle \underset{xB}{}}f(\alpha (x))`$
$`=`$ $`{\displaystyle \frac{1}{|G|}}{\displaystyle \underset{(\overline{g},h)P(G)}{}}|\overline{g}|{\displaystyle \underset{BC(g)C(h)}{}}{\displaystyle \underset{yX}{}}f(y)^{|B|}`$
$`=`$ $`{\displaystyle \frac{1}{|G|}}{\displaystyle \underset{(\overline{g},h)P(G)}{}}|\overline{g}|{\displaystyle \underset{k1}{}}|X|_{f^k}^{b_k(C(g)C(h))}`$
$`=`$ $`P_G^{orb}(|X|_f,|X|_{f^2},\mathrm{}).`$
Let $`X=[n]`$ and $`f(i)=x_i`$, then one can check directly from the definition that
$$|[n]^m/^{orb}G|_f=\underset{(i_1,\mathrm{},i_n)^n}{}c_G^{orb}(i_1,\mathrm{},i_n)x_1^{i_1}\mathrm{}x_n^{i_n},$$
where $`c_G^{orb}(i_1,\mathrm{},i_n)`$ counts colorations $`c`$ of $`[m]`$ with colors in $`[n]`$ such that:
* There are $`i_k`$ elements in $`[m]`$ of color $`k[n].`$
* $`c`$ is $`g`$-invariant for some $`\overline{g}C(G).`$
* Two $`g`$-invariant colorations are identified if they can be linked by the action of $`Z(g).`$
The orbi-analgogue of the Pólya-Redfield gives us another way to compute the coefficients $`c_G^{orb}(i_1,\mathrm{},i_n)`$, namely we have that
$$|[n]^m/^{orb}G|_f=P_G^{orb}(\underset{j=1}{\overset{n}{}}x_j,\underset{j=1}{\overset{n}{}}x_j^2,\mathrm{},\underset{j=1}{\overset{n}{}}x_j^n).$$
## 4 Orbicycle index polynomial of $`_n`$
Let $`_+`$ be the set of positive integers and let $`(x_1,x_2,\mathrm{},x_k)`$ be the greatest common divisor of $`x_1,x_2,\mathrm{},x_k^+`$. The cyclic group with $`n`$-elements is denoted by $`_n=\{1,2,\mathrm{},n\}`$. For $`n,k_+`$ we define an equivalence relation on $`_n^k`$ as follows: $`x`$ and $`y`$ are equivalent if and only if $`(x,n)=(y,n).`$ It is easy to verify that $`_n^k=_{d|n}\{x_n^k(x,n)=d\},`$ and thus we have
$$n^k=\underset{dn}{}|\{x_n^k(x,n)=d\}|.$$
For $`n_+`$ the Jordan totient function $`J_k`$, see , is given by $`J_k(n)=\left|\{x_n^k(x,n)=1\}\right|.`$ For each $`d|n`$ we have $`J_k\left(\frac{n}{d}\right)=\left|\{x_n^k(x,n)=d\}\right|,`$ therefore we get that $`n^k=_{dn}J_k(d).`$ By the M$`\ddot{\text{o}}\mathrm{bius}`$ inversion formula $`J_k(n)=_{dn}\mu \left(\frac{n}{d}\right)d^k,`$ thus $`J_k(p^r)=p^{kr}p^{k(r1)},`$ for $`p`$ prime, and for arbitrary integer $`n=p_1^{\alpha _1}\mathrm{}p_r^{\alpha _r}`$ where $`p_1,\mathrm{},p_r`$ are distinct prime numbers, we get
$$J_k(n)=n^k\left(1\frac{1}{p_1^k}\right)\mathrm{}\left(1\frac{1}{p_r^k}\right).$$
We shall need the following property, an easy consequence of the previous considerations, of the Jordan totient function:
$$\underset{x_n^k}{}f((x,n))=\underset{dn}{}J_k\left(\frac{n}{d}\right)f(d),$$
for any $`f:\{d:dn\}𝔸.`$ Recall that if $`x,y_nS_n,`$ then $`|C(x)C(y)|=(x,y,n),`$ and all blocks in $`C(x)C(y)`$ are of cardinality $`\frac{n}{(n,x,y)}.`$ Indeed if $`x_n,`$ then $`\frac{_n}{(x)}_{(n,x)},`$ since for $`a,b_n`$ we have that $`a=b`$ in $`_n/(x)`$ if and only if $`a=b`$ mod $`(n,x).`$ Thus there are $`(n,x)`$ blocks in $`_n/(x)`$ all of them of cardinality $`\frac{n}{(n,x)}.`$ Similarly if $`x,y_n,`$ then $`a,b_n`$ are in the same block of $`C(x)C(y)`$ if and only if there exist $`r,s`$ such that $`b=a+rx+sy`$ mod $`n`$, or equivalently $`a=b`$ mod $`(n,x,y).`$
###### Theorem 5.
$$P__n^{orb}(y_1,y_2,\mathrm{},y_n)=\frac{1}{n}\underset{dn}{}J_2(d)y_d^{\frac{n}{d}}.$$
###### Proof.
$`P__n^{orb}(y_1,y_2,\mathrm{},y_n)`$ $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{(\overline{g},h)P(_n)}{}}|\overline{g}|{\displaystyle \underset{k1}{}}y_k^{b_k(C(g)C(h))}`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{(x,y)_n\times _n}{}}{\displaystyle \underset{k1}{}}y_k^{b_k(C(x)C(y))}`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{(x,y)_n\times _n}{}}y_{\frac{n}{(x,y,n)}}^{(x,y,n)}`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{dn}{}}J_2\left(d\right)y_d^{\frac{n}{d}}.`$
The coefficients $`c__n(i_1,\mathrm{},i_m)`$ are computed in . As a corollary of the previous result we obtain that
$$c__n^{orb}(i_1,\mathrm{},i_m)=\frac{1}{n}\underset{d(i_1,\mathrm{},i_m)}{}J_2(d)\frac{\left(\frac{i_1}{d}+\mathrm{}+\frac{i_m}{d}\right)!}{\frac{i_1}{d}!\mathrm{}\frac{i_m}{d}!}.$$
Let $`p`$ be a prime number and $`r^+,`$ necklaces without a clasp with $`p`$ beads and $`r`$ colors may be identified with the set $`C_p([r])=[r]^p/_p.`$ As explained in we have that
$$|C_p([r])|=P__p(r,r,\mathrm{})=\frac{1}{p}\underset{dp}{}\phi (d)r^{\frac{p}{d}}=r+\frac{r^pr}{p}.$$
###### Definition 6.
The set $`C_n^{orb}(X)`$ of orbi $`n`$-cycles in $`X`$ is given by $`C_n^{orb}(X)=X^n/^{orb}_n.`$
In analogy with the example above we define the set of orbi-necklaces without a clasp with $`p`$ beads and $`r`$ colors to be $`C_p^{orb}([r])=[r]^p/^{orb}_p.`$ Its cardinality is given by
$$|C_p^{orb}([r])|=P__p^{orb}(r,r,\mathrm{})=\frac{1}{p}\underset{dp}{}J_2(d)r^{\frac{p}{d}}=rp+\frac{r^pr}{p}.$$
Next couple of results count explicitly the number of orbicycles in orbiquotient sets.
###### Theorem 7.
If $`G`$ acts on $`(X,f)`$ then
$$|C_n^{orb}(X/^{orb}G)|_f=\frac{1}{n}\underset{\alpha :[\frac{n}{d}]P(G)}{}\frac{J_2(d)}{|G|^{\frac{n}{d}}}\underset{i=1}{\overset{\frac{n}{d}}{}}|\pi _{C(G)}(\alpha (i))||X^{\alpha (i)}|_{f^d}.$$
###### Proof.
$`|C_n^{orb}(X/^{orb}G)|_f`$ $`=`$ $`P__n^{orb}(|X/^{orb}G|_f,|X/^{orb}G|_{f^2},\mathrm{})`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{dn}{}}J_2(d)\left(|X/^{orb}G|_{f^d}\right)^{\frac{n}{d}}`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{dn}{}}J_2(d)\left({\displaystyle \frac{1}{|G|}}{\displaystyle \underset{(\overline{g},h)P(G)}{}}|\overline{g}||X^{g,h}|_f\right)^{\frac{n}{d}}`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{dn}{}}{\displaystyle \frac{J_2(d)}{|G|^{\frac{n}{d}}}}{\displaystyle \underset{\alpha :[\frac{n}{d}]P(G)}{}}{\displaystyle \underset{i=1}{\overset{\frac{n}{d}}{}}}|\pi _{C(G)}(\alpha (i))||X^{\alpha (i)}|_{f^d}`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{\alpha :[\frac{n}{d}]P(G)}{}}{\displaystyle \frac{J_2(d)}{|G|^{\frac{n}{d}}}}{\displaystyle \underset{i=1}{\overset{\frac{n}{d}}{}}}|\pi _{C(G)}(\alpha (i))||X^{\alpha (i)}|_{f^d}.`$
###### Theorem 8.
Let $`(X,f)`$ be be an $`𝔸`$-weithed set and $`GS_m`$, then we have
$$|C_n^{orb}(X^{[m]}/^{orb}G)|_f=\frac{1}{n}\underset{\alpha :[\frac{n}{d}]P(G)}{}\frac{J_2(d)}{|G|^{\frac{n}{d}}}\underset{i=1}{\overset{\frac{n}{d}}{}}|X|_{f^d}^{b(\alpha (i))}$$
where $`|X|_{f^d}^{b(\alpha (i))}=_{k1}|X|_{f^{dk}}^{b_k(C(\pi _1(\alpha (i)))C(\pi _2(\alpha (i))))}`$
###### Proof.
$`|C_n^{orb}(X^{[m]}/^{orb}G)|_f`$ $`=`$ $`P__n^{orb}(|X^{[m]}/^{orb}G|_f,|X^{[m]}/^{orb}G|_{f^2},\mathrm{})`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{dn}{}}J_2(d)P_G^{orb}(|X|_{f^d},|X|_{f^{2d}},\mathrm{})^{\frac{n}{d}}`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{dn}{}}J_2(d)\left({\displaystyle \frac{1}{|G|}}{\displaystyle \underset{(\overline{g},h)P(G)}{}}{\displaystyle \underset{k1}{}}|X|_{f^{dk}}^{b_k(C(g)C(h))}\right)^{\frac{n}{d}}`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{dn}{}}{\displaystyle \frac{J_2(d)}{|G|^{\frac{n}{d}}}}{\displaystyle \underset{\alpha :[\frac{n}{d}]P(G)}{}}{\displaystyle \underset{i=1,k1}{}}|X|_{f^{dk}}^{b_k(C(\pi _1(\alpha (i)))C(\pi _2(\alpha (i))))}`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{\alpha :[\frac{n}{d}]P(G)}{}}{\displaystyle \frac{J_2(d)}{|G|^{\frac{n}{d}}}}{\displaystyle \underset{i=1}{}}|X|_{f^d}^{b(\alpha (i))}.`$
Using similar methods one can count cycles on orbiquotient sets:
$$|C_n(X/G)|_f=\frac{1}{n}\underset{\alpha :[\frac{n}{d}]G}{}\frac{J_2(d)}{|G|^{\frac{n}{d}}}\underset{i=1}{\overset{\frac{n}{d}}{}}|X^{\alpha (i)}|_{f^d},$$
and
$$|C_n(X^m/G)|_f=\frac{1}{n}\underset{\alpha :[\frac{n}{d}]G}{}\frac{J_2(d)}{|G|^{\frac{n}{d}}}\underset{i=1}{\overset{\frac{n}{d}}{}}|X|_{f^d}^{c(\alpha (i))},$$
where $`|X|_{f^d}^{c(g)}={\displaystyle \underset{k1}{}}|X|_{f^{dk}}^{c_k(g)}`$ for $`gG.`$
## 5 Orbicyle index polynomial of $`S_n`$
A partition of depth $`k`$, denoted by $`\alpha _kn`$, of $`n_+`$ is a map $`\alpha :\left(_+\right)^k`$ such that
$$\underset{(i_1,\mathrm{},i_k)_+^k}{}i_1\mathrm{}i_k\alpha (i_1,\mathrm{},i_k)=n.$$
A partition of depth $`1`$ is a partition in the usual sense. To each partition $`\alpha `$ we associate a canonical permutation of $`[n]`$ whose cycle structure is determined by $`\alpha .`$ Keeping this correspondence in mind one can check that if $`\alpha n`$ then
1. $`Z(\alpha )`$ is isomorphic to $`_{i=1}^n_i^{\alpha _i}S_{\alpha _i}.`$
2. If $`hZ(\alpha ),`$ then $`b_k(C(\alpha )C(h))=_{dk}c_{\frac{k}{d}}(\pi _d(h))`$ where $`\pi _d(h)`$ is the projection of $`h`$ into $`_{i=1}^nS_{\alpha _i}.`$
###### Theorem 9.
$$P_{S_n}^{orb}(y_1,y_2,\mathrm{},y_n)=\underset{\beta _2n}{}\underset{(i,j,k)[n]\times [\alpha _i]\times [n]}{}\frac{y_k^{_{dk}\beta (d,\frac{k}{d})}}{j^{\beta (i,j)}\beta (i,j)!}.$$
###### Proof.
By the previous remarks we have
$`P_{S_n}^{orb}(y_1,y_2,\mathrm{},y_n)`$ $`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle \underset{(\overline{g},h)P(S_n)}{}}|\overline{g}|{\displaystyle \underset{k=1}{\overset{n}{}}}y_k^{b_k(C(g)C(h))}`$
$`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle \underset{\alpha n}{}}{\displaystyle \frac{n!}{_{i=1}^ni^{\alpha _i}\alpha _i!}}{\displaystyle \underset{hZ(\alpha )}{}}{\displaystyle \underset{k=1}{\overset{n}{}}}y_k^{b_k(C(\alpha )C(h))}`$
$`=`$ $`{\displaystyle \underset{\alpha n}{}}{\displaystyle \frac{1}{_{i=1}^n\alpha _i!}}{\displaystyle \underset{h_{i=1}^nS_{\alpha _i}}{}}{\displaystyle \underset{k=1}{\overset{n}{}}}y_k^{_{dk}c_{\frac{k}{d}}(\pi _d(h))}`$
$`=`$ $`{\displaystyle \underset{\alpha n}{}}{\displaystyle \frac{1}{_{i=1}^n\alpha _i!}}{\displaystyle \underset{\beta _i\alpha _i}{}}{\displaystyle \frac{_{i=1}^n\alpha _i!}{_{(i,j)[n]\times [\alpha _i]}j^{\beta (i,j)}\beta (i,j)!}}{\displaystyle \underset{k=1}{\overset{n}{}}}y_k^{_{dk}\beta _{\frac{k}{d}}(d)}`$
$`=`$ $`{\displaystyle \underset{\beta _2n}{}}{\displaystyle \underset{(i,j,k)[n]\times [\alpha _i]\times [n]}{}}{\displaystyle \frac{y_k^{_{dk}\beta (d,\frac{k}{d})}}{j^{\beta (i,j)}\beta (i,j)!}}.`$
Above we used the fact that $`_{dk}c_{\frac{k}{d}}(\pi _d(h))`$ depends only on the cycle structure of $`\pi _d(h).`$
The orbicycle index polynomial can be use to compute the even dimensions of the orbifold cohomology groups for global orbifolds of the form $`M^n/^{\text{orb}}G,`$ where $`M`$ is a compact smooth manifold, and $`GS_n.`$ For simplicity we only consider cohomology in even dimensions. The orbifold cohomology is defined as follows:
$$\mathrm{H}^{orb}\left(M^n/G\right)=\underset{\overline{g}C(G)}{}\mathrm{H}\left(\left(M^n\right)^g\right)^{Z(g)}.$$
The following result is a direct consequence of the characterization of the centralizer of permutations previously discussed.
###### Lemma 10.
$$\mathrm{H}^{\text{orb}}\left(M^n/G\right)=\underset{\overline{g}C(G)}{}\underset{i}{}\left(\mathrm{H}\left(M\right)^{c_i(g)}\right)^{S_{c_i(g)}}.$$
###### Proof.
$`\mathrm{H}^{\mathrm{orb}}\left(M^n/G\right)=`$ $`{\displaystyle \underset{\overline{g}C(G)}{}}\mathrm{H}\left(\left(M^n\right)^g\right)^{Z(g)}`$
$`=`$ $`{\displaystyle \underset{\overline{g}C(G)}{}}\mathrm{H}({\displaystyle \underset{i}{}}M^{c_i(g)})^{Z(g)}`$
$`=`$ $`{\displaystyle \underset{\overline{g}C(G)}{}}{\displaystyle \underset{i}{}}\left(\mathrm{H}\left(M\right)^{c_i(g)}\right)^{S_{c_i(g)}}.`$
Assume that we are given a finite basis $`X`$ for $`\mathrm{H}(M)`$, then we have the following result.
###### Theorem 11.
$$dim\left(\mathrm{H}^{orb}\left(M^n/S_n\right)\right)=P_G^{orb}(|X|,\mathrm{},|X|).$$
###### Proof.
$`dim\left(\mathrm{H}^{orb}\left(M^n/S_n\right)\right)`$ $`={\displaystyle \underset{\overline{g}C(G)}{}}{\displaystyle \underset{i}{}}\left|X^{c_i(g)}/S_{c_i(g)}\right|`$
$`=\left|X^n/^{orb}S_n\right|`$
$`=P_G^{orb}(|X|,\mathrm{},|X|).`$
Notice that above we use the trivial weight on $`X`$; using a generic weight we obtain further information on the orbifold cohomology groups. Theorem $`12`$ gives a combinatorial interpretation for the orbifold cohomology groups, however we do not have a combinatorial interpretation for the orbifold product introduced by Chen and Ruan in . Until recently this problem seemed hopeless, however the alternative description of the Chen-Ruan product introduced by Jarvis, Kauffman and Kimura in could pave the way for such a combinatorial understanding. Theorem $`11`$ suggests the possibility of constructing, along the lines of , an orbi-analogue for the symmetric functions. This issue deserves further research.
## 6 Orbicycle index polynomial of $`D_n`$
The generators $`\rho `$ and $`\tau `$ of the dihedral group $`D_n=\{e,\rho ,\mathrm{},\rho ^{n1},\tau ,\mathrm{},\tau \rho ^{n1}\}`$ are such that $`\rho ^n=e`$, $`\tau ^2=e`$ and $`\tau \rho =\rho ^{n1}\tau .`$ The conjugacy classes of the dihedral groups are described in the following tables, see . For $`n`$ odd there are $`\frac{n+3}{2}`$ conjugacy classes organized in three families
| Conjugacy class | Representative | Centralizer subgroup |
| --- | --- | --- |
| $`\{e\}`$ | $`e`$ | $`D_n`$ |
| $`\{\rho ^i,\rho ^i\}`$ for $`1i\frac{n1}{2}`$ | $`\rho ^i`$ | $`\{\rho ^i0i<n\}`$ |
| $`\{\rho ^i\tau 0i<n\}`$ | $`\tau `$ | $`\{e,\tau \}`$ |
For $`n`$ even there are $`\frac{n+6}{2}`$ conjugacy classes organized in five families
| Conjugacy class | Representative | Centralizer subgroup |
| --- | --- | --- |
| $`\{e\}`$ | $`e`$ | $`D_n`$ |
| $`\{\rho ^{\frac{n}{2}}\}`$ | $`\rho ^{\frac{n}{2}}`$ | $`D_n`$ |
| $`\{\rho ^i,\rho ^i\},`$ $`1i<\frac{n}{2}`$ | $`\rho ^i`$ | $`\{\rho ^i0i<n\}`$ |
| $`\{\rho ^{2i}\tau 0i<\frac{n}{2}\}`$ | $`\tau `$ | $`\{e,\tau ,\rho ^{\frac{n}{2}},\rho ^{\frac{n}{2}}\tau \}`$ |
| $`\{\rho ^{2i+1}\tau 0i<\frac{n}{2}\}`$ | $`\rho \tau `$ | $`\{e,\rho \tau ,\rho ^{\frac{n}{2}},\rho ^{\frac{n}{2}+1}\tau \}`$ |
For real numbers $`a_11,\mathrm{},a_k1`$, and $`n^+`$, see , we set
$$\phi (a_1,\mathrm{},a_k,n)=|\{x_n^k1x_ia_i,i[k],(x,n)=1\}|.$$
For $`x`$ we denote by $`x`$ the floor function of $`x`$, that is the largest integer no greater than $`x`$. We use the fact that for a real number $`x1`$ and $`n^+`$ we have
$$\frac{x}{n}=\{a:1ax,na\}$$
Proceeding as in the previous sections one can show that
$$\phi (\frac{a_1}{d},\mathrm{},\frac{a_k}{d},\frac{n}{d})=|\{x_n^k1x_ia_i,i[k],(x,n)=d\}|.$$
and thus
$$\underset{d|n}{}\phi (\frac{a_1}{d},\mathrm{},\frac{a_k}{d},\frac{n}{d})=a_1\mathrm{}a_k$$
(1)
and this implies
$$\phi (a_1,\mathrm{},a_k,n)=\underset{dn}{}\mu \left(d\right)\frac{a_1}{d}\mathrm{}\frac{a_k}{d},$$
and also for $`f:\{d:dn\}𝔸`$ an arbitrary map we have that
$$\underset{x_{a_1}\times \mathrm{}\times _{a_k}}{}f\left((x,n)\right)=\underset{dn}{}\phi (\frac{a_1}{d},\mathrm{},\frac{a_k}{d},\frac{n}{d})f(d).$$
###### Theorem 12.
Let $`n_+`$ be odd. The orbicycle index polynomial of $`D_n`$ is given by
$$P_{D_n}^{orb}=\frac{1}{n}\underset{dn}{}\phi (\frac{n1}{2d},\frac{n1}{d},\frac{n}{d})x_d^{\frac{n}{d}}\frac{1}{2}P__n+\frac{3}{2}x_1x_2^{\frac{n1}{2}}$$
###### Proof.
$`P_{D_n}^{orb}(x_1,\mathrm{},x_n)`$ $`=`$ $`{\displaystyle \frac{1}{2n}}{\displaystyle \underset{(\overline{g},h)P(D_n)}{}}|\overline{g}|{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(g)C(h))}`$
$`=`$ $`P_{D_n}+{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{\frac{n1}{2}}{}}}{\displaystyle \underset{j=0}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^i)C(\rho ^j))}+x_1x_2^{\frac{n1}{2}}`$
$`=`$ $`P_{D_n}+{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=0}{\overset{\frac{n1}{2}}{}}}{\displaystyle \underset{j=0}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^i)C(\rho ^j))}P__n+x_1x_2^{\frac{n1}{2}}`$
$`=`$ $`{\displaystyle \frac{1}{2}}P__n+{\displaystyle \frac{3}{2}}x_1x_2^{\frac{n1}{2}}+{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=0}{\overset{\frac{n1}{2}}{}}}{\displaystyle \underset{j=0}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^i)C(\rho ^j))}P__n`$
$`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{dn}{}}\phi ({\displaystyle \frac{n1}{2d}},{\displaystyle \frac{n1}{d}},{\displaystyle \frac{n}{d}})x_d^{\frac{n}{d}}{\displaystyle \frac{1}{2}}P__n+{\displaystyle \frac{3}{2}}x_1x_2^{\frac{n1}{2}}.`$
Recall that $`\tau `$ and $`\rho `$ are given $`\tau (x)=3x`$ and $`\rho ^r(x)=x+r.`$ Our next table gives the equivalence class of $`x_n`$ under five different equivalence relations.
| Partition | Equivalence class |
| --- | --- |
| $`C(\tau )C(\rho ^{\frac{n}{2}})`$ | $`\{x,3x,x+\frac{n}{2},3x+\frac{n}{2}\}`$ |
| $`C(\tau )C(\rho ^{\frac{n}{2}}\tau )`$ | $`\{x,3x,3x+\frac{n}{2},x+\frac{n}{2}\}`$ |
| $`C(\rho \tau )C(\rho ^{\frac{n}{2}})`$ | $`\{x,4x,x+\frac{n}{2},4x+\frac{n}{2}\}`$ |
| $`C(\rho \tau )C(\rho ^{\frac{n}{2}+1}\tau )`$ | $`\{x,4x,4x+\frac{n}{2},x+\frac{n}{2}\}`$ |
| $`C(\rho ^{\frac{n}{2}})C(\tau \rho ^i)`$ | $`\{x,x+\frac{n}{2},3ix,3ix+\frac{n}{2}\}`$ |
So we see that the equivalence class of $`x[n]`$ under the equivalence relation $`C(\rho ^{\frac{n}{2}})C(\tau \rho ^i)`$ is
$$\overline{x}=\{x,x+\frac{n}{2},3ix,3ix+\frac{n}{2}\},$$
so that $`|\overline{x}|\{2,4\}.`$ It is not difficult to see that $`|\overline{x}|=2`$ if and only if either $`2x3i`$ or $`2x3i+\frac{n}{2}.`$ Therefore $`b_2(C(\rho ^{\frac{n}{2}})C(\tau \rho ^i))`$ is $`1`$ if $`\frac{n}{2}`$ is odd, $`2`$ if $`\frac{n}{2}`$ is even and $`i`$ is odd, and $`0`$ if $`\frac{n}{2}`$ is even and $`i`$ is even.
###### Theorem 13.
Let $`n_+`$ be even. According to whether $`n`$ is $`0`$ or $`2`$ mod $`4`$, the orbicycle index polynomial of $`D_n`$ is given by
$`P_{D_n}^{orb}`$ $`=`$ $`P_{D_n}P__n+{\displaystyle \frac{1}{2n}}({\displaystyle \underset{d\frac{n}{2}}{}}\phi ({\displaystyle \frac{n1}{d}},{\displaystyle \frac{n}{2d}})x_{\frac{n}{d}}^d+\left\{\begin{array}{cc}nx_2x_4^{\frac{n2}{4}}& \\ \frac{n}{2}x_4^{\frac{n}{4}}+\frac{n}{2}x_2^2x_4^{\frac{n4}{4}}& \end{array}\right)`$
$`+`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{dn}{}}\phi ({\displaystyle \frac{n2}{2d}},{\displaystyle \frac{n1}{d}},{\displaystyle \frac{n}{d}})x_d^{\frac{n}{d}}+\left({\displaystyle \frac{n+2}{2n}}\right)(x_2^{\frac{n}{2}}+x_1^2x_2^{\frac{n2}{2}}+\{\begin{array}{cc}x_2x_4^{\frac{n2}{4}}& \\ x_4^{\frac{n}{4}}& \end{array}+\left\{\begin{array}{cc}x_2x_4^{\frac{n2}{4}}& \\ x_2^2x_4^{\frac{n}{4}1}& \end{array}\right).`$
###### Proof.
$`P_{D_n}^{orb}`$ $`=`$ $`{\displaystyle \frac{1}{2n}}{\displaystyle \underset{(\overline{g},h)P(D_n)}{}}|\overline{g}|{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(g)C(h))}`$
$`=`$ $`P_{D_n}+{\displaystyle \frac{1}{2n}}{\displaystyle \underset{hD_n}{}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^{\frac{n}{2}})C(h))}`$
$`+`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=0}{\overset{\frac{n2}{2}}{}}}{\displaystyle \underset{j=0}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^i)C(\rho ^j))}P__n`$
$`+`$ $`{\displaystyle \frac{n+2}{4n}}{\displaystyle \underset{h\{e,\tau ,\rho ^{\frac{n}{2}},\rho ^{\frac{n}{2}}\tau \}}{}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\tau )C(h))}`$
$`+`$ $`{\displaystyle \frac{n+2}{4n}}{\displaystyle \underset{h\{e,\rho \tau ,\rho ^{\frac{n}{2}},\rho ^{\frac{n}{2}+1}\tau \}}{}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho \tau )C(h))}.`$
We compute the last four summands in the expression above
$`{\displaystyle \underset{hD_n}{}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^{\frac{n}{2}})C(h))}`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^{\frac{n}{2}})C(\rho ^i))}+{\displaystyle \underset{i=0}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^{\frac{n}{2}})C(\tau \rho ^i))}`$
$`=`$ $`{\displaystyle \underset{i=0}{\overset{n1}{}}}x_{\frac{n}{(\frac{n}{2},i)}}^{(\frac{n}{2},i)}+{\displaystyle \underset{i=0}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^{\frac{n}{2}})C(\tau \rho ^i))}`$
$`=`$ $`{\displaystyle \underset{d\frac{n}{2}}{}}\phi ({\displaystyle \frac{n1}{d}},{\displaystyle \frac{n}{2d}})x_{\frac{n}{d}}^d+\{\begin{array}{cc}nx_2x_4^{\frac{n2}{4}}& \\ \frac{n}{2}x_4^{\frac{n}{4}}+\frac{n}{2}x_2^2x_4^{\frac{n4}{4}}& \end{array}`$
$`{\displaystyle \underset{i=0}{\overset{\frac{n2}{2}}{}}}{\displaystyle \underset{j=0}{\overset{n1}{}}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho ^i)C(\rho ^j))}`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{\frac{n2}{2}}{}}}{\displaystyle \underset{j=0}{\overset{n1}{}}}x_{\frac{n}{(i,j,n)}}^{(i,j,n)}={\displaystyle \underset{dn}{}}\phi ({\displaystyle \frac{n2}{2d}},{\displaystyle \frac{n1}{d}},{\displaystyle \frac{n}{d}})x_d^{\frac{n}{d}}`$
$`{\displaystyle \underset{h\{e,\tau ,\rho ^{\frac{n}{2}},\rho ^{\frac{n}{2}}\tau \}}{}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\tau )C(h))}=2x_2^{\frac{n}{2}}+2\{\begin{array}{cc}x_2x_4^{\frac{n2}{4}}& \\ x_4^{\frac{n}{4}}& \end{array}`$
$`{\displaystyle \underset{h\{e,\rho \tau ,\rho ^{\frac{n}{2}},\rho ^{\frac{n}{2}+1}\tau \}}{}}{\displaystyle \underset{k=1}{\overset{n}{}}}x_k^{b_k(C(\rho \tau )C(h))}`$ $`=`$ $`2x_1^2x_2^{\frac{n2}{2}}+2\{\begin{array}{cc}x_2x_4^{\frac{n2}{4}}& \\ x_2^2x_4^{\frac{n}{4}1}& \end{array}`$
In this work we have extended Pólya theory to the context of orbiquotient sets. The main ingredient of the new theory is the orbicycle index polynomial which we computed in various cases. We expect that our construction will find applications in the study of the topology of orbifolds and also in the theory of species. I would be interesting to search for a further extension of Pólya theory within the context of rational combinatorics introduced in based on the previous work and further discussed in . One should obtain a generalization of Pólya theory in which finite sets are replaced by finite groupoids .
### Acknowledgment
Thanks to Edmundo Castillo, Federico Hernandez, Eddy Pariguan, Sylvie Paycha, and Domingo Quiroz.
[email protected]
Departamento de Matemáticas Puras y Aplicadas,
Universidad Simón Bolívar
Caracas, Venezuela
[email protected]
Instituto de Matemáticas y sus Aplicaciones
Universidad Sergio Arboleda
Bogotá, Colombia.
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# BRUSLIB and NETGEN: the Brussels nuclear reaction rate library and nuclear network generator for astrophysics
## 1 Introduction
Since around the nineteen fifties, astrophysics has advanced at a remarkable pace, and has achieved an impressive record of success. One of the factors contributing to these rapid developments is undoubtedly a series of spectacular breakthroughs in nuclear astrophysics, which embodies the special interplay between nuclear physics and astrophysics.
The close relationship between these two major scientific disciplines comes about because of the clear demonstration that the Universe is pervaded with nuclear physics imprints at all scales. Starting with the Big Bang nucleosynthesis episode, the structure, evolution and composition of a large variety of cosmic objects, including the Solar System and its various constituants (down to meteoritic grains), bear strong imprints of the properties of atomic nuclei, as well as of their interactions. Therefore, careful and dedicated experimental and theoretical studies of a large variety of nuclear processes are indispensable tools for the modeling of the ultra-macroscopic systems astrophysics has to deal with.
Over the years, an impressive body of nuclear data of astrophysics interest have been obtained through laboratory efforts, complemented with theoretical developments. The latter are indispensable as nuclear experiments conducted today or expected to be performed in any foreseeable future can and will cover only a minute fraction of the astrophysics needs. This is because an extremely large variety of highly unstable ‘exotic’ nuclei that, for long, will not be produced and studied in the laboratory are expected to be involved in the modelling of many astrophysics processes and events. Many of the basic properties of these nuclei are to be known for this purpose, and so are their interactions, in particular those with nucleons or $`\alpha `$-particles. Even when laboratory-studied nuclei are considered, theory has very often to be called for assistance. In many respects, laboratory conditions are indeed very different from stellar ones, which are highly versatile and are often charaterized by high temperatures and/or densities that are out of reach of laboratory simulations. In addition, nuclear reactions between charged particles inside non-exploding stars take place in an energy regime that is in all but a few exceptional cases out of reach of direct experimental scrutiny. Indirect methods bring some complement of information, but clearly do not cover all the needs. In explosive situations, the energies of astrophysical interest are higher, and the cross sections are correspondingly larger. However, in such events, many reactions involve unstable targets, so that the fraction of those reactions of potential interest for which experimental reaction data are lacking is even larger.
The experimental and theoretical nuclear physics achievement mentioned above in fact let nuclear astrophysics face a new and difficult challenge. The rapidly growing wealth of nuclear data becomes, ironically, less and less easily accessible to the astrophysics community. Mastering this volume of information and making it available in an accurate and usable form for incorporation into astrophysics models become urgent goals of prime necessity. The establishment of the required level in the privileged communication between nuclear physicists and astrophysicists necessarily requires the build-up of well documented and evaluated sets of experimental data or theoretical predictions of astrophysical relevance. Some years ago, the Institut d’Astronomie et d’Astrophysique of the Université Libre de Bruxelles has decided to tackle this challenge. This work has grown into the BRUSsels nuclear LIBrary for astrophysics applications, referred to as BRUSLIB, and into a nuclear NETwork GENerator called NETGEN.
One goal of this paper is to *present for the first time a detailed description of BRUSLIB and of NETGEN in a journal and in a format such that these nuclear packages become easily accessible to astrophysicists* for a large variety of applications. In particular, the BRUSLIB reaction rates presented in a tabular form that is well suited to astrophysics needs are available electronically at the address http://www-astro.ulb.ac.be.
BRUSLIB is composed of two main parts. The first one concerns the Nuclear Astrophysics Compilation of REaction rates referred to as NACRE. This compilation provides the rates of 86 thermonuclear reactions of astrophysical relevance based on an in-depth analysis of experimental data. Its description is presented in Sect. 2. The second part describes theoretical evaluations of a collection of about 100000 rates for thermonuclear reactions induced by nucleons or $`\alpha `$-particles, as well as photo-induced reactions not included in NACRE, including nuclei with $`8Z110`$ lying between the proton and the neutron drip lines (Sect. 3). The rate calculations are based on the statistical Hauser-Feshbach (hereafter HF) model. They require the knowledge of a substantial amount of data concerning basic properties of the nuclei and of their interactions. The predictions of these properties rely on the use of global and coherent microscopic nuclear models. These utterly important model characteristics make the BRUSLIB rate library unique. A confrontation between selected experimental data and the predictions used in BRUSLIB is provided in Sect. 4.
A second goal is to describe NETGEN (Sect. 5), a package for constructing nuclear reaction networks on grounds of the nuclear physics input from BRUSLIB and, when necessary, from other sources.
It has to be noted that no information is provided by the present release of BRUSLIB on the rates of non-statistical (‘direct’) or of non-thermonuclear (‘spallation’) reactions that can develop in low-temperature and low-density astrophysical media, like the interstellar or circumstellar medium.
## 2 The experimentally-based NACRE evaluation and compilation of nuclear reaction rates
The necessity of placing relevant nuclear data of high quality at the disposal of the astrophysics community has been the driving motivation for the construction of a compilation aimed at superseding the work of Fowler and collaborators (see Caughlan & Fowler 1988, hereafter CF88, and references therein to former compilations). The goal was not just to update the CF88 rates with newly available experimental data, but also to modify quite deeply different aspects of the format of the CF88 compilation. At this point, slightly more than half of the CF88 rates have been re-compiled on grounds of a careful evaluation of experimental data available by 15 June 1998 , and the results of the work make the so-called NACRE (Nuclear Astrophysics Compilation of REactions) compilation (Angulo et al. 1999). The reactions analyzed up to now are listed in Table 1. They comprise an ensemble of 86 charged particle induced reactions on stable targets up to Si involved in Big Bang nucleosynthesis and in the non-explosive H- and He-burning modes, complemented with a restricted number of reactions of special astrophysical significance on the unstable $`{}_{}{}^{7}\mathrm{Be}`$, $`{}_{}{}^{22}\mathrm{Na}`$ and $`{}_{}{}^{26}\mathrm{Al}`$ nuclides. An updated and enlarged version of NACRE is currently in preparation.
NACRE is described in detail by Angulo et al. (1999). It contains in particular
(1) the formalism that has been adopted in order to derive the Maxwellian-averaged astrophysical rates of the (exothermic and endothermic) charged particle induced reactions and of their reverse;
(2) a general description of the treatment of the data, which has sometimes to be adapted to specific cases, as described in the comments accompanying each of the compiled reactions. If necessary, tables of narrow resonances with their characteristics are provided. A careful analysis of the experimental uncertainties in the quantities involved in the evaluation of the reaction rates is carried out for each reaction. From this, adopted rates are provided, as well as low and high limits. This large-scale estimate of the uncertainties is considered to be an essential feature of the NACRE compilation;
(3) the procedure adopted for extrapolating experimental data when required in order to evaluate the reaction rates. This extrapolation to the very low energies necessary to evaluate the rates down to the lowest considered temperatures raises many difficult problems. A particular one concerns the necessity to correct the laboratory cross sections for electron screening (e.g. Rolfs & Rodney 1988) before their use in the rate calculations. This effect becomes significant for relative energies $`E`$ of the reaction partners such that $`E/U_\mathrm{e}100`$, where $`U_\mathrm{e}`$ is the so-called screening potential. An approximate procedure is devised in order to eliminate data that can be ‘polluted’ by laboratory screening. On the other hand, in stellar conditions, the nuclei are surrounded by a dense electron gas that reduces the Coulomb repulsion and makes the penetration of the Coulomb barrier easier. The cross sections are therefore enhanced in comparison with those of reactions between bare nuclei. This stellar screening effect can be evaluated by applying, for example, the Debye-Hückel theory (e.g. Cox & Giuli 1968). The NACRE rates exclude these stellar screening factors.
The evaluation of the rates for temperatures as high as $`T_9=10`$ also requires the application of special extrapolation techniques when reliable cross section data are lacking at sufficiently high energies. In those cases, a duly documented HF approach is used (see also Sect. 3.1), which smoothly connects to the experimentally-based rate estimates;
(3) the procedure used to evaluate the contribution of excited states of the target nuclei to the effective stellar reaction rates. In a stellar plasma, the excited levels of a target nucleus are indeed thermally populated, and thus contribute to the reaction mechanism. As a result, the stellar rates may differ from those obtained when the target nuclei are in their ground state. This difference is expressed in terms of a correction factor $`r_{tt}`$ that has to multiply the target ground state rate in order to obtain the stellar rate. In general, this correction cannot be derived experimentally. The very approximate treatment of this correction by Caughlan & Fowler (1988) is replaced in NACRE by a more quantitative procedure based on the use of the HF model (Sect. 3.1), and on the classical assumption of a Maxwellian population of the nuclear excited states. Note that this assumption may be invalid if the target nucleus has an isomeric state which is not in thermal equilibrium with the ground state in certain astrophysical situations. Such an example of astrophysical interest concerns $`{}_{}{}^{26}\mathrm{Al}`$, to which a special treatment is applied.
The calculation of $`r_{tt}`$ requires in particular the knowledge of the temperature-dependent partition functions of the target nuclei normalized to their ground state
$`G_i(T)={\displaystyle \frac{1}{2J_i^0+1}}{\displaystyle \underset{\mu }{}}(2J_i^\mu +1)\mathrm{exp}\left({\displaystyle \frac{ϵ_i^\mu }{kT}}\right),`$ (1)
where the summation extends over the states $`\mu `$ of target $`i`$ with excitation energy $`ϵ_i^\mu `$ and spin-parity $`J_i^\mu `$ ($`\mu =0`$ for the ground state);
(4) The values of the adopted, low and high Maxwellian-averaged unscreened rates on the ground state target are provided in tabular form for a selection of temperatures in the $`0.001T_910`$ range ($`T_9`$ is the temperature in billion K). As they may be desired by some users, analytical formulae approximating the tabulated adopted ground state rates or those including the contribution of thermalized excited states are provided as well. It has to be noted that some of these expressions have a form deviating from those classically advertised in the nuclear astrophysics literature (e.g. Caughlan & Fowler 1988; Rolfs & Rodney 1988; see e.g. Eq. (23) of Angulo et al. 1999). They are considered to provide more secure approximations to the numerically calculated Maxwellian-averaged rates. In addition, NACRE provides analytical formulae for the multiplicative factor that has to be applied to the rate of a given reaction in order to evaluate the rate of the inverse process, as well as for the nuclear partition functions for the targets involved in the compiled reactions (see Eq. 1).
NACRE is accessible electronically through the BRUSLIB library website http://www-astro.ulb.ac.be. Much more material is available there than presented above or published by Angulo et al. (1999). In particular, reaction rates are available for far more extended temperature grids. In addition, graphical material (astrophysical $`S`$-factors) can be retrieved.
Finally, it is noted that other astrophysics-oriented experimentally-based reaction rate compilations have appeared around the time of the NACRE publication (Adelberger et al. 1998), or later (Iliadis et al. 2001; Descouvemont et al. 2004). The reactions considered in the latter two compilations are identified in Table 1, which also provides references to additional data or evaluations of relevance. This additional information will of course be duly analyzed and integrated into the new version of NACRE currently in preparation.
## 3 The BRUSLIB theoretical reaction rates
Much effort has been devoted in the last decades to measurements of reaction cross sections for astrophysical purposes. As already stated in Sect. 1, difficulties related to the conditions prevailing in the astrophysical plasmas remain, however. In particular, charged-particle induced reactions at stellar energies (far below the Coulomb barrier) have extremely small cross sections that are highly difficult to measure. Specific difficulties are also raised by the measurements of photoreactions that have to be included in many astrophysics models (e.g. Arnould & Goriely 2003 for a review). In addition, a huge number (thousands) of reactions of relevance involve more or less exotic nuclides. Clearly, many of these reaction rate experiments will remain unfeasible for a long time to come. Theory has thus to supply the necessary data, which represents a major challenge of its own.
BRUSLIB provides astrophysicists with a very extended set of thermonuclear rates in the $`0.01T_910`$ temperature range for all the reactions induced by neutron, proton and $`\alpha `$-particle captures by all nuclei with $`Z,N8`$ and $`Z<110`$ located between the neutron and proton drip lines (i.e some 8000 nuclei). The rates of the ($`\gamma `$,n), ($`\gamma `$,p) and ($`\gamma `$,$`\alpha `$) photodisintegrations of all these nuclides are also tabulated for the same temperature grid. The calculations rely on the code MOST (Goriely 1998; Arnould & Goriely 2003) based on the HF model (Hauser & Feshbach 1952). This approach is valid if the density of nuclear levels of the compound systems formed as a result of these captures is high enough. This is indeed the case if the targets are heavy enough and if they are located far enough from the proton or neutron drip lines to ensure that the excitation energies of the compound systems are high enough. Roughly speaking, the reliability of the BRUSLIB rates may thus be limited to nuclides with mass numbers $`A>\mathrm{\hspace{0.17em}\hspace{0.17em}40}`$ close enough to the valley of stability, this mass limit being of course shifted to higher and higher values as one moves further and further away from the valley. If these constraints are not met, a non-statistical treatment is more suited (Goriely 1997). It has to be noticed that purely theoretical HF rates are provided for some reactions already included in NACRE (Table 1). In these cases, it is advised to adopt the NACRE rates, or those from the Iliadis et al. (2001) or Descouvemont et al. (2004) compilations, as these rely on experimental data.
A few words are in order at this point to clarify the basic philosophy governing the selection of the models necessary for the calculation of the various ingredients entering the HF BRUSLIB rates. Their evaluation is based on global and coherent microscopic (or at least semi-microscopic) models. These features have an importance in the astrophysics modellings that cannot be underestimated, and make the BRUSLIB rate library unique. They are indeed not shared by HF calculations based on other codes, like Talys (Koning et al. 2002), Empire (Herman et al. 2002), or Non-Smoker (Rauscher & Thielemann 2001).
The global character of the underlying models is made highly desirable by the fact that, for many specific applications including nuclear astrophysics, a very large body of reaction rates for which no experimental data exist have to be provided. The microscopic nature of the underlying models is essential as well. For nuclear astrophysics, as well as in other fields, a large amount of data need to be extrapolated far away from experimentally known regions. In these situations, two major features of the nuclear theories have to be considered. The first one is the accuracy of a model. In most nuclear applications, this criterion has been the main, if not the unique, one for selecting a model. The second one is the reliability of the predictions. A physically sound model that is based on first principles and is as close as possible to a microscopic description of the nuclear systems is expected to provide the best possible reliability of extrapolations far away from experimentally known regions. Of course, the accuracy of such microscopic models in reproducing experimental data may be poorer than the one obtained from more phenomenological models in which enough free parameters can guarantee a satisfactory reproduction of the data at the expense of the quality of the input physics, and consequently of the reliability. The coherence (or ‘universality’) of these microscopic models (through e.g. the use of the same basic nuclear inputs, like the effective nuclear forces) is also required as different ingredients have to be predicted in order to evaluate each reaction rate. Failure to meet this requirement could lead to bad ways in the rate evaluations.
The BRUSLIB MOST rates are based on microscopic (or at least semi-microscopic) models for the necessary nuclear structure ingredients or interaction properties (Sect. 3.1) that have been constructed to be global and coherent to the largest possible extent. In addition, these models have the major advantage of reaching a satisfactory compromise between accuracy and reliability for the relevant nuclear inputs, and consequently for the rates themselves. The level of accuracy of these models is in fact fully comparable to the one obtained from available more or less highly parametrized phenomenological approaches, while their reliability is by far better.
Finally, the BRUSLIB reaction rates duly take into account the necessary astrophysical specificities. The temperature dependence of the rates is predicted from the consideration of the Maxwell-Boltzmann distribution of the relative velocities of the reaction partners and of the target nuclear excited states states (par. (3) of Sect. 2). In contrast, the provided rates are not corrected for stellar electron screening effects.
### 3.1 Reaction rate calculations: general framework
Some basics of the HF formalism adopted in the code MOST are briefly recalled in e.g. Arnould & Goriely (2003), and are not repeated here. The limits of its validity are also reminded above. Let us just point out the following:
(1) under local thermodynamic equilibrium conditions, the effective stellar rate of $`I+jL+k`$ per pair of particles in the entrance channel at temperature $`T`$ taking due account of the contributions of the various excited states $`\mu `$ of the target is expressed in a classical notation as
$`N_A\sigma v_{jk}^{}(T)=\left({\displaystyle \frac{8}{\pi m}}\right)^{1/2}{\displaystyle \frac{N_A}{(kT)^{3/2}G(T)}}\times `$ (2)
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{\mu }{}}{\displaystyle \frac{(2J_I^\mu +1)}{(2J_I^0+1)}}\sigma _{jk}^\mu (E)E\mathrm{exp}\left({\displaystyle \frac{E+\epsilon _I^\mu }{kT}}\right)dE,`$ (3)
where $`k`$ is the Boltzmann constant, $`m`$ the reduced mass of the $`I^0+j`$ system, $`N_A`$ the Avogadro number, $`\sigma _{jk}^\mu (E)`$ the cross section at relative energy $`E`$ of the $`I^\mu +jL+k`$ reaction, and $`G(T)`$ is the partition function given by Eq. 1, where $`J_I^\mu `$ and $`ϵ_I^\mu `$ are defined. The BRUSLIB photodisintegration rates are estimated by applying the reciprocity theorem on the radiative capture rates derived from Eq. 3. This procedure leads to
$`\lambda _{(\gamma ,j)}^{}(T)`$ $`=`$ $`{\displaystyle \frac{(2J_I+1)(2J_j+1)}{(2J_L+1)}}{\displaystyle \frac{G_I(T)}{G_L(T)}}\left({\displaystyle \frac{A_IA_j}{A_L}}\right)^{3/2}\times `$ (5)
$`\left({\displaystyle \frac{kT}{2\pi \mathrm{}^2N_A}}\right)^{3/2}N_A\sigma v_{(j,\gamma )}^{}\mathrm{e}^{Q_{j\gamma }/kT},`$
where $`Q_{j\gamma }`$ is the Q-value of the $`I^0(j,\gamma )L^0`$ capture. Note that, in stellar conditions, the reaction rates for targets in thermal equilibrium obey reciprocity since the forward and reverse channels are symmetrical, in contrast to the situation which would be encountered for targets in their ground states only (Holmes et al. 1976).
The uncertainties involved in any HF prediction are dominated by those involved in the evaluation of the nuclear quantities necessary for the calculation of the cross sections, such as the masses, deformations, matter distributions, single-particle levels, and level densities of target and residual nuclei, as well as the optical potentials. Special problems are also raised by the evaluation of the photon widths $`T_\gamma `$. The interested reader is referred to e.g. Arnould & Goriely (2003) for more details.
### 3.2 Nuclear masses, level densities, and partition functions
The BRUSLIB MOST predictions rely on the experimental nuclear mass data compiled by Audi et al. (2003). When not measured (a very common situation for various astrophysics applications), use is made of the HFB-9 microscopic mass model (Goriely et al. 2005). This model also provides the necessary information on nuclear deformation, charge and matter distributions, pairing properties and single-particle spectra. The nuclear level densities are extracted from a microscopic model developed by Goriely (1996) (see also e.g. Demetriou & Goriely 2001).
The nuclear partition functions $`G_i(T)`$ entering the evaluation of the astrophysical reaction rates (Sects. 2 and 3) are calculated from (i) a summation over experimentally known levels (Eq. 1) up to an excitation energy $`ϵ^\omega `$ above which the knowledge of the energy spectrum is considered to be incomplete, and (ii) a generalization of Eq. 1 involving an integration over a level density evaluated as described above.
### 3.3 Optical potentials
Phenomenological optical potentials (OPs) (generally of the Woods-Saxon type) may not be well suited for certain applications, particularly those involving exotic nuclei. It is considered profitable to use more microscopically-based potentials, whenever possible. A semi-microscopic OP, usually referred to as the JLM potential (Jeukenne et al. 1977), is available for the description of the nucleon-nucleus case. This OP has been revised recently for nucleons incident on spherical or quasi-spherical nuclei with masses $`40A209`$ at energies ranging from 1 keV to 200 MeV (Bauge et al. 2001). The resulting new version gives a global satisfactory agreement with experimental data, even if some improvements would be most welcome, especially in the low-energy domain and in the treatment of deformed or exotic nuclei. It is adopted for the BRUSLIB rate evaluations.
The situation for the $`\alpha `$-particle-nucleus OPs is much less satisfactory, and one still has to rely on phenomenological potentials. Most of the proposed OPs are derived from fits to elastic $`\alpha `$-nucleus scattering data at energies $`E>\mathrm{\hspace{0.17em}\hspace{0.17em}80}\mathrm{MeV}`$ or, in some cases, to (n,$`\alpha `$) cross sections at lower energies. However, the OP, and in particular its imaginary component, is known to depend strongly on energy below the Coulomb barrier. As a consequence, its extrapolation to sub-Coulomb energies that is necessitated by astrophysics applications is more insecure than in the case of nucleons. Several attempts to device a global $`\alpha `$-nucleus OP for the description of the scattering and reaction cross sections at energies $`E<\mathrm{\hspace{0.17em}\hspace{0.17em}20}`$ MeV of better relevance to astrophysics have been conducted (e.g. Arnould & Goriely 2003 for more details and references). The scarcity of experimental data, particularly in the $`A>100`$ mass range, limits dramatically the predictive power of any of the constructed global OPs. This has the immediate consequence of reducing the reliability of the rate predictions as they depend sensitively on the $`\alpha `$-particle-nucleus OPs. The BRUSLIB $`\alpha `$-capture rates are calculated with the global OP III developed by Demetriou et al. (2002).
### 3.4 $`\gamma `$-ray strength function
When applied to radiative captures, the total photon transmission coefficient entering the calculation of the $`\sigma _{j\gamma }^\mu `$ cross section of Eq. 3 is dominated by the $`E1`$ transitions. The calculation of the $`E1`$-strength function necessitates the knowledge of the low-energy tail of the Giant Dipole Resonance (GDR) of the compound system formed in the reaction process. The photon transmission coefficient is most frequently described in the framework of the phenomenological generalized Lorentzian model (e.g. Goriely 1998). The Lorentzian GDR approach suffers, however, from shortcomings of various sorts. On the one hand, it is unable to predict the enhancement of the $`E1`$ strength at energies below the neutron separation energy demonstrated by nuclear resonance fluorescence experiments. This departure from a Lorentzian profile may manifest itself more clearly for neutron-rich nuclei, and especially in the form of a so-called pygmy $`E1`$ resonance (Govaert et al. 1998, Zilges et al. 2002). On the other hand, even if a Lorentzian function provides a suitable representation of the $`E1`$ strength, the location of its maximum and its width remain to be predicted from some underlying model for each nucleus. For astrophysical applications, these properties have often been obtained from a droplet-type of model (Myers et al. 1977), which clearly lacks reliability when dealing with exotic nuclei. This introduces large uncertainties in certain rate estimates (e.g. Goriely & Khan 2002, Goriely et al. 2004)
In view of this situation, it is clearly of substantial interest to develop models of the microscopic type which are hoped to provide a reasonable reliability and predictive power for the $`E1`$-strength function. Various attempts of this sort have been conducted (e.g. Arnould & Goriely 2003, and references therein). The BRUSLIB MOST rate calculations are based on the semi-microscopic QRPA $`E1`$ model developed by Goriely et al. (2004).
## 4 A confrontation between measured and calculated reaction rates
In order to evaluate the overall quality of the BRUSLIB reaction rate predictions, this section presents a confrontation between selected experimental data and calculations in which only the ground state contribution ($`\mu =0`$) is taken into account in Eqs. 3 and 5, the consideration of target excited states being irrelevant in the laboratory conditions.
Figure 1 compares the BRUSLIB predictions of the Maxwellian-averaged (n,$`\gamma `$) rates $`\sigma v`$ at $`T=3.5\times 10^8`$ K with experimental data for some 228 nuclei heavier than $`{}_{}{}^{40}\mathrm{Ca}`$ included in the compilation of Bao et al. (2000). It appears that the calculations agree with all data to within a factor of three. Figure 2 compares the experimental cross sections for some low-energy (p,$`\gamma `$) reactions on targets heavier than Fe with the corresponding BRUSLIB rates. Again, and broadly speaking, the agreement is very satisfactory for the vast majority of the (p,$`\gamma `$) data. It is difficult, however, to be much more specific, as the quality of this agreement may depend more or less drastically on temperature, in contrast to the situation encountered in the neutron capture case. In addition, the comparison is limited to targets up to Sn. Experiments with heavier nuclei would be most welcome, as they might help refining the predictions.
The situation is by far less clear for the ($`\alpha `$,$`\gamma `$) reactions. This results from the lack of a large enough body of experimental data for sub-Coulomb cross sections combined with the difficulties to construct global and reliable $`\alpha `$-nucleus OPs (Sect. 3.3). These theoretical problems are magnified by the fact that, at the sub-Coulomb energies of astrophysical relevance, the reaction rate predictions are highly sensitive to these potentials through the corresponding $`\alpha `$-particle transmission coefficients. Figure 3 displays a comparison between some low-energy measurements of ($`\alpha `$,$`\gamma `$) cross sections and MOST predictions used for evaluating the BRUSLIB rates. The quality of the agreement appears to vary from case to case, and is also highly sensitive to temperature. Uncertainties in the calculated values originate mainly from those in the nuclear level densities and $`\alpha `$-nucleus OPs (see Arnould & Goriely 2003 for details, and in particular for a discussion of the $`{}_{}{}^{144}\mathrm{Sm}(\alpha ,\gamma ){}_{}{}^{148}\mathrm{Gd}`$ reaction of special astrophysics interest).
## 5 The nuclear network generator NETGEN
The Nuclear Network Generator NETGEN is an interactive Web-based tool to generate Maxwellian-averaged nuclear reaction rates for networks and temperature grids specified by the user. It is fully documented at the web address http://www-astro.ulb.ac.be/Netgen.
NETGEN relies mainly on the BRUSLIB/NACRE library. It also makes use of
(1) the post-NACRE compilations by Iliadis et al. (2001) and Descouvemont et al. (2004) (see Table 1);
(2) experimentally-based published rates for about 70 charged-particle induced reactions not included in NACRE. Specific references are provided for each of these reactions;
(2) more than 200 experimentally-based radiative neutron capture rates. Most of these are adopted from Bao et al. (2000). Specific references are given in each case;
(3) $`\beta `$-decay and electron-capture rates, including (i) laboratory measurements compiled by Horiguchi et al. (1996), (ii) theoretical estimates based on the evaluation of individual transitions (Takahashi & Yokoi 1987), on the gross theory (Tachibana et al. 1990), or on the more microscopic ETFSI + cQRPA model (Borzov & Goriely 2000).
For the sake of completeness, the references of Table 1 that are not contained in the NACRE, Iliadis et al. (2001) or Descouvemont et al. (2004) compilations are also included in NETGEN.
For each reaction or $`\beta `$-decay rate, NETGEN selects by default the data source that is considered to be the most reliable. We select in order of preference the latest available compilation (NACRE, Iliadis et al. 2001, or Descouvemont et al. 2004), experimental data, detailed microscopic calculations, and BRUSLIB rates derived from global calculations. The user may nevertheless adopt another choice for selected cases by specifying ‘bibliographic indexes’ for each reaction, the table of rates being accompanied by a ‘log file’ listing the selected data source among all those available in the library. A FORTRAN program handling the rates is also made available. Note that all rates duly take into account the contribution from the excited states of the target nuclei, as discussed in the previous sections.
Various NETGEN options are currently offered through the web interface, as described in detail at http://www-astro.ulb.ac.be/Netgen:
(1) Generate a table of reaction rates on a temperature grid for a network that has been
* typed in reaction by reaction (offering the possibility to select non-default rates, as mentioned above), or
* generated automatically between interactively-selected boundaries on the proton and mass numbers, and involving various possible sets of reactions (p-, n-, $`\alpha `$-captures, $`\beta `$-decays and/or photodisintegrations), or
* uploaded on the server.
(2) Plot individual reaction rates, and provide .gif or .ps files.
## 6 Conclusions
This is the first release in an astronomy and astrophysics journal of the BRUSLIB nuclear reaction rate library and of the nuclear network generator NETGEN. The format of the packages is chosen in order to make them easily accessible to astrophysicists and to be well suited for a large variety of astrophysics needs. They are made available through the web site http://www-astro.ulb.ac.be.
The BRUSLIB NACRE package contains a detailed experimentally-based evaluation and compilation of the rates of 86 proton or $`\alpha `$-capture reactions on (mainly) stable targets up to Si for temperatures ranging from $`10^6`$ to $`10^{10}`$ K. The electronic files contain much more information than published by Angulo et al. (1999).
The NACRE data are complemented with about 100000 thermonuclear rates of nucleon and $`\alpha `$-captures on about 8000 $`8Z110`$ nuclei located between the proton and neutron drip lines. The calculations are based on a statistical Hauser-Feshbach model featuring a microscopic (or at least very close to microscopic) evaluation of the basic ingredients of the model. These predictions are seen to compare favourably with the limited set of experimental reaction cross section data on intermediate-mass and heavy nuclei at energies close to those of astrophysics relevance. The rates of photodisintegration of the whole set of nuclei are also provided. They are derived from the application of the reciprocity theorem.
NETGEN is an interactive web-based tool allowing the construction on a user-friendly basis of nuclear reaction networks specified by the user on a temperature grid of his/her choice. A full documentation of its use can be found at the web address http://www-astro.ulb.ac.be/Netgen.
It is hoped that the easy availability of a very large set of nuclear reaction rate evaluations and predictions, as well as of a nuclear network generator will be helful to many researchers for a large variety of applications. Quite clearly, this enterprise is of a highly dynamical and long-term nature. BRUSLIB will be continuously improved and expanded, and new releases will be made accessible to the community every time a substantial enough body of new data becomes available.
###### Acknowledgements.
The authors thank all the collaborators who have made possible the development over the years of the Brussels library of nuclear data, and in particular M. Pearson, P. Demetriou and E. Khan. This work has been supported in part by the Interuniversity Attraction Pole IAP 5/07 of the Belgian Federal Science Policy and by the Konan University - Université Libre de Bruxelles convention ‘Construction of an Extended Nuclear Database for Astrophysics’. S.G. is FNRS Research Associate. A.J. is FNRS Senior Research Associate.
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# Gravitational contribution to fermion masses
## I Introduction
Conceived as an alternative to the standard general relativistic metric approach to gravity, gauge theories of spacetime groups describe gravitational forces in close analogy to the remaining interactions Utiyama:1956sy Sciama1 Sciama2 Kibble:1961ba Hayashi:1980wj Lord:1987uq Lord:1988nd Hehl:1995ue Gronwald:1995em . The Lorentz group and the $`GL(4,R)`$ group are usual candidates proposed by different authors Sciama1 Sciama2 Ivanenko:1983vf Sardanashvily:2002mi to play the role of local symmetries. Instead, Hehl et al. Kibble:1961ba Hehl:1995ue Hehl:1974cn Hehl:1976kj consider gravity as the gauge theory of either the Poincar é or the affine group: in any case of a group including translations. Actually, the interpretation of tetrads as a certain kind of translational connections allows an uniform description of all known interactions, gravity included, in terms of gauge potentials declared as the unique force mediators Julve:1994bh Lopez-Pinto:1995qb Lopez-Pinto:1997aw Tresguerres:2000qn .
We are interested in analyzing the consequences for matter fields of considering translations included in the gauge group, as for instance in the Poincaré gauge theory (PGT) of gravity, where the full Poincaré group is treated as the local gauge group of a Yang-Mills type theory. Given that such approach is a serious candidate to become the fundamental theory of gravity, obviously we must know how the corresponding Poincaré covariant derivatives of matter fields look like, with both the homogeneous Lorentz group contributions and those of translations taken into account. The present paper is devoted to give an answer to the question how translational connections couple in particular to Dirac fields.
Independently of the interest of PGT in itself, the fact that we choose it with preference to a more general gauge theory of gravitation, such as metric-affine gravity (MAG), is partly determined by a technical reason, namely the possibility of building an explicit matrix representation of the Poincaré algebra. In fact, besides the usual spin operators $`\sigma _{\alpha \beta }`$ constituting the representation of the Lorentz generators acting on Dirac fields, one can introduce the complementary realization $`\pi _\mu `$ of the translational generators. The affine group is more problematic to deal with due to the fact that no finite dimensional spinor representation of $`GL(4,R)`$ exists Hehl:1995ue .
As an unexpected consequence of the explicit construction of covariant Poincaré derivatives with intrinsic translations, we find that the translational connections contribute to the Dirac action with a fermion mass term of PGT-gravitational nature. Such result is exclusive for a certain kind of gauge theories of gravity, having nothing to do either with ordinary General Relativity or with gauge approaches based on spacetime groups not including translations. More precisely, we derive the background fermion mass from the nonlinear approach to PGT established by us in a number of previous papers Julve:1994bh Lopez-Pinto:1995qb Lopez-Pinto:1997aw Tresguerres:2000qn rrdph . There we developed a suitable treatment of spacetime groups with translations, explaining the identification of tetrads as (nonlinear) translational connections, and one of us proposed an adapted fibre bundle description Tresguerres:2002uh .
In next section we review a few main concepts, necessary to deduce the key formula (26) expressing nonlinear connections in terms of linear ones. Then in section III we apply the nonlinear procedure to the Poincaré group, paying special attention to covariant derivatives of Dirac fields, and finally in section IV we build the matter action showing the emergence of the translation-induced mass term.
## II Generalized bundle structure of gauge theories
### II.1 Composite fiber bundles
The ordinary gauge theory of a given Lie group $`G`$ is known to have the structure of a principal bundle $`P(M,G)`$ equipped with a connection, being matter fields defined on associated bundles. However, gauge theories involving nonlinearly realized local symmetries, as for instance gauge theories of spacetime groups, require a slight modification of this bundle scheme, as discussed in Tresguerres:2002uh . The composite fiber bundles studied there are particularly suitable to highlight the underlying geometry of gauge theories of groups including translations, such as the Poincaré gauge theory of gravity, thus constituting the main support of the present paper. Let us briefly remaind the reader on its defining features. For what follows, see Sardanashvily:1992fq Sardanashvily:1994 Sardanashvily:1994fg Sardanashvily:1995ew , as much as Kobayashi:1963 , pgs. 54 and 57.
Let $`\pi _{_{PM}}:PM`$ be a principal fiber bundle with structure Lie group $`G`$, and let $`H`$ be a closed subgroup of $`G`$. The quotient space $`G/H`$ constitutes a manifold on which $`G`$ acts on the left in a natural way. Then it is possible to build the $`P`$-associated bundle $`\pi _{_{\mathrm{\Sigma }M}}:\mathrm{\Sigma }M`$ with standard fiber $`G/H`$ and with total space consisting of the quotient space $`\mathrm{\Sigma }=(P\times G/H)/G`$ of the Cartesian product $`P\times G/H`$ by the right action of $`G`$ defined as $`P\times G/H(u,\xi )(ug,g^1\xi )P\times G/H,gG`$. The total space $`\mathrm{\Sigma }`$ can be identified with the quotient space $`P/H`$ of $`P`$ by the right action of $`H`$ on $`P`$, and consequently one finds $`P(\mathrm{\Sigma },H)`$ to be a principal fiber bundle with structure group $`H`$ and with well defined projection $`\pi _{_{P\mathrm{\Sigma }}}:P\mathrm{\Sigma }`$ onto the base space $`\mathrm{\Sigma }=P/H`$, see Prop. 5.5 of Kobayashi:1963 . Indeed, each orbit $`uH`$ through $`uP`$ –diffeomorphic to the standard fiber $`H`$– projects into a single element (a left coset) of $`P/H`$.
Nonlinearly realized gauge theories to be studied here differ from ordinary gauge theories in that they are based on principal bundles $`P(M,G)`$ whose structure group $`G`$ is reducible to a closed subgroup $`H`$. According to Prop. 5.6 of Kobayashi:1963 , such reducibility of the structure group $`G`$ to $`HG`$ is guaranteed if and only if a cross section $`s_{_{M\mathrm{\Sigma }}}:M\mathrm{\Sigma }=P/H`$ of the associated bundle $`\mathrm{\Sigma }`$ exists. Furthermore, there is a one to one correspondence between sections $`s_{_{M\mathrm{\Sigma }}}`$ and the reduced subbundles of $`\pi _{_{P\mathrm{\Sigma }}}:P\mathrm{\Sigma }`$ consisting of the set of points $`uP`$ such that
$$\pi _{_{P\mathrm{\Sigma }}}(u)=s_{_{M\mathrm{\Sigma }}}\pi _{_{PM}}(u),$$
(1)
see Kobayashi:1963 . From condition (1) follows trivially
$$\pi _{_{PM}}=\pi _{_{\mathrm{\Sigma }M}}\pi _{_{P\mathrm{\Sigma }}},$$
(2)
providing a decomposition of the total projection $`\pi _{_{PM}}`$ into partial projections. Accordingly, the principal bundle $`\pi _{_{PM}}:PM`$ transforms into the composite bundle
$$\pi _{_{\mathrm{\Sigma }M}}\pi _{_{P\mathrm{\Sigma }}}:P\mathrm{\Sigma }M.$$
(3)
In (3) we distinguish two bundle sectors, characterized respectively by the partial projections
$$\pi _{_{P\mathrm{\Sigma }}}:P\mathrm{\Sigma },\pi _{_{\mathrm{\Sigma }M}}:\mathrm{\Sigma }M.$$
(4)
The latter one, with standard fiber $`G/H`$, can be seen as an intermediate space, in the sense that it is built over the primary base space $`M`$, and simultaneously plays a role as the base space of the principal bundle $`\pi _{_{P\mathrm{\Sigma }}}:P\mathrm{\Sigma }`$ with structure group $`H`$. More precisely, in the context of composite bundles one can regard the $`G`$-diffeomorphic fibers of $`P(M,G)`$ as being, say, bent into two sectors, corresponding respectively to the fibers $`H`$ of $`\pi _{_{P\mathrm{\Sigma }}}:P\mathrm{\Sigma }`$ and $`G/H`$ of $`\pi _{_{\mathrm{\Sigma }M}}:\mathrm{\Sigma }M`$. The $`H`$-diffeomorphic fiber branches are attached to points of the intermediate base space $`\mathrm{\Sigma }`$, which trivialize locally as $`(x,\xi )`$, with $`\xi `$ coordinatizing the fiber branches $`G/H`$ attached to $`xM`$.
In parallel to (2), the local sections $`s_{_{MP}}:MP`$ are decomposed as
$$s_{_{MP}}=s_{_{\mathrm{\Sigma }P}}s_{_{M\mathrm{\Sigma }}},$$
(5)
see Section V of Tresguerres:2002uh . In terms of suitable zero sections, denoting as $`\sigma _{_{MP}}`$ those corresponding to $`s_{_{MP}}`$, and so on, the sections in (5) become respectively
$$s_{_{MP}}=R_{\stackrel{~}{g}}\sigma _{_{MP}},\stackrel{~}{g}G,$$
(6)
$$s_{_{\mathrm{\Sigma }P}}=R_a\sigma _{_{\mathrm{\Sigma }P}},aH,$$
(7)
and
$$s_{_{M\mathrm{\Sigma }}}=R_b\sigma _{_{M\mathrm{\Sigma }}},bG/H.$$
(8)
Conditions
$$\stackrel{~}{g}=ba,\sigma _{_{\mathrm{\Sigma }P}}R_b=R_b\sigma _{_{\mathrm{\Sigma }P}},$$
(9)
ensure that, in analogy to (5), the relation
$$\sigma _{_{MP}}=\sigma _{_{\mathrm{\Sigma }P}}\sigma _{_{M\mathrm{\Sigma }}}$$
(10)
also holds. The usefulness of this structure will become evident in the following.
In summary, composite fiber bundles (3) provide the mathematical foundation for gauge theories involving nonlinear gauge realizations (as the generalization of induced representations). Relevant physical theories comprised among the concerned ones are on the one hand the standard model –since a correspondence between nonlinear realizations and spontaneous symmetry breaking exists Cho:1978ss – and on the other hand nonlinear gauge theories of gravity, as developed below. Nonlinear realizations characteristic for such theories take place on principal fiber bundles $`P(M,G)`$ whose structure group $`G`$ is reducible to a closed subgroup $`HG`$. While the total symmetry remains that of the gauge group $`G`$, one exploits the possibility of working with the explicit symmetry $`H`$ of the principal subbundle of $`\pi _{_{P\mathrm{\Sigma }}}:P\mathrm{\Sigma }`$ whose base space is the total space of $`\pi _{_{\mathrm{\Sigma }M}}:\mathrm{\Sigma }M`$. (The sections $`s_{_{M\mathrm{\Sigma }}}`$ defined on the latter bundle are identified as Goldstone fields Sardanashvily:1992fq .)
### II.2 Nonlinear realizations in composite bundles
In Tresguerres:2002uh it was showed that the composite bundle structure defined by (4)–(10) provides the natural framework to deal with nonlinear gauge realizations, exactly as standard principal bundles constitute the arena for the ordinary gauge treatment of groups. Actually, the main results on nonlinear realizations Coleman:1969sm Callan:1969sn Salam:1969rq Isham:1971dv Borisov74 Stelle:1980aj are easily derived. So, the nonlinear gauge transformation equation
$$L_g\sigma _{_{\mathrm{\Sigma }P}}(x,\xi )=R_h\sigma _{_{\mathrm{\Sigma }P}}(x,\xi ^{})$$
(11)
is obtained by comparing two bundle elements, both with the form (7), differing from each other by the left action $`L_g`$ of elements $`gG`$, the latter being local in the sense that $`g=g(x)`$, $`xM`$, see Tresguerres:2002uh for details. Regarding (11) referred to the base space $`M`$, it manifests itself as a vertical bundle automorphism not affecting $`xM`$, in analogy to ordinary gauge transformations. Nevertheless, when referred to the intermediate base space $`\mathrm{\Sigma }M\times G/H`$, the action of $`L_g`$ not only transforms the sections $`\sigma _{_{\mathrm{\Sigma }P}}`$ vertically along the $`H`$ fiber branches by means of $`R_h`$, $`hH`$, but simultaneously it induces a transformation affecting the points $`(x,\xi )\mathrm{\Sigma }`$, thus mapping $`H`$ fiber branches into fiber branches defined on different $`\mathrm{\Sigma }`$-points (as expected for spacetime groups, in particular for translations, see (33) below).
In order to deal with ordinary geometrical objects defined on the base manifold $`M`$, we pull back to the latter, by means of $`s_{_{M\mathrm{\Sigma }}}`$, the quantities defined on the plateau $`\mathrm{\Sigma }`$. Taking into account the property of pullbacks when applied to functions $`\phi `$, namely $`f^{}\phi =\phi f`$, we first define the pullback of $`\sigma _{_{\mathrm{\Sigma }P}}`$ as
$$\sigma _\xi (x):=(s_{_{M\mathrm{\Sigma }}}^{}\sigma _{_{\mathrm{\Sigma }P}})(x)=\sigma _{_{\mathrm{\Sigma }P}}s_{_{M\mathrm{\Sigma }}}(x).$$
(12)
Then we calculate $`s_{_{M\mathrm{\Sigma }}}^{}(L_g\sigma _{_{\mathrm{\Sigma }P}})=L_g\sigma _{_{\mathrm{\Sigma }P}}s_{_{M\mathrm{\Sigma }}}=L_g\sigma _\xi `$ and $`s_{_{M\mathrm{\Sigma }}}^{}(R_h\sigma _{_{\mathrm{\Sigma }P}}^{})=R_h\sigma _{_{\mathrm{\Sigma }P}}^{}s_{_{M\mathrm{\Sigma }}}=R_h\sigma _\xi ^{}`$, so that (11) gives rise to
$$L_g\sigma _\xi (x)=R_h\sigma _\xi ^{}(x).$$
(13)
In (13) (Eq.(6.6) of Tresguerres:2002uh ) one recognizes the fundamental equation for nonlinear realizations Coleman:1969sm Lord:1987uq Lord:1988nd .
The nonlinear gauge transformations of fields induced by (13) are deduced in Section VIII of Tresguerres:2002uh . Taking in (13) $`hI+\mu `$ to be infinitesimal, with $`\mu `$ defined on the Lie algebra of $`H`$, the fields $`\psi (\sigma _\xi (x)):=(\sigma _\xi ^{}\psi )(x)`$ of any given representation space of $`H`$ are found to transform infinitesimally under $`G`$ as
$$\delta \psi (\sigma _\xi (x)):=\sigma _\xi ^{}^{}\psi (L_g\sigma _\xi )^{}\psi \rho (\mu )\psi (\sigma _\xi (x)),$$
(14)
being $`\rho (\mu )`$ the suitable representation of the $`H`$-algebra element $`\mu `$. (See Eqs.(8.9), (8.11) of Tresguerres:2002uh .) Eqs. (13) and (14) summarize the main results of Coleman:1969sm . In the nonlinear approach, the relevant fact is that the fields $`\psi `$ of representation spaces of $`HG`$ also constitute a representation space for the nonlinear action (13) of the full group $`G`$.
### II.3 Bundle approach to nonlinear connections and covariant derivatives
Covariant derivatives of the fields in (14) require the introduction of suitable (nonlinear) connections on $`M`$. As a crucial result for this purpose, in the present paragraph we will derive equation (26) below, implicit in Tresguerres:2002uh but not explicitly given there, expressing the nonlinear connections in terms of standard (linear) gauge potentials.
Depending on the bundle base space we consider, that is, either $`M`$ or the plateau $`\mathrm{\Sigma }`$, at least two alternative expressions can be given for the Ehresmann connection form. On the one hand, taking the quantities in (6) into account,
$$\omega =\stackrel{~}{g}^1(d+\pi _{_{PM}}^{}A__M)\stackrel{~}{g},$$
(15)
involving the ordinary gauge potential $`A__M`$ on the base space $`M`$, defined as the pullback
$$A__M=\sigma _{_{MP}}^{}\omega .$$
(16)
On the other hand, with (7) at view,
$$\omega =a^1(d+\pi _{_{P\mathrm{\Sigma }}}^{}\mathrm{\Gamma }__\mathrm{\Sigma })a,$$
(17)
where we introduce the nonlinear connection on the intermediate space $`\mathrm{\Sigma }`$, turning out to be the pullback
$$\mathrm{\Gamma }__\mathrm{\Sigma }=\sigma _{_{\mathrm{\Sigma }P}}^{}\omega .$$
(18)
Since $`\stackrel{~}{g}=ba`$, see (9), comparison of (15) and (17) yields
$$\pi _{_{P\mathrm{\Sigma }}}^{}\mathrm{\Gamma }__\mathrm{\Sigma }=b^1(d+\pi _{_{PM}}^{}A__M)b.$$
(19)
From the defining condition $`\pi _{_{P\mathrm{\Sigma }}}\sigma _{_{\mathrm{\Sigma }P}}=id_\mathrm{\Sigma }`$ for sections follows $`\sigma _{_{\mathrm{\Sigma }P}}^{}\pi _{_{P\mathrm{\Sigma }}}^{}=id_{T^{}(\mathrm{\Sigma })}`$, so that (19) gives rise to
$$\mathrm{\Gamma }__\mathrm{\Sigma }=\sigma _{_{\mathrm{\Sigma }P}}^{}[b^1(d+\pi _{_{PM}}^{}A__M)b].$$
(20)
We operate on (20) taking into account that, in terms of the pulled back quantity
$$\widehat{b}(x,\xi ):=b\sigma _{_{\mathrm{\Sigma }P}}(x,\xi ),$$
(21)
the relations $`\sigma _{_{\mathrm{\Sigma }P}}^{}(b^1db)=\widehat{b}^1d\widehat{b}`$, and $`\sigma _{_{\mathrm{\Sigma }P}}^{}R_b^{}=R_{\widehat{b}}^{}\sigma _{_{\mathrm{\Sigma }P}}^{}`$ hold, being $`b^1\pi _{_{PM}}^{}A__Mb=R_b^{}\pi _{_{PM}}^{}A__M`$ and, in view of (2), $`\pi _{_{PM}}^{}=\pi _{_{P\mathrm{\Sigma }}}^{}\pi _{_{\mathrm{\Sigma }M}}^{}`$. We find
$$\mathrm{\Gamma }__\mathrm{\Sigma }=\widehat{b}^1(d+\pi _{_{\mathrm{\Sigma }M}}^{}A__M)\widehat{b}.$$
(22)
Pulling back (22), defined on $`\mathrm{\Sigma }`$, by means of $`s_{_{M\mathrm{\Sigma }}}^{}`$, compare with (12), we get
$$\mathrm{\Gamma }__M=s_{_{M\mathrm{\Sigma }}}^{}\mathrm{\Gamma }__\mathrm{\Sigma }$$
(23)
as the nonlinear connection, defined on the base space $`M`$, we will deal with in the following. Obviously, in view of (18) and (12)
$$\mathrm{\Gamma }__M=s_{_{M\mathrm{\Sigma }}}^{}\sigma _{_{\mathrm{\Sigma }P}}^{}\omega =\sigma _\xi ^{}\omega .$$
(24)
Analogous calculations to those leading from (20) to (22) allow us to find, in terms of the new pulled back quantity
$$\stackrel{~}{b}(x):=b\sigma _{_{\mathrm{\Sigma }P}}s_{_{M\mathrm{\Sigma }}}(x)=b(\sigma _\xi (x)),$$
(25)
the relation
$$\mathrm{\Gamma }__M=\stackrel{~}{b}^1(d+A__M)\stackrel{~}{b},$$
(26)
between the nonlinear connection $`\mathrm{\Gamma }_M`$ and the linear connection $`A_M`$, that is, between the alternative pullbacks (24) and (16) of the connection 1-form $`\omega `$ to $`M`$. Our deduction of (26) provides a geometrical interpretation of eqs. (19) of Coleman:1969sm , (6) of Cho:1978ss and (2.7) of Tseytlin:1982nu , while it shows the incompleteness of eqs. (2.15) of Salam:1969rq or (22) of Borisov74 . The importance of (26) for what follows becomes evident in view of the transformation properties of $`\mathrm{\Gamma }__M`$, given in Eq.(7.14) of Tresguerres:2002uh , namely
$$\delta \mathrm{\Gamma }__M=\sigma _\xi ^{}^{}\omega (L_g\sigma _\xi )^{}\omega (d\mu +[\mathrm{\Gamma }__M,\mu ]),$$
(27)
with $`\mu `$ being the same $`H`$-algebra element as in (14). Eq. (27) shows that only the $`H`$-algebra components of $`\mathrm{\Gamma }__M`$ still transform inhomogenously as $`H`$-connections, while the $`G/H`$-algebra components transform as $`H`$-tensors. According to (14) and (27), the nonlinear covariant differential defined as
$$D\psi :=(d+\rho (\mathrm{\Gamma }__M))\psi $$
(28)
transforms as an $`H`$-covariant differential
$$\delta D\psi =\rho (\mu )D\psi $$
(29)
under nonlinear gauge transformations (13) of the full group $`G`$. The general procedure established here will be applied in the next section to $`G`$ as the Poincaré group and $`H`$ as the Lorentz group in order to derive PGT.
## III Nonlinear Poincaré gauge theory of gravity
### III.1 Poincaré covariant derivatives
The main results of the previous section are summarized in the transformation law (13) and the induced field transformation (14), plus the relation (26) between the nonlinear connection (24) and the linear one (16), with the corresponding nonlinear connection transformation (27). In terms of these elements, one defines the covariant differential (28) transforming as (29).
Now, in order to perform explicit calculations, we need to transform (13) into a more manageable formula. From (12) with (8), (9) and (10), we get $`\sigma _\xi (x)=R_b\sigma _{_{MP}}(x)`$. (In the latter equation we identify $`b=\sigma _{_{MP}}^1(x)\sigma _\xi (x)=b(\sigma _\xi (x))=:\stackrel{~}{b}(x)`$ as given by (25).) Analogously, $`\sigma _\xi ^{}(x)=R_b^{}\sigma _{_{MP}}(x)`$. Replacing these values into (13), it follows $`L_gR_b\sigma _{_{MP}}(x)=R_hR_b^{}\sigma _{_{MP}}(x)`$. Finally, since $`\sigma _{_{MP}}^1(x)g\sigma _{_{MP}}(x)=g`$, we find
$$gb=b^{}h.$$
(30)
Eq. (30) is the form of (13) appearing in Coleman:1969sm , appropriate for practical computational purposes, with $`b`$ being (25) and thus identical with $`\stackrel{~}{b}`$ in (26).
Now we merely apply the general formalism mechanically to the gauge group $`G=`$ Poincaré, with $`H=`$ Lorentz. (Other choices of $`H`$ have been studied elsewhere Lopez-Pinto:1997aw .) In (30) we replace the infinitesimal group elements $`gI+iϵ^\mu P_\mu +i\beta ^{\alpha \beta }L_{\alpha \beta }`$ of the Poincaré group and $`hI+i\mu ^{\alpha \beta }L_{\alpha \beta }`$ of the homogeneous Lorentz group, and we parametrize $`b`$ and $`b^{}`$ respectively as $`b=e^{i\xi ^\mu P_\mu }`$ with finite translational parameters $`\xi ^\mu `$, and $`b^{}=e^{i\xi _{}^{}{}_{}{}^{\mu }P_\mu }`$ with $`\xi _{}^{}{}_{}{}^{\mu }\xi ^\mu +\delta \xi ^\mu `$. Then, taking into account the Poincaré commutation relations
$$[L_{\alpha \beta },L_{\mu \nu }]=i(o_{\alpha [\mu }L_{\nu ]\beta }o_{\beta [\mu }L_{\nu ]\alpha }),$$
(31)
$$[L_{\alpha \beta },P_\mu ]=io_{\mu [\alpha }P_{\beta ]},[P_\mu ,P_\nu ]=0,$$
(32)
with the help of the Hausdorff-Campbell formula, (30) yields on the one hand the value $`\mu ^{\alpha \beta }=\beta ^{\alpha \beta }`$ for the $`H`$-parameter, and on the other hand
$$\delta \xi ^\mu =\xi ^\nu \beta _\nu {}_{}{}^{\mu }ϵ^\mu .$$
(33)
Observe how the transformations (33) of the translational parameters resemble those of Cartesian coordinates.
Let us now pay attention to the connections. Starting with the ordinary linear ones for the Poincaré group, say
$$A__M=i\mathrm{\Gamma }^{\alpha \beta }L_{\alpha \beta }i\stackrel{\left(T\right)}{\mathrm{\Gamma }^\mu }P_\mu ,$$
(34)
we make use of (26) to construct, in terms of (34) and of $`b=e^{i\xi ^\mu P_\mu }`$, the nonlinear connections
$$\mathrm{\Gamma }__M=i\mathrm{\Gamma }^{\alpha \beta }L_{\alpha \beta }i\vartheta ^\mu P_\mu ,$$
(35)
where simple calculations yield for the nonlinear translational connection the structure
$$\vartheta ^\mu :=D\xi ^\mu +\stackrel{\left(T\right)}{\mathrm{\Gamma }^\mu },$$
(36)
being $`D\xi ^\mu :=d\xi ^\mu +\mathrm{\Gamma }_\nu {}_{}{}^{\mu }\xi _{}^{\nu }`$. More explicitly, since all quantities are pulled back to the base space $`M`$, (36) reads
$$\vartheta ^\mu =dx^i(_i\xi ^\mu +\mathrm{\Gamma }_{i\nu }{}_{}{}^{\mu }\xi _{}^{\nu }+\stackrel{\left(T\right)}{\mathrm{\Gamma }_i^\mu })=:dx^ie_i{}_{}{}^{\mu },$$
(37)
where we introduce the usual notation $`e_i^\mu `$ for vierbeins in order to show the identification we make of the nonlinear translational connections with the tetrads. Such interpretation of tetrads is possible since, in view of (27), they obey the gauge transformations
$$\delta \vartheta ^\mu =\vartheta ^\nu \beta _\nu {}_{}{}^{\mu }.$$
(38)
In addition we find for the Lorentz part of (35)
$$\delta \mathrm{\Gamma }_\alpha {}_{}{}^{\beta }=D\beta _\alpha {}_{}{}^{\beta }.$$
(39)
As a consistence condition for (33), (36), (38) and (39) follows the transformation of the linear translational connection
$$\delta \stackrel{\left(T\right)}{\mathrm{\Gamma }^\mu }=\stackrel{\left(T\right)}{\mathrm{\Gamma }^\nu }\beta _\nu {}_{}{}^{\mu }+Dϵ^\mu .$$
(40)
Comparison of (40) with the transformations (38) of the nonlinear translational connections (36) clarify why the latter, as a result of the nonlinear approach, can play the role of tetrads. Actually, tetrad variations (38) constitute a particular case of the above mentioned fact that the nonlinear connection components associated to generators of $`G`$ not belonging to $`H`$ behave as $`H`$-tensors.
With the previous results at hand, the main task of the present paragraph is to construct the Poincaré covariant derivatives of matter fields. As shown by (14), the gauge action of the full Poincaré group $`G`$ takes place through the representation $`\rho (\mu )=i\mu ^{\alpha \beta }\rho (L_{\alpha \beta })`$ of the algebra of the Lorentz group $`H`$, acting on fields of arbitrary representation spaces of $`H`$. In particular, for Dirac fields we take the spinor generators $`\rho (L_{\alpha \beta })=\sigma _{\alpha \beta }`$ as given by (47) below, being $`\mu ^{\alpha \beta }=\beta ^{\alpha \beta }`$ as mentioned just before (33). We find
$$\delta \psi =i\beta ^{\alpha \beta }\sigma _{\alpha \beta }\psi .$$
(41)
The covariant derivative (28) of such fields, although resembling an ordinary $`H`$-covariant differential, is build with a nonlinear connection defined on the whole $`G`$-algebra. Thus, a representation of the full Poincaré algebra is required in order to realize the nonlinear connection (35) as
$$\rho (\mathrm{\Gamma }__M)=i\mathrm{\Gamma }^{\alpha \beta }\sigma _{\alpha \beta }i\vartheta ^\mu \pi _\mu ,$$
(42)
where $`\pi _\mu =\rho (P_\mu )`$ is the finite matrix representation of translational generators to be studied below. According to the general formula (28), the Poincaré covariant derivatives of Dirac fields read
$$D\psi =d\psi i(\mathrm{\Gamma }^{\alpha \beta }\sigma _{\alpha \beta }+\vartheta ^\mu \pi _\mu )\psi ,$$
(43)
transforming in analogy to (41) as
$$\delta D\psi =i\beta ^{\alpha \beta }\sigma _{\alpha \beta }D\psi .$$
(44)
Certainly, due to the particular nonlinear Poincaré transformations (38) and (41), the contributions associated to the translational generators are not necessary to guarantee covariance of (43). Nevertheless, the general scheme requires these contributions to be present in the otherwise Lorentz covariant derivatives, as an unavoidable heritage of the gauged Poincaré group. So we need to know how the $`G`$ generators not belonging to $`H`$ act on the fields $`\psi `$ of the representation space of $`H`$. In our case, this means that, besides (47), we have to look for the already mentioned representation of the translational generators in order to complete the finite matrix realization of the abstract Poincaré algebra (31), (32).
### III.2 Intrinsic translations of fermion fields
According to our conventions, the Dirac gamma matrices are defined so that their product reads
$$\gamma _\alpha \gamma _\beta =o_{\alpha \beta }I4i\sigma _{\alpha \beta },$$
(45)
expressed in terms of the Minkowski metric
$$o_{\alpha \beta }:=diag(+++)$$
(46)
and of the spinor generators
$$\sigma _{\alpha \beta }:=\frac{i}{8}[\gamma _\alpha ,\gamma _\beta ]$$
(47)
of the Lorentz group, being $`\sigma _{\alpha \beta }=\rho (L_{\alpha \beta })`$ the usual $`4\times 4`$ matrix representation of the Lorentz algebra (31) acting on 4-dimensional Dirac bispinors $`\psi `$. Let us discuss how to extend the Lorentz algebra to the Poincaré algebra, the latter one constituting a subalgebra of the conformal algebra as shown in the appendix.
The possibility of constructing also intrinsic translational operators $`\pi _\mu =\rho (P_\mu )`$ from the gamma matrices rests on the fact that
$$[\sigma _{\alpha \beta },\gamma _\mu ]=io_{\mu [\alpha }\gamma _{\beta ]},$$
(48)
and on the properties of the $`\gamma _5`$ matrix, defined as
$$\gamma _5:=i\gamma ^0\gamma ^1\gamma ^2\gamma ^3,$$
(49)
such that $`\gamma _5^2=I`$ and satisfying the commutation relations
$$[\sigma _{\alpha \beta },\gamma _5]=0,$$
(50)
and the anticommutation relations
$$\{\gamma _\mu ,\gamma _5\}=0,$$
(51)
and
$$\{\sigma _{\alpha \beta },\gamma _\mu \}=\frac{1}{2}\eta _{\alpha \beta \mu }{}_{}{}^{\nu }\gamma _{\nu }^{}\gamma _5,$$
(52)
(where $`\eta _{\alpha \beta \gamma \delta }`$, with $`\alpha ,\beta \mathrm{}=0,\mathrm{},3`$, is defined so that $`\eta _{0abc}=ϵ_{abc}`$, with $`a,b,c=1,2,3`$). Making use of these elements, one finds operators
$$\pi _\mu \gamma _\mu (\mathrm{\hspace{0.17em}1}+\lambda \gamma _5)$$
(53)
to exist, with $`\lambda ^2=1`$, satisfying the commutation relations
$$[\sigma _{\alpha \beta },\pi _\mu ]=io_{\mu [\alpha }\pi _{\beta ]},[\pi _\mu ,\pi _\nu ]=0,$$
(54)
characteristic for translational generators, see (32). Notice that eqs. (54) do not completely determine $`\pi _\mu `$. Actually, in (53) a global factor as much as the sign $`\lambda `$ $`(=\pm 1)`$ remain unfixed. This fact reflects the existence of two inequivalent realizations of the full conformal algebra, of which the Poincaré algebra is a subalgebra. Invoking dimensionality consistence of the intrinsic linear momentum $`\pi _\mu `$ with the orbital linear momentum $`i_\mu `$ we require the former, in natural units $`\mathrm{}=c=1`$, to have dimensions $`[L]^1`$. Since the gamma matrices in (53) are dimensionless, we are enforced to introduce a dimensional constant, say $`m[L]^1`$. Let us also fix the undetermined sign in (53), see the appendix, and define
$$\pi _\mu :=\frac{m}{4}\gamma _\mu (\mathrm{\hspace{0.17em}1}+\gamma _5),$$
(55)
where the numerical factor is introduced for later convenience.
A remarkable feature of (55) is that $`\pi _\mu \pi _\nu =0`$. The resulting anticommutation relations $`\{\pi _\mu ,\pi _\nu \}=0`$ are compatible with the finite matrix realization of the Poincaré algebra given by (47) and (55). Since the commutation relations alone are responsible for the transformations (33) of the coordinate-like parameters the matter fields depend on, they suffice to induce the change from $`\psi (\sigma _\xi (x))`$ into $`\psi (\sigma _\xi ^{}(x))`$ where their gauge variation (14) is evaluated.
On the other hand, the usual Casimir characterization of mass still holds in our scheme despite the nilpotence of $`\pi _\mu `$ by considering the complete translational generators as consisting of the sum of an orbital plus an intrinsic contribution, namely $`P_\mu =iI_\mu +\pi _\mu `$. Observe that, in the limit of vanishing components of the Lorentz connections, the translational parameters $`\xi ^\mu `$ become indistinguishable from Cartesian coordinates and the covariant derivative (43) reduces to the action of such a $`P_\mu `$ on fermions as $`id\xi ^\mu P_\mu \psi `$. Since $`\pi _\mu `$ is traceless and $`\pi _\mu \pi _\nu =0`$, the Casimir relation $`Tr(P_\mu P^\mu )m^2`$ is valid for $`m0`$.
Our intrinsic translational generators (55) resemble the momentum spin introduced by Gürsey Gursey:1964 in the context of the contraction of $`O(3,2)`$ to the Poincaré group Inonu:1953sp . Indeed, such momentum spin is conceived as the intrinsic part of the pseudotranslational generators $`\mathrm{\Pi }_\mu :=(1/R)L_{5\mu }`$ whose commutation relations, in the limit $`R\mathrm{}`$, reproduce those of Poincaré translations.
## IV Poincaré gauge invariant Dirac action
The discussion of previous section guarantees the translational contributions in (43) not only to make sense, but to be an essential part of (nonlinear) Poincaré covariant derivatives. Thus we have all the elements needed to build the Dirac matter action in the presence of gravity, when the latter is described by (nonlinear) PGT. Following the notation of Hehl:1990yq , with $`\gamma :=\vartheta ^\mu \gamma _\mu `$, and $`{}_{}{}^{}\gamma `$ its Hodge dual, the Dirac Lagrange density 4-form –without explicit mass term– reads
$$L_D=\frac{i}{2}\overline{\psi }^{}\gamma D\psi +h.c.,$$
(56)
with the usual definition $`\overline{\psi }:=\psi ^{}\gamma ^0`$, and $`h.c.`$ standing for the Hermitian conjugate of the given term. Let us calculate the latter in order to make all our conventions explicit. From (45) we get $`\gamma _0^2=1`$. Provided
$$\gamma ^0\gamma _\mu ^{}\gamma ^0=\gamma _\mu ,$$
(57)
as it is the case for instance for the Dirac representation of gamma matrices in terms of Pauli matrices as
$$\gamma ^0=\left(\begin{array}{cc}I\hfill & 0\hfill \\ 0\hfill & I\hfill \end{array}\right),\gamma ^a=\left(\begin{array}{cc}0\hfill & \sigma ^a\hfill \\ \sigma ^a\hfill & 0\hfill \end{array}\right),\gamma _5=\left(\begin{array}{cc}0\hfill & I\hfill \\ I\hfill & 0\hfill \end{array}\right),$$
(58)
we realize that
$$(\frac{i}{2}\overline{\psi }^{}\gamma D\psi )^{}=\frac{i}{2}\overline{D\psi }^{}\gamma \psi ,$$
(59)
with $`\overline{D\psi }:=(D\psi )^{}\gamma ^0`$. Furthermore, (45) with (57) yields
$$\gamma ^0\sigma _{\alpha \beta }^{}\gamma ^0=\sigma _{\alpha \beta },$$
(60)
guaranteeing the invariance of (56) by enforcing $`\overline{\psi }`$ to transforms as
$$\delta \overline{\psi }=(\delta \psi )^{}\gamma ^0=i\overline{\psi }\beta ^{\alpha \beta }\sigma _{\alpha \beta },$$
(61)
and on the other hand from definition (49) with (57) we get
$$\gamma ^0\gamma _5^{}\gamma ^0=\gamma _5.$$
(62)
Applying (57) and (62) to (55), it follows
$$\gamma ^0\pi _\mu ^{}\gamma ^0=\pi _\mu ,$$
(63)
a result which was not a priori obvious. Taking (60) and (63) into account, from (43) we find
$$\overline{D\psi }:=(D\psi )^{}\gamma ^0=d\overline{\psi }+i\overline{\psi }(\mathrm{\Gamma }^{\alpha \beta }\sigma _{\alpha \beta }+\vartheta ^\mu \pi _\mu ),$$
(64)
transforming as
$$\delta \overline{D\psi }=i\overline{D\psi }\beta ^{\alpha \beta }\sigma _{\alpha \beta },$$
(65)
compare with (61). If desired, in order to take into account other forces besides gravitation, one can extend the gauge symmetry replacing the Poincaré group by the direct product of Poincaré times an internal group. To do so, one merely has to replace (43) by
$$D\psi =d\psi +i(gA\mathrm{\Gamma }^{\alpha \beta }\sigma _{\alpha \beta }\vartheta ^\mu \pi _\mu )\psi $$
(66)
(and analogously (64)) without affecting what follows. In view of the previous results, the explicit form of (56) becomes
$$L_D=\frac{i}{2}(\overline{\psi }^{}\gamma D\psi +\overline{D\psi }^{}\gamma \psi ).$$
(67)
Let us separate the translational parts, no more indispensable for the covariance of the covariant derivatives, from (43), (resp. (66)) as
$$D\psi =:\stackrel{~}{D}\psi i\vartheta ^\mu \pi _\mu \psi ,$$
(68)
and analogously
$$\overline{D\psi }=:\overline{\stackrel{~}{D}\psi }+i\overline{\psi }\vartheta ^\mu \pi _\mu ,$$
(69)
see (64), where we denote with tildes the translations-independent pieces. Replacing (68) and (69) in (67), the Lagrange density transforms into
$$L_D=\frac{i}{2}(\overline{\psi }^{}\gamma \stackrel{~}{D}\psi +\overline{\stackrel{~}{D}\psi }^{}\gamma \psi )+^{}m\overline{\psi }\psi ,$$
(70)
where we made use of the fact that $`\vartheta ^\alpha ^{}\vartheta _\beta =\delta _\beta ^\alpha \eta `$, with $`\eta =^{}1`$ as the 4-dimensional volume element, so that
$${}_{}{}^{}\gamma \vartheta ^\mu \pi _\mu =\eta \gamma ^\mu \pi _\mu =^{}m(1+\gamma _5),$$
(71)
and
$$\vartheta ^\mu \pi _\mu ^{}\gamma =\eta \pi _\mu \gamma ^\mu =^{}m(1\gamma _5).$$
(72)
Although $`\gamma _5`$ is necessary to guarantee the commutation relations (54) to hold, both contributions (71) and (72) are combined in the action in such a way that $`\gamma _5`$ cancels out. So the matter Lagrange density (70) merely retains a mass term, which is unavoidable since it derives from the translational contribution to the Poincaré connection (42). Accordingly, either one of the projections $`\psi _L`$ or $`\psi _R`$ is lacking (in which case $`\overline{\psi }\psi =0`$), or otherwise the field $`\psi `$ is necessarily massive.
## V Conclusions
Independently from other possible origins of fermion masses, a gravitational background mass contribution is predicted by PGT when treated as a nonlinear local realization of the Poincaré group. Provided both left and right projections of Dirac fields are simultaneously present, (70) prevents massless Dirac fields from existing. The irremovable fermion masses are a consequence of gravitational interaction (in particular of the underlying translational group) in the context of PGT as the fundamental theory of gravity.
As a phenomenological consequence, when considered together with the standard model, PGT gives rise to a background contribution of gravitational origin to the masses of all fermions: in particular to the quark mass parameters of the QCD sector of the Lagrangian, as much as to the neutrino masses. Neutrinos are thus predicted by PGT to be massive. Certainly, our approach does not determine the value of the universal translational mass parameter $`m`$. However, from the observed masses of neutrinos it is clear that $`m`$ (the same for all fermions) has to be very small, so that, accordingly, its contribution to the observable hadron masses is expected to be quite limited.
Matter currents corresponding to the Poincaré symmetry are the spin current $`\tau _{\alpha \beta }:=L_D/\mathrm{\Gamma }^{\alpha \beta }`$ and the energy-momentum 3-form $`\mathrm{\Sigma }_\mu :=L_D/\stackrel{\left(T\right)}{\mathrm{\Gamma }^\mu }=L_D/\vartheta ^\mu `$. The former is found to be $`\tau _{\alpha \beta }=\frac{1}{4}\overline{\psi }\vartheta _\alpha \vartheta _\beta \gamma \gamma _5\psi `$ as it is well known. Its coupling term to the Lorentz connection $`\mathrm{\Gamma }^{\alpha \beta }`$ falls off from the action in the limit of absence of gravity (that is for $`\mathrm{\Gamma }^{\alpha \beta }=0`$, $`\stackrel{\left(T\right)}{\mathrm{\Gamma }^\mu }=0`$). Instead, the mass term does not cancel out in this limit. The reason is that, according to the nonlinear approach to PGT, the tetrads have the structure (36), not vanishing for zero linear connections. Actually, ordinary Minkowskian flat spacetime may be regarded as the residual structure left by nonlinear PGT in the absence of the gravitational force carried by spin connections, that is in the limit where the components of the latter ones are chosen to vanish. The tetrads are in this case $`\vartheta ^\mu =d\xi ^\mu `$, so that the mass term associated to them still remains in the action despite translational linear connections are switched out.
*
## Appendix A The O(2 ,4) and the Poincaré algebra
The Poincaré algebra is a subalgebra of the conformal algebra Mack to be examined here. Consider the $`O(2,4)`$ generators $`L_{AB}=L_{BA}`$, $`A,B\mathrm{}=0,\mathrm{},3,5,6`$, satisfying the commutation relations
$$[L_{AB},L_{MN}]=i(g_{A[M}L_{N]B}g_{B[M}L_{N]A}),$$
(73)
where the six-dimensional metric tensor is taken to be
$$g_{AB}=diag(+++,+).$$
(74)
In order to relate (73) to the ordinary form of the conformal commutation relations, let us decompose (74) into the Minkowski metric
$$g_{\alpha \beta }=o_{\alpha \beta }:=diag(+++),$$
(75)
where $`\alpha ,\beta =0,\mathrm{},3`$, plus
$$g_{55}=1,g_{66}=1,$$
(76)
and define the translational generators
$$P_\mu :=L_{\mu 5}+L_{\mu 6},$$
(77)
the special conformal generators
$$K_\mu :=L_{\mu 5}L_{\mu 6},$$
(78)
and the dilatational generators
$$D:=2L_{56}.$$
(79)
In terms of $`L_{\alpha \beta }`$, (77), (78) and (79), the commutation relations (73) give rise to the conformal algebra
$$[L_{\alpha \beta },L_{\mu \nu }]=i(o_{\alpha [\mu }L_{\nu ]\beta }o_{\beta [\mu }L_{\nu ]\alpha }),$$
(80)
$$[L_{\alpha \beta },P_\mu ]=io_{\mu [\alpha }P_{\beta ]},$$
(81)
$$[L_{\alpha \beta },K_\mu ]=io_{\mu [\alpha }K_{\beta ]},$$
(82)
$$[P_\mu ,K_\nu ]=i(L_{\mu \nu }+\frac{1}{2}o_{\mu \nu }D),$$
(83)
$$[D,P_\mu ]=iP_\mu ,$$
(84)
$$[D,K_\mu ]=iK_\mu ,$$
(85)
$$[P_\mu ,P_\nu ]=[K_\mu ,K_\nu ]=[D,L_{\mu \nu }]=[D,D]=0.$$
(86)
As pointed out in Mack , all finite dimensional representations of the $`O(2,4)`$ algebra can be obtained by reducing out tensor products of two inequivalent fundamental 4-dimensional representations (corresponding respectively to the choices $`\lambda =\pm 1`$ in what follows) builded from the gamma matrices as
$$\rho (L_{\alpha \beta })=\sigma _{\alpha \beta }:=\frac{i}{8}[\gamma _\alpha ,\gamma _\beta ],$$
(87)
$$\rho (L_{\mu 5})=\frac{1}{2}(\pi _\mu +\kappa _\mu )=\lambda \frac{m}{4}\gamma _\mu \gamma _5,$$
(88)
$$\rho (L_{\mu 6})=\frac{1}{2}(\pi _\mu \kappa _\mu )=\frac{m}{4}\gamma _\mu ,$$
(89)
$$\rho (L_{56})=\frac{1}{2}\mathrm{\Delta }=\lambda \frac{i}{4}\gamma _5.$$
(90)
Obviously, as read out from (88), (89) and (90), the corresponding fundamental inequivalent representations of (77), (78) and (79) read
$$\pi _\mu :=\frac{m}{4}\gamma _\mu (\mathrm{\hspace{0.17em}1}+\lambda \gamma _5),$$
(91)
$$\kappa _\mu :=\frac{m}{4}\gamma _\mu (\mathrm{\hspace{0.17em}1}\lambda \gamma _5),$$
(92)
and
$$\mathrm{\Delta }:=\lambda \frac{i}{2}\gamma _5,$$
(93)
where the role of $`\pi _\mu `$ and $`\kappa _\mu `$ is interchangeable by fixing $`\lambda `$ to be either $`\pm 1`$, and by accordingly change the sign of (93). The Poincaré algebra considered in the main text is the subalgebra of the conformal algebra consisting of the spin and translational generators only, having fixed $`\lambda =1`$.
###### Acknowledgements.
The authors are very grateful to Friedrich Wilhelm Hehl and to Yuri Obukhov for helpful and clarifying discussions.
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# 1 Introduction
## 1 Introduction
### 1.1 Quantum theory on noncommutative space time
Much of the development of quantum theories on noncommutative space-times was and still is driven by the question whether noncommutative geometry might lead to an ultra-violet regularization of quantum field theory, as it was suggested by Heisenberg as early as 1938 . For the simplest conceivable examples of noncommutative geometries, where the commutator of the coordinates is constant, these hopes for regularization were met with disappointment. Such theories exhibit a mixing of ultra-violet and infra-red cutoff scales , which has been understood recently on the level of renormalization . For that reason these simple noncommutative spaces, while interesting objects on their own, hardly seem to improve the divergent behavior of quantum field theory at all. The natural next step is to turn again to more complicated noncommutative spaces, such as quantum spaces with Lie type or homogeneous commutation relations.
Quantum Minkowski space is one of the most realistic examples of a quantum deformation of space-time with a rich and fairly well understood mathematical structure: It is four dimensional in the sense that it is generated as algebra by four coordinates such that the ordered monomials are a Poincaré-Birkhoff-Witt basis. The time coordinate is central, which is important for a causal interpretation of quantum theory . By construction, it is a module algebra with respect to the quantum Poincaré algebra , so it has a well defined quantum symmetry structure . The action of the generators of the inhomogeneous part of the quantum Poincaré algebra, the momenta, on quantum Minkowski space induces a first order covariant differential calculus .
With these key mathematical structures present, much of the construction of quantum theory on commutative space-time was mimicked within an algebraic and representation theoretic approach. Sections of noncommutative space-time were given by spectra of noncommutative coordinates, which yield discrete space-time lattices . Free elementary particles were constructed as irreducible representations of the quantum Poincaré algebra, including Wigner representations of spin . The free particles were shown to obey quantum wave equations which project reducible quantum Lorentz spinor representations on their irreducible subrepresentations . The coupling of free fields to gauge fields was studied systematically within noncommutative gauge theory . Quantum statistics of the tensor representation of a second quantized free theory was studied within the framework of braided tensor categories using Drinfeld twists . Only a small selection of the numerous contributions to this program can be cited here. At the time of this writing, most aspects of free quantum theory on quantum Minkowski space are understood, and the construction of interaction terms on the level of single particle field equations is known. What is not known is how to solve these field equations by noncommutative wave functions, not even for the free case. However, this will be indispensable if free particles are to be coupled by multiplying their noncommutative wave functions within the algebra of quantum space-time.
We remark, that there are other approaches which involve wave equations on noncommutative space-time which are inequivalent to the ones considered here: In covariance arguments have been used to find wave equations for a certain deformation of the conformal symmetry algebra, which does not contain the $`q`$-Poincaré algebra we consider here. In order to construct $`q`$-deformed relativistic wave equations various other methods have been proposed, based on $`q`$-Clifford algebras , $`q`$-deformed co-spinors , or various abstract differential calculi on quantum spaces , each leading to different results.
### 1.2 Noncommutative differential calculus
The free wave equations on quantum Minkowski space can be determined uniquely on a purely representation theoretic level. To give the simplest example: Just as in the undeformed case, the momentum square $`p_\mu p^\mu `$ is central within the quantum Poincaré algebra. Therefore, it must be represented within the irreducible representation corresponding to a free particle by a multiple of the identity operator,
$$p_\mu p^\mu =m^2,$$
(1)
$`m^2`$ being the square of the particle’s mass. This line of reasoning is purely representation theoretic and does not depend in any way on how the momenta act on noncommutative wave functions, that is, on elements of quantum Minkowski space.
In order to interpret Eq. (1) as Klein-Gordon equation on noncommutative space-time we have to let the momenta act on noncommutative wave functions. Momenta then correspond to partial derivatives, $`p^\mu =\mathrm{i}^\mu `$, defined by the noncommutative first order differential calculus. This yields a wave equation given by a noncommutative partial differential equation,
$$_\mu ^\mu \psi =m^2\psi ,$$
(2)
where $`\psi `$ is an element of the algebra of quantum Minkowski space. There is a considerable amount on literature on first order differential calculi on noncommutative algebras, mostly on quantum groups and to some extent on homogeneous quantum spaces . Most mathematical work has concentrated on the construction, structural analysis, and classification of differential calculi. However, hardly any literature on the subject has dealt with concrete calculations within these differential calculi. On quantum Minkowski space a differential calculus has been known for some time, and can be deduced most elegantly from the coproduct of momenta . The coproduct defines the action of partial derivatives on noncommutative functions recursively. While the recursion relations can be used to expand the derivatives in terms of nested summations and partition functions, as it was carried out in a very detailed manner in , this does not solve the main computational problem which arises with wave equations like Eq. (2).
Trying to solve Eq. (2) by a brute force calculation with a general ansatz for the wave function ought to work in principle, but turns out to run into considerable computational complexity. Moreover, the solutions thus obtained in terms of recursion relations do not give much of structural insight in the solutions which, after all, are the noncommutative counterpart of something as simple and basic as plane waves. Why is this computation so difficult? It can be shown that just as in the commutative case the space of solutions of the massive Klein-Gordon Eq. (2) is generated, as representation of the quantum Poincaré algebra, by a rest state, which satisfies
$$^0\psi =\mathrm{i}m\psi ,^A\psi =0,$$
(3)
where $`A`$ is a three vector index of spatial coordinates. In the commutative case we could now infer that $`\psi `$ is a function of the time coordinate $`x_0`$ alone, thus reducing Eq. (3) to an ordinary differential equation in one coordinate $`x_0`$. Not so in the noncommutative case, where the partial derivatives of a function in $`x_0`$ depend on all coordinates $`x^\mu `$, which makes the solution of Eq. (3) so involved. In analogy to partial differential equations in non-cartesian coordinates we will describe this mixing of dependencies by saying that within the noncommutative calculus the standard variables are not separated.
It is now obvious to ask, whether new variables can be found which separate the differential calculus, in such a way that Eq. (3) becomes easily solvable. The main purpose of this paper is to show that this question has a positive answer.
### 1.3 Structure, main results, notation
The paper is organized as follows: In Sec. 2 the structure of the differential calculus on quantum Minkowski space is reviewed. We recall the commutation relations, general structure of the quantum Lorentz algebra, universal $``$-matrix, and the definition of the coproduct of momenta in order to make this paper self contained and fix conventions and notation unambiguously.
Sec. 3 contains the bulk of the paper: the computation of derivatives of arbitrary elements in the quantum Minkowski algebra. We first develop a notation and calculus in order to deal with functions of the algebra valued 4$`\times `$4-matrices which come from the action of the universal $`L`$-matrices appearing in the definition of the coproduct of momenta. The computation problem is here the calculation of powers of algebra valued matrices, which can then used to calculate the derivatives of arbitrary functions of the noncommutative coordinates. We do this calculation in two steps: In the first step we deal with polynomials of the non-central variables. The main result is in Proposition 1, Eqs. (69) and can be written in the concise form of a chain rule, where the outer derivatives are partial Jackson derivatives. In the second step we calculate derivatives of functions of the central coordinates, time and four length. The partial derivatives of functions in the time generator are complicated and depend on all other generators. However, we are able to give a noncommutative coordinate transformation such that the derivatives of functions in the central generators are disentangled and take the form of a chain rule, too. This is the main result of the paper, presented in Proposition 3, Eq. (97). The results of these two steps can be combined to yield a compact formula for the derivatives of arbitrary elements of the algebra, given in Eq. (101).
In Sec. 4 we apply the formulas for the derivations in order to the initial problem of finding the momentum eigenstate solutions of the massless and massive quantum Klein-Gordon equation, viewed as differential equations within the noncommutative differential calculus. We show that the separation of variables of the preceding section separates the differential equations completely. This reduces the partial differential equations to single variable equations involving Jackson derivatives, which are easily solved. The solutions are given in Eqs. (110) and (117) in terms of $`q`$-exponential functions.
In Sec. 5 we assess the results, give an outlook and references to further work in this direction of this paper.
Appendix A finally contains a few more technical results which are used in the paper. This makes this paper reasonably self-contained and will be useful for a reader who wishes to reproduce the calculations in detail.
Notation: Throughout the paper $`q`$ denotes a real deformation parameter, such that for $`q1`$ we retrieve the undeformed, commutative limit. We will often use the abbreviations $`\lambda :=qq^1`$ and symmetrized and standard quantum numbers
$$[n]=\frac{q^nq^n}{qq^1},(n)_{q^2}=\frac{q^{2n}1}{q^21},$$
(4)
where $`n`$ is a natural number. Repeated upper and lower indices are summed over $`x_\mu x^\mu _\mu x_\mu x^\mu `$.
## 2 Differential calculus on quantum Minkowski space
### 2.1 Quantum Minkowski space
Quantum Minkowski space has first been constructed as braided product of two copies of the quantum plane . We will need the definition in terms of generators which are decomposed into the scalar time coordinate and the three vector of space coordinates.
###### Definition 1.
The complex algebra generated by the four generators $`x_0`$, $`x_{}`$, $`x_+`$, and $`x_3`$, divided by the relations
$$\begin{array}{c}x_{}x_0=x_0x_{},x_+x_0=x_0x_+,x_3x_0=x_0x_3,\\ q^1x_{}x_3qx_3x_{}=\lambda x_{}x_0,q^1x_3x_+qx_+x_3=\lambda x_+x_0,\\ x_{}x_+x_+x_{}\lambda x_3x_3=\lambda x_3x_0,\end{array}$$
(5)
is called the algebra of quantum Minkowski space, denoted by $`𝒳`$.
Here, $`x_0`$ is the time coordinate, while $`x_{}`$, $`x_+`$, and $`x_3`$ are the space coordinates with $`\mathrm{su}_2`$-weight indices. The commutation relations (5) can be written in more sophisticated forms, e.g., in terms of $`q`$-Clebsch-Gordon coefficients or in terms of an $`R`$-matrix as we will recall in the next section.
The center of $`𝒳`$ is generated by $`x_0`$ and the four-square
$$x^2:=(x_0)^2+q^1x_{}x_++qx_+x_{}(x_3)^2.$$
(6)
Since we will not use 2 as index, $`x^2`$ cannot be confused with any one of the coordinates. The four-square can also be written as
$$x^2=x_\mu x_\nu \eta ^{\mu \nu }$$
(7)
in terms of the $`q`$-Minkowski metric defined by
$$\eta ^{00}=1\eta ^+=q^1,\eta ^+=q\eta ^{33}=1,$$
(8)
all other components vanishing. The metric defines upper four-vector indices by
$$x^\mu :=\eta ^{\mu \nu }x_\nu .$$
(9)
Note, that $`x^2=x_\mu x^\mu `$ but $`x^2x^\mu x_\mu `$ because $`\eta ^{\mu \nu }`$ is not symmetric.
### 2.2 Quantum Lorentz algebra
By construction $`𝒳`$ is a module algebra with respect to the quantum Lorentz algebra $`𝒰_q(\mathrm{sl}_2())`$. There are several essentially isomorphic forms of $`𝒰_q(\mathrm{sl}_2())`$ in use . In the chiral decomposition form $`𝒰_q(\mathrm{sl}_2())`$ is given, as algebra, by the tensor square of quantum $`𝒰(\mathrm{su}_2)`$,
$$𝒰_q(\mathrm{sl}_2())=𝒰_q(\mathrm{su}_2)𝒰_q(\mathrm{su}_2).$$
(10)
The definition of $`𝒰_q(\mathrm{su}_2)`$ is recalled in Appendix A. The chiral decomposition implies that the irreducible representations of $`𝒰_q(\mathrm{sl}_2())`$ can be labelled by the highest weights of the two tensor factors. For example, the structure morphism of the four-vector representation is given on this chiral form by
$$\rho ^{(\frac{1}{2},\frac{1}{2})}:=\rho ^{\frac{1}{2}}\rho ^{\frac{1}{2}},$$
(11)
where $`\rho ^{\frac{1}{2}}`$ is the fundamental representation of $`𝒰_q(\mathrm{su}_2)`$ with highest weight or spin $`j=\frac{1}{2}`$.
The quantum Lorentz algebra is quasi-triangular. In the chiral form (10) the universal $``$-matrix of $`𝒰_q(\mathrm{sl}_2())`$ can be expressed in terms of the well known universal $``$-matrix of $`𝒰_q(\mathrm{su}_2)`$. There are two inequivalent universal $``$-matrices ,
$$_\mathrm{I}=_{41}^1_{31}^1_{24}_{23},_{\mathrm{II}}=_{41}^1_{13}_{24}_{23},$$
(12)
where $``$ is the universal $``$-matrix of $`𝒰_q(\mathrm{su}_2)`$ given in Eq. (124), and where we have used the tensor leg notation $`_{23}=11`$ etc. The four-vector representations of the universal $``$-matrices are denoted by
$`R_\mathrm{I}^{\mu \nu }_{\sigma \tau }`$ $`:=(\rho ^{(\frac{1}{2},\frac{1}{2})}\rho ^{(\frac{1}{2},\frac{1}{2})})(_\mathrm{I})^{\mu \nu }{}_{\sigma \tau }{}^{},`$ (13)
$`R_{\mathrm{II}}^{\mu \nu }_{\sigma \tau }`$ $`:=(\rho ^{(\frac{1}{2},\frac{1}{2})}\rho ^{(\frac{1}{2},\frac{1}{2})})(_{\mathrm{II}})^{\mu \nu }{}_{\sigma \tau }{}^{}.`$
With the $`R`$-matrices the commutation relations (5) can be written in the compact form
$$x_\sigma x_\tau =x_\mu x_\nu R_\mathrm{I}^{\mu \nu }{}_{\tau \sigma }{}^{}.$$
(14)
This implies that the $`𝒳`$-algebra has the Poincaré-Birkhoff-Witt property, because the $`R`$-matrices satisfy the Yang-Baxter equation. From the decomposition of the $`R`$-matrices into eigenspaces we can, furthermore, derive the relation
$$(1+q\widehat{R}_{\mathrm{II}})(1\widehat{R}_\mathrm{I})=0,$$
(15)
where the hat denotes swapping the two upper indices $`\widehat{R}^{\mu \nu }{}_{\tau \sigma }{}^{}=R^{\nu \mu }_{\tau \sigma }`$. We will need this formula later.
The disadvantage of the chiral form (10) is that the inner tensor factors of the coproduct are twisted with the $``$-matrix of $`𝒰_q(\mathrm{su}_2)`$. Since the $``$-matrix is given by an infinite series, which exists only as formal power series, some of the generators of $`𝒰_q(\mathrm{sl}_2())`$ have a complicated coproduct which then exists only as formal power series, too. Therefore, it is also convenient to use the Drinfeld double form , where $`𝒰_q(\mathrm{sl}_2())`$ is represented as Drinfeld double of the Hopf-dually paired $`𝒰_q(\mathrm{su}_2)`$ and $`SU_q(2)^{\mathrm{op}}`$.
### 2.3 Noncommutative differential calculus
The generators of $`q`$-momenta are required to transform as four-vector with respect to $`q`$-Lorentz transformations. The fact that $`q`$-momenta transform in the same way as the coordinate generators of $`q`$-Minkowski space implies that they have to satisfy the same commutation relations as well. This leads to the definition of the $`q`$-momentum algebra as the algebra generated by $`p_0`$, $`p_{}`$, $`p_+`$, $`p_3`$ with commutation relations (5), replacing $`x`$ with $`p`$ everywhere.
Noncommutative partial derivatives
$$_\mu :=\mathrm{i}p_\mu $$
(16)
are defined by the action of the momentum generators on the $`q`$-Minkowski algebra of space functions. Partial derivatives of the coordinates have to be dimensionless numbers from the base field of complex numbers (including the deformation parameter $`q`$). Just as in the undeformed case the only tensor in the base field with two four-vector indices is the metric. We conclude that the action of partial derivatives on coordinates is given by
$$_\mu x_\nu =\eta _{\mu \nu }^\mu x_\nu =\delta _\nu ^\mu .$$
(17)
This action on the generators extends to the entire algebra of $`q`$-Minkowski space by the coproduct of the momenta ,
$$\begin{array}{cc}\hfill \mathrm{\Delta }(p^\mu )& :=p^\mu 1+(\kappa 1)_{\mathrm{II}}{}_{}{}^{1}(1p^\mu )_{\mathrm{II}}\hfill \\ & =p^\mu 1+\kappa L^\mu {}_{\nu }{}^{}p^\nu .\hfill \end{array}$$
(18)
Here $`_{\mathrm{II}}`$ is the second of the universal $``$-matrices of the $`q`$-Lorentz algebra given in Eq. (12). The $`L`$-matrix is the four-vector half-representation of this $``$-matrix
$$L^\mu {}_{\nu }{}^{}:=_{\mathrm{II}}{}_{[1]}{}^{}\rho _{}^{(\frac{1}{2},\frac{1}{2})}(_{\mathrm{II}}{}_{[2]}{}^{})^\mu {}_{\nu }{}^{},$$
(19)
where we have used a Sweedler-like notation for $`_{\mathrm{II}}=_{\mathrm{II}}{}_{[1]}{}^{}_{\mathrm{II}}_{[2]}`$. Finally, $`\kappa `$ is a group-like scaling operator,
$$p_\mu \kappa =q\kappa p_\mu ,x_\mu \kappa =q^1\kappa x_\mu ,\mathrm{\Delta }(\kappa )=\kappa \kappa .$$
(20)
The action (17) of partial $`q`$-derivatives together with the coproduct (18) defines a $`q`$-Lorentz-covariant differential calculus on $`q`$-Minkowski space. Letting $`p^\mu `$ act on $`x_\nu f`$, $`f𝒳`$, we obtain
$$p^\mu (x_\nu f)=\mathrm{i}\delta _\nu ^\mu f+qR_{\mathrm{II}}^{\nu ^{}\mu }{}_{\nu \mu ^{}}{}^{}x_{\nu ^{}}^{}(p^\mu ^{}f),$$
(21)
from which we can deduce the $`q`$-deformed Heisenberg commutation relations
$$p^\mu x_\nu qR_{\mathrm{II}}^{\nu ^{}\mu }{}_{\nu \mu ^{}}{}^{}x_{\nu ^{}}^{}p^\mu ^{}=\mathrm{i}\delta _\nu ^\mu ,$$
(22)
which are of the type originally introduced by Wess and Zumino .
In order to show that Eqs. (17) and (18) define a differential calculus which is well defined on the algebra $`𝒳`$, that is, which does not depend on the ordering of the coordinate generators, we have to check if the action of the momenta is consistent with commutation relations (5) defining $`𝒳`$. This is best done using the $`R`$-matrix form (14),
$$\begin{array}{cc}\hfill p^\lambda (x_\sigma x_\tau x_\mu x_\nu R_\mathrm{I}^{\mu \nu }{}_{\tau \sigma }{}^{})& =(p^\lambda x_\nu x_\mu )(\delta _\tau ^\nu \delta _\sigma ^\mu \widehat{R}_\mathrm{I}^{\nu \mu }{}_{\tau \sigma }{}^{})\hfill \\ & =x_\rho (\delta _\nu ^\lambda \delta _\mu ^\rho +q\widehat{R}_{\mathrm{II}}^{\lambda \rho }{}_{\nu \mu }{}^{})(\delta _\tau ^\nu \delta _\sigma ^\mu \widehat{R}_\mathrm{I}^{\nu \mu }{}_{\tau \sigma }{}^{})\hfill \\ & =0,\hfill \end{array}$$
(23)
where in the last step we have used Eq. (15).
In principle, the action (17) of the partial derivatives on generators together with the coproduct (18) of momenta defines the action of partial derivatives on arbitrary elements of quantum Minkowski space. Note, however, that this definition is recursive. It does not yield formulas for the partial derivatives of a basis of the noncommutative space algebra, such as the Poincaré-Birkhoff-Witt basis of ordered monomials,
$$_{\mathrm{PBW}}:=\{(x_0)^{n_0}(x_{})^n_{}(x_+)^{n_+}(x_3)^{n_3}|n_0,n_{},n_+,n_3_0\}.$$
(24)
To derive such a formula is one of the goals of the next section.
## 3 Computation in noncommutative calculus
### 3.1 L-matrix calculus
Let us first consider the simple case of a power of one single generator, $`f=x_\alpha ^n`$, $`\alpha \{0,,+,3\}`$, $`n`$. Using the coproduct (18) we get by induction
$$^\mu x_\alpha ^n=\underset{k=0}{\overset{n1}{}}q^k(L^\mu {}_{\nu }{}^{}x_\alpha ^k)(^\nu x_\alpha )x_\alpha ^{nk1},$$
(25)
where the factor $`q^k`$ comes from the action of the scaling operator (20),
$$\kappa x_\alpha ^k=\kappa _{(1)}x_\alpha ^kS(\kappa _{(2)})=\kappa x_\alpha ^k\kappa ^1=q^kx_\alpha ^k.$$
(26)
By construction, the $`L`$-matrix is a comultiplicative quantum matrix, $`\mathrm{\Delta }(L^\mu {}_{\nu }{}^{})=L^\mu {}_{\sigma }{}^{}L^\sigma _\nu `$. So we get for the action of $`L`$ on a power
$$L^\mu {}_{\nu }{}^{}x_\alpha ^k=(L^\mu {}_{\sigma _1}{}^{}x_\alpha )(L^{\sigma _1}{}_{\sigma _2}{}^{}x_\alpha )\mathrm{}(L^{\sigma _{k1}}{}_{\nu }{}^{}x_\alpha ).$$
(27)
Since each $`L^{\sigma _k}{}_{\sigma _{k+1}}{}^{}x_\alpha `$ is an $`𝒳`$-valued 4$`\times `$4-matrices, the right hand side is the matrix product of such matrices. This suggests to introduce an index-free notation:
###### Definition 2.
Let $`L^\mu _\nu `$ be the $`L`$-matrix defined in Eq. (19), $`^\mu `$ the partial derivatives defined in (16), $`f𝒳`$ an element of the algebra of quantum Minkowski space. Then we denote by $`L_f`$ the $`𝒳`$-valued 4$`\times `$4-matrix and by $`f`$ the $`𝒳`$-valued four-vector with entries
$$(L_f)^\mu {}_{\nu }{}^{}:=L^\mu {}_{\nu }{}^{}f\text{and}(f)^\mu :=^\mu f,$$
(28)
respectively.
To illustrate this, we have for example
$$(L_f^2f)^\mu =(L^\mu {}_{\nu }{}^{}f)(L^\nu {}_{\sigma }{}^{}f)(^\sigma f).$$
(29)
In this index-free notation Eq. (27) is written as $`Lx_\alpha ^k=L_{x_\alpha }^k`$, such that Eq. (25) becomes
$$x_\alpha ^n=\underset{k=0}{\overset{n1}{}}q^kL_{x_\alpha }^k(x_\alpha )x_\alpha ^{nk1}.$$
(30)
We recall that by definition (17) of the action of partial derivatives on the generators we have for the components of the gradient
$$(x_\alpha )^\mu =^\mu x_\alpha =\delta _\alpha ^\mu ,$$
(31)
where $`\mu ,\alpha \{0,,+,3\}`$.
The computationally difficult part of (30) is the evaluation of powers of the algebra-valued matrices $`L_{x_\alpha }`$. For example, for $`(L_{x_0})^\mu _\nu `$ we get with respect to four-vector indices $`\mu ,\nu \{0,,+,3\}`$
$$L_{x_0}=\frac{1}{[2]}\left(\begin{array}{cccc}\frac{[4]}{[2]}x_0& q\lambda x_{}& q\lambda x_+& q\lambda x_3\\ \frac{\lambda }{q^2}x_+& \frac{[4]}{[2]}x_0\frac{\lambda ^2}{q}x_3& 0& \lambda x_+\\ \lambda x_{}& 0& \frac{[4]}{[2]}x_0q\lambda ^2x_3& \lambda x_{}\\ \frac{\lambda }{q}x_3& q\lambda x_{}& \frac{\lambda }{q}x_+& \frac{[4]}{[2]}x_0\lambda ^2x_3\end{array}\right).$$
(32)
Little is known about the calculation of powers of general algebra valued matrices. Since the entries of this matrix do not commute we cannot resort to the usual methods of linear algebra. And in a brute force approach the number of terms in $`L_{x_0}^n`$ which have to be reordered to the Poincaré-Birkhoff-Witt form (or any other ordering) increases exponentially with $`n`$.
### 3.2 Block decomposition of L-matrices
The computational problem of computing the powers of the $`𝒳`$-valued $`L_{x_\alpha }`$-matrices can be simplified by considering the $`L`$-matrix (19) not with respect to four-vector indices running through $`\{0,,+,3\}`$ as in Eq. (32), but with respect to pairs of indices of the two spin-$`\frac{1}{2}`$ representations of the chiral decomposition $`𝒰_q(\mathrm{sl}_2())=𝒰_q(\mathrm{su}_2)𝒰_q(\mathrm{su}_2)`$ of the $`q`$-Lorentz algebra, each labelled by the weights $`\{\frac{1}{2},+\frac{1}{2}\}=\{,+\}`$. We recall that the generators of quantum Minkowski space in this $`(\frac{1}{2},\frac{1}{2})`$-spinor basis labelled by indices $`\{,+,+,++\}`$ are related to those in the four-vector basis by
$`x_{}`$ $`=[2]^{\frac{1}{2}}x_{},`$ (33)
$`x_+`$ $`=q^{\frac{1}{2}}(x_3x_0),`$
$`x_+`$ $`=q^{\frac{1}{2}}x_3+q^{\frac{3}{2}}x_0,`$
$`x_{++}`$ $`=[2]^{\frac{1}{2}}x_+.`$
Using formula (12) for the universal $`_{\mathrm{II}}`$-matrix in terms of the $``$-matrix of $`𝒰_q(\mathrm{su}_2)`$ and definition (19) of the $`L`$-matrix, we calculate the $`L`$-matrix with respect to the spinor basis,
$$\begin{array}{cc}\hfill (L)^{ij}_{kl}& =(\mathrm{id}\mathrm{id}\rho ^{\frac{1}{2}}\rho ^{\frac{1}{2}})(_{\mathrm{II}})^{ij}_{kl}\hfill \\ & =\left(L_{}^{\frac{1}{2}}\right)^j{}_{j^{}}{}^{}\left(L_+^{\frac{1}{2}}\right)_{}^{i}{}_{i^{}}{}^{}\left(L_+^{\frac{1}{2}}\right)^j^{}{}_{l}{}^{}\left(L_+^{\frac{1}{2}}\right)_{}^{i^{}}_k\hfill \\ & =B^j{}_{l}{}^{}(L_+^{\frac{1}{2}})_{}^{i}{}_{k}{}^{},\hfill \end{array}$$
(34)
where all indices run through $`\{\frac{1}{2},+\frac{1}{2}\}=\{,+\}`$. Here $`B=(B^j{}_{l}{}^{})`$ is the matrix of “boosts” which generate the $`SU_q(2)^{\mathrm{op}}`$ Hopf subalgebra of the Drinfeld double form of the $`q`$-Lorentz algebra , and $`L_+^{1/2}`$ the $`L_+`$-matrix of the $`𝒰_q(\mathrm{su}_2)`$ subalgebra of rotations,
$$B=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right),L_+^{\frac{1}{2}}=\left(\begin{array}{cc}K^{\frac{1}{2}}& q^{\frac{1}{2}}\lambda K^{\frac{1}{2}}E\\ 0& K^{\frac{1}{2}}\end{array}\right),$$
(35)
with respect to the $`\{,+\}`$ basis. The explicit form of the boosts in the chiral decomposition of the $`q`$-Lorentz algebra is given in Eq. (125).
The decomposition (34) of the $`L`$-matrix into boosts and rotations yields a natural block matrix decomposition of the $`𝒳`$-valued 4$`\times `$4-matrices $`L_{x_\alpha }`$ as defined in Eq. (28). For the four-vector index $`\alpha `$ it turns out to be convenient to use instead of $`x_3`$ the light cone coordinate
$$x_{30}:=x_3x_0,$$
(36)
which is essentially the $`x_+`$-coordinate in the spinor basis (33) of quantum Minkowski space. And for the block decomposition it is convenient to define $`𝒳`$-valued 2$`\times `$2-matrices by the action of the boosts
$$(B_f)^i{}_{j}{}^{}:=B^i{}_{j}{}^{}f,f𝒳,$$
(37)
using the same index-free notation as introduced in Eq. (28). Since $`B^i_j`$ is a comultiplicative quantum matrix, $`\mathrm{\Delta }(B^i{}_{j}{}^{})=B^i{}_{k}{}^{}B^k_j`$, the matrices $`B_{x_\alpha }`$ satisfy the same commutation relations (5) as $`x_\alpha `$. Explicitly, the $`B_{x_\alpha }`$-matrices can be computed to
$`B_{x_0}`$ $`=\left(\begin{array}{cc}\frac{[4]}{[2]^2}x_0+\frac{\lambda }{q[2]}x_3& q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}x_+\\ q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}x_{}& \frac{[4]}{[2]^2}x_0\frac{q\lambda }{[2]}x_3\end{array}\right)`$ (38)
$`B_x_{}`$ $`=\left(\begin{array}{cc}x_{}& q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}x_{30}\\ 0& x_{}\end{array}\right)`$
$`B_{x_+}`$ $`=\left(\begin{array}{cc}x_+& 0\\ q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}x_{30}& x_+\end{array}\right)`$
$`B_{x_{30}}`$ $`=\left(\begin{array}{cc}q^1x_{30}& 0\\ 0& qx_{30}\end{array}\right),`$
with respect to spin-$`\frac{1}{2}`$ indices running through $`\{,+\}`$. In terms of these matrices the $`L_{x_\alpha }`$-matrices can be expressed as
$`L_{x_0}`$ $`=\left(\begin{array}{cc}B_{x_0}& 0\\ 0& B_{x_0}\end{array}\right)`$ (39)
$`L_x_{}`$ $`=\left(\begin{array}{cc}qB_x_{}& q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}B_{x_3}\\ 0& q^1B_x_{}\end{array}\right)`$
$`L_{x_+}`$ $`=\left(\begin{array}{cc}q^1B_{x_+}& 0\\ 0& qB_{x_+}\end{array}\right)`$
$`L_{x_{30}}`$ $`=\left(\begin{array}{cc}B_{x_{30}}& q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}B_{x_+}\\ 0& B_{x_{30}}\end{array}\right),`$
with respect to the index structure $`\{,+,+,++\}`$.
### 3.3 Calculating powers of L-matrices
In the form of Eq. (39) the $`L_{x_\alpha }`$ matrices are upper block triangular. This makes the computation of their powers a lot easier. $`L_{x_0}`$ and $`L_{x_+}`$ are even block diagonal, so we immediately get
$$L_{x_0}^n=\left(\begin{array}{cc}B_{x_0}^n& 0\\ 0& B_{x_0}^n\end{array}\right),L_{x_+}^n=\left(\begin{array}{cc}q^nB_{x_+}^n& 0\\ 0& q^nB_{x_+}^n\end{array}\right).$$
(40)
For the calculation of $`L_x_{}^n`$ and $`L_{x_{30}}^n`$ we have to use the commutation relations (5) for the $`B_{x_\alpha }`$ matrices. We thus get by induction
$`L_x_{}^n`$ $`=\left(\begin{array}{cc}q^nB_x_{}^n& q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}\left(q^{n1}\frac{[2n]}{[2]}B_{x_{30}}+[n]B_{x_0}\right)B_x_{}^{n1}\\ 0& q^nB_x_{}^n\end{array}\right)`$
$`L_{30}^n`$ $`=\left(\begin{array}{cc}B_{x_{30}}^n& q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}q^{n1}[n]B_{x_+}B_{x_{30}}^{n1}\\ 0& B_{x_{30}}^n\end{array}\right).`$ (41)
Now we have expressed the powers of $`L_{x_\alpha }`$ matrices in terms of powers of the $`B_{x_\alpha }`$ matrices. It remains to calculate the powers of the $`B_{x_\alpha }`$-matrices. Since $`B_{x_{30}}`$ is diagonal we immediately get
$$B_{x_{30}}^n=\left(\begin{array}{cc}q^nx_{30}^n& 0\\ 0& q^nx_{30}^n\end{array}\right).$$
(42)
For the powers of the matrices $`B_{x_\pm }`$, which are block triangular, we use again the commutation relations (5) to obtain by induction
$`B_x_{}^n`$ $`=\left(\begin{array}{cc}x_{}^n& q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}q^{n1}[n]x_{30}x_{}^{n1}\\ 0& x_{}^n\end{array}\right)`$ (43)
$`B_{x_+}^n`$ $`=\left(\begin{array}{cc}x_+^n& 0\\ q^{\frac{1}{2}}\lambda [2]^{\frac{1}{2}}q^{1n}[n]x_{30}x_+^{n1}& x_+^n\end{array}\right).`$
The calculation of the powers of $`B_{x_0}`$ is more difficult because it is not triangular. Using $`q`$-Pauli matrices (123), $`B_{x_0}`$ of Eqs. (38) can be written more compactly as
$$B_{x_0}:=\frac{[4]}{[2]^2}x_0\frac{\lambda }{[2]}\stackrel{~}{\sigma }_Ax^A,$$
(44)
where the index $`A`$ is summed over $`\{,3,+\}`$. The $`q`$-Pauli matrices satisfy the relation
$$\stackrel{~}{\sigma }_A\stackrel{~}{\sigma }_B=g_{BA}\stackrel{~}{\sigma }_C\epsilon _B{}_{}{}^{C}{}_{A}{}^{},$$
(45)
where the quantum metric and epsilon tensor are defined by quantum Clebsch-Gordon coefficients as
$`g^{AB}`$ $`:=\sqrt{[3]}C_q(1,1,0;A,B,0)`$ (46)
$`\epsilon ^{AB}_C`$ $`:=\sqrt{{\displaystyle \frac{[4]}{[2]}}}C_q(1,1,1;A,B,C).`$
Contracting Eq. (45) with $`x^Ax^B`$ and using the commutation relations (5) of the coordinates, which can be written as $`x_Ax_B\epsilon ^{AB}{}_{C}{}^{}=\lambda x_Cx_0`$, we derive
$$\begin{array}{cc}\hfill \stackrel{~}{\sigma }_A\stackrel{~}{\sigma }_Bx^Ax^B& =(g_{BA}\stackrel{~}{\sigma }_C\epsilon _B{}_{}{}^{C}{}_{A}{}^{})x^Ax^B\hfill \\ & =x_Ax^A+\lambda x_0\stackrel{~}{\sigma }_Ax^A.\hfill \end{array}$$
(47)
This can be used to compute the square of Eq. (44),
$$\begin{array}{cc}\hfill B_{x_0}^2& =\frac{[4]^2}{[2]^2}(x_0)^22\frac{\lambda [4]}{[2]^3}x_0\stackrel{~}{\sigma }_Ax^A+\frac{\lambda ^2}{[2]^2}\stackrel{~}{\sigma }_A\stackrel{~}{\sigma }_Bx^Ax^B\hfill \\ & =\frac{[4]^2}{[2]^4}(x_0)^22\frac{\lambda [4]}{[2]^3}x_0\stackrel{~}{\sigma }_Ax^A+\frac{\lambda ^2}{[2]^2}(x_Ax^A+\lambda x_0\stackrel{~}{\sigma }_Ax^A)\hfill \\ & =[2]x_0B_{x_0}(x_0)^2\frac{\lambda ^2}{[2]^2}x^2.\hfill \end{array}$$
(48)
Together with Eq. (40) we conclude that $`L_{x_0}`$ satisfies the polynomial equation
$$\begin{array}{cc}\hfill 0& =L_{x_0}^2[2]x_0L_{x_0}+\left((x_0)^2+\frac{\lambda ^2}{[2]^2}x^2\right)\hfill \\ & =L_{x_0}^22bL_{x_0}+c,\hfill \end{array}$$
(49)
where we have introduced the abbreviations
$$b:=\frac{[2]}{2}x_0,c:=(x_0)^2+\frac{\lambda ^2}{[2]^2}x^2,$$
(50)
which are both in the center of $`𝒳`$. Eq. (49) enables us to derive formulas for arbitrary powers of $`L_{x_0}`$. This is most elegantly done by considering the generating function of the powers $`(1L_{x_0}z)^1=1+L_{x_0}z+L_{x_0}^2z^2+\mathrm{}`$, where $`z`$ is a formal parameter. First we rewrite (49) as
$$(1L_{x_0}^2z^2)2bz(1L_{x_0}z)(12bz+cz^2)=0$$
(51)
and divide it by $`(1L_{x_0}z)`$ and $`(12bz+cz^2)`$, yielding
$$\frac{1}{1L_{x_0}z}=\frac{1+(L_{x_0}2b)z}{12bz+cz^2}.$$
(52)
The right hand side can be expanded in powers of $`z`$, observing that it contains the generating function for Chebyshev polynomials of the second kind,
$$\frac{1}{12yt+t^2}=\underset{n=0}{\overset{\mathrm{}}{}}U_n(y)t^n,$$
(53)
where we have to set $`y=b/\sqrt{c}`$ and $`t=\sqrt{c}z`$. We thus obtain
$$\begin{array}{cc}\hfill L_{x_0}^n& =U_n\left(b/\sqrt{c}\right)c^{\frac{n}{2}}+U_{n1}\left(b/\sqrt{c}\right)c^{\frac{n1}{2}}(L_{x_0}2b)\hfill \\ & =U_{n2}\left(b/\sqrt{c}\right)c^{\frac{n}{2}}+U_{n1}\left(b/\sqrt{c}\right)c^{\frac{n1}{2}}L_{x_0}\hfill \end{array}$$
(54)
for $`n2`$. The fact that $`U_n(y)`$ is a polynomial of degree $`n`$ which contains either only even or only odd powers of $`y`$ implies that $`U_n(1/y)y^n`$ is a quadratic polynomial in $`y`$. Therefore, the right hand side of Eq. (54) contains only even positive powers of $`\sqrt{c}`$, that is, does not contain inverses or square roots of $`c`$ and, hence, yields well-defined elements of the algebra $`𝒳`$. Together with Eqs. (40), (3.3), (42), and (43) this completes our calculation of powers of the $`L_{x_\alpha }`$ matrices.
### 3.4 Calculating derivatives
#### 3.4.1 Choice of basis and partial Jackson derivatives
With the formulas for the powers of the $`L_\alpha `$ matrices the derivatives of arbitrary elements $`f𝒳`$ could be computed by calculating the derivatives of the Poincaré-Birkhoff-Witt basis of $`𝒳`$. However, it turns out to be more convenient to work in a slightly different basis. Observe that from definition (7) of the $`q`$-Lorentz invariant length and the commutation relations (5) we can deduce the relation
$$[2]x_{}x_+=x^2+q^2(x_{30})^2+q\lambda x_0x_{30}.$$
(55)
Using this relation, all products of equal powers of $`x_{}`$ and $`x_+`$ in the Poincaré-Birkhoff-Witt basis can be substituted by powers of $`x^2`$. We thus obtain a basis of ordered monomials which contains powers of $`x^2`$ and either powers of $`x_{}`$ or $`x_+`$. To be more precise, we have a basis of $`𝒳`$,
$$=_{}_+,$$
(56)
where
$`_{}`$ $`:=\{(x_\mu x^\mu )^i(x_0)^j(x_{30})^k(x_{})^l|i,j,k,l_0\}`$ (57)
$`_+`$ $`:=\{(x_\mu x^\mu )^i(x_+)^l(x_0)^j(x_{30})^k|i,j,k,l_0\}.`$
Note, that $`_{}`$ and $`_+`$ have a non-empty intersection, the monomials for which $`l=0`$. But this will not be a problem here. Let us introduce a more suggestive notion for sums of monomials in a particular order:
###### Definition 3.
Let $`x_1,\mathrm{},x_n𝒳`$ be elements of the algebra of quantum Minkowski space. Then we will denote by $`f(x_1,\mathrm{},x_n)`$ linear combinations of ordered products of monomials in $`x_1,\mathrm{},x_n𝒳`$. That is,
$$f(x_1,\mathrm{},x_n)\mathrm{Span}\{(x_1)^{k_1}\mathrm{}(x_n)^{k_n}|k_1,\mathrm{},k_n_0\}.$$
(58)
With this notation the statement that $`_{}_+`$ is a basis of quantum Minkowski space can be written as
$$f(x_0,x_+,x_{30},x_{})=f_{}(x^2,x_0,x_{30},x_{})+f_+(x^2,x_0,x_+,x_{30}),$$
(59)
decomposing an arbitrary element of $`𝒳`$ into two parts containing powers in either $`x_{}`$ or $`x_+`$.
It is important to have a notation which keeps track of the ordering of generators because the formulas for derivatives are most elegantly expressed in terms of partial Jackson derivatives which depend on the ordering. Recall, that the Jackson derivative or $`q`$-derivative of a function $`f=f(x_\alpha )`$ in a single variable $`x_\alpha `$ is defined as difference quotient,
$$\frac{_{q^2}f}{_{q^2}x_\alpha }:=\frac{f(q^2x_\alpha )f(x_\alpha )}{q^2x_\alpha x_\alpha }.$$
(60)
For monomials $`f(x_\alpha )=(x_\alpha )^k`$, this yields
$$\frac{_{q^2}(x_\alpha )^k}{_{q^2}x_\alpha }=(k)_{q^2}(x_\alpha )^{k1},$$
(61)
where $`(k)_{q^2}`$ is the usual $`q`$-number defined in Eq. (4). This formula naturally generalizes to a partial Jackson derivative on monomials of several noncommutative variables.
###### Definition 4.
Let $`x_1,\mathrm{},x_n𝒳`$ be elements of the algebra of quantum Minkowski space (or any other algebra). We define the partial Jackson derivative with respect to one of these elements $`x_\alpha `$ on ordered monomials by
$$\frac{_{q^2}}{_{q^2}x_\alpha }\left((x_1)^{k_1}\mathrm{}(x_n)^{k_n}\right):=(k_\alpha )_{q^2}(x_1)^{k_1}\mathrm{}(x_\alpha )^{k_\alpha 1}\mathrm{}(x_n)^{k_n},$$
(62)
and extend it linearly to arbitrary linear combinations of such monomials. This defines partial Jackson derivatives on general functions $`f=f(x_1,\mathrm{},x_n)`$ in the sense of Definition 3.
It must be emphasized that the partial Jackson derivatives depend on the particular order and are not well-defined on the algebra. For example, within the algebra we have the commutation relation $`x_{30}x_+q^2x_+x_{30}=0`$, but
$$\frac{_{q^2}}{_{q^2}x_+}(x_{30}x_+q^2x_+x_{30})=x_{30}q^2x_{30}0.$$
(63)
This is why we have defined the notation $`f=f(x_+,x_{30})`$ to not just denote elements of the algebra but to denote in addition a specific ordering of the variables. With this notation the action of partial Jackson derivatives is defined unambiguously.
We break up the calculation of derivatives in two steps. First, we calculate the derivatives of non-central variables, $`x_{}`$, $`x_{30}`$, $`x_+`$. Working within the basis $``$ we can limit the considerations in this first step to polynomial functions $`f=f(x_{30},x_{})`$ and $`f=f(x_+,x_{30})`$ which depend either on $`x_{}`$ or on $`x_+`$. This calculation is rather straightforward. In a second step we consider polynomial functions $`f=f(x^2,x_0)`$ of the center of $`𝒳`$. This calculation is much more involved.
#### 3.4.2 Functions of non-central coordinates
We begin the calculation of the derivatives of the basis $``$ by reading off Eqs. (38) and (39) that the $`L_\alpha `$ matrices possess some algebra valued eigenvalues,
$`L_Ax_A`$ $`=qx_Ax_A`$ (64)
$`L_{30}x_{}`$ $`=q^1x_{30}x_{}`$
$`L_+x_{30}`$ $`=q^1x_+x_{30},`$
for $`A\{,+,30\}`$ (no summation over $`A`$). We recall that $`x_\alpha `$ is the index free notation for $`(x_\alpha )^\mu =^\mu x_\alpha =\delta _\alpha ^\mu `$.
Inserting the first of these eigenvalue equations into Eq. (30) we obtain for powers of the generators
$$x_A^n=\underset{k=0}{\overset{n}{}}q^{2k}x_A^k(x_\alpha )x_\alpha ^{nk1}=(n)_{q^2}x_A^{n1}x_\alpha ,$$
(65)
for $`A\{,+,30\}`$, where $`(n)_{q^2}`$ denotes the quantum number (4). From Eq. (65) we can deduce that the derivative of any polynomial $`f=f(x_A)`$ in a single one of the generators $`x_A`$ can be expressed in terms of the Jackson derivative by
$$f(x_A)=\frac{_{q^2}f}{_{q^2}x_A}x_A^\mu f(x_A)=\frac{_{q^2}f}{_{q^2}x_A}\delta _A^\mu .$$
(66)
Using the second of Eqs. (64) we get
$$\begin{array}{cc}\hfill (x_{30}^kx_{}^n)& =(x_{30}^k)x_{}^n+q^kL_{30}^k(x_{}^n)\hfill \\ & =(k)_{q^2}x_{30}^{k1}x_{}^nx_{30}+x_{30}^k(n)_{q^2}x_{}^{n1}x_{},\hfill \end{array}$$
(67)
and analogously for the derivative of $`x_+^kx_{30}^n`$,
$$(x_+^kx_{30}^n)=(k)_{q^2}x_+^{k1}x_{30}^nx_++x_+^k(n)_{q^2}x_{30}^{n1}x_{30}.$$
(68)
From Eqs. (67) and (68) we can deduce the following result:
###### Proposition 1.
Let $`f=f(x_{30},x_{})`$ and $`f=f(x_+,x_{30})`$ be ordered polynomials in the notation of Definition 3. Then their partial derivatives are expressed in terms of partial Jackson derivatives as
$`f(x_{30},x_{})`$ $`={\displaystyle \frac{_{q^2}f}{_{q^2}x_{30}}}x_{30}+{\displaystyle \frac{_{q^2}f}{_{q^2}x_{}}}x_{}`$ (69a)
$`f(x_+,x_{30})`$ $`={\displaystyle \frac{_{q^2}f}{_{q^2}x_+}}x_++{\displaystyle \frac{_{q^2}f}{_{q^2}x_{30}}}x_{30}.`$ (69b)
#### 3.4.3 Functions of four-length
Next, we calculate the derivative of powers of the coordinate four-square $`x^2`$. Since $`x^2`$ is a Lorentz scalar the $`L`$-matrix acts with the antipode $`\epsilon `$,
$$L^\mu {}_{\nu }{}^{}x^2=\epsilon (L^\mu {}_{\nu }{}^{})x^2=\delta _\nu ^\mu x^2.$$
(70)
By the same reasoning which led to Eq. (30) we obtain for powers of the four-square
$$\begin{array}{cc}\hfill (x^2)^n& =\underset{k=0}{\overset{n1}{}}q^{2k}(x^2)^{n1}x^2\hfill \\ & =(n)_{q^2}(x^2)^{n1}x^2.\hfill \end{array}$$
(71)
Again, for a general function $`f=f(x^2)`$ this can be written in terms of a Jackson derivative by
$$f(x^2)=\frac{_{q^2}f}{_{q^2}(x^2)}x^2.$$
(72)
It remains to compute the derivative of $`x^2`$, which for representation theoretic reasons we expect to be proportional to the coordinate four vector $`x^\mu `$. After lengthy calculations we, indeed, find
$$(x^2)^\mu =q^1[2]x^\mu .$$
(73)
#### 3.4.4 Functions of time
Like the calculation of powers of $`L_{x_0}`$, the calculation of derivatives of powers of the time coordinate is more difficult. In order to calculate the derivative of $`(x_0)^n`$ we first observe that, since $`x_0`$ is central in the space algebra we can write Eq. (30) in the form
$$x_0^n=\frac{(qL_{x_0})^n(x_0)^n}{qL_{x_0}x_0}x_0,$$
(74)
which is similar to a $`q`$-difference quotient. We can get rid of the matrix in the denominator,
$$\begin{array}{cc}\hfill \frac{(qL_{x_0})^n(x_0)^n}{qL_{x_0}x_0}& =\frac{(q^nL_{x_0}^nx_0^n)(q^1L_{x_0}x_0)}{(qL_{x_0}x_0)(q^1L_{x_0}x_0)}\hfill \\ & =\frac{[2](q^{n1}L_{x_0}^{n+1}q^nL_{x_0}^nx_0q^1L_{x_0}x_0^nx_0^{n+1})}{\lambda ^2x^2},\hfill \end{array}$$
(75)
where in the second step we have used Eq. (49). The powers of $`L_{x_0}`$ in the numerator have been calculated in Eq. (54). From explicit expression (32) of $`L_{x_0}`$ we deduce
$$(L_{x_0}x_0)^\mu =qx_0\delta _0^\mu \frac{\lambda }{q[2]}x^\mu .$$
(76)
Putting things together we obtain
$$\begin{array}{cc}\hfill (x_0^n)^\mu & =\frac{[2]}{\lambda }\left(\frac{U_nc^{\frac{n}{2}}q^{n2}U_{n1}c^{\frac{n1}{2}}q^{n1}x_0x_0^n}{x^2}\right)x^\mu \hfill \\ & +U_{n1}c^{\frac{n1}{2}}q^{n1}\delta _0^\mu ,\hfill \end{array}$$
(77)
where $`U_n=U_n(b/\sqrt{c})`$ denotes the Chebyshev polynomials of the second kind with the same argument as in Eq. (54).
In order to generalize the formulas to general functions $`f=f(x_0)`$ we first note that we can deduce from Eq. (74)
$$f(x_0)=\frac{f(qL_{x_0})f(x_0)}{qL_{x_0}x_0}x_0,$$
(78)
which could be viewed as matrix valued generalization of the Jackson derivative. In order to compute this expression we need to generalize formula (54) for the powers of $`L_{x_0}`$ to general functions. The result is given by:
###### Proposition 2.
Let $`L_{x_0}`$ be the $`𝒳`$-valued 4$`\times `$4-matrix given by Eq. (32). Let $`b`$ and $`c`$ be defined as in Eq. (50). Define
$$\tau _\pm :=b\pm \sqrt{b^2c}=\frac{1}{2}\left([2]x_0\pm \lambda \sqrt{(x_0)^2\frac{4}{[2]^2}(x^2)}\right),$$
(79)
and
$$\mathrm{\Pi }_\pm :=\frac{1}{2}\left(1\pm \frac{L_{x_0}b}{\sqrt{b^2c}}\right).$$
(80)
Then for any polynomial function $`f`$ in one variable we have
$$f(L_{x_0})=f(\tau _+)\mathrm{\Pi }_++f(\tau _{})\mathrm{\Pi }_{}.$$
(81)
###### Proof.
First, we note that Chebyshev polynomials of the second kind can be written as
$$U_{n1}(x)=\frac{(x+\sqrt{x^21})^n(x\sqrt{x^21})^n}{2\sqrt{x^21}}.$$
(82)
For $`x=b/\sqrt{c}`$ this identity takes the form
$$U_{n1}(b/\sqrt{c})c^{\frac{n1}{2}}=\frac{\tau _+^n\tau _{}^n}{\tau _+\tau _{}},$$
(83)
with $`\tau _\pm `$ defined as in the proposition. Inserting (83) into (54) yields for an arbitrary function
$$f(L_{x_0})=\frac{f(\tau _+)f(\tau _{})}{\tau _+\tau _{}}(L_{x_0}b)+\frac{1}{2}\left(f(\tau _+)+f(\tau _{})\right),$$
(84)
which can be written using $`\mathrm{\Pi }_\pm `$ in the form of Eq. (81). ∎
The $`𝒳`$-valued matrices $`\mathrm{\Pi }_\pm `$ arise here naturally because they are orthogonal, complementary idempotents,
$$\mathrm{\Pi }_\pm ^2=\mathrm{\Pi }_\pm ,\mathrm{\Pi }_+\mathrm{\Pi }_{}=0=\mathrm{\Pi }_{}\mathrm{\Pi }_+,\mathrm{\Pi }_++\mathrm{\Pi }_{}=1.$$
(85)
This property ensures that Eq. (81) is consistent with the algebra structure of functions, because it can be immediately verified that $`(fg)(L_{x_0})=f(L_{x_0})g(L_{x_0})`$.
We can now insert Eq. (84) into Eq. (78) to obtain
$$f(x_0)=\frac{f(q\tau _+)f(x_0)}{q\tau _+x_0}\mathrm{\Pi }_+x_0+\frac{f(q\tau _{})f(x_0)}{q\tau _{}x_0}\mathrm{\Pi }_{}x_0$$
(86)
for the derivative of an arbitrary function of the time coordinate. For this formula to be explicit it remains to calculate
$$(\mathrm{\Pi }_\pm x_0)^\mu =\frac{\frac{1}{2}\left(\pm x_0+\sqrt{(x_0)^2\frac{4}{[2]^2}(x^2)}\right)\delta _0^\mu \frac{1}{q[2]}x^\mu }{\sqrt{(x_0)^2\frac{4}{[2]^2}(x^2)}}.$$
(87)
Finally, we can combine the previous result (72) for the derivative of functions of the four-length $`x^2=x_\mu x^\mu `$ and (86) for the derivative of functions of the time coordinate $`x_0`$ into a formula for functions $`f=f(x^2,x_0)`$ of both variables,
$$\begin{array}{cc}\hfill f(x^2,x_0)& =\frac{f(q^2x^2,x_0)f(x^2,x_0)}{q^2x^2x^2}x^2\hfill \\ & +\frac{f(q^2x^2,q\tau _+)f(q^2x^2,x_0)}{q\tau _+x_0}\mathrm{\Pi }_+x_0\hfill \\ & +\frac{f(q^2x^2,q\tau _{})f(q^2x^2,x_0)}{q\tau _{}x_0}\mathrm{\Pi }_{}x_0.\hfill \end{array}$$
(88)
This formula can be further simplified to
$$\begin{array}{cc}\hfill f(x^2,x_0)=& \frac{f(q^2x^2,q\tau _+)f(x^2,x_0)}{q\tau _+x_0}\mathrm{\Pi }_+x_0\hfill \\ \hfill +& \frac{f(q^2x^2,q\tau _{})f(x^2,x_0)}{q\tau _{}x_0}\mathrm{\Pi }_{}x_0.\hfill \end{array}$$
(89)
#### 3.4.5 Separation of variables
We have seen that the derivatives of functions of $`x_\pm `$, $`x_{30}`$, and $`x^2`$ can be written in terms of Jackson $`q`$-derivatives. Derivatives of functions of the time coordinate, however, do not have this property but depend on the four-length $`x^2`$, as well. In the language of partial differential equations, the partial derivative $`^\mu `$ is not separated with respect to the standard time coordinate of quantum Minkowski space. We will now show that there is a remarkable non-linear transformation of coordinates such that the partial derivatives are separated in the new coordinates.
These new coordinates $`\xi _\pm =\xi _\pm (x^2,x_0)`$ are given by
$$\xi _\pm =\frac{1}{2}\left(x_0\pm \sqrt{(x_0)^2\frac{4}{[2]^2}(x^2)}\right).$$
(90)
In terms of the new variables $`\tau _\pm `$ is expressed as
$`\tau _+`$ $`=q\xi _++q^1\xi _{}`$ (91)
$`\tau _{}`$ $`=q^1\xi _++q\xi _{}.`$
The back transform is given by
$`x_0`$ $`=\xi _++\xi _{}`$ (92)
$`x^2`$ $`=[2]^2\xi _+\xi _{}.`$
A function $`f=f(x^2,x_0)`$ is expressed in terms of the new variables by the transformed function
$$\begin{array}{cc}\hfill \stackrel{~}{f}(\xi _+,\xi _{})& :=f(x^2(\xi _+,\xi _{}),x_0(\xi _+,\xi _{}))\hfill \\ & =f([2]^2\xi _+\xi _{},\xi _++\xi _{}),\hfill \end{array}$$
(93)
for which the quotients of Eq. (89) take the form of Jackson derivatives
$`{\displaystyle \frac{f(q^2x^2,q\tau _+)f(x^2,x_0)}{q\tau _+x_0}}`$ $`={\displaystyle \frac{\stackrel{~}{f}(q^2\xi _+,\xi _{})\stackrel{~}{f}(\xi _+,\xi _{})}{q^2\xi _+\xi _+}}`$ (94)
$`{\displaystyle \frac{f(q^2x^2,q\tau _+)f(x^2,x_0)}{q\tau _+x_0}}`$ $`={\displaystyle \frac{\stackrel{~}{f}(\xi _+,q^2\xi _{})\stackrel{~}{f}(\xi _+,\xi _{})}{q^2\xi _{}\xi _{}}}.`$
In order to completely rewrite Eq. (89) in $`q`$-derivative form we still have to rewrite the expressions $`\mathrm{\Pi }_\pm x_0`$. Towards this end, we must understand on what the projection operators $`\mathrm{\Pi }_\pm `$ actually project.
Key to understanding the operators $`\mathrm{\Pi }_\pm `$ is the calculation of the derivatives of the new coordinates viewed as functions of $`x^2`$ and $`x_0`$. Inserting the functions $`\xi _\pm =\xi _\pm (x^2,x_0)`$ into formula (89) and using the formula
$$\sqrt{\tau _\pm ^2\frac{4}{[2]^2}(x^2)}=\frac{1}{2}\left(\pm \lambda x_0+[2]\sqrt{(x_0)^2\frac{4}{[2]^2}(x^2)}\right),$$
(95)
we obtain after long calculations the compact result
$$\xi _\pm =\mathrm{\Pi }_\pm x_0.$$
(96)
In this sense $`\mathrm{\Pi }_\pm `$ can be viewed as generalized Jacobian of the coordinate transformation (90). We thus arrive at the main result of this paper:
###### Proposition 3.
Let $`f=f(\xi _+,\xi _{})`$ be a general ordered polynomial function in $`\xi _\pm `$. Then
$$f(\xi _+,\xi _{})=\frac{_{q^2}f}{_{q^2}\xi _+}\xi _++\frac{_{q^2}f}{_{q^2}\xi _{}}\xi _{}.$$
(97)
#### 3.4.6 General functions
Finally, we can combine this result with formulas (69) for derivations of functions of the non-central variables. A general polynomial in $`𝒳`$ can be written as a sum of products of functions $`g=f(\xi _+,\xi _{})`$ and $`h=h_+(x_+,x_{30})+h(x_{30},x_{})`$. Using the coproduct (18) we obtain for the derivative of the product
$$(fg)=(f)g+(\kappa Lf)g,$$
(98)
where we can derive from Proposition 2 the formula
$$\kappa Lg(\xi _+,\xi _{})=\mathrm{\Pi }_+g(q^2\xi _+,\xi _{})+\mathrm{\Pi }_{}g(\xi _+,q^2\xi _{}).$$
(99)
Inserting Eq. (99) into (98) we have to observe that $`\mathrm{\Pi }_\pm `$ does not commute with the non-central observables. In order to write down the end result, we first define for a general linear combination of the basis $``$ the matrix valued object
$$\delta f:=\mathrm{\Pi }_+\left(f|_{\xi _+q^2\xi _+}f\right)+\mathrm{\Pi }_{}\left(f|_{\xi _{}q^2\xi _{}}f\right).$$
(100)
This definition does not depend on the ordering because $`\xi _+`$ and $`\xi _{}`$ are both central. Putting things together we obtain for a general $`f\mathrm{Span}`$
$$\begin{array}{cc}\hfill f& =\frac{_{q^2}f}{_{q^2}\xi _+}\xi _++\frac{_{q^2}f}{_{q^2}\xi _{}}\xi _{}\hfill \\ & +\frac{_{q^2}f}{_{q^2}x_+}x_++\frac{_{q^2}f}{_{q^2}x_{30}}x_{30}+\frac{_{q^2}f}{_{q^2}x_{}}x_{}\hfill \\ & +\frac{_{q^2}(\delta f)}{_{q^2}x_+}x_++\frac{_{q^2}(\delta f)}{_{q^2}x_{30}}x_{30}+\frac{_{q^2}(\delta f)}{_{q^2}x_{}}x_{}.\hfill \end{array}$$
(101)
The first two lines look like a generalized chain rule with partial Jackson derivatives as outer derivatives. The partial Jackson derivatives in the last line are to be understood not to act on the projections $`\mathrm{\Pi }_\pm `$ contained in $`\delta f`$. By definition of $`\delta f`$ the last line vanishes for $`q1`$ and the first two lines reproduce the usual chain rule.
## 4 Solutions of quantum wave equations
In order to give an application of the separation of variables and the formulas for derivations, let us now come back to the initial problem of calculating solutions of the quantum Klein-Gordon equation (2). It can be shown that, just as in the commutative case, the space of solutions is generated as representation of the quantum algebra by a momentum eigenstate . However, unlike in the commutative case, there is only a finite number of such momentum eigenstates. In the massless case it is the light cone state, defined by
$$p_0\psi =k\psi ,p_3\psi =k\psi ,p_\pm \psi =0,$$
(102)
where $`k`$ is related to the helicity. In the massive case it is the rest state
$$p_0\psi =m\psi ,p_3\psi =0,p_\pm \psi =0,$$
(103)
where $`m`$ is the mass. The difficult part is to calculate these states within the noncommutative differential calculus as solutions of quantum differential equations. All other solutions of the quantum Klein-Gordon equation can then be obtained by quantum Lorentz transformations, which are well known and straightforward to calculate.
### 4.1 The massless case
Let us start with the massless case which turns out to be quite simple. As $`q`$-differential equation Eq. (102) reads
$$^0\psi =\mathrm{i}k\psi ,^3\psi =\mathrm{i}k\psi ,^\pm \psi =0,$$
(104)
where we have raised the indices with the $`q`$-Minkowski metric, e.g., $`^3=_3`$. Using the coordinate free notation and the lightcone coordinate $`x_{30}`$ to be a function of the light cone coordinate $`x_{30}=x_3x_0`$, we can write Eq. (104) as
$$\psi =\mathrm{i}k\psi x_{30}.$$
(105)
Since $`x_{30}`$ is one of the variables separating the noncommutative differential calculus, we can conclude that (105) can be solved by a function $`\psi =\psi (x_{30})`$ of $`x_{30}`$ alone. Using the formulas (69) for the derivatives we obtain
$$\psi (x_{30})=\frac{_{q^2}\psi }{_{q^2}x_{30}}x_{30}=\mathrm{i}k\psi x_{30},$$
(106)
which is equivalent to
$$\frac{_{q^2}\psi }{_{q^2}x_{30}}=\mathrm{i}k\psi (x_{30}).$$
(107)
This equation is solved by the well-known $`q`$-exponential function
$$\mathrm{exp}_q(z)\mathrm{e}_q^z:=\underset{n=1}{\overset{\mathrm{}}{}}\frac{z^n}{(n)_{q^2}!},$$
(108)
where the $`q`$-factorial is defined in the obvious way as
$$(n)_{q^2}!=(n)_{q^2}(n1)_{q^2}\mathrm{}(1)_{q^2}.$$
(109)
With the $`q`$-exponential the general solution of Eq. (107) and, hence, of Eq. (104) can be written as
$$\psi =C\mathrm{e}_q^{\mathrm{i}k(x_3x_0)},$$
(110)
where $`C`$ is a normalization constant.
### 4.2 The massive case
Let us now turn to the massive case where calculation of the rest state is more involved. As $`q`$-differential function the defining eigenvalue equation (103) reads
$$^0\psi =\mathrm{i}m\psi ,^A\psi =0,$$
(111)
where the index runs through $`A=,3,+`$. In index free notation we can write this as
$$\psi =\mathrm{i}m\psi x_0.$$
(112)
Since the time coordinate $`x_0`$ is not one of the variables which separates the noncommutative differential calculus, we cannot conclude that Eq. (112) can be solved by a function in $`x_0`$ alone. Since the eigenvalue equation is invariant with respect to quantum rotations we can conclude, though, that $`\psi `$ must be a function of the scalars with respect to quantum rotations $`x_0`$ and $`x^2`$. In order to solve Eq. (112) we have to separate it with the separating coordinates $`\xi _{}`$ and $`\xi _+`$. Using Eq. (97) and $`x_0=\xi _{}+\xi _{}`$ we obtain
$$\psi (\xi _+,\xi _{})=\frac{_{q^2}\psi }{_{q^2}\xi _+}\xi _++\frac{_{q^2}\psi }{_{q^2}\xi _{}}\xi _{}=\mathrm{i}m(\psi \xi _++\psi \xi _{}),$$
(113)
which can be rewritten as
$$(\frac{_{q^2}\psi }{_{q^2}\xi _+}\mathrm{i}m\psi )\xi _++(\frac{_{q^2}\psi }{_{q^2}\xi _{}}\mathrm{i}m\psi )\xi _{}=0.$$
(114)
The differential equation is now separated so we can make the ansatz
$$\psi (\xi _+,\xi _{})=\psi _+(\xi _+)\psi _{}(\xi _{}),$$
(115)
which solves Eq. (114) if
$$\frac{_{q^2}\psi _+}{_{q^2}\xi _+}=\mathrm{i}m\psi _+,\frac{_{q^2}\psi _{}}{_{q^2}\xi _{}}=\mathrm{i}m\psi _{}.$$
(116)
We see that the separation of variables of the differential calculus by the new coordinates $`\xi _\pm `$ leads to a complete separation of the wave equation. Eqs. (116) are again solved by the $`q`$-exponential function (108). As end result we obtain
$$\begin{array}{cc}\hfill \psi & =C\mathrm{e}_q^{\mathrm{i}m\xi _+}\mathrm{e}_q^{\mathrm{i}m\xi _{}}\hfill \\ & =C\mathrm{e}_q^{\frac{\mathrm{i}m}{2}\left(x_0+\sqrt{(x_0)^2{\scriptscriptstyle \frac{4}{[2]^2}}(x^2)}\right)}\mathrm{e}_q^{\frac{\mathrm{i}m}{2}\left(x_0\sqrt{(x_0)^2{\scriptscriptstyle \frac{4}{[2]^2}}(x^2)}\right)},\hfill \end{array}$$
(117)
where $`C`$ is a normalization constant. This rest state is a solution of the massive quantum Klein-Gordon equation (2). All other solutions can be obtained from this one by quantum Lorentz transformations .
## 5 Conclusion
The separation of variables (97) by a nonlinear coordinate transformation (90) seems to be the most intriguing result presented in this paper. We do not yet understand on a fundamental level what property of quantum Minkowski space is responsible for the existence of such a transformation. Nor do we know if and how this can be generalized to other quantum spaces.
In definition (90) of the separating variables a square root expression appears which, strictly speaking, is not an element of the algebra of Minkowski space proper. In order to make all statements completely rigorous we have to enlarge $`𝒳`$ by a central square root element $`\alpha `$ satisfying
$$\alpha ^2=[2]^2(x_0)^24(x^2).$$
(118)
Alternatively, one could think of the coordinates being represented by operators on a Hilbert space as in . In this case the square root would be defined through functional calculus for normal operators. In the expansion of solution (117) of the quantum Klein-Gordon equation the square root drops out. There, it can be seen as an auxiliary object which makes the notation of a certain generating function more compact. Anyway, it is clear that this mathematical subtlety does not affect the results of this paper.
Finally, we would like to note that expressions for derivatives within the noncommutative calculus can also be derived in a straightforward manner using the recursion relations defined by the commutation relations of momenta and coordinates. While this approach circumvents to some degree the mathematical machinery we have introduced here, it produces lengthy formulas which do not give the structural insight desirable for solving noncommutative differential equations. Even when free field equation can be solved by brute force, using computer algebra for example, the results turn to be quite complicated. But if already something as basic as the wave function of a free particle were described by complicated expressions, some fundamental questions of quantum field theory, such as independence of in and out states, would become very hard to address. In this respect we believe results of the type presented here to be not just a matter of computational convenience but an indispensable requirement for the further development of quantum theory on quantum Minkowski space.
## Appendix A Basic definitions and formulas
###### Definition 5.
The Hopf $``$-algebra generated by $`E`$, $`F`$, $`K`$, and $`K^1`$ with relations
$$\begin{array}{c}KK^1=1=K^1K,KEK^1=q^2E,\\ KFK^1=q^2F,[E,F]=\lambda ^1(KK^1),\end{array}$$
(119)
Hopf structure
$$\begin{array}{c}\mathrm{\Delta }(E)=EK+1E,\mathrm{\Delta }(F)=F1+K^1F,\\ \mathrm{\Delta }(K)=KK,\epsilon (E)=0=\epsilon (F),\epsilon (K)=1,\\ S(E)=EK^1,S(F)=KF,S(K)=K^1,\end{array}$$
(120)
and $``$-structure
$$E^{}=FK,F^{}=K^1E,K^{}=K$$
(121)
is called $`𝒰_q(\mathrm{su}_2)`$, the $`q`$-deformation of the enveloping algebra $`𝒰(\mathrm{su}_2)`$ .
Another useful set of generators is given by angular momentum vector $`\{J_A\}=\{J_{},J_3,J_+\}`$ which is a vector operator with respect to the Hopf adjoint action of $`𝒰(\mathrm{su}_2)`$ on itself ,
$`J_{}`$ $`:=q[2]^{\frac{1}{2}}KF,`$ (122)
$`J_3`$ $`:=[2]^1(q^1EFqFE),`$
$`J_+`$ $`:=[2]^{\frac{1}{2}}E.`$
The spin-$`\frac{1}{2}`$ representation of the antipode of this vector operator yields by $`\stackrel{~}{\sigma }_A:=[2]\rho ^{\frac{1}{2}}(SJ_A)`$ a variant of the $`q`$-Pauli matrices which we will need here. Explicitly, we obtain
$$\stackrel{~}{\sigma }_{}=[2]^{\frac{1}{2}}\left(\begin{array}{cc}0& q^{\frac{1}{2}}\\ 0& 0\end{array}\right),\stackrel{~}{\sigma }_+=[2]^{\frac{1}{2}}\left(\begin{array}{cc}0& 0\\ q^{\frac{1}{2}}& 0\end{array}\right),\stackrel{~}{\sigma }_3=\left(\begin{array}{cc}q^1& 0\\ 0& q\end{array}\right)$$
(123)
with respect to the $`\{,+\}`$ basis. $`𝒰_q(\mathrm{su}_2)`$ is quasitriangular with universal $``$-matrix
$$=\mathrm{e}^{\mathrm{}(HH)/2}\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{e}^{\mathrm{}n(n1)/2}\frac{(\mathrm{e}^{\mathrm{}}\mathrm{e}^{\mathrm{}})^n}{[n]!}(E^nF^n),$$
(124)
where $`\mathrm{}:=\mathrm{ln}q`$. The quantum Lorentz algebra contains a $`SU_q(2)^{\mathrm{op}}`$ Hopf sub-algebra , the generators of which are given explicitly by
$`a`$ $`:=K^{\frac{1}{2}}K^{\frac{1}{2}}`$ (125)
$`b`$ $`:=q^{\frac{1}{2}}\lambda K^{\frac{1}{2}}K^{\frac{1}{2}}E`$
$`c`$ $`:=q^{\frac{1}{2}}\lambda FK^{\frac{1}{2}}K^{\frac{1}{2}}`$
$`d`$ $`:=K^{\frac{1}{2}}K^{\frac{1}{2}}\lambda ^2FK^{\frac{1}{2}}K^{\frac{1}{2}}E.`$
#### Acknowledgements
The work of F.B. was supported by a PhD fellowship of the Max-Planck Society. C.B. was supported by the European Union with an outgoing international Marie-Curie fellowship under contract MOIF-CT-2005-8559.
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# A Detailed Study of Molecular Clouds toward the TeV Gamma-Ray SNR G347.3–0.5
## 1 Introduction
Supernovae are generally thought to be the most energetic events in the Galactic disk and supernova remnants (SNRs) have profound effects on the dynamics and physical/chemical processes in the interstellar space through direct interactions with the surroundings, providing a unique laboratory to test high energy processes related to strong shocks. Because of the large energy release of 10<sup>51</sup> erg, SNRs have also been considered as the origin of Galactic cosmic rays (e.g., Shklovsky 1953; Hayakawa 1956; Ginzburg 1957). Cosmic rays themselves have brought crucial information on elementary particles and in addition play a crucial role in the evolution of the interstellar medium mainly via heating and ionization, although detailed physics of particle acceleration in SNRs remains unclear. It is therefore one of the key issues in astrophysics to elucidate the interaction of SNRs with their ambient matter and the probable production of cosmic rays therein.
Observations of the interstellar gas have revealed the interactions between SNRs and the surroundings under various physical conditions including ionized hot gas, neutral atomic gas and dense molecular gas (e.g., Rho et al. 1994 (optical and X-ray), Koo & Heiles 1991 (H I),van Dishoeck et al. 1993 (molecular line)). In particular, millimeter and sub-mm wave observations of dense molecular gas have proved to be a powerful probe of the shocked gas and the interactions; broad wings of several 10 km s<sup>-1</sup> accelerated by the shocks are observed in several SNRs in the mm and sub-mm spectra of interstellar CO and other molecules (e.g., W44; Seta et al. 1998, W28; Arikawa et al. 1999, IC443; White et al. 1987). It is also to be noted that observations of molecular gas are extremely useful in constraining the kinematic distance of SNRs owing to the better angular resolutions of mm-wave telescopes and the intrinsic smaller velocity dispersions of the molecular gas compared to atomic gas.
On the other hand, observational indications for the acceleration of high-energy particles in SNRs have been quite poor particularly for cosmic ray protons. Cosmic rays consist of protons as the major constituent as well as of the other minor constituents including electrons and heavier atomic nuclei. For years, observational indications for the cosmic ray particles in SNRs were available only for the electrons that emit synchrotron radiation in the radio band. The highest energy of such electrons are however in the order of MeV, much below the highest energy of the Galactic cosmic rays around 10<sup>15</sup> eV. The detection of nonthermal X-ray emission from the shell-type SNR SN 1006 provided evidence for such very high energy electrons (e.g., Koyama et al. 1995). Considering the very short lifetime of such high energy electrons the site of their prediction is definitely constrained in the SNR itself. We are however still lacking of understanding of the origin of the major constituent, cosmic ray protons, as there is not yet clear evidence for proton acceleration in SNRs.
Two shell-type remnants with nonthermal-dominant X-ray spectra have been identified subsequently; they are G347.3–0.5 (Koyama et al. 1997; Slane et al. 1999) and G266.2–1.2 (Slane et al. 2001). G347.3–0.5 (RX J1713.7–3946) was first discovered in the ROSAT All-Sky Survey by Pfeffermann & Aschenbach (1996). Later, ASCA observations revealed that the X-ray emission from the remnant is predominantly nonthermal (Koyama et al. 1997; Slane et al. 1999). The remnant is $``$ 1 in diameter and appears to be of a shell-type morphology with the brightest emission in the western region (We shall use the equatotial coordinates to identify directions in the sky throughout the present paper). Uchiyama et al. (2003) showed that the synchrotron cutoff energy is unusually high beyond 10 keV, corresponding to 10<sup>14</sup> eV as the highest particle energies. G347.3–0.5 was also detected at the TeV $`\gamma `$-ray range with the CANGAROO telescope (Muraishi et al. 2000; Enomoto et al. 2002), and models of broadband emission point to IC scattering as the origin of the TeV $`\gamma `$-ray photons (Muraishi et al. 2000; Ellison et al. 2001). More recently, follow-up TeV $`\gamma `$-ray observations have led to a different conclusion, suggesting pion decay as the source of energetic photons (Enomoto et al. 2002), but the nature of this emission is still uncertain (see Butt et al. 2002; Reimer & Pohl 2002; Aharonian et al. 2004).
The SNR is located at l $``$ 347$`\mathrm{°}`$, towards the direction of the central part of the Galaxy, a very crowded region with various galactic objects in the line of sight. This makes it particularly confusing and uncertain to identify physically associated objects and accordingly, the distance was rather poorly determined at best. The determination of the distance is crucial in all the interpretation of observations. Originally, Pfeffermann & Aschenbach (1996) adopted a distance to the SNR of $``$ 1 kpc based on the estimation of the column density derived from the spectral analysis of the ASCA data (see also Koyama et al. 1997). Later, the possible association of a molecular cloud (later refered to as ”cloud A”) was used to derive a distance to be 6.3 $`\pm `$ 0.4 kpc from observations of the 2.6 mm CO($`J`$=1–0) line emission at 9 arc min resolution by Slane et al. (1999). In their discussion, these authors argued that an enhanced value of CO($`J`$=2–1)/($`J`$=1–0) line ratio in cloud A indicates the physical interaction with the SNR shock (see also Butt et al. 2001). This estimation, $``$ 6 kpc, was often used in the following studies of the physical properties of G347.3–0.5 (e.g., Lazendic et al. 2004; Pannuti et al. 2003, Cassam-Chenai et al. 2004a).
Most recently, Fukui et al. (2003, hereafter Paper I) performed new CO observations at 2.6 mm wavelength with NANTEN, a 4-meter mm and sub-mm telescope located at Las Campanas Observatory in Chile. These new observations have a resolution of 2$`\mathrm{}`$.6, a factor of 3.4 better than that of the CO data used by Slane et al. (1999). The authors discovered molecular gas at a distance of $``$ 1 kpc which shows hole-like distribution having a striking correlation with the X-ray image of G347.3–0.5, in addition to the presence of a broad CO line in one of the CO peaks that likely represents dynamical interaction between the SNR shock and the ambient molecular gas. They also find that cloud A shows worse spatial contact with the edge of the SNR than in the lower resolution CO data. In addition, Koo et al. (2003), Koo et al. (2004) and Cassam-Chenai et al. (2004b) derive a similar distance $``$ 1 kpc by analysing and comparing it with X-ray absorption. Most recently, the H.E.S.S. collaboration made a better resolution image of TeV $`\gamma `$-rays and has supported the interaction of the SNR with the CO clouds at $``$ 1 kpc (Aharonian et al. 2004).
Subsequent to Paper I, we present in this paper a detailed analysis of the CO distributions and the physical properties of molecular clouds associated with G347.3–0.5. We present the observations in section 2, and the results of the observations in section 3. We discuss the distance, the evolutionary phase, and the surrounding environments of the SNR in section 4 and give conclusions in section 5.
## 2 Observations
Observations were made in the <sup>12</sup>CO ($`J`$=1–0) line by using a 4-meter mm/sub-mm telescope of Nagoya University, NANTEN, at Las Campanas Observatory, Chile. The telescope had a half-power beam width (HPBW) of 2.6 at a frequency of 115 GHz and was equipped with a 4 K cryogenically cooled Nb superconductor-insulator-superconductor (SIS) mixer receiver (Ogawa et al. 1990) that provided a typical system temperature of $``$ 250 K in the single-side band, including the atmosphere toward the zenith. NANTEN was equipped with two acousto-optical spectrometers (AOS) with 2048 channels. The total bandwidth and the effective resolution were 250 MHz and 250 kHz for a wide-band mode, corresponding to a velocity coverage of 600 km s<sup>-1</sup> and a velocity resolution of 0.65 km s<sup>-1</sup>, respectively. Present observations used the wide-band mode. The data is part of the NANTEN Galactic Plane Survey. The survey carried out from 1999 to 2003 observed about 1.1 million points in total, covering a galactic longitude range of 240 degrees (from $`l`$ = 180 to 60) and galactic latitude coverage of 10–20 degrees at 4$`\mathrm{}`$–8$`\mathrm{}`$ (partly 2$`\mathrm{}`$) grid-spacing. The grid spacing for the longitude coverage is basically 4 for $`|b|`$ 5, 8 for $`|b|`$ 5, respectively. (see Matsunaga et al. and the other papers in PASJ Vol. 53 No.6 2001)
After the 4$`\mathrm{}`$ grid unbiased survey, we carried out a sensitive 2$`\mathrm{}`$ grid survey, by covering the whole area of G347.3–0.5 in April 2003. An area of $``$ 1.91 square degrees in the region of 346$`{}_{}{}^{}.7l`$ 348.0 and $``$1$`{}_{}{}^{}.2b`$ 0.2 were mapped. The integration time per point was typically $``$ 10 s/point, resulting in a typical rms noise fluctuations of 0.2–0.3 K at a velocity resolution of 0.65 km s<sup>-1</sup>. In total, 1720 positions were observed. For the intensity calibration, the room-temperature chopper wheel method was employed and the absolute intensity calibration was made by observing Orion KL \[$`\alpha `$(1950) = 5<sup>h</sup>32<sup>m</sup>47.<sup>s</sup>0, $`\delta `$(1950) = $``$52421<sup>′′</sup>\] and $`\rho `$ Oph East \[$`\alpha `$(1950) = 16<sup>h</sup>29<sup>m</sup>20<sup>s</sup>.9, $`\delta `$(1950) = $``$242213<sup>′′</sup>\] every a few hours. We assumed the absolute radiation temperatures, $`T_\mathrm{R}^{}`$, of Orion KL and $`\rho `$ Oph East to be 65 K and 15 K, respectively.
Following the NANTEN survey, we have carried out a <sup>12</sup>CO($`J`$=3–2) one-point observation toward one of the strong <sup>12</sup>CO($`J`$=1–0) peaks identified with NANTEN, peak C discussed in Paper I, using the 10.4m telescope at the Caltech Submilimeter Observatory (CSO) in April 2004. CSO is installed at the summit of Mt. Mauna Kea at an altitude of 4,000m. The observed position is ($`l`$, $`b`$) = (347$`\mathrm{°}`$.07, $``$0$`\mathrm{°}`$.400) and the beam size at the <sup>12</sup>CO($`J`$=3–2) line was 22$`\mathrm{}`$. CSO was equipped with four acousto-optical spectrometers (AOS) with differenet bandwidths, 50 MHz, 500 MHz (two), and 1.5 GHz, respectively. The number of the channels is 2048 for 1.5 GHz and 500MHz, 1024 for 50 MHz, respectively. We used two of these AOS’s, 1.5GHz and 50MHz. The effective ferequency resolution were 700 kHz for the 1.5GHz AOS and 48 kHz for the 50MHz one, respectively. The velocity coverage and resolution were 1200 km s<sup>-1</sup> and 0.61 km s<sup>-1</sup> for 1.5GHz, 40 km s<sup>-1</sup> and 0.04 km s<sup>-1</sup> for 50 MHz, respectively. The system temperature for the 1.5 GHz was $``$ 1200 K for an elevation angle of $``$ 26$`\mathrm{°}`$ in the double-side band including the atmosphere. The total integration time is 120 s and the rms noise temperature per channel is 0.25 K for 1.5GHz and 0.55 K for 50 MHz, respectively. We adopted a beam efficiency $``$ 75% for 345 GHz, the recommended value calculated by past observations of planets.
Successively, we have carried out high resolution-extensive <sup>12</sup>CO($`J`$=3–2) survey toward G347.3–0.5 using ASTE 10-meter submillimeter telescope of NAOJ (Ezawa et al. 2004) located at the site of Atacama desert in Chile, with an altitude of 4800m, in November 2004. The beam size at the <sup>12</sup>CO($`J`$=3–2) line was 23$`\mathrm{}`$. We have observed the three regions containing the CO peaks selected from Paper I with a grid spacing of 30$`\mathrm{}`$ (partially 1$`\mathrm{}`$), including 780 observed points in total. ASTE was equipped with four digital backend system (auto-correlator) with 2048 channels and we could select either of the two observational modes (wide-band or narrow-band). The total bandwidths of the wide-band mode and the narrow-band one were 512 MHz and 128 MHz corresponding to an effective velocity coverages of 450 km s<sup>-1</sup> and 110 km s<sup>-1</sup>, respectively. The effective resolutions were 500 kHz and 125 kHz, corresponding to a velocity resolution of 0.43 km s<sup>-1</sup> and 0.12 km s<sup>-1</sup>, respectively. We used the wide-band mode. The system temperature was 300-400 K for an elevation angle of $``$ 40-70$`\mathrm{°}`$ in the double-side band including the atmosphere. Total integration time per point is 30 s and the rms noise temperature per channel is 0.4-0.9 K, respectively. The beam efficiency was $``$ 50-60% for 345 GHz in daytime. For system check and absolute intensity calibration, M17 SW \[$`\alpha `$(1950) = 18<sup>h</sup>17<sup>m</sup>30.<sup>s</sup>0, $`\delta `$(1950) = $``$16136<sup>′′</sup>\] was observed every day.
## 3 Results
### 3.1 Large Scale CO Distribution
Figure 1(a) shows a CO map covering from $`l`$ = 343$`\mathrm{°}`$ to $`l`$ = 352$`\mathrm{°}`$ taken with NANTEN as part of the galactic plane CO survey (Fukui et al. 2005, in preparation) superposed on the ROSAT X-ray boundary of SNR G347.3–0.5. The SNR happens to be located toward a hole of the integrated CO intensity. Figure 1(b) indicates a velocity-longitude diagram in the same longitude range with Figure 1(a) for a velocity ( = radial velocity with respect to the Local Standard of Rest = $`V_{\mathrm{LSR}}`$) range from -150 km s<sup>-1</sup> to 20 km s<sup>-1</sup>. The velocity range of the molecular gas argued to be interacting with the SNR is close to $``$ 0 km s<sup>-1</sup>, from $``$12 to $``$3 km s<sup>-1</sup> (Paper I). It is seen that the lower negative velocity edge at $`V_{\mathrm{LSR}}`$ $``$ from 0 km s<sup>-1</sup> to $``$10 km s<sup>-1</sup> of the Sagittarius arm exhibits an intensity depression in $`l`$ $``$ 347$`\mathrm{°}`$-348$`\mathrm{°}`$. This corresponds to a cavity-like CO distribution delineating the SNR G347.3–0.5 in Paper I, having a size similar to that of the SNR. More significant is a hole of CO emission at $`V_{\mathrm{LSR}}`$ from $``$ $``$10 km s<sup>-1</sup> to $``$ $``$35 km s<sup>-1</sup> in the same longitude range of l $``$ 347$`\mathrm{°}`$-348$`\mathrm{°}`$, associated with the Sagittarius arm. This hole corresponds to a molecular supershell named as SG 347.3–0.0–21 at a distance of $``$ 3 kpc (Matsunaga et al. 2001), which is likely due to multiple supernovae/stellar winds by OB stars. This chance coincidence with the supershell in the line of sight causes the apparent CO depression toward G347.3–0.5 in Figure 1(a), while the size of the supershell, $``$ 1.5$`\mathrm{°}`$, is much larger than that of G347.3–0.5 and they are not likely physically connected. Another feature of past concern is cloud A located at $`l`$ $``$ 348$`\mathrm{°}`$ and $`V_{\mathrm{LSR}}`$ $``$ $``$90 km s<sup>-1</sup>, which was ascribed to be associated with the SNR G347.3–0.5 by Slane et al. (1999). This cloud is rather isolated on the far side of the Norma arm.
### 3.2 CO Clouds towards the SNR: Coarse Velocity Channel Maps
In Figure 2, we present a montage of CO velocity channel maps taken every 10 km s<sup>-1</sup> from $`V_{\mathrm{LSR}}`$ =-160 km s<sup>-1</sup> to 20 km s<sup>-1</sup>. The directions seen from the center of the SNR in the equatrical coordinates are denoted in Figure 2(a), such as N (north), NE (northeast), and so on. The central position is denoted by the cross, where we adopt ($`l`$, $`b`$) = (347.3, $``$0.5), the catalogued position originally given in Pfeffermann & Aschenbach (1996).
The kinematic distance corresponding to the velocity centroid at each panel is indicated as calculated by using the Galactic rotation curve model (Brand & Blitz 1993). The outer boundary of the ROSAT X-ray image is superposed in each panel for reference. At various velocities, the SNR appears to be in contact with CO features as naturally expected in the direction of the galactic center, but most of the apparent coincidence must be fortuitous.
At $`V_{\mathrm{LSR}}`$ $``$ $``$160 to $``$110 km s<sup>-1</sup> (a)–(e) the association of CO with the SNR is not obvious. At $``$100 to $``$90 km s<sup>-1</sup> (g), a CO cloud towards ($`l`$, $`b`$) = (347$`\mathrm{°}`$.9, $``$0$`\mathrm{°}`$.3) appears close to the boundary of the SNR, which is cloud A. The present data taken at a higher angular resolution than that of Slane et al. (1999) shows that the western side of cloud A is not well delineating the SNR boundary compared to the lower resolution data at 9$`\mathrm{}`$ particularly an a latitude range greater than $``$0$`\mathrm{°}`$.3. Another CO feature is seen at $`V_{\mathrm{LSR}}`$ $``$ $``$80 to $``$70 km s<sup>-1</sup> (i) towards ($`l`$, $`b`$) = (347$`\mathrm{°}`$.30, 0$`\mathrm{°}`$.05), which was also suggested to be interacting with the SNR (Slane et al. 1999), although the kinematic distances of cloud A and the $``$70 km s<sup>-1</sup> cloud differ by $``$ 3 kpc, making it highly unlikely that both of the clouds are a connected single cloud interacting with the SNR. At $`V_{\mathrm{LSR}}`$ $``$ $``$60 to $``$50 km s<sup>-1</sup> (k), weak CO features appear to surround the SNR on the northwest, while this was not noted before. The panel at $`V_{\mathrm{LSR}}`$ $``$ $``$50 to $``$40 km s<sup>-1</sup> (l) shows little CO emission and these velocities correspond to the inter arm. At $`V_{\mathrm{LSR}}`$ $``$ $``$40 to $``$20 km s<sup>-1</sup> (m and n) the supershell SG 347.3–0.5–21 mentioned above are seen. This supershell has a significantly larger radius than that of the SNR G347.3–0.5 while some of the CO features appear to overlap with the SNR perhaps by chance. What is most remarkable a panel (p) at $``$10 to $``$0 km s<sup>-1</sup> which shows that the CO distribution appears to well delineate the SNR boundary. For over three quarters of the SNR boundary except for the southeast edge, the X-ray image seems to be bound by the CO emission. This is the velocity component ascribed to be interacting with the SNR in Paper I. More details of the features will be shown in the following sections (section 3.3 and 3.4). We also note a strong CO peak at ($`l`$, $`b`$) = (347$`\mathrm{°}`$.1, $``$0$`\mathrm{°}`$.4) in panel (o) is likely associated with the SNR (see section 3.4). Beyond $`V_{\mathrm{LSR}}`$ $``$ 0 km s<sup>-1</sup>, the CO emission shows a fairly good correspondence with the depression of the X-ray image towards $`l`$ $``$ 347$`\mathrm{°}`$.0 to 347$`\mathrm{°}`$.4 and $`b`$ $``$ $``$0.$`\mathrm{°}`$6 to $``$0$`\mathrm{°}`$.7, which is elongated by half a degree from the southwest to east. This depression is ascribed to the absorption by a local molecular gas at 0 to 10 km s<sup>-1</sup> in section 3.6. We note that most of the CO emission seen in panel (q) is due to the local clouds independent of the others at minus velocities. The local clouds have narrow linewidths of $``$ 1.5 km s<sup>-1</sup> and peak velocity of $``$ 6 to 7 km s<sup>-1</sup>, clearly discerned from the distant clouds in $`V_{\mathrm{LSR}}`$ $`<`$ 0 km s<sup>-1</sup>, as seen in Figure 1(b).
Figure 1(a) shows that the molecular distribution towards G347.3–0.5 generally shows a depression toward the SNR and Figure 2 indicated the same trend. We shall here introduce a factor to describe the degree of geometrical contact between the CO and X-ray images, which may help one to test how the molecular gas is in contact with the SNR at each distance. This factor, a ”covering factor”, is defined as a ratio of a covering angle of the CO distribution surrounding the SNR boundary to the whole angle of the SNR in each velocity range every 10 km s<sup>-1</sup>. We adopt the center of the angular measurement to be ($`l`$, $`b`$) = (347$`\mathrm{°}`$.3, $``$0$`\mathrm{°}`$.5). For simplicity, this covering angle is calculated only for the overlapping region of the SNR boundary and CO lowest contours in the channel maps shown in Figure 2. The contributions from the isolated clouds located inside the X-ray boundary are not included. Figure 3(a)$``$(b) are schematics for estimating the covering factor. The factor is $``$ 10% in the case of Figure 3(a) and 80% in Figure 3(b), respectively.
Since the distribution of molecular clouds toward G347.3–0.5 generally tends to surround the SNR and to avoid the interior region of the SNR in any velocity range, except for the positive velocity corresponding to the local cloud components, the covering factor gives a quantitative measure of the overall spatial coincidence between CO and X-ray around the boundary of the SNR. The covering factor is shown as a function of $`V_{\mathrm{LSR}}`$ in Figure 4. For most velocity ranges, the factor is less than $``$ 40%, but at around $``$10 km s<sup>-1</sup> it becomes largest, $``$ 80%, indicating that the correlation of CO with the SNR is the highest at that velocity range in panel (p) of Figure 2. We suggest this lends another support for the assignment of the low velocity CO emission as interacting with the SNR in Paper I.
### 3.3 Detailed Comparisons with X-rays
In order to make a more detailed examination of the spatial correlation between the X-ray features and the interacting molecular gas at velocity range from $``$ $``$20 km s<sup>-1</sup> to $``$ 0 km s<sup>-1</sup>, we show in Figure 5 a superposition of ROSAT image for the whole SNR and CO and in Figure 6 a superposition of the XMM image at harder energy range (2-7 keV) and at higher angular resolution (15$`\mathrm{}`$) (Hiraga et al. 2004). Figure 5(a) shows the gray scale and contour map of <sup>12</sup>CO($`J`$=1–0) integrated intensity and Figure 5(b) shows a superposition of the ROSAT X-ray image at an energy range of 0.1-2.0 keV on the CO for a velocity range from $``$12 km s<sup>-1</sup> to $``$3 km s<sup>-1</sup> which as we find shows the best correlation with the SNR.
In order to estimate the physical parameters of the molecular gas interacting with the SNR, we need to identify molecular clouds. As shown in Figures 5 and 6, we named CO peaks from A through Y, while we admit the complicated CO distribution makes it difficult to identify clouds uniquely. A, B, C, and D, have been used in Paper I, and in addition to these, we indicate other peaks of weaker intensity as E through Y. Basic physical properties of the clouds tentatively identified are listed in Table 1. We note the present cloud definition has not signifacant influence on later discussion of the physical aspects of the interaction.
The method of defining a CO cloud is as follows:
(1) In Figure 5, an area enclosed by the contour of 8.5 K km s<sup>-1</sup> (equal to 5 $`\sigma `$ level for the integration for a velocity range of 8 km s<sup>-1</sup>), which is significantly above the extended emission, is defined as an individual cloud.
(2) In case that a CO peak has relatively weak intensity ($``$ 6.5 K km s<sup>-1</sup>), the area enclosed by the contour of 4.5 K km s<sup>-1</sup> is identified as an individual cloud (O, Q, R, S, T, U, and V).
(3) For a cloud for which it is hard to define a closed boundary with the contour of 8.5 K (because the area enclosed by cotour is too small or extended filamentarily), a circular area with a radius of 4$`\mathrm{}`$, slightly larger than the beam size, around the peak position is defined as a size of the cloud for simplicity (F, I, H, K1, and K2).
(4) When a peak position is separated by 0.3$`\mathrm{°}`$ or larger from the nearest boundary of the ROSAT X-ray image in Figure 5(b), the cloud is regarded as not associated directly with the SNR and is not taken into account (e.g., ($`l`$, $`b`$) $``$ (346$`\mathrm{°}`$.73, 0$`\mathrm{°}`$.02), (347$`\mathrm{°}`$.71, 0$`\mathrm{°}`$.10)).
(5) When a CO spectrum has a double-peaked profile, each velocity component is defined as an independent cloud in case that the velocity difference between the two peaks is larger than 5 km s<sup>-1</sup>, $``$ 2.5 times larger than the typical linewidth of a cloud, (K1, K2, N1, and N2). The two components are devided at the velocity of the intensity minumum between the peaks.
(6) When two or more peaks are contained in one area enclosed by a contour, the area devided at the minimum in the intensity map are defined as independent peaks.
In Figure 5, The stronger CO peaks tend to be located close to the X-ray intensity peaks on the west and northwest. The most prominent X-ray peaks at ($`l`$, $`b`$) = (347$`\mathrm{°}`$.1, $``$0$`\mathrm{°}`$.3) and (347$`\mathrm{°}`$.3, $``$0$`\mathrm{°}`$.05) are well correlated with CO peaks; X-ray peak (347$`\mathrm{°}`$.1, $``$0$`\mathrm{°}`$.3) is associated with peaks A and C. B is not obvious here but as shown later B is associated with an X-ray peak at harder energy band. The X-ray peak at (347$`\mathrm{°}`$.3, $``$0$`\mathrm{°}`$.05) is associated with D and L, where D is coincident with the CANGAROO and H.E.S.S. TeV $`\gamma `$-ray peak (Paper I, Aharonian et al. 2004). Most of these CO peaks including E and G are located just on the outer edge of the SNR, as is consistent with a picture that the SNR shock is impacting these molecular clumps from the inward. We also note that peaks O, Q and R correspond to the X-ray peaks at ($`l`$, $`b`$) = (347$`\mathrm{°}`$.52, $``$0$`\mathrm{°}`$.30) , (347$`\mathrm{°}`$.72, $``$0$`\mathrm{°}`$.50), and (347$`\mathrm{°}`$.60, $``$0$`\mathrm{°}`$.73), respectively. On the south, the molecular gas including CO peaks S-Y also well delineates the SNR boundary. To summarize, at scales of arc min, the X-ray image shows a good spatial correlation with NANTEN CO image.
Figure 6 shows an even more striking correspondence between CO and X-ray at 10$`\mathrm{}`$ scales taken with XMM at the energy range harder than that in Figure 5. This X-ray image should be less affected by foreground absorption than the ROSAT image. The CO peaks A and B are well delineated by the thin X-ray filament, which becomes clear in this harder energy range, extending by 15$`\mathrm{}`$ in the northwest-southeast. This thin filament is not well traced in Figure 5 perhaps due to the absorption effect in the southeastern half of the filament facing CO peak B and also in part due to lower resolution of ROSAT. Peak C appears surrounded by two thin X-ray features of a few to several arcmin lengths on the north-west and south-east. Peak D is also X-ray bright on the south. The other peaks are all located just outside the X-ray features. The inside of the SNR around the point source (1WGA J1713.4–3949) are almost empty in the CO distribution. To summarize, this comparison further confirms the remarkable spatial correlation seen in Figure 5 at a finer angular resolution by a factor of three.
### 3.4 Sub-mm Results; CO $`J`$=3–2 Emission
Three regions of probable sites of the interaction, <sup>12</sup>CO($`J`$=1–0) peaks of A, C, and D, have been mapped in the <sup>12</sup>CO($`J`$=3–2) emission with ASTE. The same transition has also been observed toward peak C with the CSO 10.4m telescope. The intensity distributions of the CO ($`J`$=3–2) emission toward peak A, C, and D are shown in Figure 7. The ratios of the integrated intensities of the $`J`$=3–2 and $`J`$=1–0 emission are $``$ 0.7, $``$ 0.8, and $``$ 0.5, towards peaks A, C, and D, respectively, where each of the CO intensity has been integrated within a contour at half-intensity levels (11 $`\sigma `$, 8 $`\sigma `$, and 6 $`\sigma `$ level in the $`J`$ = 3–2 distribution, respectively.) We also note that the $`J`$=3–2 line is generally not excited to be observable with the current typical sensitivity in most of the cold molecular clouds in the Galaxy; Towards peak C, the $`J`$ = 3–2 emission is seeen at $`V_{\mathrm{LSR}}`$ = $``$11 km s<sup>-1</sup> and $``$90 km s<sup>-1</sup>, only toward two of the four components seen in the $`J`$ = 1–0 emission (Figure 8).
Here we shall use the large velocity gradient (LVG) radiative transfer model to obtain constraints on the physical conditions in the molecular gas. The LVG model is a useful method to solve the equation of radiation and collision for molecular excitation states. In the LVG model, it is assumed that the molecular cloud has the velocity gradient, and that all photons radiated in the cloud escape without absorption because a different position has a different Doppler velocity shift. In this analysis, we use a LVG calculation method developed by Kim et al. (2002) based on Goldreich & Kwan (1974). In the model, the radiative transfer equation gives population of each excitation state as a function of the kinetic temperature of the gas, $`T_{\mathrm{kin}}`$, for the molecular hydrogen number density $`n(\mathrm{H}_2)`$. For each observed point of CO ($`J`$=3–2) and ($`J`$=1–0), we invert these functions to determine the $`T_{\mathrm{kin}}`$ and $`n(\mathrm{H}_2)`$ corresponding to the observed line ratios (see Figure 9). The model assumes a plane-parallel cloud geometry. The parameter needed for the model is the ratio $`X`$(CO)/$`V`$, where $`X`$(CO) is the fractional CO abundance and $`V`$ is the velocity gradient. We adopt $`X`$(CO)/$`V`$ = 10<sup>-4.5</sup> pc km<sup>-1</sup> s in this case.
With the LVG analysis, we find that peaks A and C require rather high density and temperature as listed in Table 2. The lowest temperature for peak A must exceed 30 K if we adopt typical density less than $``$ 10<sup>3</sup> cm<sup>-3</sup>; we note that the typical density and temperature in local <sup>12</sup>CO emitting dark clouds are 10<sup>3</sup> cm<sup>-3</sup> and 10 K, respectively (e.g., Snell et al. 1981). This indicates that the high-$`J`$ transition of CO is appreciably excited in the peaks A and C, and that the gas is warm, being consistent with the shocked molecular gas (e.g., Arikawa et al. 1999, Irwin & Avery 1992, Zhu et al. 2003). In particular, peak A showing the good correlation with the X-ray bright filament requires higher temperature than in typical quiescent molecular gas in the solar vicinity. In the excitation condition of typical cold molecular gas in nearby dark clouds with low mass star formation ($`T_{\mathrm{kin}}`$ $``$ 10-20 K, $`n(\mathrm{H}_2)`$ $``$ 10<sup>3-4</sup> cm<sup>-3</sup> ), the $`J`$=3–2 transition is only weakly excited and the line intensity ratio between the $`J`$=3–2 and $`J`$=1–0 lines is likely below 0.3, significantly less than what is observed in the three peaks of G347.3–0.5 (see Figure 9). Unfortunately, the observational data in the CO($`J`$=3–2) emission is quite limited over the Galaxy and only a limited number of areas have been observed so far (e.g., Hunter et al. 1997; Yamaguchi et al. 2003). In order to make an estimate on the general CO $`J`$=3–2/$`J`$=1–0 ratio in typical molecular clouds in the solar vicinity which are not subject to extreme excitation conditions, we shall here integrate the CO emission of the present $`J`$=3–2 data and NANTEN $`J`$=1–0 data of the three regions toward G347.3–0.5 (shown in Figure 7) by eliminating the regions of probable interaction in velocity space. The results are as follows; the $`J`$=3–2/$`J`$=1–0 ratios are $`<`$ 4.2 K km s<sup>-1</sup>/130 K km s<sup>-1</sup> = 0.03 in the local clouds in $`V_{\mathrm{LSR}}`$ = 0-10 km s<sup>-1</sup> (4.2 K km s<sup>-1</sup> corresponds to a 3 $`\sigma `$ detection limit of $`J`$=3–2 emission over the integrated velocity range), and 23 K km s<sup>-1</sup>/91 K km s<sup>-1</sup> = 0.25 in the Sagittarius arm ($`V_{\mathrm{LSR}}`$ from $``$35 km s<sup>-1</sup> to $``$20 km s<sup>-1</sup>), respectively. Thus, the ratio is below 0.3 in the other quiescent regions as is consistent with the above estimate for the nearby dark clouds.
We infer therefore that the present high ratios of 0.7-0.8, in particular, between the $`J`$=3–2 and $`J`$=1–0 lines in G347.3–0.5 lend a support for highly excited states rather uncommon among the galactic molecular gas, keeping in mind that a more extensive survey in the CO $`J`$=3–2 emission over the galactic molecular gas is highly desirable to better establish this in future. This temperature rise is perhaps due to local enhanced heating by the SNR, where possible mechanisms may include X-rays, $`\gamma `$-rays and/or cosmic ray protons. We shall postpone to deal with more details on the heating mechanisms until the possible contribution by another young source becomes reasonably well understood (see below).
In order to test an alternative possibility to explain the higher temperature we looked for embedded infrared sources towards the CO peaks in the IRAS point sources catalogue (1988). It turns out that all the three regions are associated with IRAS sources as listed in Table 3 (see also peak A, C, and D in Figure 7) and all of them appear to have steeply-rising far-infrared spectra explicable as embedded young stars. It is not certain if all these are really at the same distance of the molecular peaks, 1 kpc, and the physical association needs to be more carefully checked. If we tentatively assume that they are at 1 kpc, their radiation luminosities are estimated as 140-560 $`L_{}`$, which may contribute to heat up the molecular gas to some degree. Thus, we shall postpone to affirm the shock heating/compression in the two peaks A and C until better estimates of density and temperature become available and until better establishing of the association of the IRAS sources.
In conclusion, the strong sub-mm emission of CO shown in the present work certainly demonstrate highly excited conditions of the molecular gas interacting with the SNR, while it remains open if the conditions are solely ascribed to the shock interaction or if an additional heating may be caused by young stars embedded. Further efforts to better constrain the physical parameters of the molecular gas and the nature of the IRAS sources are essentially important.
### 3.5 Broad Wings
Observations made by both of ASTE and CSO indicate that peak C shows fairly strong CO $`J`$=3–2 emission with some wing-like feature. The CO $`J`$=1–0 wings are already discovered and are discussed in terms of the shock acceleration by the blast wave in Paper I. It is possible that the sub-mm wings indicate the accelerated gas due to the interaction with the SNR.
Figure 10 shows a CO $`J`$ = 1–0 profile map of the peak C region. The broad wings have $``$ 20 km s<sup>-1</sup> extent showing a peak at $``$11 km s<sup>-1</sup>, and are localized within 2-4 arc min ($``$ 0.6-1.2 pc at 1 kpc). The blue-shifted component is seen from $``$ $``$23 km s<sup>-1</sup> to $``$14.5 km s<sup>-1</sup>, and the possible red-shifted component from $``$ $``$9.0 km s<sup>-1</sup> to $``$7 km s<sup>-1</sup>, respectively. For the red-shifted side, we note that confusion with the other clouds at $`V_{\mathrm{LSR}}`$ from $``$ $``$6 km s<sup>-1</sup> to 0 km s<sup>-1</sup> makes it difficult to separate the wing component clearly. There also remains a possibility that the $`J`$=1–0 red-shifted wing-like profile is not due to the shock acceleration, but just a superposition of unshocked clouds.
As shown in Figure 11, a CO $`J`$ = 3–2 profile map, peak C exhibits broad wings. These sub-mm wings are more intense in the blue-shifted side and is more localized than $`J`$=1–0 wings within $``$ 0.3 pc of the intensity maximum of peak C. In Paper I, it is argued that the CO broad wings in peak C may represent the accelerated molecular gas due to the impacting blast wave of the SNR. For the blue-shifted wing, we confirm that this is a viable interpretation since the higher density/temperature of the wings indicated by the sub-mm spectrum are consistent with this interpretation.
The present wings are seen both in the red- and blue-shifted sides of the quiescent gas. Such a trend is not odd in shocked gas and is seen in the other shock excited cases (e.g., IC 443; van Dishoeck et al. 1993, W28; Arikawa et al. 1999). These two velocity components may be due to complicated geometry of the shock fronts; a possibility here is that the far and near sides of peak C are being shocked to produce the two components. It is also remarkable that the size of the $`J`$=3–2 wing is much smaller than that of the $`J`$=1–0 wings. This suggests that the higher excitation condition in density is localized toward the central part of peak C.
Here, we should note an alternative possibility that the wings may represent molecular outflow driven by young protostars (for molecular outflows see e.g., Lada 1985; Fukui et al. 1986; 1989; 1993). The above-mentioned IRAS source IRAS 17089–3951 toward peak C (Table 3) may lend support to this alternative. The total molecular mass of peak C is estimated to be $``$ 400 $`M_{}`$ by assuming the X-factor for <sup>12</sup>CO($`J`$=1–0), 2 $`\times `$ 10<sup>20</sup> cm<sup>-2</sup> K<sup>-1</sup> (km s<sup>-1</sup>)<sup>-1</sup> (Bertsch et al. 1993), and the luminosity of the IRAS source $``$ 300 $`L_{}`$ seems consistent with a young stellar object embedded in it. By assuming that it is a protostellar molecular outflow, the physical quantities, kinetic energy, momentum and mass of outflow, are estimated (Table 4) and it seems that these quantities are not unreasonably different from typical parameters of outflow (Fukui et al. 1989; 1993).
If this interpretation is correct, peak C may represent a pre-existing dense star-forming cloud core which happened to be close to the SN and the enhanced density in it may have allowed the core to survive against the blast wave. The possibility of the protostellar outflow in peak C should be further pursued by mid-infrared imaging and spectroscopy of the IRAS source as well as by higher resolution molecular observations.
In summary, we confirm the CO broad wings in peak C both in the $`J`$=1–0 and $`J`$=3–2 CO transitions and present two alternative interpretations; one is the shock acceleration by the SNR blast waves as in Paper I and the other is the molecular outflow driven by a protostar. We need to accumulate more observations to discern these alternatives.
### 3.6 Local Cloud Absorption
The previous comparison of 0$``$10 km s<sup>-1</sup> velocity window (Figure 2(q)) suggests absorption of soft X-rays may be affecting the apprearance of X-ray distribution. Generally speaking, the southern region of the SNR is relatively weak in X-ray, suggesting possible absorption; a most notable soft X-ray depression in the ROSAT image is a straight feature extending from ($`l`$, $`b`$) $``$ (347$`\mathrm{°}`$.0, $``$0$`\mathrm{°}`$.7) to (347$`\mathrm{°}`$.4, $``$0$`\mathrm{°}`$.9) over half a degree which may be due to absorption. In the following, we make a detailed comparison of CO with soft X-ray and show that the soft X-ray image may be considerably affected by the foreground absorption by molecular gas.
Figure 12 is for comparison with the molecular clouds in a velocity range from 6.5 km s<sup>-1</sup> to 7.5 km s<sup>-1</sup> that is probably responsible for the absorption of X-rays.
The main component of the molecular clouds is located in the southern part of the SNR, from $`l`$ = 346$`\mathrm{°}`$.9 to 347$`\mathrm{°}`$.4 and from $`b`$ = $``$0$`\mathrm{°}`$.4 to $``$1$`\mathrm{°}`$.0. Including this, we have named seven CO peaks from a through g as listed in Table 5.
We list the likely candidates for absorption below.
1) The main cloud including peaks a, b, c, and e corresponds to the weaker X-ray intensity. In particular we note that the most intense X-ray arc in the northwest toward the southwest appears bound by the main cloud at ($`l`$, $`b`$) $``$ (347$`\mathrm{°}`$.05, $``$0$`\mathrm{°}`$.45). This suggests that the sharp cut off of the X-ray arc may be due to absorption at least in part.
2) The straight X-ray depression feature from ($`l`$, $`b`$) $``$ (347$`\mathrm{°}`$.0, $``$0$`\mathrm{°}`$.5) to (347$`\mathrm{°}`$.4, $``$0$`\mathrm{°}`$.9) well corresponds to the CO ridge in the same direction.
3) A few CO peaks are also located towards the regions of weaker X-ray intensity; f at ($`l`$, $`b`$) $``$ (347$`\mathrm{°}`$.55, $``$0$`\mathrm{°}`$.2 - +0$`\mathrm{°}`$.2) and g at ($`l`$, $`b`$) $``$ (347$`\mathrm{°}`$.55, $``$0$`\mathrm{°}`$.75). The exception is d at ($`l`$, $`b`$) = (347$`\mathrm{°}`$.20, $``$0$`\mathrm{°}`$.25) where a weak CO peak is overlapped with the part of the northwestern X-ray rim.
These are promising candidates for absorption features and could be more firmly established through detailed spectral analysis in the X-ray energy spectra. The CO integrated intensities of these peaks are 6.0-12.6 K km s<sup>-1</sup>, corresponding to the atomic column density of $``$ 2.8-5.0 $`\times `$ 10<sup>21</sup> cm<sup>-2</sup> (X-factor of 2 $`\times 10^{20}`$ cm<sup>-2</sup> K<sup>-1</sup> (km s<sup>-1</sup>)<sup>-1</sup> is adopted). The absorption column density based on the recent X-ray observations (Cassam-Chenai 2004b, Hiraga et al. 2004) ranges 4-10 $`\times `$ 10<sup>21</sup> cm<sup>-2</sup>, fairly consistent with our results within a factor of $``$ 2. This column density of several $`\times `$ 10<sup>21</sup> cm<sup>-2</sup> is large enough to cause the absorption (e.g., Tatematsu et al. 1990), while the probable variation in the X-ray intensity may make it difficult to estimate the amount of absorption quantitatively. In conclusion, we have shown that the soft X-ray image of the G347.3–0.5, particularly in the south, is likely affected by the foreground absorption due to the local molecular cloud within several 100 pc of the sun, and any detailed analysis of the soft X-ray distribution is to be made by taking this into account. More quantitative study of the absorption toward the molcular clouds is desirable by using new X-ray datasets in future. We also need to consider further the possible contribution of H I. This is to be done in future with appropriate high resolution H I measurements.
Anyway, it is likely that most of these clouds in positive velocity are local components within several 100 pc of the sun. The fact that the clouds can be recognized as dark extinction features in optical pictures (Digitized Sky Survey archive data from ESO/ST-ECF Science Archive) supports this suggestion.
It may be worthy to discuss on the origin of the small positive velocities of the local clouds. One of the possible explanations is that these clouds are parts of the Gould belt. The Gould belt passes through the area of $`b`$ $``$ 0$`\mathrm{°}`$-10$`\mathrm{°}`$ at $`l`$ $``$ 347$`\mathrm{°}`$ (Taylor et al. 1987), close to the location of the clouds in issue. Taylor et al. (1987) also estimated the distance and the expansion velocity of the Gould’s belt ring to be $``$ 300 pc and $``$ 5 km s<sup>-1</sup> respectively, which are not inconsistent with the properties of these clouds. We can speculate other explanations for the positive velocities, such as the random motion or the streaming motion of the clouds. Dispersion due to random motion of local clouds is measured in order of seveal km s<sup>-1</sup> (e.g., $``$ 4 km s <sup>-1</sup> by Liszt et al. 1984, $``$ 8 km s <sup>-1</sup> by Stark & Brand 1989), and the streaming motion due to the velocity field of the Galactic density wave is estimated to be $``$ 4 km s<sup>-1</sup> in the solar vicinity (e.g. Burton & Bania 1974). Each of these hypotheses is capable explaining of the positive velocity of $``$ 6–7 km s<sup>-1</sup> in this case.
## 4 Discussion
### 4.1 Distance
Based on the good spatial correlation between CO and X-ray the authors in Paper I argued that the distance of the SNR G347.3–0.5 is most likely 1 kpc instead of 6 kpc previously favoured. The present detailed analysis of the same CO datasets also confirms this small distance.
The peak velocity range of the molecular clouds named A-Y in section 3.3. is $``$12 - $``$3 km s<sup>-1</sup> (Table 1). Most of the clouds are highly likely to be in almost the same distance and associated with the SNR, since they show good spatial coincidences with the X-ray distribution. The velocity centroid of these clouds is $``$ $``$6 km/s as remarked in Paper I, giving the kinematic distance of $``$ 1 kpc. Here we adopt the Galactic rotation curve model by Brand & Blitz (1993) for deriving kinematic distances. If one assumes that the distance of these clouds in the velocity range of $``$12 - $``$3 km s<sup>-1</sup> are uniformly 1 kpc, they are localized in an area of $``$ 20 pc (= 1$`\mathrm{°}`$ in angular size) and have a dispersion of the peak velocities of $``$ 9 km/s in maximum each other. A simple and feasible explanation of this cloud-cloud velocity dispersion is the initial proper motion. Generally, it is not unusual that localized molecular clouds have velocity dispersions of $``$ several - 10 km s<sup>-1</sup>, being due to random and/or streaming motion. Thus, in this case it is also natural to regard these clouds in the velocity range of $``$12 - $``$3 km s<sup>-1</sup> as having almost the same distance of $``$ 1 kpc.
We note that the range of the kinematic distance corresponding the velocity range of -12 - -3 km s<sup>-1</sup> is about 0.5 - 2.0 kpc. However, such an uncertainty does not conflict with our supposition that the G347.3–0.5 is at the small distance basically.
It is also noteworthy that recent work, independent of Paper I, have provided support for the small distance; these are studies of the X-ray and interstellar absorption (Uchiyama et al. 2003; Koo et al. 2003; Cassam-Chenai et al. 2004b), H.E.S.S. TeV $`\gamma `$-ray image (Aharonian et al. 2004). Koo et al. (2003) and Cassam-Chenai et al. (2004b) made analyses of the interstellar absorption by using low-resolutional H I (15–30) and CO (8) data and concluded that the X-ray absorption is consistent with the small distance around 1 kpc as opposed to the large distance around 6 kpc. Uchiyama et al. (2003) posed a constraint that the shock velocity must be $`>`$ 5000 km s<sup>-1</sup> from Chandra observations of the northwestern shell, implying that G347.3-0.5 is a quite young SNR. The physical properties of the SNR derived assuming the distance of 1 kpc is consistent with this demand for the shock velocity (see next subsection). The H.E.S.S. TeV $`\gamma `$-ray image has revealed a detailed distribution of the highest energy photons. Their major finding is the shell-like distribution of TeV $`\gamma `$-ray similar to the X-ray distribution. These authors support the distance 1 kpc by noting that the TeV $`\gamma `$-ray peaks are stronger in the northwestern region where the interaction between the local molecular gas and the SNR is suggested in Paper I and that the TeV $`\gamma `$-ray flux is weak toward the eastern regions where cloud A was suggested to be interacting.
The small distance 1 kpc is therefore acquiring even stronger supports since the publication of Paper I.
### 4.2 Evolution of the SNR
In this section we discuss the evolution of the SNR specifically for the case of 1 kpc distance which is strongly supported observationally by the present work as well as by others (e.g., Koo et al. 2003, Cassam-Chenai et al. 2004b, Aharonian et al. 2004). Slane et al. (1999) already discussed that the two distances, 1 kpc and 6 kpc, are both possible in explaining the evolution of the SNR though the time scales etc. should be quite different between the two cases. We shall here highlight the main aspects of the SNR evolution at 1 kpc; we do not intend to use the evolutionary models to discern the distance since it is no more required.
There are seven historical supernovae and G347.3–0.5 likely corresponds to the historical SN AD 393, the only one that has escaped identification among the seven recorded in the historical literature (Wang et al. 1997). Thus, we confirm that G347.3–0.5 corresponds to SN AD 393 and adopt its age as 1600 yrs.
The small distance 1 kpc drastically changes many physical parameters of the SNR from those expected in case of 6 kpc. Various physical parameters are compared for the two distances 1 kpc and 6 kpc in Table 6. At 1 kpc, the radius of the SNR is 8.7 pc and the age is as small as $``$ 1600 yrs. Because of lack of thermal X-ray emission, the emission measure of the X-ray emitting hot gas should be very small, like $``$ 10<sup>15</sup> cm<sup>-5</sup> (e.g., Slane et al. 1999). Such a very low emission measure should require very low densities of the post-shock region. It also indicates that the total mass of swept-up matter is not large enough to achieve the adiabatic phase and the remnant is still in free-expansion phase with thermal non-equilibrium between electrons and ions. If we assume that the progenitor is a massive star and the ejecta mass is 3 $`M_{}`$ (Borkowski et al. 1996), the explosion energy of 10<sup>51</sup> ergs gives the ejecta velocity of 5,800 km. This velocity is nearly equal to average velocity of the shock front $``$ 5,500 km s<sup>-1</sup> as derived from X-ray measurements (e.g., Koyama et al 1997, Uchiyama et al. 2003), indicating that the blast-wave has not been decelerated yet. If we apply a simple model that the free expansion phase ends when the swept-up mass becomes equal to the mass of ejecta, we infer that the explosion occurred in low-density surroundings having density of $`<`$ 0.01 cm<sup>-3</sup> and that the ejecta is non-radiative at present. We then obtain emission measure of $`n_\mathrm{e}n_\mathrm{H}`$ $`dl`$ $`<`$ 3 $`\times `$ 10<sup>15</sup> cm<sup>-5</sup>. This value is nearly the same as that in case of 6 kpc distance (Table 6) and is small enough to explain the lack of thermal emission. In addition, the electron temperature is estimated to be $`<`$ 0.5 keV by applying the analytic formula developed by Masai (1994) to the X-ray data (Slane et al. 1999). This low value of electron temperature is consistent with the observed X-ray spectrum. We confirm that the basic X-ray properties of the SNR are explained consistently at 1 kpc.
Slane et al. (1999) reported that the spectrum of the central object 1WGA J1713.4–3949 can be fitted with 0.38 keV blackbody model and the radius of emitting region is 0.5$`{}_{0.11}{}^{}{}_{}{}^{+0.16}`$ $`\times `$ diameter at 1 kpc. The temperature of the blackbody model is significantly higher than that of a cooling neutron star whose age is $``$ 1600 year (e.g., Slane et al. 2002). The polar cap heating model seems to be good in this case, and is able to explain the observed temperature and the radius of the emitting region plausibly at 1 kpc. Here, we have to notice another possibility still remaining that the central object is not associated with the SNR. Slane et al. (1999) for instance pointed out that the X-ray spectra of the central object can be fitted either by a power-law or thin thermal emission models with absorption having the column density comparable to the total absorption through the Galaxy in this direction. To summarize, the physical parameters are well explained in the framework of very young SNR in the free expansion phase at 1 kpc. Another detailed discussion on the SNR evolution may be found elsewhere (e.g., Cassam-Chenai et al. 2004b).
### 4.3 Origin of Gamma-Rays and Cosmic Ray Protons
Higher quality multi-wavelength data including radio, X-ray and $`\gamma `$-ray are needed to obtain definitive conclusions on cosmic ray acceleration in G347.3–0.5, the most unique object known to date in testing the proton acceleration. Higher energy X-ray imaging observations will provide synchrotron cut-off energy and more sensitive X-ray spectroscopy is needed to detect thermal X-ray emission. $`\gamma `$-ray observations with higher spatial resolution will allow us to better constrain the acceleration mechanism.
It is in this context a very important issue to understand precisely the mechanism of TeV $`\gamma `$-ray emission in G347.3–0.5. For this purpose, it is essential to discern the contributions of the cosmic ray proton component from the electronic contribution via inverse Compton process and this should be achieved by making a detailed comparison between spatially resolved images of CO and TeV $`\gamma `$-ray.
The recent H.E.S.S. image of TeV $`\gamma `$-ray (Aharonian et al. 2004) is certainly encouraging in that it is beginning to resolve the TeV-$`\gamma `$ distribution at a few arc min, similar to the angular resolution of NANTEN CO dataset. It is particularly intriguing that the TeV $`\gamma `$-ray distribution largely resembles that of the CO distribution in that they are enhanced in the northwestern rim of the SNR as noted by Aharonian et al. (2004). While the statistical significance in H.E.S.S. results may not yet be high enough to make a detailed comparison of the individual peaks, we make a preliminary estimate of the efficiency for proton acceleration over the entire SNR by using a formula from Enomoto et al. (2002), ($`E`$/10<sup>48</sup>)($`M_{\mathrm{cloud}}`$/200)($`l`$/3)<sup>-3</sup>($`d/1)^5`$ = 1.35 , where $`E`$ (erg) is the total energy of cosmic ray protons, $`M_{\mathrm{cloud}}`$ ($`M_{}`$) the molecular cloud mass interacting with them, $`l`$ (pc) the typical length of the cloud, and $`d`$ (kpc) the distance, same as that in Paper I. We assume that the interacting molecular cloud mass amounts to $``$ 2500 $`M_{}`$ by summing up the major CO peaks in contact with the SNR (from Table 1). The equation gives the total energy of cosmic ray protons of $``$ 2.4 $`\times `$ 10<sup>49</sup> erg, and then, the efficiency of the proton acceleration is estimated to be $``$ 0.024 by dividing by the total energy of the SN explosion, $``$10<sup>51</sup> erg.
## 5 Conclusions
We conclude the main results of the present work as follows;
1. The present detailed analysis of the NANTEN CO data has established that the SNR G347.3–0.5 is interacting with the molecular gas at a distance of 1 kpc instead of 6 kpc previously suggested. This is supported by subsequent H I and TeV $`\gamma `$-ray studies. The physical picture of the SNR evolution is shown to be consistent with the distance 1 kpc if the SNR is in the free expansion phase.
2. The intense CO $`J`$=3–2 emission has been detected toward two of the CO peaks, strongly indicating that the interacting molecular gas is in a highly excited states in density and temperature unusual in typical local dark clouds. This provides another support for the interaction.
3. We note an alternative possibility that IRAS point sources, candidates for embedded young stars in the CO peaks, may provide additional heating in these CO peaks, although their actual association needs to be tested by further observations in the infrared and others.
4. Peak C showing strong interaction with the SNR blast wave is established to exhibit broad CO wings in the $`J`$=1–0 and $`J`$=3–2 transitions. Two alternative interpretations for the wings are presented; one is the shock acceleration (Paper I) and the other protostellar outflow. Further observational efforts are desirable to discern them.
5. It is demonstrated that the southern part of the SNR may suffer from significant absorption by a foreground molecular gas within several 100 pc, while quantitative estimate for the absorption is not available at present.
We would like to thank all the stuff members of the NANTEN project. We also acknowledge Prof. Tadayuki Takahashi, Dr. Yasunobu Uchiyama, Dr. Hiraga Junko, for their helpful comments. We appreciate the hospitality of all the people of Las Campanas Observatory of the Carnegie Institution of Washington, and thank Prof. Tom Phillips for kindly allocating the observing time of the CSO (The Caltech Submillimeter Observatory). The authors greatly acknowledge the hospitality of all members of NAOJ (National Astronomical Observatory Japan), University of Tokyo, and Osaka Prefecture University who are working for the ASTE project, for thier dedicated supports. The NANTEN project was based on the mutual agreements between Nagoya University and Carnegie Institution of Washington. We also acknowledge that NANTEN project was realized by the contributions from many Japanese public donors and companies. This reserch was financially supported by JSPS Grant-in-Aid for Scientific Research (B) No. 14403001 and MEXT Grant-in-Aid for Scientific Research on Priority Areas No. 15071202.
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# Performance of various quantum key distribution systems using 1.55 𝜇m up-conversion single-photon detectors
## I Introduction
Quantum key distribution (QKD) allows two parties to share an unconditionally secure secret key. Security is guaranteed by the laws of quantum mechanics, ensuring that the key can be used afterwards to encrypt and decrypt secret messages as a one-time pad. The most common QKD protocols, which have been implemented in experiments over the last years Gisin et al. (2002), are the BB84 protocol, which uses single photons as information carriers Bennett and Brassard (IEEE, New York, 1984), and the entanglement-based BBM92 protocol Bennett et al. (1992). A security analysis for these protocols under realistic system parameters and against individual attacks has been performed L$`\ddot{\text{u}}`$tkenhaus (2000); Waks et al. (2002). This analysis shows that the performance of a quantum cryptography system, in terms of communication distance and secure communication rate, is determined by the characteristics of the source of single or entangled photons, and of the single-photon detectors. In addition to the BB84 and BBM92 protocols, we consider the recently proposed differential phase shift keying (DPSK) protocol, which uses a weak coherent pulse train as the information carrier Inoue et al. (2002, 2003). To this end, we develop a security analysis against certain types of hybrid attacks.
To date, fiber-optic QKD systems have invariably used InGaAs/InP avalanche photodiodes (APDs) as single-photon detectors. Recently, an alternative technology for very efficient single-photon detection at 1.55 $`\mu `$m, based on the principle of frequency up-conversion, was presented Langrock et al. (2005). Using realistic experimental parameters, we perform comparisons for the various types of sources and protocols, and show that longer communication distances and higher communication rates can be achieved using the up-conversion detector in all cases.
## II 1.55 $`\mu 𝐦`$ single-photon detectors
### II.1 InGaAs/InP avalanche photodiode
The InGaAs/InP avalanche photodiodes have been the subject of thorough investigation over the last decade due to their importance as single-photon detectors in fiber-optic QKD implementations. Although considerable progress has been achieved in the performance of these detectors Yoshizawa et al. (2004); Bethune et al. (2004); Stucki et al. (2001); Bourennane et al. (2001); Gobby et al. (2004), they exhibit low quantum efficiencies (typically on the order of 0.1), and, most seriously, they suffer from after-pulse effects caused by trapped charge carriers, which produce large dark count rates during a relatively long time. The high dark count probability imposes *gated-mode operation*, which limits their capabilities significantly. In particular, when operated in gated mode, the APD device is raised above breakdown threshold for a few nsec, which ensures low probability of a dark count and high efficiency for detecting light. Subsequently, the device is returned to below breakdown for a time long enough for any trapped charge carrier to leak away. Given that the trapping lifetime is on the order of a $`\mu `$sec, this mode allows operation at MHz rates, while the after-pulse probability is reduced by the ratio of the gate width to the time separation between gates. In a QKD application, this gate frequency determines the repetition rate of the signal pulse and, therefore, limits the attainable communication rate. Furthermore, the dark count rate, which is critical for the communication distance, is determined by the gate width, limited by the response time of the semiconductor material. Typically, gate widths of $`12`$ nsec at $`1`$ MHz repetition frequency are used with resulting dark count rates on the order of $`10^4`$/sec.
### II.2 Up-conversion detector
In the 1.55 $`\mu `$m up-conversion single-photon detector Langrock et al. (2005), a single photon at 1.55 $`\mu `$m interacts with a strong pump at 1.32 $`\mu `$m in a periodically poled lithium niobate (PPLN) waveguide, designed for sum-frequency generation at these wavelengths Roussev et al. (2004). Due to the quasi-phase-matching and the tight mode confinement over long interaction lengths achieved in a guided-wave structure, this device allows for very high conversion efficiency of the signal to the $``$ 0.7 $`\mu `$m sum frequency output. The converted photon is subsequently detected by a silicon APD. Contrary to InGaAs/InP APDs, Si APDs have high quantum efficiencies in the near-infrared (typically on the order of $`0.60.7`$), very low dark count rates, and very small after-pulse effects. The last characteristic enables *Geiger (non-gated) mode operation* of the Si APD, which does not impose any severe limitation to the attainable communication rate in a QKD system. In practice, however, the rate is limited by the dead time of Si APD detectors, which is on the order of 50 nsec for commercial devices. During this time period that follows a photo-detection event, the photodiode cannot respond to subsequent events, and, eventually, a very large photon flux saturates the device. This effect is taken into account in the calculations of Section IV.
The main characteristics of the up-conversion detector, such as the quantum efficiency, $`\eta _{\text{up}}`$, and the dark count rate, $`D_{\text{up}}`$, depend on the pump power, $`p`$ Langrock et al. (2005). When the phase-matching condition in the waveguide is met and sufficient pump power is available to achieve almost 100% photon conversion, a maximum overall quantum efficiency of 0.46 is achieved, as shown in Fig. 1. In agreement with the coupled mode theory for three-wave interactions in a waveguide, which predicts a $`\text{sin}^2`$ dependence of $`\eta _{\text{up}}`$ on $`p`$, the fitting curve of the experimental results is given by the following expression:
$$\eta _{\text{up}}(p)=a_1\mathrm{sin}^2(\sqrt{a_2p})$$
(1)
where $`a_1=0.465,a_2=79.75`$, and $`p`$ is given in mW.
On the other hand, the dark count rate is dominated by a combined nonlinear process: Initially, the pump photons are scattered by the phonons of both the PPLN waveguide and the fiber via a spontaneous Raman scattering process. This process scales linearly with the pump power, and generates a spectrum of Stokes photons, which includes the signal wavelength of 1.55 $`\mu `$m. Subsequently, the noise photons interact with the pump photons in the waveguide via the phase-matched sum-frequency generation process, and create dark counts. The combined process results in an approximately quadratic dependence of the dark counts on the pump power, as shown in Fig. 2. A more accurate polynomial fitting curve is given by the following expression:
$$D_{\text{up}}(p)=b_0+b_1p+b_2p^2+b_3p^3+b_4p^4\text{ }\text{ }\text{ }\text{(/sec)}$$
(2)
where $`b_0=50`$, $`b_1=826.4`$, $`b_2=110.3`$, $`b_3=0.403`$, $`b_4=0.00065`$, and $`p`$ is again given in mW.
An important feature of the up-conversion detector stems from the fact that the dark counts depend on the bandwidth of the waveguide, as this determines the number of noise photons. We can define a quantity, $`D_{\text{up-Hz}}=\frac{D_{\text{up}}}{B_\text{d}}`$ $`(\text{sec}^1\text{Hz}^1)`$, for a detector with bandwidth $`B_\text{d}`$, which corresponds to the dark counts per mode. Then, we can think of the ideal communication system shown in Fig. 3 with a matched filter with bandwidth equal to the bit rate $`B`$. In such a system, the dark counts per time window, $`d_{\text{up}}`$, a parameter of great importance in QKD applications, is equal to $`D_{\text{up-Hz}}`$. Note that $`d_{\text{up}}`$ is independent of the bit rate $`B`$ (or measurement time window $`1/B`$) under this optimum filtering. An InGaAs/InP APD operated in gated mode has dark counts per gate, $`d_{\text{APD}}`$, calculated by $`D_{\text{APD}}\frac{1}{B}`$, where $`D_{\text{APD}}`$ (/sec) is the dark count rate of the InGaAs/InP APD. In Fig. 4, the quantity $`d`$ is plotted for the two types of detectors as a function of the bit rate. For the InGaAs/InP APD, the typical value $`D_{\text{APD}}=10^4`$/sec is used. For the up-conversion detector, we calculate the quantity $`D_{\text{up-Hz}}`$ at the operating point of the detector, where the normalized Noise Equivalent Power (NEP), $`\sqrt{2D}/\eta `$, is minimized, which corresponds to $`D_{\text{up}}=6.4\times 10^3`$/sec and $`\eta _{\text{up}}=0.075`$. Given a bandwidth of $`B_\text{d}=50`$ GHz for the up-converter, we find that the optimum $`d_{\text{up}}`$ is $`1.3\times 10^7`$, as shown in Fig. 4. This result illustrates the significant advantage of the up-conversion detector for most practical system bit rates.
The dependence of the dark counts on the waveguide bandwidth, together with the non-gated mode operation of the Si APD and the pump power dependence of the detector characteristics, have a significant effect on the performance of a quantum cryptography system employing up-conversion detectors, as we will see in the following sections.
## III Communication rate equations
In this paper, we will consider only individual attacks, that is Eve is restricted to attack only individual bits; she is not allowed to perform a coherent attack consisting of collective quantum operations and measurements of many qubits with quantum computers. In a QKD system, the raw transmission of random bits is followed by a public exchange of information on the time of single-photon detection and the bases used by the two parties, which results in the *sifted key*. The steps of classical error-correction and privacy amplification follow. The first step serves the dual purpose of correcting all erroneously received bits and giving an estimate of the error rate. Privacy amplification is then used to distill a shorter key, the *final key*, which can be made as secure as desired. The security analyses of L$`\ddot{\text{u}}`$tkenhaus (2000); Waks et al. (2002) take all the above steps into account and derive the communication rate equations that are re-stated in Sections III.1 and III.2. In Section III.3, we derive the corresponding equation for the DPSK protocol, based on the security analysis against certain types of hybrid attacks.
### III.1 BB84 protocol
In the BB84 protocol, Alice sends Bob single photons randomly modulated in two non-orthogonal bases. Bob measures the polarization states of the single photons in a randomly chosen polarization basis. The secure communication rate of this protocol against an arbitrary individual attack, including the most commonly considered intercept-resend and photon-number splitting (PNS) attacks L$`\ddot{\text{u}}`$tkenhaus (2000), is given by the following expression:
$`R_{\text{BB84}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\nu p_{\text{click}}\{\tau (e,\beta )+f(e)[e\mathrm{log}_2e+`$ (3)
$`(1e)\mathrm{log}_2(1e)]\}`$
In the above equation, the factor $`\frac{1}{2}`$ is called the sifting parameter and is due to the fact that half of the times Alice’s and Bob’s polarization bases are not the same. The repetition rate of the transmission is given by $`\nu `$. The probability that Bob detects a photon is
$$p_{\text{click}}=p_{\text{signal}}+p_{\text{dark}}$$
(4)
Simultaneous signal and dark counts are ignored in the above expression, and the two components are given by
$`p_{\text{signal}}`$ $`=`$ $`\mu \eta 10^{(\alpha L+L_\text{r})/10}`$ (5)
$`p_{\text{dark}}`$ $`=`$ $`4d`$ (6)
where $`\mu `$ is the average number of photons per pulse, $`\eta `$ the quantum efficiency of the detector, $`\alpha `$ the loss coefficient of the optical fiber in dB/km, $`L`$ the distance in km, $`L_\text{r}`$ the loss of the receiver unit in dB, and $`d`$ the dark counts per measurement time window of the system. The coefficient 4 in Eq. (6) is due to the assumption of a passive detection unit involving four detectors at Bob’s site, as in Waks et al. (2002). For an ideal single-photon source, $`\mu =1`$, while for a Poisson source, which corresponds to the common weak laser pulse implementations Gisin et al. (2002), $`\mu `$ becomes a free variable which should be optimized.
The error rate is given by the expression:
$$e=\frac{\frac{1}{2}p_{\text{dark}}+bp_{\text{signal}}}{p_{\text{click}}}$$
(7)
where $`b`$ is the baseline system error rate, which cannot be distinguished from tampering. The last term in Eq. (3) corresponds to the additional shrinking of the sifted key due to the leakage of information to Eve during classical error correction. The function $`f(e)`$ depends on the error-correction algorithm and its values are given in Table 1 for the bi-directional algorithm developed in Brassard and Salvail (Springer, Berlin, 1994).
Finally, the main shrinking factor $`\tau (e,\beta )`$ in the privacy amplification step is related through the expression
$$\tau =\mathrm{log}_2p_\text{c}$$
(8)
to the average collision probability, $`p_\text{c}`$. This is a measure of Eve’s mutual information with Alice and Bob. In L$`\ddot{\text{u}}`$tkenhaus (2000) the following result is derived for $`\tau `$:
$$\tau (e,\beta )=\beta \mathrm{log}_2\left[\frac{1}{2}+2\frac{e}{\beta }2\left(\frac{e}{\beta }\right)^2\right]$$
(9)
The parameter $`\beta `$ is defined as the fraction of single-photon states emitted by the source:
$$\beta =\frac{p_{\text{click}}p_\text{m}}{p_{\text{click}}}$$
(10)
where $`p_\text{m}`$ is the probability that the source emits a multi-photon state. For an ideal single-photon source, $`p_\text{m}=0`$ (i.e., $`\beta =1`$), while for a Poisson source,
$$p_\text{m}=1(1+\mu )e^\mu $$
(11)
Essentially, the parameter $`\beta `$ accounts for the PNS attacks, with which Eve can obtain full information without causing any error in the communication between Alice and Bob by performing a quantum non-demolition (QND) measurement of the photon number in each pulse, keeping one photon in her quantum memory when she detects multiple photons, and applying a delayed measurement on her photon after the public announcement of the bases by Bob. This attack is a major restricting factor in the performance of a weak laser pulse implementation of the BB84 protocol. The secure communication rate decreases quadratically with the transmission of the quantum channel, $`10^{\alpha L/10}`$, for small error rate and $`p_{\text{dark}}p_{\text{signal}}1`$. On the contrary, for an ideal single-photon source implementation, under the same conditions we find $`R_{\text{BB84}}\frac{1}{2}\nu p_{\text{signal}}`$, i.e., the rate decreases only linearly with the fiber transmission.
The above security analysis is based on the assumption that Eve has a quantum memory with an infinitely long coherence time because Alice and Bob can delay the public announcement for an arbitrarily long time. If Eve is not equipped with such a quantum memory, she must perform the polarization measurement with a randomly chosen basis. In this realistic case, Eq. (9) must be modified to:
$$\tau (e,\beta )=\frac{1+\beta }{2}\mathrm{log}_2\left[\frac{1}{2}+4\frac{e}{1+\beta }8\left(\frac{e}{1+\beta }\right)^2\right]$$
(12)
### III.2 BBM92 protocol
The BBM92 protocol is the two-photon variant of BB84. Alice and Bob each share a photon of an entangled photon-pair, for which they measure the polarization state in a randomly-chosen basis out of two non-orthogonal bases. It was shown in Waks et al. (2002) that the average collision probability, $`p_\text{c}`$, for this protocol is the same as that of the BB84 with a single-photon source, i.e., with $`\beta =1`$. The shrinking factor $`\tau `$ becomes:
$$\tau (e)=\mathrm{log}_2\left[\frac{1}{2}+2e2e^2\right]$$
(13)
This indicates that there is no analog to a photon-number splitting attack in BBM92. In general, the nature of this entanglement-based protocol renders it more robust than BB84; for example it is less vulnerable to errors caused by dark counts, since one dark count alone cannot produce an error in this protocol. The equation for the secure communication rate against any individual attack is given by the following expression Waks et al. (2002):
$`R_{\text{BBM92}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\nu p_{\text{coin}}\{\tau (e)+f(e)[e\mathrm{log}_2e+`$ (14)
$`(1e)\mathrm{log}_2(1e)]\}`$
The sifting parameter is the same as in BB84, while the probability of a coincidence between Alice and Bob is
$$p_{\text{coin}}=p_{\text{true}}+p_{\text{false}}$$
(15)
The expressions for the probability of a true coincidence, $`p_{\text{true}}`$, and the probability of a false coincidence, $`p_{\text{false}}`$, are different for a deterministic entangled-photon source and a Poissonian entangled-photon source, such as a parametric down converter (PDC). They are given below, under the assumption that the source is placed halfway between the two parties Waks et al. (2002).
1. Deterministic entangled-photon source
$`p_{\text{true}}`$ $`=`$ $`\eta ^210^{(\alpha L+2L_\text{r})/10}`$ (16)
$`p_{\text{false}}`$ $`=`$ $`8d\eta 10^{(\alpha L+2L_\text{r})/20}+16d^2`$ (17)
2. Poissonian entangled-photon source
$`p_{\text{true}}`$ $`=`$ $`c_1`$ (18)
$`p_{\text{false}}`$ $`=`$ $`16d^2c_2+8dc_3+c_4`$ (19)
where
$`c_1`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}^4\chi }}{\displaystyle \frac{2t_L^2\mathrm{tanh}^2\chi }{\left[1\mathrm{tanh}^2\chi (1t_L)^2\right]^4}}`$ (20)
$`c_2`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}^4\chi }}{\displaystyle \frac{1}{\left[1\mathrm{tanh}^2\chi (1t_L)^2\right]^2}}`$ (21)
$`c_3`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}^4\chi }}{\displaystyle \frac{2t_L(1t_L)\mathrm{tanh}^2\chi }{\left[1\mathrm{tanh}^2\chi (1t_L)^2\right]^3}}`$ (22)
$`c_4`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}^4\chi }}{\displaystyle \frac{4t_L^2(1t_L)^2\mathrm{tanh}^4\chi }{\left[1\mathrm{tanh}^2\chi (1t_L)^2\right]^4}}`$ (23)
and
$$t_L=\eta 10^{(\alpha L+2L_\text{r})/20}$$
(24)
All the parameters in the above equations are defined as in the previous section. The parameter $`\chi `$, which appears in the case of the Poissonian entangled-photon source, is a free variable that depends on the average photon-pair number per pulse, i.e. the nonlinear coefficient, the pump energy and the interaction time of the down conversion process. Finally, the error rate is given by the expression:
$$e=\frac{\frac{1}{2}p_{\text{false}}+bp_{\text{true}}}{p_{\text{coin}}}$$
(25)
For small error rate and $`p_{\text{false}}p_{\text{true}}`$, the secure communication rate of BBM92 decreases linearly with the transmission of the quantum channel, similarly to the case of the BB84 protocol with a single photon source. Note that Eve does not need a quantum memory to attack the BBM92 protocol. Equation (14) is solely determined by the intercept and resend attack.
### III.3 DPSK protocol
Instead of using two non-orthogonal bases as in BB84 and BBM92, the differential phase shift keying (DPSK) protocol uses many non-orthogonal states consisting of many pulses Inoue et al. (2002, 2003). In particular, it is based on the fact that highly attenuated coherent states of many pulses with random $`\{0,\pi \}`$ phase modulation are mutually non-orthogonal. The idea of encoding the information in the phase of highly attenuated coherent pulses was first presented by Bennett in 1992 Bennett (1992). The DPSK protocol is a simpler but more efficient protocol compared to the B92 protocol. A similar protocol has also recently been proposed Gisin et al. (2004).
In the DPSK protocol, shown in Fig. 5, all pulses are highly attenuated and randomly phase-modulated by $`\{0,\pi \}`$. Each photon coherently spreads over many pulses with a fixed phase modulation pattern. In the receiver side, Bob randomly modulates the delay time $`N\times \tau `$ in his interferometer by randomly choosing a positive integer $`N`$, as shown in Fig. 5, where $`\tau `$ is the inverse of the clock frequency. After passing through Bob’s interferometer, the pulses interfere at Bob’s output beam-splitter, and which detector clicks depends on the phase difference of the two pulses separated by a time $`N\times \tau `$. Bob announces publicly the time instances at which a photon was detected and the randomly chosen positive integer $`N`$. From her modulation data Alice knows which detector recorded the event. Thus, they form a secret key by assigning a bit value to each detector. The shifting parameter is 1 since all bits are utilized during the key formation.
The security of the DPSK protocol stems from the fact that the information is encoded on the differential phase of two nonlocal pulses. This renders the protocol robust against any type of individual photon splitting attack Inoue and Honjo (2005); Honjo and Inoue (2005). In order to derive the communication rate equation, we need to calculate the privacy amplification shrinking factor, $`\tau `$, defined in Eq. (8) as a function of the average collision probability, $`p_\text{c}`$. Our analysis takes into account a hybrid attack, which consists of two types of collective attacks:
1. Beam-splitter attack
Eve uses a beam-splitter with transmission $`\eta _{\text{BS}}`$ to obtain coherent copies of the quantum state of many pulses that Alice sends to Bob. She also replaces the lossy optical fiber with a loss-less one, and the inefficient detectors at Bob’s receiver unit with ideal ones. Without Eve’s intervention, Bob’s probability of detecting a signal photon, $`p_{\text{signal}}`$, is identical to the one given in Eq. (5). In order to leave this probability unaltered, Eve has to set the beam-splitter transmission, $`\eta _{\text{BS}}`$, to:
$$\eta _{\text{BS}}=\eta 10^{(\alpha L+L_\text{r})/10}$$
(26)
where all the parameters are defined as in Section III.1. One possibility for Eve is to measure the pulses that she picks up with an interferometer with delay time $`M\times \tau `$ chosen independently from Bob’s. In this case, her information gain is calculated as follows: the probability of a detection event at Eve’s and Bob’s site at a given time slot is given by $`\mu (1\eta _{\text{BS}})`$ and $`\mu \eta _{\text{BS}}`$ respectively, where $`\mu `$ is the average number of photons per pulse. Thus, the probability of a detection event at the same time instance is equal to $`\mu ^2\eta _{\text{BS}}(1\eta _{\text{BS}})`$. On the other hand, the probability that Eve’s randomly chosen $`M`$ matches Bob’s $`N`$ is equal to $`1/N`$. Then, the probability that Eve gains bit information relative to Bob is $`\mu ^2\eta _{\text{BS}}(1\eta _{\text{BS}})/(\mu \eta _{\text{BS}}N)=\mu (1\eta _{\text{BS}})/N`$. This is true if we assume that Eve is not equipped with a quantum memory with an infinitely long coherence time or if Alice and Bob encrypt their public channel communication. However, if we allow Eve to have a quantum memory and the two parties do not encrypt their public exchange of information, Eve’s strategy can be changed in order to increase her information gain. In this case, she keeps the pulses in her quantum memory and waits for Bob’s announcement. Note that Alice and Bob can delay the public announcement for an arbitrarily long time, so Eve’s quantum memory must have an infinitely long coherence time. Then, Eve uses an optical interferometer with an active switch that allows her to interfere only the pulses for which she is aware that Bob has obtained the differential phase information. This strategy increases Eve’s probability of gaining bit information to $`2\mu (1\eta _{\text{BS}})`$. The beam-splitter attack does not cause any error in the communication between Alice and Bob, hence it gives full information, i.e., $`p_\text{c}=1`$, to Eve for a fraction of bits equal to $`\mu (1\eta _{\text{BS}})/N`$ or $`2\mu (1\eta _{\text{BS}})`$. The remaining fraction of the bits is given by:
$$\gamma =\{\begin{array}{cc}1\frac{\mu (1\eta _{\text{BS}})}{N}=1\frac{\mu }{N}+\frac{p_{\text{signal}}}{N}\hfill & \\ \text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{:without quantum memory}\hfill & \\ 12\mu (1\eta _{\text{BS}})=12\mu +2p_{\text{signal}}\hfill & \\ \text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\text{:with quantum memory}\hfill & \end{array}$$
(27)
2. Intercept-resend attack
Eve also applies an intercept and resend attack to some of the pulses that are sent to Bob after her beam splitter. In particular, Eve intercepts two pulses with a time interval $`M\times \tau `$, lets them pass through an interferometer with an identical delay $`M\times \tau `$, measures the differential phase, and according to her measurement result she sends an appropriate state to Bob. We assume that in the case of an inconclusive or vacuum outcome she sends the vacuum state, while when she measures a single photon she sends a photon split into two pulses with the correct phase difference applied between them. In this case, when Bob picks up an identical delay, $`N=M`$, and measures the central time slot, he does not detect the eavesdropping because he obtains the correct answer. However, with probability $`1\frac{1}{2N}`$ he chooses another delay, $`NM`$, or measures the side time slots, which yield random, uncorrelated results, and with probability $`\frac{1}{2}`$ these lead in error. Hence, this attack causes a bit error of $`\frac{1}{2}\left(1\frac{1}{2N}\right)`$ in the communication between Alice and Bob. If the error rate of the system is $`e`$, Eve is allowed to apply her attack to a fraction $`\frac{2e}{11/2N}`$ of the pulse-pairs in order not to exceed this error rate. With probability $`\frac{1}{2N}`$, she obtains full information for these intercepted pulse-pairs.
In summary, taking into account the hybrid attack consisting of the beam-splitter and intercept-resend attacks, we find that the fraction of bits for which Eve has no information, i.e., for which $`p_\text{c}=\frac{1}{2}`$, is equal to $`\gamma \frac{e}{N(11/2N)}`$. Thus, we have calculated the privacy amplification shrinking factor,
$$\tau (e,\gamma )=\gamma \frac{e}{N\left(1\frac{1}{2N}\right)}$$
(28)
where $`\gamma `$ is given by Eq. (27). We can now write the equation for the secure communication rate of the DPSK protocol against the hybrid attack we considered:
$`R_{\text{DPSK}}`$ $`=`$ $`\nu p_{\text{click}}\{\tau (e,\gamma )+f(e)[e\mathrm{log}_2e+`$ (29)
$`(1e)\mathrm{log}_2(1e)]\}`$
In the above equation, $`\nu `$ is the repetition rate of the transmission. The probability that Bob detects a photon, $`p_{\text{click}}`$, is defined in Eq. (4). The probability of a signal count, $`p_{\text{signal}}`$, is given by Eq. (5), while the probability of a dark count, $`p_{\text{dark}}`$, in this case is given by the expression:
$$p_{\text{dark}}=2d$$
(30)
because there are two detectors at the receiver unit. Finally, the error rate is defined in Eq. (7), and the values of $`f(e)`$ are given in Table 1.
In the case of small error rate and $`p_{\text{dark}}p_{\text{signal}}1`$, Eq. (29) gives $`R_{\text{DPSK}}\nu (1\frac{\mu }{N})p_{\text{signal}}`$ without a quantum memory, or $`R_{\text{DPSK}}\nu (12\mu )p_{\text{signal}}`$ with a quantum memory. This means that the secure rate for the DPSK protocol decreases linearly with the fiber transmission. This is in agreement with the results of Gisin et al. (2004) and Koashi (2005), who have considered a protocol similar to DPSK and a slightly modified B92 protocol respectively.
## IV Numerical results
We compare the performance of quantum key distribution systems implementing the BB84, BBM92 and DPSK protocols, when the up-conversion single-photon detector is used. In order to do that, we calculate the secure communication rate as a function of distance for fiber-optic implementations of the three protocols, based on Eqs. (3),(14) and (29) respectively. In the case of BB84 and BBM92, both ideal and realistic sources of single and entangled photons are considered. Some parameters are fixed in all simulations: the channel loss is set to $`\alpha =0.2`$ dB/km at 1.55 $`\mu `$m, the baseline system error rate is set to $`b=0.01`$, and in addition to the fiber losses we assume an extra loss of $`L_\text{r}=1`$ dB at the receiver site. As mentioned in Section III.1, in the case of a weak laser pulse implementation of the BB84 protocol, the average number of photons per pulse, $`\mu `$, is an adjustable parameter, with respect to which the rate is numerically optimized at each distance. Intuitively, such optimization is necessary because when this parameter is too low the dark counts dominate, while when it is too high the probability of multi-photon pulses becomes very large. In both cases, secure communication quickly becomes impossible. The corresponding adjustable parameter is $`\chi `$ in the case of the BBM92 protocol with a Poissonian entangled-photon source.
It is clear from the analysis of Section III that the critical parameters for the performance of a quantum cryptography system related to the single-photon detector employed are the dark counts per measurement time window, $`d`$, the quantum efficiency, $`\eta `$, and the repetition rate of the transmission that it allows, $`\nu `$. In the case of the up-conversion single-photon detector, due to the non-gated mode operation of the Si APD there is no severe limitation to the repetition rate of the experiment. In practice, the limit is set by the speed of the electronic equipment as well as by the timing jitter of the Si APD (typically $`0.50.7`$ nsec). A realistic value, compatible with currently available components, is $`\nu _{\text{up}}=1`$ GHz. As was explained in Section II.2, the limiting factor for the attainable communication rate is the dead time of the Si APD, $`t_\text{d}`$. Assuming that the photo-detection events follow a Poisson process, the probability of two events occurring in a time period larger than $`t_\text{d}`$ is given by the exponential factor $`e^{\delta \nu p_{\text{click}}t_\text{d}}`$, where $`\delta `$ depends on the number of detectors in the receiver unit. For the typical value $`t_\text{d}=50`$ nsec, this saturation factor becomes rather small at rates greater than a few MHz, limiting the final rate at small fiber losses. Using Eqs. (1) and (2), we numerically optimize the communication rate for each protocol with respect to the pump power, $`p`$, at each distance. Such optimization is intuitively necessary because depending on the communication distance an equilibrium between the values of the quantum efficiency and the dark counts of the up-conversion detector has to be established. The result of this optimization indicates the optimal regime of operation of the detector at each distance. Finally, the optimum filtering configuration, shown in Fig. 3, is assumed, which sets the measurement time window to 1 nsec.
The simulation results are shown in Figs. 6, 7, and 8 for the BB84, BBM92, and DPSK protocols respectively. Each curve features a cut-off distance, which is due to the increasing contribution of the dark counts with fiber length. The saturation effect, related to the dead time of the Si APD, is apparent for small fiber losses and high bit rates.
In the case of BB84 with a Poisson single-photon source, we observe in Fig. 6 that not allowing Eve to possess a quantum memory with an infinitely long coherence time does not have a major effect on the performance of the system. The quadratic decrease of the rate of the communication rate with the fiber length, a consequence of the PNS attacks, is a dominant factor, making this implementation unsuitable for long-distance quantum cryptography. On the contrary, the use of an ideal single-photon source allows for a significantly longer communication distance with high communication rates. However, such a source does not exist today at 1.55 $`\mu `$m, although efforts towards this goal are underway Fasel et al. (2004a).
As shown in Fig. 7, the inherently more robust entanglement-based BBM92 protocol allows for even longer communication distances, having the capability to achieve a practical 1 Hz secure key generation rate at more than 300 km with a deterministic entangled-photon source. However, technological difficulties related to entanglement generation and coincidence detection at 1.55 $`\mu `$m have limited until today this distance to 30 km Fasel et al. (2004b).
The DPSK protocol features characteristics very similar to BB84 with a single-photon source, due to its robustness to PNS attacks, as was shown in the security analysis of Section III.3. In this case, when a realistic scenario is assumed, where Eve does not possess a quantum memory with an infinitely long coherence time, or Alice and Bob encrypt their public communication, we observe in Fig. 8 a significant effect on the performance of the system. Indeed, introducing a time delay parameter $`N`$ greater than 1 enhances both the secure communication rate and the communication distance of the system considerably. Nevertheless, the advantage becomes comparatively smaller as $`N`$ increases to values greater than 10. This result shows that the DPSK protocol is a very practical and appealing alternative for a long-distance QKD system, with the potential of 1 kHz secure key generation rate over distances longer than 200 km.
For all the QKD protocols, if instead of the up-conversion detector we assume an InGaAs/InP APD with $`\nu _{\text{APD}}=10`$ MHz, which is the best gate frequency achieved to date Yoshizawa et al. (2004), and the typical values $`\eta _{\text{APD}}=0.1`$ and $`d_{\text{APD}}=10^5`$/gate Gisin et al. (2004), we find that the maximum communication distance is about half of the one achieved with an up-conversion detector, while the communication rate is two orders of magnitude lower than with the up-conversion detector, due to the gated-mode operation of the InGaAs/InP APD. Clearly, the up-conversion detector offers a great advantage over the InGaAs/InP APD as a single-photon detector in a QKD system, both in terms of secure communication rate and communication distance.
Finally, in Fig. 9 we compare the performance of quantum key distribution systems implementing the three protocols, under the assumptions that Eve is equipped with an ideal quantum memory and that the dark counts of the up-conversion detector, caused by parasitic non-linear processes in the PPLN waveguide, are eliminated. This means that the detector’s performance is ideally limited by the Si APD characteristics, which corresponds to $`d_{\text{up}}=5\times 10^8`$. Operation at the maximum quantum efficiency regime is also assumed, i.e. $`\eta _{\text{up}}=0.46`$. We observe that, ultimately, 250 km of secure communication distance is possible with the DPSK protocol and an ideal single-photon source implementation of BB84, while BBM92 has the potential of extending this distance to 350 km with a deterministic entangled-photon source.
## V Conclusions
In this paper, we studied the main characteristics of two types of 1.55 $`\mu `$m single-photon detectors, the InGaAs/InP APD and the up-conversion detector, which combines frequency up-conversion in a PPLN waveguide and detection by a silicon APD. We presented the communication rate equations for the BB84 and the BBM92 QKD protocols, and we derived a corresponding equation for the DPSK protocol, developing a security analysis of this protocol against certain types of hybrid attacks. Based on these equations, we compared the performance of fiber-optic quantum key distribution systems employing the protocols under consideration, with realistic experimental parameters. In all cases, we found that a secure communication rate of two orders of magnitude higher than before is possible, while the use of the up-conversion detector enables quantum key distribution over communication distances longer by a factor of 2 than with an InGaAs/InP APD. Furthermore, the importance of the implemented protocol was illustrated, and the impact of Eve’s allowed capabilities was investigated. We concluded that the simple and efficient DPSK protocol allows for more than 200 km of secure communication distance with high communication rates, in the realistic case that Eve does not possess a quantum memory with an infinitely long coherence time, and the time delay parameter $`N`$ is greater than 1. The BBM92 protocol can extend this distance to 300 km with a reasonably high secure key generation rate. It is clear that improving the performance of the Si APDs with respect to their dead time and timing jitter and reducing the dark counts of the up-converter will extend the capabilities of fiber-optic QKD systems employing these protocols even further.
###### Acknowledgements.
The authors would like to thank Edo Waks for his helpful comments and suggestions. Financial support was provided by the MURI Center for Photonic Quantum Information Systems (ARO/ARDA DAAD19-03-1-0199), and the Quantum Entanglement Project, SORST, JST.
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# KEK Preprint 2005-101Belle Preprint 2006-5 Observation of 𝒃→𝒅𝜸 and Determination of |𝑽_{𝒕𝒅}/𝑽_{𝒕𝒔}|
The Belle Collaboration
(Feb. 14, 2006)
## Abstract
We report the observation of the flavor-changing neutral current process $`bd\gamma `$ using a sample of $`386\times 10^6`$ $`B`$ meson pairs accumulated by the Belle detector at the KEKB $`e^+e^{}`$ collider. We measure branching fractions for the exclusive modes $`B^{}\rho ^{}\gamma `$, $`\overline{B}{}_{}{}^{0}\rho ^0\gamma `$ and $`\overline{B}{}_{}{}^{0}\omega \gamma `$. Assuming that these three modes are related by isospin, we find $`(\overline{B}(\rho ,\omega )\gamma )=(1.32_{0.31}^{+0.34}(\mathrm{stat}.)_{0.09}^{+0.10}(\mathrm{syst}.))\times 10^6`$ with a significance of $`5.1\sigma `$. This result is used to determine the ratio of CKM matrix elements $`|V_{td}/V_{ts}|`$ to be $`0.199_{0.025}^{+0.026}(\mathrm{exp}.)_{0.015}^{+0.018}(\mathrm{theo}.)`$.
The $`bd\gamma `$ process, which proceeds via a loop diagram (Fig. 1(a)) in the Standard Model (SM), is suppressed with respect to $`bs\gamma `$ by the Cabibbo-Kobayashi-Maskawa (CKM) factor bib:ckm $`|V_{td}/V_{ts}|^20.04`$, with large uncertainty due to the lack of precise knowledge of $`|V_{td}|`$. The exclusive modes $`\overline{B}\rho \gamma `$ and $`\overline{B}{}_{}{}^{0}\omega \gamma `$ are presumably the easiest modes to search for; no evidence for the decays has been previously reported bib:belle-rhogam ; bib:babar-rhogam . The predicted branching fractions are $`(0.9\text{}2.7)\times 10^6`$ bib:ali-cdlu ; bib:bosch-buchalla based on the measured rate for the $`bs\gamma `$ process $`\overline{B}\overline{K}{}_{}{}^{}\gamma `$ and the $`|V_{td}/V_{ts}|^2`$ factor with corrections due to form factors, $`SU(3)`$ breaking effects, and, for the $`B^{}`$ decay, inclusion of an annihilation diagram (Fig. 1(b)). Measurement of these exclusive branching fractions allows one to determine the value of $`|V_{td}/V_{ts}|`$ in the context of the SM and to search for physics beyond the SM bib:rhogam-bsm . In this Letter, we report the observation of the $`bd\gamma `$ process using a sample of $`(386\pm 5)\times 10^6`$ $`B`$ meson pairs accumulated at the $`\mathrm{{\rm Y}}(4S)`$ resonance. With a larger data sample and an improved analysis procedure, the results supersede those of our previous publication bib:belle-rhogam .
The data are produced in $`e^+e^{}`$ annihilation at the KEKB energy-asymmetric (3.5 on 8 GeV) collider bib:kekb and collected with the Belle detector bib:belle-detector , which includes a silicon vertex detector (SVD), a central drift chamber (CDC), aerogel threshold Cherenkov counters (ACC), time-of-flight (TOF) scintillation counters, and an electromagnetic calorimeter (ECL) of CsI(Tl) crystals located inside a 1.5 T superconducting solenoid coil.
We reconstruct three signal modes, $`B^{}\rho ^{}\gamma `$, $`\overline{B}{}_{}{}^{0}\rho ^0\gamma `$ and $`\overline{B}{}_{}{}^{0}\omega \gamma `$, and two control samples, $`B^{}K^{}\gamma `$ and $`\overline{B}{}_{}{}^{0}\overline{K}{}_{}{}^{0}\gamma `$. Charge conjugate modes are implicitly included throughout this Letter. The following decay modes are used to reconstruct the intermediate states: $`\rho ^{}\pi ^{}\pi ^0`$, $`\rho ^0\pi ^+\pi ^{}`$, $`\omega \pi ^+\pi ^{}\pi ^0`$, $`K^{}K^{}\pi ^0`$, $`\overline{K}{}_{}{}^{0}K^{}\pi ^+`$, and $`\pi ^0\gamma \gamma `$.
Photon candidates are reconstructed from ECL energy clusters with a photon-like shape and no associated charged track. A photon in the barrel ECL ($`33^{}<\theta _\gamma <128^{}`$ in the laboratory frame polar angle) with a center-of-mass (c.m.) energy in the range $`1.8\text{ GeV}<E_\gamma <3.4\text{ GeV}`$ is selected as the primary photon candidate. To suppress backgrounds from $`\pi ^0/\eta \gamma \gamma `$ decays, we apply a veto algorithm based on likelihoods to be and not to be a $`\pi ^0/\eta `$. The likelihoods are calculated for every combination of the primary photon and another photon in the event using the energy of the other photon and the invariant mass of the pair. We also reject the primary photon candidate if the ratio of the energy in the central $`3\times 3`$ ECL cells to that in the central $`5\times 5`$ cells is less than 0.95.
Neutral pions are formed from photon pairs with invariant masses within $`\pm 16\text{ MeV}/c^2`$ ($`3\sigma `$) of the $`\pi ^0`$ mass. The photon momenta are then recalculated with a $`\pi ^0`$ mass constraint. We require the energy of each photon to be greater than 50 (100) MeV inside (outside) the barrel ECL. We also require the cosine of the angle between the two photons in the laboratory frame to be greater than 0.7; this requirement suppresses the copious combinatorial background with momenta below $`0.6\text{ GeV}/c`$.
Charged pions and kaons are selected from tracks in the CDC and SVD. Each track is required to have a transverse momentum greater than $`100\text{ MeV}/c`$ and a distance of closest approach to the interaction point of less than $`0.5\text{ cm}`$ in radius and $`\pm 3.0\text{ cm}`$ along the $`z`$-axis, which is parallel to the positron beam. We do not use a track to form the signal candidate if, when it is combined with an oppositely charged track, the resulting pair has an invariant mass within $`\pm 30\text{ MeV}/c^2`$ of the $`K_S^0`$ mass and a displaced vertex that is consistent with that of a $`K_S^0`$. We determine pion ($`_\pi `$) and kaon ($`_K`$) likelihoods from ACC, CDC and TOF information and form a likelihood ratio $`_\pi /(_\pi +_K)`$ to separate pions from kaons. The criteria for pions have efficiencies of $`83\%`$, $`81\%`$ and $`91\%`$ for $`\rho ^{}`$, $`\rho ^0`$ and $`\omega `$, respectively; the corresponding kaon misidentification rates are $`5.8\%`$, $`6.3\%`$ and $`8.4\%`$. For $`K^{}`$ candidates, we select kaons with an efficiency of $`90\%`$.
Invariant masses for the $`\rho `$, $`\omega `$ and $`K^{}`$ candidates are required to be within windows of $`\pm 150`$, $`\pm 30`$ and $`\pm 75\text{ MeV}/c^2`$, respectively, around their nominal values.
Candidate $`B`$ mesons are reconstructed by combining a $`\rho `$ or $`\omega `$ candidate with the primary photon and calculating two variables: the beam-energy constrained mass $`M_{\mathrm{bc}}=\sqrt{(E_{\mathrm{beam}}^{}/c^2)^2|\stackrel{}{p}{}_{B}{}^{}/c|^2}`$, and the energy difference $`\mathrm{\Delta }E=E_B^{}E_{\mathrm{beam}}^{}`$. Here, $`\stackrel{}{p}_B^{}`$ and $`E_B^{}`$ are the c.m. momentum and energy of the $`B`$ candidate, and $`E_{\mathrm{beam}}^{}`$ is the c.m. beam energy. To improve resolution, the magnitude of the photon momentum is replaced by $`(E_{\mathrm{beam}}^{}E_{\rho /\omega }^{})/c`$ when the momentum $`\stackrel{}{p}_B^{}`$ is calculated.
To optimize the event selection, we study Monte Carlo (MC) events in a signal box defined as $`5.273\mathrm{GeV}/c^2<M_{\mathrm{bc}}<5.285\mathrm{GeV}/c^2`$ and $`0.10\mathrm{GeV}<\mathrm{\Delta }E<0.08\mathrm{GeV}`$. We choose selection criteria to maximize $`N_S/\sqrt{N_B}`$, where $`N_S`$ and $`N_B`$ are the expected signal and the sum of the background yields.
The dominant background arises from continuum events ($`e^+e^{}q\overline{q}(\gamma )`$, $`q=u,d,s,c`$), where a random combination of a $`\rho `$ or $`\omega `$ candidate with a photon forms a $`B`$ candidate. We suppress this background using the following quantities: (1) $``$, a Fisher discriminant constructed from 16 modified Fox-Wolfram moments bib:belle-pi0pi0 ; bib:fox-wolfram and the scalar sum of the transverse momenta of all charged tracks and photons. (2) $`\mathrm{cos}\theta _B^{}`$, where $`\theta _B^{}`$ is the c.m. polar angle of the $`B`$ candidate direction: true $`B`$ mesons follow a $`1\mathrm{cos}^2\theta _B^{}`$ distribution, while candidates in the continuum background are almost uniformly distributed. (3) $`\mathrm{\Delta }z`$, the separation along the $`z`$-axis between the decay vertex of the candidate $`B`$ meson and the fitted vertex of the remaining tracks in the event. Discrimination is provided due to the displacement of the signal $`B`$ decay vertex from the other $`B`$, as tracks from continuum events typically have a common vertex. For each of the quantities $``$, $`\mathrm{cos}\theta _B^{}`$ and $`\mathrm{\Delta }z`$, we construct likelihood distributions for signal and continuum events. The $``$, $`\mathrm{cos}\theta _B^{}`$ and signal $`\mathrm{\Delta }z`$ distributions are determined from MC samples; the continuum $`\mathrm{\Delta }z`$ distribution is determined from the data sideband $`5.20\text{ GeV}/c^2<M_{\mathrm{bc}}<5.24\text{ GeV}/c^2`$, $`0.1\text{ GeV}<\mathrm{\Delta }E<0.5\text{ GeV}`$.
We form product likelihoods $`_s`$ and $`_c`$ for signal and continuum background, respectively, from the likelihood distributions for $``$, $`\mathrm{cos}\theta _B^{}`$ and (where available) $`\mathrm{\Delta }z`$. In addition, we use a tagging quality variable $`r`$ that indicates the level of confidence in the $`B`$-flavor determination as described in Ref. bib:hamlet . In the $`(r,)`$ plane defined by the tagging quality $`r`$ and the likelihood ratio $`=_s/(_s+_c)`$, signal tends to populate the edges at $`r=1`$ and $`=1`$, while continuum preferentially populates the edges at $`r=0`$ and $`=0`$. We divide the events into six bins of $`r`$ (two bins between 0 and 0.5, and four between 0.5 and 1) and determine the minimum $``$ requirement for each bin. In the $`\rho ^{}\gamma `$ mode, we also assign events to the bin $`0r<0.25`$ if the tagging-side flavor is the same as the signal-side. The signal efficiency is $`40\%`$, and $`95\%`$ of continuum background is rejected. For the $`K^{}\gamma `$ ($`\overline{K}{}_{}{}^{0}\gamma `$) mode we use the selection criteria for the $`\rho ^{}\gamma `$ ($`\rho ^0\gamma `$) mode.
We consider the following backgrounds from $`B`$ decays: $`\overline{B}\overline{K}{}_{}{}^{}\gamma `$, other $`BX_s\gamma `$ processes, decays with a $`\pi ^0/\eta `$ ($`B\rho \pi ^0`$, $`\omega \pi ^0`$, $`\rho \eta `$ and $`\omega \eta `$), other charmless hadronic $`B`$ decays, and $`bc`$ decay modes. We find the $`bc`$ background to be negligible. The $`\overline{B}\overline{K}{}_{}{}^{}\gamma `$ background can mimic the $`\overline{B}\rho \gamma `$ signal if the kaon from the $`K^{}`$ is misidentified as a pion. To suppress $`\overline{B}\overline{K}{}_{}{}^{}\gamma `$ events we calculate $`M_{K\pi }`$, where the kaon mass is assigned to one of the charged pion candidates, and reject the candidate if $`M_{K\pi }<0.95`$ ($`0.92`$)$`\text{ GeV}/c^2`$ for the $`\rho ^0\gamma `$ ($`\rho ^{}\gamma `$) mode. This requirement removes 82% (64%) of the $`K^{}\gamma `$ background while retaining 63% (87%) of the signal. The decay chain $`\overline{B}{}_{}{}^{0}\overline{K}{}_{}{}^{0}\gamma `$, $`\overline{K}{}_{}{}^{0}K_S^0\pi ^0`$, $`K_S^0\pi ^+\pi ^{}`$ has a small contribution to $`\overline{B}{}_{}{}^{0}\omega \gamma `$ due to the tail of the $`K^{}`$ Breit-Wigner lineshape. In addition, $`\overline{B}\overline{K}{}_{}{}^{}\gamma `$ and other $`BX_s\gamma `$ decays contribute to the background when the $`\rho `$ and $`\omega `$ candidates are formed from random combinations of particles.
Hadronic decays with a $`\pi ^0/\eta `$ can mimic the signal if a photon from the $`\pi ^0`$ or $`\eta \gamma \gamma `$ decay is soft and passes the $`\pi ^0/\eta `$ veto. To suppress this background, we reject the candidate if $`|\mathrm{cos}\theta _{\mathrm{hel}}|>0.75`$, $`0.70`$ and $`0.80`$ for the $`\rho ^{}\gamma `$, $`\rho ^0\gamma `$ and $`\omega \gamma `$ modes, respectively, where the helicity angle $`\theta _{\mathrm{hel}}`$ is the angle between the $`\pi ^{}`$ track (normal to the $`\omega `$ decay plane) and the $`B`$ momentum vector in the $`\rho `$ ($`\omega `$) rest frame (similarly for the $`K^{}\gamma `$ modes). Other hadronic decays make smaller contributions.
The reconstruction efficiency for each mode is defined as the fraction of the signal remaining after all selection criteria are applied, where the signal yield is determined from a fit to the sum of the signal and continuum MC samples using the procedure described below. The total efficiencies are listed in Table 1. The systematic error on the efficiency is the quadratic sum of the following contributions, estimated using control samples: the uncertainty in the photon detection efficiency (2.2%) as measured in radiative Bhabha events; charged tracking efficiency (1.0% per track) from partially reconstructed $`D^+D^0\pi ^+`$, $`D^0K_S^0\pi ^+\pi ^{}`$, $`K_S^0\pi ^+(\pi ^{})`$; charged pion and kaon identification (0.7–1.7% per track) and misidentification (15–17%) from $`D^+D^0\pi ^+`$, $`D^0K^{}\pi ^+`$; neutral pion detection (4.6%) from $`\eta `$ decays to $`\gamma \gamma `$, $`\pi ^+\pi ^{}\pi ^0`$ and $`3\pi ^0`$; $``$-$`r`$ and $`\pi ^0/\eta `$ veto requirements (2.8–5.5%) from $`B^{}D^0\pi ^{}`$, $`D^0K^{}\pi ^+`$ and $`\overline{B}{}_{}{}^{0}D^+\pi ^{}`$, $`D^+K^{}\pi ^+\pi ^+`$; the $`\omega \pi ^+\pi ^{}\pi ^0`$ branching fraction (0.8%); and uncertainty due to MC statistics (0.5–0.7%).
We perform an unbinned extended maximum likelihood fit to candidates satisfying $`|\mathrm{\Delta }E|<0.5\text{ GeV}`$ and $`M_{\mathrm{bc}}>5.2\text{ GeV}/c^2`$, individually and simultaneously for the three signal modes. In the latter case we assume isospin symmetry, and we also simultaneously fit the two $`\overline{B}\overline{K}{}_{}{}^{}\gamma `$ modes. We describe the events in the fit region using a sum of functions for the signal, continuum, $`K^{}\gamma `$ (for the three signal modes only), and other background hypotheses. The signal distribution is modeled as the product of a Crystal Ball lineshape bib:cbls in $`\mathrm{\Delta }E`$ to reproduce the asymmetric ECL energy response, and a Gaussian (another Crystal Ball lineshape) in $`M_{\mathrm{bc}}`$ for the mode without (with) a $`\pi ^0`$ in the final state. The signal parameters for $`M_{\mathrm{bc}}`$ and $`\mathrm{\Delta }E`$ are determined from separate fits to the $`B^{}K^{}\gamma `$ and $`\overline{B}{}_{}{}^{0}\overline{K}{}_{}{}^{0}\gamma `$ samples for the modes with and without a neutral pion, respectively. The branching fraction is the only parameter that is allowed to float for the signal component. The continuum background component is modeled as the product of a linear function in $`\mathrm{\Delta }E`$ and an ARGUS function bib:argus-function in $`M_{\mathrm{bc}}`$. The continuum shape parameters and normalizations are mode dependent and allowed to float. We use the distributions of MC events to model the shapes of other background components. The size of the $`K^{}\gamma `$ background component in each signal mode is constrained using the fit to the $`K^{}\gamma `$ events and the known misidentification probability. Other radiative and charmless decays are considered as an additional background component when we extract the signal yield. The levels of the other backgrounds are fixed using known branching fractions or upper limits bib:hfag2005 . We constrain branching fractions in the simultaneous fit using the isospin relations bib:ali-1994 ; bib:ali-cdlu $`(\overline{B}(\rho ,\omega )\gamma )(B^{}\rho ^{}\gamma )=2\frac{\tau _{B^+}}{\tau _{B^0}}(\overline{B}{}_{}{}^{0}\rho ^0\gamma )=2\frac{\tau _{B^+}}{\tau _{B^0}}(\overline{B}{}_{}{}^{0}\omega \gamma )`$ and $`(\overline{B}\overline{K}{}_{}{}^{}\gamma )(B^{}K^{}\gamma )=\frac{\tau _{B^+}}{\tau _{B^0}}(\overline{B}{}_{}{}^{0}\overline{K}{}_{}{}^{0}\gamma )`$, where $`\frac{\tau _{B^+}}{\tau _{B^0}}=1.076\pm 0.008`$ bib:hfag2005 .
The results of the fits are shown in Fig. 2 and listed in Table 1. The simultaneous fit gives
$$(\overline{B}(\rho ,\omega )\gamma )=(1.32_{0.31}^{+0.34}{}_{0.09}{}^{+0.10})\times 10^6,$$
(1)
where the first and second errors are statistical and systematic, respectively. The result is consistent with previous results bib:belle-rhogam ; bib:babar-rhogam and in agreement with SM predictions bib:ali-cdlu ; bib:bosch-buchalla . The significance of the simultaneous fit is $`5.1\sigma `$, where the significance is defined as $`\sqrt{2\mathrm{ln}(_0/_{\mathrm{max}})}`$, and $`_{\mathrm{max}}`$ ($`_0`$) is the value of the likelihood function when the signal branching fraction is floated (set to zero). Here, the likelihood function from the fit is convolved with a Gaussian systematic error function in order to include the systematic uncertainty. The invariant $`\pi \pi (\pi )`$ mass and helicity angle distributions for the events in the signal box are consistent with those expected from the sum of the signal and background components. The fit also gives $`(\overline{B}\overline{K}{}_{}{}^{}\gamma )=(41.1_{1.3}^{+1.4})\times 10^6`$ (statistical error only), which is consistent with the world average value bib:hfag2005 . The individual fit results are in marginal agreement with the isospin relation. We test our fitting procedure using MC simulation and find no statistically significant bias. We perform MC pseudo-experiments where events are generated according to the isospin relation; from the two-dimensional distribution of the deviation between the $`B^{}`$ and averaged $`\overline{B}^0`$ rates and that between the $`\rho ^0\gamma `$ and $`\omega \gamma `$ rates, we find the probability to observe an isospin violation equal to or larger than our measurement to be 4.9%. The expected level of isospin violation is within $`\pm 10\%`$ bib:ali-cdlu .
The systematic error is estimated by varying each of the fixed parameters by $`\pm 1\sigma `$ and then taking the quadratic sum of the deviations in the branching fraction from the nominal value. We note that the ARGUS background shape in the fit to the $`\omega \gamma `$ mode is steeper than those for the other two modes. Therefore we also vary the ARGUS shape parameter for the $`\omega \gamma `$ mode by $`2\sigma `$ and include the deviation in the systematic error.
The ratio $`(\overline{B}(\rho ,\omega )\gamma )/(\overline{B}\overline{K}{}_{}{}^{}\gamma )=0.032\pm 0.008(\mathrm{stat}.)\pm 0.002(\mathrm{syst}.)`$, which we obtain from a separate fit, can be used to determine $`|V_{td}/V_{ts}|`$. The fit takes into account the correlated systematic errors between the signal and $`K^{}\gamma `$ modes and thus gives a reduced total error. Using the relation bib:ali-2004 $`\frac{(\overline{B}(\rho ,\omega )\gamma )}{(\overline{B}\overline{K}{}_{}{}^{}\gamma )}=\left|\frac{V_{td}}{V_{ts}}\right|^2\frac{(1m_{(\rho ,\omega )}^2/m_B^2)^3}{(1m_K^{}^2/m_B^2)^3}\zeta ^2[1+\mathrm{\Delta }R]`$, where the form factor ratio $`\zeta =0.85\pm 0.10`$ and the $`SU(3)`$-breaking correction $`\mathrm{\Delta }R=0.1\pm 0.1`$, we obtain
$$|V_{td}/V_{ts}|=0.199_{0.025}^{+0.026}(\mathrm{exp}.)_{0.015}^{+0.018}(\mathrm{theo}.).$$
(2)
We obtain a 95% confidence level interval of $`0.142<|V_{td}/V_{ts}|<0.259`$ using an ensemble of MC samples in which the experimental error is a quadratic sum of the asymmetric Gaussian statistical and systematic errors, and the theory error is a flat distribution in the given range. This result is in agreement with the range favored by a fit to the unitarity triangle bib:pdg2004 assuming $`|V_{ts}|=|V_{cb}|`$.
In conclusion, we observe the process $`bd\gamma `$ using the $`B\rho \gamma `$ and $`\omega \gamma `$ modes. The resulting branching fractions are consistent with SM predictions bib:ali-cdlu ; bib:bosch-buchalla . The ratio of the $`\overline{B}(\rho ,\omega )\gamma `$ branching fraction to that for $`\overline{B}\overline{K}{}_{}{}^{}\gamma `$ is used to determine $`|V_{td}/V_{ts}|`$.
We thank the KEKB group for the excellent operation of the accelerator, the KEK cryogenics group for the efficient operation of the solenoid, and the KEK computer group and the NII for valuable computing and Super-SINET network support. We acknowledge support from MEXT and JSPS (Japan); ARC and DEST (Australia); NSFC (contract No. 10175071, China); DST (India); the BK21 program of MOEHRD and the CHEP SRC program of KOSEF (Korea); KBN (contract No. 2P03B 01324, Poland); MIST (Russia); MHEST (Slovenia); SNSF (Switzerland); NSC and MOE (Taiwan); and DOE (USA).
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# Dynamic inverse problem in a weakly laterally inhomogeneous medium.
## 1. Introduction
The inverse problems of geo-exploration implies investigation of domains inside the earth’s crust which contain gas, oil or other minerals as well as determination of their physical properties such as density, porosity, pressure, and so on. The most important problem is determination of the boundaries of the domains which gives information about its location, volume and makes possible to predict costs and outcomes in exploitation. In practice, in e.g. oil exploration and seismology this inverse problem is mainly solved by means of the map migration method. The map migration method assumes the knowledge of the velocity profile above the interface and provides robust, effective numerical algorithms which are quite stable with respect to variations of the velocity and complicated geometry of the interface. Therefore, it has attracted much attention of mathematicians and geophysicists, see e.g. , , , , with new approaches continuing to appear e.g. . However, an absence of an independent way to recover the velocity profile above the interface may hinder the map migration techniques. This makes crucial for geosciences to develop algorithms to solve an inverse problem of the velocity reconstruction from the measurements of the wave field on the ground surface, $`z=0`$.
This gives rize to a fully non-linear dynamic multidimensional inverse problem. To our knowledge, the only method valid for arbitrary inhomogeneous velocity profiles is the Boundary Control method (BC-method), see e.g. , . However, at the moment, this method is developed only for smooth velocity profiles, in the absence of discontinuities. Moreover, existing variants of the BC-method do not possess good stability properties and have a rather poor numerical implementation, e.g. .
On the other hand, in the case of a pure layered structure, with the properties depending only on depth, $`z`$, inverse problems of geophysics are often reduced to one-dimensional inverse problems. The theory of one-dimensional inverse problems which goes back to the classical results obtained by Borg, Levinson, Gel’fand-Levitan, Marchenko and Krein of late 40-th - 50-th is still an area of active research with new theoretical results and numerical algorithms appearing regularly in mathematical and geophysical literature, see e.g. , , , , , to mention just a few.
Returning to geophysical applications, a pure layered structure occurs rare. However, in many cases in seismology and oil exploration the parameters of the earth are not described by arbitrary functions of 3 dimensions having jumps across arbitrary 2 dimensional surfaces. Rather the properties of the medium depend mainly on the depth, $`z`$, with only slow dependence on horizontal coordinates, $`x=(x_1,x_2)`$, i.e. depend on $`ϵx,ϵ<<1`$ rather than $`x`$ (in seismology, often $`ϵ0.10.2`$) and we deal with a weakly laterally inhomogeneous medium (WLIM). The importance of WLIM is now well-understood in theoretical and mathematical geophysics. There are currently numerous results on the direct problem of the wave propagation in WLIM, see e.g. , , , , , . However, to our knowledge, except , there is currently no special inversion techniques for WLIM. In this paper we develop and numerically implement an inversion method for WLIM based on different ideas than those in . This method provides, for inverse problems occurring in important practical applications, a technique which inherits some practically useful properties of the one-dimensional inverse problems, e.g. robustness and fast convergence rate of numerical algorithms.
In the acoustical approximation we use in this paper, the wave propagation in WLIM is described by the wave equation,
(1)
$$u_{tt}c^2(z,x)\mathrm{\Delta }_{z,x}u=0,z>0$$
with
(2)
$$c(z,x)=c_0(z)+ϵ<x,c_1(z)>+\mathrm{}.$$
Due to (2), the response of the media, measured at $`z=0`$, may also be decomposed in power series with respect to $`ϵ`$. Analyzing this decomposition, we obtain a recurrent system of inverse problems for subsequent terms $`c_p(z)`$. The inverse problem for $`c_0(z)`$ reduces to a family of one-dimensional inverse problems. We employ as a basis to solve this problem the method coupled non-linear Volterra integral equations developed by Blagovestchenskii, e.g. , which goes back to . Note that , provide a local reconstruction of $`c_0`$ near the ground surface $`z=0`$. In this paper, we pursue the method of coupled non-linear Volterra integral equations further making it global, i.e. providing reconstruction of $`c_0(y)`$ from inverse data measured for $`t(0,2T)`$ up to the depth $`T`$ in the travel-time coordinates, $`y`$ (see formula (16) and obtain conditional Lipschitz stability estimates for this procedure. The results of the zero-order reconstruction serve as a starting point to develop a recurrent procedure to determine, from the inverse data, further terms in expansion (2). Such procedure, for the inverse problem of the reconstruction of an unknown potential $`q(z,ϵx)`$ rather than $`c(z,ϵx)`$,was suggested by Blagovestchenskii , . The method of is based on the use of polynomial moments with respect to $`x`$. We suggest another approach based on the Fourier transform, $`x=(x_1,x_2)\xi =(\xi _1,\xi _2)`$, of the equation (1). This makes possible to utilize the dependence of the resulting problems on $`\xi `$ in order to solve the recurrent system of 1D inverse problems. The coefficients $`c_1(z),c_2(z),\mathrm{}`$ may then be obtained as solutions to some linear Volterra integral equations. Having said so, we note that the integral equations for higher order unknowns contain higher and higher order derivatives of the previously found terms, thus increasing ill-posedness of the inverse problem. This is hardly surprising taking into account well-known strong ill-posedness of multi-dimensional inverse problems , . What is, however, interesting is that within the model considered there is just a gradual increase of instability adding two derivations at each stage of the reconstruction algorithm.
In this paper we confine ourself to the reconstruction only of $`c_0,c_1`$. Reconstruction of the higher order terms $`c_p,p>1`$ is, in principle, possible using the same ideas as for $`c_1`$, although is more technically involved. However, in practical applications in geophysics the measured data make possible to find inverse data only for $`c_0,c_1`$. Indeed, (1), (2) imply that, for the type of the boundary sources we consider,
$$u(z,x,t)=u_0(z,x,t)+ϵu_1(z,x,t)+ϵ^2u_2(z,x,t)+\mathrm{},$$
so that the measured data at $`z=0`$,
$$u_z|_{z=0}(x,t)=R(x,t)=R_0(x,t)+ϵR_1(x,t)+ϵ^2R_2(x,t)+\mathrm{}.$$
What is more, $`R_p(x,t)`$ are even with respect to $`x`$ for even $`p`$ and $`R_p(x,t)`$ are odd for odd $`p`$. Clearly, in real measurements $`R`$ is not given as a power series with respect to $`ϵ`$. However, using the fact that $`R_0`$ is even, and $`R_1`$ is odd, we can find $`R_0`$ and $`R_1`$ up to $`O(ϵ^2)`$ from the measured $`R`$.
The plan of the paper is as follows. In the next section we give a rigorous formulation of the problem and provide a general outline of the perturbation scheme for this problem in WLIM. In section 3 a modified method of coupled non-linear Volterra integral equations to reconstruct $`c_0(z)`$ is described. In section 4 we derive a coupled system of linear Volterra integral equations for $`c_1(z)`$. Section 5 is devoted to the global solvability of the non-linear system for $`c_0(z)`$ and conditional stability estimates. In section 6 we test the method numerically for the two dimensional case (inhomogeneous half-plane). As we stop with $`c_0,c_1`$, in practical applications this would result in an error of the magnitude $`O(ϵ^2)`$. The final section is devoted to some concluding remarks.
## 2. Formulation of the problem
Let us consider the wave propagation into the inhomogeneous acoustic half-space
(3)
$$n^2(z,ϵx)u_{tt}\frac{1}{\rho }\mathrm{div}(\rho u)=0,x=(x_1,x_2),$$
where $`n`$ is the refractive index, $`n=c^1`$. Above and below the interface $`z=h(ϵx)`$ the refractive index is described by
$$n(z,ϵx)=\{\begin{array}{ccc}\hfill n^{(1)}(z,ϵx),& 0<z<h(ϵx),& \\ \hfill n^{(2)}(z,ϵx),& z>h(ϵx),& \end{array},$$
where $`ϵ`$ is a small parameter ($`0<ϵ1`$) which characterizes the ratio of the horizontal and vertical gradients of $`n`$. We assume that the density is piece-wise constant,
$$\rho =\{\begin{array}{ccc}\hfill \rho _1,& 0<z<h(ϵx),& \\ \hfill \rho _2,& z>h(ϵx),& \end{array}$$
with known $`\rho _1,\rho _2`$. As we deal with WLIM, the following notations are used,
(4)
$$n^2(z,ϵx)=n_0^2(z)+ϵ<x,\overline{n}(z)>+O(ϵ^2),\overline{n}(z)=(n_1(z),n_2(z)),$$
(5)
$$h(ϵx)=h_0+ϵ<x,\overline{h}>+O(ϵ^2),\overline{h}=(h_1,h_2),$$
where $`<,>`$ means the scalar product.
We assume that $`u=0`$ for $`t0`$, and the boundary condition is
(6)
$$u|_{z=0}=\delta (x)f(t)\sqrt{\frac{\rho _1}{n_0(0)}},f(t)=\delta (t)\mathrm{or}f(t)=\theta (t),$$
where $`\delta (x)`$ and $`\theta (t)`$ are $`\delta `$-function and Heaviside function, correspondingly. On the interface between these two layers we assume the usual continuity conditions,
(7)
$$[u]|_{z=h(ϵx)}=0,[\rho ^1\frac{u}{n}]|_{z=h(ϵx)}=0,$$
where $`[\mathrm{}]`$ stand for a jump across the interface. Using the Fourier transform with respect to $`x`$
(8)
$$u(z,x,t)=\frac{1}{4\pi ^2}\underset{R^2}{}e^{i<\xi ,x>}U(z,\xi ,t)𝑑\xi ,\xi =(\xi _1,\xi _2),.$$
$`U`$ can be expanded into asymptotic series,
(9)
$$U(z,\xi ,t)\underset{n=0}{\overset{\mathrm{}}{}}ϵ^ni^nU^{(n)}(z,\xi ,t),$$
where all functions $`U^{(n)}(z,\xi ,t)`$ are real. Using decompositions (4), (9), it is easily seen that $`U^{(n)}(z,\xi ,t)`$ are even functions with respect to $`\xi `$ for even $`n`$ and odd for odd $`n`$.
Our goal is to reconstruct the refractive index and the shape of interface, namely, the functions $`n_0(z),n_{1,2}(z)`$ and constants $`h_0,h_{1,2}`$, from the response data collected during time $`0<t<2T`$,
$$\frac{u}{z}|_{z=0}=R(x,t,ϵ),$$
with
$$R(x,t,ϵ)\underset{n=0}{\overset{\mathrm{}}{}}ϵ^nR_n(x,t).$$
Decomposing the wave equation (1) and interface conditions (7) with respect to $`ϵ`$, we obtain initial-boundary value problems for $`U^{(0)}`$ and $`U^{(1)}`$. The zero-order problem is
(10)
$$n_0^2(z)U_{tt}^{(0)}U_{zz}^{(0)}+|\xi |^2U^{(0)}=0,|\xi |^2=\xi _1^2+\xi _2^2,U^{(0)}|_{z=0}=f(t)\sqrt{\frac{\rho _1}{n_0(0)}}$$
with the interface continuity conditions
(11)
$$[U^{(0)}]|_{z=h_0}=0,[\frac{1}{\rho }U_z^{(0)}]|_{z=h_0}=0.$$
and inverse data of the form
(12)
$$U_z^{(0)}|_{z=0}=r_0(t,\xi )=\underset{R^2}{}\mathrm{cos}(<\xi ,x>)R_0(x,t)dx.$$
The first-order problem is
(13)
$$n_0^2(z)U_{tt}^{(1)}U_{zz}^{(1)}+|\xi |^2U^{(1)}=<\overline{n},_\xi U_{tt}^{(0)}>,U^{(1)}|_{z=0}=0.$$
with the interface continuity conditions
(14)
$$[U^{(1)}<\overline{h},_\xi U_z^{(0)}>]|_{z=h_0}=0,\left[\frac{1}{\rho }\{U_z^{(1)}<\overline{h},_\xi U_{zz}^{(0)}>+U^{(0)}<\xi ,\overline{h}>\}\right]|_{z=h_0}=0$$
and inverse data of the form
(15)
$$U_z^{(1)}|_{z=0}=r_1(t,\xi )=\underset{R^2}{}\mathrm{sin}(<\xi ,x>)R_1(x,t)dx.$$
In this paper we confine ourselves to the reconstruction of only $`n_0,\overline{n}`$ and $`h_0,\overline{h}`$. The reconstruction of the higher order terms is, in principal, possible using the same ideas as for $`\overline{n}`$ and $`\overline{h}`$, although is more technically involved. Moreover, in practical applications in geophysics the measured data make possible to find the inverse data only for $`n_0,\overline{n}`$ and $`h_0,\overline{h}`$. Using the experimental data of the response $`R(x,t)`$, one cannot expand it into power series with respect to $`ϵ`$. However, as $`R_0`$ is even and $`R_1`$ is odd with respect to $`x`$, we have
$$r_0(t,\xi )=\underset{R^2}{}\mathrm{cos}(<\xi ,x>)R(x,t)dx+O(ϵ^2),r_1(t,\xi )=ϵ^1\underset{R^2}{}\mathrm{sin}(<\xi ,x>)R(x,t)dx+O(ϵ^2).$$
So, the only quantities we may observe from the experiment are $`r_0(\xi ,t)`$ $`r_1(\xi ,t)`$ up to error $`ϵ^2`$. Thus, although the inverse problems for the higher-order terms in (4) can be considered analytically, inverse data for them are not available from the measurements. In this connection, we do not discuss higher approximations in this paper. Also the exact value for the parameter $`ϵ`$ is not known from the experiment. However, this quantity may be evaluated numerically using inverse data of response $`R(x,t,ϵ)`$. For example, one of the ways to calculate $`ϵ`$ is as follows
$$ϵ=sup\left|\underset{R^2}{}\mathrm{sin}(<\xi ,x>)R(x,t)dx\right|,0t2T,|\xi |\xi _{max},$$
where $`\xi _{max}`$ is a positive real number of order 1. It is clear that $`ϵ`$ is qualitatively related to an extent of the medium being a weakly laterally inhomogeneous one.
## 3. Inverse problem in the zero-order approximation
### 3.1. Half-space
In this section we describe an algorithm to solve the inverse problem in the zero-order approximation for an inhomogeneous half-space. Namely, we will describe an algorithm to determine $`n_0(z)`$ from the response $`r_0(t,\xi )`$. This is a generalization of the approach developed first by Blagovestchenskii , see also . The crucial point of the approach is the derivation of a non-linear Volterra-type system of integral equations to solve the inverse problem.
Let us introduce two new independent variables
(16)
$$y=\underset{0}{\overset{z}{}}n_0(z)𝑑z,\sigma (y)=\frac{n_0(z(y))}{\rho }.$$
The function $`\sigma (y)`$ is called acoustic admittance while $`y`$ is the vertical travel-time. Then,
(17)
$$U_{tt}^{(0)}\frac{1}{\sigma (y)}\frac{}{y}\left(\sigma (y)U_y^{(0)}\right)+\frac{|\xi |^2}{n_0^2}U^{(0)}=0,$$
Let $`f(t)=\delta (t)`$. Then, the boundary condition and response are
$$U^{(0)}|_{y=0}=\delta (t)\sqrt{\frac{\rho _1}{n_0(0)}},U_y^{(0)}|_{y=0}=\frac{r_0(\xi ,t)}{n_0(0)}.$$
Let us change the dependent variable
(18)
$$\psi _1(y,t)=\sqrt{\sigma (y)}U^{(0)},\psi _2(y,t)=\frac{\psi _1}{t}+\frac{\psi _1}{y},$$
thus reducing the second order PDE for $`U^{(0)}`$ to a system of two PDE of the first order
(19)
$$\{\begin{array}{ccc}\hfill \psi _{1t}+\psi _{1y}& =& \psi _2\hfill \\ \hfill \psi _{2t}\psi _{2y}& =& q(y)\psi _1,\hfill \end{array}$$
where
(20)
$$q(y,\xi )=\frac{(\sqrt{\sigma (y)})^{\prime \prime }}{\sqrt{\sigma (y)}}\frac{|\xi |^2}{n_0^2}.$$
Expressions in the left-hand side of (19) are the total derivatives of $`\psi _1`$ and $`\psi _2`$ along corresponding characteristics (see Fig.1). Integrating the first equation along the lower characteristics, and the second equation - along the upper one, we obtain for $`t>y`$
(21)
$$\{\begin{array}{ccc}\hfill \psi _1(y,t)& =& \underset{0}{\overset{y}{}}\psi _2(\eta ,t+\eta y)𝑑\eta ,\hfill \\ \hfill \psi _2(y,t)& =& \underset{0}{\overset{y}{}}q(\eta ,\xi )\psi _1(\eta ,t\eta +y)𝑑\eta +g(t+y,\xi ),\hfill \end{array}$$
where
(22)
$$g(t,\xi )=\delta ^{}(t)+\frac{n_0^{}(0)}{2n_0(0)}\delta (t)+\frac{r_0(t,\xi )}{\sqrt{\rho _1n_0(0)}}$$
is bounded as $`t0`$ due to the cancellation of $`\delta ^{}(t)`$ and $`\delta (t)`$ singularities in the right-hand side of (22).
The system (21) is not complete to solve the inverse problem as it has three unknown functions $`\psi _1(y,t)`$, $`\psi _2(y,t)`$ and $`q(y,\xi )`$ and just two equations. We use the progressive wave expansion, see e.g. , ,
(23)
$$U^{(0)}(y,t)=\underset{n=0}{}f_n(ty)U_n^{(0)}(y),f_{n+1}(t)=\underset{0}{\overset{t}{}}f_n(t)𝑑t,f_0(t)=\delta (t)$$
to derive the third equation. The amplitudes $`U_n^{(0)}`$ are found from the transport equations with e.g.
(24)
$$U_0^{(0)}=\frac{1}{\sqrt{\sigma (y)}},U_1^{(0)}=\frac{1}{2\sqrt{\sigma (y)}}\underset{0}{\overset{y}{}}q(\eta )𝑑\eta ,$$
so that
$$\psi _2(y,t)=\frac{1}{2}\theta (ty)q(y)+\mathrm{}.$$
Here and later we often skip the dependence of $`q`$ and other functions on $`\xi `$ when our considerations are independent of its value.
Substituting the second equation of (21) into this relation, and taking $`ty+0`$, we obtain the third desired equation. We summarize the above results in the follwing
###### Proposition 3.1.
Let $`\psi _{1,2}(y,t)`$ is obtained from the solution $`U^{(0)}(y,t)`$ of the wave equation (17) by formula (18). Then $`\psi _{1,2}(y,t)`$ together with $`q(y)`$ satisfy the following system of non-linear Volterra-type equations,
(25)
$$\{\begin{array}{ccc}\hfill \psi _1(y,t)& =& \underset{0}{\overset{y}{}}\psi _2(\eta ,t+\eta y)𝑑\eta ,\hfill \\ \hfill \psi _2(y,t)& =& \underset{0}{\overset{y}{}}q(\eta )\psi _1(\eta ,t\eta +y)𝑑\eta +g(t+y),\hfill \\ \hfill q(y)& =& 2\underset{0}{\overset{y}{}}q(\eta )\psi _1(\eta ,2y\eta )𝑑\eta +2g(2y).\hfill \end{array}$$
This non-linear Volterra-type system of integral equations, in the sense of Goh’berg-Krein , allows us to determine, at least locally, the function $`q(y)`$ from the response $`g(t)`$ for $`0<t<2T`$. In section 5 we return to the question of the global solvability of system (25). Having solved this system for two different values $`\xi _1`$ and $`\xi _2`$, we determine the refractive index in the zero-approximation,
(26)
$$n_0(y)=\sqrt{\frac{|\xi |_2^2|\xi |_1^2}{q(y,\xi _1)q(y,\xi _2)}},$$
which must be used together with
$$z=\underset{0}{\overset{y}{}}\frac{dy}{n_0(y)}.$$
The refractive index $`n_0(y)`$ may also be determined by solving system (25) just for a single value of $`\xi `$. It can be done by integration the differential equation for $`w(y)=\sqrt{\sigma (y)}`$ which follows from (20)
$$w^{\prime \prime }+wq(y,\xi )=\frac{|\xi |^2}{\rho _1^2w^3}.$$
The general solution of this equation may be represented in the form, see .
$$w(y)=\sqrt{\underset{i,j=1,2}{}A_{i,j}w_i(y)w_j(y)},$$
where $`w_{1,2}(y)`$ make a fundamental system of the corresponding homogeneous equation. Constant real symmetric matrix $`A_{i,j}`$ satisfies
$$\mathrm{det}A_{i,j}W^2(w_1,w_2)=\frac{|\xi |^2}{\rho _1^2},$$
where
$$W(w_1,w_2)=w_1w_2^{}w_1^{}w_2.$$
Given $`w(0)`$ and $`w^{}(0)`$ we are then able to determine all entries of the matrix $`A`$.
### 3.2. Two layers problem
Consider a wave propagation through the interface, $`y=y(h_0)=L`$, using the singularity analysis of the incident, reflected and transmitted waves. Then, in the layer $`0<y<L`$, the singularities of $`U^{(0)}`$ given by the incident and reflected waves, are
$`U^{(0)}={\displaystyle \frac{\delta (ty)}{\sqrt{\sigma (y)}}}+d_0{\displaystyle \frac{\delta (t+y2L)}{\sqrt{\sigma (y)}}}+`$
$`+{\displaystyle \frac{\theta (ty)}{2\sqrt{\sigma (y)}}}Q_0(y)+{\displaystyle \frac{\theta (t+y2L)}{2\sqrt{\sigma (y)}}}(d_1+Q_1(y))+\mathrm{},0<y<L,L<t<2L.`$
For the transmitted wave, the WKB asymptotics is given by
$$U^{(0)}=s_0\frac{\delta (ty)}{\sqrt{\sigma (y)}}+\frac{\theta (ty)}{2\sqrt{\sigma (y)}}(s_1+Q_1(y))+\mathrm{},y>L,L<t<2L,$$
where
$$Q_0(y)=\underset{0}{\overset{y}{}}q(\eta )𝑑\eta ,Q_1(y)=\underset{L}{\overset{y}{}}q(\eta )𝑑\eta .$$
Unknown quatities $`d_0`$, $`d_1`$ and $`s_0`$, $`s_1`$ are the reflection and transmission coefficients. Using the interface continuity conditions (11), leads to a linear system of equations for $`d_0`$, $`d_1`$ and $`s_0`$, $`s_1`$ so that
(27) $`d_0={\displaystyle \frac{1\frac{\sigma _+}{\sigma _{}}}{1+\frac{\sigma _+}{\sigma _{}}}},d_1={\displaystyle \frac{\sigma _{}^{}(1+d_0)+Q_0(L)(\sigma _{}\sigma _+)s_0\sigma _+^{}\sqrt{\frac{\sigma _{}}{\sigma _+}}}{\sigma _++\sigma _{}}},`$
$`s_0=\sqrt{{\displaystyle \frac{\sigma _+}{\sigma _{}}}}{\displaystyle \frac{2}{1+\frac{\sigma _+}{\sigma _{}}}},s_1=\sqrt{{\displaystyle \frac{\sigma _+}{\sigma _{}}}}{\displaystyle \frac{\sigma _{}^{}(1+d_0)+2Q_0(L)\sigma _{}s_0\sigma _+^{}\sqrt{\frac{\sigma _{}}{\sigma _+}}}{\sigma _++\sigma _{}}},`$
where $`\sigma _\pm `$ and $`\sigma _\pm ^{}`$ are the limiting values of $`\sigma (y)`$ and $`\sigma ^{}(y)`$ at $`yL0`$ and $`yL+0`$, correspondingly.
Analyzing the incoming singularities at $`y=0`$, we can find $`L`$ and, therefore, $`h0`$ and also $`d_0`$ and $`d_1`$. This will give us $`\sigma _+`$ and $`\sigma _+^{}`$. The next step is reconstruction of the velocity at large depth $`L<y<T`$, which is illustrated on Fig. 2. This reconstruction requires the data observed for $`0<t<2T`$. At this stage we again apply the non-linear Volterra system of integral equations (25) in the new frame $`y^{},\eta ^{},t^{}`$ ($`y=y^{}+L,\eta =\eta ^{}+L,t=t^{}+L`$)
(28)
$$\{\begin{array}{ccc}\hfill \psi _1(y^{},t^{})& =& \underset{0}{\overset{y^{}}{}}\psi _2(\eta ^{},t^{}+\eta ^{}y^{})𝑑\eta ^{}+\psi _1^+(0,t^{}y^{}),\hfill \\ \hfill \psi _2(y^{},t^{})& =& \underset{0}{\overset{y^{}}{}}q(\eta ^{})\psi _1(\eta ^{},t^{}\eta ^{}+y^{})𝑑\eta ^{}+\psi _2^+(0,t^{}+y^{}),\hfill \\ \hfill q(y^{})& =& 2\underset{0}{\overset{y^{}}{}}q(\eta ^{})\psi _1(\eta ^{},2y^{}\eta ^{})𝑑\eta ^{}+2\psi _2^+(0,2y^{}).\hfill \end{array}$$
These require knowledge of $`\psi _1^+(0,t^{}y^{})`$ and $`\psi _2^+(0,t^{}+y^{})`$ at the vertical segment $`AB`$ as the limits $`\psi _{1,2}^+(0,t)=lim\psi _{1,2}(y,t)`$ for $`yL+0`$. To this end, we first determine $`\psi _{1,2}^{}(0,t)=lim\psi _{1,2}(y,t)`$ for $`yL0`$ using system (21). Employing the interface continuity conditions
$$\frac{\psi _1^{}}{\sqrt{\sigma _{}}}=\frac{\psi _1^+}{\sqrt{\sigma _+}},\psi _{1y}^{}\sqrt{\sigma _{}}\psi _1^{}\frac{\sigma _{}^{}}{2\sqrt{\sigma _{}}}=\psi _{1y}^+\sqrt{\sigma _+}\psi _1^+\frac{\sigma _+^{}}{2\sqrt{\sigma _+}},$$
we obtain the required data, $`\psi _1^+(0,t^{}y^{})`$ and $`\psi _2^+(0,t^{}+y^{})`$ at $`AB`$.
## 4. Inverse problem in the first-order approximation
### 4.1. Half-space
In this section we describe an algorithm to determine $`\overline{n}(z)`$ in an inhomogeneous half-space from the response $`r_1(t,\xi )`$. For convenience, we integrate all the functions of the zero-order problem with respect to time $`t`$ so that
$$U^{(0)}|_{y=0}=\theta (t)\sqrt{\frac{\rho _1}{n_0(0)}}.$$
Using the variables $`y`$ and $`\sigma (y)`$, we rewrite the problem (13), (14) in the form
(29)
$$U_{tt}^{(1)}\frac{1}{\sigma (y)}\frac{}{y}\left(\sigma (y)U_y^{(1)}\right)+\frac{|\xi |^2}{n_0^2}U^{(1)}=\frac{1}{n_0^2}<\overline{n},_\xi U_{tt}^{(0)}>,$$
$$U^{(1)}|_{y=0}=0,U_y^{(1)}|_{y=0}=\frac{\widehat{r}_1(t,\xi )}{n_0(0)},\widehat{r}_1(t,\xi )=\underset{0}{\overset{t}{}}r_1(t,\xi )𝑑t.$$
We start with the progressive wave expansion for $`U^{(1)}`$. As by (23), (24),
$$<\overline{n},_\xi U_{tt}^{(0)}>=<\overline{n},\xi >\frac{\delta (ty)}{\sqrt{\sigma (y)}}\underset{0}{\overset{y}{}}\frac{d\eta }{n_0^2(\eta )}+\mathrm{}.$$
the progressive wave expansion for $`U^{(1)}(y,t)`$ has the form
$$U^{(1)}(y,t)=\theta (ty)A_0(y)+(ty)_+A_1(y)+\mathrm{}.$$
Therefore, equation (29) implies, in particular, that
(30)
$$A_0(y)=\frac{1}{\sqrt{\sigma (y)}}\underset{0}{\overset{y}{}}<\overline{n}(\eta ),\xi >p(\eta )d\eta ,p(y)=\frac{1}{2n_0^2(y)}\underset{0}{\overset{y}{}}\frac{d\eta }{n_0^2(\eta )}.$$
Let $`G(y,\eta ,t;\xi )`$ be the space-causal Green’s function,
(31)
$$G_{tt}\frac{1}{\sigma (y)}\frac{}{y}\left(\sigma (y)G_y\right)+\frac{|\xi |^2}{n_0^2}G=\delta (t)\delta (y\eta ),G=0\mathrm{if}y<\eta .$$
Then
(32) $`U^{(1)}(y,y+0;\xi )=`$
$`{\displaystyle \underset{0}{\overset{y}{}}}𝑑\eta {\displaystyle \underset{\eta }{\overset{2y\eta }{}}}𝑑\tau G(y,\eta ,y\tau ;\xi ){\displaystyle \frac{<\overline{n},_\xi U_{tt}^{(0)}(\eta ,\tau ;\xi )>}{n_0^2(\eta )}}{\displaystyle \frac{1}{n_0(0)}}{\displaystyle \underset{0}{\overset{2y}{}}}𝑑\tau G(y,0,y\tau ;\xi )\widehat{r}_1(\tau ,\xi ),`$
where we now write down explicitly the dependence on $`\xi `$.
Introducing new unknown functions,
(33)
$$\overline{\phi }(y)=(\phi _1(y),\phi _2(y))=\overline{n}(y)p(y)$$
and employing the progressive wave expansion for $`U^{(1)}`$, we see that
$$<\overline{\phi }(y),\xi >=\frac{d}{dy}\{\sqrt{\sigma (y)}\underset{0}{\overset{y}{}}d\eta \underset{\eta }{\overset{2y\eta }{}}d\tau G(y,\eta ,y\tau ,\xi )\frac{<\overline{\phi }(\eta ),_\xi U_{tt}^{(0)}(\eta ,\tau ,\xi )>}{n_0^2(\eta )p(\eta )}$$
(34)
$$\frac{\sqrt{\sigma (y)}}{n_0(0)}\underset{0}{\overset{2y}{}}d\tau G(y,0,y\tau ,\xi )\widehat{r}_1(\tau ,\xi )\}.$$
Finally, the desired system of linear Volterra integral equations may be obtained from (34) by differentiation and setting $`\xi `$ equal to e.g. $`\xi _1=(a,0)`$ and $`\xi _2=(0,a)`$, where $`a0`$. Namely, we obtain
###### Proposition 4.1.
Let $`\overline{\phi }(y)=(\phi _1(y),\phi _2(y))`$ be given by (33) where the scalar factor $`p(y)`$ depends only on the already found $`n_0(y)`$. Then $`\overline{\phi }(y)`$ satisfies the system of linear Volterra equations
$`a\phi _i(y)=2{\displaystyle \underset{0}{\overset{y}{}}}𝑑\eta G_1(y,\eta ,\eta y,\xi _i){\displaystyle \frac{<\overline{\phi }(\eta ),_\xi U_{tt}^{(0)}(\eta ,2y\eta ,\xi _i)>}{n_0^2(\eta )p(\eta )}}`$
$`{\displaystyle \frac{2}{n_0(0)}}G_1(y,0,y,\xi _i)\widehat{r}_1(2y,\xi _i)+{\displaystyle \underset{0}{\overset{y}{}}}𝑑\eta {\displaystyle \underset{\eta }{\overset{2y\eta }{}}}𝑑\tau G_2(y,\eta ,y\tau ,\xi _i){\displaystyle \frac{<\overline{\phi }(\eta ),_\xi U_{tt}^{(0)}(\eta ,\tau ,\xi _i)>}{n_0^2(\eta )p(\eta )}}`$
$`{\displaystyle \frac{1}{n_0(0)}}{\displaystyle \underset{0}{\overset{2y}{}}}𝑑\tau G_2(y,0,y\tau ,\xi _i)\widehat{r}_1(\tau ,\xi _i),i=1,2.`$
where
$$G_1(y,\eta ,t\tau ;\xi )=\sqrt{\sigma (y)}G(y,\eta ,t\tau ;\xi )G_2(y,\eta ,t\tau ;\xi )=G_{1t}(y,\eta ,t\tau ;\xi )+G_{1y}(y,\eta ,t\tau ;\xi ).$$
Observe that, due to the first equation in (25),
$$\underset{\eta 0}{lim}_\xi U_{tt}^{(0)}(\eta ,\tau )=0,$$
so that the system of integral equtions in Proposition 4.1 is not singular.
### 4.2. Two layers problem
Let us consider the inverse problem for the first-order approximation in the case of a layer and a semi-infinite bottom separated by an interface. Our goal is to derive a coupled system of linear Volterra integral equations similar to that in Proposition 4.1 to reconstruct $`\overline{n}(z)`$ and also to determine constants $`\overline{h}`$ characterizing this interface. Recall the formulation of the corresponding problem using $`y`$ variable, see (13), (14),
(35)
$$U_{tt}^{(1)}\frac{1}{\sigma (y)}\frac{}{y}\left(\sigma (y)U_y^{(1)}\right)+\frac{\xi ^2}{n_0^2}U^{(1)}=\frac{1}{n_0^2}<\overline{n},_\xi U_{tt}^{(0)}>,$$
$$U^{(1)}|_{y=0}=0,U_y^{(1)}|_{y=0}=\frac{\widehat{r}_1(t,\xi )}{n_0(0)}.$$
$$[U^{(1)}n_0<\overline{h},_\xi U_y^{(0)}>]|_{y=L}=0,$$
$$[\sigma (y)U_y^{(1)}(\rho \sigma ^2(y)<\overline{h},_\xi U_{yy}^{(0)}>+\rho \sigma (y)\sigma ^{}(y)<\overline{h},_\xi U_y^{(0)}>U^{(0)}\frac{<\xi ,\overline{h}>}{\rho })]|_{y=L}=0.$$
As before, we start with the singularity analysis of the transmitted and reflected waves. The progressive wave expansion of the inhomogeneous term in (35) is given by
$$\frac{}{\xi _i}U_{tt}^{(0)}=\frac{\xi _i\delta (ty)}{\sqrt{\sigma (y)}}\underset{0}{\overset{y}{}}\frac{d\eta }{n_0^2(\eta )}+\frac{\delta (t+y2L)}{\sqrt{\sigma (y)}}\left(\frac{1}{2}\frac{}{\xi _i}d_1\xi _i\underset{L}{\overset{y}{}}\frac{d\eta }{n_0^2(\eta )}\right)+\mathrm{},$$
$$0<y<L,L<t<2L,i=1,2,$$
$$\frac{}{\xi _i}U_{tt}^{(0)}=\frac{\delta (ty)}{\sqrt{\sigma (y)}}\left(\frac{1}{2}\frac{}{\xi _i}s_1\xi _i\underset{L}{\overset{y}{}}\frac{d\eta }{n_0^2(\eta )}\right)+\mathrm{},y>L,L<t<3L.$$
Then the sum of the incident and reflected waves is given by
(36)
$$U^{(1)}=\frac{\theta (ty)}{\sqrt{\sigma (y)}}\underset{0}{\overset{y}{}}<\overline{n},\xi >p(\eta )d\eta +\frac{\theta (t+y2L)}{\sqrt{\sigma (y)}}\left(\underset{L}{\overset{y}{}}<\overline{n},\overline{p}_1>d\eta +D\right)+\mathrm{},$$
$$0<y<L,L<t<2L,$$
where
$$\overline{p}_1^{(i)}(y)=\frac{1}{2n_0^2(y)}\left(\frac{1}{2}\frac{}{\xi _i}d_1\xi _i\underset{L}{\overset{y}{}}\frac{d\eta }{n_0^2(\eta )}\right),0<y<L.$$
The transmitted wave is then
(37)
$$U^{(1)}=\frac{\theta (ty)}{\sqrt{\sigma (y)}}\left(\underset{L}{\overset{y}{}}<\overline{n},\overline{p}_2>d\eta +S\right)+\mathrm{},y>L,L<t<3L,$$
where
$$\overline{p}_2(y)=\frac{1}{2n_0^2(y)}\left(\frac{1}{2}\frac{}{\xi _i}s_1\xi \underset{L}{\overset{y}{}}\frac{d\eta }{n_0^2(\eta )}\right),y>L$$
and coefficients $`d_1,s_1`$ are defined in (27).
Substituting these expressions for $`U^{(1)}`$ into the interface continuity relations in (35), and solving the corresponding system of linear equations, we find that
(38) $`D(\xi )={\displaystyle \frac{(1\frac{\sigma _+}{\sigma _{}})\underset{0}{\overset{L}{}}<\overline{n},\xi >p(y)dy+\frac{\sigma _+}{\sigma _{}}<\overline{h},\overline{\mathrm{\Gamma }}_1>+<\overline{h},\overline{\mathrm{\Gamma }}_2>}{1+\frac{\sigma _+}{\sigma _{}}}},`$
$`S(\xi )={\displaystyle \frac{2\underset{0}{\overset{L}{}}<\overline{n},\xi >p(y)dy+<\overline{h},\overline{\mathrm{\Gamma }}_2\overline{\mathrm{\Gamma }}_1>}{\sqrt{\frac{\sigma _{}}{\sigma _+}}+\sqrt{\frac{\sigma _+}{\sigma _{}}}}}.`$
Here the vectors $`\overline{\mathrm{\Gamma }}_1`$ and $`\overline{\mathrm{\Gamma }}_2`$ are given by
$`\overline{\mathrm{\Gamma }}_1^{(i)}=\rho _1\sigma _{}\left({\displaystyle \frac{}{2\xi _i}}d_1+\xi _i{\displaystyle \underset{0}{\overset{L}{}}}{\displaystyle \frac{d\eta }{n_0^2(\eta )}}\right)+\sqrt{{\displaystyle \frac{\sigma _{}}{\sigma _+}}}\rho _2\sigma _+{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{\xi _i}}s_1,`$
$`\overline{\mathrm{\Gamma }}_2^{(i)}=\rho _1\sigma _{}\left({\displaystyle \frac{}{2\xi _i}}d_1\xi _i{\displaystyle \underset{0}{\overset{L}{}}}{\displaystyle \frac{d\eta }{n_0^2(\eta )}}\right)\sqrt{{\displaystyle \frac{\sigma _+}{\sigma _{}}}}\rho _2\sigma _+{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{\xi _i}}s_1,i=1,2.`$
From the jump of $`U^{(1)}`$ at $`y=0`$ at the time $`t=2L`$, we can find $`D`$. Now equation (38) with $`\xi _1=(a,0)`$ and $`\xi _2=(0,a)`$ make possible to find $`\overline{h}`$,
(39)
$$h_i=\left[D(\xi _i)\left(1+\frac{\sigma _+}{\sigma _{}}\right)\left(1\frac{\sigma _+}{\sigma _{}}\right)\underset{0}{\overset{L}{}}<\overline{n},\xi _i>p(y)dy\right]\left[2a\rho _1\frac{\sigma _+\sigma _{}}{\sigma _++\sigma _{}}\underset{0}{\overset{L}{}}\frac{d\eta }{n_0^2(\eta )}\right]^1.$$
Clearly, at this stage we can also find $`S`$.
Let us describe an algorithm to determine $`\overline{n}(z)`$ in the second layer. The boundary conditions $`U^{(1)}|_{y=L+0}`$ and the response $`U_y^{(1)}|_{y=L+0}`$ for the transmitted wave can be found using $`\widehat{r}_1(t,\xi )`$ and already known, in the upper layer, $`n_0(y),\overline{n}(y)`$ together with the interface continuity conditions in (35). Using the coordinates $`y^{},\eta ^{},t^{},y=y^{}+L,\eta =\eta ^{}+L,t=t^{}+L`$, see Fig. 1, we have the progressive wave expansion for $`U^{(1)}`$,
$$U^{(1)}(y^{},t^{})=\frac{\theta (ty^{})}{\sqrt{\sigma (y^{})}}\left(\underset{0}{\overset{y^{}}{}}<\overline{n},\xi >p_2(\eta )d\eta +S\right)+\mathrm{},\overline{p}_2(y^{})=\xi p_2(y^{}).$$
Employing the Green function (31) in the second layer, we obtain
$`p_2(y^{})<\overline{n}(y^{}),\xi >=`$
$`{\displaystyle \frac{d}{dy^{}}}\{\sqrt{\sigma (y^{})}({\displaystyle \underset{0}{\overset{y^{}}{}}}d\eta ^{}{\displaystyle \underset{\eta ^{}}{\overset{2y^{}\eta ^{}}{}}}d\tau G(y^{},\eta ^{},y^{}\tau ){\displaystyle \frac{<\overline{n}(\eta ^{}),_\xi U_{tt}^{(0)}(\eta ^{},\tau ,\xi )>}{n_0^2(\eta ^{})}}`$
$`{\displaystyle \underset{0}{\overset{2y^{}}{}}}𝑑\tau G(y^{},0,y^{}\tau )\left(\chi _2(\tau )+{\displaystyle \frac{\sigma _+^{}}{\sigma _+}}\chi _1(\tau )\right)+{\displaystyle \underset{0}{\overset{2y^{}}{}}}𝑑\tau G_\eta ^{}(y^{},0,y^{}\tau )\chi _1(\tau )`$
$`+[G(y^{},0,y^{})\chi _1(0)G(y^{},0,y^{})\chi _1(2y^{})])\},`$
where
$`\chi _1(t)=U^{(1)}|_{y=L},\chi _2(t)=U_y^{(1)}|_{y=L}`$
Finally, the desired coupled system of linear Volterra integral equations determining $`\overline{n}(y)`$ below the interface may be obtained from this equation by setting e.g. $`\xi _1=(a,0)`$ and $`\xi _2=(0,a)`$ (compare with Proposition 4.1). It is worth noting that these integral equations are not singular.
## 5. Global reconstruction and stability
As mentioned in section 3, the non-linear Volterra-type system of integral equations (25) may, in principle, be solvable only locally. In this section we show that, assuming a priori bounds for $`q(t)`$ in $`C(0,T)`$, it can be found on the whole interval $`(0,T)`$ using a layer-stripping method based on (25). What is more, we will prove Lipschitz stability of this procedure with respect to the variation of the response data $`g(t)`$ in $`C(0,2T)`$.
We first note that if $`q_{C(0,T)}<M`$ then solving the direct problem gives
(40)
$$\psi _{1,2}_{C(\mathrm{\Delta }_T)},q_{C(0,T)}<A=A(T,M),$$
where we can assume $`A\mathrm{max}(1,M)`$. Here $`\mathrm{\Delta }_T`$ and generally $`\mathrm{\Delta }_b,\mathrm{\hspace{0.17em}0}<b<T`$, is the triangle in $`R^2`$ bounded by the characteristics $`\eta =\tau ,\eta +\tau =2b`$ and the axis $`\eta =0`$. Assume that we have already found $`\psi _{1,2}`$ in $`\mathrm{\Delta }_b`$ and $`q`$ in $`(0,b)`$.
###### Lemma 5.1.
Let $`\lambda >0`$ satisfies
(41)
$$\lambda <\frac{1}{16(1+A)^2}.$$
Then system (25) uniquely determines $`\psi _{1,2}(y,t),q(y)`$ for $`(y,t)\mathrm{\Delta }_{b+\lambda }`$, $`y(0,b+\lambda )`$, respectively. Moreover, these functions can be found from system (25) by Picard iterations.
Proof Utilizing new coordinates $`y^{}=yb,\eta ^{}=\eta b,t^{}=tb`$ (compare with the end of section 3), we obtain a system of non-linear Volterra-type equations similar to (28),
(42)
$$\{\begin{array}{ccc}\hfill \psi _1(y^{},t^{})& =& \underset{0}{\overset{y^{}}{}}\psi _2(\eta ^{},t^{}+\eta ^{}y^{})𝑑\eta ^{}+g_1(y^{},t^{}),\hfill \\ \hfill \psi _2(y^{},t^{})& =& \underset{0}{\overset{y^{}}{}}q(\eta ^{},\xi )\psi _1(\eta ^{},t^{}\eta ^{}+y^{})𝑑\eta ^{}+g_2(y^{},t^{}),\hfill \\ \hfill q(y^{})& =& 2\underset{0}{\overset{y^{}}{}}q(\eta ^{})\psi _1(\eta ^{},2y^{}\eta ^{})𝑑\eta ^{}+g_3(y^{}).\hfill \end{array}$$
Here $`g_1(y^{},t^{})=\psi _1(0,t^{}y^{}),g_2(y^{},t^{})=\psi _2(0,t^{}=y^{}),g_3(y^{})=2\psi _2(0,2y^{})`$ may be found from the "direct" system (21) and satisfy $`|g_i|<A`$.
Using the first equation in (42), we eliminate $`\psi _1`$ to get
(43)
$$\{\begin{array}{ccc}\hfill \psi _2(y^{},t^{})& =& \underset{0}{\overset{y^{}}{}}q(\eta ^{},\xi )\left[g_1(\eta ^{},t^{}\eta ^{}+y^{})+_0^\eta ^{}\psi _2(\xi ^{},t^{}\eta ^{}+\xi ^{})𝑑\xi ^{}\right]𝑑\eta ^{}+g_2(y^{},t^{}),\hfill \\ \hfill q(y^{})& =& 2\underset{0}{\overset{y^{}}{}}q(\eta ^{})\left[g_1(\eta ^{},2y^{}\eta ^{})+_0^\eta ^{}\psi _2(\xi ^{},y^{}\eta ^{}+\xi ^{})𝑑\xi ^{}\right]𝑑\eta ^{}+g_3(y^{}).\hfill \end{array}$$
We rewrite this system in the form,
$$(\psi _2,q)=K_\lambda (\psi _2,q),$$
where $`K_\lambda `$ is a non-linear operator in $`C(\mathrm{\Delta }_\lambda \times C(0,\lambda )`$ determined by the right-hand side (43). Let us show that $`K_\lambda `$ is a contraction in
$$𝐃_{2A}=B_{2A}(C(\mathrm{\Delta }_\lambda ))\times B_{2A}(C(0,\lambda )),$$
where $`B_\rho ()`$ is the ball of radius $`\rho `$ in the corresponding function space and we use the norm
$$(\psi _2,q)=\mathrm{max}(\psi _2,q).$$
Considering
$$K_\lambda (\psi _2,q)K_\lambda (\widehat{\psi }_2,\widehat{q})=\left(K_\lambda (\psi _2,q)K_\lambda (\psi _2,\widehat{q})\right)+\left(K_\lambda (\psi _2,\widehat{q})K_\lambda (\widehat{\psi }_2,\widehat{q})\right),$$
it follows from (43) that in $`𝐃_{2A}`$
$$K_\lambda (\psi _2,q)K_\lambda (\widehat{\psi }_2,\widehat{q})2A\lambda \left(q\widehat{q}+2\lambda \psi _2\widehat{\psi }_2+2\lambda q\widehat{q}\right).$$
As $`\lambda `$ satisfies (41), this implies that
$$K_\lambda (\psi _2,q)K_\lambda (\widehat{\psi }_2,\widehat{q})\frac{3}{8}(\psi _2,q)(\widehat{\psi }_2,\widehat{q}).$$
Similar arguments show that $`K_\lambda :𝐃_{2A}𝐃_{2A}`$. Thus (43) has a unique solution in $`𝐃_{2A}`$ and, since (28) is a Volterra-type system, in $`C(\mathrm{\Delta }_\lambda )\times C(0,\lambda )`$.
QED
Next we prove the Lipschitz stability of the non-linear Volterra system (25). Let us denote by $`\stackrel{~}{\psi }_{1,2}(y,t)`$, $`\stackrel{~}{q}(y,\xi )`$ and $`\stackrel{~}{g}(t)`$ small variations of the corresponding functions.
###### Lemma 5.2.
Let $`\psi _{1,2},q`$ satisfy (40). There exist $`c_0=c_0(T,A),ϵ_0=ϵ0(T,A)`$ such that for $`\stackrel{~}{g}_{C(0,2T)}<ϵ<ϵ_0`$ the corresponding system (25) with $`g+\stackrel{~}{g}`$ instead of $`g`$ has a unique solution $`\psi _{1,2}+\stackrel{~}{\psi }_{1,2}`$, $`q+\stackrel{~}{q}`$, and
(44)
$$\stackrel{~}{\psi }_{1,2}_{C(\mathrm{\Delta }_T)},\stackrel{~}{q}_{C(0,2T)}<c_0ϵ.$$
Proof Substituting expressions for $`\psi _{1,2}(y,t)+\stackrel{~}{\psi }_{1,2}(y,t)`$, $`q(y)+\stackrel{~}{q}(y)`$ and $`g(t)+\stackrel{~}{g}(t)`$ into (25) , we obtain the corresponding system in variations
$$\{\begin{array}{ccc}\hfill \stackrel{~}{\psi }_1(y,t)& =& \underset{0}{\overset{y}{}}\stackrel{~}{\psi }_2(\eta ,t+\eta y)𝑑\eta ,\hfill \\ \hfill \stackrel{~}{\psi }_2(y,t)& =& \underset{0}{\overset{y}{}}q(\eta )\stackrel{~}{\psi }_1(\eta ,t\eta +y)𝑑\eta \hfill \\ & & \underset{0}{\overset{y}{}}\stackrel{~}{q}(\eta )(\psi _1(\eta ,t\eta +y)+\stackrel{~}{\psi }_1(\eta ,t\eta +y))𝑑\eta +\stackrel{~}{g}(t+y),\hfill \\ \hfill \stackrel{~}{q}(y)& =& 2\underset{0}{\overset{y}{}}\stackrel{~}{q}(\eta )\psi _1(\eta ,2y\eta )𝑑\eta 2\underset{0}{\overset{y}{}}(q(\eta )+\stackrel{~}{q}(\eta ))\stackrel{~}{\psi }_1(\eta ,2y\eta )𝑑\eta +2\stackrel{~}{g}(2y).\hfill \end{array}$$
Let $`\stackrel{~}{g}<ϵ`$. Then, we have
$$\{\begin{array}{ccc}\hfill |\stackrel{~}{\psi }_1(y,t)|& & \underset{0}{\overset{y}{}}|\stackrel{~}{\psi }_2(\eta ,t+\eta y)|𝑑\eta ,\hfill \\ \hfill |\stackrel{~}{\psi }_2(y,t)|& & 2A\underset{0}{\overset{y}{}}|\stackrel{~}{\psi }_1(\eta ,t\eta +y)|𝑑\eta +A\underset{0}{\overset{y}{}}|\stackrel{~}{q}(\eta )|𝑑\eta +ϵ,\hfill \\ \hfill |\stackrel{~}{q}(y)|& & 4A\underset{0}{\overset{y}{}}|\stackrel{~}{q}(\eta )|𝑑\eta +2A\underset{0}{\overset{y}{}}|\stackrel{~}{\psi }_1(\eta ,2y\eta )|𝑑\eta +2ϵ.\hfill \end{array}$$
Substituting the upper inequality into the second one to replace $`|\stackrel{~}{\psi }_1(\eta ,t\eta +y)|`$, we obtain
$$|\stackrel{~}{\psi }_2(y,t)|2A\underset{\mathrm{\Delta }_{y,t}}{}|\stackrel{~}{\psi }_2(\eta ,\tau )|𝑑\eta 𝑑\tau +2A\underset{0}{\overset{y}{}}|\stackrel{~}{q}(\eta )|𝑑\eta +ϵ,$$
where $`\mathrm{\Delta }_{y,t}`$ is the the triangle in $`R^2`$ bounded by the characteristics $`\tau =\eta +ty,\tau =t+y\eta `$ and the axis $`\eta =0`$. Similarly,
$$|\stackrel{~}{q}(y)|4A\underset{0}{\overset{y}{}}|\stackrel{~}{q}(\eta )|𝑑\eta +4A\underset{\mathrm{\Delta }_{y,t}}{}|\stackrel{~}{\psi }_2(\eta ,\tau )|𝑑\eta 𝑑\tau +2ϵ.$$
Let
$$p(\eta )=\mathrm{max}_\tau |\stackrel{~}{\psi }_2(\eta ,\tau )|,(\eta ,\tau )\mathrm{\Delta }_T.$$
Then, we have
$$|\stackrel{~}{\psi }_2(y,t)|4AT\underset{0}{\overset{y}{}}p(\eta )𝑑\eta +2A\underset{0}{\overset{y}{}}|\stackrel{~}{q}(\eta )|𝑑\eta +ϵ,$$
$$|\stackrel{~}{q}(y)|4A\underset{0}{\overset{y}{}}|\stackrel{~}{q}(\eta ,\xi )|𝑑\eta +8AT\underset{0}{\overset{y}{}}p(\eta )𝑑\eta +2ϵ,$$
as
$$\underset{\mathrm{\Delta }_{y,t}}{}p(\eta )𝑑\eta 𝑑\tau =2\underset{0}{\overset{y}{}}p(\eta )(y\eta )𝑑\eta 2T\underset{0}{\overset{y}{}}p(\eta )𝑑\eta .$$
Moreover, it holds that
$$p(y)4AT\underset{0}{\overset{y}{}}p(\eta )𝑑\eta +2A\underset{0}{\overset{y}{}}|\stackrel{~}{q}(\eta )|𝑑\eta +ϵ.$$
Introduce
$$\rho (y)=p(y)+|\stackrel{~}{q}(y)|.$$
Adding the last two inequalities for $`p(y)`$ and $`|\stackrel{~}{q}(y)|`$ yields
$$\rho (y)\underset{0}{\overset{y}{}}(12ATp(\eta )+6A|\stackrel{~}{q}(\eta )|)𝑑\eta +3ϵ,$$
or, with $`C=6A\mathrm{max}(1,2T)`$,
$$\rho (y)C\underset{0}{\overset{y}{}}\rho (\eta )𝑑\eta +3ϵC\underset{0}{\overset{y}{}}\left(C\underset{0}{\overset{\eta }{}}\rho (\eta _1)𝑑\eta _1+3ϵ\right)𝑑\eta +3ϵ=$$
$$\frac{C^2}{1!}\underset{0}{\overset{y}{}}\rho (\eta )(y\eta )𝑑\eta +(Cy+1)3ϵ.$$
Continuing this process, we come to the estimate
$$\rho (y)\frac{C^{n+1}}{n!}\underset{0}{\overset{y}{}}\rho (\eta )(y\eta )^n𝑑\eta +\left(y^n\frac{C^n}{n!}+y^{n1}\frac{C^{n1}}{(n1)!}+\mathrm{}+1\right)3ϵ.$$
Let now $`\rho _0=\mathrm{max}\rho (y)`$ for $`y[0,T]`$. Then,
$$\rho (y)\rho _0\frac{C^{n+1}y^{n+1}}{(n+1)!}+3e^{Cy}ϵ,$$
i.e.
$$\rho _0\rho _0\frac{(CT)^{n+1}}{(n+1)!}+3e^{CT}ϵ.$$
Clearly, for sufficiently large $`n`$ it holds that
$$\frac{(CT)^{n+1}}{(n+1)!}\frac{1}{2},$$
and then,
$$\rho _06e^{CT}ϵ.$$
Thus, assuming $`\stackrel{~}{g}ϵ`$, we obtain that
(45)
$$\stackrel{~}{q}6e^{CT}ϵ,\stackrel{~}{\psi }_26e^{CT}ϵ,\stackrel{~}{\psi }_16Te^{CT}ϵ,$$
which implies (44) with
$$c_0=6\mathrm{max}(1,T)e^{CT},ϵ_0=\frac{1}{6\mathrm{max}(1,T)}e^{CT}\left(A\mathrm{max}(\psi _{1,2},q)\right).$$
QED
Observe that when $`\stackrel{~}{g}=O(1)`$ estimate (45) does not imply, for sufficiently large $`T`$, that $`\stackrel{~}{\psi }_{1,2}+\psi _{1,2},\stackrel{~}{q}+q`$ satisfy (40). This is a manifestation of a possible non-convergence of the Picard method for large $`\stackrel{~}{g}`$.
Observe also that the conditional stability, in $`C(0,T)`$, of the inversion method based on the Volterra-type system (25) implies the conditional stability, in $`C^2(0,T)`$, of the original inverse problem of the reconstruction of $`n_0(y)`$. Indeed, if $`n_0(y)`$ is a priori bounded in $`C^2(0,T)`$, i.e. $`q(y,\xi )`$ is a priori bounded in $`C^2(0,T)`$ for bounded $`\xi `$, the inverse map, $`g(t,\xi )n_0(y)`$, is Lipschitz stable from $`C(0,2T)`$ to $`C(0,T)`$. Due to (20), this implies the Lipschitz stability of the above map from $`C(0,2T)`$ to $`C^2(0,T)`$.
Regarding Volterra system in Proposition 4.1, we first note that, if $`n_0C^3(0,T)`$, then $`\mathrm{}_\xi U_{tt}^{(0)}C(\mathrm{\Delta }_T)`$. This implies the Lipschitz stability of this system in $`C(0,T)`$. It is clear from (30), (33) and can be also checked directly using the fact that
$$\widehat{r}_1(t)=n_0(0)U^{(0)}(\eta ,t\tau )n_0^2(\eta )<\overline{n}(\eta ),\mathrm{}_\xi U_{tt}^{(0)}(\eta ,\tau >d\eta d\tau ,$$
that $`\widehat{r}_1(t)=tf(t),fC(0,2T)`$. Therefore, the inversion method to find $`\overline{n}`$ using Proposition 4.1 is Lipschitz stable into $`C(0,T)`$ with respect to the variation of $`\widehat{r}_1`$ in the norm $`t^1\widehat{r}_1_{C(0,2T)}`$.
## 6. Numerical results
In this section we demonstrate efficiency of the described method in numerical solution of the inverse problem. The computational results were obtained for the two-dimensional problem with coordinates $`(z,x)`$. In this case,
$$n^2=n_0^2(z)+ϵxn_1(z)+O(ϵ^2),h(ϵx)=h_0+ϵxh_1+O(ϵ^2).$$
The constant $`h_0`$ is determined by the arrival time of the $`\delta `$type singularity reflected from the interface, while $`h_1`$ may be evaluated by measuring, at $`z=0`$, the amplitudes of the reflected waves, see (39). In the present numerical implementation, we assume that their values are known. The described algorithms for solving the zero-order and first-order inverse problems are implemented into a computer code to reconstruct $`n_0(z)`$ and $`n_1(z)`$. The first part of the computer code generates responses for both orders approximations $`r^{(0)}(t,\xi )`$ and $`r^{(1)}(t,\xi )`$. The chosen profiles are described by
$$\sigma _0(z)=p_0+p_1z+p_2z^2+q\mathrm{sin}f_0z$$
for the zero-order problem, and by the trigonometric polynomial
$$n_1(z)=r_0+r_1\mathrm{cos}f_1z+r_2\mathrm{cos}2f_1z+q_1\mathrm{sin}f_1z+q_2\mathrm{sin}2f_1z$$
for the first-order problem. Taking various coefficients in the above representations below and above the interface, we present on Fig.3 and 4 the numerically computed values of $`\sigma _0(z)`$ and $`n_1(z)`$ against original data. Here we chose $`\rho _1=1`$ and $`\rho _2=1.5`$. As we can see there is a good agreement of exact and computed profiles in both cases.
On Fig.5 we demonstrate the results in the case of the response data are corrupted by a $`2\%`$ noise for $`R_0`$ and a $`10\%`$ noise for $`R_1`$. The ratio of 2 and 10 may be explained by assuming that $`ϵ=\frac{1}{5}`$. The data with $`5\%`$ noise for $`R_0`$ and $`25\%`$ for $`R_1`$ are presented in Fig.6. To clarify the results of numerical reconstructions one should take into the account that, with respect to the complete refractive index, the error in the first-order approximation should be multiplied by $`ϵ1`$.
The method has shown to be quite stable, fast and accurate. When solving Volterra-type integral equations, both non-linear and linear, the iteration processes need just a few iterations (for all graphs the number of iterations was chosen 10). Clearly, the number of iterations and accuracy in the reconstruction depend on the scale of discretization. On Fig. 7 we demonstrate the error dependence on the scale of discretization.
Numerous computer experiments have shown that for a better accuracy and fast convergence of the iteration process it is reasonable to use for the chosen profiles the segment $`|\xi |<0.5`$. It is worth noting that the parameter $`|\xi |`$, the maximum depth $`T`$ and the maximum of $`n_0^{\prime \prime }(z)`$ are interconnected. For example, for the larger values of $`T`$ and the maximum of $`n_0^{\prime \prime }(z)`$, while computing the profiles we were forced to take smaller values of $`|\xi |`$. Moreover, due to the non-linearity of (25), when we increase $`T`$ and/or $`n_0^{\prime \prime }(z)`$ and $`|\xi |`$, a blow up effect can occur, i.e. the iterations stop to converge. This may be remedied, using the results of section 5, by a variant of the layer-stripping as it was used in the sections 3.2 and 4.2 but for an interface without discontinuity.
In applications to geophysics, the unity of the refractive index corresponds to the average speed of the wave propagation $`c=`$2-2.5km/sec. Thus, the dimensionless depth coordinate $`z`$ must be multiplied by $`22.5km`$.
## 7. Concluding remarks
### 7.1.
Typical distances of interest in seismology/oil exploration are few kilometers. As a typical velocity of the wave propagation is around $`22.5km/sec`$, in the travel-time coordinates, $`x,z=O(1)`$. As we have already mentioned in Introduction, the method described in the paper can, in principle, reconstruct the velocity profile in this region up to an error of the order $`O(ϵ^2)`$. Methods based on an approximation of WLIM by a purely layered medium would give rize to an error of the order $`O(ϵ^2+ϵ|x|)`$. A natural way to improve the result when using inversion techniques for purely layered medium is to increase the number of sources placing them at distance $`O(ϵ)`$. However, in applications to seismology/oil exploration this is not always possible. Indeed, a typical structure of the earth contains, in addition to WLIM, various inclusion of different nature with domains of interest often lying below these inclusions. In map migration method, the rays used often propagate oblique to the surface $`y=0`$ with their substantial part lying in WLIM making it desirable to know well the properties of this medium. Taking into account that, in order to determine the velocity profile in WLIM up to depth $`T=O(1)`$ it is necessary to make measurements during the time interval $`0<t<2T`$, the sources should be located at a distance $`O(1)`$ from the inclusion not to be contaminated by its influence. Therefore, using inversion methods based on an approximation by a purely layered medium, we would end up with a reconstruction error, near inclusion, of the order $`O(ϵ)`$.
### 7.2.
Another observation, partly related to the above one, concerns with the case when it is necessary/desirable to make measurements only on a part of the ground surface, $`y=0`$, near the origin. Observe that, although integrals (12), (15) is taken over $`R^2`$, due to the finite velocity of the wave propagation, $`R(x,t)=0`$ for $`x`$ with $`d(x,0)>t`$, where $`d((x,y),(\stackrel{~}{x},\stackrel{~}{y})`$ are the distance in the metric $`dl^2=c^2dx^2+dy^2`$. Consider e.g. an inclusion located near $`y=0`$ at the distance less then $`2T`$ from the origin which would contaminate the measurements. If we, however, make measurements near the origin, the inclusions starts to affect our measurements only when distance goes down to $`T`$. According to a result obtained by the BC-method and valid for a general multidimensional medium, making measurements on a subdomain, $`\mathrm{\Gamma }`$, on the surface during time $`2T`$ makes possible, in principle, to recover the velocity profile in the $`T`$neighbourhood of $`\mathrm{\Gamma }`$ . In the case of a layered medium, due to it is even sufficient to make measurements in a single point on $`y=0`$.
Let us show that a simple modification of the procedure described in the paper makes it possible to determine $`c_0(y),c_1(y),\mathrm{\hspace{0.17em}0}<y<T`$ given $`R(x,t)`$ for $`|x|<2a,t<2T`$, where $`a>0`$ is arbitrary. Denote by $`\widehat{c}=\mathrm{max}c(x,y),`$ over $`(x,y)`$ lying at the distance less than $`2T`$ from the origin. Observe that, for any $`b>0`$, the layer $`y>b`$ affects $`R(x,t)`$ with $`|x|>2a`$ only when $`t>2t(a,b)`$. Here $`2t(a,b)`$ is the time needed for a wave from the origin which propagates through a medium with the length element $`dl^2=\widehat{c}^2dx^2+dy^2`$ to reach the layer $`y=b`$ and return to the surface $`y=0`$ at a point $`|x|>2a`$, i.e.
$$t(a,b)=\sqrt{b^2+a^2\widehat{c}^2}.$$
Therefore, if we know $`c_0(y),c_1(y)`$ for $`y<b`$ and $`R(x,t)`$ for $`|x|<2a,t<2T`$, we can determine, up to an error of the order $`O(ϵ^2)`$, $`R(x,t)`$ for all $`xR^2`$ and $`t<2t(a,b)`$. Clearly, the procedure described makes it possible to determine $`c_0(y),c_1(y)`$ for $`y<t(a,b)`$. Iterating this process, we reach the level $`y=T`$ in a finite number of steps.
Acknowledgements The authors would like to acknowledge the financial support from EPSRC grants GR/R935821/01 and GR/S79664/01. they are grateful to Prof. C. Chapman for numerous consultations on the geophysical background of the problem, Dr. K. Peat for the assistance with numerics for the direct problem and Prof. A.P.Katchalov for stimulating discussions.
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# CPT Violation: What and where to look for
## 1 Introduction: CPT Theorem
A century after Einstein’s annus mirabilis, and ninety years after his revolutionary proposal on the dynamical nature of space time, which is the basis for the classical theory of General Relativity, we are still lacking a consistent theory that would describe the quantum nature of space time at short-distance scales, of order of the Planck length $`10^{35}`$ m. Any complete theory of quantum gravity is bound to address fundamental issues, directly related to the emergence of space-time and its structure at energies beyond the Planck energy scale $`M_P10^{19}`$ GeV. From our relatively low energy experience so far, we are lead to expect that a theory of quantum gravity should respect most of the fundamental symmetries of particle physics, that govern the standard model of electroweak and strong interactions: Lorentz symmetry and CPT invariance, that is invariance under the combined action of Charge Conjugation (C), Parity (reflection P) and Time Reversal Symmetry (T). Actually the latter invariance is a theorem of any local quantum field theory that we can use to describe the standard phenomenology of particle physics to date. The CPT theorem can be stated as follows cpt : Any quantum theory, formulated on flat space time is symmetric under the combined action of CPT transformations, provided the theory respects (i) Locality, (ii) Unitarity (i.e. conservation of probability) and (iii) Lorentz invariance.
## 2 Potential CPT Violation
If such a theorem exists, then why do we have to bother to test CPT invariance, given that all our phenomenology up to now has been based on such quantum theories ? The answer to this question is intimately linked with our understanding of quantum gravity.
### 2.1 CPT Violation through Decoherence: Quantum Gravity as a ‘medium’, opening up the matter quantum system
First of all, the theorem is not valid (at least in its strong form) in highly curved (singular) space times, such as black holes, or in general in space-time backgrounds of some quantum gravity theories involving the so-called quantum space-time foam backgrounds wheeler , that is singular quantum fluctuations of space time geometry, such as black holes etc, with event horizons of microscopic Planckian size ($`10^{35}`$ meters). Such backgrounds result in apparent violations of unitarity in the following sense: there is part of information (quantum numbers of incoming matter) “disappearing” inside the microscopic event horizons, so that an observer at asymptotic infinity will have to trace over such “trapped” degrees of freedom. This would define a “space time foam” situation wheeler , in which the quantum gravity ground state resembles that of a decoherening “medium” in open system quantum mechanics ehns ; lopez ; mavromatosdecohe . Barring the exciting possibility of holographic properties of such microscopic black holes maldacena , one may face, under such circumstances, a situation in which an initially pure state evolves in time to get mixed: the asymptotic states are described by density matrices, defined as follows: $`\rho _{\mathrm{out}}=\mathrm{Tr}_M|\psi ><\psi |`$, where the trace is over trapped (unobserved) quantum states, that disappeared inside the microscopic event horizons in the foam. Such a non-unitary evolution results in the impossibility of defining a standard quantum-mechanical scattering matrix, connecting asymptotic states in a scattering process: $`|\mathrm{out}>=S|\mathrm{in}>`$, $`S=e^{iH(t_it_f)}`$, where $`t_it_f`$ is the duration of the scattering (assumed much longer than other time scales in the problem). Instead, in foamy situations, one can define an operator that connects asymptotic density matrices hawking : $`\rho _{\mathrm{out}}\mathrm{Tr}_M|\mathrm{out}><\mathrm{out}|=\$\rho _{\mathrm{in}},\$SS^{},`$ where the lack of factorization is attributed to the apparent loss of unitarity of the effective low-energy theory, defined as the part of the theory accessible to low-energy observers who perform scattering experiments. This defines what we mean by particle phenomenology in such situations. The $ matrix is not invertible, and this reflects the effective unitarity loss. It is this property, actually, that leads to a violation of CPT invariance, at least in its strong form, in the sense that the generator of CPT symmetry is ill-defined. This has been shown rigorously in ref. wald . However, as cautiously pointed out there, despite the strong violation of CPT in such a situation, one cannot exclude the possibility that a weak form of CPT invariance remains, which is reflected in the possibility of the “experimentalist” to prepare pure initial quantum states, $`|\psi `$, that evolve into pure asymptotic states, $`|\varphi `$ (defining, in some sense, a “decoherence-free subspace” in the language of open systems), in such a way that CPT is preserved in the respective probabilities: $`P(\psi \varphi )=P(\theta ^1\varphi \theta \psi )`$, where $`\theta `$: $`_{\mathrm{in}}_{\mathrm{out}}`$, with $``$ denoting the appropriate Hilbert state spaces. The notation is such that the CPT operator $`\mathrm{\Theta }`$ acting on density matrices is: $`\mathrm{\Theta }\rho =\theta \rho \theta ^{},\theta ^{}=\theta ^1(\mathrm{anti}\mathrm{unitary})`$.
As we shall discuss in this article, such issues can, in principle, be resolved experimentally, provided of course the sensitivity of the experiments is appropriate to the order expected from theoretical models of quantum gravity.
### 2.2 Cosmological CPT Violation?
Another type of possible violation of CPT mavromatosdecohe , which falls within the remit of the theorem of wald , may be associated with the recent experimental evidence from supernovae and temperature fluctuations of the cosmic microwave background on a current era acceleration of the Universe, and the fact that more than 70 % of its energy budget consists of Dark Energy. If this energy substance is a positive cosmological constant (de Sitter Universe), $`\mathrm{\Lambda }>0`$, then there is a cosmic horizon, in the sense that in a flat Universe, as the data seem to indicate we live in, light emitted at some future moment $`t_0`$ in the cosmic time takes an infinite time to cover the finite horizon radius. The presence of a cosmic future horizon implies impossibility of defining proper asymptotic states, and in particular quantum decoherence of a matter quantum field theory in this de Sitter geometry, as a result of an environment of modes crossing the horizon. In view of wald , then, this would imply an analogous situation with the foam, discussed above, i.e. a strong form of CPT violation. In this particular case, as the horizon is macroscopic, one would definitely have the evolution of an initial pure quantum state to a mixed one, and probably no weak form of CPT invariance would exist, in contrast to the black-hole case mavromatosdecohe .
### 2.3 CPT Violation in the Hamiltonian, consistent with closed system quantum mechanics
Another reason for CPT violation (CPTV) in quantum gravity is spontaneous breaking of Lorentz symmetry (LV), without necessarily implying decoherence. This has also been argued to occur in string theory and other models of quantum gravity sme ; kostelecky , but, in my opinion, no concrete microscopic model has been given as yet. In string theory, for instance, where a LV vacuum has been argued to exist sme , it was not demonstrated that this vacuum is the energetically preferred one. So far, most of the literature in this respect is concentrating on consistent parameterizations of extension of the standard model (SME) sme , which can be used to bound the relevant Lorentz and/or CPT violating parameters. An important difference of this approach, as compared with the decoherence one is the fact that in SME or other spontaneous Lorentz symmetry breaking approaches, the CPT Violation is linked merely with non-commutativity of the CPT operator, which otherwise is a well-defined quantum mechanical operator, with the Hamiltonian of the system. In such models there is a well defined scattering matrix, and the usual phenomenology applies. As we shall argue below, experimentally one can in principle disentangle these two types of violation, by means of appropriate observables and their time evolution properties.
### 2.4 Order of Magnitude Estimates of Quantum Gravity Effects
At first sight, the CPT violating effects can be estimated to be strongly suppressed. Indeed, naively, Quantum Gravity (QG) has a dimensionful coupling constant: $`G_N1/M_P^2`$, where $`M_P=10^{19}`$ GeV is the Planck scale. Hence, CPT violating and decoherening effects may be expected to be suppressed by terms (with dimensions of energy) of order $`E^3/M_P^2`$, where $`E`$ is a typical energy scale of the low-energy probe. If such is the case, the current facilities seem far from reaching this sensitivities. But, as we shall mention below, high energy cosmic neutrino observations may well reach such sensitivities in the future. However, there may be cases where loop resummation and other effects in theoretical models may result in much larger CPT-violating effects of order: $`\frac{E^2}{M_P}`$. This happens, for instance, in some loop gravity approaches to Quantum Gravity, or some non-equilibrium stringy models of space-time foam involving open string excitations mavromatosdecohe . Such large effects can lie within the sensitivities of current or immediate future experimental facilities (terrestrial and astrophysical), and hence can be, or are already, excluded reviews .
## 3 Phenomenology
The phenomenology of CPT Violation is complicated, as there is no single figure of merit, and the relevant sensitivities are highly model dependent. Below I outline briefly, the main phenomenological searches for CPT violation, and the respective sensitivities, in the various major approaches to CPT Violation outlined above.For further details see refs. sme ; kostelecky ; mavromatosdecohe ; reviews .
### 3.1 CPT Violation in the Hamiltonian
The main activities in this area concern: (i) Lorentz violation in extensions of the standard model sme ; kostelecky , which provide an exhaustive phenomenological study of various Lorentz and CPT Violating effects in a plethora of atomic, nuclear and particle (neutrinos) physics experiments. Modified Dirac equation in the presence of external gauge fields can be used for constraining CPTV and Lorentz violating parameters by means of, say, antimatter factories factories . I will not discuss further such tests here. For details I refer the reader to the literature sme ; kostelecky ; mavro\_leap03 . Specific tests for CPTV in the antihydrogen system, via precision measurements of its hyperfine structure, can be found in hayano . At present, the most sensitive of the parameters of CPT and Lorentz violation, $`b`$ in the notation of kostelecky can be constrained to be smaller than $`b<10^{27}`$ GeV (or $`b<10^{31}`$ GeV in masers). Since there is no microscopic model underlying the standard model extension at present it is difficult to interpret such small bounds, as far as sensitivity at Planckian energy scales is concerned. In the naive dimensional estimate that $`[b]=E^2/M_{QG}`$, where $`M_{QG}`$ is a quantum gravity scale, one obtains sensitivities to scales up to $`M_{QG}=10^{27}`$ GeV, thereby tending to excluding linear suppression by the Planck mass. However, if, as expected in such models, there are quadratic (or higher) suppressions by $`M_P`$, then we are some 10 orders of magnitude away from Planck scale physics at present.
(ii) Tests of modified dispersion relations for matter probes. Stringent tests for charged fermions (e.g. electrons) are provided by synchrotron radiation measurements from astrophysical sources, e.g. Crab nebula crab ; reviews , whilst the most accurate tests of modified dispersion relations for photons are at present provided by Gama-Ray-Burst observations on arrival times of radiation at various energy channels nature . The present observations of photons from Gamma Ray Bursts and Active Galactic Nuclei (AGN) concern energy scales up to a few MeV at most. If the quantum gravity effects are then linearly suppressed by the quantum gravity scale, then the sensitivity is $`M_{QG}>10^{16}`$ GeV, otherwise is much less. There is an issue here concerning the universality of quantum gravity effects. As argued in equiv , studies in certain microscopic models of modified dispersion relations in string theory, have indicated that only photons and standard model gauge bosons may be susceptible to quantum gravity effects in some models. This is due to some stringy gauge symmetry protection of chiral matter particles against such effects, to all orders in perturbation theory in a low-energy quantum field theory context. It is in the above sense that it is important to obtain limits on such effects from various sources and for various systems.
(iii) As a final comment in this type of CPT violation we also mention the possibility that quantum gravity may induce non hermiticity of the effective low energy Hamiltonian, in the sense of complex coupling constants appearing in low energy theories okun . This may lead to interesting phenomena, for instance complex anomalous magnetic moments of, say, protons, from which one may place stringent constraints in such effects. Since in all microscopic models we have in our disposition so far such imaginary effects do not appear we shall not pursue this discussion further inhere. However, we stress again, we cannot exclude this possibility from appearing in future microscopic models of quantum gravity.
### 3.2 Quantum Gravity Decoherence and CPT Violation
*Neutral Mesons (Single Kaon states)*
Quantum Gravity may induce decoherence and oscillations among neutral mesons, such as kaons $`K^0\overline{K}^0`$ ehns ; lopez . The modified evolution equation for the respective density matrices of neutral kaon matter can be parametrized as follows ehns :
$$_t\rho =i[\rho ,H]+\delta H/\rho ,$$
where $`H`$ is the standard Kaon Hamiltonian, and the “entanglement” matrix is given by
$$\delta H/_{\alpha \beta }=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 2\alpha & 2\beta \\ 0& 0& 2\beta & 2\gamma \end{array}\right).$$
Positivity of $`\rho `$ requires: $`\alpha ,\gamma >0,\alpha \gamma >\beta ^2`$. Notice that $`\alpha ,\beta ,\gamma `$ violate CPT in the sense of an induced microscopic time irreversibility of wald , as being associated with decoherence, but also they violate CP since they do not commute with the CP operator lopez : $`CP=\sigma _3\mathrm{cos}\theta +\sigma _2\mathrm{sin}\theta `$,$`[\delta H/_{\alpha \beta },CP]0`$.
An important remark is now in order. We should distinguish two types of CPTV: (i) CPTV within Quantum Mechanics: $`\delta M=m_{K^0}m_{\overline{K}^0}`$, $`\delta \mathrm{\Gamma }=\mathrm{\Gamma }_{K^0}\mathrm{\Gamma }_{\overline{K}^0}`$. This could be due to (spontaneous) Lorentz violation (c.f. above).
(ii) CPTV through decoherence $`\alpha ,\beta ,\gamma `$ (entanglement with QG ‘environment’, leading to modified evolution for $`\rho `$ and $`\$SS^{}`$).
The important point is that the two types of CPTV can be disentangled experimentally lopez . Experimentally, the best available bounds of $`\alpha ,\beta ,\gamma `$ parameters to date for single neutral Kaon states come from CPLEAR measurements cplear $`\alpha <4.0\times 10^{17}\mathrm{GeV},|\beta |<2.3.\times 10^{19}\mathrm{GeV},\gamma <3.7\times 10^{21}\mathrm{GeV}`$, which are not much different from theoretically expected values in the most optimistic of the scenaria, involving linear Planck-scale suppression $`\alpha ,\beta ,\gamma =O(\xi \frac{E^2}{M_P})`$. For more details we refer the reader to the literature lopez ; mavro\_leap03 .
*Neutral Meson factories: entangled meson states*
The above-described decoherence formalism can be used to derive the (non-unitary) evolution of the entangled products of the $`\varphi `$ or $`\mathrm{{\rm Y}}`$ decays in meson factories. The requirement of *complete positivity* of the respective density matrices imposes further restrictions among the decoherence parameters in that case, amounting to setting $`\beta =0`$, $`\alpha =\gamma 0`$ benatti .
An entirely novel observable, exclusive to a breakdown of CPT Violation through decoherence, in which case the CPT operator is not well defined, pertains to observations of the modifications of Einstein-Podolsky-Rosen entangled states of neutral mesons in meson factories ($`\varphi `$ or B-factories) bernabeu . These modifications concern the nature of the products of the decay of the neutral mesons in a factory on the two sides of the detector. For instance, for neutral kaons, if the CPT operator is a well defined operator, even if it does not commute with the Hamiltonian of the system, the products of the decay contain states $`K_LK_S`$ only. On the contrary, in the case of CPT breakdown through decoherence (ill defined CPT operator), one obtains in the final state, in addition to $`K_LK_S`$, also $`K_SK_S`$ and/or $`K_LK_L`$ states. In a similar manner, when this formalism of CPT breaking through decoherence is applied to B-systems, one observes that flavor tagging fails in B-factories as a result of such CPTV modifications bernabeu .
*Neutrinos*
Similarly to neutral mesons, one could have decoherening modifications of the oscillation probabilities for neutrinos. Due to lack of space we shall not describe the details here. We refer the interested reader in the literature mavromatosdecohe . We only mention that at present there is no experimental evidence in neutrino physics on quantum gravity decoherence effects, only very stringent bounds exist. Although, at first sight, it appears that neutrino anomalies, such as the LSND result lsnd indicating asymmetric rates for antineutrino oscillations as compared to neutrino ones, can be fit in the context of a decoherence model with mixed energy dependences, $`E`$ and $`1/E`$, and with different orders of the decoherence parameters between particle and antiparticle sectors barenboim , nevertheless recent spectral distortions from KamLand kamland indicate that standard oscillations are capable of explaining these distortions, excluding the order of decoherence in the antineutrino sector necessary to match the LSND result with the rest of the available neutrino data. This leaves us with the following bounds as far as decoherence coefficients $`\gamma `$ are concerned for neutrinos lizzi : in a parametrization $`\gamma \gamma _0(\frac{E}{\mathrm{GeV}})^n`$, with $`n=0,2,1`$, we have: (a) $`n=0`$, $`\gamma _0<3.5\times 10^{23}`$ GeV, (b) $`n=2`$, $`\gamma _0<0.9\times 10^{27}`$ GeV (compare with the neutral Kaon case above), and (c) $`n=1`$, $`\gamma _0<2\times 10^{21}`$ GeV. Very stringent limits on quantum gravity decoherence may be placed from future observations of high-energy neutrinos from extragalactic sources mavromatosdecohe (Supernovae, AGN), if, for instance, Quantum Gravity induces lepton number violation and/or flavor oscillations. We therefore see that neutrinos provide the most sensitive probe at present for tests of quantum gravity decoherence effects, provided the latter pertain to these particles.
*Disentangling Matter Effects from Genuine Quantum Gravity induced Decoherence effects*
However, there is an important aspect in all such decoherence searches, which should not be overlooked. This concerns “fake” decoherence effects as a result of the passage of the particle probe, e.g. neutrino, through matter media. Due to the apparent (“extrinsic”) breaking of CPT in such a case, one obtains fake violations of the symmetry, which have nothing to do with genuine microscopic gravity effects. Matter effects, for instance, include standard CPTV differences in neutrino oscillation probabilities of the form: $`P_{\nu _\alpha \nu _\beta }P_{\overline{\nu }_\beta \overline{\nu }_\alpha }`$, non-zero result for the so-called $`A_{CPT}^r`$ asymmetry to leading order, in the case of Kaons through a regenerator lopez , as well as uncertainties in the energy and oscillation length of neutrino beams, which result in damping factors in front of terms in the respective oscillation probabilities. Such factors are of similar nature to those induced by decoherence $`\gamma `$ coefficients olsson . Nevertheless, the energy and oscillation length dependence of the effects is different mavromatosdecohe ; barenboim , and allows disentanglement from genuine effects, if the latter are present and are of comparable order to the matter effects. Namely, quantum gravity effects are in general expected to increase with the energy of the neutrino probe, as a result of the fact that the higher the energy the stronger the back reaction onto the space time. This should be contrasted with ordinary matter effects, which are expected to decrease with the energy of the matter probe. Also, the damping factors in the case of neutrinos with uncertainty in their energy have a different oscillation length dependence, as compared with quantum gravity decoherence effects. Note that matter decoherence effects may be tiny olsson , with the appropriate coefficients of order smaller than $`\gamma _{\mathrm{fake}}<10^{24}`$ GeV, thereby leading to the temptation of identifying (incorrectly) such effects with genuine microscopic effects.
## 4 Conclusions
There are various ways by means of which Quantum Gravity (QG)-induced CPT breaking can occur, which are in principle independent of each other. For instance, quantum decoherence and Lorentz Violation are in general independent effects. There is no single figure of merit of CPT violation, and the associated sensitivity depends on the way CPT is broken, and on the relevant observable. The pertinent phenomenology is not simple.
Neutrino (astro)physics may provide some of the most stringent (to date) constraints on QG CPT Violation icecube . There are interesting theoretical issues on Quantum Gravity decoherence and neutrinos, which go as far as the origin of neutrino mass differences (which could even be due to space time foam effects, thus explaining their smallness), and their contributions to the Dark Energy barenboim2 .
Quantum-Gravity induced decoherence CPTV, in which the CPT generator is ill defined, may lead to interesting novel observables with high sensitivity in the near future, associated with entangled state modifications in neutral meson factories (B, $`\varphi `$).
Thus, it seems that a century after annus mirabilis, there is still a long way to go before an understanding of Quantum Gravity is achieved. But the challenge is there, and we think that there may be pleasant experimental surprises in the near future. This is why experimental searches of CPT violation and other quantum gravity effects along the lines presented here are worth pursuing.
The author would like to thank Prof. W. Oelert and the other organizers of the LEAP05 conference in Bonn/GSI for the invitation, and for creating such a stimulating meeting.
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# Critical potentials of the eigenvalues and eigenvalue gaps of Schrödinger operators
## 1. Introduction and Statement of main Results
Let $`M`$ be a compact connected Riemannian manifold of dimension $`d`$, possibly with nonempty boundary $`M`$, and let $`\mathrm{\Delta }`$ be its Laplace-Beltrami operator acting on functions with, in the case where $`M\mathrm{}`$, Dirichlet or Neumann boundary conditions. In all the sequel, as soon as the Neumann Laplacian will be considered, the boundary of $`M`$ will be assumed to be sufficiently regular (e.g. $`C^1`$, but weaker regularity assmptions may suffice, see ) in order to guarantee the compactness of the embedding $`H^1(M)L^2(M)`$ and, hence, the compactness of the resolvent of the Neumann Laplacian (note that it is well known, using standard arguments like in \[14, p.89\], that compactness results for Sobolev spaces on Euclidean domains remain valid in the Riemannian setting).
For any bounded real valued potential $`q`$ on $`M`$, the Schrödinger type operator $`\mathrm{\Delta }+q`$ has compact resolvent (see \[16, Theorem IV.3.17\] and observe that a bounded $`q`$ leads to a relatively compact operator with respect to $`\mathrm{\Delta }`$). Therefore, its spectrum consists of a nondecreasing and unbounded sequence of eigenvalues with finite multiplicities:
$$Spec(\mathrm{\Delta }+q)=\{\lambda _1(q)<\lambda _2(q)\lambda _3(q)\mathrm{}\lambda _i(q)\mathrm{}\}.$$
Each eigenvalue $`\lambda _i(q)`$ can be considered as a (continuous) function of the potential $`qL^{\mathrm{}}(M)`$ and there are both physical and mathematical motivations to study existence and properties of extremal potentials of the functionals $`\lambda _i`$ as well as of the differences, called gaps, between them. A very rich literature is devoted to the existence and the determination of maximizing or minimizing potentials for the eigenvalues (especially the fundamental one, $`\lambda _1`$) and the eigenvalue gaps (especially the first one, $`\lambda _2\lambda _1`$) under various constraints often motivated by physical considerations (see, for instance, and the references therein). Note that, since the function $`\lambda _i`$ commutes with constant translations, that is, $`\lambda _i(q+c)=\lambda _i(q)+c`$, such constraints are necessary.
Our aim in this paper is to investigate critical points, including ”local minimizers” and ”local maximizers”, of the eigenvalue functionals $`q\lambda _i(q)`$ and the eigenvalue gap functionals $`q\lambda _j(q)\lambda _i(q)`$, the potentials $`q`$ being subjected to the constraint that their mean value (or, equivalently, their integral) over $`M`$ is fixed. All along this paper, the mean value of an integrable function $`q`$ will be denoted $`\overline{q}`$, that is,
$$\overline{q}=\frac{1}{V(M)}_Mq𝑑v,$$
$`V(M)`$ and $`dv`$ being respectively the Riemannian volume and the Riemannian volume element of $`M`$.
Actually, most of the results below can be extended, modulo some slight changes, to the case where this constraint is replaced by the more general one
$$_MF(q)𝑑v=\text{constant},$$
where $`F:`$ is a continuous function such that $`F^{}(x)0`$ if $`x0`$, like $`F(x)=|x|^\alpha `$ or $`F(x)=x|x|^{\alpha 1}`$ with $`\alpha 1`$. However, for simplicity and clarity reasons, we preferred to focus only on the mean value constraint. Therefore, we fix a constant $`c`$ and consider the functionals
$$\lambda _i:qL_c^{\mathrm{}}(M)\lambda _i(q),$$
where $`L_c^{\mathrm{}}(M)=\left\{qL^{\mathrm{}}(M)\right|\overline{q}=c\}`$. The tangent space to $`L_c^{\mathrm{}}(M)`$ at any point $`q`$ is given by
$$L_{}^{\mathrm{}}(M):=\left\{uL^{\mathrm{}}(M)\right|_Mu𝑑v=0\}.$$
### 1.1. Critical potentials of the eigenvalue functionals
Since it is always nondegenerate, the first eigenvalue gives rise to a differentiable functional in the sense that, for any $`qL_c^{\mathrm{}}(M)`$ and any $`uL_{}^{\mathrm{}}(M)`$, the function $`t\lambda _1(q+tu)`$ is differentiable in $`t`$. A potential $`qL_c^{\mathrm{}}(M)`$ will be termed critical for this functional if $`\frac{d}{dt}\lambda _1(q+tu)|_{t=0}=0`$ for any $`uL_{}^{\mathrm{}}(M)`$.
In the case of empty boundary or of Neumann boundary conditions, the constant function 1 belongs to the domain of the operator $`\mathrm{\Delta }+q`$ and one obtains, as a consequence of the min-max principle, that the constant potential $`c`$ is a global maximizer of $`\lambda _1`$ over $`L_c^{\mathrm{}}(M)`$ (see also and ). Constant potential $`c`$ is actually the only critical one for $`\lambda _1`$. On the other hand, under Dirichlet boundary conditions, the functional $`\lambda _1`$ admits no critical potentials in $`L_c^{\mathrm{}}(M)`$. Indeed, we have the following
###### Theorem 1.1.
1. Assume that either $`M=\mathrm{}`$ or $`M\mathrm{}`$ and Neumann boundary conditions are imposed. Then, for any potential $`q`$ in $`L_c^{\mathrm{}}(M)`$, we have
$$\lambda _1(q)\lambda _1(c)=c,$$
where the equality holds if and only if $`q=c`$. Moreover, the constant potential $`c`$ is the only critical one of the functional $`\lambda _1`$ over $`L_c^{\mathrm{}}(M)`$.
2. Assume that $`M\mathrm{}`$ and that Zero Dirichlet boundary conditions are imposed. Then the functional $`\lambda _1`$ does not admit any critical potential in $`L_c^{\mathrm{}}(M)`$.
Higher eigenvalues are continuous but not differentiable in general. Nevertheless, perturbation theory enables us to prove that, for any function $`uL^{\mathrm{}}(M)`$, the function $`t\lambda _i(q+tu)`$ admits left and right derivatives at $`t=0`$ (see section 2.2). A generalized notion of criticality can be naturally defined as follows :
###### Definition 1.1.
A potential $`q`$ is said to be critical for the functional $`\lambda _i`$ if, for any $`uL_{}^{\mathrm{}}(M)`$, the left and right derivatives of $`t\lambda _i(q+tu)`$ at $`t=0`$ have opposite signs, that is
$$\frac{d}{dt}\lambda _i(q+tu)|_{t=0^+}\times \frac{d}{dt}\lambda _i(q+tu)|_{t=0^{}}0.$$
It is immediate to check that $`q`$ is critical for $`\lambda _i`$ if and only if, for any $`uL_{}^{\mathrm{}}(M)`$, one of the two following inequalities holds :
$$\lambda _i(q+tu)\lambda _i(q)+o(t)\text{as}t0$$
or
$$\lambda _i(q+tu)\lambda _i(q)+o(t)\text{as}t0.$$
In all the sequel, we will denote by $`E_i(q)`$ the eigenspace corresponding to the $`i`$-th eigenvalue $`\lambda _i(q)`$ whose dimension coincides with the number of indices $`j`$ such that $`\lambda _j(q)=\lambda _i(q)`$.
As for the first eigenvalue, the functionals $`\lambda _i`$, $`i2`$, admit no critical potentials under Dirichlet boundary conditions.
###### Theorem 1.2.
Assume that $`M\mathrm{}`$ and that Zero Dirichlet boundary conditions are imposed. Then, $`i^{}`$, the functional $`\lambda _i`$ does not admit any critical potential in $`L_c^{\mathrm{}}(M)`$.
Under the two remaining boundary conditions, the following theorem gives a necessary condition for a potential $`q`$ to be critical for the functional $`\lambda _i`$. This condition is also sufficient for the indices $`i`$ such that $`\lambda _i(q)>\lambda _{i1}(q)`$ or $`\lambda _i(q)<\lambda _{i+1}(q)`$, which means that $`\lambda _i(q)`$ is the first one or the last one in a cluster of equal eigenvalues.
###### Theorem 1.3.
Assume that either $`M=\mathrm{}`$ or $`M\mathrm{}`$ and Neumann boundary conditions are imposed. Let $`i`$ be a positive integer.
If $`qL_c^{\mathrm{}}(M)`$ is a critical potential of the functional $`\lambda _i`$, then $`q`$ is smooth and there exists a finite family of eigenfunctions $`f_1,\mathrm{},f_k`$ in $`E_i(q)`$ such that $`_{1jk}f_j^2=1`$.
Reciprocally, if $`\lambda _i(q)>\lambda _{i1}(q)`$ or $`\lambda _i(q)<\lambda _{i+1}(q)`$, and if there exists a family of eigenfunctions $`f_1,\mathrm{},f_kE_i(q)`$ such that $`_{1jk}f_j^2=1`$, then $`q`$ is a critical potential of the functional $`\lambda _i`$.
Note that the identity $`_{1jk}f_j^2=1`$, with $`f_1,\mathrm{},f_kE_i(q)`$, immediately implies another one (that we obtain from $`\mathrm{\Delta }_{1jk}f_j^2=0`$):
$$q=\lambda _i(q)\underset{1jk}{}|f_j|^2,$$
from which we can deduce the smoothness of $`q`$.
###### Remark 1.1.
1. The identity $`_{1jk}f_j^2=1`$ with $`\mathrm{\Delta }f_j+qf_j=\lambda _i(q)f_j`$, means that the map $`f=(f_1,\mathrm{},f_k)`$ from $`M`$ to the Euclidean sphere $`𝕊^{k1}`$ is harmonic with energy density $`|f|^2=\lambda _i(q)q`$ (see ). Hence, a necessary (and sometime sufficient) condition for a potential $`q`$ to be critical for the functional $`\lambda _i`$ is that the function $`\lambda _i(q)q`$ is the energy density of a harmonic map from $`M`$ to a Euclidean sphere.
2. If one replaces the constraint on the mean value $`\frac{1}{V(M)}_Mq𝑑v=c`$ by the general constraint $`_MF(q)𝑑v=c,`$ then the necessary and sufficient condition $`_{1jk}f_j^2=1`$ of Theorem 1.3 becomes (even under Dirichlet boundary conditions) $`_{1jk}f_j^2=F^{}(q)`$. In particular, $`q`$ is a critical potential of the functional $`\lambda _1`$ if and only if $`F^{}(q)0`$ and $`F^{}(q)^{\frac{1}{2}}`$ is a first eigenfunction of $`\mathrm{\Delta }+q`$, see for a discussion of the case $`F(q)=|q|^\alpha `$.
Under each one of the boundary conditions we consider, a constant function can never be an eigenfunction associated to an eigenvalue $`\lambda _i(q)`$ with $`i2`$. Hence, an immediate consequence of Theorem 1.3 is the following
###### Corollary 1.1.
If $`qL_c^{\mathrm{}}(M)`$ is a critical potential of the functional $`\lambda _i`$ with $`i2`$, then the eigenvalue $`\lambda _i(q)`$ is degenerate, that is $`\lambda _i(q)=\lambda _{i1}(q)`$ or $`\lambda _i(q)=\lambda _{i+1}(q)`$
If $`\{f_1,\mathrm{},f_k\}`$ is an $`L^2`$-orthonormal basis of $`E_i(\mathrm{\Delta })`$, then the function $`_{1jk}f_j^2`$ is invariant under the isometry group of $`M`$. Indeed, for any isometry $`\rho `$ of $`M`$, $`\{f_1\rho ,\mathrm{},f_k\rho \}`$ is also an orthonomal basis of $`E_i(\mathrm{\Delta })`$ and then, there exists a matrix $`AO(d)`$ such that $`(f_1\rho ,\mathrm{},f_d\rho )=A.(f_1,\mathrm{},f_d)`$. In particular, if $`M`$ is homogeneous, that is, the isometry group acts transitively on $`M`$, then $`_{1jk}f_j^2`$ would be constant. Another consequence of Theorem 1.3 is then the following
###### Corollary 1.2.
If $`M`$ is homogeneous, then constant potentials are critical for all the functionals $`\lambda _i`$ such that $`\lambda _i(\mathrm{\Delta })<\lambda _{i+1}(\mathrm{\Delta })`$ or $`\lambda _i(\mathrm{\Delta })>\lambda _{i1}(\mathrm{\Delta })`$.
Recall that Euclidean spheres, projective spaces and flat tori are examples of homogeneous Riemannian spaces.
A potential $`qL_c^{\mathrm{}}(M)`$ is said to be a local minimizer (resp. local maximizer) of the functional $`\lambda _i`$ (in a weak sense) if, for any $`uL_{}^{\mathrm{}}(M)`$, the function $`t\lambda _i(q+tu)`$ admits a local minimum (resp. maximum) at $`t=0`$. The result of Corollary 1.1 takes the following more precise form in the case of a local minimizer or maximizer.
###### Theorem 1.4.
Let $`qL_c^{\mathrm{}}(M)`$ and $`i2`$.
1. If $`q`$ is a local minimizer of the functional $`\lambda _i`$, then $`\lambda _i(q)=\lambda _{i1}(q)`$.
2. If $`q`$ is a local maximizer of the functional $`\lambda _i`$, then $`\lambda _i(q)=\lambda _{i+1}(q)`$.
Since the first eigenvalue is simple, we always have $`\lambda _2(q)>\lambda _1(q)`$. The previous results, applied to the functional $`\lambda _2`$ can be summarized as follows.
###### Corollary 1.3.
Assume that either $`M=\mathrm{}`$ or $`M\mathrm{}`$ and Neumann boundary conditions are imposed. A potential $`qL_c^{\mathrm{}}(M)`$ is critical for the functional $`\lambda _2`$ if and only if, $`q`$ is smooth, $`\lambda _2(q)=\lambda _3(q)`$ and there exist eigenfunctions $`f_1,\mathrm{},f_k`$ in $`E_2(q)`$ such that $`_{1jk}f_j^2=1`$.
Moreover, the functional $`\lambda _2`$ admits no local minimizers in $`L_c^{\mathrm{}}(M)`$.
In , Ilias and the first author have proved that, under some hypotheses on $`M`$, satisfied in particular by compact rank-one symmetric spaces, irreducible homogeneous Riemannian spaces and some flat tori, the constant potential $`c`$ is a global maximizer of $`\lambda _2`$ over $`L_c^{\mathrm{}}(M)`$. In , they studied the critical points of $`\lambda _i`$ considered as a functional on the set of Riemannian metrics of fixed volume on $`M`$.
### 1.2. Critical potentials of the eigenvalue gaps functionals
We consider now the eigenvalue gaps functionals $`qG_{ij}(q)=\lambda _j(q)\lambda _i(q)`$, where $`i`$ and $`j`$ are two distinct positive integers, and define their critical potentials as in Definition 1.1. These functionals are invariant under translations, that is $`G_{ij}(q+c)=G_{ij}(q)`$. Therefore, critical potentials of $`G_{ij}`$ with respect to fixed mean value deformations are also critical with respect to arbitrary deformations.
###### Theorem 1.5.
If $`qL_c^{\mathrm{}}(M)`$ is a critical potential of the gap functional $`G_{ij}=\lambda _j\lambda _i`$, then there exist a finite family of eigenfunctions $`f_1,\mathrm{},f_k`$ in $`E_i(q)`$ and a finite family of eigenfunctions $`g_1,\mathrm{},g_l`$ in $`E_j(q)`$, such that $`_{1pk}f_p^2=_{1pl}g_p^2`$.
Reciprocally, if $`\lambda _i(q)<\lambda _{i+1}(q)`$ and $`\lambda _j(q)>\lambda _{j1}(q)`$, and if there exist $`f_1,\mathrm{},f_k`$ in $`E_i(q)`$ and $`g_1,\mathrm{},g_l`$ in $`E_j(q)`$ such that $`_{1pk}f_p^2=_{1pl}g_p^2`$, then $`q`$ is a critical potential of $`G_{ij}`$.
In the particular case of the gap between two consecutive eigenvalues, we have the following
###### Corollary 1.4.
A potential $`qL_c^{\mathrm{}}(M)`$ is critical for the gap functional $`G_{i,i+1}=\lambda _{i+1}\lambda _i`$ if and only if, either $`\lambda _{i+1}(q)=\lambda _i(q)`$, or there exist a family of eigenfunctions $`f_1,\mathrm{},f_k`$ in $`E_i(q)`$ and a family of eigenfunctions $`g_1,\mathrm{},g_l`$ in $`E_{i+1}(q)`$, such that $`_{1pk}f_p^2=_{1pl}g_p^2`$.
###### Remark 1.2.
The characterization of critical potentials of $`G_{ij}`$ given in Theorem 1.5 remains valid under the constraint $`_MF(q)𝑑v=c`$.
An immediate consequence of Theorem 1.5 is the following
###### Corollary 1.5.
Let $`qL_c^{\mathrm{}}(M)`$ be a critical potential of the gap functional $`G_{ij}=\lambda _j\lambda _i`$. If $`\lambda _i(q)`$ (resp. $`\lambda _j(q)`$) is nondegenerate, then $`\lambda _j(q)`$ (resp. $`\lambda _i(q)`$) is degenerate.
The following is an immediate consequence of the discussion above concerning homogeneous Riemannian manifolds.
###### Corollary 1.6.
If M is a homogeneous Riemannian manifold, then, for any positive integer $`i`$, constant potentials are critical points of the gap functional $`G_{i,i+1}=\lambda _{i+1}\lambda _i`$.
Potentials $`q`$ such that $`\lambda _{i+1}(q)=\lambda _i(q)`$ are of course global minimizers of the gap functional $`G_{i,i+1}`$. These potentials are also the only local minimizers of $`G_{i,i+1}`$. Indeed, we have the following
###### Theorem 1.6.
If $`qL_c^{\mathrm{}}(M)`$ is a local minimizer of the gap functional $`G_{ij}=\lambda _j\lambda _i`$, then, either $`\lambda _i(q)=\lambda _{i+1}(q)`$, or $`\lambda _j(q)=\lambda _{j1}(q)`$. If $`q`$ is a local maximizer of $`G_{ij}`$, then, either $`\lambda _i(q)=\lambda _{i1}(q)`$, or $`\lambda _j(q)=\lambda _{j+1}(q)`$.
In particular, $`q`$ is a local minimizer of the gap functional $`G_{i,i+1}=\lambda _{i+1}\lambda _i`$ if and only if $`G_{i,i+1}(q)=0`$.
Finally, let us apply the results of this section to the first gap $`G_{1,2}`$.
###### Corollary 1.7.
A potential $`qL_c^{\mathrm{}}(M)`$ is critical for the gap functional $`G_{1,2}=\lambda _2\lambda _1`$ if and only if $`\lambda _2(q)`$ is degenerate and there exists a family of eigenfunctions $`g_1,\mathrm{},g_l`$ in $`E_2(q)`$ such that $`_{1jl}g_j^2=f^2`$, where $`f`$ is a basis of $`E_1(q)`$.
The functional $`G_{1,2}`$ does not admit any local minimizer in $`L_c^{\mathrm{}}(M)`$.
The authors wish to thank the referee for his valuable remarks.
## 2. Proof of Results
### 2.1. Variation Formula and proof of Theorem 1.1
Given on $`M`$ a potential $`q`$ and a function $`uL^{\mathrm{}}(M)`$, we consider the family of operators $`\mathrm{\Delta }+q+tu`$. Suppose that $`\mathrm{\Lambda }(t)`$ is a differentiable family of eigenvalues of $`\mathrm{\Delta }+q+tu`$ and that $`f_t`$ is a differentiable family of corresponding normalized eigenfunctions, that is, $`t`$,
$$(\mathrm{\Delta }+q+tu)f_t=\mathrm{\Lambda }(t)f_t,$$
and
$$_Mf_t^2𝑑v=1,$$
with $`f_t|_M=0`$ or $`\frac{f_t}{\nu }|_M=0`$ if $`M\mathrm{}`$. The following formula, giving the derivative of $`\mathrm{\Lambda }`$, is already known at least in the case of Euclidean domains with Dirichlet boundary conditions.
###### Proposition 2.1.
$$\mathrm{\Lambda }^{}(0)=_Muf_0^2𝑑v.$$
###### Proof.
First, we have, for all $`t`$,
$$\mathrm{\Lambda }(t)=\mathrm{\Lambda }(t)_M(f_t)^2𝑑v=_Mf_t(\mathrm{\Delta }+q+tu)f_t𝑑v.$$
Differentiating at $`t=0`$, we get
$$\mathrm{\Lambda }^{}(0)=\frac{d}{dt}\left(_Mf_t(\mathrm{\Delta }+q)f_t𝑑v+t_Mu(f_t)^2𝑑v\right)|_{t=0}.$$
Now, noticing that the function $`\frac{d}{dt}f_t|_{t=0}`$ satisfies the same boundary conditions as $`f_0`$ in case $`M\mathrm{}`$, and using integration by parts, we obtain
$`{\displaystyle \frac{d}{dt}}{\displaystyle _M}f_t(\mathrm{\Delta }+q)f_t𝑑v|_{t=0}`$ $`=`$ $`2{\displaystyle _M}(\mathrm{\Delta }+q)f_0{\displaystyle \frac{d}{dt}}f_t|_{t=0}dv`$
$`=`$ $`2\mathrm{\Lambda }(0){\displaystyle _M}f_0{\displaystyle \frac{d}{dt}}f_t|_{t=0}dv`$
$`=`$ $`\mathrm{\Lambda }(0){\displaystyle \frac{d}{dt}}{\displaystyle _M}f_t^2𝑑v|_{t=0}=0.`$
On the other hand, we have
$`{\displaystyle \frac{d}{dt}}\left(t{\displaystyle _M}uf_t^2𝑑v\right)|_{t=0}`$ $`=`$ $`{\displaystyle _M}uf_0^2𝑑v+\left(t{\displaystyle _M}u{\displaystyle \frac{d}{dt}}f_t^2𝑑v\right)|_{t=0}`$
$`=`$ $`{\displaystyle _M}uf_0^2v_g.`$
Finally, $`\mathrm{\Lambda }^{}(0)=_Muf_0^2𝑑v`$. ∎
###### Proof.
(of Theorem 1.1.) (i) First, let us show that constant potentials are maximizing for $`\lambda _1`$. Indeed, let $`c`$ be a constant potential and let $`q`$ be an arbitrary one in $`L_c^{\mathrm{}}(M)`$. From the variational characterization of $`\lambda _1(\mathrm{\Delta }+q)`$ in the case $`M=\mathrm{}`$ as well as in the case of Neumann boundary conditions, we get
$`\lambda _1(\mathrm{\Delta }+q)`$ $`=`$ $`\underset{fH^1(M)}{inf}{\displaystyle \frac{_M(|f|^2+qf^2)𝑑v}{f_{L^2(M)}^2}}`$
$``$ $`{\displaystyle \frac{_M(|1|^2+q1^2)𝑑v}{1_{L^2(M)}^2}}={\displaystyle \frac{_Mq𝑑v}{V(M)}}=c.`$
Hence, $`\lambda _1(q)\lambda _1(c)`$ and the constant potential $`c`$ maximizes the functional $`\lambda _1`$ on $`L_c^{\mathrm{}}(M)`$. In particular, constant potentials are critical for this functional.
Now, suppose that $`qL_c^{\mathrm{}}(M)`$ is a critical potential for $`\lambda _1`$. For any $`uL_{}^{\mathrm{}}(M)`$, we consider a differentiable family $`f_t`$ of normalized eigenfunctions corresponding to the first eigenvalue of $`(\mathrm{\Delta }+q+tu)`$ and apply the variation formula above to obtain:
$$\frac{d}{dt}\lambda _1(q+tu)|_{t=0}=_Muf_0^2𝑑v.$$
Hence, $`_Muf_0^2𝑑v=0`$ for any $`uL_{}^{\mathrm{}}(M)`$, which implies that $`f_0`$ is constant on $`M`$. Since $`(\mathrm{\Delta }+q)f_0=qf_0=\lambda _1(q)f_0`$, the potential $`q`$ must be constant on $`M`$.
(ii) Let $`f_0`$ be the first nonnegative Dirichlet eigenfunction of $`\mathrm{\Delta }+q`$ satisfying $`_Mf_0^2𝑑v=1`$. The function $`u=V(M)f_0^21`$ belongs to $`L_{}^{\mathrm{}}(M)`$ and we have
$$\frac{d}{dt}\lambda _1(q+tu)|_{t=0}=_Muf_0^2𝑑v=V(M)_Mf_0^4𝑑v1>0,$$
where the last inequality comes from Cauchy-Schwarz inequality and the fact that $`f_0`$ is not constant (recall that $`f_0|_M=0`$). Therefore, the potential $`q`$ is not critical for $`\lambda _1`$. ∎
### 2.2. Characterization of critical potentials
Let $`i`$ be a positive integer and let $`m1`$ be the dimension of the eigenspace $`E_i(q)`$ associated to the eigenvalue $`\lambda _i(q)`$. For any function $`uL_{}^{\mathrm{}}(M)`$, perturbation theory of unbounded self-adjoint operators (see for instance Kato’s book ) that we apply to the one parameter family of operators $`\mathrm{\Delta }+q+tu`$, tells us that, there exists a family of $`m`$ eigenfunctions $`f_{1,t},\mathrm{},f_{m,t}`$ associated with a family of $`m`$ (non ordered) eigenvalues $`\mathrm{\Lambda }_1(t),\mathrm{},\mathrm{\Lambda }_m(t)`$ of $`\mathrm{\Delta }+q+tu`$, all depending analytically in $`t`$ in some interval $`(\epsilon ,\epsilon )`$, and satisfying
* $`\mathrm{\Lambda }_1(0)=\mathrm{}=\mathrm{\Lambda }_m(0)=\lambda _i(q)`$,
* $`t(\epsilon ,\epsilon )`$, the $`m`$ functions $`f_{1,t},\mathrm{},f_{m,t}`$ are orthonormal in $`L^2(M)`$.
From this, one can easily deduce the existence of two integers $`km`$ and $`lm`$, and a small $`\delta >0`$ such that
$$\lambda _i(q+tu)=\{\begin{array}{c}\mathrm{\Lambda }_k(t)\text{if}t(\delta ,0)\hfill \\ \\ \mathrm{\Lambda }_l(t)\text{if}t(0,\delta ).\hfill \end{array}$$
Hence, the function $`t\lambda _i(q+tu)`$ admits a left sided and a right sided derivatives at $`t=0`$ with
$$\frac{d}{dt}\lambda _i(q+tu)|_{t=0^{}}=\mathrm{\Lambda }_k^{}(0)=_Muf_{k,0}^2𝑑v$$
and
$$\frac{d}{dt}\lambda _i(q+tu)|_{t=0^+}=\mathrm{\Lambda }_l^{}(0)=_Muf_{l,0}^2𝑑v.$$
To any function $`uL_{}^{\mathrm{}}(M)`$ and any integer $`i`$, we associate the quadratic form $`Q_u^i`$ on $`E_i(q)`$ defined by
$$Q_u^i(f)=_Muf^2𝑑v.$$
The corresponding symmetric linear transformation $`L_u^i:E_i(q)E_i(q)`$ is given by
$$L_u^i(f)=P_i(uf),$$
where $`P_i:L^2(M)E_i(q)`$ is the orthogonal projection of $`L^2(M)`$ onto $`E_i(q)`$.
It follows immediately that
###### Proposition 2.2.
If the potential $`q`$ is critical for the functional $`\lambda _i`$, then, $`uL_{}^{\mathrm{}}(M)`$, the quadratic form $`Q_u^i(f)=_Muf^2𝑑v`$ is indefinite on the eigenspace $`E_i(q)`$.
The following lemma enables us to establish a converse to this proposition.
###### Lemma 2.1.
$`k,lm`$, we have
$$_Muf_{k,0}f_{l,0}𝑑v=\{\begin{array}{c}0\text{if}kl\hfill \\ \\ \mathrm{\Lambda }_k^{}(0)\text{if}k=l.\hfill \end{array}$$
In other words, $`\mathrm{\Lambda }_1^{}(0),\mathrm{},\mathrm{\Lambda }_m^{}(0)`$ are the eigenvalues of the symmetric linear transformation $`L_u^i:E_i(q)E_i(q)`$ and the functions $`f_{1,0},\mathrm{},f_{m,0}`$ constitute an orthonormal eigenbasis of $`L_u^i`$.
###### Proof.
Differentiating at $`t=0`$ the equality $`(\mathrm{\Delta }+q+tu)f_{k,t}=\mathrm{\Lambda }_k(t)f_{k,t}`$, we obtain
$$uf_{k,0}+(\mathrm{\Delta }+q)\frac{d}{dt}f_{k,t}|_{t=0}=\mathrm{\Lambda }_k^{}(0)f_{k,0}+\mathrm{\Lambda }_k(0)\frac{d}{dt}f_{k,t}|_{t=0},$$
and then,
$`{\displaystyle _M}uf_{k,0}f_{l,0}𝑑v=\mathrm{\Lambda }_k^{}(0){\displaystyle _M}f_{k,0}f_{l,0}𝑑v`$ $`+`$ $`\mathrm{\Lambda }_k(0){\displaystyle _M}f_{l,0}{\displaystyle \frac{d}{dt}}f_{k,t}|_{t=0}dv`$
$``$ $`{\displaystyle _M}f_{l,0}(\mathrm{\Delta }+q){\displaystyle \frac{d}{dt}}f_{k,t}|_{t=0}dv.`$
Integration by parts gives, after noticing that $`\mathrm{\Lambda }_k(0)=\mathrm{\Lambda }_l(0)=\lambda _i(q)`$ and that the functions $`\frac{d}{dt}f_{k,t}|_{t=0}`$ satisfy the considered boundary conditions,
$`{\displaystyle _M}f_{l,0}(\mathrm{\Delta }+q){\displaystyle \frac{d}{dt}}f_{k,t}|_{t=0}dv`$ $`=`$ $`{\displaystyle _M}{\displaystyle \frac{d}{dt}}f_{k,t}|_{t=0}(\mathrm{\Delta }+q)f_{l,0}dv`$
$`=`$ $`\mathrm{\Lambda }_k(0){\displaystyle _M}f_{l,0}{\displaystyle \frac{d}{dt}}f_{k,t}|_{t=0}dv,`$
and finally,
$$_Muf_{k,0}f_{l,0}𝑑v=\mathrm{\Lambda }_k^{}(0)_Mf_{k,0}f_{l,0}𝑑v=\mathrm{\Lambda }_k^{}(0)\delta _{kl}.$$
###### Proposition 2.3.
Assume that $`\lambda _i(q)>\lambda _{i1}(q)`$ or $`\lambda _i(q)<\lambda _{i+1}(q)`$. Then the following conditions are equivalent:
* the potential $`q`$ is critical for $`\lambda _i`$
* $`uL_{}^{\mathrm{}}(M)`$, the quadratic form $`Q_u^i(f)=_Muf^2𝑑v`$ is indefinite on the eigenspace $`E_i(q)`$.
* $`uL_{}^{\mathrm{}}(M)`$, the linear transformation $`L_u^i`$ admits eigenvalues of both signs.
###### Proof.
Conditions (ii) and (iii) are clearly equivalent and the fact that (i) implies (ii) was established in Proposition 2.2. Let us show that (iii) implies (i). Assume that $`\lambda _i(q)>\lambda _{i1}(q)`$ and let $`uL_{}^{\mathrm{}}(M)`$ and $`\mathrm{\Lambda }_1(t),\mathrm{},\mathrm{\Lambda }_m(t)`$ be as above. For small $`t`$, we will have, for continuity reasons, $`km`$, $`\mathrm{\Lambda }_k(t)>\lambda _{i1}(q+tu)`$ and then, $`\lambda _i(q+tu)\mathrm{\Lambda }_k(t)`$. Since $`\lambda _i(q+tu)\{\mathrm{\Lambda }_1(t),\mathrm{},\mathrm{\Lambda }_m(t)\}`$, we get
$$\lambda _i(q+tu)=\underset{km}{\mathrm{min}}\mathrm{\Lambda }_k(t).$$
It follows that
$$\frac{d}{dt}\lambda _i(q+tu)|_{t=0^{}}=\underset{km}{\mathrm{max}}\mathrm{\Lambda }_k^{}(0)$$
and
$$\frac{d}{dt}\lambda _i(q+tu)|_{t=0^+}=\underset{km}{\mathrm{min}}\mathrm{\Lambda }_k^{}(0).$$
Thanks to Lemma 2.1, Condition (iii) implies that $`\mathrm{min}_{km}\mathrm{\Lambda }_k^{}(0)0\mathrm{max}_{km}\mathrm{\Lambda }_k^{}(0)`$ which implies the criticality of $`q`$.
The case $`\lambda _i(q)<\lambda _{i+1}(q)`$ can be treated in a similar manner. ∎
### 2.3. Proof of Theorems 1.2 and 1.3
Let $`q`$ be a potential in $`L_c^{\mathrm{}}(M)`$. To prove Theorem 1.2 we first notice that, since $`f|_M=0`$ for any $`fE_i(q)`$, the constant function 1 does not belong to the vector space $`F`$ generated in $`L^2(M)`$ by $`\{f^2;fE_i(q)\}`$. Hence, there exists a function $`u`$ orthogonal to $`F`$ and such that $`u,1_{L^2(M)}<0`$. The function $`u_0=u\overline{u}`$ belongs to $`L_{}^{\mathrm{}}(M)`$ and the quadratic form $`Q_{u_0}^i(f)=_Mu_0f^2𝑑v=\overline{u}f_{L^2(M)}^2`$ is positive definite on $`E_i(q)`$. Hence, the potential $`q`$ is not critical for $`\lambda _i`$ (see Proposition 2.2).
The proof of Theorem 1.3 follows directly from the two propositions above and the following lemma.
###### Lemma 2.2.
Let $`i`$ be a positive integer. The two following conditions are equivalent:
* $`uL_{}^{\mathrm{}}(M)`$, the quadratic form $`Q_u^i(f)=_Muf^2𝑑v`$ is indefinite on the eigenspace $`E_i(q)`$.
* there exists a family of eigenfunctions $`f_1,\mathrm{},f_k`$ in $`E_i(q)`$ such that $`_{1jk}f_j^2=1`$.
###### Proof.
To see that (i) implies (ii) we introduce the convex cone $`C`$ generated in $`L^2(M)`$ by the set $`\{f^2;fE_i(q)\}`$, that is $`C=\{_{jJ}f_j^2;f_jE_i(q),J,J\text{ is finite}\}`$. Condition (ii) is then equivalent to the fact that the constant function 1 belongs to $`C`$. Let us suppose, for a contradiction, that $`1C`$. Then, applying classical separation theorems (in the finite dimensional vector subspace of $`L^2(M)`$ generated by $`\{f^2;fE_i(q)\}`$ and $`1`$, see ), we prove the existence of a function $`uL^2(M)`$ such that $`\overline{u}=\frac{1}{V(M)}_Mu1𝑑v<0`$ and $`_Muf^2𝑑v0`$ for any $`fC`$. Hence, the function $`u_0=u\overline{u}`$ belongs to $`L_{}^{\mathrm{}}(M)`$ and satisfies, $`fE_i(q)`$,
$$Q_{u_0}^i(f)=_Muf^2𝑑v\frac{1}{V(M)}_Mu𝑑v_Mf^2𝑑v\overline{u}f_{L^2(M)}^2.$$
The quadratic form $`Q_{u_0}^i`$ is then positive definite which contradicts (i) (see Proposition 2.2).
Reciprocally, the existence of $`f_1,\mathrm{},f_k`$ in $`E_i(q)`$ satisfying $`_{1jk}f_j^2=1`$ implies that, $`uL_{}^{\mathrm{}}(M)`$,
$$\underset{jk}{}Q_u^i(f_j)=\underset{jk}{}_Muf_j^2𝑑v=_Mu=0,$$
which implies that the quadratic form $`Q_u^i`$ is indefinite on $`E_i(q)`$. ∎
Finally, let us check that the condition $`_{1jk}f_j^2=1`$, with $`f_jE_i(q)`$, implies that $`q`$ is smooth. Indeed, since $`qL^{\mathrm{}}(M)`$, we have, for any eigenfunction $`fE_i(q)`$, $`\mathrm{\Delta }fL^2(M)`$ and then, $`fH^{2,2}(M)`$. Using standard regularity theory and Sobolev embeddings (see, for instance, ), we obtain by an elementary iteration, that $`fH^{2,p}(M)`$ for some $`p>n`$, and, then, $`fC^1(M)`$. From $`_{1jk}f_j^2=1`$ and $`\mathrm{\Delta }_{1jk}f_j^2=0`$, we get
$$q=\lambda _i(q)\underset{1jk}{}|f_j|^2,$$
which implies that $`q`$ is continuous. Again, elliptic regularity theory tells us that the eigenfunctions of $`\mathrm{\Delta }+q`$ are actually smooth, and, hence, $`q`$ is smooth.
### 2.4. Proof of Theorem 1.4
Assume that the potential $`q`$ is a local minimizer of the functional $`\lambda _i`$ on $`L_c^{\mathrm{}}(M)`$ and let us suppose for a contradiction that $`\lambda _i(q)>\lambda _{i1}(q)`$. Let $`u`$ be a function in $`L_{}^{\mathrm{}}(M)`$ and let $`\mathrm{\Lambda }_1(t),\mathrm{},\mathrm{\Lambda }_m(t)`$ be a family of $`m`$ eigenvalues of $`\mathrm{\Delta }+q+tu`$, where $`m`$ is the multiplicity of $`\lambda _i(q)`$, depending analytically in $`t`$ and such that $`\mathrm{\Lambda }_1(0)=\mathrm{}=\mathrm{\Lambda }_m(0)=\lambda _i(q)`$. For continuity reasons, we have, for sufficiently small $`t`$ and any $`km`$, $`\mathrm{\Lambda }_k(t)>\lambda _{i1}(q+tu)`$. Hence, $`km`$ and $`t`$ sufficiently small,
$$\mathrm{\Lambda }_k(t)\lambda _i(q+tu)\lambda _i(q)=\mathrm{\Lambda }_k(0).$$
Consequently, $`km`$, $`\mathrm{\Lambda }_k^{}(0)=0`$. Applying Lemma 2.1 above we deduce that the symmetric linear transformation $`L_u^i`$ and then the quadratic form $`Q_u^i`$ is identically zero on the eigenspace $`E_i(q)`$. Therefore, $`uL_{}^{\mathrm{}}(M)`$ and $`fE_i(q)`$, we have $`_Muf^2v_g=0`$. In conclusion, $`fE_i(q)`$, $`f`$ is constant on $`M`$ which is impossible for $`i2`$. The same arguments work to prove Assertion (ii).
### 2.5. Proof of Theorem 1.5
Let $`q`$ be a potential and let $`i`$ and $`j`$ be two distinct positive integers such that $`\lambda _i(q)\lambda _j(q)`$. We denote by $`m`$ (resp. $`n`$) the dimension of the eigenspace $`E_i(q)`$ (resp. $`E_j(q)`$). Given a function $`u`$ in $`L_{}^{\mathrm{}}(M)`$, we consider, as above, $`m`$ (resp. $`n`$) $`L^2(M)`$-orthonormal families of eigenfunctions $`f_{1,t},\mathrm{},f_{m,t}`$ (resp. $`g_{1,t},\mathrm{},g_{n,t}`$) associated with $`m`$ (resp. $`n`$) families of eigenvalues $`\mathrm{\Lambda }_1(t),\mathrm{},\mathrm{\Lambda }_m(t)`$ (resp. $`\mathrm{\Gamma }_1(t),\mathrm{},\mathrm{\Gamma }_n(t)`$) of $`\mathrm{\Delta }+q+tu`$, all depending analytically in $`t(\epsilon ,\epsilon )`$, such that $`\mathrm{\Lambda }_1(0)=\mathrm{}=\mathrm{\Lambda }_m(0)=\lambda _i(q)`$ (resp. $`\mathrm{\Gamma }_1(0)=\mathrm{}=\mathrm{\Gamma }_n(0)=\lambda _j(q)`$). Hence, there exist four integers $`km`$, $`k^{}m`$, $`ln`$ and $`l^{}n`$, such that
$`{\displaystyle \frac{d}{dt}}(\lambda _j\lambda _i)(q+tu)|_{t=0^{}}`$ $`=`$ $`\mathrm{\Gamma }_l^{}(0)\mathrm{\Lambda }_k^{}(0)`$
$`=`$ $`{\displaystyle _M}u(g_{l,0}^2f_{k,0}^2)𝑑v`$
and
$`{\displaystyle \frac{d}{dt}}(\lambda _j\lambda _i)(q+tu)|_{t=0^+}`$ $`=`$ $`\mathrm{\Gamma }_l^{}^{}(0)\mathrm{\Lambda }_k^{}^{}(0)`$
$`=`$ $`{\displaystyle _M}u(g_{l^{},0}^2f_{k^{},0}^2)𝑑v.`$
Recall that (Lemma 2.1) the eigenfunctions $`f_{1,0},\mathrm{},f_{m,0}`$ (resp. $`g_{1,0},\mathrm{},g_{n,0}`$) constitutes an $`L^2(M)`$-orthonormal basis of $`E_i(q)`$ (resp. $`E_j(q)`$) which diagonalizes the quadratic form $`Q_u^i`$ (resp. $`Q_u^j`$). Therefore, the family $`(f_{k,0}g_{l,0})_{km,ln}`$ constitutes a basis of the space $`E_i(q)E_j(q)`$ which diagonalizes the quadratic form $`S_u^{i,j}`$ given by
$`S_u^{i,j}(fg)`$ $`=`$ $`f_{L^2(M)}^2Q_u^j(g)g_{L^2(M)}^2Q_u^i(f)`$
$`=`$ $`{\displaystyle _M}u(f_{L^2(M)}^2g^2g_{L^2(M)}^2f^2)𝑑v.`$
The corresponding eigenvalues are $`(\mathrm{\Gamma }_l^{}(0)\mathrm{\Lambda }_k^{}(0))_{km,ln}`$. The criticality of $`q`$ for $`\lambda _j\lambda _i`$ then implies that this quadratic form admits eigenvalues of both signs, which means that it is indefinite.
On the other hand, in the case where $`\lambda _i(q)<\lambda _{i+1}(q)`$ and $`\lambda _j(q)>\lambda _{j1}(q)`$, we have, as in the proof of Proposition 2.3, for sufficiently small $`t`$, $`\lambda _i(q+tu)=\mathrm{max}_{km}\mathrm{\Lambda }_k(t)`$ and $`\lambda _j(q+tu)=\mathrm{min}_{ln}\mathrm{\Gamma }_l(t)`$, which yields
$`{\displaystyle \frac{d}{dt}}(\lambda _j\lambda _i)(q+tu)|_{t=0^{}}`$ $`=`$ $`\underset{ln}{\mathrm{max}}\mathrm{\Gamma }_l^{}(0)\underset{km}{\mathrm{min}}\mathrm{\Lambda }_k^{}(0)`$
$`=`$ $`\underset{km,ln}{\mathrm{max}}(\mathrm{\Gamma }_l^{}(0)\mathrm{\Lambda }_k^{}(0))`$
and
$`{\displaystyle \frac{d}{dt}}(\lambda _j\lambda _i)(q+tu)|_{t=0^+}`$ $`=`$ $`\underset{ln}{\mathrm{min}}\mathrm{\Gamma }_l^{}(0)\underset{km}{\mathrm{max}}\mathrm{\Lambda }_k^{}(0)`$
$`=`$ $`\underset{km,ln}{\mathrm{min}}(\mathrm{\Gamma }_l^{}(0)\mathrm{\Lambda }_k^{}(0)).`$
One deduces the following
###### Proposition 2.4.
If the potential $`qL_c^{\mathrm{}}(M)`$ is critical for the functional $`G_{ij}=\lambda _j\lambda _i`$, then, $`uL_{}^{\mathrm{}}(M)`$, the quadratic form $`S_u^{i,j}`$ is indefinite on $`E_i(q)E_j(q)`$.
Reciprocally, if $`\lambda _i(q)<\lambda _{i+1}(q)`$ and $`\lambda _j(q)>\lambda _{j1}(q)`$, and if, $`uL_{}^{\mathrm{}}(M)`$, the quadratic form $`S_u^{i,j}(g)`$ is indefinite on $`E_i(q)E_j(q)`$, then $`q`$ is a critical potential of the functional $`G_{ij}`$.
The following lemma will completes the proof of Theorem 1.5
###### Lemma 2.3.
The two following conditions are equivalent:
* $`uL_{}^{\mathrm{}}(M)`$, the quadratic form $`S_u^{i,j}`$ is indefinite on $`E_i(q)E_j(q)`$.
* there exist a finite family of eigenfunctions $`f_1,\mathrm{},f_k`$ in $`E_i(q)`$ and a finite family of eigenfunctions $`g_1,\mathrm{},g_l`$ in $`E_j(q)`$, such that $`_{1pk}f_p^2=_{1pl}g_p^2`$.
The proof of this lemma is similar to that of Lemma 2.2. Here, we consider the two convex cones $`C_i`$ and $`C_j`$ in $`L^2(M)`$ generated respectively by $`\left\{f^2;fE_i(q),f0\right\}`$ and $`\left\{g^2;gE_j(q),g0\right\}`$. Condition (ii) is then equivalent to the fact that these two cones admit a nontrivial intersection. As in the proof of Lemma 2.2, separation theorems enable us to prove that, if $`C_iC_j=\mathrm{}`$, then there exists a function $`u`$ such that $`_Muf^2𝑑v<0`$ for any $`fE_i(q)`$, and $`_Mug^2𝑑v0`$ for any $`fE_j(q)`$, which implies that $`S_u^{i,j}`$ is positive definite on $`E_i(q)E_j(q)`$. Since $`S_1^{i,j}=0`$, we have, $`S_u^{i,j}=S_{u_0}^{i,j}`$ with $`u_0=u\overline{u}L_{}^{\mathrm{}}(M)`$. Proposition 2.4 enables us to conclude.
Reciprocally, assume the existence of $`f_1,\mathrm{},f_kE_i(q)`$ and $`g_1,\mathrm{},g_lE_j(q)`$ satisfying $`_{1pk}f_p^2=_{1pl}g_p^2`$. Then, $`uL_{}^{\mathrm{}}(M)`$,
$$\underset{1pk}{}\underset{1p^{}l}{}S_u^{i,j}(f_pg_p^{})=\mathrm{}=0,$$
which implies that $`S_u^{i,j}`$ is indefinite on $`E_i(q)E_j(q)`$.
### 2.6. Proof of Theorem 1.6
Let $`q`$ be a local minimizer of $`G_{ij}=\lambda _j\lambda _i`$ and let us suppose for a contradiction that $`\lambda _i(q)<\lambda _{i+1}(q)`$ and $`\lambda _j(q)>\lambda _{j1}(q)`$. Given a function $`u`$ in $`L_{}^{\mathrm{}}(M)`$, we consider, as above, $`m`$ (resp. $`n`$) families of eigenvalues $`\mathrm{\Lambda }_1(t),\mathrm{},\mathrm{\Lambda }_m(t)`$ (resp. $`\mathrm{\Gamma }_1(t),\mathrm{},\mathrm{\Gamma }_n(t)`$) of $`\mathrm{\Delta }+q+tu`$, with $`m=dimE_i(q)`$ and $`n=dimE_j(q)`$, such that $`\mathrm{\Lambda }_1(0)=\mathrm{}=\mathrm{\Lambda }_m(0)=\lambda _i(q)`$ and $`\mathrm{\Gamma }_1(0)=\mathrm{}=\mathrm{\Gamma }_n(0)=\lambda _j(q)`$. As in the proof of Theorem 1.4, we will have for sufficiently small $`t`$, $`\lambda _i(q+tu)=\mathrm{max}_{km}\mathrm{\Lambda }_k(t)`$ and $`\lambda _j(q+tu)=\mathrm{min}_{ln}\mathrm{\Gamma }_l(t)`$. Hence, $`km`$ and $`ln`$,
$`\mathrm{\Gamma }_l(t)\mathrm{\Lambda }_k(t)`$ $``$ $`\lambda _j(q+tu)\lambda _i(q+tu)=G_{ij}(q+tu)`$
$``$ $`G_{ij}(q)=\mathrm{\Gamma }_l(0)\mathrm{\Lambda }_k(0).`$
It follows that, $`km`$ and $`ln`$, $`\mathrm{\Gamma }_l^{}(0)\mathrm{\Lambda }_k^{}(0)=0`$ and, then, the quadratic form $`S_u^{i,j}`$ is identically zero on $`E_i(q)E_j(q)`$ (recall that $`\mathrm{\Gamma }_l^{}(0)\mathrm{\Lambda }_k^{}(0)`$ are the eigenvalues of $`S_u^{i,j}`$). This implies that, $`fE_i(q)`$ and $`gE_j(q)`$, the function $`f_{L^2(M)}^2g^2g_{L^2(M)}^2f^2`$ is constant equal to zero (since its integral vanishes) which is clearly impossible unless $`i=j`$.
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# The metal abundance distribution of the oldest stellar component in the Sculptor dwarf spheroidal galaxyBased on data collected at the European Southern Observatory, proposal number 71.B-0621
## 1 Introduction
In hierarchical merging scenarios dwarf galaxies are thought to be the bricks from which larger galaxies were assembled. The study of the star formation history (SFH) and of the chemical evolution of presently existing dwarf galaxies is then extremely important for a proper understanding of the formation and evolution of the larger systems that this type of galaxies may have contributed to build in the past. Colour magnitude diagrams (CMDs) reaching the faint main sequence turn-off (TO) of the oldest stellar components are the most traditional and reliable way to derive the SFH of any individual dwarf galaxy. However, these detailed studies, possible so far mainly for the dwarf members of the Local Group (LG), require very time consuming observations.
The RR Lyrae variables, being about 3 magnitudes brighter than coeval TO stars ($`t>10`$ Gyrs), are much easier to observe. They offer an excellent tool for tracing the oldest stellar populations, and therefore the epoch of galaxy formation, in composite systems such as the resolved LG dwarf galaxies. Recent work by several groups has lead to the discovery of RR Lyrae stars in increasing numbers of LG dwarf galaxies (e.g. Leo I and II, IC 1613, Fornax, And VI, NGC6822: Held et al. 2001; Siegel & Majewski 2000; Dolphin et al. 2001; Bersier & Wood 2002; Pritzl et al. 2002; Clementini et al. 2003a, just to mention a few of them). An early stellar population, nearly coeval to the old Galactic globular clusters, has been found in the majority of LG galaxies, irrespective of their star formation histories. This indicates that all LG dwarfs started forming stars at an early epoch, $``$13 Gyr ago (e.g. Held et al. 2000).
The Sculptor dwarf spheroidal galaxy is no exception to this general trend. 226 RR Lyrae stars and 3 Anomalous Cepheids have been detected in a 15$`\mathrm{}\times 15\mathrm{}`$ area around the galaxy centre by Kaluzny et al. (1995) who published multi-epoch photometry for all of them. Kaluzny et al. (1995) also found that the period distribution of the RRab stars shows a sharp cut-off at $`P=0.475`$ d implying a metallicity of \[Fe/H\]$`1.7`$ (on the Zinn & West 1984 scale), and that the dispersion of the average $`V`$ magnitudes is most likely due to the metallicity spread exhibited by the stars in this galaxy. Similarly, Kovács (2001) found $``$\[Fe/H\]$`1.5`$ with a large dispersion from about $`2.0`$ to $`0.8`$ dex (on the metallicity scale by Jurcsik 1995), from the Fourier decomposition of the light curves of the RRab stars.
Indeed, Sculptor dwarf spheroidal (dSph) has long been known, from photometric studies, to have a large metallicity spread and bimodality in the metallicity and spatial distribution of its horizontal branch (HB) stars (Majewski et al., 1999; Hurley-Keller et al., 1999; Harbeck et al., 2001; Rizzi et al., 2004; Babusiaux et al., 2005). Very few spectroscopic determinations of the metal abundance of Sculptor stars existed so far (Norris & Bessell 1998; Tolstoy et al. 2001, 2003). Geisler et al. (2005), from high resolution spectroscopy of red giants, find \[Fe/H\] values in the range from $`2.10`$ to $`0.97`$ dex, confirming a large metallicity spread in Sculptor. Two distinct ancient populations showing an abrupt change in the \[Fe/H\] distribution at about 12 arcmin from the galaxy center are found in Sculptor by Tolstoy et al. (2004) based on FLAMES@VLT low resolution spectroscopy and WFI imaging of the galaxy. The two components have also different spatial distribution and velocity dispersion.
In this paper we present metal abundance determinations based on multi-slit low resolution spectroscopy obtained with FORS2 at the VLT for more than a hundred RR Lyrae stars in Sculptor. The variables are spread over a 15$`\mathrm{}\times 15\mathrm{}`$ area around the galaxy centre, thus being almost coincident with the internal region of Sculptor where Tolstoy et al. (2004) find segregation of red HB stars. The knowledge of the metal distribution of the RR Lyrae population allows to put important constraints on the early star formation and chemical evolution histories of the host galaxy, by removing the age-metallicity degeneracy: the earliest measurable data point for the chemical enrichment history of Sculptor can be determined with accuracy.
Observations and data reduction are discussed in Section 2. Metal abundances are presented in Section 3. The luminosity-metallicity relation and the radial velocities we determined for the RR Lyrae stars in Sculptor are discussed in Sections 4 and 5, respectively. In Section 6 the metallicity distribution of the Sculptor RR Lyrae stars is compared with that of other old stellar populations in the galaxy. A summary and some final remarks in Section 7 close the paper.
## 2 Observations and data reductions
Observations of 110 variable stars (107 RR Lyrae, 1 Anomalous Cepheid, 1 binary system, and 1 variable of unknown type) in Sculptor and of RR Lyrae stars in 4 Galactic globular clusters (namely: M 2, M 15, NGC 6171, and NGC 6441) were carried out using FORS2 (FOcal Reducer/low dispersion Spectrograph 2), mounted on the ESO Very Large Telescope (Paranal, Chile). The data were collected in service mode during the period July 29 to August 5, 2003. Typical seeing values during the observations were in the range 0$`\stackrel{}{.}`$7–1$`\stackrel{}{.}7`$ and on average of about 1$`\stackrel{}{.}2`$. We used the MXU (Mask eXchange Unit) configuration, that allows to observe simultaneously many objects with more freedom in choosing the location, size and shape of individual slitlets with respect to the standard MOS mode. The detector is a mosaic of two MIT CCDs with 15 $`\mu `$m pixel size.
Spectra were collected using the blue grism GRIS\_600B covering the 3450-5900 Å wavelength range, with a dispersion of 50 Åmm<sup>-1</sup> with slits 1<sup>′′</sup> wide, and usually 14<sup>′′</sup> long to allow for sky subtraction. With this configuration, each pixel corresponds to 0.75 Å. An effort was made to cover for each star the relevant wavelength range ($``$ 3900-5100 Å) containing both the CaII K and the hydrogen Balmer lines up to H$`\beta `$. We have used an instrumental set-up similar (i.e. same spectral range, resolution, and typical S/N) to that employed in our study of the RR Lyrae stars in the LMC (Gratton et al., 2004), so that RR Lyrae variables in some calibrating GCs are already available (namely in clusters NGC 1851, NGC 3201, and M 68). Exposure times on the Sculptor variables were of 31 min, as an optimal compromise between S/N and time resolution of the light curve of the RR Lyrae targets. We employed 9 masks in Sculptor, 2 in NGC 6441, 1 in M2, 1 in NGC 6171 and 1 in M15. The 9 Sculptor fields were slightly overlapped, so that for 25 variables we have more than 1 spectrum. A detailed log of the observations is given in Table 1 where N is the number of variable stars observed in each mask. The complete listing of the variables observed in Sculptor is provided in Table 2 where we have adopted Kaluzny et al. (1995) identification numbers. Their location on a $`16\mathrm{}\times 16\mathrm{}`$ map of the central region of Sculptor dSph galaxy is shown in Fig. 1. Finding charts corresponding to the nine $`6.8\mathrm{}\times 6.8\mathrm{}`$ FORS2 subfields are given in Appendix A. Centre of field coordinates are provided in Table 1. Equatorial coordinates for all our targets can be found in table 2 of Kaluzny et al. (1995).
Data reduction was performed using the standard IRAF<sup>1</sup><sup>1</sup>1 IRAF is distributed by the NOAO, which are operated by AURA, under contract with NSF routines. Images have been trimmed, corrected for bias and for the normalized flat field. Then we used the IRAF command lineclean to reduce the contamination by cosmic rays. Up to 19 spectra were present in each pointing, and were extracted with the optimal extraction and automated cleaning options switched on. The sky contribution was subtracted making use of the slit length. The contamination of targets from nearby stars was reduced to a minimum, except for a few objects. For each science mask a HeCdHg lamp was acquired, and used to calibrate in wavelength the spectra, each one covering a different spectral range, depending on the target position. Not less than 10 lines of the calibration lamp were visible for each aperture, and the resulting dispersion solutions have r.m.s. of about 0.03 Å. Further cleaning of cosmic rays hits and bad sky subtractions was done using the clipping option in the IRAF splot task. Fig. 2 shows examples of the final spectra.
## 3 Derivation of the metal abundances
We obtained spectra for 110 of the variable stars identified in Sculptor by Kaluzny et al. (1995), and for 25 of them we have multiple observations. These authors published photometry in the $`V`$ band for all our targets. Periods and epochs of maximum light were determined from their time series data (kindly made available by Dr. J. Kaluzny), using the period search package GRaphycal Analyzer of TIme Series (GRATIS, Di Fabrizio 1999, Clementini et al. 2000). The new ephemerides are provided in Table 2. We found that our periods are in general slightly different from those published by Kaluzny et al. (1995); differences are in most cases around the fourth or fifth digit. However, there are a number of cases where our periods and type classifications significantly differ from those of Kaluzny et al. (1995) who published aliases of the periodicities preferred here. All these objects have been flagged in the last column of Table 2, where we provide comments on individual stars. We used our new periods to phase the spectra of our target stars in Sculptor since the scatter in the light curves appears to be significantly reduced than using the published values. Phases corresponding to the Heliocentric Julian Day (HJD) at half exposure are listed in Column 3 of Table 2; for the double-mode RR Lyrae stars they correspond to the first overtone pulsation period. We also computed intensity-averaged luminosities for all the variables in our study, that are given in Column 5 of Table 2. Based on our study of the light curves our sample contains 62 ab-type, 40 c-type, 3 confirmed and 2 suspected d-type RR Lyrae stars, 1 Anomalous Cepheid, 1 suspected binary system, and 1 variable of unknown type.
### 3.1 Metallicities and metal abundance distribution of the variable stars in Sculptor
Precise and homogeneous metal abundances for the target stars in Sculptor were measured using the revised version of the $`\mathrm{\Delta }S`$ method (Preston, 1959) devised by Gratton et al. (2004). We do not actually measure $`\mathrm{\Delta }`$S values, but rather estimate metallicities for individual variables by comparing the strength of the H lines and of the K Ca II line with analogous data for variables in GCs of known metallicity. A detailed description of our method can be found in Gratton et al. (2004). A summary of the technique, and an update of the calibration procedures are provided in Appendix B, to which the interested reader is referred to for details. Spectral line indices measured for the variables in Sculptor following Gratton et al. (2004), are given in Columns from 8 to 12 of Table 2.
The correlation between $`K`$ and $`H`$ spectral indices is shown in Fig. 3. In this figure the solid lines represent the mean relations of the calibrating clusters M 68 and NGC 1851 (see Appendix B). The variables in Sculptor fall almost entirely below the mean line of NGC 1851 indicating that they are metal-poorer than this cluster (\[Fe/H\]<sub>ZW</sub>=$``$1.36), and also extend below the mean line of M 68, showing that there is a number of variables in Sculptor metal-poorer than \[Fe/H\]<sub>ZW</sub>=$``$2.09. Metallicity indices $`M.I.`$ defined according to equations (1), (2) and (3) in Appendix B, are listed in Column 13 of Table 2. Individual metallicities were derived from the calibration equations (4) and (5) in Appendix B in the Zinn & West (1984, hereafter ZW84) and in the Carretta & Gratton (1997, hereafter CG97) metallicity scales separately; they are listed in Columns 14 and 15 of Table 2. For different observations of the same star the \[Fe/H\] values are averages of all the available measurements. From the objects with multiple observations we estimate that errors of individual abundance determinations are 0.15 and 0.16 dex in ZW84 and CG97 metallicity scales, respectively.
The average metal abundance of the 107 RR Lyrae stars in our sample is \[Fe/H\]=$``$1.83$`\pm `$0.03 (r.m.s.=0.26 dex) with a total metallicity range of $`2.42<`$\[Fe/H\]$`<0.85`$ in ZW84 scale, and an offset of about 0.2 dex to higher metallicity in CG97 scale. The star-to-star scatter inferred from the r.m.s. dispersions is about 0.19-0.23 dex, hence larger than typical measurement errors. If we adopt 0.23 dex as the measured metallicity spread, and quadratically subtract a \[Fe/H\] measurement error of 0.16 dex, we obtain $`0.16`$ dex as our estimate of the intrinsic spread in the metal abundance of the Sculptor RR Lyrae stars.
The observed metallicity distribution of the Sculptor RR Lyrae stars (uncorrected for the measurement errors) in the ZW84 metallicity scale is shown in Fig. 4. Variable stars are divided by type. We find that the different types follow the same metallicity distributions. As well known, CG97 scale provides systematically higher metallicities than ZW84. In the following we will adopt ZW84 values, unless explicitly noted. However, independently of the adopted metallicity scale there are only very few metal-rich stars in our RR Lyrae sample. Based on the period distribution of the ab-type pulsators, Kaluzny et al. (1995) conclude that the bulk of the RR Lyrae stars in Sculptor have metallicities equal to or lower than \[Fe/H\]$`=1.7`$. Indeed, there are only 26 stars with \[Fe/H\]$`>1.7`$ in our sample. This is also the metallicity at which Tolstoy et al. (2004) find the break between metal-poor and metal-rich old populations in Sculptor. We will come back to this point in Section 6.
### 3.2 Comparison with metallicities from the Fourier parameters of the light curve and the pulsation equations
Metal abundances for the ab-type RR Lyrae stars in Sculptor have been estimated by Kovács (2001) using the parameters of the Fourier decomposition of the light curves and the Jurcsik & Kovács (1996) method. Derived metallicities are in the range from $`0.8`$ to $`2.0`$ dex, with an average value of \[Fe/H\]$``$1.5 dex in Jurcsik (1995) metallicity scale. Based on the frequency analysis of Kaluzny et al. (1995) data, Kovács (2001) also identified 15 confirmed and 3 suspected double mode pulsators in Sculptor and estimated their metal abundance using the pulsation equations. He found that, similarly to the RRab’s, the bulk of the double-mode RR Lyrae stars in Sculptor have \[Fe/H\]$`1.5`$, with only two RRd’s (stars 1168 and 5354 in Kaluzny et al. 1995) at \[Fe/H\]=$`1.9`$. The metallicity distributions of RRab and RRd stars separately are shown in fig. 5 of Kovács (2001), and can be compared with our distribution in Fig. 4. We recall that the Jurcsik (1995) metallicity scale is on average 0.2 dex more metal rich than the ZW84 scale at \[Fe/H\]$`1.5`$. This partially accounts for the difference between average values in Kovács (2001) and in our distribution in Fig. 4, but there still is a residual difference of 0.17 dex between the average values, and our distribution appears to cover a metallicity range larger than in Kovács (2001).
We determined metallicities for six of the RRd variables discovered by Kovács (2001), three of which are actually only suspected RRd’s. From our re-analysis of the light curves we think that these stars are monoperiodic c-type RR Lyraes with noisy light curves. The comparison between individual metallicity values is shown in Table 3, where the range given in Column 4 was evaluated from the metallicity distribution of the RRd stars in fig. 5 of Kovács (2001). Unfortunately we did not observe the two most metal poor RRd’s in Kovács (2001).
As with the RRab’s the metallicity range spanned by the RRd stars we analyzed is larger than that obtained by Kovács (2001) for the same stars.
## 4 The luminosity-metallicity relation
The luminosity-metallicity relation followed by the RR Lyrae stars in Sculptor was derived using the intensity-averaged mean magnitudes and the metal abundances in Table 2. We first discarded all objects that are not RR Lyrae stars or have incomplete light curves (stars 3302, 3710, 4780, and 5724). We also eliminated the most metal rich variable in our sample (star 4263), even if this star falls extremely well on the mean relations we derive. Following the procedure applied by Gratton et al. (2004) to the LMC RR Lyrae stars, we divided the Sculptor variables into 6 metallicity bins 0.1 dex wide; the corresponding average apparent magnitudes are given in Table 4 in the two metallicity scales, respectively. A least square fit weighted by the errors in both variables gives:
$$V=(0.092\pm 0.027)([\mathrm{Fe}/\mathrm{H}]_{\mathrm{ZW}}+1.5)+(20.158\pm 0.009)$$
where the errors in the slopes were evaluated via Monte Carlo simulations. The same slope is found for metallicities in CG97 scale.
The luminosity-metallicity relation of the Sculptor RR Lyrae stars is shown in Fig. 5. It is based on 105 stars covering the metallicity range \[Fe/H\] from $`1.36`$ to $`2.42`$, and we used different symbols for the various types of RR Lyrae stars. All variables seem to follow the same luminosity-metallicity relation independent of type. Gratton et al. (2004) found that the LMC double-mode RR Lyrae stars are offset to brighter luminosities in the luminosity-metallicity plane and explain this evidence with the LMC RRd’s being more evolved than the single-mode pulsators. The lack of a similar difference in luminosity between single and double-mode Sculptor RR Lyrae stars suggests that in this galaxy also the single-mode variables are evolved.
A more striking difference is the slope of the luminosity-metallicity relation, significantly shallower than that obtained for the LMC variables by Gratton et al. (2004)<sup>2</sup><sup>2</sup>2 However, we note that in both studies the slopes of the RR Lyrae luminosity-metallicity relations could be slightly underestimated by a few hundredths of magnitude, because of the relatively large errors of the individual \[Fe/H\] values.. On the other hand, we note that also the scatter in the average apparent luminosity of the RR Lyrae stars in Sculptor is less than half of that observed for the variables in the LMC: $`\sigma _V`$(Sculptor)=0.07 mag, compared to the 0.15-0.16 mag in the LMC (Clementini et al. 2003b). This scatter is almost entirely accounted for by photometric errors and the dispersion in metallicity of the Sculptor RR Lyrae stars, thus indicating that these variables have a very similar degree of evolution off the zero age horizontal branch. The shallow slope of the luminosity-metallicity relation in Sculptor could thus be explained with the Sculptor RR Lyrae’s being all rather evolved stars arising from the old metal poor population that also gives origin to the blue HB in the galaxy. Some support to this suggestion arises from the comparison of the average luminosity of the RR Lyrae stars with that of non variable horizontal branch stars in Sculptor.
From the color-magnitude diagram of Sculptor, values of $`V_{\mathrm{HB}}`$ 20.18-20.20 mag, $`V_{\mathrm{blue}\mathrm{HB}}`$ 20.20 and $`V_{\mathrm{red}\mathrm{HB}}`$ 20.29 mag are obtained by Kaluzny et al. (1995), Babusiaux, Gilmore, & Irwin (2005), and Rizzi (2003), respectively. This is at least $``$0.06-0.08 mag fainter than the average luminosity of the RR Lyrae stars which, for the 106 RR Lyrae with full light curve coverage, is $`V`$=20.127 with a dispersion of $`\sigma `$=0.072 (this value agrees well with the estimate by Kaluzny et al. 1995, based on the total sample of 226 RR Lyrae stars in Sculptor).
## 5 Radial velocity determinations
Radial velocities were measured from the spectra of the Sculptor RR Lyrae stars, and are given in Column 16 of Table 2. Multiple observations of the same stars show that our estimates have typical errors of about $`\pm 15`$ km s<sup>-1</sup>, with no systematic differences for different masks. This error includes the contribution of measurement uncertainties, errors related to the centering of the stars in the slit, and uncertainties due to the poor sampling of the radial velocity curves of the variable stars.
The average heliocentric radial velocity of all our variable stars in Sculptor is $`V_r`$=$`109.1\pm `$ 1.9 km s<sup>-1</sup> (r.m.s.=19.9 km s<sup>-1</sup>, 110 stars, having preliminarily averaged individual values for stars with multiple observations). The average radial velocity of the RR Lyrae stars alone is $`V_r`$=$`109.6\pm 1.9`$ km s<sup>-1</sup> (r.m.s.=19.8 km s<sup>-1</sup>, 107 stars). These values are in excellent agreement with the estimate obtained from K-giants in Sculptor by Queloz, Dubath & Pasquini (1995: $`V_r`$=$`109.9\pm 1.4`$ km s<sup>-1</sup>) and are consistent with the value of $`V_r`$=$`107\pm 2.0`$ km s<sup>-1</sup> previously derived by Armandroff & Da Costa (1986).
The good agreement with the literature values suggests that the offcentering problems noted for the calibrating cluster variables (see Appendix B) do not seem to affect the variables in Sculptor and that the undersampling of the variable star’s radial velocity curves does not significantly bias our estimates of the average $`V_r`$ value. To further check this point we extracted from the database of the Galactic field RR Lyrae stars analyzed with the Baade-Wesselink method (Liu & Janes 1990; Jones et al. 1992; Cacciari, Clementini & Fernley 1992; Skillen et al. 1993; Fernley 1994) template radial velocity curves of ab- and c-type RR Lyrae stars with metal abundance comparable to that of the variable stars in Sculptor, and estimated phase-dependent radial velocity corrections for each spectrum of RR Lyrae star. Using this procedure we found that the average correction to apply to the radial velocity mean value of the 107 RR Lyrae stars in our sample is less than $``$0.1 km s<sup>-1</sup> and can be safely neglected.
The difference between r.m.s. scatter of the RR Lyrae stars (19.8 km s<sup>-1</sup>) and typical measurement errors (15 km s<sup>-1</sup>) implies an intrinsic radial velocity dispersion of $`\sigma `$=12.9 km s<sup>-1</sup> for the RR Lyrae stars. This value is larger than the $`6.3+1.1,1.3`$ km s<sup>-1</sup> and $`6.2\pm 1.1`$ km s<sup>-1</sup> found for the K-giants by Armandroff & Da Costa (1986) and Queloz et al. (1995), and for Sculptor metal rich red giant stars by Tolstoy et al. (2004; $`\sigma _{metalrich}`$=$`7\pm 1`$ km s<sup>-1</sup>), but is consistent with the $`11\pm 1`$ km s<sup>-1</sup> dispersion observed by these same authors for the metal poor red giants in Sculptor. This result gives further support to the hypothesis that the RR Lyrae stars in Sculptor arise from the old, metal-poor population giving origin to the galaxy blue-horizontal branch, although our value for their velocity dispersion needs to be confirmed by higher resolution spectroscopy.
## 6 Metallicity distributions of the different old stellar components in Sculptor
The average metallicity, metal abundance distribution and range in metal abundance spanned by the RR Lyrae stars can be compared with the analogous quantities for other old stellar components in Sculptor, namely with the metallicity spread inferred from the width of the red giant branch (e.g. Kaluzny et al. 1995; Majewski et al. 1999; Rizzi 2003; Babusiaux et al. 2005), and with the abundances directly measured for red giants in Sculptor by Geisler et al. (2005) and Tolstoy et al. (2004), respectively.
The metallicity distribution in Fig. 4 shows that the RR Lyrae stars in Sculptor cover a full metallicity range of about 1.6 dex, which however reduces to $``$ 1 dex if the single most metal rich star in the sample is discarded. This range is larger than inferred from the spread of the red giant stars ($`\mathrm{\Delta }[\mathrm{Fe}/\mathrm{H}]`$ 0.6 dex, Kaluzny et al. 1995; Rizzi 2003; $``$0.8 dex, Majewski et al. 1999; $``$0.7 dex, Babusiaux et al. 2005), and is consistent with the spectroscopic study of red giants in Sculptor by Geisler et al. (2005: $``$1.1 dex) and Tolstoy et al. (2004) in the galaxy inner region (see upper panel of their fig. 3).
Tolstoy et al. (2004) find that the ancient stellar component ($``$ 10 Gyr old) in Sculptor is divided into two distinct groups having different metal abundance, kinematics and spatial distribution. The metal-rich population is concentrated within the $`r=0.2`$ degree central region, has metallicities in the range $`1.7<`$\[Fe/H\]$`<0.9`$, velocity dispersion of $`\sigma _{metalrich}`$=7$`\pm `$1 km s<sup>-1</sup> and is related to the Sculptor red horizontal branch. The metal-poor population is more spatially extended, has metallicities in the range $`2.8<`$\[Fe/H\]$`<1.7`$, velocity dispersion of $`\sigma _{metalpoor}`$=11$`\pm `$1 km s<sup>-1</sup> and is related to the Sculptor blue horizontal branch.
Our RR Lyrae stars are located in the central region of Sculptor, virtually coincident with the region where Tolstoy et al. (2004) find segregation of red HB stars. However, their metallicity distribution is dominated by the metal-poor objects with an average value of \[Fe/H\]=$``$1.83. The 26 stars with \[Fe/H\]$`>`$1.7 are marked by boxes in Fig. 1. Their average luminosity $`V_{[\mathrm{Fe}/\mathrm{H}]>1.7}`$=20.154 ($`\sigma `$=0.060, 26 stars) is marginally fainter that the average of the remaining 80 stars with \[Fe/H\]$``$1.7 ($`V_{[\mathrm{Fe}/\mathrm{H}]1.7}`$=20.118, $`\sigma `$=0.075, 80 stars), as expected given the shallow slope of the Sculptor RR Lyrae luminosity-metallicity relation, and is anyway brighter than both the blue and red HBs of the non variable stars, thus indicating that they are evolved objects, like their metal-poor counterpart. The average radial velocity of the two samples is only marginally different: $`v_r`$=117.5 (r.m.s.=18.5, 26 stars) for the metal-rich sample and $`v_r`$=107.0 (r.m.s.=19.6, 80 stars) for the metal-poor stars. The r.m.s. scatters are very similar and, once deconvolved for the measurement errors (15 km s<sup>-1</sup>), lead to very similar velocity dispersions of 10.8 and 12.6 km s<sup>-1</sup> in agreement with the velocity dispersion measured by Tolstoy et al. (2004) for the metal-poor component associated to Sculptor blue HB.
## 7 Summary and conclusions
Low resolution spectra obtained with FORS2 at the VLT have been used to measure individual metal abundances \[Fe/H\] and radial velocities for 107 RR Lyrae stars in the Sculptor dwarf spheroidal galaxy. Metallicities were derived using a revised version of the $`\mathrm{\Delta }`$S method (Gratton et al., 2004). The RR Lyrae stars in Sculptor are predominantly metal-poor with an average metal abundance of \[Fe/H\]=$`1.83\pm 0.03`$ (r.m.s.=0.26 dex) on the ZW84 metallicity scale, and only a few outliers having metallicities larger than $`1.4`$ dex. The observed metallicity dispersion is larger than the observational errors, thus showing that these variables have a real metallicity spread.
The RR Lyrae stars in Sculptor are found to follow a luminosity-metallicity relation with a slope of 0.09 mag dex<sup>-1</sup>, which is shallower than in the LMC (Gratton et al., 2004). This is explained with the Sculptor variable stars being rather evolved from the zero age HB, as also supported by their brighter luminosity compared to the non variable HB stars.
From our spectra we measured an intrinsic velocity dispersion 12.9 km s<sup>-1</sup> for the RR Lyrae stars, which appears to be in agreement with the dispersion derived by Tolstoy et al. (2004) for metal-poor red giants associated to the blue-HB stars in Sculptor.
All these evidences suggest that our RR Lyrae sample, and the RR Lyrae stars in Sculptor in general, are connected to the blue-HB population and arise from the first burst of star formation that produced the galaxy blue metal-poor HB. They allow to trace and distinguish this older component in the internal regions of the galaxy which are otherwise dominated by the metal-rich and younger red-HB stars. Indeed, this component only produced a few, if any, of the Sculptor RR Lyrae stars since average luminosity and kinematic properties of the few metal-rich objects in our sample suggest that they are more likely the high metallicity tail of the metal-poor star population rather than objects associated with Sculptor red-HB stars.
## Acknowledgments
A special thanks goes to J. Kaluzny for providing us the time series photometry of the Sculptor variables. We thank the anonymous referee for his/her comments and suggestions. This research was funded by MIUR, under the scientific projects: 2002028935, “Stellar Populations in the Local Group” (P.I.: Monica Tosi) and 2003029437, “Continuity and Discontinuity in the Milky Way Formation” (P.I.: Raffaele Gratton).
## 8 Appendix A - Finding charts
We present in this Section finding charts (Fig.s 6a,b,c) for all the variable stars we observed in Sculptor. They correspond to the nine $`6.8\mathrm{}\times 6.8\mathrm{}`$ FORS2 subfields used to map the central $`15\mathrm{}\times 15\mathrm{}`$ area of the galaxy. Centre of field coordinates are provided in Table 1. In each map North is up and East to the left, and variables are identified according to Kaluzny et al. (1995) identifiers. The RR Lyrae stars are marked by open circles (in red in the electronic edition of the journal), other types of variables by (blue) open squares. The complete listing of the variables observed in Sculptor is provided in Table 2. Equatorial coordinates for all our targets can be found in table 2 of Kaluzny et al. (1995).
## 9 Appendix B - The metallicity calibration
Following the procedure devised by Gratton et al. (2004), line indices for RR Lyrae stars are computed from the spectra shifted to rest wavelength by directly integrating the instrumental fluxes in spectral bands centered on the Ca II K, H$`\delta `$, H$`\gamma `$, and H$`\beta `$ lines (see table 2 and fig. 10 of Gratton et al. 2004 for the definition of the spectral bands). Then a $`H`$ index is defined as the average of the indices of the 3 hydrogen lines, and $`K`$ as the index of the Ca II K line. Metallicities are derived by comparing the $`H`$ and $`K`$ indices measured for the target stars to the same quantities for variables in a number of globular clusters of known metal abundance.
The calibration of the line indices of the Sculptor variables in terms of metal abundances \[Fe/H\] was obtained using RR Lyrae stars in the newly observed clusters M 15, M 2 and NGC 6171, and in the calibrating clusters used in Gratton et al. (2004), namely M 68, NGC 1851 and NGC 3201. For all these clusters, precise metal abundances are available on both the ZW84 and CG97 metallicity scales. NGC 6441, a metal rich cluster (\[Fe/H\]=$`0.59`$ according to ZW84 metallicity scale) having RR Lyrae stars with anomalously long periods, was not used for calibration purposes and will be discussed elsewhere (see Clementini et al. 2005).
Line indices measured for RR Lyrae stars in the calibrating clusters analyzed in the present paper are provided in Tables 56, and 7 for M 15, M 2 and NGC 6171, respectively.
The calibrating clusters define mono-metallic correlations in the $`K`$ versus $`H`$ plane. These relations are shown in Fig. 7 where filled squares represent stars in M 15, M 2 and NGC 6171, and solid lines are the mean relations defined by M 68, NGC 1851 and NGC 3201 taken from Gratton et al. (2004). RR Lyrae stars in M 15 fall precisely on the mean line of M 68 and those in M 2 closely follow the mean line of NGC 3201, confirming the good agreement between metal abundances of these two pairs of clusters. Stars in NGC 6171 generally fall slightly above the mean line of NGC 1851, in agreement with the slightly higher metallicity of NGC 6171 with respect to NGC 1851.
According to Gratton et al. (2004) the mean relations drawn in Fig. 7 by M 68 and NGC 1851 are:
$$K_1=0.30931.2815H+1.3045H^2$$
(1)
(valid for $`H`$ between 0.12 and 0.40) for M 68, and:
$$K_2=0.64322.6043H+3.0820H^2$$
(2)
(valid for $`H`$ between 0.04 and 0.34) for NGC 1851. We thus defined metallicity index $`M.I.`$ the quantity:
$$M.I.=(KK_1)/(K_2K_1)$$
(3)
where $`K`$ is the Ca II K line index of the star, and $`K_1,K_2`$ are derived entering the $`H`$ index measured for the star into equations (1) and (2). $`M.I.`$ values derived by this procedure for the RR Lyrae stars in the calibrating clusters analyzed in the present paper are listed in Column 10 Tables 5, 6, 7.
The calibration of the metallicity index in terms of metal abundance \[Fe/H\] was derived by computing average $`M.I.`$ values for the M 15, M 2 and NGC 6171 variables from the individual $`M.I.`$’s in Tables 5, 6, 7 and from the $`M.I.`$ values in Column 9 of tables 3, 4, and 5 of Gratton et al. (2004) for M 68, NGC 1851 and NGC 3201, and correlating these $`M.I.`$’s with the metal abundances of the clusters on the ZW and CG metallicity scales, respectively. The average $`M.I.`$ values and their dispersions are summarized in Columns 6 and 7 of Table 8 along with the cluster metallicities (and their uncertainties) in the two metallicity scales (Columns from 2 to 5).
The correlation between \[Fe/H\] and $`M.I.`$ values is very well represented by linear regressions with scatter typically within the error of measure. These linear regressions are described by the following equations:
$$[\mathrm{Fe}/\mathrm{H}]_{\mathrm{ZW}}=0.882M.I.2.170$$
(4)
$$[\mathrm{Fe}/\mathrm{H}]_{\mathrm{CG}}=0.941M.I.2.000$$
(5)
and are shown in Fig. 8 for the ZW84 and CG97 metallicity scales, respectively.
Individual metal abundances for the RR Lyrae stars in the calibrating clusters analyzed in the present paper were derived from the above calibration equations. They are listed in Tables 5, 6, 7 respectively, while the mean metallicities derived from the averages of these individual values and their respective dispersions are given in Table 8.
The last column in each of Tables 56, and 7 gives individual radial velocities measured from the spectra of the stars observed in the calibrating clusters. Averages of these values are listed in Table 8. We note that these average radial velocities differ somewhat from the literature values, particularly for NGC 6171, suggesting the presence of systematic offsets possibly caused by offcentering of the cluster variables in the slit. These systematic differences are not found in the case of the Sculptor variables, for which the average radial velocity we measured is perfectly consistent with the literature values (see Section 5). This suggests that we are simply seeing an effect of the small samples in the Galactic clusters, while the one in Sculptor is large enough to average away the offcentering effects.
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# Chandra Observations of A Galactic Supernova Remnant Vela Jr.: A New Sample of Thin Filaments Emitting Synchrotron X-Rays
## 1 Introduction
Supernova remnants (SNRs) play crucial roles for heating and chemical evolution of galaxies. Their shocks are also famous as cosmic-ray accelerators. Koyama et al. (1995) discovered synchrotron X-rays from shells of SN 1006, which is the first observational result indicating that SNRs accelerate electrons up to $``$TeV. At present, several SNRs have been identified as synchrotron X-ray emitters (e.g. RX J1713.7$``$3946, Koyama et al. 1997; Slane et al. 1999; RCW 86, Bamba, Tomida, & Koyama 2000; Borkowski et al. 2001; Rho et al. 2002; Cas A., Vink & Laming 2003; see also Bamba et al. 2005; Bamba 2004).
The most plausible acceleration mechanism is the diffusive shock acceleration (DSA) (e.g. Bell, 1978; Drury, 1983; Blandford & Eichler, 1987; Jones & Ellison, 1991; Malkov & Drury, 2001), which can accelerate charged particles on the shock into a power-law distribution, similar to the observed spectrum of cosmic rays on the earth. However, since there is still little observational information and theoretical understanding, there are many unresolved problems, such as the injection efficiency, the maximum energy of accelerated particles, the configurations of magnetic field, and so on. The magnetic field amplification is also pointed out as an remarkable process (Bell & Lucek, 2001; Lucek & Bell, 2000). Recently, it is found that nonthermal X-rays from the shells of historical SNRs concentrate on very narrow filamentary regions (Bamba et al., 2003b, 2005). Together with the roll-off frequency ($`\nu _{rolloff}`$) of synchrotron X-rays (Reynolds, 1998; Reynolds & Keohane, 1999), Bamba et al. (2005) found that an empirical function, $`\nu _{rolloff}/w_d^2`$, decreases with the age of an SNR ($``$—age relation). The relation may reflect the time evolution of the magnetic field around the shock front (Bamba et al., 2005). However, the $``$—age relation has still many uncertainties because of poor statistics. Clearly other new samples with synchrotron X-ray filaments are desired.
Vela Jr. (RX J0852.0$``$4622, G266.6$``$1.2) is one of the most curious galactic SNR with large radius of about 60 arcmin (Green, 2004). It was discovered by ROSAT (Aschenbach, 1998). and well known that the COMPTEL team reported a possible detection of $`1.157`$ MeV $`\gamma `$-ray line emitted by <sup>44</sup>Ti (Iyudin et al., 1998). This fact strongly indicates that Vela Jr. is a nearby and very young SNR because the half lifetime of <sup>44</sup>Ti is only about 60 years, however, there remains some doubt about the detection (Shonfelder et al., 2000). Tsunemi et al. (2000) reported the presence of Ca K line, suggesting that there are a plenty of Ca produced from <sup>44</sup>Ti. On the other hand, Slane et al. (2001) reported that the 1$`\sigma `$ upper-limit of Sc-K line is $`4.4\times 10^6`$ cm<sup>-2</sup>s<sup>-1</sup>. Recently, Iyudin et al. (2005) found the possible detection of Ti and Sc K emission lines with XMM-Newton. For the 78.4 keV line emission from <sup>44</sup>Ti, von Keinlin et al. (2004) obtained the upper-limit of $`1.1\times 10^4`$ cm<sup>-2</sup>s<sup>-1</sup> with INTEGRAL/SPI. The presence of a central source has been also suggested (Aschenbach, 1998), which was confirmed with Chandra (Pavlov et al., 2001). The spectrum of the source, CXOU J085201.4$``$461753, together with the absence of any optical counter part (Mereghetti, Pellizzoni, & de Luca, 2002), reminds us of a neutron star. These observational reports might indicate that the progenitor is a core-collapse supernova (SN).
The other fact which makes this SNR famous is the existence of nonthermal X-rays from the rims (Slane et al., 2001), which implies that the rims of Vela Jr. are cosmic ray accelerators like SN 1006 (Koyama et al., 1995) and other SNRs with synchrotron X-rays (e.g. Bamba, 2004, and references therein). Recently, TeV $`\gamma `$-rays with energies greater than 500 GeV are also detected at the 6$`\sigma `$ level from the north-western rim of the SNR with CANGAROO-II (Katagiri et al., 2005) and H.E.S.S. (Aharonian et al., 2005). These $`\gamma `$-rays are produced by accelerated electrons and/or protons with energies more than $``$ TeV.
Despite such an interesting source, the precise age of the remnant is still unknown. There are some spikes in nitrate concentration measured in an Antarctic ice core, which might be regarded as signals of SNe (Burgess & Zuber, 2000). One of them is not yet identified with historical SNe ever known. If the spike really associated with the progenitor of Vela Jr., the SN may have occurred around AD 1320 and the age of the remnant is about 680 yrs. However, there are many uncertainties, hence further studies are necessary. In addition, an unclear issue is the distance to the remnant, which brings us numerous information, such as the total luminosity of the <sup>44</sup>Ti emission line, the average expansion velocity, and so on. It may also give us some answer on the current debate whether or not Vela (the distance of 250 pc; Cha et al., 1999) and Vela Jr. interact with each other. Iyudin et al. (1998) estimated the distance to be $``$200 pc with an assumption that the age is 680 years mentioned above, while Slane et al. (2001) suggested that the distance is 1–2 kpc from the relatively large absorption column. Although Moriguchi et al. (2001) observed the molecular distribution using <sup>12</sup>CO (J = 1—0) emission with NANTEN, they could not determine the precise distance because the remnant is located near the tangential point of a galactic arm. The upper limit of the distance they estimated is $``$1 kpc.
In this paper, using the Chandra data for the first time, we report on the analysis of hard X-ray filaments associated with Vela Jr. and discuss on the acceleration mechanism at the shock of the SNR, its age, and the distance to the SNR. This paper is organized as follows. We summarize the observation details of Vela Jr. with Chandra in § 2. Section 3 describes the results of the observations. We discuss on <sup>44</sup>Sc emission line reported by Iyudin et al. (2005) in § 4, the origin of nonthermal X-rays in § 5, on the interaction between the SNR and molecular clouds (§ 6), and on the origin of TeV $`\gamma `$-rays (§ 7). In § 8, we argue the distance and the age of the SNR using the correlation between a function $``$ and the SNR age.
## 2 Observations
We used the Chandra archival data of the ACIS of the north-western (NW) rim of Vela Jr. (ObsID = 3846 and 4414). Figure 1 shows the ASCA GIS image (Tsunemi et al., 2000) with the field of view of Chandra observations. The satellite and the instrument are described by Weisskopf et al. (2002) and Garmire et al. (2000), respectively. Data acquisition from the ACIS was made in the Timed-Exposure Faint mode. The data reductions and analysis were made using the Chandra Interactive Analysis of Observations (CIAO) software version 3.0.2. Using the Level 2 processed events provided by the pipeline processing at the Chandra X-ray Center, we selected ASCA grades 0, 2, 3, 4, and 6, as the X-ray events. The “streak” events on the CCD chip S4 were removed by using the program destreak<sup>1</sup><sup>1</sup>1See http://asc.harvard.edu/ciao2.3/ahelp/destreak.html. in CIAO. In order to make statistics better, we improved the astrometory of the data following the CIAO data analysis threads and combined these data. The exposure time of each observation is 39 ks (ObsID = 3846) and 35 ks (ObsID = 4414), respectively. Table 1 summarizes about these observations. Hereinafter, all results are from the combined data so far as there is no mention.
## 3 Results
Figure 2 shows images of the NW rim in the 0.5–2.0 keV (a) and 2.0–10.0 keV (b) bands, binned with 4 arcsec scale. Neither the subtraction of background photons, smoothing process, nor correction of the exposure time did not be performed. The difference of background level among the chips are caused by that of the CCD chips (back-illuminated and front-illuminated). These images look alike, and show very straight filamentary structures clearly on the outer edge of the rim. Two filaments can be well recognized especially in the upper panel. Iyudin et al. (2005) suggests that the inner filament may be the reverse shock emission. However, we have no clue to distinguish whether it indicates the reverse shock or apparently overlaid forward shock via the projection effect. The filaments are very similar to those in SN 1006 (Long et al., 2003; Bamba et al., 2003b), indicating that they may be efficient accelerators of electrons like SN 1006 filaments. Then we made analysis in the same way to our previous one for SN 1006 (Bamba et al., 2003b).
For the first step, we made a spectrum from the whole emitting region of this rim. Since the rim is covered by two CCD chips (chip 6 and 7), source spectra were made separately from each chip. Background photons were accumulated from outer regions of the rim. All the spectra extend up to $`>`$2 keV, indicating that the hard X-ray emission is nonthermal. There is some line-like feature below 2 keV, which probably comes from thermal emission of Vela and/or Vela Jr. as already mentioned (Slane et al., 2001). Since our data does not have sufficient statistics to examine the properties of thermal plasma, we ignored the component. In Hiraga et al. (2005), detailed analysis for the thermal plasma is done with the XMM-Newton deep observation. We fitted the spectra with an absorbed power-law model only with photons above 2 keV, in order to avoid the contamination of thermal photons. The absorption column was subsequently calculated using the cross sections by Morrison & McCammon (1983) with the solar abundances (Anders & Grevesse, 1989). The best-fit parameters are shown in Table 2.
The best-fit values of the photon index and absorption column are consistent with those of the nonthermal component in previous results (Slane et al., 2001; Iyudin et al., 2005; Hiraga et al., 2005), indicating that almost all photons above 2 keV are nonthermal origin. Thus, we regarded all the photons above 2 keV as nonthermal in the following. As shown in Figure 2(b), three filaments were selected (filament 1–3) in order to study their spatial and spectral characteristics, which are straight and free from other structures, in the same way to the analysis of SN 1006 case (Bamba et al., 2003b). Since all the filaments are located within 4 arcmin from the aim-point (see Table 1), the size of a point spread function (PSF) is about 0.5 arcsec. Although we can see the filaments more clearly in the soft X-ray band in Figure 2, Bamba et al. (2003b) showed us that the contamination of thermal photons makes filaments broader. Therefore, we conducted the spatial analysis only with photons above 2 keV. Figure 3 shows spatial profiles of the filaments in the 2.0–10.0 keV (nonthermal) band binned to a resolution of 1 arcsec. We can see profiles with clear decays on the downstream sides and sharp edges on the upstream sides in all profiles. The filament 2 and 3 are double peaked, as suggested by Iyudin et al. (2005). To estimate the scale width of these filaments on upstream ($`w_u`$) and downstream ($`w_d`$) sides, we fitted them with exponential function as a simple and empirical fitting model for the apparent profiles, which is the same as used in the previous analysis in the SN 1006 (Bamba et al., 2003b) and other historical SNR (Bamba et al., 2005) cases. Background photons were not subtracted and treated as a spatially independent component. For the filament 3, the $`w_u`$ could not be determined due to the lack of statistics, then it was frozen to be the PSF size (0.5 arcsec). For each filament, fitting region was limited around the outermost peak in order to avoid the influence of the inner structure. The model well reproduced the data. The best-fit models and parameters are shown in Figure 3 and Table 3, respectively.
For the next step, we made spectra of the filaments. The background regions were selected from the downstream sides of the filaments, where there was no other structure. The spectra were so hard and fitted with an absorbed power-law model again. The absorption column was fixed to be $`6.1\times 10^{21}\mathrm{cm}^2`$ in order to make statistics better, which is the best-fit value for the total region (Table 2) and consistent with previous results. The index of the power-law component is similar to those of nonthermal shells in other SNRs (SN 1006, Bamba et al. 2003b; RX J1713.7$``$3946, Koyama et al. 1997; Slane et al. 1999), which is considered to be synchrotron emission. Therefore, we concluded that the nonthermal emission in Vela Jr. shell is synchrotron, and applied the SRCUT model, which is one of the model representing the synchrotron emission from electrons of power-law distribution with exponential roll-off in a homogeneous magnetic field (Reynolds, 1998; Reynolds & Keohane, 1999). The spectral index at 1 GHz was fixed to be $`0.3`$ according to a report by Combi, Romero, & Benaglia (1999). Since the value $`p=0.3`$ is rather smaller than those of ordinary SNRs, about 0.5, we also tested the model under the assumption of $`p=0.5`$. All models can reproduce the data with similar values of reduced $`\chi ^2`$ as shown in Table 4.
## 4 Comments on <sup>44</sup>Sc fluorescent line
Iyudin et al. (2005) reported a significant emission line at 4.45$`\pm `$0.05 keV from the NW rim of Vela Jr. with XMM-Newton. We checked the spectrum of each filament as well as the averaged one so as to investigate the presence of some line-like feature. As a result, we found that the spectrum of filament 1 has an excess around 4 keV as shown in Figure 4, while the spectrum averaged for whole rim does not show the line feature. Since one can see the excess in the data of both observations, it may be a real line. Then, the power-law plus a narrow line model were applied and accepted, although the reduced $`\chi ^2`$ does not improve significantly ($`\chi ^2`$/d.o.f.=69.5/79; see also Table 4). The center energy is $`4.1\pm 0.2`$ keV, which is slightly lower than the averaged one by Iyudin et al. (2005) but consistent with the ASCA result (Tsunemi et al., 2000). The total flux, $`7.3_{4.5}^{+5.1}\times 10^7`$ photons cm<sup>-2</sup>s<sup>-1</sup>, is about 10% of that derived by Iyudin et al. (2005) and Tsunemi et al. (2000), and consistent with the upper limit by Slane et al. (2001). The difference of the intensity may be because of our smaller region the spectra were accumulated than those by other authors. Moreover, the center energy and significance depend on the choice of the background region and the model of the continuum component. The lack of statistics prevents us from reliable conclusion. Detailed analysis with excellent statistics and spectral resolution are encouraged.
## 5 Origin of the Filaments
We found that the scale width of the filaments in Vela Jr. is much smaller than its radius. Filament 3 has a small $`w_u`$ in the order of Chandra PSF size (0.5 arcsec), indicating that the $`w_u`$ may be similar to or smaller than the PSF size. Although the other filaments have $`w_u`$ significantly larger than the PSF size, it may be caused by the projection effect or other structures, such as the second peak and/or other filaments. Therefore, we concluded that the scale width on the upstream side is similar to or smaller than the PSF size. On the other hand, the value of $`w_d`$ is much larger than the PSF size. These results are similar to other case of nonthermal filaments in young SNRs (Bamba et al., 2005). The power-law index is similar to the results with previous observations (Slane et al., 2001; Hiraga et al., 2005). These thin filaments with hard X-rays remind us the nonthermal X-rays from filaments in young SNRs. Wide band spectra of most of these SNRs show that the nonthermal X-rays are synchrotron radiation from highly accelerated electrons. Therefore, we consider that the filaments emit X-rays via synchrotron radiation. In order to confirm our conclusion, we need more information in the radio continuum band to make the wide band spectrum of synchrotron emission.
## 6 Interaction with Molecular Clouds?
The filaments in Vela Jr. are rather straight and their lengths are comparable to the radius of the SNR, that are similar to those in SN 1006 (Bamba et al., 2003b) and Tycho (Hwang et al., 2002), and are unlike clumpy filaments in Cas A (Vink & Laming, 2003), RCW 86 (Rho et al., 2002), and RX J1713.7$``$3946 (Uchiyama et al., 2003; Lazendic et al., 2004; Cassam-Chenaï et al., 2004). The former samples are located in tenuous interstellar space, whereas for the two SNRs in the latter cases, there are some reports about the interaction between shocks and molecular clouds (RCW 86: Moriguchi et al., private communication; RX J1713.7$``$3946: Fukui et al. 2003). Interactions with molecular clouds may distort beautiful filaments via turbulence and so on. There are molecular clouds around Vela Jr. (Moriguchi et al., 2001). However, it is not clear whether the shocks are interacting with them or not, because there is no observation with excited molecular cloud lines (CO(2$``$1) and so on). The straight filaments may indicate that there is no interaction between the shock and the molecular cloud in this region. Further observations with excited molecular lines are needed.
## 7 Origin of TeV $`\gamma `$-rays
Recently, two instruments have independently detected TeV $`\gamma `$-rays significantly (Katagiri et al., 2005; Aharonian et al., 2005). The reported differential flux is consistent with each other, but the photon index is slightly different. In this section, we examine the origin of the TeV $`\gamma `$-rays, considering the preferable values of the maximum energy of electrons ($`E_{max}`$) and the downstream magnetic field ($`B_d`$) (Yamazaki et al., 2004). Bamba et al. (2005) suggested that the width of the filament on the downstream side ($`w_d`$) and the roll-off frequency of synchrotron emission ($`\nu _{rolloff}`$) reflect the value of $`E_{max}`$ and $`B_d`$, then we can use these two observational value for the study on $`E_{max}`$ and $`B_d`$. Hereinafter, we adopt flux-averaged mean values, such as
$`w_d`$ $`=`$ $`49.5\mathrm{arcsec}=0.24_{0.07}^{+0.19}\left({\displaystyle \frac{D}{1\mathrm{kpc}}}\right)\mathrm{pc},`$ (1)
$`\nu _{rolloff}`$ $`=`$ $`6.6_{1.6}^{+2.1}\times 10^{16}\mathrm{Hz},`$ (2)
where $`D`$ is the distance to Vela Jr. Considering the projection effect, $`w_d`$ can be written as
$`w_d`$ $`=`$ $`\alpha r^1u_st_{loss},`$
$`t_{loss}`$ $`=`$ $`1.25\times 10^3\mathrm{yrs}\left({\displaystyle \frac{E_{max}}{100\mathrm{TeV}}}\right)^1\left({\displaystyle \frac{B_d}{10\mu \mathrm{G}}}\right)^2,`$
where $`\alpha `$, $`r`$, $`u_s`$, and $`t_{loss}`$ are a correction factor of the projection effect, the compression ratio, the shock velocity, and synchrotron loss time scale (Bamba et al., 2005; Yamazaki et al., 2004). Although $`\alpha `$ is an unknown parameter depending on the shape of profiles and the curvature radius, it becomes $``$1–10 as far as the width of filaments is negligible to the SNR radii (Berezhko & Völk, 2004). Then the above equations lead
$$\left(\frac{E_{max}}{100\mathrm{TeV}}\right)^1\left(\frac{B_d}{10\mu \mathrm{G}}\right)^2=6.3_{1.8}^{+5.0}\times 10^2\alpha ^1r\left(\frac{u_s}{3000\mathrm{km}\mathrm{s}^1}\right)^1\left(\frac{D}{1\mathrm{kpc}}\right).$$
(3)
On the other hand, the roll-off frequency $`\nu _{rolloff}`$ is represented as (Reynolds & Keohane, 1999; Reynolds, 1998)<sup>2</sup><sup>2</sup>2See http://heasac.gsfc.nasa.gov/docs/xanadu/xspec for the erratum of coefficients.
$`\nu _{rolloff}`$ $`=`$ $`1.6\times 10^{18}\mathrm{Hz}\left({\displaystyle \frac{E_{max}}{100\mathrm{TeV}}}\right)^2\left({\displaystyle \frac{B_d}{10\mu \mathrm{G}}}\right).`$ (4)
In the following, we adopt typical values $`\alpha =5`$ and $`r=4`$. Then, using Eqs. (3) and (4), we derived
$`E_{max}`$ $`=`$ $`4.1_{0.9}^{+0.9}\mathrm{TeV}\left({\displaystyle \frac{u_s}{3000\mathrm{km}\mathrm{s}^1}}\right)^{1/3}\left({\displaystyle \frac{D}{1\mathrm{kpc}}}\right)^{1/3}\left({\displaystyle \frac{\alpha }{5}}\right)^{1/3}\left({\displaystyle \frac{r}{4}}\right)^{1/3},`$ (5)
$`B_d`$ $`=`$ $`2.2_{0.2}^{+0.2}\times 10^2\mu \mathrm{G}\left({\displaystyle \frac{u_s}{3000\mathrm{km}\mathrm{s}^1}}\right)^{2/3}\left({\displaystyle \frac{D}{1\mathrm{kpc}}}\right)^{2/3}\left({\displaystyle \frac{\alpha }{5}}\right)^{2/3}\left({\displaystyle \frac{r}{4}}\right)^{2/3}.`$ (6)
Katagiri et al. (2005) examined whether the TeV $`\gamma `$-rays arise from the inverse Compton emission via accelerated electrons or $`\pi ^0`$ decay caused by accelerated protons. The former requires very high size ratio of X-ray and TeV $`\gamma `$-ray emission regions ($`V_{\mathrm{TeV}}/V_{\mathrm{X}\mathrm{ray}}10^5`$) and strong magnetic field of $`B_d`$1.6 mG. Then we obtain the shock velocity with eq.(6) as
$$u_s6\times 10^8\mathrm{cm}\mathrm{s}^1\left(\frac{D}{100\mathrm{pc}}\right)\left(\frac{\alpha }{5}\right)^1\left(\frac{r}{4}\right),$$
(7)
so that the small distance, $`D100`$ pc (the very small physical radius and the very young age, in other words) may be required, which is somewhat doubtful. This result is consistent with the discussion in Aharonian et al. (2005). On the other hand, $`\pi ^0`$-decay model, which requires the interaction of the SNR and molecular clouds, can naturally explain the multi-band spectrum (Katagiri et al., 2005). However, there might be some discrepancy between the observed straight filaments as discussed in § 6.
## 8 Estimation of the Age and the Distance
More determinative arguments on the distance ($`D`$) can be possible, and at this time, the age of the SNR ($`t_{age}`$) can be also discussed simultaneously. As a tool of the estimation, $``$—age relation (Bamba et al., 2005) is used, where
$``$ $``$ $`\nu _{rolloff}/w_d{}_{}{}^{2}=Ct_{age}{}_{}{}^{\alpha },`$ (8)
$`C`$ $`=`$ $`2.6_{1.4}^{+1.2}\times 10^{27}\mathrm{Hzpc}^2,`$ (9)
$`\alpha `$ $`=`$ $`2.96_{0.06}^{+0.11}.`$ (10)
Let $`\theta _R`$ ($`R_s/D`$) and $`\theta _d`$ ($`w_d/D`$) be the angular radius of the SNR and the observed angular scale width of synchrotron X-ray filaments on the downstream sides, respectively. Here $`R_s`$ is the radius of the SNR in the unit of pc. From the equation defining $``$, we obtain
$`D`$ $`=`$ $`0.65\mathrm{kpc}{\displaystyle \frac{10\mathrm{arcsec}}{\theta _d}}\left({\displaystyle \frac{\nu _{rolloff}}{10^{17}\mathrm{Hz}}}\right)^{1/2}\left({\displaystyle \frac{(t_{age})}{10^{20}\mathrm{Hz}\mathrm{pc}^2}}\right)^{1/2}.`$ (11)
The thin solid and dashed lines in Figure 5 represent the relation between the age and the distance using eq.(11), (1), and (2).
On the other hand, $`D=R_s/\theta _R`$ is rewritten as
$`D`$ $`=`$ $`0.34\mathrm{kpc}{\displaystyle \frac{10\mathrm{arcmin}}{\theta _R}}{\displaystyle \frac{R_s(t_{age})}{1\mathrm{pc}}},`$ (12)
$`R_s(t_{age})`$ $`=`$ $`\{\begin{array}{cc}1.1\mathrm{pc}(\frac{E_{51}}{M_{ej}\rho _0}t_{age}{}_{}{}^{4})^{1/7}\hfill & t<t_{ST}\hfill \\ 1.2\mathrm{pc}\left(\frac{E_{51}}{\rho _0}\right)^{1/5}(t_{age}0.22E_{51}{}_{}{}^{1/2}M_{ej}^{}{}_{}{}^{5/6}\rho _{0}^{}{}_{}{}^{1/3})^{2/5}\hfill & t>t_{ST}\hfill \end{array}.`$ (15)
The function $`R_s(t_{age})`$ is given by Truelove & McKee (1999) (see eq.(1) and (2), and table 6 and 7), which studies the shock dynamics of SNRs, where $`E_{51}`$, $`M_{ej}`$, $`\rho _0`$, and $`t_{ST}`$ is the explosion energy in the unit of $`10^{51}`$ ergs, the ejecta mass in the unit of $`M_{}`$, the ambient density in the unit of g cm<sup>-3</sup>, and the time scale that SNRs enter the Sedov-Taylor phase (in which shock begins to deaccelerate due to the interstellar medium), respectively (Truelove & McKee, 1999). For simplicity, we consider the case of the constant interstellar medium, $`n_0`$,
$$n_0\frac{\rho _0}{\mu _H}$$
where $`\mu _H`$ is the mean mass per hydrogen nucleus, $`1.4\times 1.67\times 10^{24}`$ g. In the following, we choose the kinetic energy of the ejecta, the ejecta mass, and $`n_0`$ to be $`10^{51}`$ ergs, 1.4$`M_{}`$, and 0.1 cm<sup>-3</sup>, respectively. Solving these equations, one can estimate both $`D`$ and $`t_{age}`$.
Figure 5 represents the relation between the age and the distance. We derived the allowed range of the parameters to be
$`t_{age}`$ $`=`$ $`660(4201400)\mathrm{yrs},`$ (16)
$`D`$ $`=`$ $`0.33(0.260.50)\mathrm{kpc}.`$ (17)
When we vary the ambient number density into $`n_0=0.5\mathrm{cm}^3`$ or the ejecta mass into $`M_{ej}=10M_{}`$, the result remains basically unchanged as can be seen in Figure 5.
Although the allowed regions are too wide to derive some conclusion, we may be able to say that Vela Jr. is a nearby and relatively young SNR. The results are consistent with previous reports, which may imply that this is an indirect confirmation of $``$—age relation (Bamba et al., 2005) and our distance/age indicator may become a useful tool in the future. Furthermore, derived distance is significantly larger than that of Vela SNR (250 pc; Cha et al., 1999), then there may be no interaction between these two SNRs.
The most common distance indicator of SNR is an application of the relation between the surface brightnesses at 1 GHz and the diameters of SNRs ($`\mathrm{\Sigma }`$$`D`$ relation; Case & Bhattacharya, 1998), however, the relation has still large uncertainties. Recent X-ray surveys recognize a new type of SNRs, which are dim in radio band and bright in hard X-ray band, such as RX J1713.7$``$3946 (Koyama et al., 1997), G28.6$``$0.1 (Bamba et al., 2001; Ueno et al., 2003), and so on (e.g. Combi et al., 2005; Yamaguchi et al., 2004). The number of such SNRs may be more than $``$20 (Bamba et al., 2003a). They must be significant accelerators of cosmic rays because most of them emit synchrotron X-rays, Unfortunately, they are dim in radio bands, so that the $`\mathrm{\Sigma }`$$`D`$ relation cannot directly applied. Our method may also be an useful tool to estimate both $`D`$ and $`t_{age}`$ of such synchrotron X-ray emitting SNRs.
Again let us consider $`E_{max}`$ and $`B_d`$ using Eqs. (5) and (6). Assuming that $`D300\mathrm{pc}`$, that is the most preferable value of our distance indicator, we obtain $`E_{max}3`$ TeV and $`B_d5\times 10^2\mu `$G. Although the uncertainty is very large, our result may indicate that the magnetic field is highly amplified.
These results does not include the uncertainty of the model of SNR dynamics and the function $``$. Considered these uncertainty, the allowed regions for $`t_{age}`$ and $`D`$ become more larger. It is the uncertainty of $``$’s normalization which makes error regions the most widest, about 50%. In order to improve this method for $`t_{age}`$ and $`D`$, more samples are needed to make the uncertainty of $``$ smaller.
## 9 Summary
We have conducted systematic spectral and spatial analysis of filamentary structures in Vela Jr. NW rim for the first time. A summary of our results is as follows:
1. We found that nonthermal X-rays from Vela Jr. NW rim are concentrated on very thin filamentary structures. The average scale width on the upstream side is similar to or smaller than the PSF size of Chandra, whereas that on the downstream side is 49.5 (36.0–88.8) arcsec.
2. The spectra of filaments are hard and have no line-like structure, which is well reproduced with an absorbed power-law model of $`\mathrm{\Gamma }`$ = 2.67 (2.55–2.77), or SRCUT model with $`\nu _{rolloff}`$ = 4.3 (3.4–5.3)$`\times 10^{16}`$ Hz under the assumption of $`p=0.3`$.
3. We tried to estimate the distance $`(D)`$ and the age $`(t_{age})`$ of Vela Jr. using the function $``$ and estimated that $`t_{age}`$ = 660 (420–1400) yrs and $`D`$ = 0.33 (0.26–0.50) kpc, which is consistent with previous reports. These results may suggest that there is no interaction between Vela SNR and Vela Jr.
4. Using the estimated $`D`$, we derived the most preferable values $`E_{max}3\mathrm{TeV}`$ and $`B_d500\mu `$G. Our result may imply that the magnetic field on the filament is highly amplified.
Our particular thanks are due to the anonymous referee, K. Makishima, F. Takahara, Y. Mochizuki, Y. Moriguchi, Y. Uchiyama, M. Tsujimoto, J. Vink, and K. Ebisawa, for their fruitful discussions and comments. R.Y. and J.S.H. are supported by JSPS Research Fellowship for Young Scientists.
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# Correlations at intermediate 𝑝_𝑇
## 1 Introduction
The study of the correlations at intermediate $`p_T`$ in heavy-ion collisions at high energy is important for the understanding of the interaction between a hard parton and the hot dense medium that it traverses. By intermediate $`p_T`$ we mean the region that stands between the soft region ($`p_T<2`$ GeV/c) where the recombination of thermal partons is most important and the hard region ($`p_T>9`$ GeV/c) where the fragmentation of partons is dominant. Note that this classification of regions is determined by the modes of hadronization, which we shall review briefly, rather than by the nature of scattering, soft or hard. In the intermediate $`p_T`$ region the recombination of the thermal and shower partons is more important than any other component of hadronization and clearly conveys the medium effect on hard partons. And it is in that region where the recent analysis of the data from RHIC reveals a wealth of information on jet structure. We shall examine the properties of correlation in the framework of parton recombination, which is the only viable hadronization scheme that can account for the species dependence of the particles produced. Our emphasis will be on near-side correlation, which depends mainly on the correlation among shower partons in a jet. The away-side correlation involves other issues besides hadronization and will be the subject of a future investigation.
## 2 Single-particle distribution
Before discussing two-particle correlation, it is fitting to review first the single-particle distribution as determined in the recombination model . In that model the shower partons in a jet play a crucial role. They are semi-hard and can recombine with soft thermal partons on the one hand, and also with one another on the other hand. Their distributions cannot be calculated in perturbative QCD, but can be determined phenomenologically from the fragmentation functions (FF), which are themselves determined by fitting the fragmentation processes in the collisions of simple systems. In the framework of parton recombination the shower parton distributions (SPD) can be extracted from the fragmentation function $`D(x)`$ by use of the equation
$`xD_i^\pi (x)={\displaystyle \frac{dx_1}{x_1}\frac{dx_2}{x_2}\{S_i^j(x_1),S_i^j^{}\left(\frac{x_2}{1x_1}\right)\}R_\pi (x_1,x_2,x)},`$ (1)
where $`i`$ specifies the type of hard parton that fragments, $`j`$ and $`j^{}`$ denotes the types of two partons that recombine, and $`R_\pi `$ is the recombination function (RF) for the formation of a pion. The two SPDs, $`S_i^j`$ and $`S_i^j`$, are symmetrized in the order of emission with momentum fractions $`x_1`$ and $`x_2`$ \[see (12) below for the details\]. Five such parton distributions have been determined from five types of $`D(x)`$ functions . The RF for pion is
$`R_\pi (x_1,x_2,x)={\displaystyle \frac{x_1x_2}{x}}\delta (x_1+x_1x),`$ (2)
and is inferred from pion-induced Drell-Yan process; for proton formation the details are given also in .
In heavy-ion collisions the probability of finding a shower parton $`j`$ is
$`𝒮^j(q)=\xi {\displaystyle \underset{i}{}}{\displaystyle 𝑑kkf_i(k)S_i^j(q/k)},`$ (3)
where $`f_i(k)`$ is the probability of producing a hard parton of species $`i`$ at transverse momentum $`k`$ . $`\xi `$ is the average fraction of the number of hard partons that emerge from the bulk medium to hadronize in vacuum.
The thermal parton distribution (TPD) is determined by fitting the soft pion distribution at $`p_T<2`$ GeV/c by use of the recombination formula
$`{\displaystyle \frac{dN_\pi }{pdp}}={\displaystyle \frac{1}{p^2}}{\displaystyle \frac{dq_1}{q_1}\frac{dq_2}{q_2}F_{q\overline{q}}(q_1,q_2)R_\pi (q_1,q_2,p)}`$ (4)
where $`p_T`$ is denoted by $`p`$, for brevity. For TPD we use the factorizable form
$`F_{q\overline{q}}(q_1,q_2)=𝒯(q_1)𝒯(q_2),`$ (5)
where
$`𝒯(q)=Cqe^{q/T}.`$ (6)
It is found from the low-$`p_T`$ data of pions that
$`C=23.2(\mathrm{GeV}/\mathrm{c})^1,T=0.317\mathrm{GeV}/\mathrm{c}`$ (7)
for central Au-Au collisions. For non-central collisions the parameters are given in .
With these basic quantities specified we can now describe how the pion distribution can be determined for any $`p_T`$ by use of the same equation (4), but with the two-parton distribution generalized to include the shower partons.
$`F_{q\overline{q}}(q_1,q_2)=𝒯(q_1)𝒯(q_2)+𝒯(q_1)𝒮(q_2)+𝒮𝒮(q_1,q_2)`$ (8)
where the last term is written in that way to emphasize that it is not factorizable, i.e.,
$`(𝒮𝒮)(q_1,q_2)=\xi {\displaystyle \underset{i}{}}{\displaystyle 𝑑kkf_i(k)\{S_i\left(\frac{q_1}{k}\right),S_i\left(\frac{q_2}{kq_1}\right)\}}.`$ (9)
In view of (1) and (4) it should be clear that the $`𝒮𝒮`$ terms in $`F_{q\overline{q}}`$ leads to fragmentation
$`{\displaystyle \frac{dN_\pi ^{𝒮𝒮}}{pdp}}={\displaystyle \frac{\xi }{p}}{\displaystyle \underset{i}{}}{\displaystyle 𝑑kf_i(k)D_i^\pi \left(\frac{p}{k}\right)}.`$ (10)
What is new is the $`𝒯𝒮`$ term in (8); it dominates in the intermediate $`p_T`$ region, as evidenced in Fig. 1, in which the overall normalization is adjusted to fit the data by letting $`\xi `$ be 0.07. The shape of the $`p_T`$ dependence is a prediction of the model. In that figure the shower-shower (2 jet) line corresponds to the recombination of shower partons arising from two different jets, and should be ignored for collisions at RHIC energies. The dominance of $`𝒯𝒮`$ recombination in the $`3<p_T<9`$ GeV/c region cannot be reproduced by fragmentation even if the FF used is medium-modified because the momentum fraction $`x`$ in the FF requires the parton momentum to be greater than the pion momentum $`p`$, whereas the RF requires the coalescing parton momenta to be less than $`p`$. Since the parton momenta are damped by a power law, the latter process always wins.
The contrast between the two processes of hadronization becomes even more pronounced in the case of proton production. Since three quarks recombine to form a proton, the average parton momentum is $`p/3`$, so they are even more abundantly available. To form a proton by fragmentation, one pays a heavy penalty to produce a high $`k`$ parton, and then pays an even heavier penalty to require that it fragments into a proton, the FF for which is an order of magnitude smaller than $`D^\pi `$. This is why the $`p/\pi `$ ratio can be high in the recombination model but very small in the fragmentation model. The production of proton in central AuAu collisions has been calculated in the recombination model where $`𝒯𝒯𝒮`$ and $`𝒯𝒮𝒮`$ components have been found to be more important than $`𝒮𝒮𝒮`$ component (i.e., fragmentation) for $`p_T<9`$ GeV/c . This is shown in Fig. 2, where the data exist only up to 4 GeV/c. But that is enough to exhibit the large $`p/\pi `$ ratio , as shown in Fig. 3, where the dashed line takes the proton mass into account at low $`p_T`$ . Actually, the large $`p/\pi `$ ratio was obtained in an earlier paper on recombination even before the shower parton distributions were obtained. The parton distributions there were inferred from the pion distribution. Two other groups have also obtained similar results using recombination/coalescence model .
The Cronin effect has for thirty years been referred to as the manifestation of $`k_T`$ broadening by multiple scattering in the initial state of pA collisions. That relationship does not take into account of the fact that the experimental $`p_T`$ spectrum in $`p+Ah+X`$ depends on $`A`$ as $`A^{\alpha _h}`$, where $`\alpha _p>\alpha _\pi `$ . If the effect of the nuclear medium on hard scattering is before fragmentation, then the exponent $`\alpha _h`$ should be independent of whether the hadron $`h`$ is a pion or a proton. In reality, not only is $`\alpha _p>\alpha _\pi `$ experimentally, the FF for proton $`D^p`$ is much smaller than that for pion, $`D^\pi `$, by roughly an order of magnitude. This failure in interpreting the data has been corrected by use of parton recombination as the hadronization mechanism. We have studied the production of hadrons (pion and proton) at intermediate $`p_T`$ in d-Au collisions at all centralities in the recombination model . Fig. 4 shows our results on $`R_{CP}`$ for pion and proton. Evidently, we obtain $`R_{CP}^p>R_{CP}^\pi `$ in the range $`1<p_T<3`$ GeV/c, in good agreement with the data . This result may be regarded as the strongest support for the recombination model, since no other approaches have indicated the possibility of attaining the same.
One last feature of single-particle distributions that we choose to mention here is the suppression of $`R_{CP}`$ in forward production. BRAHMS data on d-Au collision have shown that $`R_{CP}`$ is as low as 0.5 at $`\eta =3.2`$. Such a suppression of central production has been interpreted as suggestive evidence for color glass condensate, since at large $`\eta `$ the small-$`x`$ nuclear partons are presumed to be important and their high density there reveals saturation physics. We have, however, calculated the $`p_T`$ spectra at large $`\eta `$ in the recombination model without incorporating any exotic physics, and found results in agreement with the data . Our input is the data on $`dN_{ch}/d\eta `$ which decreases with increasing $`\eta `$ much more rapidly for central d-Au collisions than for peripheral collisions. Since $`dN_{ch}/d\eta `$ is dominated by soft partons, the $`𝒯𝒮`$ recombination results in the corresponding decrease of the hadron distributions at large $`\eta `$. There is no change of the underlying physics as $`\eta `$ is carried from backward to forward direction.
## 3 Parton and hadron correlations in jets
Having established some degree of reliability of parton recombination in the treatment of single-particle distributions, we now consider correlation of particles in jets at high $`p_T`$. There are two ways to study correlation: one is to use trigger particles to select events in which the associated particles reveal near-side and away-side characteristics; the other is to treat the two particles on equal footing and study the two-particle correlation function. Both approaches have been adopted by different groups within the STAR collaboration. We have some recent results on both that can be reported here, starting with the latter.
In general, the two-particle correlation is defined by
$`C_2(1,2)=\rho _2(1,2)\rho _1(1)\rho _1(2).`$ (11)
where $`\rho _2(1,2)`$ is the two-particle distribution, and $`\rho _1(1)`$ is the one-particle distribution, which for pion production at intermediate $`p_T`$ is given by (4) in the recombination model. It is, however, more appropriate to discuss first, not particle correlation, but parton correlation in a jet. Consider a hard parton with a fixed momentum $`k`$ in vacuum, as in $`e^+e^{}`$ annihilation. Since we shall discuss the correlation in terms of momentum fractions $`x_i`$, it does not matter what $`k`$ is so long as it is high enough. Thus for a single shower parton, we have $`\rho _1(1)=S_i^j(x_1)`$.
For two shower partons we have
$`\rho _1(1,2)=\{S_i^j(x_1),S_i^j^{}\left({\displaystyle \frac{x_2}{1x_1}}\right)\}={\displaystyle \frac{1}{2}}\left[S_i^j(x_1)S_i^j^{}\left({\displaystyle \frac{x_2}{1x_1}}\right)+S_i^j\left({\displaystyle \frac{x_1}{1x_2}}\right)S_i^j^{}(x_2)\right]`$ (12)
which guarantees that $`x_1+x_21`$ and symmetrizes the order of emission. Evidently, the two partons are correlated by virtue of the form in (12). If we define the normalized distribution by the ratio
$`r_2(1,2)={\displaystyle \frac{\rho _2(1,2)}{\rho _1(1)\rho _1(2)}},`$ (13)
then the result of our calculation for $`r_2(1,2)`$ is shown in Fig. 5 . It is clear that there is correlation for almost all $`x_1`$ and $`x_2`$ except when they are very small, where $`r_2(1,2)1`$.
Now, consider two shower partons in a jet in heavy-ion collisions. In that case the hard parton momentum $`k`$ is not fixed, so the corresponding $`\rho _1`$ and $`\rho _2`$ involve integrals over $`k`$, i.e.,
$`\rho _1(1)=𝒮(q_1),\rho _2(1,2)=(𝒮𝒮)^{jj^{}}(q_1,q_2),`$ (14)
where $`𝒮`$ and $`𝒮𝒮`$ are given in (3) and (9). The calculated results for $`r_2(1,2)`$ in that case are shown in Fig. 6 for both central and peripheral collisions . They become very large at large $`q_1`$ and $`q_2`$ because each $`\rho _1`$ is power damped at large $`k`$, as is $`\rho _2`$. The drastic difference between Figs. 5 and 6 underscores the effect of hard scattering in heavy-ion collision even when the only correlation in the problem is the same in both cases.
The correlation between pions in jets is far more complicated to calculate because of the many ways that partons can recombine. The two-pion distribution is
$`\rho _2(1,2)={\displaystyle \frac{dN_{\pi _1\pi _2}}{p_1p_2dp_1dp_2}}={\displaystyle \frac{1}{(p_1p_2)^2}}{\displaystyle \left(\underset{i}{}\frac{dq_i}{q_i}\right)F_4(q_1,q_2,q_3,q_4)R(q_1,q_3,p_1)R(q_2,q_4,p_2)}.`$ (15)
where
$`F_4=(𝒯𝒯+𝒮𝒯+𝒮𝒮)_{13}(𝒯𝒯+𝒮𝒯+𝒮𝒮)_{24}.`$ (16)
While many parts of $`F_4`$ are factorizable, and therefore make no contribution to $`C_2(1,2)`$, there are non-factorizable parts that involve at least one $`𝒮`$ in each of $`(\mathrm{})_{13}`$ and $`(\mathrm{})_{24}`$, the most important example of which is $`(𝒮𝒯)_{13}(𝒮𝒯)_{24}`$. The two $`𝒮`$ terms in that component, involving the shower partons $`S(q_1)`$ and $`S(q_2)`$ are correlated because they are in the same jet. Using (15) and (4) in (11), we obtain $`C_2(1,2)`$ which is shown in Fig. 7. There is not too much difference in the shapes of $`C_2`$ for the central and peripheral cases for most of $`p_1`$ and $`p_2`$, except when the momenta are small where $`C_2`$ becomes negative for central collisions and therefore cannot be exhibited in the log plot . The rapid decrease at large momenta is due to the power-law damping of $`f_i(k)`$ in both $`\rho _1`$ and $`\rho _2`$. In the lower part of the intermediate $`p_T`$ region various components of $`C_2`$ become negative, as shown in Fig. 8. That occurs because of the competition for the shower parton momenta in a jet, when the hard parton momentum $`k`$ can be low enough to avoid the severe suppression of $`f_i(k)`$. Since the correlation of the shower partons is negative, as we have seen in Fig. 5, it is not surprising that $`C_2`$ for the hadrons also becomes negative when $`p_1`$ and $`p_2`$ are not too high.
To compensate for that suppression at high $`p_T`$ let us define
$`G_2(1,2)={\displaystyle \frac{C_2(1,2)}{\left[\rho _1(1)\rho _1(2)\right]^{1/2}}},`$ (17)
whose dependence on $`p_1`$ and $`p_2`$ can now be shown in linear plots. Fig. 9(a) shows that for central collisions $`G_2`$ becomes negative for $`p_1`$ and/or $`p_2\stackrel{<}{}4`$ GeV/c. The ratio, $`R_{CP}^{G_2}`$, of $`G_2`$ for the two extreme centralities exhibits a minimum at $`p_1p_22`$ GeV/c, as shown in Fig. 9(b) . Data on that ratio, thus far not analyzed, would be able to provide information on whether there exist any dynamical correlations that we have not incorporated in our calculation.
So far we have only considered the correlations in the momentum variables $`p_1`$ and $`p_2`$ of $`C_2(1,2)`$. We can also study the autocorrelation in $`\mathrm{\Delta }\eta `$ and $`\mathrm{\Delta }\varphi `$, for which there are data at low $`p_T`$ . For such studies we need information on the angular distribution of shower partons. With that goal in mind we turn next to the investigation of correlations with trigger particles selected to serve as reference.
## 4 Associated particle distributions
If we use $`p_1`$ to denote the transverse momentum of the trigger, and $`p_2`$ that of the associated particle, then the per-trigger distribution of the latter for pions is
$`{\displaystyle \frac{dN_\pi ^{AP}}{p_2dp_2}}={\displaystyle \frac{𝑑p_1p_1𝑑N_{\pi \pi }/p_1p_2𝑑p_1𝑑p_2}{𝑑p_1p_1𝑑N_\pi /p_1𝑑p_1}},`$ (18)
where $`p_1`$ is integrated over a range that corresponds to the experimental cut on the trigger momentum. The integrands in the numerator and denominator are, respectively, (15) and (4). The associated particle distribution (APD) has been calculated for both dAu and AuAu collisions at various centralities . It is found that the $`p_T`$ distribution of the APD for dAu collisions has negligible dependence on centrality, although the dependence is quite significant for AuAu collisions. Those results are summarized in Fig. 10. They are different from the results reported in because different quantities are calculated: whereas (18) corresponds to the ratio of integrals, that in is the integral of the ratio. Comparison with the experimental data should be made only when the appropriate quantity is chosen that corresponds to what is measured and analyzed.
Let us now go on to the angular dependence of the APD. To exhibit the structures of the near- and away-side jets, it is necessary to make background subtraction of the data. F. Wang has presented data on $`\mathrm{\Delta }\varphi `$ and $`\mathrm{\Delta }\eta `$ distributions, where the subtraction scheme used results in the vanishing of the APD in $`\mathrm{\Delta }\varphi `$ at $`\left|\mathrm{\Delta }\varphi \right|=1`$ . The corresponding distribution in $`\mathrm{\Delta }\eta `$ shows a pedestal on top of which sits a peak at $`\mathrm{\Delta }\eta =0`$. To address these features found in the data it is necessary for us to generalize the formalism that we have used for recombination. So far our consideration has only been one-dimensional (1D), where the parton and hadron momenta are all collinear. Now we must consider shower partons in a jet cone that has 3D characteristics. Furthermore, we must take into account the energy loss of the hard partons and the subsequent hadronization of the medium that has absorbed the radiated energy. These aspects of generalization have been considered in .
The fact that the APD in $`\mathrm{\Delta }\eta `$ has a pedestal, not found in the $`\mathrm{\Delta }\varphi `$ distribution, suggests the basic lack of symmetry between the longitudinal and azimuthal directions. Indeed, whereas there is longitudinal expansion of the compressed medium, there is no azimuthal expansion in the transverse plane, only radial expansion. That means that there is no mixing of the various $`\varphi `$ sectors, making possible the implication that the particles detected in the peak region with $`\left|\mathrm{\Delta }\varphi \right|<1`$ arise from partons, soft or hard, that are originally in the same $`\varphi `$ sector, i.e., where the trigger is measured in the azimuth. That is not the case with the $`\eta `$ variable due to longitudinal expansion, and therein lies the possibility of a pedestal outside the peak region in $`\mathrm{\Delta }\eta `$ where the trigger is.
Endowing the parton momenta with vectorial properties in 3D, we use $`\psi `$ to denote the angle between $`\stackrel{}{q}_2`$ and $`\stackrel{}{k}`$, assuming for simplicity that $`\stackrel{}{q}_1`$ is along $`\stackrel{}{k}`$, based on the recognition that it is the relative angle between $`\stackrel{}{q}_1`$ and $`\stackrel{}{q}_2`$ that matters. We further assume that there is enough dispersion of the thermal partons around the average direction at any $`(\eta ,\varphi )`$ such that the hadron momenta $`\stackrel{}{p}_1`$ and $`\stackrel{}{p}_2`$ due to $`𝒯𝒮`$ recombination can be directed along $`\stackrel{}{q}_1`$ and $`\stackrel{}{q}_2`$, respectively. That means that $`\psi `$ is also the angle between the measured pion momenta. Consequently, it is possible to relate $`\psi `$ to the pseudorapidities $`\eta _1`$ and $`\eta _2`$. Let us describe the angular distribution of the shower partons around the jet axis by a Gaussian
$`G(\psi ,x)=exp\left[\psi ^2/2\sigma ^2(x)\right],`$ (19)
where the width depends on the momentum fraction $`x`$ as
$`\sigma (x)=\sigma _0(1x),`$ (20)
which is a simple way to capture the property that the jet cone is wider for softer partons. $`\sigma _0`$ is a free parameter that is to be determined phenomenologically.
With the angular variable described above we can now write down the contribution to $`F_4`$ that gives rise to the trigger at $`\stackrel{}{p}_1`$ and the AP at $`\stackrel{}{p}_2`$ within the peak in $`\mathrm{\Delta }\eta `$ through $`𝒯𝒮`$ recombination for both pions
$`F_4^{TSTS}=\xi {\displaystyle \underset{i}{}}{\displaystyle 𝑑kkf_i(k)𝒯(q_3)\{S(q_1),S(q_2)\}𝒯(q_4)G(\psi ,q_2/k)}`$ (21)
where $`q_1`$ and $`q_3`$ form the trigger at $`p_1`$, and $`q_2`$ and $`q_4`$ form the AP at $`\stackrel{}{p}_2`$ at an angle $`\psi `$ relative to $`\stackrel{}{p}_1`$. $`𝒯(q)`$ has the same form as in (6), but the inverse slope $`T`$ is now allowed to be higher than the value used in the past in order to take into account the enhanced thermal distribution due to the loss of energy of the hard parton while traversing the medium. Since the enhanced thermal partons are in the immediate vicinity of the hard parton, they are the ones that recombine with the shower partons, as expressed in (21). How $`T`$ differs from the value $`T_0`$ determined from the soft hadron distribution for $`p_T<2`$ GeV/c without being in the presence of any jet is another parameter $`\mathrm{\Delta }T`$ in the problem. We determine $`\mathrm{\Delta }T`$ from the pedestal height. In the recombination model we are able to attribute the pedestal effect to the difference of $`𝒯𝒯`$ recombination when there is a jet and $`𝒯_0𝒯_0`$ recombination of the background in the absence of a jet. That is, for the pedestal we have
$`F_4^{ped}=\xi {\displaystyle \underset{i}{}}{\displaystyle 𝑑kkf_i(k)S(q_1)𝒯(q_3)\left[𝒯(q_2)𝒯(q_4)𝒯_0(q_2)𝒯_0(q_4)\right]}`$ (22)
where the part before the square brackets is for the trigger as in (21), while the quantity inside the square brackets is what remains of the thermal partons for recombination after background subtraction. Using the value $`𝒯_0=0.317`$ GeV/c we can determine $`T=T_0+\mathrm{\Delta }T`$ by varying $`\mathrm{\Delta }T`$ to fit the data on the pedestal.
The result of our calculation for the APD in $`\mathrm{\Delta }\eta `$ is shown in Fig. 11 . The two free parameters $`\sigma _0`$ and $`\mathrm{\Delta }T`$ are adjusted to get the good fit of the STAR data ; they are
$`\sigma _0=0.22,\mathrm{\Delta }T=15\mathrm{MeV}/\mathrm{c}.`$ (23)
The change from $`T_0`$ to $`T`$ is only 5%, so the difference is insignificant in the calculation of single-particle distribution. However, the difference is sufficient to give rise to a pedestal in $`\mathrm{\Delta }\eta `$, whose origin is therefore the feedback from the lost energy of the hard parton to the hadrons through the enhanced thermal medium. It should be recognized that the good fit in Fig. 11 is not a trivial consequence of the use of two parameters, since the height of the peak is the result of multidimensional integrals involving many components of recombination.
With the parameters in (23) fixed the APD in $`\mathrm{\Delta }\varphi `$ can be calculated; the result is shown in Fig. 12 . The dashed line corresponds to the pedestal in $`\mathrm{\Delta }\eta `$, but is forced to vanish at $`\left|\mathrm{\Delta }\varphi \right|=1`$ because of the subtraction scheme. The solid line includes the $`𝒯𝒮`$ contribution on top of the mount and exhibits good agreement with the data . The yield of the AP has been studied in on the basis of some correlation among the soft partons.
Although we have not put in by hand any short- or long-range correlation, one may interpret the correlation of the shower partons exhibited in Fig. 5 as an intrinsic short-range correlation in a jet, and the feedback mechanism of the energy loss of hard partons to the enhancement of thermal partons as a form of long-range correlation. Through parton recombination these correlations are transmitted from the partons to the hadrons that are measured.
The advantage of studying the particles associated with triggers is that the details of the jet structure become manifest and allow us to determine the properties of the shower partons, such as $`\sigma _0`$ in (23). The drawback is the necessity of making subtraction of the background that may involve some ambiguity. A way to avoid the drawback is to study autocorrelation in $`\mathrm{\Delta }\eta `$ and $`\mathrm{\Delta }\varphi `$ starting with either (11) or (17), in which the two particles are treated on equal footing and no subtraction beyond the definition of $`C_2(1,2)`$ is needed. That will be our next project, in which we can use the results of as the basis for the calculation of the autocorrelation in $`\mathrm{\Delta }\eta `$ and $`\mathrm{\Delta }\varphi `$ within a definite $`p_T`$ range.
## 5 Concluding remarks
In conclusion, it is worth stressing that parton recombination provides a framework to describe correlation at intermediate $`p_T`$ range. So far we have not assumed any exotic correlation among the partons, since none seems necessary. However, some dynamical correlation may be present when probed properly, in which case our formalism may be well suited to decipher its characteristics. Of course, it is the correct hadronization process that must first be established. After that we can then investigate not only the detail properties of the correlation in the near-side jet, but also the nature of jet quenching on the away side.
I am grateful to Gunther Roland and Tom Trainor for inviting me to this very stimulating workshop on Correlations and Fluctuations. I also want to thank my collaborators, C. B. Chiu, R. Fries, Z. Tan and C. B. Yang, who have been crucial in the completion of the work reviewed here, so the credit for the success of our approach should go mostly to them, without whom there are no results to show. This work was supported in part, by the U. S. Department of Energy under Grant No. DE-FG02-96ER40972.
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# A solution of the Schottky-Type problem for curves with automorphisms
## 1. Introduction
The objective of the Schottky problem is to characterize the principally polarized abelian varieties (p.p.a.v.) which are Jacobians of smooth algebraic curves. This problem was solved, in the framework of the theory of KP equations, by Shiota (\[Sh\]). The analogous problem for Prym varieties was studied and partially solved by Shiota (\[Sh2\]) in terms of the BKP hierarchy.
The Schottky problem for Pryms is also related to the characterization of Jacobians of algebraic curves which admit a non-trivial involution. The moduli space of curves with non-trivial automorphisms was studied in the previous paper (\[GMP\]), and the points of the Sato Grassmannian defined by those curves were characterized. Moreover, an explicit set of algebraic equations defining the moduli space of curves with automorphisms as a subscheme of the Sato Grassmannian were obtained.
In §3 of the present paper, these equations are given as an explicit hierarchy of differential equations for the $`\tau `$-functions (Theorems 3.14 and 3.15). The first non-trivial equations are studied in detail and related to the KP and KdV hierarchies.
These results do not solve the Schottky problem for Riemann surfaces with automorphisms. For solving this problem one must characterize the conditions that a theta function of a p.p.a.v. should satisfy in order to be a jacobian theta function of a curve with a non-trivial automorphism. Such a characterization is obtained in this paper, as follows.
Let $`X_\mathrm{\Omega }`$ be an irreducible p.p.a.v. of dimension $`g`$. For a natural number $`r`$, the $`r`$-th Baker-Akhiezer functions are defined from the theta function of $`X_\mathrm{\Omega }`$ and some auxiliary data; namely, an $`r`$-tuple of $`g\times \mathrm{}`$-matrices $`(A^{(1)},\mathrm{},A^{(r)})`$ of rank $`g`$ and $`r`$ symmetric quadratic forms $`Q^{(1)},\mathrm{},Q^{(r)}`$ (see subsection 4.B for precise definitions).
###### Theorem (Characterization).
Let $`X_\mathrm{\Omega }`$ be an irreducible p.p.a.v. of dimension $`g>1`$.
Then the following conditions are equivalent.
1. There exists a projective irreducible smooth genus $`g`$ curve $`C`$ with a non-trivial automorphism $`\sigma _C:CC`$ such that $`X_\mathrm{\Omega }`$ is isomorphic as a p.p.a.v. to the Jacobian of $`C`$.
2. There exist a prime number $`p`$, $`p`$ matrices $`A^{(1)},\mathrm{},A^{(p)}`$ ($`A^{(j)}`$ being a $`g\times \mathrm{}`$-matrix of rank $`g`$) and $`p`$ symmetric quadratic forms $`Q^{(1)},\mathrm{},Q^{(p)}`$, such that
1. for some $`\xi _0^g`$ the corresponding BA-functions satisfy the $`(1,\stackrel{𝑝}{\mathrm{}},1)`$-KP hierarchy
$$\mathrm{Res}\left(\underset{j=1}{\overset{p}{}}z^{\delta _{ju}\delta _{jv}}\psi _{u,\xi _0}^{(j)}(z,t)\psi _{v,\xi _0}^{(j)}(z,s)\right)dz=\mathrm{\hspace{0.17em}0},\text{ and }$$
2. there exists $`\xi _1^g`$ (depending on $`\xi _0`$) such that
$$\mathrm{Res}\left(\underset{j=1}{\overset{p}{}}z^{\delta _{ju}\delta _{jv}}\psi _{v+1,\xi _0}^{(j+1)}(z,\sigma ^{}(t))\psi _{u,\xi _1}^{(j)}(z,s)\right)dz=\mathrm{\hspace{0.17em}0}$$
where $`\sigma ^{}(t):=(t^{(p)},t^{(1)},t^{(2)},\mathrm{},t^{(p1)})`$.
This theorem can be understood as a translation into equations of the characterization theorems of §5, where they are stated in terms of orbits of finite dimension as in the approach of Mulase (\[M\]) and Shiota (\[Sh\]). The idea of our proof is similar to the first part of the paper of Shiota; in fact, our main ingredient is a generalization of Theorem 6 of Shiota to the case of the $`(1,\stackrel{𝑟}{\mathrm{}},1)`$-KP hierarchy (see Theorem 6.2).
Standard arguments allow us to reinterpret the above identities as a hierarchy of PDE for the associated tau functions. The final step of our program should consist of proving that the “first” non-trivial equations of the hierarchy suffice to characterize the theta functions of Riemann surfaces with non-trivial automorphisms. This problem has an analytical counterpart which is the analogue of Shiota’s proof of Novikov conjecture. We hope to study it in a future paper.
## 2. Preliminary results
### 2.A. Formal group schemes
We establish the notation and recall some of the results proved in the papers \[MP2\] and \[GMP\] that will be used in subsequent sections.
In what follows, $`V`$ will denote a finite $`((z))`$–algebra endowed with an action of the group $`/p`$, where $`p`$ is a prime number, such that its fixed subset $`V^{/p}`$ is equal to $`((z))`$. Let $`\sigma `$ denote a (fixed) generator of this subgroup of $`\mathrm{Aut}_{((z))}V`$.
Since $`V`$ is a finite $`((z))`$-algebra, there are canonical maps given by the trace and the norm
$`\mathrm{Tr}:V`$ $`((z))`$
$`\mathrm{Nm}:V`$ $`((z))`$
which map an element $`g`$ in $`V`$ to the trace (respectively norm) of the homothety of $`V`$ defined by $`g`$ (as a $`((z))`$-vector space).
Throughout the paper, we will consider only the following two cases.
1. Ramified case: $`V=V_\mathrm{R}=((z_1))`$, where the $`((z))`$–algebra structure is given by mapping $`z`$ to $`z_1^p`$ and $`\sigma `$ is given by $`\sigma (z_1)=\omega z_1`$, where $`\omega `$ is a primitive $`p`$-th root of $`1`$ in $``$. In this case we set $`V_\mathrm{R}^+=[[z_1]]`$ and $`V_\mathrm{R}^{}=z_1^1[z_1^1]`$.
2. Non-ramified case: $`V=V_{\mathrm{NR}}=((z_1))\times \mathrm{}\times ((z_p))`$, where the $`((z))`$–algebra structure is given by mapping $`z`$ to $`(z_1,\mathrm{},z_p)`$ and $`\sigma `$ is given by $`\sigma (z_i)=z_{i+1}`$ (for $`i<p`$) and $`\sigma (z_p)=z_1`$. In this case we set $`V_{\mathrm{NR}}^+=[[z_1]]\times \mathrm{}\times [[z_p]]`$ and $`V_{\mathrm{NR}}^{}=z_1^1[z_1^1]\times \mathrm{}\times z_p^1[z_p^1]`$.
###### Example 2.1.
Let $`C`$ be a projective irreducible smooth curve with an order $`p`$ automorphism, denoted by $`\sigma _C:CC`$.
If $`\pi :C\overline{C}:=C/<\sigma _C>`$ denotes the quotient map and $`p\overline{C}`$ is a smooth point, then the $`\widehat{𝒪}_{\overline{C},p}`$-algebra $`\widehat{𝒪}_{C,\pi ^1(p)}`$ is isomorphic to one of the above types.
The formal base curve is $`\widehat{C}:=\mathrm{Spf}[[z]]`$ and the formal spectral cover is $`\widehat{C}_V:=\mathrm{Spf}V^+`$. Let $`\mathrm{\Gamma }_V`$ be the formal group scheme representing the functor
$`\left\{\text{category of formal }\text{-schemes}\right\}`$ $``$ $`\left\{\text{category of groups}\right\}`$
$`S`$ $``$ $`\left(V\widehat{}_{}H^0(S,𝒪_S)\right)_0^{}`$
where the subscript $`0`$ denotes the connected component of the identity and the superscript denotes the invertible elements. Replacing $`V`$ by $`V^+`$ (respectively by $`1+V^{}`$) we define the subgroup $`\overline{\mathrm{\Gamma }}_V^+`$ (respectively $`\mathrm{\Gamma }_V^{}`$) and thus obtain the decomposition
$$\mathrm{\Gamma }_V=\mathrm{\Gamma }_V^{}\times \overline{\mathrm{\Gamma }}_V^+.$$
The formal Jacobian of the formal spectral cover is the formal group scheme $`𝒥(\widehat{C}_V):=\mathrm{\Gamma }_V^{}`$. A straightforward calculation shows that $`𝒥(\widehat{C}_V)`$ is the formal spectrum of the ring
$$𝒪\left(𝒥\left(\widehat{C}_V\right)\right)=\left\{\{t_1^{\left(1\right)},t_2^{\left(1\right)},\mathrm{}\}\right\}\widehat{}\mathrm{}\widehat{}\left\{\{t_1^{\left(r\right)},t_2^{\left(r\right)},\mathrm{}\}\right\}$$
where $`t_i^{(j)}`$ are indeterminates, $`\{\{t_1,t_2,\mathrm{}\}\}`$ denotes the inverse limit $`\underset{}{\mathrm{lim}}[[t_1,\mathrm{},t_n]]`$, and $`r=1`$ for the ramified case (in which case the superscript $`(1)`$ is dropped) and $`r=p`$ for the non-ramified one.
Replacing $`V`$, $`V^+`$ and $`V^{}`$ by $`((z))`$, $`[[z]]`$ and $`z^1[z^1]`$ respectively in the previous constructions, one obtains formal schemes $`\mathrm{\Gamma }`$, $`\overline{\mathrm{\Gamma }}^+`$ and $`𝒥(\widehat{C}):=\mathrm{\Gamma }^{}`$. It is straightforward that the canonical morphism $`((z))V`$ gives rise to $`\mathrm{\Gamma }\mathrm{\Gamma }_V`$ and that the trace and the norm yield corresponding morphisms $`\mathrm{\Gamma }_V\mathrm{\Gamma }`$.
Recall that the Abel map $`\varphi _V:\widehat{C}_V𝒥(\widehat{C}_V)`$ is the morphism corresponding to the $`\widehat{C}_V`$-valued point of $`\mathrm{\Gamma }_V`$ associated to the $`r`$-tuple of series
(2.2)
$$(\left(1\frac{\overline{z}_1}{z_1}\right)^1,\mathrm{},\left(1\frac{\overline{z}_r}{z_r}\right)^1)$$
where $`\widehat{C}_V\mathrm{Spf}\left([[\overline{z}_1]]\times \mathrm{}\times [[\overline{z}_r]]\right)`$, with $`r=1`$ in the ramified case and $`r=p`$ in the non ramified one.
###### Proposition 2.3.
The Albanese variety of $`\widehat{C}_V`$ is the pair $`(𝒥(\widehat{C}_V),\varphi _V)`$.
###### Proof.
This statement is proved in \[MP\] for the case $`V=((z))`$. The present case follows from that result and from the fact that the Albanese variety of a disjoint union is the product of the corresponding Albanese varieties. ∎
The natural morphism $`\widehat{C}_V\stackrel{\pi }{}\widehat{C}`$ induces two group homomorphisms; namely, the pull-back
$$\pi ^{}:𝒥(\widehat{C})𝒥(\widehat{C}_V)$$
and the Albanese
$$𝒥(\widehat{C}_V)𝒥(\widehat{C}).$$
The Albanese map coincides with the restriction of the $`\mathrm{Nm}`$ to $`𝒥(\widehat{C}_V)`$; it will also be denoted by $`\mathrm{Nm}`$.
Observe that the automorphism $`\sigma `$ of $`V`$ gives rise to automorphisms of $`\mathrm{\Gamma }_V`$ and $`𝒥(\widehat{C}_V)`$ (also denoted by $`\sigma `$) such that $`\mathrm{\Gamma }_V^\sigma =((z))`$ and $`𝒥(\widehat{C}_V)^\sigma =𝒥(\widehat{C})`$. In particular, it follows that $`\mathrm{Tr}`$ (respectively $`\mathrm{Nm}`$) maps an element $`g\mathrm{\Gamma }_V`$ to $`_{i=0}^{p1}\sigma ^i(g)`$ (respectively $`_{i=0}^{p1}\sigma ^i(g)`$).
### 2.B. Infinite Grassmannians
Recall that the infinite Grassmannian $`\mathrm{Gr}(V)`$ of the pair $`(V,V^+)`$ is a $``$-scheme, which is not of finite type, whose set of rational points is
$$\left\{\begin{array}{c}\text{subspaces }UV\text{ such that }UV/V^+\\ \text{has finite dimensional kernel and cokernel}\end{array}\right\}.$$
This scheme is equipped with the determinant bundle, $`\mathrm{Det}_V`$, which is the determinant of the complex of $`𝒪_{\mathrm{Gr}(V)}`$-modules
$$V/V^+\widehat{}_{}𝒪_{\mathrm{Gr}(V)},$$
where $``$ is the universal submodule of $`\mathrm{Gr}(V)`$ and the morphism is the natural projection. The connected components of the Grassmannian are indexed by the Euler–Poincaré characteristic of the complex. The connected component of index $`m`$ will be denoted by $`\mathrm{Gr}^m(V)`$.
The group $`\mathrm{\Gamma }_V`$ acts by homotheties on $`V`$, and this action gives rise to a natural action on $`\mathrm{Gr}(V)`$
$$\mathrm{\Gamma }_V\times \mathrm{Gr}(V)\mathrm{Gr}(V).$$
Furthermore, this action preserves the determinant bundle.
These facts allow us to introduce $`\tau `$-functions and Baker-Akhiezer functions of points of $`\mathrm{Gr}(V)`$. Let us recall the definition and some properties of these functions (\[MP2\], §3).
The determinant of the morphism $`V/V^+\widehat{}_{}𝒪_{\mathrm{Gr}(V)}`$ gives rise to a canonical global section
$$\mathrm{\Omega }_+H^0(\mathrm{Gr}^0(V),\mathrm{Det}_V^{}).$$
In order to extend this section to $`\mathrm{Gr}(V)`$ (in a non-trivial way), we fix elements $`v_mV^+`$ for $`m>0`$ such that $`dimV^+/v_mV^+=m`$. Setting $`v_m:=v_m^1`$ for $`m>0`$, we define $`\mathrm{\Omega }_+(U):=\mathrm{\Omega }_+(v_m^1U)`$ for $`U\mathrm{Gr}^m(V)`$.
Now, the $`\tau `$-function and BA functions will be introduced following \[MP2\]. The $`\tau `$-function of $`U`$, $`\tau _U(t)`$, is a function on $`𝒥(\widehat{C}_V)`$ and is introduced as a suitable trivialization of the function $`g\mathrm{\Omega }_+(gU)`$ for $`g𝒥(\widehat{C}_V)`$,
$$\tau _U(g)=\frac{\mathrm{\Omega }_+(gU)}{g\delta _U}$$
where $`\delta _U`$ is a non-zero element in the fibre of $`\mathrm{Det}_V^{}`$ over $`U`$.
Let $`t`$ be the set of variables $`(t^{(1)},\mathrm{},t^{(r)})`$ (where $`t^{(j)}=(t_1^{(j)},t_2^{(j)},\mathrm{})`$) and $`z_{}`$ denote $`(z_1,\mathrm{},z_r)`$. For $`1u,vr`$, let $`[z_v]:=(z_v,\frac{z_v^2}{2},\frac{z_v^3}{3},\mathrm{})`$, $`t+[z_v]:=(t^{(1)},\mathrm{},t^{(v)}+[z_v],\mathrm{},t^{(r)})`$, $`U_{uu}=U`$ and, if $`uv`$, $`U_{uv}:=(1,\mathrm{},z_u,\mathrm{},z_v^1,\mathrm{},1)U`$.
The $`u`$-th Baker-Akhiezer function of a point $`U\mathrm{Gr}(V)`$ is the $`V`$-valued function
$$\psi _{u,U}(z_{},t):=(\mathrm{exp}\left(\underset{i1}{}\frac{t_i^{\left(1\right)}}{z_1^i}\right)\frac{\tau _{U_{u1}}\left(t+\left[z_1\right]\right)}{\tau _U\left(t\right)},\mathrm{},\mathrm{exp}\left(\underset{i1}{}\frac{t_i^{\left(r\right)}}{z_r^i}\right)\frac{\tau _{U_{ur}}\left(t+\left[z_r\right]\right)}{\tau _U\left(t\right)}).$$
From the decomposition $`V=_{i=1}^r((z_i))`$ we can write
$$\psi _{u,U}(z_{},t)=(\psi _{u,U}^{\left(1\right)}(z_1,t),\mathrm{},\psi _{u,U}^{\left(r\right)}(z_r,t))\underset{i=1}{\overset{r}{}}\left(\left(\left(z_i\right)\right)\widehat{}_{}𝒪\left(𝒥\left(\widehat{C}_V\right)\right)\right).$$
Let $`p_n(t)`$ be the Schur polynomials and $`\stackrel{~}{}_{t^{(j)}}=(_{t_1^{(j)}},\frac{1}{2}_{t_2^{(j)}},\frac{1}{3}_{t_3^{(j)}},\mathrm{})`$. Using the identity $`\mathrm{exp}\left(_{i1}z_j^i\stackrel{~}{}_{t_i^{(j)}}\right)\tau _{U_{uj}}(t)=\tau _{U_{uj}}(t+[z_j])`$, we obtain
(2.4) $`\psi _{u,U}^{\left(j\right)}(z_j,t)`$ $`=\mathrm{exp}\left({\displaystyle \underset{i1}{}}{\displaystyle \frac{t_i^{\left(j\right)}}{z_j^i}}\right){\displaystyle \frac{\tau _{U_{uj}}\left(t+\left[z_j\right]\right)}{\tau _U\left(t\right)}}=`$
$`=\left({\displaystyle \underset{i0}{}}{\displaystyle \frac{p_i\left(t^{\left(j\right)}\right)}{z_j^i}}\right){\displaystyle \frac{\left(_{i0}p_i\left(\stackrel{~}{}_{t^{\left(j\right)}}\right)z_j^i\right)\tau _{U_{uj}}\left(t\right)}{\tau _U\left(t\right)}}.`$
The main property of these Baker-Akhiezer functions is that they can be understood as generating functions for $`U`$ as a subspace of $`V`$, as we recall next.
###### Theorem 2.5 (\[MP2\]).
Let $`U\mathrm{Gr}^m(V)`$. Then
$$\psi _{u,U}(z_{},t)=v_m^1(1,\mathrm{},z_u,\mathrm{},1)\underset{i>0}{}(\psi _{u,U}^{(i,1)}(z_1),\mathrm{},\psi _{u,U}^{(i,r)}(z_r))p_{ui,U}(t)$$
where
$$\left\{(\psi _{u,U}^{(i,1)}(z_1),\mathrm{},\psi _{u,U}^{(i,r)}(z_r))\right|i>0,\mathrm{\hspace{0.17em}1}ur\}$$
is a basis of $`U`$ and $`p_{ui,U}(t)`$ are functions in $`t`$.
Consider the following pairing
$`V\times V`$ $``$
$`(a,b)`$ $`\mathrm{Res}_{z=0}\mathrm{Tr}(a,b)dz.`$
Since it is non-degenerate, there is an involution of $`\mathrm{Gr}(V)`$ which maps any point $`U`$ to its orthogonal complement $`U^{}`$. This involution sends the connected component $`\mathrm{Gr}^m(V)`$ to $`\mathrm{Gr}^{1mp}(V)`$ in the ramified case and to $`\mathrm{Gr}^m(V)`$ in the non-ramified one.
Finally, the adjoint Baker-Akhiezer functions of $`U`$ are defined by
$$\psi _{u,U}^{}(z_{},t):=\psi _{u,U^{}}(z_{},t),$$
whose components are given by
(2.6)
$$\psi _{v,U}^{\left(j\right)}(z_j,s)=\left(\underset{i0}{}\frac{p_i\left(s^{\left(j\right)}\right)}{z_j^i}\right)\frac{\left(_{i0}p_i\left(\stackrel{~}{}_{s^{\left(j\right)}}\right)z_j^i\right)\tau _{U_{jv}}\left(s\right)}{\tau _U\left(s\right)}.$$
## 3. Moduli of curves with an order $`p`$ automorphism
The aim of this section is to give explicit differential equations that characterize the tau functions coming from curves with an order $`p`$ automorphism among the tau functions coming from curves.
Recall from §5 of \[GMP\] that $`^{\mathrm{}}(p,\mathrm{R})`$ (respectively $`^{\mathrm{}}(p,\mathrm{NR})`$) is the subscheme of $`\mathrm{Gr}(V_\mathrm{R})`$ (respectively $`\mathrm{Gr}(V_{\mathrm{NR}})`$) parametrizing $`(C,\sigma _C,x,t_x)`$ where $`C`$ is a curve, $`\sigma _C`$ is an order $`p`$ automorphism of $`C`$, $`x`$ is a smooth point of $`C`$ fixed under $`\sigma _C`$ (respectively an orbit consisting of $`p`$ pairwise distinct points) and $`t_x`$ is a formal parameter at $`x`$.
More precisely, if $`^{\mathrm{}}(1)`$ (respectively $`^{\mathrm{}}(p)`$) is the moduli space parameterizing $`(C,x,t_x)`$ where $`C`$ is a curve, $`x`$ is a smooth point (respectively $`p`$ pairwise distinct smooth points) and a formal parameter at $`x`$, then the Krichever map induces an immersion
$$\mathrm{Kr}:^{\mathrm{}}(1)\mathrm{Gr}(V_\mathrm{R})$$
(respectively $`\mathrm{Kr}:^{\mathrm{}}(p)\mathrm{Gr}(V_{\mathrm{NR}})`$) such that
(3.1)
$$^{\mathrm{}}(p,\mathrm{R})=^{\mathrm{}}(1)\mathrm{Gr}(V_\mathrm{R})^\sigma $$
(respectively $`^{\mathrm{}}(p,\mathrm{NR})=^{\mathrm{}}(p)\mathrm{Gr}(V_{\mathrm{NR}})^\sigma `$) where $`\mathrm{Gr}(V)^\sigma `$ denotes the set of points in $`\mathrm{Gr}(V)`$ fixed under the action of $`\sigma `$.
Therefore, our task consists of writing down the hierarchies corresponding to these subschemes.
### 3.A. Hierarchies for invariant subspaces
The automorphism of $`\mathrm{Gr}(V)`$ induced by $`\sigma :VV`$ will also be denoted by $`\sigma `$; it preserves the determinant bundle. If we denote by $`\mathrm{Gr}(V)^\sigma `$ the set of points in $`\mathrm{Gr}(V)`$ fixed under the action of $`\sigma `$, then it is known that $`\mathrm{Gr}(V)^\sigma `$ is a closed subscheme.
The action of $`\sigma `$ on $`\mathrm{Gr}(V)`$ can be easily described in terms of the Baker-Akhiezer and the tau functions, as we show next.
Let us begin by studying the ramified case. We denote
$$\lambda t=(\lambda t_1,\lambda ^2t_2,\mathrm{},\lambda ^jt_j,\mathrm{})\text{ for any }\lambda $$
and
$$\sigma ^{}(t):=\omega ^1t.$$
Then
$$\tau _{\sigma (U)}(t)=\tau _U(\omega ^1t)$$
(up to a constant) and
(3.2)
$$\psi _{\sigma (U)}(z_1,t)=\psi _U(\omega ^1z_1,\omega ^1t)=\sigma ^1\left(\psi _U(z_1,\sigma ^{}(t))\right)$$
where the expression on the r.h.s. corresponds to the action of $`\sigma `$ on $`V`$; that is, since $`\psi _U(z_1,t)`$ is $`V`$-valued, $`\sigma `$ acts on $`z_1`$ and acts trivially on $`t`$.
It is also known that a point $`U\mathrm{Gr}(V)`$ lies in $`\mathrm{Gr}(V)^\sigma `$ if and only if its BA functions satisfy the following identities (\[MP2\], Theorem 3.15).
(3.3)
$$\mathrm{Res}_{z=0}\mathrm{Tr}\left(\frac{1}{z_1}\psi _{\sigma (U)}(z_1,t)\psi _U^{}(z_1,s)\right)\frac{dz}{z}=\mathrm{\hspace{0.17em}0}$$
Then, one has the following
###### Theorem 3.4 (ramified case).
Let $`V=V_\mathrm{R}`$ and let $`U`$ be a closed point of $`\mathrm{Gr}(V)`$.
Then $`U`$ is a point of $`\mathrm{Gr}(V)^\sigma `$ if and only its $`\tau `$-function $`\tau _U`$ satisfies the following differential equations
(3.5)
$$\underset{\begin{array}{c}\beta _1+\beta _2\alpha _1\alpha _2=1\\ \alpha _1\beta _1j(modp)\end{array}}{}\left(D_{\lambda _1,\alpha _1}\left(\stackrel{~}{}_x\right)p_{\beta _1}\left(\stackrel{~}{}_x\right)D_{\lambda _2,\alpha _2}\left(\stackrel{~}{}_s\right)p_{\beta _2}\left(\stackrel{~}{}_s\right)\right)\tau _U\left(x\right)\tau _U\left(s\right)=0$$
for all Young diagrams $`\lambda _1,\lambda _2`$ and all integers $`j\{0,\mathrm{},p1\}`$, where the definition of $`D_{\lambda ,\alpha }`$ is
$$D_{\lambda ,\alpha }(\stackrel{~}{}_y)(f(y)):=\underset{\lambda \mu =(\alpha )}{}\chi _\mu (\stackrel{~}{}_y)(f(y))_{|_{y=0}}.$$
###### Proof.
Let $`U`$ be a closed point of $`\mathrm{Gr}(V)`$. The residue condition given by (3.3) is equivalent to the vanishing of the constant term of
$$\underset{j=1}{\overset{p}{}}\frac{1}{\omega ^jz_1}\psi _{\sigma \left(U\right)}(\omega ^jz_1,t)\psi _U^{}(\omega ^jz_1,s)=\underset{j=1}{\overset{p}{}}\frac{1}{\omega ^jz_1}\psi _U(\omega ^{j1}z_1,\omega ^1t)\psi _U^{}(\omega ^jz_1,s)$$
Using (2.4) and (2.6) this coefficient turns out to be
$$\begin{array}{c}\underset{j=1}{\overset{p}{}}\underset{\begin{array}{c}\beta _1+\beta _2\alpha _1\alpha _2=1\end{array}}{}\omega ^{j(\beta _2+\beta _1\alpha _1\alpha _21)+\alpha _1\beta _1}\hfill \\ \hfill p_{\alpha _1}(\omega ^1t)p_{\beta _1}(\stackrel{~}{}_{\omega ^1t})\tau _U(\omega ^1t)p_{\alpha _2}(s)p_{\beta _2}(\stackrel{~}{}_s)\tau _U(s).\end{array}$$
Since this coefficient must vanish for each $`p`$-th rooth $`\omega ^j`$ of $`1`$ (note that the case $`\omega ^0=1`$ is precisely the KP-hierarchy, see \[MP\]), we obtain the equivalent conditions
$$\underset{\begin{array}{c}\beta _1+\beta _2\alpha _1\alpha _2=1\\ \alpha _1\beta _1j(modp)\end{array}}{}p_{\alpha _1}(\omega ^1t)p_{\beta _1}(\stackrel{~}{}_{\omega ^1t})\tau _U(\omega ^1t)p_{\alpha _2}(s)p_{\beta _2}(\stackrel{~}{}_s)\tau _U(s)=\mathrm{\hspace{0.17em}0}$$
for all $`j`$ in $`\{0,\mathrm{},p1\}`$ and all $`t`$ and $`s`$.
Equivalently (substituting $`\omega ^1t`$ by $`x`$), we obtain
(3.6)
$$\underset{\begin{array}{c}\beta _1+\beta _2\alpha _1\alpha _2=1\\ \alpha _1\beta _1j(modp)\end{array}}{}p_{\alpha _1}(x)p_{\beta _1}(\stackrel{~}{}_x)\tau _U(x)p_{\alpha _2}(s)p_{\beta _2}(\stackrel{~}{}_s)\tau _U(s)=\mathrm{\hspace{0.17em}0}$$
for all values of $`x`$ and $`s`$, and for all $`j`$ in $`\{0,\mathrm{},p1\}`$.
Since the vanishing of a function $`f(x,s)`$ (such as the left hand side of (3.6)) for all values of $`x`$ and $`s`$ is equivalent to the vanishing of $`\chi _{\lambda _1}(\stackrel{~}{}_x)\chi _{\lambda _2}(\stackrel{~}{}_s)f(x,s)_{|_{x=0,s=0}}`$ for all Young diagrams $`\lambda _1`$ and $`\lambda _2`$, a calculation shows that (3.6) is equivalent to (3.5), thus proving the theorem. ∎
###### Remark 3.7.
We now compute the first equations in the previous statement and relate them to the KP equations. Observe that if $`E_j`$ denotes the l.h.s. of equation (3.5) for $`j\{0,\mathrm{},p1\}`$, then the KP hierarchy is $`E_0+\mathrm{}+E_{p1}=0`$. Let us make this explicit.
The first non trivial equation in the KP hierarchy, which corresponds to the Young diagrams $`\lambda _1=(1,1,1)`$ and $`\lambda _2=0`$, is the celebrated KP equation
$$\begin{array}{c}\tau _U(0)\chi _{(2,2)}(\stackrel{~}{}_s)\tau _U(s)_{|_{s=0}}p_1(\stackrel{~}{}_x)\tau _U(x)_{|_{x=0}}\chi _{(2,1)}(\stackrel{~}{}_s)\tau _U(s)_{|_{s=0}}\hfill \\ \hfill +p_2(\stackrel{~}{}_x)\tau _U(x)_{|_{x=0}}\chi _{(1,1)}(\stackrel{~}{}_s)\tau _U(s)_{|_{s=0}}=0.\end{array}$$
On the other hand, equations (3.5) for the same Young diagrams are
$$\{\begin{array}{cc}& p_1\left(\stackrel{~}{}_x\right)\tau _U\left(x\right)_{|_{x=0}}\chi _{(2,1)}\left(\stackrel{~}{}_s\right)\tau _U\left(s\right)_{|_{s=0}}=\mathrm{\hspace{0.17em}0}\hfill \\ & \tau _U\left(0\right)\chi _{(2,2)}\left(\stackrel{~}{}_s\right)\tau _U\left(s\right)_{|_{s=0}}+p_2\left(\stackrel{~}{}_x\right)\tau _U\left(x\right)_{|_{x=0}}\chi _{(1,1)}\left(\stackrel{~}{}_s\right)\tau _U\left(s\right)_{|_{s=0}}=\mathrm{\hspace{0.17em}0},\hfill \end{array}$$
for $`p=2`$ and
$$\{\begin{array}{cc}& p_1\left(\stackrel{~}{}_x\right)\tau _U\left(x\right)_{|_{x=0}}\chi _{(2,1)}\left(\stackrel{~}{}_s\right)\tau _U\left(s\right)_{|_{s=0}}=\mathrm{\hspace{0.17em}0}\hfill \\ & \tau _U\left(0\right)\chi _{(2,2)}\left(\stackrel{~}{}_s\right)\tau _U\left(s\right)_{|_{s=0}}=\mathrm{\hspace{0.17em}0}\hfill \\ & p_2\left(\stackrel{~}{}_x\right)\tau _U\left(x\right)_{|_{x=0}}\chi _{(1,1)}\left(\stackrel{~}{}_s\right)\tau _U\left(s\right)_{|_{s=0}}=\mathrm{\hspace{0.17em}0}.\hfill \end{array}$$
for $`p2`$
###### Example 3.8.
As an example we give some more equations from (3.5) for other pairs of Young diagrams. Note that the corresponding equations in the KP hierarchy are all trivial.
1. Consider the Young diagrams $`\lambda _1=0`$ and $`\lambda _2=(1)`$ and the corresponding $`p`$ equations given by (3.5).
If $`p=2`$, both equations are trivial.
However, if $`p2`$, two of the equations (3.5) (the cases $`\alpha _1\beta _10(modp)`$ and $`\alpha _1\beta _12(modp)`$) are equivalent to the following
$$\tau _U(0)p_2(\stackrel{~}{}_x)\tau _U(x)_{|_{x=0}}=0.$$
Similarly, the consideration of $`\lambda _1=(1)`$ and $`\lambda _2=0`$ in (3.5) yields trivial equations for $`p=2`$, whereas for $`p2`$ we obtain
$$\tau _U(0)p_2(\stackrel{~}{}_x)\tau _U(x)_{|_{x=0}}=0.$$
2. More generally, for $`n2`$ and not divisible by $`p`$ we obtain
$`\tau _U(0)p_n(\stackrel{~}{}_t)\tau _U(t)_{|_{t=0}}`$ $`=0`$
$`\tau _U(0)p_n(\stackrel{~}{}_t)\tau _U(t)_{|_{t=0}}`$ $`=0`$
when considering (3.5) with $`\lambda _i=(n1)`$ and $`\lambda _j=0`$.
3. When $`n2`$ is not divisible by $`p`$, we obtain
$`p_1(\stackrel{~}{}_t)\tau _U(t)_{|_{t=0}}p_{n+1}(\stackrel{~}{}_s)\tau _U(s)_{|_{s=0}}`$ $`=0`$
$`p_1(\stackrel{~}{}_t)\tau _U(t)_{|_{t=0}}p_{n+1}(\stackrel{~}{}_s)\tau _U(s)_{|_{s=0}}`$ $`=0`$
by considering $`\lambda _i=(n1,1)`$ and $`\lambda _j=0`$.
4. When $`n1`$ is not divisible by $`p`$, the equations become
$`p_2(\stackrel{~}{}_t)\tau _U(t)_{|_{t=0}}p_{n+2}(\stackrel{~}{}_s)\tau _U(s)_{|_{s=0}}`$ $`=0`$
$`p_2(\stackrel{~}{}_t)\tau _U(t)_{|_{t=0}}p_{n+2}(\stackrel{~}{}_s)\tau _U(s)_{|_{s=0}}`$ $`=0`$
for $`\lambda _i=(n+1,2)`$ and $`\lambda _j=0`$.
###### Remark 3.9.
Let us study the relation with the KdV hierarchy. Consider an invariant point $`U\mathrm{Gr}(V)^\sigma `$ (in the ramified case). If we impose the condition that $`[z^1]U=U`$, then we obtain the $`p`$-KdV hierarchy. To see this, recall that the $`\tau `$-function of $`U`$, $`\tau _U`$, is (up to a constant) the pullback of the global section $`\mathrm{\Omega }_+`$ to $`\mathrm{\Gamma }_V`$ by
$$\mathrm{\Gamma }_V\times \{U\}\mathrm{Gr}(V).$$
Since the condition means that $`\mathrm{\Gamma }^{}U=U`$, we obtain the following diagram
where
$$𝒫:=\left\{g𝒥\left(\widehat{C}_V\right)\right|\mathrm{Nm}\left(g\right)=1\}=\left\{\mathrm{exp}\left(t_iz_1^i\right)\mathrm{\Gamma }_V\right|t_i=0\text{ for }i=\dot{p}\}.$$
Therefore, for the ramified case and for $`p=2`$, it makes sense to write $`\tau _U=\tau _U(t_1,t_3,\mathrm{})`$ as an element of $`𝒪(𝒫)=\{\{t_1,t_3,\mathrm{}\}\}`$ so that the $`\tau `$-function only depends on $`t_i`$ with $`i`$ odd. The resulting hierarchy is the classical KdV hierarchy (see also \[SW\], Proposition 5.11).
Now we focus in the non-ramified case. Denote
$$\sigma ^{}(t):=(t^{(p)},t^{(1)},t^{(2)},\mathrm{},t^{(p1)});$$
then the corresponding relations are as follows.
$`\tau _{\sigma \left(U\right)}\left(t\right)`$ $`=\tau _{\sigma \left(U\right)}\left((t^{\left(1\right)},\mathrm{},t^{\left(p\right)})\right)=`$
$`=\tau _U\left((t^{\left(p\right)},t^{\left(1\right)},t^{\left(2\right)},\mathrm{},t^{\left(p1\right)})\right)=\tau _U\left(\sigma ^{}\left(t\right)\right)`$
(up to a constant) and
(3.10) $`\psi _{u,\sigma \left(U\right)}(z_{},t)`$ $`=(\psi _{u+1,U}^{\left(2\right)}(z_1,\sigma ^{}\left(t\right)),\psi _{u+1,U}^{\left(3\right)}(z_2,\sigma ^{}\left(t\right)),\mathrm{},\psi _{u+1,U}^{\left(1\right)}(z_p,\sigma ^{}\left(t\right)))=`$
$`=\sigma ^1\left(\psi _{u+1,U}(z_{},\sigma ^{}\left(t\right))\right)`$
where the action of $`\sigma `$ on the r.h.s. is the action on $`V`$-valued functions.
It is also known that a point $`U\mathrm{Gr}(V)`$ lies in $`\mathrm{Gr}(V)^\sigma `$ if and only if its BA functions satisfy the following identities (\[MP2\])
(3.11)
$$\mathrm{Res}_{z=0}\mathrm{Tr}\left(\frac{\psi _{u,\sigma (U)}(z_{},t)}{(1,\mathrm{},z_u,\mathrm{},1)}\frac{\psi _{v,U}^{}(z_{},s)}{(1,\mathrm{},z_v,\mathrm{},1)}\right)dz=\mathrm{\hspace{0.17em}0}$$
for all $`u,v\{1,\mathrm{},p\}`$ where $`(1,\mathrm{},z_u,\mathrm{},1)`$ denotes the element of $`V`$ with entries equal to $`1`$ except the $`u`$-th, which is $`z_u`$.
In a similar manner to the ramified case but this time using (3.11), we obtain the following result.
###### Theorem 3.12 (non-ramified).
Let $`V=V_{\mathrm{NR}}`$ and let $`U`$ be a closed point of $`\mathrm{Gr}(V)`$.
Then $`U`$ is a point of $`\mathrm{Gr}(V)^\sigma `$ if and only the $`\tau `$-functions of $`U`$, $`\tau _{U_{uv}}`$, satisfy the following differential equations.
$$\begin{array}{c}\underset{j=1}{\overset{p}{}}\underset{\begin{array}{c}\beta _1+\beta _2\alpha _1\alpha _2\\ =\delta _{ju}+\delta _{jv}1\end{array}}{}(D_{\lambda _j,\alpha _1}(\stackrel{~}{}_{x^{\left(j+1\right)}})p_{\beta _1}\left(\stackrel{~}{}_{x^{\left(j+1\right)}}\right)\text{Ď}_\lambda _{}^{j+1}\left(\stackrel{~}{}_x\right)\tau _{U_{u+1,j+1}}\left(x\right)\hfill \\ \hfill D_{\mu _j,\alpha _2}\left(\stackrel{~}{}_{s^{\left(j\right)}}\right)p_{\beta _2}(\stackrel{~}{}_{s^{\left(j\right)}})\text{Ď}_\mu _{}^j\left(\stackrel{~}{}_s\right)\tau _{U_{jv}}\left(s\right))=0\end{array}$$
for all Young diagrams $`\lambda _1,\mu _1,\mathrm{},\lambda _p,\mu _p`$, where the differential operator $`\text{Ď}_\mu _{}^k`$ is given by
$$\text{Ď}_\mu _{}^k(\stackrel{~}{}_s)(f(s))=\underset{\mathrm{}k}{}\chi _\mu _{\mathrm{}}(\stackrel{~}{}_s^{\mathrm{}})(f(y))_{|_{s^{\mathrm{}}=0}}.$$
###### Proof.
Let $`U`$ be a closed point of $`\mathrm{Gr}(V)`$. The residue condition given by (3.11) is equivalent to the vanishing of
$$\begin{array}{c}\mathrm{Res}\left(\underset{j=1}{\overset{p}{}}z^{\delta _{ju}\delta _{jv}}\psi _{u,\sigma \left(U\right)}^{\left(j\right)}(z,t)\psi _{v,U}^{\left(j\right)}(z,s)\right)dz=\hfill \\ \hfill =\mathrm{Res}\left(\underset{j=1}{\overset{p}{}}z^{\delta _{ju}\delta _{jv}}\psi _{u+1,U}^{\left(j+1\right)}(z,\sigma ^{}\left(t\right))\psi _{v,U}^{\left(j\right)}(z,s)\right)dz\end{array}$$
Using (2.4) and (2.6) and denoting $`x=\sigma ^{}(t)`$, the coefficient of $`z^1`$ in the latter sum turns out to be
(3.13)
$$\begin{array}{c}\underset{j=1}{\overset{p}{}}\underset{\begin{array}{c}\beta _1+\beta _2\alpha _1\alpha _2\\ =\delta _{ju}+\delta _{jv}1\end{array}}{}p_{\alpha _1}(x^{(j+1})p_{\beta _1}(\stackrel{~}{}_{x^{(j+1)}})\tau _{U_{u+1,j+1}}(x)\hfill \\ \hfill p_{\alpha _2}(s^{(j)})p_{\beta _2}(\stackrel{~}{}_{s^{(j)}})\tau _{U_{jv}}(s).\end{array}$$
Since the vanishing of the function $`f(x,s)`$ one given by (3.13) for all values of $`x`$ and $`s`$ is equivalent to the vanishing of
$$\underset{1a,bp}{}\chi _{\lambda _a}(\stackrel{~}{}_{x^{(a)}})\chi _{\mu _b}(\stackrel{~}{}_{s^{(b)}})f(x,s)_{|_{x=0,s=0}}$$
for all Young diagrams $`\lambda _1,\mu _1,\mathrm{},\mu _p`$, a calculation shows that the vanishing of (3.13) is equivalent to (3.11), thus proving the theorem. ∎
### 3.B. Equations of the moduli space
Recalling the relation 3.1 and the Theorem 3.4 we can now write down the characterization of $`^{\mathrm{}}(p,\mathrm{R})`$ in terms of differential equations for the $`\tau `$-functions.
###### Theorem 3.14 (ramified case).
Let $`V=V_\mathrm{R}`$ and let $`U`$ be a closed point of $`\mathrm{Gr}^m(V)`$.
Then $`U`$ is a point of $`^{\mathrm{}}(p,\mathrm{R})`$ if and only if its $`\tau `$-function $`\tau _U(t)`$ satisfies the following set of differential equations, for all Young diagrams $`\lambda _1`$, $`\lambda _2`$, $`\lambda _3`$, $`\lambda `$, and all integers $`j`$ in $`\{0,1,\mathrm{},p1\}`$.
1. The equations (3.5):
$$\underset{\begin{array}{c}\beta _1+\beta _2\alpha _1\alpha _2=1\\ \alpha _1\beta _1j(modp)\end{array}}{}\left(D_{\lambda _1,\alpha _1}\left(\stackrel{~}{}_x\right)p_{\beta _1}\left(\stackrel{~}{}_x\right)D_{\lambda _2,\alpha _2}\left(\stackrel{~}{}_s\right)p_{\beta _2}\left(\stackrel{~}{}_s\right)\right)\tau _U\left(x\right)\tau _U\left(s\right)=0$$
2. $$\underset{\beta \alpha =m}{}\left(D_{\lambda ,\alpha }\left(\stackrel{~}{}_t\right)p_\beta \left(\stackrel{~}{}_t\right)\right)\tau _U\left(t\right)=0$$
3. $$\begin{array}{c}\underset{\beta _1+\beta _2+\beta _3\alpha _1\alpha _2\alpha _3=2m}{}(D_{\lambda _1,\alpha _1}(\stackrel{~}{}_t)p_{\beta _1}\left(\stackrel{~}{}_t\right)D_{\lambda _2,\alpha _2}(\stackrel{~}{}_s)p_{\beta _2}\left(\stackrel{~}{}_s\right)\hfill \\ \hfill D_{\lambda _3,\alpha _3}\left(\stackrel{~}{}_u\right)p_{\beta _3}(\stackrel{~}{}_3)\left)\tau _U\right(x\left)\tau _U\right(s\left)\tau _U\right(t)=0.\end{array}$$
Similarly, we can write down differential equations for $`^{\mathrm{}}(p,\mathrm{NR})`$ for the non ramified case. In this case the non-trivial curves with automorphisms appear when $`m0`$; if $`m=qp+f`$ with $`0f<p`$, then $`v_m=z_1^{q+1}\mathrm{}z_f^{q+1}z_{f+1}^q\mathrm{}z_p^q`$, and $`v_m=v_m^1`$.
###### Theorem 3.15 (non ramified case).
Let $`V=V_{\mathrm{NR}}`$ and let $`U`$ be a closed point of $`\mathrm{Gr}^m(V)`$.
Then $`U`$ is a point of $`^{\mathrm{}}(p,\mathrm{NR})`$ if and only if the following set of differential equations is satisfied by its $`\tau `$-functions $`\tau _{U_{uv}}(t)`$, for all Young diagrams $`\lambda _{}=\{\lambda _1,\mathrm{},\lambda _p\}`$, $`\mu _{}=\{\mu _1,\mathrm{},\mu _p\}`$, $`\nu _{}=\{\nu _1,\mathrm{},\nu _p\}`$ and all integers $`u,v,w`$ in $`\{1,2,\mathrm{},p\}`$.
1. The equations of Theorem 3.12
$$\begin{array}{c}\underset{j=1}{\overset{p}{}}\underset{\begin{array}{c}\beta _1+\beta _2\alpha _1\alpha _2\\ =\delta _{ju}+\delta _{jv}1\end{array}}{}(D_{\lambda _j,\alpha _1}(\stackrel{~}{}_{x^{\left(j+1\right)}})p_{\beta _1}\left(\stackrel{~}{}_{x^{\left(j+1\right)}}\right)\text{Ď}_\lambda _{}^{j+1}\left(\stackrel{~}{}_x\right)\tau _{U_{u+1,j+1}}\left(x\right)\hfill \\ \hfill D_{\mu _j,\alpha _2}\left(\stackrel{~}{}_{s^{\left(j\right)}}\right)p_{\beta _2}(\stackrel{~}{}_{s^{\left(j\right)}})\text{Ď}_\mu _{}^j\left(\stackrel{~}{}_s\right)\tau _{U_{jv}}\left(s\right))=0\end{array}$$
2. $$\begin{array}{c}\underset{j=1}{\overset{f}{}}\underset{\begin{array}{c}\beta \alpha \\ =\delta _{ju}2q\end{array}}{}\left(D_{\lambda _j,\alpha }\left(\stackrel{~}{}_{t^{\left(j\right)}}\right)p_\beta \left(\stackrel{~}{}_{t^{\left(j\right)}}\right)\text{Ď}_\lambda _{}^j\left(\stackrel{~}{}_t\right)\right)\tau _{U_{ju}}\left(t\right)+\hfill \\ \hfill +\underset{j=f+1}{\overset{p}{}}\underset{\begin{array}{c}\beta \alpha \\ =\delta _{ju}1q\end{array}}{}\left(D_{\lambda _j,\alpha }\left(\stackrel{~}{}_{t^{\left(j\right)}}\right)p_\beta \left(\stackrel{~}{}_{t^{\left(j\right)}}\right)\text{Ď}_\lambda _{}^j\left(\stackrel{~}{}_t\right)\right)\tau _{U_{ju}}\left(t\right)=0\end{array}$$
3. $$\begin{array}{c}\underset{j=1}{\overset{f}{}}\underset{\begin{array}{c}\beta _1+\beta _2+\beta _3\alpha _1\alpha _2\alpha _3\\ =q+\delta _{ju}+\delta _{jv}+\delta _{jw}\end{array}}{}(D_{\lambda _j,\alpha _1}(\stackrel{~}{}_{t^{\left(j\right)}})p_{\beta _1}\left(\stackrel{~}{}_{t^{\left(j\right)}}\right)\text{Ď}_\lambda _{}^j(\stackrel{~}{}_t)D_{\mu _j,\alpha _2}(\stackrel{~}{}_{s^{\left(j\right)}})p_{\beta _2}\left(\stackrel{~}{}_{s^{\left(j\right)}}\right)\text{Ď}_\mu _{}^j(\stackrel{~}{}_s)\hfill \\ \hfill D_{\nu _j,\alpha _3}\left(\stackrel{~}{}_{x^{\left(j\right)}}\right)p_{\beta _3}(\stackrel{~}{}_{x^{\left(j\right)}})\text{Ď}_\nu _{}^j(\stackrel{~}{}_x)\left)\tau _{U_{uj}}\right(t\left)\tau _{U_{vj}}\right(s\left)\tau _{U_{jw}}\right(x)+\\ \hfill +\underset{j=f+1}{\overset{p}{}}\underset{\begin{array}{c}\beta _1+\beta _2+\beta _3\alpha _1\alpha _2\alpha _3=\\ =q+\delta _{ju}+\delta _{jv}+\delta _{jw}1\end{array}}{}(D_{\lambda _j,\alpha _1}(\stackrel{~}{}_{t^{\left(j\right)}})p_{\beta _1}\left(\stackrel{~}{}_{t^{\left(j\right)}}\right)\text{Ď}_\lambda _{}^j(\stackrel{~}{}_t)D_{\mu _j,\alpha _2}(\stackrel{~}{}_{s^{\left(j\right)}})p_{\beta _2}\left(\stackrel{~}{}_{s^{\left(j\right)}}\right)\text{Ď}_\mu _{}^j(\stackrel{~}{}_s)\\ \hfill D_{\nu _j,\alpha _3}\left(\stackrel{~}{}_{x^{\left(j\right)}}\right)p_{\beta _3}(\stackrel{~}{}_{x^{\left(j\right)}})\text{Ď}_\nu _{}^j(\stackrel{~}{}_x)\left)\tau _{U_{uj}}\right(t\left)\tau _{U_{vj}}\right(s\left)\tau _{U_{jw}}\right(x)=0.\end{array}$$
## 4. $`\tau `$-functions and theta functions of Jacobians
### 4.A. Geometrical meaning of “formal”objects
This section aims at motivating the relation between $`\tau `$-functions and theta functions of Jacobians. In particular, $`\tau `$-functions attached to a Riemann surface with marked points will be defined following the works of Fay, Krichever, Shiota and Adler–Shiota–van Moerbecke (\[F, K, Sh, ASvM\]).
Throughout this section, $`C`$ will be an integral complete curve over $``$ of genus $`g`$. Let $`J_{g1}(C)`$ denote the scheme parametrizing invertible sheaves of degree $`g1`$. For the sake of clarity, we will assume $`C`$ to be smooth although most of the results established here hold in greater generality.
We let $`r=1`$ in the ramified case and $`r=p`$ in the non-ramified case.
Let us fix data $`(C,\overline{x},t_{\overline{x}})`$, where $`\overline{x}=\{x_1,\mathrm{},x_r\}`$ are $`r`$ pairwise distinct points of $`C`$ and $`t_{\overline{x}}`$ is a collection of formal parameters $`\{t_{x_1},\mathrm{},t_{x_r}\}`$ giving corresponding isomorphisms $`t_{x_j}:\widehat{𝒪}_{C,x_j}\stackrel{}{}[[z_j]]`$.
###### Proposition 4.1.
For each invertible sheaf $`LJ_{g1}(C)`$, there is a canonical morphism $`\gamma :𝒥(\widehat{C}_V)J_{g1}(C)`$ such that the diagram
is commutative. Here $`\varphi _V`$ is the Abel map and $`\varphi _L`$ sends a point $`x^{}C`$ to $`L(rx^{}\overline{x})`$.
###### Proof.
The data $`(C,\overline{x},t_{\overline{x}})`$ gives rise to a canonical morphism:
$$𝒪_C\widehat{𝒪}_{C,\overline{x}}\stackrel{}{}V^+.$$
Then we define $`\alpha `$ to be the induced morphism between the corresponding schemes.
Because of the Albanese property given in Proposition 2.3, the composition $`\widehat{C}_VJ_{g1}(C)`$ factors through the Abel map $`\varphi _V:\widehat{C}_V𝒥(\widehat{C}_V)`$, thus defining $`\gamma `$. ∎
Let $`J_{g1}^{\mathrm{}}(C,\overline{x})`$ be the scheme parametrizing pairs $`(L,\varphi )`$ where $`LJ_{g1}(C)`$ and $`\varphi :\widehat{L}_{\overline{x}}\stackrel{}{}\widehat{𝒪}_{C,\overline{x}}`$. It carries a canonical action of $`\overline{\mathrm{\Gamma }}_V^+`$ since this group acts by homotheties on the trivialization of $`L`$. Summing up, there is an exact sequence of group schemes
$$0\overline{\mathrm{\Gamma }}_V^+J_{g1}^{\mathrm{}}(C,\overline{x})J_{g1}(C)\mathrm{\hspace{0.17em}0}.$$
This sequence has also a formal counterpart
$$0\overline{\mathrm{\Gamma }}_V^+\mathrm{\Gamma }_V𝒥(\widehat{C}_V)\mathrm{\hspace{0.17em}0}.$$
Since this last sequence splits, $`\mathrm{\Gamma }_V𝒥(\widehat{C}_V)\times \overline{\mathrm{\Gamma }}_V^+`$ and every element $`g\mathrm{\Gamma }_V`$ can be written as $`(g_{},g_+)`$ with $`g_{}𝒥(\widehat{C}_V)`$, $`g_+\overline{\mathrm{\Gamma }}_V^+`$ and $`g=g_{}g_+`$.
###### Proposition 4.2.
Let $`(L,\varphi )`$ be a point in $`J_{g1}^{\mathrm{}}(C,\overline{x})`$. Then the Abel maps $`\varphi _V`$ and $`\varphi _L`$ have canonical lifts to $`\mathrm{\Gamma }_V`$ and $`J_{g1}^{\mathrm{}}(C,\overline{x})`$ respectively, and the map $`\gamma `$ has a lift
compatible with those of $`\varphi _V`$ and $`\varphi _L`$.
###### Proof.
Let $`E(u,v)`$ denote the prime form of $`C`$ as a meromorphic function on $`C\times C`$. Then the meromorphic function $`\frac{E(x,x_i)}{E(x,x^{})}`$ has a zero at $`x=x_i`$ and a pole at $`x=x^{}`$. Moreover, if $`z_i`$ is a coordinate at $`x_i`$ such that $`z_i(x_i)=0`$ and $`\overline{z}_i:=z_i(x^{})`$, then the expansion of this function at $`x_i`$ is the following series.
$$t_{x_i}\left(\frac{E(x,x_i)}{E(x,x^{})}\right)\left(1\frac{\overline{z}_i}{z_i}\right)^1(1+z_i[[z_i]]).$$
Therefore, the $`r`$-tuple consisting of the expansions
$$t_{\overline{x}}(\frac{E(x,x_1)}{E(x,x^{})},\mathrm{},\frac{E(x,x_r)}{E(x,x^{})})$$
corresponds to a morphism $`\widehat{C}_V\mathrm{\Gamma }_V`$ which lifts the Abel morphism $`\varphi _V`$ defined by the $`r`$-tuple (2.2).
To define the lift of $`\varphi _L`$, it is enough to observe that if the line bundle $`L`$ carries a formal trivialization $`\varphi :\widehat{L}_{\overline{x}}\stackrel{}{}\widehat{𝒪}_{C,\overline{x}}`$, then $`L(rx^{}\overline{x})`$ is canonically endowed with the trivialization given by
$$\varphi \underset{i=1}{\overset{r}{}}\frac{E(x,x_i)}{E(x,x^{})}.$$
Finally, the lift of $`\gamma `$ is defined by
$`\widehat{\gamma }:\mathrm{\Gamma }_V`$ $`J_{g1}^{\mathrm{}}(C,\overline{x})`$
$`g`$ $`(LL_g_{},\varphi g)`$
where $`L_g_{}`$ is given as follows: let $`D_i`$ be a small disk around $`x_i`$ such that $`z_i`$ defines a coordinate in $`D_i`$, then $`L_g_{}`$ consists of gluing the trivial bundles on $`C\overline{x}`$ and on $`\stackrel{}{D}_1,\mathrm{},\stackrel{}{D}_r`$ by the transition functions $`(g_{})_1,\mathrm{},(g_{})_r`$ (see \[SW\], Remark 6.8 and \[Sh\], Lemma 4). ∎
Recall that the Krichever map associated to $`(C,\overline{x},t_{\overline{x}})`$ is the morphism
$$\begin{array}{cc}\hfill J_{g1}^{\mathrm{}}(C,\overline{x})& \mathrm{Gr}(V)\hfill \\ \hfill (L,\varphi )& (t_{\overline{x}}\varphi )\left(H^0(C\overline{x},L)\right)\hfill \end{array}.$$
Then, we obtain the following
###### Theorem 4.3.
Let $`\mathrm{Kr}`$ denote the Krichever map associated to $`(C,\overline{x},t_{\overline{x}})`$. For $`(L,\varphi )J_{g1}^{\mathrm{}}(C,\overline{x})`$, let $`U=\mathrm{Kr}(L,\varphi )\mathrm{Gr}(V)`$.
Then the composition
$$\mathrm{\Gamma }_V\stackrel{\widehat{\gamma }}{}J_{g1}^{\mathrm{}}(C,\overline{x})\stackrel{\mathrm{Kr}}{}\mathrm{Gr}(V)$$
coincides with the morphism $`\mu _U:\mathrm{\Gamma }_V\mathrm{\Gamma }_V\times \{U\}\mathrm{Gr}(V)`$ mapping $`g`$ to $`gU`$.
###### Proof.
We have to check that $`(\mathrm{Kr}\widehat{\gamma })(g)=gU`$, but this is a straightforward consequence of the definition of $`\widehat{\gamma }`$. ∎
In particular, we have obtained morphisms
As a straightforward consequence of the determinantal construction of the theta divisor in $`J_{g1}(C)`$ and of the determinant bundle on $`\mathrm{Gr}(V)`$, it follows that
$$\mathrm{Kr}^{}\mathrm{Det}_V^{}\pi ^{}𝒪_{J_{g1}}(\mathrm{\Theta })$$
###### Remark 4.4.
It is easy to check that there is a canonical isomorphism of formal schemes
$$\mathrm{\Gamma }_V(U)/\overline{\mathrm{\Gamma }}_V^+J_{g1}(C)_L^\widehat{},$$
where $`J_{g1}(C)_L^\widehat{}`$ denotes the formal completion of $`J_{g1}(C)`$ at the point $`L`$. If $`A_U=\{vV:vUU\}`$ is the stabilizer of $`U`$, then the above isomorphism yields, at the level of tangent spaces,
$$T_1\left(\mathrm{\Gamma }_V(U)/\overline{\mathrm{\Gamma }}_V^+\right)=V/(A_U+V^+)H^1(C,𝒪_C)$$
since $`A_U=t_{\overline{x}}(H^0(C\overline{x},𝒪_C))`$. This is related to Mulase’s characterization of Jacobian varieties (\[M\]).
### 4.B. $`\tau `$-function attached to a Jacobian
The relations just described between the determinant bundle and the theta line bundle induce an explicit relation between $`\tau `$-functions and theta functions, which we now study.
Let $`J(C)`$ denote the Jacobian variety of $`C`$. Choose a symplectic basis $`\{\alpha _1,\mathrm{},\alpha _g,\beta _1,\mathrm{},\beta _g\}`$ for $`H_1(C,)`$ and let $`\{\omega _1,\mathrm{},\omega _g\}`$ denote the corresponding canonical basis of holomorphic $`1`$-forms on $`C`$; that is,
$$_{\alpha _i}\omega _j=\delta _{ij}\text{ and }_{\beta _i}\omega _j=\mathrm{\Omega }_{ij}$$
where $`\mathrm{\Omega }=\left(\mathrm{\Omega }_{ij}\right)`$ is the period matrix of $`C`$. Then, as a complex torus, we have
$$J(C)=^g/^g+\mathrm{\Omega }^g.$$
We also recall that the Riemann theta function associated to $`\mathrm{\Omega }`$ is the quasi-periodic function on $`^g`$ (the universal cover of $`J(C)`$) given by
$$\theta (z)=\underset{m^g}{}\mathrm{exp}\left(2\pi im^tz+\pi im^t\mathrm{\Omega }m\right).$$
By the Riemann-Kempf Theorem, there is an identification of $`J(C)`$ and $`J_{g1}(C)`$ such that the zero divisor of $`\theta (z)`$ corresponds to the divisor $`\mathrm{\Theta }J_{g1}(C)`$. Therefore, a point $`\xi ^g`$ gives rise to a invertible sheaf $`L_\xi J_{g1}(C)`$.
Given $`(C,\overline{x}=\{x_1,\mathrm{},x_r\},t_{\overline{x}})`$ there is a canonical map
$`H^0(C,\mathrm{\Omega }_C)`$ $`\mathrm{\Omega }_{\widehat{C}_V}[[z_1]]dz_1\times \mathrm{}\times [[z_r]]dz_r`$
$`\omega _i`$ $`(\mathrm{},(a_{i1}^{(j)}+a_{i2}^{(j)}z_j+a_{i3}^{(j)}z_j^2+\mathrm{})dz_j,\mathrm{})`$
which maps each $`\omega _i`$ to its local expansions at the points $`x_1,\mathrm{},x_r`$. Since we have chosen bases on both sides, for each $`j\{1,\mathrm{},r\}`$ we have a $`g\times \mathrm{}`$ matrix over $``$, $`A^{(j)}`$, associated to the map $`H^0(C,\mathrm{\Omega }_C)[[z_j]]dz_j`$. Moreover, $`\mathrm{rank}A^{(j)}=g`$ for each $`j\{1,\mathrm{},r\}`$.
###### Remark 4.5.
The transpose of the above map is related to the map induced by Proposition 4.1 at the level of tangent spaces.
We will now define the $`\tau `$-function of $`(C,x_1,\mathrm{},x_r,z_1,\mathrm{},z_r,\xi )`$, where $`\{x_1,\mathrm{},x_r\}`$ are $`r`$ different points in $`C`$, with respective local holomorphic coordinates $`z_1,\mathrm{},z_r`$, and $`\xi ^g`$ is a point in the universal cover of $`J(C)`$.
For each $`j\{1,\mathrm{},r\}`$ and each natural number $`n`$, we let $`\eta _n^{(j)}`$ denote the normalized meromorphic $`1`$-form on $`C`$ with a unique pole of order $`n+1`$ at $`x_j`$ of the form $`d(z_j^{(n+1)})+O(1)`$ and such that
$$^x\eta _n^{(j)}=z_j^n+O(z_j)\text{ at }x=x_j.$$
We also consider the complex numbers $`q_{nm}^{(j)}`$ defined by the following identities.
$$^x\eta _n^{(j)}=z_j^n2\underset{m=1}{\overset{\mathrm{}}{}}q_{nm}^{(j)}\frac{z_j^m}{m}\text{ at }x=x_j.$$
If we let $`t`$ be the $`r`$-tuple $`(t^{(1)},\mathrm{},t^{(r)})`$ where each $`t^{(j)}`$ is a family of variables $`(t_1^{(j)},t_2^{(j)},\mathrm{})`$, then we define the following quadratic form
$$Q(t)=\underset{n,m1}{}q_{nm}^{(1)}t_n^{(1)}t_m^{(1)}+\mathrm{}+\underset{n,m1}{}q_{nm}^{(r)}t_n^{(r)}t_m^{(r)}$$
and the point $`A(t)`$ of $`J(C)`$ with values in $`\{\{t^{(1)},\mathrm{},t^{(r)}\}\}`$ given by
$$A(t)=A^{(1)}t^{(1)}+\mathrm{}+A^{(r)}t^{(r)}.$$
Finally, the $`\tau `$-function of $`(C,x_1,\mathrm{},x_r,z_1,\mathrm{},z_r,\xi )`$ is defined as follows
(4.6)
$$\tau (\xi ,t):=\mathrm{exp}(Q(t))\theta (A(t)+\xi )$$
for $`t=(t^{(1)},\mathrm{},t^{(r)})`$.
It is worth pointing out that standard calculations shows that the $`\tau `$-function $`\tau (\xi ,t)`$ coincides with the $`\tau `$-function $`\tau _U(t)`$ of the point $`U`$ defined through the Krichever map (\[SW, ASvM\]).
Now, let us introduce the associated BA function for the ramified case ($`r=1`$) by
(4.7)
$$\psi _\xi (z,t):=\mathrm{exp}(\underset{i1}{}\frac{t_i}{z^i})\frac{\tau (\xi ,t+[z])}{\tau (\xi ,t)}$$
and the adjoint BA function by
(4.8)
$$\psi _\xi ^{}(z,t):=\mathrm{exp}(\underset{i1}{}\frac{t_i}{z^i})\frac{\tau (\xi ,t[z])}{\tau (\xi ,t)}$$
where $`t=(t_1,t_2,\mathrm{})`$ and $`t+[z]=(t_1+z,t_2+z^2/2,\mathrm{},t_n+z^n/n,\mathrm{})`$.
Then it follows from \[DJKM, F\] that this BA function coincides with the BA function of $`U`$ (the corresponding point under the Krichever map) and that the following equality holds
(4.9)
$$\mathrm{Res}_{z=0}\psi _\xi (z,t)\psi _\xi ^{}(z,s)\frac{dz}{z^2}=0$$
for all $`t`$ and $`s`$.
In the non-ramified case ($`r=p`$), BA functions will be introduced as follows. Let $`u,v`$ be two integers in $`\{1,\mathrm{},p\}`$ and let $`T_{uv}`$ denote the homothety of $`V`$ given by the element $`(1,\mathrm{},z_u,\mathrm{},z_v^1,\mathrm{},1)`$ if $`uv`$, and the identity if $`u=v`$. Then it is easy to check that $`T_{uv}`$ induces an automorphism of each connected component of $`\mathrm{Gr}(V)`$ and that these automorphisms lift canonically to automorphisms of the determinant bundle, which will also be denoted by $`T_{uv}`$. In particular, we have automorphisms of $`H^0(\mathrm{Gr}(V),\mathrm{Det}^{})`$.
Recalling the bosonization isomorphism and the fact that $`𝒪(\mathrm{\Gamma }_V^+)^{}𝒪(\mathrm{\Gamma }_V^{})`$, we obtain isomorphisms
$$T_{uv}^{}:𝒪(\mathrm{\Gamma }_V^{})\stackrel{}{}𝒪(\mathrm{\Gamma }_V^{})$$
and we define
$$\tau _{uv}(\xi ,t):=T_{uv}^{}(\tau (\xi ,t))$$
where $`\tau (\xi ,t)`$ is given by (4.6). Note that $`\tau _{U_{uv}}`$ coincides with $`T_{uv}^{}(\tau _U)`$ for any $`U\mathrm{Gr}(V)`$.
We will also consider the corresponding formal $`u`$BA functions as follows
(4.10)
$$\psi _{u,\xi }(z_{},t):=(\psi _{u,\xi }^{(1)}(z_1,t),\mathrm{},\psi _{u,\xi }^{(r)}(z_r,t))$$
where
(4.11)
$$\psi _{u,\xi }^{(j)}(z_j,t):=\mathrm{exp}\left(\underset{i1}{}\frac{t_i^{(j)}}{z_j^i}\right)\frac{\tau _{uj}(\xi ,t+[z_j])}{\tau (\xi ,t)},$$
and the formal adjoint $`u`$BA functions by
(4.12)
$$\psi _{u,\xi }^{}(z_{},t):=(\psi _{u,\xi }^{(1)}(z_1,t)\mathrm{},\psi _{u,\xi }^{(r)}(z_r,t))$$
where
(4.13)
$$\psi _{u,\xi }^{(j)}(z_j,t):=\mathrm{exp}\left(\underset{i1}{}\frac{t_i^{(j)}}{z_j^i}\right)\frac{\tau _{ju}(\xi ,t[z_j])}{\tau (\xi ,t)}.$$
Then, from \[MP2\], we obtain that the BA functions satisfy the $`(1,\stackrel{𝑝}{\mathrm{}},1)`$-KP-hierarchy
$$\mathrm{Res}_{z=0}\mathrm{Tr}\left(\frac{\psi _{u,U}(z_{},t)}{(1,\mathrm{},z_u,\mathrm{},1)}\frac{\psi _{v,U}^{}(z_{},s)}{(1,\mathrm{},z_v,\mathrm{},1)}\right)dz=\mathrm{\hspace{0.17em}0}$$
for all $`u,v\{1,\mathrm{},p\}`$. Equivalently, in terms of $`\tau `$-functions we have
(4.14)
$$\begin{array}{c}\underset{j=1}{\overset{p}{}}\underset{\begin{array}{c}\beta _1+\beta _2\alpha _1\alpha _2\\ =\delta _{ju}+\delta _{jv}1\end{array}}{}p_{\alpha _1}(t^{(j})p_{\beta _1}(\stackrel{~}{}_{t^{(j)}})\tau _{U_{u,j}}(t)\hfill \\ \hfill p_{\alpha _2}(s^{(j)})p_{\beta _2}(\stackrel{~}{}_{s^{(j)}})\tau _{U_{jv}}(s)=0.\end{array}$$
Moreover, it follows from \[MP2\] that these hierarchies characterize when the BA-functions (respectively $`\tau `$-function) are the BA-functions (respectively $`\tau `$-function) of a point of $`\mathrm{Gr}(V)`$.
## 5. Geometric characterization of Jacobian varieties with automorphisms in terms of Sato Grassmannian
In this section we characterize the points of the Sato Grassmannian that arise from geometric data over an algebraic curve with automorphisms via the Krichever construction. We prove that these points are those whose orbit under the action of $`\mathrm{\Gamma }_V`$ is finite dimensional (up to the action of $`\overline{\mathrm{\Gamma }}_V^+`$). A related result has been established in \[GMP\]. This type of characterizations dates back to the approach of Mulase (\[M\]).
Let us denote by $`\overline{\mathrm{Pic}}(C)`$ the moduli space of rank $`1`$ torsion free sheaves on $`C`$.
We say that $`(C,\overline{x},L)`$ is maximal for a curve $`C`$, a divisor $`\overline{x}`$ composed by $`r`$ pairwise distinct points $`x_1,\mathrm{},x_r`$ in $`C`$, and $`L\overline{\mathrm{Pic}}(C)`$ when one of the following equivalent condition holds (see page 38 of \[SW\])
* let $`(C^{},\overline{x}^{},L^{})`$ be another such triple, and suppose there exists $`\psi :C^{}C`$ a birational morphism such that $`\psi (\overline{x}^{})=\overline{x}`$ and $`\psi _{}L^{}\stackrel{}{}L`$; then $`\psi `$ is an isomorphism.
* The canonical map $`𝒪_C\mathrm{End}(L)`$ is an isomorphism.
We begin with a detailed study of the ramified case and then generalize the results to the non-ramified case.
### 5.A. Ramified case
For the sake of notation, in this subsection the subscript $`\mathrm{R}`$ in $`V_\mathrm{R}`$, $`V_\mathrm{R}^+`$, $`\mathrm{\Gamma }_{V_\mathrm{R}}`$, $`\mathrm{}`$ will be omitted.
###### Definition 5.1.
Let $`\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{R})`$ be the contravariant functor from the category of $``$-schemes to the category of sets defined by
$$S\{(C,\sigma _C,x,t_x,L,\varphi _x)\}$$
where
1. $`p_C:CS`$ is a proper and flat morphism whose fibres are geometrically integral curves.
2. $`\sigma _C:CC`$ is an order $`p`$ automorphism (over $`S`$).
3. $`x:SC`$ is a smooth section of $`p_C`$, fixed under $`\sigma _C`$, such that the Cartier divisor $`x(s)`$ is a smooth point of $`C_s:=p_C^1(s)`$ for all closed points $`sS`$.
4. $`t_x`$ is an equivariant formal parameter along $`x(S)`$; that is, an equivariant isomorphism of $`𝒪_S`$-modules $`t_x:\widehat{𝒪}_{C,x(S)}\stackrel{}{}\widehat{V}_S^+`$.
5. $`L\overline{\mathrm{Pic}}(C)`$ satisfies that $`(C_s,x(s),L|_{C_s})`$ is maximal for all closed point $`sS`$.
6. $`\varphi _x`$ is a formal trivialization of $`L`$ along $`x(S)`$; that is, an isomorphism $`\varphi _x:\widehat{L}_{x(S)}\widehat{𝒪}_{C,x(S)}`$.
7. $`(C,\sigma _C,x,t_x,L,\varphi _x)`$ and $`(C^{},\sigma _C^{},x^{},t_x^{},L^{},\varphi _x^{})`$ are said to be equivalent when there is an isomorphism of $`S`$-schemes $`C\stackrel{}{}C^{}`$ compatible with all the data.
The Krichever morphism for the functor $`\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{R})`$ is the morphism of functors
$$\mathrm{Kr}:\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{R})\underset{¯}{\mathrm{Gr}(V)}$$
which sends the $`S`$-valued point $`(C,\sigma _C,x,t_x,L,\varphi _x)`$ to the following submodule of $`\widehat{V}_S:=V\widehat{}𝒪_S`$
$$(t_x\varphi _x)\left(\underset{𝑚}{\underset{}{\mathrm{lim}}}(p_C)_{}L(mx)\right)\widehat{V}_S.$$
###### Theorem 5.2.
The functor $`\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{R})`$ is representable by a subscheme $`\overline{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{R})`$ of $`\mathrm{Gr}(V)`$.
###### Proof.
Consider the morphism from $`\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{R})`$ to $`\mathrm{Gr}(V)\times \mathrm{Gr}(V)`$ which sends the $`S`$-valued point $`(C,\sigma _C,x,t_x,L,\varphi _x)`$ to the following pair of submodules
$$(t_x\left(\underset{𝑚}{\underset{}{\mathrm{lim}}}\left(p_C\right)_{}𝒪_C\left(mx\right)\right),\left(t_x\varphi _x\right)\left(\underset{𝑚}{\underset{}{\mathrm{lim}}}\left(p_C\right)_{}L\left(mx\right)\right))\mathrm{Gr}\left(V\right)\times \mathrm{Gr}\left(V\right)$$
where $`p_C:C\times SC`$ is the projection.
From the inverse construction of the Krichever map (\[K, SW\]) one has that this map is injective and that the image is contained in the set $`Z`$ of those pairs $`(𝒜,)`$ in $`\mathrm{Gr}(V)\times \mathrm{Gr}(V)`$ such that
$$𝒪_S𝒜,𝒜𝒜𝒜,𝒜,\sigma (𝒜)=𝒜.$$
Let us examine the maximality condition. For $`(𝒜,)`$ satisfying the above conditions, let $`A_{}`$ denote the stabilizer of $``$
$$A_{}:=\{v\widehat{V}_S\text{ such that }v\},$$
and let $`(C^{},\sigma _C^{},x^{},t_x^{},L^{},\varphi _x^{})`$ be the geometric data defined by the pair $`(A_{},)`$. Then the inclusion $`𝒜A_{}`$ gives rise to an equivariant morphism of $`S`$-schemes $`\psi :C^{}C`$ such that $`\psi (x^{})=x`$ and $`\psi _{}L^{}L`$. The maximality condition says that $`\psi _s`$ is an isomorphism for every closed point $`sS`$. That is, $`A_{}`$ is a finite $`𝒜`$-module such that $`𝒜_s=(A_{})_s`$ for all $`s`$. Therefore we have that $`𝒜=A_{}`$. Summing up, we are interested on the subset $`Z_0`$ of $`Z`$ consisting of those pairs $`(𝒜,)`$ such that $`𝒜=𝒜_{}`$.
From the proof of Theorem 6.5 of \[MP\] we know that the condition $`𝒜_{}𝒜`$, and hence $`Z_0`$, is closed. The Krichever construction implies that $`Z_0`$ represents $`\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{R})`$. Finally, $`p_2|_{Z_0}:Z_0\mathrm{Gr}(V)`$ is a closed immersion (where $`p_2`$ is the projection onto the second factor), and the theorem is proved. ∎
Let us consider the following action of $`\mathrm{\Gamma }_V`$ on $`\mathrm{Gr}(V)^p:=\mathrm{Gr}(V)\times \stackrel{𝑝}{\mathrm{}}\times \mathrm{Gr}(V)`$
$`\mu ^p:\mathrm{\Gamma }_V\times \mathrm{Gr}(V)^p`$ $`\mathrm{Gr}(V)^p`$
$`(g,(U_1,\mathrm{},U_p))`$ $`(gU_1,\mathrm{},gU_p).`$
Then the morphism
$`\mathrm{Gr}(V)`$ $`\mathrm{Gr}(V)^p`$
$`U`$ $`U_\sigma :=(U,\sigma (U),\mathrm{},\sigma ^{p1}(U))`$
is a $`\mathrm{\Gamma }_V`$-equivariant closed immersion.
The orbit of $`U_\sigma \mathrm{Gr}(V)^p`$ under the action of $`\mathrm{\Gamma }_V`$ is the schematic image of $`\mu _{U_\sigma }^p:=\mu ^p|_{\mathrm{\Gamma }_V\times \{U_\sigma \}}`$; it will be denoted by $`\mathrm{\Gamma }_V(U_\sigma )`$. Section 4 of \[GMP\] implies that $`\mathrm{\Gamma }_V(U_\sigma )/\overline{\mathrm{\Gamma }}_V^+`$ is a formal scheme whose tangent space is
$$T_{U_\sigma }\left(\mathrm{\Gamma }_V(U_\sigma )/\overline{\mathrm{\Gamma }}_V^+\right)T_1\mathrm{\Gamma }_V/(\mathrm{Ker}d\mu _{U_\sigma }^p+T_1\overline{\mathrm{\Gamma }}_V^+),$$
where
$$d\mu _{U_\sigma }^p:T_1\mathrm{\Gamma }_VT_{U_\sigma }\mathrm{Gr}(V)^p$$
is the map induced by $`\mu _{U_\sigma }^p`$ on the respective tangent spaces.
###### Theorem 5.3 (ramified).
Let $`U`$ be a closed point of $`\mathrm{Gr}(V)`$. Then the following conditions are equivalent.
1. $`dim_{}T_{U_\sigma }(\mathrm{\Gamma }_V(U_\sigma )/\overline{\mathrm{\Gamma }}_V^+)<\mathrm{}`$ where $`V=V_\mathrm{R}`$, and
2. there exists $`(C,\sigma _C,x,t_x,L,\varphi _x)\overline{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{R})`$ such that its image by the Krichever morphism is $`U`$.
###### Proof.
The result follows straightforwardly from similar arguments to those of Theorems 4.13 and 4.15 of \[GMP\]. ∎
###### Remark 5.4.
If the conditions of the theorem hold, it follows that $`\mathrm{Ker}d\mu _{U_\sigma }^p=_{i=0}^{p1}\sigma ^iA_U`$, with $`A_U`$ the stabilizer of $`U`$. It is then straightforward to show that a similar characterization exists in terms of two copies of $`\mathrm{Gr}(V)`$ instead of $`p`$ copies.
###### Theorem 5.5.
Let $`\underset{¯}{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{R})`$ denote the subfunctor of $`\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{R})`$ consisting of those data $`(C,\sigma _C,x,t_x,L,\varphi _x)`$ such that the fibres $`C_s`$ are smooth curves for all closed points $`sS`$.
Then $`\underset{¯}{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{R})`$ is representable by an open subscheme $`\mathrm{Pic}^{\mathrm{}}(p,\mathrm{R})`$ of $`\overline{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{R})`$.
###### Proof.
Let $`𝒞\overline{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{R})`$ be the universal curve. Then $`\mathrm{Pic}^{\mathrm{}}(p,\mathrm{R})`$ is given by the open subscheme of $`\overline{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{R})`$ consisting of the points $`s`$ such that $`𝒞_s`$ is a smooth curve.
###### Remark 5.6.
Observe that we have a forgetful map
which sends $`(C,\sigma _C,x,t_x,L,\varphi _x)`$ to $`(C,\sigma _C,x,t_x)`$. The fibre of $`(C,\sigma _C,x,t_x)`$ when $`C`$ is a smooth curve is essentially the scheme studied in Section 4.A.
### 5.B. Non-ramified case
Similarly to the ramified case we now consider the functor $`\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{NR})`$. To define it, it suffices to replace conditions ($`1`$) and ($`3`$) in Definition 5.1 respectively by
1. $`p_C:CS`$ is a proper and flat morphism whose fibres are geometrically reduced curves.
2. $`x_i:SC`$ are disjoint smooth sections of $`p_S`$ ($`i=1,\mathrm{},p`$) such that $`\sigma _C(x_i)=x_{i+1}`$ for $`i<p`$ and $`\sigma _C(x_p)=x_1`$. Also, for every closed point $`sS`$ and each irreducible component of $`C_s`$, there is at least one $`i`$ such that $`x_i(s)`$ lies on that component.
We also take $`V`$, $`V^+`$, $`\mathrm{\Gamma }_V`$, $`\mathrm{}`$ to be $`V_{\mathrm{NR}}`$, $`V_{\mathrm{NR}}^+`$, $`\mathrm{\Gamma }_{V_{\mathrm{NR}}}`$, $`\mathrm{}`$ in this case.
With arguments similar to those in the previous section, we can prove the following results.
###### Theorem 5.7.
The functor $`\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{NR})`$ is representable by a subscheme $`\overline{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{NR})`$ of $`\mathrm{Gr}(V)`$.
The subfunctor $`\underset{¯}{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{N}R)`$ of $`\overline{\underset{¯}{\mathrm{Pic}}}^{\mathrm{}}(p,\mathrm{NR})`$ consisting of those data $`(C,\sigma _C,\overline{x},t_x,L,\varphi _x)`$ such that the fibres $`C_s`$ are smooth curves for all closed points $`sS`$ is representable by an open subscheme $`\mathrm{Pic}^{\mathrm{}}(p,\mathrm{NR})`$ of $`\overline{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{NR})`$.
###### Theorem 5.8 (non-ramified).
Let $`U`$ be a closed point of $`\mathrm{Gr}(V)`$. Then the following conditions are equivalent.
1. $`dim_{}T_{U_\sigma }(\mathrm{\Gamma }_V(U_\sigma )/\overline{\mathrm{\Gamma }}_V^+)<\mathrm{}`$ where $`V=V_{\mathrm{NR}}`$,
2. there exists $`(C,\sigma _C,\overline{x},t_x,L,\varphi _x)\overline{\mathrm{Pic}}^{\mathrm{}}(p,\mathrm{NR})`$ such that its image by the Krichever morphism is $`U`$.
###### Remark 5.9.
Analogously to the ramified case, it follows that there is a similar statement in terms of two copies of $`\mathrm{Gr}(V)`$ instead of $`p`$ copies.
## 6. Characterization of Jacobian theta functions of Riemann surfaces with non-trivial automorphisms
In this section we give the conditions that a theta function of a p.p.a.v. should satisfy in order to be the theta function of the Jacobian of a smooth irreducible projective curve.
We begin with the proof of a generalization of the known theorems of Mulase (\[M\]) and Shiota (\[Sh\]) in terms of the Sato Grassmannian.
We will use the following notation.
Let $`\mathrm{\Omega }^g`$ be a point in the Siegel upper half space such that the principally polarized abelian variety $`X_\mathrm{\Omega }:=^g/(^g+\mathrm{\Omega }^g)`$ is irreducible. Let $`\theta (z)=\theta (z,\mathrm{\Omega })`$ denote the Riemann theta function of $`X_\mathrm{\Omega }`$.
Let $`r`$ be a natural number, let $`A^{(j)}=(a_1^{(j)},a_2^{(j)},\mathrm{})(^g)^{\mathrm{}}`$ be a $`g\times \mathrm{}`$-matrix of rank $`g`$ for each $`j\{1,\mathrm{},r\}`$, let $`Q^{(j)}(t^{(j)})=_{i,k=1}^{\mathrm{}}q_{ik}^{(j)}t_i^{(j)}t_k^{(j)}`$, with $`q_{ik}^{(j)}`$, be a quadratic form for each $`j`$ in $`\{1,\mathrm{},r\}`$, and let $`\xi `$ be in $`^g`$.
###### Definition 6.1.
The $`\tau `$-function and BA-functions associated to the data $`(X_\mathrm{\Omega },\mathrm{\Omega },\{A^{(1)},\mathrm{},A^{(r)}\},\{Q^{(1)},\mathrm{},Q^{(r)}\},\xi )`$ are given formally by formulae (4.6) and (4.10)-(4.13).
The $`\tau `$-function will be denoted by $`\tau (\xi ,t)`$ while the BA-function (respectively adjoint BA-function) will be denoted by $`\psi _{u,\xi }(z_{},t)`$ (respectively $`\psi _{u,\xi }^{}(z_{},t)`$).
###### Theorem 6.2.
Let $`X_\mathrm{\Omega }`$ be an irreducible p.p.a.v. of dimension $`g`$.
Then the following conditions are equivalent.
1. There exists a triple $`(C,\overline{x},t_{\overline{x}})`$, where $`C`$ is a projective irreducible smooth curve of genus $`g`$, $`\overline{x}=(x_1,\mathrm{},x_r)`$ is an $`r`$-tuple of distinct points in $`C`$ and $`t_{\overline{x}}=(t_{x_1},\mathrm{},t_{x_r})`$ is an $`r`$-tuple of local parameters at the corresponding $`x_j`$, such that $`X_\mathrm{\Omega }`$ is isomorphic as a p.p.a.v. to the Jacobian of $`C`$.
2. For each $`j\{1,\mathrm{},r\}`$ there exist a $`g\times \mathrm{}`$-matrix of rank $`g`$, $`A^{(j)}=(a_1^{(j)},a_2^{(j)},\mathrm{})`$, with $`a_i^{(j)}^g`$, and a quadratic form $`Q^{(j)}(t^{(j)})=_{i,k=1}^{\mathrm{}}q_{ik}^{(j)}t_i^{(j)}t_k^{(j)}`$, with $`q_{ik}^{(j)}`$, such that for every $`\xi ^g`$, the corresponding $`\tau `$-function $`\tau (\xi ,t)`$ is a $`\tau `$-function of the $`(1,\stackrel{r}{\mathrm{}},1)`$KP-hierarchy (4.14).
Moreover, if one of the conditions is fulfilled, the matrices $`A^{(j)}`$ and the quadratic forms $`Q^{(j)}`$ are equal to the data associated to the triple $`(C,\overline{x},t_{\overline{x}})`$ in Section 4.B.
###### Proof.
$`1.2.`$ It follows from Section 4.
$`2.1.`$ We denote
$$A(t)=\underset{j=1}{\overset{r}{}}A^{(j)}t^{(j)}\text{ and }Q(t)=\underset{j=1}{\overset{r}{}}\underset{i,k1}{}q_{ik}^{(j)}t_i^{(j)}t_k^{(j)}.$$
Since $`\tau (\xi ,t)`$ is a $`\tau `$-function of the $`(1,\stackrel{r}{\mathrm{}},1)`$KP-hierarchy for every $`\xi ^g`$, it follows that $`\tau (\xi ,t)`$ defines a point $`U_\xi \mathrm{Gr}(V)`$ (with $`V=((z))\times \stackrel{r}{\mathrm{}}\times ((z))`$) such that $`\tau (\xi ,t)=\tau _{U_\xi }(t)`$ (up to a constant). From Theorem 3.12 in \[MP2\], we have that
$$U_\xi =(p_i\left(\stackrel{~}{_t}\right)\psi _{u,\xi }^{\left(1\right)}(z_1,t)_{|_{t=0}},\mathrm{},p_i\left(\stackrel{~}{_t}\right)\psi _{u,\xi }^{\left(r\right)}(z_r,t)_{|_{t=0}}),i0,1ur.$$
Therefore, we have obtained a morphism
$`\phi :^g`$ $`\mathrm{Gr}(V)`$
$`\xi `$ $`U_\xi .`$
We claim that this morphism induces an injection
(6.3)
$$X\mathrm{Gr}(V).$$
Indeed, given $`\xi _1`$ and $`\xi _2`$ in $`^g`$, the condition $`U_{\xi _1}=U_{\xi _2}`$ is equivalent to $`\tau (\xi _1,t)=\tau (\xi _2,t)`$ for all $`t`$ (up to a constant), which is in turn equivalent to $`\theta (A(t)+\xi _1)=\theta (A(t)+\xi _2)`$ for all $`t`$ (up to a constant), and therefore equivalent to $`\xi _1\xi _2^g+\mathrm{\Omega }^g`$, since $`\mathrm{\Theta }`$ is a principal polarization on $`X`$.
Now the function $`A`$ can be interpreted as a surjective linear map
$$^{\mathrm{}}\times \stackrel{r}{\mathrm{}}\times ^{\mathrm{}}\stackrel{A}{}^g$$
and, with the identifications $`^{\mathrm{}}\times \stackrel{r}{\mathrm{}}\times ^{\mathrm{}}=T_1\mathrm{\Gamma }_V^{}`$ and $`T_\xi X=^g`$, $`A`$ corresponds to a surjective morphism of formal group schemes
$$\mathrm{\Gamma }_V^{}\stackrel{A_\xi }{}\widehat{X}_\xi .$$
We claim now that the surjective morphism
$`\mu _\xi :\mathrm{\Gamma }_V^{}`$ $`\mathrm{\Gamma }(U_\xi )/\overline{\mathrm{\Gamma }}_V^+`$
$`g`$ $`gU_\xi `$
factorizes by $`A_\xi :\mathrm{\Gamma }_V^{}\widehat{X}_\xi `$. Observe that if $`s=(s^{(1)},\mathrm{},s^{(r)})\mathrm{Ker}A`$ then
$$\tau _{U_\xi }\left(t+s\right)=\tau (\xi ,t+s)=q_\xi (t,s)\mathrm{exp}\left(Q\left(t\right)\right)\theta \left(A\left(t\right)+\xi \right)=q_\xi (t,s)\tau (\xi ,t)$$
where $`q_\xi (t,s)`$ is an exponential of a linear function in $`t`$. Generalizing Lemma 3.8 of \[SW\] to the case of $`\mathrm{\Gamma }_V`$, there exists $`g\overline{\mathrm{\Gamma }}_V^+`$ (which depends on $`s`$) such that
$$\tau _{U_\xi }(t+s)=q_\xi (t,s)\tau (\xi ,t)=\tau _{gU_\xi }(t).$$
Hence there is a factorization
In particular, it follows that $`dimT_{U_\xi }\mathrm{\Gamma }(U_\xi )/\overline{\mathrm{\Gamma }}_V^+`$ is finite, and applying the results of \[M\] one has that there exists data $`(C_\xi ,\overline{x}_\xi ,t_\xi ,L_\xi ,\varphi _\xi )`$ associated to $`U_\xi `$ under the Krichever map.
Let us check that the piece of data $`(C_\xi ,\overline{x}_\xi ,t_\xi )`$ does not depend on $`\xi `$. Indeed, for $`\xi ,\xi ^{}X`$ let $`s^{\mathrm{}}\times \stackrel{r}{\mathrm{}}\times ^{\mathrm{}}`$ be such that $`A(s)=\xi ^{}\xi `$. Then
$`\tau _\xi (t+s)`$ $`=\mathrm{exp}(Q(t+s))\theta (A(t+s)+\xi )=`$
$`=q_\xi ^{}(t)\mathrm{exp}(Q(t))\theta (A(t)+\xi ^{})=q_\xi ^{}(t)\tau _\xi ^{}(t)`$
The generalization of Lemma 3.8 of \[SW\] implies that there exists $`g_\xi ^{}(z)V`$ such that $`U_\xi ^{}=g_\xi ^{}(z)U_\xi `$. From this fact it follows that $`U_\xi ^{}`$ and $`U_\xi `$ have the same stabilizer and that, therefore, $`(C_\xi ,\overline{x}_\xi ,t_\xi )`$ does not depend on $`\xi `$. It will be denoted by $`(C,\overline{x},t_{\overline{x}})`$.
The latter fact does have further consequences. It implies that the map (6.3) takes values into $`\overline{\mathrm{Pic}}^{\mathrm{}}(C,\overline{x})`$ (the subscheme of $`\mathrm{Gr}(V)`$ parameterizing torsion free sheaves of rank $`1`$ on $`C`$ with a formal trivialization along $`\overline{x}`$). Furthermore, it says that the composite map
takes values in $`\mathrm{Pic}^0(C)L_\xi `$, the orbit of $`L_\xi \overline{\mathrm{Pic}}(C)`$ under the action of $`\mathrm{Pic}^0(C)`$. Using the surjectivity of $`A`$ we can show that the induced map
$$X\mathrm{Pic}^0(C)L_\xi $$
is surjective. Since $`(C,\overline{x},t_{\overline{x}},L_\xi ,\varphi _\xi )`$ is maximal (see §5), the action of $`\mathrm{Pic}^0(C)`$ on $`\overline{\mathrm{Pic}}(C)`$ is free. So $`\mathrm{Pic}^0(C)`$ is a quotient of an abelian variety and, therefore, $`C`$ is a smooth complete curve of genus at most $`g`$.
To finish the implication $`2.1.`$, one has only to show that $`X\mathrm{Pic}^0(C)`$ is an isomorphism of p.p.a.v.’s: Given $`(X,\xi )`$ and $`(J(C),\xi )`$ we consider the tau-functions $`\tau _X=\tau (\xi ,t)`$ associated to $`X`$ as in (6.1) and $`\tau _J=\tau (\xi ,t)`$ associated to $`J(C)`$ as in (4.6). By the construction of the data $`(C_\xi ,\overline{x}_\xi ,t_\xi ,L_\xi ,\varphi _\xi )`$, it follows that $`\tau _X=\tau _J`$ (up to a constant) and hence
$$\theta _X(A(t)+\xi )=\mathrm{exp}(q(t))\theta _J(A_J(t)+\xi )\text{ (up to a constant)},$$
where $`q(t)`$ is a quadratic function.
Therefore
$$\mathrm{\Theta }_X=\mathrm{\Theta }_J\text{ (up to translation)}.$$
In particular, the curve $`C`$ is irreducible and of genus $`g`$. ∎
###### Remark 6.4.
Considering $`r=1`$ in the previous theorem we obtain the characterization of Jacobian varieties given by Shiota (\[Sh\], Theorem 6).
We will now apply this result to give a sufficient and necessary condition for a theta function of a p.p.a.v. to be the theta function of a curve with an automorphism of prime order $`p`$ with a fixed point.
###### Theorem 6.5 (ramified case).
Let $`X_\mathrm{\Omega }`$ be an irreducible p.p.a.v. of dimension $`g`$.
Then the following conditions are equivalent.
1. There exists a quadruple $`(C,\sigma _C,x,t_x)`$, where $`C`$ is a projective irreducible smooth curve of genus $`g`$, $`\sigma _C`$ is an automorphism of order $`p`$ of $`C`$, $`x`$ is a fixed point of $`\sigma _C`$ in $`C`$, and $`t_x`$ is a local parameter at $`x`$, such that $`X_\mathrm{\Omega }`$ is isomorphic as a p.p.a.v. to the Jacobian of $`C`$.
2. There exist a $`g\times \mathrm{}`$-matrix $`A`$ of rank $`g`$ and a symmetric quadratic form $`Q(t)`$ such that for each $`\xi _0`$ in $`^g`$ there exists $`\xi _1`$ in $`^g`$ so that the corresponding BA-functions satisfy
1. the KP hierarchy
$$\mathrm{Res}_{z=0}\psi _{\xi _0}(z,t)\psi _{\xi _0}^{}(z,s)\frac{dz}{z^2}=\mathrm{\hspace{0.17em}0},\text{and}$$
2. the identity
$$\mathrm{Res}_{z=0}\psi _{\xi _0}(\omega ^1z,\omega ^1t)\psi _{\xi _1}^{}(z,s)\frac{dz}{z^2}=\mathrm{\hspace{0.17em}0}$$
for all $`t`$ and $`s`$, where $`\omega `$ is a primitive $`p`$-th root of 1.
###### Proof.
$`1.2.`$
From the results given in Section 4.B with $`r=1`$, we know that given $`(C,\sigma _C,x,t_x)`$ as in condition 1., there exist $`A`$ and $`Q`$ as in condition 2) such that its associated BA-functions $`\psi _{\xi _0}(z,t)`$ and $`\psi _{\xi _0}^{}(z,s)`$, defined by (4.7) and (4.8) respectively, satisfy (4.9) for each $`\xi _0`$ in $`^g`$. That is, there exists a point $`U_{\xi _0}\mathrm{Gr}(V_\mathrm{R})`$ such that $`\psi _{\xi _0}=\psi _{U_{\xi _0}}`$ and 2.a) follows.
Furthermore, the induced embedding $`J(C)\mathrm{Gr}(V_\mathrm{R})`$ (given by (6.3)) is compatible with the actions of $`\sigma _C`$ in $`J(C)`$ and of $`\sigma `$ on $`\mathrm{Gr}(V_\mathrm{R})`$; that is, $`\sigma (U_\xi )=U_{\sigma _C^{}\xi }`$. Having this in mind and taking into account the relation between the BA-function of $`U_{\xi _0}`$ and $`U_{\xi _1}=\sigma (U_{\xi _0})`$ given by (3.2), 2.b) follows.
$`2.1.`$
We know by Theorem 6.2 (with $`r=1`$) that the first identity implies that there exists a triple $`(C,x,t_x)`$ such that $`X_\mathrm{\Omega }J(C)`$ as p.p.a.v.’s.
The second identity implies that $`U_{\xi _1}=\sigma (U_{\xi _0})`$. In particular, we have that $`A_{\xi _1}=\sigma (A_{\xi _0})`$ where $`A_\xi `$ denotes the stabilizer of $`U_\xi `$. Then the orbit $`\mathrm{\Gamma }_{V_\mathrm{R}}(U_{\xi _0},\sigma (U_{\xi _0})/\overline{\mathrm{\Gamma }}_{V_\mathrm{R}}^+`$ is finite dimensional and, by Theorem 5.3 (see Remark 5.4), we obtain that $`\sigma `$ induces an automorphism of $`A_{\xi _0}`$. Since $`A_{\xi _0}=t_{\overline{x}}(H^0(C\overline{x},𝒪_C))`$, because of the Krichever construction, the result follows. ∎
###### Remark 6.6.
Under the hypotheses of the above Theorem, suppose there are $`\xi _0`$ and $`\xi _1`$ in $`^g`$ such that $`\xi _0\xi _1^g+\mathrm{\Omega }^g`$ and the equations of the theorem are satisfied. Then there exists a line bundle $`L`$ on $`C`$ such that $`\sigma ^{}(L)L`$.
###### Theorem 6.7 (non-ramified).
Let $`X_\mathrm{\Omega }`$ be an irreducible p.p.a.v. of dimension $`g`$.
Then the following conditions are equivalent.
1. There exists a quadruple $`(C,\sigma _C,\overline{x},t_{\overline{x}})`$, where $`C`$ is a projective irreducible smooth curve of genus $`g`$, $`\sigma _C`$ is an automorphism of order $`p`$ of $`C`$, $`\overline{x}=\{x_1,\mathrm{},x_p\}`$ is an orbit of $`\sigma _C`$ consisting of $`p`$ different points in $`C`$, and $`t_{\overline{x}}=\{t_{x_1},\mathrm{},t_{x_p}\}`$ is a collection of local parameters $`t_{x_j}`$ at each respective $`x_j`$, such that $`X_\mathrm{\Omega }`$ is isomorphic as a p.p.a.v. to the Jacobian of $`C`$.
2. There exist $`p`$ matrices $`A^{(1)},\mathrm{},A^{(p)}`$, where $`A^{(j)}`$ is a $`g\times \mathrm{}`$-matrix of rank $`g`$, and $`p`$ symmetric quadratic forms $`Q^{(1)},\mathrm{},Q^{(p)}`$, such that for each $`\xi _0^g`$ there exists $`\xi _1^g`$ so that their BA-functions satisfy
1. the $`(1,\stackrel{r}{\mathrm{}},1)`$-KP hierarchy
$$\mathrm{Res}\left(\underset{j=1}{\overset{p}{}}z^{\delta _{ju}\delta _{jv}}\psi _{u,\xi _0}^{(j)}(z,t)\psi _{v,\xi _0}^{(j)}(z,s)\right)dz=\mathrm{\hspace{0.17em}0},\text{and}$$
2. the identity
$$\mathrm{Res}\left(\underset{j=1}{\overset{p}{}}z^{\delta _{ju}\delta _{jv}}\psi _{v+1,\xi _0}^{(j+1)}(z,\sigma ^{}(t))\psi _{u,\xi _1}^{(j)}(z,s)\right)dz=\mathrm{\hspace{0.17em}0}$$
where $`\sigma ^{}(t):=(t^{(p)},t^{(1)},t^{(2)},\mathrm{},t^{(p1)})`$.
###### Proof.
$`1.2.`$
Given $`(C,\sigma _C,\overline{x},t_{\overline{x}})`$ satisfying condition 1., we construct $`A(t)`$ and $`Q(t)`$ as in Section 4.B. By Theorem 6.2, for each $`\xi ^g`$ the corresponding BA-functions satisfy 2.a).
Similarly to the ramified case, the actions of $`\sigma _C`$ in $`J(C)`$ and of $`\sigma `$ on $`\mathrm{Gr}(V_{\mathrm{NR}})`$ are compatible; that is, $`\sigma (U_\xi )=U_{\sigma _C^{}\xi }`$. Taking into account the relation (3.10) between $`\psi _{u,\xi _0}`$ and $`\psi _{v,\xi _1}`$ for $`L_{\xi _1}=\sigma _C^{}L_{\xi _0}`$, the second part of 2. follows.
$`2.1.`$
From 2.a) and Theorem 6.2, we know that there exists a triple $`(C,\overline{x},t_{\overline{x}})`$ such that $`X_\mathrm{\Omega }`$ is isomorphic as a p.p.a.v. to the Jacobian of $`C`$.
Now 2.b) shows that $`U_{\xi _1}=\sigma (U_{\xi _0})`$, so the orbit $`\mathrm{\Gamma }_V(U_{\xi _0},U_{\xi _1})/\overline{\mathrm{\Gamma }}_V^+`$ is finite dimensional. Theorem 5.8 (see also Remark 5.9) implies that $`\sigma `$ induces an automorphism of $`C`$ satisfying the conditions of 1. ∎
###### Remark 6.8.
Observe that if the condition 2. of Theorem 6.2 holds for one $`\xi ^g`$, then it holds for every $`\xi `$ (see \[Sh\], Theorem 6). Therefore, if the condition 2. of Theorem 6.5 or 6.7 holds for a given $`\xi _0^g`$, then it holds for every $`\xi _0`$.
Finally, we obtain a solution of the Schottky problem for curves with automorphisms.
###### Theorem 6.9 (Characterization).
Let $`X_\mathrm{\Omega }`$ be an irreducible p.p.a.v. of dimension $`g>1`$.
Then the following conditions are equivalent.
1. There exists a projective irreducible smooth curve $`C`$ of genus $`g`$ with a non-trivial automorphism $`\sigma _C:CC`$ such that $`X_\mathrm{\Omega }`$ is isomorphic as a p.p.a.v. to the Jacobian of $`C`$.
2. There exist a prime number $`p`$, $`p`$ matrices $`A^{(1)},\mathrm{},A^{(p)}`$ ($`A^{(j)}`$ being a $`g\times \mathrm{}`$-matrix of rank $`g`$) and $`p`$ symmetric quadratic forms $`Q^{(1)},\mathrm{},Q^{(p)}`$, such that
1. for some $`\xi _0^g`$, the corresponding BA-functions satisfy the $`(1,\stackrel{p}{\mathrm{}},1)`$-KP hierarchy
$$\mathrm{Res}\left(\underset{j=1}{\overset{p}{}}z^{\delta _{ju}\delta _{jv}}\psi _{u,\xi _0}^{(j)}(z,t)\psi _{v,\xi _0}^{(j)}(z,s)\right)dz=\mathrm{\hspace{0.17em}0}$$
2. there exist $`\xi _1^g`$ (depending on $`\xi _0`$) such that
$$\mathrm{Res}\left(\underset{j=1}{\overset{p}{}}z^{\delta _{ju}\delta _{jv}}\psi _{v+1,\xi _0}^{(j+1)}(z,\sigma ^{}(t))\psi _{u,\xi _1}^{(j)}(z,s)\right)dz=\mathrm{\hspace{0.17em}0}$$
where $`\sigma ^{}(t):=(t^{(p)},t^{(1)},t^{(2)},\mathrm{},t^{(p1)})`$.
###### Proof.
Observe that any curve with non-trivial automorphism group admits an automorphism of prime order $`p`$ with an orbit consisting of $`p`$ pairwise distinct points. By the previous theorem, we conclude.
###### Remark 6.10.
Recall that standard arguments allow us to express the above equations as an infinite system of partial differential equations for the $`\tau `$-function.
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# Contents
## 1 Introduction
In this paper we continue to describe the computer system SANC Support of Analytic and Numerical Calculations for experiments at Colliders intended for semi-automatic calculations of realistic and pseudo-observables for various processes of elementary particle interactions at the one-loop precision level. This is done in the spirit of the first description of the SANC system (see and references therein which we recommend as a first acquaintance with the system).
Here we consider the implementation of several processes of $`ffbb0`$ kind (where $`f`$ stands for a fermion, $`b`$ for a boson of the Standard Model (SM), while for concrete bosons we use $`A`$ for the photon and $`Z,W^\pm ,H`$) One should emphasize also that the notation $`ffbb0`$ means that all external 4-momenta flow inwards; this is the standard SANC convention which allows to compute one-loop covariant amplitude (CA) and form factors (FF) only once and obtain it for a concrete channel by means of a crossing transformation. The present level of the system is realized in version v.1.10. Compared to version v.1.00, it is upgraded both physics-wise and computer-wise. As far as physics is concerned it contains an upgraded treatment of $`u\overline{d}l^+\nu _l`$ and $`d\overline{u}l^{}\overline{\nu }_l`$ processes (see Ref. ) and a complete implementation of $`Ff+f_1+\overline{f_1^{^{}}}`$ CC decays up to numbers and MC generators. (Here $`F`$ and $`f`$ stand for massive fermions and $`f_1`$ and $`\overline{f_1^{^{}}}`$ for massless fermions of the first generation.) Although the version 1.10 tree literally contains only $`tb+l^++\nu _l`$ decay , any decay of the kind $`Ff+f_1+\overline{f_1^{^{}}}`$ may be treated in a similar manner and we are going to implement them into the next versions. The complete description of these CC decays will be given elsewhere . Version 1.10 contains also the process $`Hf_1\overline{f}_1A`$ in three cross channels in EW branch, $`\gamma \gamma \gamma \gamma `$ scattering and $`ll\gamma \gamma ^{}`$ in QED branch, as well as a new QCD branch .
New in version 1.10 are also several $`ffbb0`$ processes, to whose implementation this paper is devoted. We describe here two of them: $`f_1\overline{f}_1ZZ0`$ and $`f_1\overline{f}_1HZ0`$, the latter one being used in two channels — annihilation and decay.
In the annihilation channel, these processes were considered in the literature extensively (see, for instance, and ), however, we are not aware of publications devoted to the $`HZf_1\overline{f}_1`$ decay.
These processes are relevant for $`H`$ search at LHC: the processes $`f_1\overline{f}_1ZZ`$ are one of the backgrounds while the one-loop calculations of the decay $`HZf_1\overline{f}_1`$ was also used for an improved treatment of the decay $`H4\mu `$ for an intermediate Higgs mass interval 130 GeV $`M_H`$ 150 GeV, see section 7.
Furthermore, in the spirit of the adopted SANC approach, all $`2f2b0`$ processes can be computed with off shell bosons thereby allowing their use also as building blocks for future studies of $`50`$ processes.
The process $`f_1\overline{f}_1ZZ`$ is very similar to the processes $`ffb\gamma `$ ($`b=\gamma ,Z,H`$) whose precomputation was described in detail in Ref. . Its tree level amplitude is represented by two diagrams in $`t`$ and $`u`$ channels, Fig. 1.
Note, that in this and in the following figures, fu=typeIU etc. denote “types” of external particles, see Table 2 of as well as the discussion in the beginning of section 4.
For the process $`f_1\overline{f}_1HZ`$ these two diagrams do not contribute in the tree approximation since fermion $`f_1`$ is considered to be massless. However, in this case there exists an $`s`$ channel amplitude, Fig. 2,
All $`ffbb`$ processes are fully implemented at Level 1 of analytical calculations. Several new modules which compute the contribution of the $`bbb`$ vertices to $`ffbb0`$ processes are added to the precomputation tree, as well as three other modules relevant for the $`s`$ channel diagram.
The modified “Precomputation” tree is shown in Fig. 5 and discussed in section 3.
The modified branch 2f2b for the “Processes” tree is shown in Fig. 3. It contains four new sub-menus $`f_1\overline{f}_1ZZ`$, $`f_1\overline{f}_1HZ`$, $`Hf_1\overline{f}_1Z`$ and $`Hf_1\overline{f}_1A`$ which in turn are branched into scalar Form Factors (FF) and Helicity Amplitudes (HA) (two for the process $`f_1\overline{f}_1HZ`$ corresponding to two annihilation and decay channels) and the accompanying bremsstrahlung contributions (BR).
These processes are implemented also at Level 2, where the s2n.f package produces the results in the “Semi Analytic” mode (see Fig. 20 of Ref. ). For the three decays we have relevant “Monte Carlo” generators which, however, are not yet implemented into the system.
The paper is organized as follows:
In section 2 we describe the covariant (CA) and helicity amplitudes for three of four new $`ffbb0`$ processes available in version 1.10.
Section 3 contains a brief description of new precomputation modules. In section 4 we describe in some more detail the renormalization procedure for the $`f_1\overline{f}_1HZ0`$ process, i.e. calculation of FFs. Section 5 contains the results for the accompanying bremsstrahlung in the semi-analytic mode for two $`f_1\overline{f}_1HZ0`$ channels. Section 6 contains numerical results for the processes $`f_1\overline{f}_1HZ(ZZ)`$ and decay $`HZf_1\overline{f}_1`$. Finally, section 7 contains a brief description a Monte Carlo generator for process $`H4\mu `$ in the single resonance approximation. The first results of numerical comparison with those of Prophecy4f are also presented.
## 2 Amplitude Basis, Scalar Form Factors, Helicity Amplitudes
### 2.1 Introduction
In this section we continue the presentation of formulae for the amplitudes of $`ffbb0`$ processes started in section 2 of Ref. . As usual, we begin with the calculation of CAs corresponding to a result of the straightforward computation of all diagrams contributing to a given process at the one-loop level. It is represented in a certain basis of structures, made of strings of Dirac matrices and external momenta, contracted with polarization vectors of vector bosons. The amplitude is parameterized by a number of FFs, which we denote by $``$ with an index labeling the corresponding structure. The number of FFs is by construction equal to the number of structures, however for the cases presented below some of the FFs can be equal, so the number of independent FFs may be less than the number of structures. For the existing tree level structures the corresponding FFs have the form
$`=1+{\displaystyle \frac{\alpha }{4\pi s_W^2}}\stackrel{~}{},`$ (1)
where “1” is due to the Born level and $`\stackrel{~}{}`$ is due to the one-loop level. As usual, we use various coupling constants:
$`Q_f,I_f^{(3)},\sigma _f=v_f+a_f,\delta _f=v_fa_f,s_W={\displaystyle \frac{e}{g}},c_W={\displaystyle \frac{M_W}{M_Z}},\text{etc.}`$ (2)
Given a CA, SANC computes a set of HAs, denoted by $`_{\lambda _1\lambda _2\lambda _3\mathrm{}}`$, where $`\lambda _1\lambda _2\lambda _3\mathrm{}`$ denote the signs of particle spin projections onto a quantization axis.
### 2.2 $`f\overline{f}ZZ`$ process
Here we present the CA of the process $`f(p_2)\overline{f}(p_1)Z(p_3)Z(p_4)`$ in the annihilation channel <sup>1</sup><sup>1</sup>1The other channels are unphysical in this case., see Fig. 1.
It contains 10 left $`(\gamma _+)`$ and 10 right $`(\gamma _{})`$ structures:
$`𝒜_{ffZZ}`$ $`=`$ $`k_0\{[\overline{v}\left(p_1\right)(/p_3\gamma _+(p_1)_\mu (p_1)_\nu _1^+(s,t)+/p_3\gamma _+(p_1)_\mu (p_2)_\nu _2^+(s,t)`$ (3)
$`+/p_3\gamma _+(p_1)_\nu (p_2)_\mu _3^+(s,t)+/p_3\gamma _+(p_2)_\mu (p_2)_\nu _4^+(s,t)+/p_3\gamma _+\delta _{\mu \nu }_5^+(s,t)`$
$`+\gamma _\mu /p_3\gamma _\nu \gamma _+_6^+(s,t)+\gamma _\mu \gamma _+(p_1)_\nu _7^+(s,t)+\gamma _\mu \gamma _+(p_2)_\nu _8^+(s,t)`$
$`+\gamma _\nu \gamma _+(p_1)_\mu _9^+(s,t)+\gamma _\nu \gamma _+(p_2)_\mu _{10}^+(s,t))u\left(p_2\right)\epsilon _\nu ^Z(p_3)\epsilon _\mu ^Z(p_4)]`$
$`+[\gamma _+\gamma _{},_i^+_i^{}]\},`$
where
$$k_0=\frac{ig^2}{8c_W^2}\text{and}\gamma _\pm =I\pm \gamma _5.$$
(4)
Furthermore,
$$\left(p_1+p_2\right)^2=s,\left(p_2+p_3\right)^2=t,\left(p_2+p_4\right)^2=u.$$
(5)
Now we give the explicit form of the CA in the tree (Born) approximation:
$`𝒜_{ffZZ}^{Born}`$ $`=`$ $`k_0\{[\sigma _f^2\overline{v}\left(p_1\right)({\displaystyle \frac{t+u}{tu}}\gamma _\mu /p_3\gamma _\nu \gamma _++{\displaystyle \frac{2}{t}}\gamma _\mu \gamma _+(p_2)_\nu +{\displaystyle \frac{2}{u}}(/p_3\gamma _+\delta _{\mu \nu }`$ (6)
$`\gamma _\mu \gamma _+(p_1)_\nu +\gamma _\nu \gamma _+(p_1)_\mu +\gamma _\nu \gamma _+(p_2)_\mu ))u\left(p_2\right)\epsilon _\nu ^Z(p_3)\epsilon _\mu ^Z(p_4)]`$
$`+[\sigma _f^2\delta _f^2,\gamma _+\gamma _{}]\}.`$
Note that this is decomposed into 12 structures of 20 and is highly asymmetric in $`t`$ and $`u`$. This is due to our choice of the 4-momentum $`p_3`$ and of the ordering of Lorentz indices $`\mu `$ and $`\nu `$ in Eq. (6).
Equation 6 may be parameterized by only two FFs if one introduces two “Born-like structures (BLS)” given by expressions in big round brackets by means of eliminating the 5th structure $`/p_3\gamma _+\delta _{\mu \nu }`$ in favor of BLS; to this structure and to the corresponding $`_0^\pm (s,t)`$ we assign the subindex “0”:
$`/p_3\gamma _\pm \delta _{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{u}{2}}\left[\text{BLS}_0^\pm +\left({\displaystyle \frac{1}{t}}+{\displaystyle \frac{1}{u}}\right)\gamma _\mu /p_3\gamma _\nu \gamma _\pm \right]{\displaystyle \frac{u}{t}}\gamma _\mu \gamma _\pm (p_2)_\nu +\gamma _\mu \gamma _\pm (p_1)_\nu \gamma _\nu \gamma _\pm (p_1)_\mu +\gamma _\nu \gamma _\pm (p_2)_\mu .`$ (7)
Moreover, between the 20 FFs there are four identities:
$$_4^\pm (s,t)=_1^\pm (s,t),_{10}^\pm (s,t)=_7^\pm (s,t).$$
(8)
Therefore, there are 16 independent FFs but 18 independent non-zero HAs for process $`f_1\overline{f}_1ZZ`$:
$`_{+\pm }`$ $`=`$ $`k_0^sc_\pm \{2\sigma _e^2({\displaystyle \frac{1}{t}}+{\displaystyle \frac{1}{u}})_0^+(s,t)+{\displaystyle \frac{s}{4}}c_{}\beta [2_1^+(s,t)_2^+(s,t)_3^+(s,t)]`$
$`2_7^+(s,t)\pm _8^+(s,t)_9^+(s,t)\},`$
$`_{+\pm }`$ $`=`$ $`k_0^sc_{}\{2\delta _e^2({\displaystyle \frac{1}{t}}+{\displaystyle \frac{1}{u}})_0^{}(s,t){\displaystyle \frac{s}{4}}c_\pm \beta [2_1^{}(s,t)_2^{}(s,t)_3^{}(s,t)]`$
$`2_7^{}(s,t)\pm _8^{}(s,t)_9^{}(s,t)\},`$
$`_{+\pm \pm }`$ $`=`$ $`k_0^s\{2\sigma _e^2({\displaystyle \frac{\beta _{}^c}{t}}{\displaystyle \frac{\beta _+^c}{u}})_0^+(s,t){\displaystyle \frac{s}{4}}\mathrm{sin}^2\vartheta _Z\beta [2_1^+(s,t)_2^+(s,t)_3^+(s,t)]`$
$`+2[\beta _{}_6^+(s,t)+\mathrm{cos}\vartheta _Z_7^+(s,t)]\pm c_{}_8^+(s,t)\pm c_\pm _9^+(s,t)\},`$
$`_+`$ $`=`$ $`k_0^s\{2\delta _e^2({\displaystyle \frac{\beta _{}^c}{t}}{\displaystyle \frac{\beta _+^c}{u}})_0^{}(s,t)+{\displaystyle \frac{s}{4}}\mathrm{sin}^2\vartheta _Z\beta [2_1^{}(s,t)_2^{}(s,t)_3^{}(s,t)]`$
$`2[\beta _{}_6^{}(s,t)+\mathrm{cos}\vartheta _Z_7^{}(s,t)]{\displaystyle \frac{1}{2}}c_{}_8^{}(s,t){\displaystyle \frac{1}{2}}c_\pm _9^{}(s,t)\},`$
$`_{+\pm 0}`$ $`=`$ $`k_1^\pm \{2\sigma _{el}^2[{\displaystyle \frac{\beta _{}^c}{t}}{\displaystyle \frac{\beta _+^c}{u}}\pm {\displaystyle \frac{2M_Z^2}{s}}({\displaystyle \frac{1}{t}}+{\displaystyle \frac{1}{u}})]_0^+(s,t)`$
$`{\displaystyle \frac{s}{4}}c_{}\beta \left[2\mathrm{cos}\vartheta _Z_1^+(s,t)\beta _{}^c_2^+(s,t)\pm \beta _+^c_3^+(s,t)\right]`$
$`\pm \beta _{}^2_6^+(s,t)(\beta _+^cc_{})_7^+(s,t)c_{}_8^+(s,t)+\beta _{}^c_9^+(s,t)\},`$
$`_{+0\pm }`$ $`=`$ $`k_1^{}\{2\sigma _e^2[{\displaystyle \frac{\beta _{}^c}{t}}{\displaystyle \frac{\beta _+^c}{u}}{\displaystyle \frac{2M_Z^2}{s}}({\displaystyle \frac{1}{t}}+{\displaystyle \frac{1}{u}})]_0^+(s,t)`$
$`{\displaystyle \frac{s}{4}}c_\pm \beta \left[\pm 2\mathrm{cos}\vartheta _Z_1^+(s,t)\pm \beta _{}^c_2^+(s,t)\beta _+^c_3^+(s,t)\right]`$
$`\pm 4{\displaystyle \frac{M_Z^2}{s}}_6^+(s,t)(\beta _+^c\pm c_\pm )_7^+(s,t)\beta _{}^c_8^+(s,t)c_\pm _9^+(s,t)\},`$
$`_{+\pm 0}`$ $`=`$ $`k_1^{}\{2\delta _{el}^2[{\displaystyle \frac{\beta _{}^c}{t}}{\displaystyle \frac{\beta _+^c}{u}}{\displaystyle \frac{2M_Z^2}{s}}({\displaystyle \frac{1}{t}}+{\displaystyle \frac{1}{u}})]_0^{}(s,t)`$
$`+{\displaystyle \frac{s}{4}}c_\pm \beta \left[2\mathrm{cos}\vartheta _Z_1^{}(s,t)\beta _{}^c_2^{}(s,t)\pm \beta _+^c_3^{}(s,t)\right]`$
$`\beta _+^2_6^{}(s,t)(\beta _+^c\pm c_\pm )_7^{}(s,t)\pm c_\pm _8^{}(s,t)+\beta _{}^c_9^{}(s,t)\},`$
$`_{+0\pm }`$ $`=`$ $`k_1^\pm \{2\delta _{el}^2[{\displaystyle \frac{\beta _{}^c}{t}}{\displaystyle \frac{\beta _+^c}{u}}\pm {\displaystyle \frac{2M_Z^2}{s}}({\displaystyle \frac{1}{t}}+{\displaystyle \frac{1}{u}})]_0^{}(s,t)`$
$`\pm {\displaystyle \frac{s}{4}}c_{}\beta \left(2\mathrm{cos}\vartheta _Z_1^{}(s,t)+\beta _{}^c_2^{}(s,t)\beta _+^c_3^{}(s,t)\right)`$
$`4{\displaystyle \frac{M_Z^2}{s}}_6^{}(s,t)(\beta _+^cc_{})_7^{}(s,t)\beta _{}^c_8^{}(s,t)\pm c_{}_9^{}(s,t)\},`$
$`_{\pm 00}`$ $`=`$ $`{\displaystyle \frac{k_0^s}{2}}{\displaystyle \frac{s}{M_Z^2}}\{2\left(\begin{array}{c}\sigma _{el}^2\\ \delta _{el}^2\end{array}\right)[{\displaystyle \frac{\beta _{}^c}{t}}{\displaystyle \frac{\beta _+^c}{u}}+2\beta _f{\displaystyle \frac{M_Z^2}{s}}({\displaystyle \frac{1}{t}}{\displaystyle \frac{1}{u}})]_0^\pm (s,t)`$ (12)
$`{\displaystyle \frac{1}{4}}s\beta \left[2\beta _+^c\beta _{}^c_1^\pm (s,t)+(\beta _{}^c)^2_2^\pm (s,t)+(\beta _+^c)^2_3^\pm (s,t)\right]`$
$`+4\beta {\displaystyle \frac{M_Z^2}{s}}_6^\pm (s,t)\pm 2\beta _+^c_7^\pm (s,t)\pm \beta _{}^c[_8^\pm (s,t)_9^\pm (s,t)]\}.`$
Here we use the following shorthand notation:
$`k_0^s`$ $`=`$ $`{\displaystyle \frac{k_0s}{2}}\mathrm{sin}\vartheta _Z,k_1^\pm ={\displaystyle \frac{k_0s}{2}}{\displaystyle \frac{\sqrt{s}}{\sqrt{2}M_Z}}c_\pm ,c_\pm =1\pm \mathrm{cos}\vartheta _Z,`$
$`\beta _\pm `$ $`=`$ $`\beta \pm 1,\beta _\pm ^c=\beta \pm \mathrm{cos}\vartheta _Z,\beta ={\displaystyle \frac{\sqrt{\lambda (s,M_Z^2,M_Z^2)}}{s}},`$ (13)
and $`\vartheta _Z`$ is the CMS angle between $`\stackrel{}{p}_2`$ and $`\stackrel{}{p}_3`$. The invariant $`t`$ and the cosine $`\mathrm{cos}\vartheta _Z`$ are related by
$$t=M_Z^2\frac{1}{2}s(1\beta \mathrm{cos}\vartheta _Z).$$
(14)
The number 18 is the product of 2 initial massless helicity states and 3$`\times `$3 states for the final $`Z`$ bosons.
### 2.3 $`f_1\overline{f}_1HZ`$ process
There are six structures for the $`f_1\overline{f}_1HZ`$ process if the fermion mass is neglected
$`𝒜_{ffHZ}`$ $`=`$ $`k\{\left[\overline{v}\left(p_1\right)(\gamma _\nu \gamma _+\sigma _f_0^+(s,t)+/p_3\gamma _+(p_1)_\nu _1^+(s,t)+/p_3\gamma _+(p_2)_\nu _2^+(s,t))u\left(p_2\right)\epsilon _\nu ^Z(p_3)\right]`$ (15)
$`+[\sigma _f\delta _f,\gamma _+\gamma _{},_i^+(s,t)_i^{}(s,t)]\},`$
where
$`k={\displaystyle \frac{ig^2}{4c_W^2}}{\displaystyle \frac{M_Z}{M_Z^2s}}.`$ (16)
The structures for the decay $`Hf_1\overline{f}_1Z`$ may be obtained by simple replacement of 4-momenta $`p_1p_3,p_2p_4,p_4p_1(p_3p_2)`$ of the structures (15).
Note, that the two terms $`\gamma _\nu \gamma _+\sigma _f_0^\pm (s,t)`$ correspond to the Born level.
As far as HAs are concerned, we present them in both channels: annihilation and decay.
#### 2.3.1 HAs in annihilation channel $`f_1\overline{f}_1HZ`$
There are 6 HAs in this case:
$`_{++}`$ $`=`$ $`k_0^sc_+\left\{k_1^{}\left[_2^+(s,t)_1^+(s,t)\right]4\sigma _e_0^+(s,t)\right\},`$
$`_{++}`$ $`=`$ $`k_0^sc_{}\left\{k_1^+\left[_1^{}(s,t)_2^{}(s,t)\right]+4\delta _e_0^{}(s,t)\right\},`$
$`_+`$ $`=`$ $`k_0^sc_{}\left\{k_1^+\left[_1^+(s,t)_2^+(s,t)\right]+4\sigma _e_0^+(s,t)\right\},`$
$`_+`$ $`=`$ $`k_0^sc_+\left\{k_1^{}\left[_2^{}(s,t)_1^{}(s,t)\right]4\delta _e_0^{}(s,t)\right\},`$
$`_{+0}`$ $`=`$ $`k_0^sk_2\left\{\sqrt{\lambda (s,M_Z^2,M_H^2)}\left[\beta _+^c_1^+(s,t)+\beta _{}^c_2^+(s,t)\right]+4\sigma _e_0^+(s,t)\right\},`$
$`_{+0}`$ $`=`$ $`k_0^sk_2\left\{\sqrt{\lambda (s,M_Z^2,M_H^2)}\left[\beta _+^c_1^{}(s,t)+\beta _{}^c_2^{}(s,t)\right]+4\delta _e_0^{}(s,t)\right\}.`$ (17)
where
$`k_0^s`$ $`=`$ $`k_0{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{\sqrt{s}M_Z}{sM_Z^2}},k_1^\pm =\sqrt{\lambda (s,M_Z^2,M_H^2)}c_\pm ,k_2={\displaystyle \frac{s+M_Z^2M_H^2}{\sqrt{2}\sqrt{s}M_Z}}\mathrm{sin}\vartheta _Z,`$
$`c_\pm `$ $`=`$ $`1\pm \mathrm{cos}\vartheta _Z,\beta _\pm ^c=\beta \pm \mathrm{cos}\vartheta _Z,\beta ={\displaystyle \frac{\sqrt{\lambda (s,M_Z^2,M_H^2)}}{s+M_Z^2M_H^2}},`$
$`t`$ $`=`$ $`M_Z^2{\displaystyle \frac{1}{2}}(s+M_Z^2M_H^2)(1\beta \mathrm{cos}\vartheta _Z).`$ (18)
#### 2.3.2 HAs in the decay channel $`Hf_1\overline{f}_1Z`$
The six HAs in this case are somewhat different from the previous case:
$`_{++}`$ $`=`$ $`k_0^s\left\{k_1\left[_1^{}(s,t)_2^{}(s,t)\right]4\delta _fc_{}_0^{}(s,t)\right\},`$
$`_{++}`$ $`=`$ $`k_0^s\left\{k_1\left[_2^+(s,t)_1^+(s,t)\right]4\sigma _fc_+_0^+(s,t)\right\},`$
$`_+`$ $`=`$ $`k_0^s\left\{k_1\left[_1^{}(s,t)_2^{}(s,t)\right]+4\delta _fc_+_0^{}(s,t)\right\},`$
$`_+`$ $`=`$ $`k_0^s\left\{k_1\left[_2^+(s,t)_1^+(s,t)\right]+4\sigma _fc_{}_0^+(s,t)\right\},`$
$`_{0+}`$ $`=`$ $`k_0^sk_2\left\{\sqrt{\lambda (M_H^2,M_Z^2,s)}\left[\beta _+^c_1^{}(s,t)+\beta _{}^c_2^{}(s,t)\right]4\delta _e_0^{}(s,t)\right\},`$
$`_{0+}`$ $`=`$ $`k_0^sk_2\left\{\sqrt{\lambda (M_H^2,M_Z^2,s)}\left[\beta _+^c_1^+(s,t)+\beta _{}^c_2^+(s,t)\right]+4\sigma _e_0^+(s,t)\right\}.`$ (19)
Here,
$`k_0^s`$ $`=`$ $`k_0{\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{\sqrt{s}M_Z}{M_Z^2s}},k_1=\sqrt{\lambda (M_H^2,M_Z^2,s)}\mathrm{sin}^2\vartheta _f,k_2={\displaystyle \frac{\left(M_H^2M_Z^2s\right)}{\sqrt{2}\sqrt{s}M_Z}}\mathrm{sin}\vartheta _f,`$
$`c_\pm `$ $`=`$ $`1\pm \mathrm{cos}\vartheta _f,\beta _\pm ^c=\beta \pm \mathrm{cos}\vartheta _f,\beta ={\displaystyle \frac{\sqrt{\lambda (M_H^2,M_Z^2,s)}}{M_H^2M_Z^2s}}.`$ (20)
The number 6 is the product of 2 initial massless helicity states and 3 states of the final $`Z`$ boson.
Furthermore, $`s=M_{f_1\overline{f}_1}^2`$ is the invariant mass of the two fermions $`f`$, varying in the limits $`4m_f^2s(M_HM_Z)^2`$; and $`t`$ is another independent kinematical variable, depending on $`s`$ and an angle $`\vartheta _f`$, varying in the limits $`0\vartheta _f\pi `$
$$t=M_Z^2+\frac{1}{2}\left[M_H^2M_Z^2s\sqrt{\lambda (M_H^2,M_Z^2,s)}\mathrm{cos}\vartheta _f\right].$$
(21)
The kinematical diagram of the process is shown in Fig. 4.
The Higgs boson with momentum $`p_2`$ at rest, decays back-to-back into a $`Z`$ boson with momentum $`p_1`$ and a fermionic compound with 4-momentum $`p_3+p_4`$ and invariant mass $`s`$. This compound decays in its own rest frame into two back-to-back fermions with $`\vartheta _f`$ being the angle between $`p_3`$ in the compound rest frame and the direction of flight of the $`Z`$ boson in the $`H`$ boson rest frame.
## 3 Precomputation news
The “Precomputation” tree of version 1.10 is shown in Fig. 5 with all modified sub-menus open, and all sub-menus closed which were not changed compared to version 1.00. In this section we briefly discuss what every new module does.
First, we added a new folder accessible via menu sequence EW $``$ Precomputation $``$ Vertex $``$ bbb $``$ ffbb bbb with five modules ffXX bbb Vertex, XX=AA, ZA, HA, ZZ, HZ which compute three boson vertices of four topologies (see Fig. 11 of Ref. ) to the corresponding $`ffbb`$ processes, Fig. 6. The results of their calculations are saved to ffXX\*.sav files to be loaded by corresponding modules computing FFs via chains EW $``$ Processes $``$ 4-legs $``$ 2f2b $``$ Neutral Current for the $`ffXX0`$ processes.
These three boson diagrams contain both bosonic and fermionic components. The latter are precomputed by the modules bbb Vertex in Boson and Fermion folders of the same level on the tree. Five modules of ffbb bbb folder load them and then apply tedious calculations involving in some cases the Schouten identity. There are many peculiarities in these calculations, forcing us to have an individual module for each $`ffbb`$ process. Note also that if the corresponding process has a Born-level $`s`$ channel exchange as in Fig. 2, then the contribution of one-loop vertices is supplemented by the relevant counterterm cross.
In the modules under discussion a summation over the exchanged boson $`B`$ is performed. In general, four neutral bosons $`B=\gamma ,Z,\varphi ^0`$ and $`H`$ can contribute if the fermion mass is not neglected, otherwise, only $`\gamma `$ and $`Z`$ contribute.
For the processes $`ffAA0`$, $`ffZA0`$, $`ffHA0`$ the “left” $`bff`$ vertex diagram, shown in Fig. 7 does not contribute, since in these cases the “right” vertex does not exist at the tree level. In general, and this is indeed the case for the processes $`ffZZ0`$ and $`ffHZ0`$, the “right” vertex exists at the tree level for $`B=Z,\varphi ^0`$, therefore, the dressed “left” vertex has to be added to the precomputation tree. Note that it does not contribute for massless fermions if $`B=\varphi ^0`$. For $`ffHZ0`$ this vertex is accessible via menu sequence EW $``$ Precomputation $``$ Vertex $``$ bff $``$ ffHZ bff Vertex. For the $`ffZZ0`$ process, only $`B=H`$ contributes, but then, for the massless $`f`$, the dressed “left” vertex vanishes. This is why we do not add the corresponding module in the latter case.
The presence of an $`s`$ channel tree level diagram in the process $`f_1\overline{f}_1HZ0`$ (Fig. 2) forces us to take into account two more self energy diagrams, Fig. 8.
The first one is accessible via menu chain: Self $``$ Boson $``$ Ren Blocks $``$ ffHZ Ren Self. Here $`B_1=\gamma ,Z,\varphi ^0`$ and $`B_2=Z,\varphi ^0`$; again $`\varphi ^0`$’s do not contribute for massless fermions. As will be explained in the next section, the second diagram is better to be combined with the “right” vertex, Fig. 6.
Note that nothing is changed, compared to $`ffXA`$ processes, as far as boxes are concerned. So, the Box sub-menu is as in version 1.00 .
The new modules contain calls to several new intrinsic procedures which will be described elsewhere.
## 4 Renormalizaton for $`ffHZ0`$ process
In this section we describe how to use the FORM module which computes FFs for the process $`ffHZ0`$. This description is supposed to help a user to understand the other modules computing FFs for any $`ffbb0`$ process.
First of all, to use our basic declaration and notation we begin the file with
```
#include Declar.h
#call Globals()
```
and define types of external particles, see Figs. 1 and 2.
```
#ifdef ‘typeIU’; * ‘fu’
#ifdef ‘typeID’; * ‘fd’
#ifdef ‘typeFU’; * ‘vu’
#ifdef ‘typeFD’; * ‘vd’
#define typeIDp "{2*(‘typeID’%2)-1+‘typeID’}"; * ‘fdp’
```
Secondly, we fix four main steering flags to define the calculation Al scheme:
1. define xi: xi = 0 to test gauge invariance in $`R_\xi `$, or xi = 1 to work in $`\xi =1`$ gauge;
2. define on: on = 0 external photons are off mass-shell, or on = 1 photons are on mass-shell;
3. define mf: mf = 0 zero external fermion mass (i.e. pm(‘fd’)=0), or mf = 1 it is not zeroed;
4. define mp: mp = 0 zero mass of the weak isospin partner of the fermion $`f`$, pm(‘fdp’)=0, or mp = 1 it is not zeroed.
Actually, for the process under consideration, $`f_1\overline{f}_1HZ0`$, only the $`\xi `$ and mp definitions are meaningful since there are no external photons and we ignore the masses of external fermions throughout the calculations. But the mass of a weak isospin partner of an external fermion that appears in the internal loop may be kept nonzero.
The ideology of building blocks (BB) is the key element for SANC development. The information about the main precomputed BB is stored in basic \*.sav files (BSF). Note that for 4-particle $`ffbb`$ processes all the BBs are 4-legs by construction. This trick will greatly simplify the procedure of projection of the covariant amplitude onto an independent basis of structures.
Moreover, for the future development it is necessary to upgrade the database of the collected information area SANC: fields of program modules, fields of procedures and the bank of BSF.
Any module computing FFs starts from loading of the calculated BB from the bank of BSF. These BSF contain the precomputed objects: self energies, vertices and boxes typically with off-shell bosons.
They are precomputed not only to accelerate the calculations. Although all BSF may be, in principle, precomputed online, we remind that in some cases the CPU time for calculating off-shell boxes in $`R_\xi `$ gauge takes many hours, see section 3.4 of Ref. . In such cases precomputation is strictly prohibited and the user must use already precomputed BSFs.
We recall also that our precomputation procedure has indeed several levels; in the modules computing FFs we tend to use the results of the last level which contains already renormalized BBs: propagators and vertices, i.e. taking into account relevant counterterms and special vertices. However, they are full of residual UV poles and $`\xi `$ dependent terms, which cancel in the sum for a one-loop CA of a physical process. This is why we still use to word “renormalization” in connection with modules computing FFs rather than a simple “summation”.
Typically, the loading of BSFs is organized in several steps. Let us consider the example of H $``$ f1 f1 Z (FF) module, see the tree in Fig. 5.
* step self
Here we manipulate objects from BSF ffHZSelfschxi‘xi’‘fu’‘fd’‘vu’‘vd’.sav
We extract from its volume the BB of bosonic self energy in the s-channel, BSEsch‘fu’‘fd’‘vu’‘vd’ see left diagram in Fig. 8.
* step vertex
Moving further over the renormalization procedure at this step we load three BSFs:
ffHZVertbffxi‘xi’‘fu’‘fd’‘vu’‘vd’.sav;
ffHZVertbbbxi‘xi’‘fu’‘fd’‘vu’‘vd’.sav;
ffbbVertxi‘xi’on‘on’mf‘mf’mp‘mp’‘fu’‘fd’‘vu’‘vd’.sav.
From these BSFs we extract various types of vertices correspondingly:
VertBff‘i’‘fu’‘fd’‘vu’‘vd’ with i=1,2,3,4 standing for $`\xi _A`$, $`\xi _Z`$, $`\xi _W`$, and no $`\xi `$ vertex clusters originating from the diagram of Fig. 7;
Vertbbbbos‘fu’‘fd’‘vu’‘vd’ and Vertbbbfer‘fu’‘fd’‘vu’‘vd’ — the bosonic and fermionic components of three-boson vertices shown in the diagram Fig. 6, where the former contains counterterms, the special vertex and the right diagram of Fig. 8. These tadpoles cancel the $`\xi _Z`$ dependence of the three-boson vertices, giving an opportunity to assign this contribution to i=3. Finally, the fermionic component should be naturally assigned to i=4;
— abelian Vert‘I’‘i’ and non-abelian vert‘I’‘i’ vertex clusters in $`t`$ and $`u`$ channels I=t,u with cluster index i=1,2,3,4 and k=22,33,24,42,44, see section 3.4.2 of Ref. for a description of the latter.
* step boxes
Here the most complex building blocks — off-shell boxes are loaded from four BSFs:
ffbb3T1xi‘xi’on‘on’mf‘mf’mp‘mp’‘fu’‘fd’‘vd’‘vu’.sav;
ffbb3T3xi‘xi’on‘on’mf‘mf’mp‘mp’‘fu’‘fd’‘vu’‘vd’.sav;
ffbb22T5xi‘xi’on‘on’mf‘mf’mp‘mp’‘vd’‘fd’‘vu’‘fu’.sav;
ffbb33T5xi‘xi’on‘on’mf‘mf’mp‘mp’‘vd’‘fd’‘vu’‘fu’.sav.
They contain precomputed boxes of topology T1 from the expression S3T1‘xi’‘on’‘mf’‘mp’‘fu’‘fd’ ‘vd’‘vu’, the boxes of topology T3 from the expression S3T3‘xi’‘on’‘mf’‘mp’‘fu’‘fd’‘vu’‘vd’, and of topology T5 with cluster index $`k1=2,3`$ , i.e. with virtual $`Z`$ and $`W`$ bosons from the expression S‘k1’‘k1’T5‘xi’‘on’‘mf’‘mp’‘vd’‘fd’‘vu’‘fu’, see section 3.4.2 of Ref. .
Only those box topologies and clusters are loaded which give a non-zero contribution for $`m_f=0`$.
* step Sum
Finally, we sum all contributions. Four expressions, Sum‘i’, corresponding to cluster index i=1,2,3,4 are being constructed here. The first three of them may carry only one gauge parameter each, the latter carries none.
After construction of four Sum‘i’, the module continues with various kinds of transformations (in particular, involving an algebra of Gram determinants) which prove the cancellation of gauge parameter dependences in first three Sum1,2,3 and the cancellation of the residual UV poles between FFs with cluster indices i=3 and i=4.
After the comment Preparing Structures for obtaining FFs: the six basis elements of the CA, Eq (15), are created and the 6$`\times 4`$ FFs are projected out of this CA.
The final step is formatting of the BSF FFf1f1HZ.sav with $`24`$ FFs --- FFgp‘k’‘i’ and FFgm‘k’‘i’ ($`k=0,1,2`$, $`i=1,4`$) --- for subsequent processing by s2n.f software.
## 5 Bremsstrahlung in $`f_1\overline{f}_1HZ0`$ processes
In this section we present the list of short final results for the contribution of accompanying bremsstrahlung processes.
### 5.1 Bremsstrahlung in $`f_1\overline{f}_1HZ`$ annihilation channel
The tree level diagram of this channel is shown in Fig. 2. The corresponding total Born cross section reads:
$`\sigma ^{\mathrm{Born}}={\displaystyle \frac{G_F^2(v_f^2+a_f^2)}{12\pi }}{\displaystyle \frac{M_Z^4\sqrt{\lambda (s,M_Z^2,M_H^2)}}{(M_Z^2s)^2+M_Z^2\mathrm{\Gamma }_Z^2}}\left[{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{s}}\left(5M_Z^2M_H^2\right)+{\displaystyle \frac{1}{2s^2}}\left(M_H^2M_Z^2\right)^2\right].`$ (22)
There are only two initial state (ISR) bremsstrahlung diagrams:
QED corrections due to virtual and soft photons are proportional to the Born cross section:
$`\sigma ^{\mathrm{Virt}}=\sigma ^{\mathrm{Born}}{\displaystyle \frac{\alpha }{\pi }}Q_f^2\left\{{\displaystyle \frac{1}{2}}\left[\mathrm{ln}\left({\displaystyle \frac{s}{m_f^2}}\right)1\right]^2+\left[{\displaystyle \frac{3}{2}}\mathrm{ln}\left({\displaystyle \frac{s}{\lambda ^2}}\right)\right]\left[\mathrm{ln}\left({\displaystyle \frac{s}{m_f^2}}\right)1\right]1+4\mathrm{Li}_2(1)\right\},`$ (23)
$`\sigma ^{\mathrm{Soft}}=\sigma ^{\mathrm{Born}}{\displaystyle \frac{\alpha }{\pi }}Q_f^2\left\{{\displaystyle \frac{1}{2}}\left[\mathrm{ln}\left({\displaystyle \frac{s}{m_f^2}}\right)1\right]^2+\mathrm{ln}\left({\displaystyle \frac{4\overline{\omega }^2}{\lambda ^2}}\right)\left[\mathrm{ln}\left({\displaystyle \frac{s}{m_f^2}}\right)1\right]+{\displaystyle \frac{1}{2}}2\mathrm{Li}_2(1)\right\},`$ (24)
with infrared divergence $`\mathrm{ln}\lambda ^2`$ being canceled out<sup>2</sup><sup>2</sup>2Note, that the “Soft” contribution to the process $`f_1\overline{f}_1ZZ`$ is also described by Eq.(24).. The hard photon contribution of the differential cross section in $`s^{}=(p_3+p_4)^2`$ has the following factorization property:
$`{\displaystyle \frac{d\sigma ^{\mathrm{Hard}}}{ds^{}}}={\displaystyle \frac{\alpha }{\pi }}Q_f^2{\displaystyle \frac{s^2+s_{}^{}{}_{}{}^{2}}{s^2(ss^{})}}\left[\mathrm{ln}\left({\displaystyle \frac{s}{m_f^2}}\right)1\right]\sigma ^{\mathrm{Born}}(s^{}).`$ (25)
It may be integrated over $`s^{}`$ leading to a rather compact expression for $`\sigma ^{\mathrm{tot}}`$ (see relevant module at SANC tree), yielding a $`\mathrm{ln}\overline{\omega }`$ which cancels against corresponding term in Eq. (24).
### 5.2 Bremsstrahlung in $`Hf_1\overline{f}_1Z`$ decay channel
Here we consider the decay channel $`Hf_1\overline{f}_1Z`$. We begin with the tree level diagram, Fig. 10:
The corresponding tree level double differential width, depending on two kinematical variables $`s,\vartheta _f`$ discussed in section 2.3.2 and with kinematics shown in Fig. 4, reads
$`{\displaystyle \frac{d^2\mathrm{\Gamma }^{\mathrm{Born}}}{dsd\mathrm{cos}\vartheta _f}}`$ $`=`$ $`k_B\{((v_f^2+a_f^2)[\mathrm{sin}^2\theta _f(14{\displaystyle \frac{m_f^2}{s}})+4{\displaystyle \frac{m_f^2}{s}}]8a_f^2{\displaystyle \frac{m_f^2}{M_Z^2}})`$ (26)
$`\left[\left(1{\displaystyle \frac{s}{M_H^2}}\right)^22{\displaystyle \frac{sM_Z^2}{M_H^4}}1+\left(1{\displaystyle \frac{M_Z^2}{M_H^2}}\right)^2\right]+8\left(v_f^2+a_f^2\right){\displaystyle \frac{sM_Z^2}{M_H^4}}\left(1+2{\displaystyle \frac{m_f^2}{s}}\right)`$
$`+4a_f^2{\displaystyle \frac{m_f^2}{M_Z^2}}[{\displaystyle \frac{sM_Z^2}{M_H^4}}(1{\displaystyle \frac{s}{M_Z^2}})^2+{\displaystyle \frac{s}{M_H^2}}(2{\displaystyle \frac{s}{M_Z^2}}2+{\displaystyle \frac{M_H^2}{M_Z^2}})12{\displaystyle \frac{M_Z^4}{M_H^4}}]\},`$
$`\text{where}k_B={\displaystyle \frac{1}{128}}{\displaystyle \frac{G_F^2}{\pi ^3}}{\displaystyle \frac{\sqrt{\lambda (M_H^2,M_Z^2,s)}M_Z^4M_H}{\left|M_Z^2iM_Z\mathrm{\Gamma }_Zs\right|^2}}.`$
There are only two final state bremsstrahlung diagrams, Fig. 11:
The fully differential phase space is characterized by five kinematical variables which we choose as follows:
$`d\mathrm{\Phi }^{(3)}={\displaystyle \frac{ds}{2\pi }}{\displaystyle \frac{d\tau }{2\pi }}\mathrm{\Phi }_1^{(2)}d\mathrm{\Phi }_2^{(2)}d\mathrm{\Phi }_3^{(2)},`$ (27)
where $`\tau =(p_4+p_5)^2`$ is the lepton--photon invariant mass.
The 3-step kinematical cascade develops as a sequence of three 2-body decays shown in Fig. 12,
with three corresponding two body phase spaces:
$`\mathrm{\Phi }_1^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{\sqrt{\lambda (M_H^2,M_Z^2,s)}}{s}},`$
$`d\mathrm{\Phi }_2^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{\sqrt{\lambda (s,\tau ,m_f^2)}}{s}}{\displaystyle \frac{1}{2}}d\mathrm{cos}\vartheta _3,`$
$`d\mathrm{\Phi }_3^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{\sqrt{\lambda (\tau ,m_f^2,0)}}{\tau }}{\displaystyle \frac{1}{2}}d\mathrm{cos}\vartheta _4{\displaystyle \frac{1}{2\pi }}d\phi _4.`$ (28)
The gauge invariant QED part of the complete one-loop EW correction is subdivided into virtual, soft and hard photon contributions. The virtual one comes from the two Born-like FFs with cluster index i=1. It is proportional to the Born width and contains the infrared divergence parameterized by the photon mass $`\lambda `$:
$$\frac{d^2\mathrm{\Gamma }^{\mathrm{Virt}}}{dsd\mathrm{cos}\vartheta _f}=\frac{d^2\mathrm{\Gamma }^{\mathrm{Born}}}{dsd\mathrm{cos}\vartheta _f}\frac{\alpha }{\pi }Q_f^2\{\mathrm{ln}\left(\frac{m_f^2}{\lambda ^2}\right)[\mathrm{ln}\left(\frac{s}{m_f^2}\right)1]+\frac{1}{2}\mathrm{ln}\left(\frac{s}{m_f^2}\right)[3\mathrm{ln}\left(\frac{s}{m_f^2}\right)]+4\mathrm{Li}_2(1)2\}.$$
(29)
The soft photon contribution is also proportional to the Born one; its infrared divergence cancels against the virtual contribution. It contains also a logarithm with soft-hard separator $`\overline{\omega }`$:
$$\frac{d^2\mathrm{\Gamma }^{\mathrm{Soft}}}{dsd\mathrm{cos}\vartheta _f}=\frac{d^2\mathrm{\Gamma }^{\mathrm{Born}}}{dsd\mathrm{cos}\vartheta _f}\frac{\alpha }{\pi }Q_f^2\left\{\left[\mathrm{ln}\left(\frac{m_f^2}{\lambda ^2}\right)+2\mathrm{ln}\left(\frac{2\overline{\omega }}{m_f}\right)\mathrm{ln}\left(\frac{s}{m_f^2}\right)\right]\left[\mathrm{ln}\left(\frac{s}{m_f^2}\right)1\right]\mathrm{Li}_2\left(1\right)+1\right\}.$$
(30)
The hard photon contribution after integration over three kinematical variables $`d\phi _4,d\mathrm{cos}\vartheta _4`$ and $`d\tau `$ (the first two vary together in full angular $`4\pi `$ limits and $`m_f^2\tau (\sqrt{s}m_f)^2`$) is
$`{\displaystyle \frac{d^2\mathrm{\Gamma }^{\mathrm{Hard}}}{dsd\mathrm{cos}\vartheta _f}}`$ $`=`$ $`{\displaystyle \frac{d^2\mathrm{\Gamma }^{\mathrm{Born}}}{dsd\mathrm{cos}\vartheta _f}}{\displaystyle \frac{\alpha }{\pi }}Q_f^2\{2\mathrm{ln}\left({\displaystyle \frac{2\overline{\omega }}{m_f}}\right)[\mathrm{ln}\left({\displaystyle \frac{s}{m_f^2}}\right)1]{\displaystyle \frac{1}{2}}\mathrm{ln}\left({\displaystyle \frac{s}{m_f^2}}\right)[53\mathrm{ln}\left({\displaystyle \frac{s}{m_f^2}}\right)]`$ (31)
$`3\mathrm{Li}_2(1)+{\displaystyle \frac{1}{4}}\}+k_B{\displaystyle \frac{\alpha }{\pi }}Q_f^2[(1{\displaystyle \frac{s}{M_H^2}})^2+10s{\displaystyle \frac{M_Z^2}{M_H^4}}1+(1{\displaystyle \frac{M_Z^2}{M_H^2}})^2].`$
The total QED correction, sum of above three contributions, is free not only of infrared divergence and of soft-hard separator, but also free of final fermion mass singularity in accordance with the KLN theorem --.
$$\frac{d^2\mathrm{\Gamma }^{\mathrm{Total}}}{dsd\mathrm{cos}\vartheta _f}=\frac{d^2\mathrm{\Gamma }^{\mathrm{Born}}}{dsd\mathrm{cos}\vartheta _f}\frac{\alpha }{\pi }Q_f^2+k_B\frac{\alpha }{\pi }Q_f^2\left[\left(1\frac{s}{M_H^2}\right)^2+10s\frac{M_Z^2}{M_H^4}1+\left(1\frac{M_Z^2}{M_H^2}\right)^2\right].$$
(32)
Finally, if one integrates over $`d\mathrm{cos}\vartheta _f`$, the well known $`Z`$ decay correction factor restores:
$$\frac{d\mathrm{\Gamma }^{\mathrm{Total}}}{ds}=\frac{d\mathrm{\Gamma }^{\mathrm{Born}}}{ds}\left[1+\frac{3}{4}\frac{\alpha }{\pi }Q_f^2\right].$$
(33)
Therefore, the QED part of the correction is small, $`0.2\%`$.
## 6 Numerical results and comparison
In the numerical calculations by s2n package we use two precompiled libraries: SancLib\_v1.00 and looptools 2.1 .
### 6.1 Numerical results for Electroweak corrections
#### 6.1.1 Process $`f_1\overline{f}_1HZ`$
For this process we present in Table 1 the results of a tuned triple comparison of the one-loop electroweak corrections, excluding the gauge-invariant QED subset of diagrams (vertex and electron self-energy) and real bremsstrahlung. The input parameters are taken as in . Table 1 shows 6-7 digits agreement between the three calculations.
Beside the input given in , $`M_W`$ is crucial for a precise comparison. The following $`M_W`$ masses have been used: $`M_W=80.231815`$GeV ($`M_H=100`$GeV), $`M_W=80.159313`$ GeV ($`M_H=300`$GeV), $`M_W=80.081409`$GeV ($`M_H=800`$GeV), following A. Denner private communication, as referred to in .
#### 6.1.2 Process $`f_1\overline{f}_1ZZ`$
For this process we compared only ‘‘virtual + soft’’ corrections in the conditions of Tables 1,2 of Ref. with input parameters tuned carefully. In this case we do not find a good agreement: SANC numbers happened to lie about 10% lower. We shall go back to searching for the origin of this discrepancy after a tuned comparison of the very similar process $`f_1\overline{f}_1ZA`$ with the results of Ref.. The implementation of the latter process into the SANC system is nearly finished.
### 6.2 Numerical results for real and complete corrections
#### 6.2.1 Hard bremsstrahlung in $`f_1\overline{f}_1HZ`$ annihilation channel
In Table 2 we present typical results of a triple comparison of the Born cross section and the cross section of hard photon bremsstrahlung between two calculations within SANC (semi-analytic, Eq.25, and MC) and those of CompHEP for $`E_\gamma 1`$GeV. Here we used $`m_H`$=130 GeV and the other parameters as in CompHEP.
The two SANC results perfectly agree within statistical errors. The agreement with CompHEP also looks quite good.
#### 6.2.2 Hard bremsstrahlung in $`f_1\overline{f}_1ZZ`$ annihilation channel
In Table 3 we present the results of a double comparison of the Born cross section and the cross section of hard photon bremsstrahlung between the calculation within SANC (MC) and those of CompHEP for $`E_\gamma 1`$GeV. <sup>3</sup><sup>3</sup>3For this channel we don’t have the semi-analytic result. Here we used $`m_H`$=130 GeV and the other parameters as in CompHEP.
Again, we see a very good within larger statistical errors in CompHEP.
#### 6.2.3 $`Hf_1\overline{f}_1Z`$ decay channel
For this decay we present the complete one-loop correction.
We present numbers, collected in the $`G_\mu `$ scheme for the standard SANC INPUT: PDG(2006)
$`\begin{array}{ccccccccccccccc}G_F\hfill & =& 1.1663710^5\mathrm{GeV}^2,\hfill & \alpha (0)\hfill & =& 1/137.03599911,\hfill & \alpha _s(M_Z)\hfill & =& 0.1187,\hfill & & & & & & \\ M_W\hfill & =& 80.403\mathrm{GeV},\hfill & \mathrm{\Gamma }_W\hfill & =& 2.141\mathrm{GeV},\hfill & & & & & & & & & \\ M_Z\hfill & =& 91.1876\mathrm{GeV},\hfill & \mathrm{\Gamma }_Z\hfill & =& 2.4952\mathrm{GeV},\hfill & & & & & & & & & \\ M_H\hfill & =& 120\mathrm{GeV},\hfill & & & & & & & & & & & & \\ m_e\hfill & =& 0.5109989210^3\mathrm{GeV},\hfill & m_u\hfill & =& 62\mathrm{MeV}\hfill & m_d\hfill & =& 83\mathrm{MeV},\hfill & & & & & & \\ m_\mu \hfill & =& 0.105658369\mathrm{GeV},\hfill & m_c\hfill & =& 1.5\mathrm{GeV},\hfill & m_s\hfill & =& 215\mathrm{MeV},\hfill & & & & & & \\ m_\tau \hfill & =& 1.77699\mathrm{GeV},\hfill & m_b\hfill & =& 4.7\mathrm{GeV},\hfill & m_t\hfill & =& 174.2\mathrm{GeV},\hfill & & & & & & \\ \mathrm{\Gamma }_t\hfill & =& 1.551\mathrm{GeV}.\hfill & & & & & & & & & & & & \end{array}`$ (42)
The only exception is the Higgs boson mass for which we again take $`M_H=130\mathrm{GeV}`$ in this section.
Table 4 shows the double and singly differential decay width for the decays $`He^+e^{}Z`$ and $`H\mu ^+\mu ^{}Z`$ for a set of $`s`$ and $`\mathrm{cos}\vartheta _l`$.
From Table 4 it is seen that at the edges of $`\mathrm{cos}\vartheta _l`$ and near the fermionic threshold the double differential width shows a $`1/s`$ behaviour, typical for Coulomb interaction. The origin of the Coulomb peak at the one-loop level may be easily understood. First we note that $`HZ\gamma `$ width does not vanish for an on-shell photon with $`Q_\gamma ^2=0`$, see first Fig. 13:
Therefore, the one-loop amplitude for $`HZf_1\overline{f}_1`$ with virtual photon exchange will show a $`1/s`$ behaviour (with $`s=Q_\gamma ^2`$). This, in turn, will lead to the $`1/s`$ behaviour of both the double and single decay differential widths. This conclusion is fully confirmed by the numbers in Table 4.
Recalling now the limits of $`s,\mathrm{\hspace{0.33em}4}m_f^2s(M_HM_Z)^2`$, one might expect the appearance after integration over $`s`$ of the big logarithm $`\mathrm{ln}((M_HM_Z)^2/m_f^2)`$, with a final state fermion mass singularity. However, the $`1/s`$ region is very narrow and it is largely washed out not only by a soft cut on the variable $`s`$ but even by the plain integration over $`s`$.
Finally let us discuss the total width for muon channel. In Born approximation it is $`\mathrm{\Gamma }^{\mathrm{Born}}=5.59210^6`$ GeV, while with the complete EW corrections it is $`\mathrm{\Gamma }^{\mathrm{Born}+1\mathrm{loop}}=5.77410^6`$ GeV. So, the correction in $`G_\mu `$ scheme amounts to $`3.2\%`$.
#### 6.2.4 Hard photon radiation in $`Hf_1\overline{f}_1Z`$ decay
Here we present the triple comparison of the hard photon bremsstrahlung contribution with a cut on the lepton--photon invariant mass $`\tau =(p_4+p_5)^2`$ again between two calculations of SANC and one of CompHEP.
In Table 5 ‘‘$`\sqrt{\tau _{min}}_{|\overline{\omega }=0.1}`$GeV’’ denotes $`\tau _{min}`$ derived by the formula $`\tau _{min}=m_l^2+2m_lE_{\gamma ,min}`$.
As seen, two SANC numbers agree within MC errors and there is a reasonable agreement with CompHEP everywhere but upper left corner (soft radiation by electrons) where CompHEP show a tendency to be unstable.
## 7 A Monte - Carlo generator for $`H4\mu `$
A Monte Carlo generator of unweighted events for process $`H4\mu `$ is the first example of a possible application of the building blocks ideology of SANC.
In our generator we implement two building blocks at the one-loop level: $`HffZ(\gamma )`$ and $`Zff(\gamma )`$. We merge these two blocks and create a link between them by means of the $`Z^{}`$ (resonating boson) line with Breit-Wigner mass distribution. In this spirit we create the generator in the single resonance approximation. We also built a double resonance generator, where we incorporated resonance approximation in two $`Z^{}`$ lines.
The range for application of the single resonance approximation was found to be $`120`$ GeV$`M_H160`$ GeV and for the double resonance approximation $`M_H`$ 180 GeV.
These conclusions are illustrated by Fig. 15 where we present the results of calculations of the tree level width of the decay $`H4\mu `$ for the three cases:
1) solid line: results of complete tree level calculations neglecting effects of identical final state muons;
2) dash dotted line: single resonance approximation,
$`\mathrm{\Gamma }_{H4\mu }^{1\mathrm{res}}={\displaystyle \frac{\mathrm{\Gamma }_{H2\mu Z}\mathrm{\Gamma }_{Z2\mu }}{\mathrm{\Gamma }_Z}};`$ (43)
3) dashed line: double resonance approximation,
$`\mathrm{\Gamma }_{H4\mu }^{2\mathrm{res}}={\displaystyle \frac{\mathrm{\Gamma }_{H2Z}\mathrm{\Gamma }_{Z2\mu }^2}{\mathrm{\Gamma }_Z^2}}.`$ (44)
The loop-corrected result is the linear combination of three types of events: Born with Identity events minus Resonanse Born events plus One-loop Resonance events. We feel that this is the main shortcoming of the generator.
Let us consider each type of events:
* ‘‘Born with Identity’’ events means a branch which computes the distributions without radiative corrections but with effects of identical muons;
* ‘‘Resonance Born’’ events means resonance approximation for one of the $`Z`$ bosons, i.e. $`HZZ^{}\mu ^+\mu ^{}\mu ^+\mu ^{}`$, where $`Z^{}`$ is the resonating $`Z`$. Here, the two building blocks in Fig. 14 are calculated at the Born level;
* ‘‘Resonance One-loop’’ events is implemented in the same spirit as ‘‘Resonance Born’’, but building blocks $`HffZ(\gamma )`$ and $`Zff(\gamma )`$ are calculated at the one-loop level.
The codes of both MC generators can be obtained from the authors by request.
Recently appeared a new MC code Prophecy4f, see Refs.--, realizing calculation of the complete one-loop corrected partial widths of the $`H4l`$ channels. We present a preliminary comparison between MC Prophecy4f and SANC in Table 6.
As seen from the Table, there is $`\pm `$1% agreement in the mass range 130--140 GeV, degrading at the edges of the interval \[120--160\], that finds its natural explanation in Fig.15. Moreover, Prophecy4f uses the complex-mass scheme and takes into account several higher order corrections. One has to emphasize, however, that SANC calculations in $`\alpha `$ and $`G_\mu `$ schemes differ by about 2%. This can be considered as a rough estimate of the theoretical error. Prophecy4f numbers lie basically inside the range of SANC predictions.
The generator in the single resonance approximation described in this section was used for a MC simulation of $`H4\mu `$ decay in the ATLAS detector and the results were compared with simulation by PYTHIA, showing notable deviations, see . This fact demonstrates the importance of higher order corrections and the necessity to reduce the theoretical error.
## 8 User Guide
### 8.1 Benchmark case 3: the process $`Hf_1\overline{f_1}Z`$
Here we consider the $`2f2b`$ NC process $`Hf_1\overline{f}_1Z`$.
One can open the relevant branch of the SANC tree as follows:
EW $``$ Processes $``$ 4 legs $``$ 2f2b $``$ Neutral Current $``$ H$``$f1f1Z
For this process there are three FORM programs: (FF) Form Factors, (HA) Helicity Amplitudes, and (BR) Bremsstrahlung. Each of them in turn is opened, compiled and run as described in Section 6 of Ref..
For the process $`He^+e^{}Z`$ we have in the Console window the particle indices shown in Table 7.
These can be changed to typeID (typeFU) = 13,14 for up- and down-quarks in the final state of the processes $`H(u\overline{u},d\overline{d})Z`$ by editing the particle numbers as explained in Section 6 of Ref.. <sup>4</sup><sup>4</sup>4See Table 2 for definitions of particle types typeXX..
Next bring the Fortran Editor sheet of the Editors List and the Numeric Form panel to the foreground. Shown in the Numeric Parameter sheet are the particle masses in GeV and the invariant mass of $`f\overline{f}Z`$ compaund in GeV, also the cosine of the angle $`\vartheta _l`$ defined in Fig.4.
Click on the Rehash button at the bottom of the Numeric Form panel: the main module of FORTRAN code appears in the Fortran Editor sheet of the Editors List. Then click on Compile. The final answer appears in the Output field. It consists of the parameters used ($`\alpha `$, $`G_F`$, particle masses, the ’t Hooft scale $`\mu `$ and the invariant mass of compaund), and the resulting differential width $`d^2\mathrm{\Gamma }/dsd\mathrm{cos}\vartheta _l`$ in the Born approximation and Born+one-loop. The results for the default parameters and for several scattering angles are summarised in Table 8.
Input parameters can be changed by editing the appropriate field of the Numeric Form panel and pressing the Rehash button. Again the Rehash button must be pressed before pressing Compile.
To produce whole Table 8 one can set flag tbprint = 1 in the Fortran Editors sheet. After editing the code just press Compile; there is no need to press the Rehash button.
One can also produce the differential width $`d\mathrm{\Gamma }/ds`$ and total width $`\mathrm{\Gamma }`$ in GeV by integrating the above differential width. To produce these numbers one can set flag inflag = 1,2, respectevely, in the Fortran Editors sheet. After editing the code just press Compile; again, there is no need to press the Rehash button.
Acknowledgments
The authors are grateful to P. Christova for a valuable discussion of bremsstrahlung issues, to A. Arbuzov for discussion of the stability of numerical calculations and physical results and also to W. Hollik for providing us with useful references. Three of us (D.B., L.K. and G.N.) are cordially indebted to S. Jadach and Z. Was for offering us an opportunity of encouraging common work at IFJ Krakow in April--May 2005 and to the IFJ directorate for hospitality which was extended to us in this period, when the major part of this study was done. We are thankful to the authors of the Prophecy4f generator for providing us with their numbers for Table 6.
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# Hydrogen atom and one-electron molecular systems in a strong magnetic field: are all of them alike?
## I The $`\mathrm{H}`$-atom
Let us take hydrogen atom - the simplest one-electron system - placed in a constant uniform magnetic field $`B`$ directed along the $`z`$-axis. Due to the Lorentz force, the spherical symmetric electron cloud (in the absence of a magnetic field) is deformed to a cigar-like form. The size of the electron cloud $`r_\mathrm{t}`$ in transversal direction to $`z`$-axis shrinks drastically $`B^{1/2}`$ at large magnetic fields, being close the value of the Larmor radius. As to the longitudinal size $`r_\mathrm{l}`$ it also contracts at large magnetic fields but at a much more moderate rate, $`(\mathrm{log}B)^1`$ (see, e.g., Ruderman:1971 ; LL-QM:1977 ; Hasegawa:1961 ). An interplay of these two types of behavior explains the cigar-type form of the electron cloud. At very large magnetic fields, the cigar-type form evolves to a needle-like form known as the Ruderman needle. In Fig. 1a the form of the electron cloud is illustrated for $`B=10^{12}`$ G <sup>4</sup><sup>4</sup>4Calculations were made using the trial function (7) from Potekhin:2001 . In particular, the longitudinal size of the electron cloud shrinks in comparison with the zero-magnetic-field case about four times. Therefore, the apparent classical (electrostatic) appearance of the magnetic field influence is characterized by a change of the form of the electron cloud, which can be roughly approximated by the ratio of two classical parameters $`r_\mathrm{t},r_\mathrm{l}`$. In fact, it is the major assumption of the present approximation scheme. We also assume that these parameters $`r_\mathrm{t},r_\mathrm{l}`$ are defined by the expectation values,
$$r_\mathrm{t}<\rho >,r_\mathrm{l}2<|z|>.$$
(4)
If a definition of the transversal size $`r_\mathrm{t}`$ looks natural from the physical point of view and rather unambiguous, definition (4) of the longitudinal size is not so obvious. It can be chosen as $`\sqrt{<z^2>}`$, or as a linear combination of $`<|z|>`$ and $`\sqrt{<z^2>}`$. So far it is not so clear what would be physical arguments which allow to specify a definition. Eventually, it turns out it is not very important what quantity is used to define $`r_\mathrm{l}`$. The results of the fit remain very similar although there can be some difference in the parameters.
The binding energy $`E_\mathrm{b}`$ is by definition the difference between the energy of free electron in magnetic field (the Larmor energy) $`B`$ and the total energy of the atom, $`E_\mathrm{b}=BE_\mathrm{T}`$. It is known that $`E_\mathrm{b}`$ in the weak-field regime is represented by the Taylor expansion in powers of $`B^2`$, while for large $`B`$ it behaves $`(\mathrm{log}^2B)`$ (see, for example, LL-QM:1977 and discussion in Karnakov:2003 ). Following the above assumption the binding energy depends on the ratio $`𝒳=r_\mathrm{t}/r_\mathrm{l}`$,
$$E_\mathrm{b}=E_\mathrm{b}(𝒳).$$
(5)
It is quite natural to approximate the transverse size $`r_\mathrm{t}<\rho >`$ as follows
$$r_\mathrm{t}=\frac{r_\mathrm{t}^0}{(1+\alpha _\mathrm{t}^2B^2)^{1/4}}\left(\frac{1+a_\mathrm{t}B^2}{1+b_\mathrm{t}B^2}\right),$$
(6)
where $`r_\mathrm{t}^0,\alpha _\mathrm{t},a_\mathrm{t},b_\mathrm{t}`$ are parameters, which are found by fitting the calculated expectation values for $`<\rho >`$. The formula (6) is written in such a way as to reproduce a functionally-correct perturbative expansion of $`<\rho >`$ at $`B=0`$ (in powers $`B^2`$) and $`r_\mathrm{t}^0=\sqrt{2}a_\mathrm{B}`$, where $`a_\mathrm{B}=1`$ a.u. is the Bohr radius. At large $`B`$ the right-hand side of Eq.(6) behaves as $`B^{1/2}`$ simulating the Larmor radius behavior.
In Fig. 2 one can see that Eq.(6) fits data on $`<\rho >`$ calculated using the formalism developed in Potekhin:2001 with accuracy better than one percent at $`B10^9`$ G. The parameters of the fit are given in Table 1. The parameter $`r_\mathrm{t}^0`$ is also found from the fit, it deviates from $`\sqrt{2}`$ (see above) by $`8\%`$. It reflects the fact that the accuracy provided by formula (6) diminishes as the magnetic field decreases (see the discussion below).
At first sight, it is a much more complicated task to describe the longitudinal size, $`r_\mathrm{l}2<|z|>`$. The approximation we propose to use is
$$r_\mathrm{l}=\frac{r_\mathrm{l}^0}{1+\alpha _\mathrm{l}\mathrm{log}(1+\beta _\mathrm{l}^2B^2+\gamma _\mathrm{l}^2B^4)}\left(\frac{1+a_\mathrm{l}B^2}{1+b_\mathrm{l}B^2}\right),$$
(7)
where $`r_\mathrm{l}^0,\alpha _\mathrm{l},\beta _\mathrm{l},\gamma _\mathrm{l},a_\mathrm{l},b_\mathrm{l}`$ are parameters, which are found by fitting the calculated expectation values for $`2<|z|>`$. Formula (7) has the perturbative expansion in powers $`B^2`$, which agrees with perturbation theory results and $`r_\mathrm{l}^0=3/2a_B`$, where $`a_\mathrm{B}=1`$a.u. is the Bohr radius. At large $`B`$, the right-hand side (7) behaves as $`(\mathrm{log}B)^1`$ as should be in accordance with the qualitative arguments.
In Fig. 3 one can see that (7) fits data on $`2<|z|>`$ obtained in the formalism developed in Potekhin:2001 with accuracy better than 1% at $`B10^9`$ G. The parameters of the fit are given in Table 2. The parameter $`r_\mathrm{l}^0`$ is also found from fit. Surprisingly, it deviates from $`3/2`$ (see above) insignificantly, by $`1\%`$, in contrast to what happened for the parameter $`r_\mathrm{t}^0`$.
In Fig. 4 a comparison of the ratio $`𝒳=r_\mathrm{t}/r_\mathrm{l}`$ (see Eqs. (6)-(7)) with parameters taken from Table 1 with results of calculations is presented. One can clearly see that both data and fitted curves of Figs. 3 and 4 demonstrate a certain irregularity in the range $`(550)\times 10^{10}`$ G. It is a transition region from the Coulomb regime, where the Coulombic forces dominate over magnetic forces to the Landau regime where in the $`(x,y)`$ plane Coulombic forces become subdominant.
Following the assumption (5) let us approximate the binding energy
$$E_b=A𝒳_\mathrm{l}^2+B𝒳_\mathrm{l}+C,𝒳_\mathrm{l}=\mathrm{log}𝒳,$$
(8)
where $`A,B,C`$ are parameters, which are found by making fit of the results of calculations of the binding energy. It is worth emphasizing that the parameters of $`𝒳(𝒳_\mathrm{l})`$ are already fixed by the fits (6), (7) of $`r_\mathrm{t},r_\mathrm{l}`$, respectively. The formula (8) agrees with the perturbative expansion in powers $`B^2`$ (at small $`B`$) and gives a correct asymptotic expansion at large $`B`$.
In Fig. 5 it is shown the fit using the formula (8) of the best known results for the binding energies from Kravchenko:1996 combined with those from Potekhin:2001 . The parameters $`A,B,C`$ are given in Table 3. In the whole range of explored magnetic fields $`10^94.414\times 10^{13}`$ G, formula (8) approximates the binding energies with a relative accuracy, which does not exceed few percent and becomes more accurate with growing magnetic field. It is worth mentioning that when the parameter $`A=4`$ in the approximation (8) its asymptotic coincides with the exact asymptotic (see, e.g., LL-QM:1977 , §112),
$$E_\mathrm{b}\mathrm{log}^2\frac{B}{B_0},B\mathrm{},$$
(9)
where $`E_\mathrm{b}`$ is in Ry. In fact, the deviation $`|A/41|`$ gives a feeling about the quality of our approximation. Clearly, this estimate is very rough ca. 20 % (see Table III), while a real accuracy of approximating the binding energy is a few percent.
One of the important characteristics of the magnetic field influence on the $`\mathrm{H}`$-atom is the appearance of the quadrupole moment
$$QQ_{zz}=2<z^2><\rho ^2>.$$
(10)
Recently, the first quantitative study of the quadrupole moment was carried out Potekhin:2001 . The formula (10) suggests immediately the following approximation
$$Q=2r_\mathrm{l}^2(A_qa_q𝒳_\mathrm{l})r_\mathrm{t}^2(B_q+b_q𝒳_\mathrm{l}),𝒳_\mathrm{l}=\mathrm{log}𝒳,$$
(11)
where $`A_q=0.325447,a_q=0.049432,B_q=1.32012,b_q=0.955362`$ are dimensionless parameters, which are found by fitting the quadrupole moment. The parameters of $`𝒳_\mathrm{l}`$ are already fixed in the fit of parameters $`r_\mathrm{t},r_\mathrm{l}`$ using (6) and (7), respectively. Formula (11) describes correctly the expansion at small and large $`B`$ (see Ruderman:1971 ; Turbiner:1987 ; Potekhin:2001 ). It fits the results of calculations in Potekhin:2001 with an accuracy of few percents (see Fig. 6).
We made an analysis of the expectation values $`<|z|^n>`$ at $`n=2,3,4,5`$. It turns out that the calculated expectation values admit a very accurate polynomial approximation in terms of a single expectation value $`<|z|>`$,
$$<|z|^n>=P_n(<|z|>),$$
(12)
where $`P_n`$ is a $`n`$-th degree polynomial. It seems natural to assume that (12) holds for any $`n`$, hence any expectation value is defined by $`<|z|>`$. This leads to a striking hypothesis that the ground state eigenfunction integrated over $`\rho `$ can be viewed as a one-parametric probability distribution (!).
## II The $`\mathrm{H}_2^+`$ molecular ion
In this Section we consider the molecular ion $`\mathrm{H}_2^+`$ in parallel configuration, when the protons are situated along the magnetic line. The form of the electron cloud is illustrated in Fig. 1b for the magnetic field $`B=10^{12}`$ G. The transversal size of the electron cloud $`r_\mathrm{t}`$ shrinks drastically, $`B^{1/2}`$, at large magnetic fields, being close to the value of the Larmor radius similarly to what happens for the hydrogen atom. As to the longitudinal size $`r_\mathrm{l}`$ it also shrinks but at a much slower rate $`(\mathrm{log}B)^1`$. In particular, the longitudinal size of the electron cloud shrinks in comparison with vanishing magnetic field about five times (see Fig. 1b).
Following the same arguments which were used earlier for the $`\mathrm{H}`$ atom, we again assume that the dynamic characteristics of the $`\mathrm{H}_2^+`$ in a magnetic field depend on the expectation values of transversal $`(r_\mathrm{t})`$ and longitudinal $`(r_\mathrm{l})`$ sizes. The dependence of them on magnetic field is approximated by similar formulas (6)-(7). The binding energy at equilibrium distance between protons depends on the ratio $`𝒳=r_\mathrm{t}/r_\mathrm{l}`$. Eventually, the binding energy is written in the same form (8) with the same expressions (6) and (7) as is done for $`\mathrm{H}`$ atom but with different parameters. For the fit we use the results of recent calculations of the binding energy which were carried out in Turbiner:2003 ; Turbiner:2004 . These parameters of the fit are presented in Tables 1-3 and the fit is illustrated by Figs. 710.
In order to approximate the equilibrium distance $`R_{\mathrm{eq}}`$ we assume that $`R_{\mathrm{eq}}`$ is proportional to the longitudinal distance $`r_\mathrm{l}`$ with a small correction in $`𝒳_\mathrm{l}`$
$$R_{\mathrm{eq}}=r_\mathrm{l}(c_0+c_1𝒳_\mathrm{l}+c_2𝒳_\mathrm{l}^2)=\frac{r_\mathrm{l}^0}{1+\alpha _\mathrm{l}\mathrm{log}(1+\beta _\mathrm{l}B^2)}\left(\frac{1+a_\mathrm{l}B^2}{1+b_\mathrm{l}B^2}\right)(c_0+c_1𝒳_\mathrm{l}+c_2𝒳_\mathrm{l}^2),$$
(13)
where the parameters $`c_0,c_1,c_2`$ are found from the fit of the results of calculations of the equilibrium distance which were carried out in Turbiner:2003 ; Turbiner:2004 . The parameters of the fit are given in Table 4. The fit is illustrated in Fig. 8. It is worth mentioning that the parameters $`c_i,i=0,1,2`$ decrease very fast with $`i`$ (see Table 4). This can be considered as an indication of adequateness of the approximation formula (13).
Similarly to what happened for $`\mathrm{H}`$ atom, the plot of the ratio $`𝒳=r_\mathrm{t}/r_\mathrm{l}`$ (Fig. 9) reveals a certain irregularity in behavior of the calculation results as well as the fit in the range $`(550)\times 10^{10}`$ G. We assign these irregularities to a transition from the Coulomb to the Landau regime. An overall quality of the fit for the domain $`10^94\times 10^{13}`$ G is very high, about 1-2 % except for the above-mentioned region where the accuracy drops to 5-10 %.
Similar to what was done for the $`\mathrm{H}`$-atom we carried out a calculation of expectation values $`<|z|^n>`$, $`n=2,3,4,5`$. It turns out that these expectation values admit very accurate polynomial approximation in terms of the expectation value $`<|z|>`$ (see (12)). It seems natural to assume that (12) holds for any $`n`$. This leads to the hypothesis that the ground state eigenfunction integrated over $`\rho `$ defines a one-parametric distribution similar to what appears for the $`\mathrm{H}`$ atom (see previous Section).
## III The $`\mathrm{H}_3^{++}`$ molecular ion
Now we consider the exotic system $`\mathrm{H}_3^{++}`$ theoretically predicted in Turbiner:1999 , which is made out of three protons situated along the magnetic line and one electron (parallel configuration). This system appears as a quasi-stationary state at $`B10^{10}`$ G Turbiner:2004 . The form of the electron cloud for $`B=10^{12}`$ G is shown in Fig. 1c. It is clearly seen that the transversal size of the electron cloud $`r_\mathrm{t}`$ shrinks drastically, $`B^{1/2}`$ at large magnetic fields similar to what happens for the hydrogen atom and the $`\mathrm{H}_2^+`$ molecular ion which is of the order of the Larmor radius. As to the longitudinal size $`r_\mathrm{l}`$ it also contracts but in much slower rate, $`(\mathrm{log}B)^1`$, at large magnetic fields.
We follow the same idea of approximation as for $`\mathrm{H}`$ and $`\mathrm{H}_2^+`$ assuming that the physics is governed by a single parameter $`𝒳=r_\mathrm{t}/r_\mathrm{l}`$. The same approximation formulas (6) and (7) are used for the transverse ($`r_\mathrm{t}`$) and longitudinal ($`r_\mathrm{l}`$) sizes, respectively, as it is done for $`\mathrm{H}`$-atom and $`\mathrm{H}_2^+`$. Their parameters are found by fitting the results of calculations. The data for $`r_\mathrm{t},r_\mathrm{l}`$ are obtained using a strategy described in Turbiner:2004 . The parameters of the fit are given in Tables 12. The fit of $`r_\mathrm{t}`$ and $`r_\mathrm{l}`$ is illustrated in Figs. 1112. Fig. 13 demonstrates the behavior of the $`𝒳`$. The binding energy $`E_b`$ which is calculated in Turbiner:2004 is approximated using the formula (8) (see Table 3 for parameters of the fit). The fit is illustrated in Fig. 14.
In the same way as it is done for $`\mathrm{H}_2^+`$, we assume that the equilibrium distance between protons are mostly defined by the longitudinal size of the electron cloud (see (7)), which are slightly modified by including the terms depending on $`𝒳_\mathrm{l}=\mathrm{log}𝒳`$. Finally, the equilibrium distance is approximated by Eq.(13) as was done for $`\mathrm{H}_2^+`$ (see Fig. 12). The parameters of the fit are given in Table 4. It is worth mentioning that the parameters $`c_i,i=0,1,2`$ decrease very fast with $`i`$. This might be considered as an indication of adequateness of the approximation formula (13).
In the fit, some irregularities can be seen in the region $`(550)\times 10^{10}`$ G, near the threshold of appearance of the $`\mathrm{H}_3^{++}`$ ion (see Figs. 1114) similarly to those that were observed for the $`\mathrm{H}`$-atom and for the $`\mathrm{H}_2^+`$-ion. One of the reasons for these irregularities can be related to highly increased technical difficulties we encountered exploring this region. This could lead to a loss of accuracy. The overall quality of the fit for the range $`10^{11}4.414\times 10^{13}`$ G is very high, 1-5 %.
Similarly to what was done for the $`\mathrm{H}`$ atom and the $`\mathrm{H}_2^+`$ molecular ion, we calculate the expectation values $`<|z|^n>,n=2,3,4,5`$ for the $`\mathrm{H}_3^{++}`$ ion. It turns out that these expectation values admit a very accurate polynomial approximation in terms of the expectation value $`<|z|>`$, see Eq.(12). It seems natural to assume that (12) holds for any $`n`$. The ground state eigenfunction integrated over $`\rho `$ seems to define a certain one-parametric distribution. A similar phenomenon occurs for the $`\mathrm{H}`$ atom and the $`\mathrm{H}_2^+`$ molecular ion.
## IV The $`(\mathrm{HeH})^{++}`$ molecular ion
Recently, it was theoretically predicted that the exotic molecular ion $`(\mathrm{HeH})^{++}`$ can exist for $`B10^{12}`$ G Turbiner:2004He . Following the same idea of approximation as it was implemented for the $`\mathrm{H}`$ atom and for the $`\mathrm{H}_2^+,\mathrm{H}_3^{++}`$ molecular ions, we can construct high-accuracy approximations for the exotic $`(\mathrm{HeH})^{++}`$ ion. Transversal $`(r_\mathrm{t})`$ and longitudinal $`(r_\mathrm{l})`$ sizes <sup>5</sup><sup>5</sup>5The $`(\mathrm{HeH})^{++}`$ molecular ion is characterized by the asymmetric electronic cloud. Therefore, the longitudinal size is defined $`r_\mathrm{l}<(zz_{\mathrm{max}})>|_{zz_{\mathrm{max}}}<(zz_{\mathrm{max}})>|_{z<z_{\mathrm{max}}}`$ where $`z_{\mathrm{max}}`$ corresponds to the $`z`$-position of the maximum of the electronic distribution. of the electron cloud as a function of the magnetic field are approximated by the expressions (6) and (7) (see Figs. 15 and 16). The parameters of the approximations (6)-(7) obtained through fitting the data from Turbiner:2004He are presented in Tables 1-2. In Fig. 17 the ratio $`𝒳`$ is compared with the calculated data from Turbiner:2004He . The fit of the binding energy was performed using the formula (8) (see Fig. 18). The parameters of the fit are presented in Table 3. For the equilibrium distance $`R_{\mathrm{eq}}`$, the approximation (13) is used (see Fig. 16); the parameters are presented in Table 4. The overall quality of the fit for the range $`10^{12}4.414\times 10^{13}`$ G is very high, around 1 %.
## V $`\mathrm{He}_2^{3+}`$ molecular ion
Recently, it was theoretically predicted that for $`B100`$ a.u. the exotic molecular ion $`\mathrm{He}_2^{3+}`$ can exist Turbiner:2004He . Following the same idea of approximation as for the $`\mathrm{H}`$ atom and the $`\mathrm{H}_2^+,\mathrm{H}_3^{++},(\mathrm{HeH})^{++}`$ molecular ions (see previous Sections), we would like to construct accurate approximations for the exotic $`\mathrm{He}_2^{3+}`$ ion. Transversal $`(r_\mathrm{t})`$ and longitudinal $`(r_\mathrm{l})`$ sizes of the electron cloud as a function of the magnetic field are approximated by the expressions (6) and (7), respectively (see Figs. 19 \- 20). The parameters of the approximations (6) and (7) obtained through the fit of the data obtained in Turbiner:2004He are presented in Tables 12, respectively. In Fig. 21 the ratio $`𝒳`$ is compared with the calculated data from Turbiner:2004He . The fit of the binding energy was performed using the formula (8) (see Fig. 22). The parameters of the fit are presented in Table 3. For the equilibrium distance $`R_{\mathrm{eq}}`$ the approximation (13) is used (see Fig. 20) with parameters presented in Table 4.
Some irregularities can be seen in the fit in the region $`(25)\times 10^{11}`$ G, near the threshold of appearance of the $`\mathrm{He}_2^{3+}`$ ion (see Figs. 1922) similar to those which were observed for the $`\mathrm{H}`$ atom and for the $`\mathrm{H}_2^+,\mathrm{H}_3^{++}`$ ions. The overall quality of the fit for the region $`10^{12}4.414\times 10^{13}`$ G is very high, around 1 %.
## VI Conclusion
In this work we presented a phenomenological model of the behavior of different one-electron atomic-molecular systems in a strong magnetic field. The model is based on a surprisingly simple physical idea that the ground state depends on a ratio of transverse to longitudinal size of a system placed in a strong magnetic field only. Since accurate numerical studies in a strong magnetic field are very tedious from a technical point of view a construction of a phenomenological model which provides approximate expressions for basic characteristics of a system for any value of a magnetic field strength can be quite useful for applications.
One of the motivations of the present work is related to the fact that the neutron star atmosphere is characterized by strong magnetic fields, $`10^{12}10^{13}`$ G. It seems natural to anticipate a wealth of new physical phenomena there. However, for many years the observational data did not indicate anything unusual, corresponding to the black-body radiation. On 2002, the CHANDRA $`X`$-ray observatory collected data on an isolated neutron star 1E1207.4-5209 which led to the discovery of two clearly-seen absorption features at $`0.7`$ keV and $`1.4`$ keV Sanwal:2002 . It is necessary to mention that the XMM-Newton $`X`$-ray observatory recently confirmed the results of Chandra/ACIS related to absorption features at 0.7 and 1.4 keV Bignami:2003 . We proposed a model of hydrogen atmosphere with main abundance of the exotic $`\mathrm{H}_3^{++}`$ molecular ion which explains these absorption features assuming that the surface magnetic field is $`5\times 10^{14}`$ G Turbiner:2004m . For other neutron stars, observational indications of the existence of absorption lines in their spectra were already foundKerkwijk:2004 ; Vink:2004 . It seems natural to anticipate forthcoming observations of other neutron stars which will likely reveal absorption features. The study presented here can be of certain use in identifying possible absorption features.
## VII Acknowledgement
One of us (ABK) is grateful to the Instituto de Ciencias Nucleares, UNAM, where the present work was initiated, for kind hospitality extended to him. This work was supported in part by CONACyT grant 36650-E and DGAPA grant IN124202 (Mexico), and by the RFBR grant 04-02-17263 and a grant of the leading scientific schools 1774.2005.2 (Russia). AVT thanks the University Program FENOMEC (UNAM) for partial financial support.
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# 1 An array of Dirac monopoles aligned along the 𝑥₂ axis with separation 𝛽. In the large 𝐽 limit, within the central slab region -{𝛽/2}≤𝑥₂≤𝛽/2, the magnetic potential goes over from the 3D 1/𝑟 form to a 2D ln(𝜌) form within the region of radius of order 𝐽𝛽.
## Acknowledgments
We thank Alex Kovner and Kimyeong Lee for discussions, and we acknowledge the support of the US DOE through grant DE-FG02-92ER40716. We are very grateful to Paul Sutcliffe for helpful comments and correspondence.
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# 1 Fuzzy spaces
## 1 Fuzzy spaces
### 1.1 Definition and construction of $`_N`$
Fuzzy spaces have been an area of research for a number of years by now . They have proved to be useful in some physical problems. Part of the motivation for this has been the discovery that noncommutative spaces, and more specifically fuzzy spaces, can arise as solutions in string and $`M`$-theories . For example, in the matrix model version of $`M`$-theory, noncommutative spaces can be obtained as $`(N\times N)`$-matrix configurations whose large $`N`$-limit will give smooth manifolds. Fluctuations of branes are described by gauge theories and, with this motivation, there has recently been a large number of papers dealing with gauge theories, and more generally field theories, on such spaces . There is also an earlier line of development, motivated by quantum gravity, using the Dirac operator to characterize the manifold and using ‘spectral actions’ .
Even apart from their string and $`M`$-theory connections, fuzzy spaces are interesting for other reasons. Because these spaces are described by finite dimensional matrices, the number of possible modes for fields on such spaces is limited by the Cayley-Hamilton theorem, and so, one has a natural ultraviolet cutoff. We may think of such field theories as a finite-mode approximation to commutative continuum field theories, providing, in some sense, an alternative to lattice gauge theories. Indeed, this point of view has been pursued in some recent work .
Analysis of fuzzy spaces and particle dynamics on such spaces are also closely related to the quantum Hall effect. The dynamics of charged particles in a magnetic field can be restricted to the lowest Landau level, if the field is sufficiently strong, and this is equivalent to dynamics on a fuzzy version of the underlying spatial manifold. (The fact that the restriction to the lowest Landau level gives noncommutativity of coordinates has been known for a long time; for a recent review focusing on the fuzzy aspects, see .)
The main idea behind fuzzy spaces is the standard correspondence principle of the quantum theory, which is as shown below.
| $`\underset{¯}{QuantumTheory}`$ | $`\mathrm{}0`$ | $`\underset{¯}{ClassicalTheory}`$ |
| --- | --- | --- |
| Hilbert space $``$ | | Phase space $`M`$ |
| | $``$ | |
| Operators on $``$ | | Functions on $`M`$ |
This correspondence suggests a new paradigm. Rather than dealing with theories on a continuous manifold $`M`$, we take the Hilbert space $``$ and the algebra of operators on it as the fundamental quantities and obtain the continuous manifold $`M`$ as an approximation. Generally, instead of $`\mathrm{}`$, we use an arbitrary deformation parameter $`\theta `$, so that the continuous manifold emerges not as the classical limit in the physical sense, but as some other limit when $`\theta 0`$, which will mathematically mimic the transition from quantum mechanics to classical mechanics. The point of view where the spacetime manifold is not fundamental can be particularly satisfying in the context of quantum gravity, and in fact, it was in this context that the first applications of noncommutative spaces to physics was initiated .
Now passing to more specific details, by a fuzzy space, we mean a sequence $`(_N,Mat_N,𝒟_N)`$, where $`_N`$ is an $`N`$-dimensional Hilbert space, $`Mat_N`$ is the matrix algebra of $`N\times N`$-matrices which act on $`_N`$, and $`𝒟_N`$ is a matrix analog of the Dirac operator or, in many instances, just the matrix analog of the Laplacian. The inner product on the matrix algebra is given by $`A,B=\frac{1}{N}\mathrm{Tr}(A^{}B)`$. The Hilbert space $`_N`$ leads to some smooth manifold $`M`$ as $`N\mathrm{}`$. The matrix algebra $`Mat_N`$ approximates to the algebra of functions on $`M`$. The operator $`𝒟_N`$ is needed to recover metrical and other geometrical properties of the manifold $`M`$. For example, information about the dimension of $`M`$ is contained in the growth of the number of eigenvalues. More generally, noncommutative spaces are defined in a similar way, with a triple $`(,𝒜,𝒟)`$, where $``$ can be infinite-dimensional, $`𝒜`$ is the algebra of operators on $``$ and $`𝒟`$ is a Dirac operator on $``$ . For fuzzy spaces the dimensionality of $``$ is finite.
Rather than discuss generalities, we will consider the construction of some noncommutative and fuzzy spaces. Consider the flat $`2k`$-dimensional space $`𝐑^{2k}`$. We build up the coherent state representation buy considering the particle action
$$\frac{𝒮}{\theta }=\frac{i}{\theta }𝑑t\overline{Z}_\alpha \dot{Z}^\alpha $$
(1)
where $`\alpha =1,2,\mathrm{},k`$. Evidently, the time-evolution of the variables $`Z^\alpha `$, $`\overline{Z}_\alpha `$ is trivial, and so, the theory is entirely characterized by the phase space, or upon quantization, by the specification of the Hilbert space. From the action, we can identify the canonical commutation rules as
$$[\overline{Z}_\beta ,Z^\alpha ]=\theta \delta _\beta ^\alpha $$
(2)
It is then possible to choose states, which are eigenstates of $`Z^\alpha `$, defined by $`z|Z^\alpha =z|z^\alpha `$, so that wave functions $`f(z)=z|f`$ can be taken to be holomorphic. The operators $`Z^\alpha ,\overline{Z}_\beta `$ are realized on these by
$$\begin{array}{cc}\hfill Z^\alpha f(z)& =z^\alpha f(z)\hfill \\ \hfill \overline{Z}_\beta f(z)& =\theta \frac{}{z^\beta }f(z)\hfill \end{array}$$
(3)
The inner product for the wave functions should be of the form
$$\begin{array}{cc}\hfill f|h& =𝑑\mu C(z,\overline{z})\overline{f}h\hfill \\ \hfill d\mu & =\underset{\alpha }{}\frac{dz^\alpha d\overline{z}_\alpha }{(2i)}\underset{\alpha }{}d^2z_\alpha \hfill \end{array}$$
(4)
By imposing the adjointness condition $`f|Zh=\overline{Z}f|h`$, we get
$$\theta \frac{C}{z^\alpha }=\overline{z}_\alpha C$$
(5)
which can be solved to yield
$$f|h=\underset{\alpha }{}\frac{d^2z_\alpha }{\pi \theta }\mathrm{exp}\left[\frac{\overline{z}_\alpha z^\alpha }{\theta }\right]\overline{f}h$$
(6)
The overall normalization has been chosen so that the state $`f=1`$ has norm equal to $`1`$.
Since $`f(z)`$ is holomorphic, a basis of states can be given by $`f=1,z^\alpha ,`$ $`z^{\alpha _1}z^{\alpha _2},`$ $`\mathrm{}`$. The Hilbert space is infinite-dimensional and can be used for the noncommutative version of $`𝐑^{2k}`$.
### 1.2 Star products
The star product is very helpful in discussing the large $`N`$ limit. (Star products have along history going back to Moyal and others. The books and reviews quoted, and others, contain expositions of the star product.) We shall consider the two-dimensional case first, generalization to arbitrary even dimensions will be straightforward. A basis for the Hilbert space is given by $`1,z,z^2,`$ etc., and using this basis, we can represent an operator as a matrix $`A_{mn}`$. Associated to such a matrix, we define a function $`A(z,\overline{z})=(A)`$, known as the symbol for $`A`$, by
$`(A)=A(z,\overline{z})`$ $``$ $`{\displaystyle \underset{mn}{}}A_{mn}{\displaystyle \frac{z^m\overline{z}^n}{\sqrt{m!n!}}}e^{z\overline{z}/\theta }`$ (7)
$`=`$ $`{\displaystyle \underset{mn}{}}A_{mn}\psi _m\psi _n^{}`$
where $`\psi _n`$ are given by
$$\psi _n=e^{z\overline{z}/2\theta }\frac{z^n}{\sqrt{n!}}$$
(8)
These are normalized functions obeying the equation
$$\frac{d^2z}{\theta \pi }\psi _n^{}\psi _m=\delta _{mn}$$
(9)
The symbol corresponding to the product of two operators (or matrices) $`A`$ and $`B`$ may be written as
$`(AB)`$ $`=`$ $`{\displaystyle \psi _mA_{mn}B_{nk}\psi _k^{}}`$
$`=`$ $`{\displaystyle \psi _m(z)A_{mn}\left[\frac{d^2w}{\theta \pi }\psi _n^{}(z+w)\psi _r(z+w)\right]B_{rk}\psi _k^{}(z)}`$
$`=`$ $`{\displaystyle \frac{d^2w}{\theta \pi }e^{w\overline{w}/\theta }A(z,\overline{z}+\overline{w})B(z+w,\overline{z})}`$
$``$ $`(A)(B)`$
$`=`$ $`(A)(B)+\theta {\displaystyle \frac{(A)}{\overline{z}}}{\displaystyle \frac{(B)}{z}}+\mathrm{}`$ (11)
Functions on $`M=𝐂`$, under the star product, form an associative but noncommutative algebra. As $`\theta `$ becomes small, we may approximate the star product by the first two terms, giving
$$([A,B])=(A)(B)(B)(A)=\theta \left(\frac{(A)}{\overline{z}}\frac{(B)}{z}\frac{(B)}{\overline{z}}\frac{(A)}{z}\right)$$
(12)
The right hand side is the Poisson bracket of $`A`$ and $`B`$, and this relation is essentially the standard result that the commutators of operators tend to ($`i`$ times) the Poisson bracket of the corresponding functions (symbols) for small values of the deformation parameter. In particular, we find
$$\begin{array}{cc}\hfill Z\overline{Z}& =z\overline{z}+\mathrm{}\hfill \\ \hfill \overline{Z}Z& =\overline{z}z+\theta +\mathrm{}\hfill \\ \hfill Z\overline{Z}\overline{Z}Z& =\theta +\mathrm{}\hfill \end{array}$$
(13)
We can interpret $`Z`$, $`\overline{Z}`$ as the coordinates of the space; they are noncommuting. The noncommutativity is characterized by the parameter $`\theta `$, as we can use the equations given above in terms of symbols to obtain the small $`\theta `$-limit.
These considerations can be generalized in an obvious way to $`M=𝐂^k`$.
### 1.3 Complex projective space $`CP^k`$
We shall now discuss the fuzzy version of $`\mathrm{𝐂𝐏}^k`$. Unlike the case of flat space, we will get a finite number of states for $`\mathrm{𝐂𝐏}^k`$, say, $`N`$, so this will be a truly fuzzy space, rather than just noncommutative. The continuous manifold $`\mathrm{𝐂𝐏}^k`$ can be obtained as $`N\mathrm{}`$. In the previous discussion, we started with continuous $`𝐑^{2k}`$, set up the quantum theory for the action (1), and the resulting Hilbert space could be interpreted as giving the noncommutative version of $`𝐑^{2k}`$. We can follow the same strategy for $`\mathrm{𝐂𝐏}^k`$. In fact, we can adapt the coherent state construction to obtain the fuzzy version of $`\mathrm{𝐂𝐏}^k`$. For a more detailed and group theoretic approach, see .
Continuous $`\mathrm{𝐂𝐏}^k`$ is defined as the set of $`k+1`$ complex variables $`Z^\alpha `$, with the identification of $`Z^\alpha `$ and $`\lambda Z^\alpha `$ where $`\lambda `$ is any nonzero complex number, i.e., we start with $`𝐂^{k+1}`$ and make the identification $`Z^\alpha \lambda Z^\alpha `$, $`\lambda 𝐂\{0\}`$. Based on the fact that there is natural action of $`SU(k+1)`$ on $`Z^\alpha `$ given by
$$Z^\alpha Z^\alpha =g_\beta ^\alpha Z^\beta ,gSU(k+1),$$
(14)
we can show that $`\mathrm{𝐂𝐏}^k`$ can be obtained as the coset
$$\mathrm{𝐂𝐏}^k=\frac{SU(k+1)}{U(k)}$$
(15)
In fact, we may take (15) as the definition of $`\mathrm{𝐂𝐏}^k`$. The division by $`U(k)`$ suggests that we can obtain $`\mathrm{𝐂𝐏}^k`$ by considering a “gauged” version of the action (1), where the gauge group is taken to be $`U(k)`$. Replacing the time-derivative by the covariant derivative, the action becomes
$$𝒮=i𝑑t\overline{Z}_\alpha (_0Z^\alpha iA_0Z^\alpha )n𝑑tA_0$$
(16)
where we have also included a term for the gauge field. (From now on, we will not display $`\theta `$ explicitly.) This action is easily checked to be invariant under the $`U(1)`$ gauge transformation
$$Z^\alpha e^{i\phi }Z^\alpha ,A_0A_0+_0\phi $$
(17)
The pure gauge field part of the action is the one-dimensional Chern-Simons term. The coefficient $`n`$ has to be quantized, following the usual arguments. For example, we can consider the transformation where $`\phi (t)`$ obeys $`\phi (\mathrm{})\phi (\mathrm{})=2\pi `$. The action then changes by $`2\pi n`$, and since $`\mathrm{exp}(i𝒮)`$ has to be single-valued to have a well-defined quantum theory, $`n`$ has to be an integer.
The variation of the action with respect to $`A_0`$ leads to the Gauss law for the theory,
$$\overline{Z}_\alpha Z^\alpha n0$$
(18)
where the weak equality (denoted by $``$) indicates, as usual, that this conditon is to be imposed as a constraint. This is a first class constraint in the Dirac sense, and hence it removes two degrees of freedom. Thus, from $`𝐂^{k+1}`$, we go to a space with $`k`$ complex dimensions. Given the $`U(k)`$ invariance, this can be identified as $`\mathrm{𝐂𝐏}^k`$.
The time-evolution of $`Z^\alpha `$ is again trivial and we are led to the complete characterization of the theory by the Hilbert space, which must be obtained taking account of the constraint (18). In the quantum theory, the allowed physical states must be annihilated by the Gauss law. Using the realization of the $`Z^\alpha `$, $`\overline{Z}_\alpha `$ given in (2), this becomes
$$\left(z^\alpha \frac{}{z^\alpha }n\right)f(z)=0$$
(19)
Thus the allowed functions $`f(z)`$ must have $`n`$ powers of $`z`$’s. They are of the form
$$f(z)=\frac{1}{\sqrt{n!}}z^{\alpha _1}z^{\alpha _2}\mathrm{}z^{\alpha _n}$$
(20)
There are $`N=(n+k)!/n!k!`$ independent functions. The Hilbert space of such functions form the carrier space of a completely symmetric rank $`n`$ irreducible representation of $`SU(k+1)`$. A simple parametrization in terms off local coordinates on $`\mathrm{𝐂𝐏}^k`$ can be obtained by writing $`z^{k+1}=\lambda `$, $`z^i=\lambda \xi ^i`$, where $`\xi ^i=z^i/z^{k+1}=z^i/\lambda `$, for $`i=1,2,\mathrm{},k`$. Correspondingly, the wave functions have the form $`f(z)=\lambda ^nf(\xi )`$. The inner product for two such wave functions can be obtained from the inner product (6). We get
$`f|h`$ $`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle \frac{d^2\lambda }{\pi }\underset{i}{}\frac{d^2\xi _i}{\pi }e^{\lambda \overline{\lambda }(1+\overline{\xi }\xi )}(\lambda \overline{\lambda })^{k+n}\overline{f}h}`$ (21)
$`=`$ $`{\displaystyle \frac{(n+k)!}{n!k!}}{\displaystyle \left[\frac{k!d^2\xi _i}{\pi ^k(1+\overline{\xi }\xi )^{k+1}}\right]\frac{\overline{f(\xi )}}{(1+\overline{\xi }\xi )^{n/2}}\frac{h(\xi )}{(1+\overline{\xi }\xi )^{n/2}}}`$
$`=`$ $`N{\displaystyle 𝑑\mu (\mathrm{𝐂𝐏}^k)\frac{\overline{f(\xi )}}{(1+\overline{\xi }\xi )^{n/2}}\frac{h(\xi )}{(1+\overline{\xi }\xi )^{n/2}}}`$
Here $`N`$ is the dimension of the Hilbert space and $`d\mu (\mathrm{𝐂𝐏}^k)`$ is the standard volume element for $`\mathrm{𝐂𝐏}^k`$ in the local coordinates $`\xi ^i`$, $`\overline{\xi }_i`$. Now, an $`SU(k+1)`$ matrix $`g`$ can be parametrized in such a way that the last column $`g_{k+1}^\alpha `$ is given in terms of $`\xi ^i`$, and the factor $`\sqrt{1+\overline{\xi }\xi }`$, as
$$g=\left[\begin{array}{ccccc}.& .& .& .& \xi ^1\\ .& .& .& .& \xi ^2\\ .& .& .& .& .\\ .& .& .& .& \xi ^k\\ .& .& .& .& 1\end{array}\right]\frac{1}{\sqrt{1+\overline{\xi }\xi }}$$
(22)
The states $`f`$ are thus of the form
$$\frac{f(\xi )}{(1+\overline{\xi }\xi )^{n/2}}=g_{k+1}^{\alpha _1}g_{k+1}^{\alpha _2}\mathrm{}g_{k+1}^{\alpha _n}$$
(23)
Let $`|n,r`$, $`r=1,2,\mathrm{},N`$, denote the states of the rank $`n`$ symmetric representation of $`SU(k+1)`$. Then the Wigner $`𝒟`$-function corresponding to $`g`$ in this representation is defined by $`𝒟_{rs}^{(n)}(g)=n,r|\widehat{g}|n,s`$; it is the matrix representative of the group element $`g`$ in this representation. One can then check easily that
$$g_{k+1}^{\alpha _1}g_{k+1}^{\alpha _2}\mathrm{}g_{k+1}^{\alpha _n}=𝒟_{r,w}^{(n)}=n,r|\widehat{g}|n,w$$
(24)
where the state $`|n,w`$ is the lowest weight state obeying
$$\begin{array}{cc}\hfill T_{k^2+2k}|n,w& =n\frac{k}{\sqrt{2k(k+1)}}|n,w\hfill \\ \hfill T_a|n,w& =0\hfill \end{array}$$
(25)
Here $`T_a`$ are the generators of the $`SU(k)`$ subalgebra and $`T_{k^2+2k}`$ is the generator of the $`U(1)`$ algebra, both for the subgroup $`U(k)`$ of $`SU(k+1)`$. The normalized wave functions for the basis states are thus $`\mathrm{\Psi }_r=\sqrt{N}𝒟_{r,w}^{(n)}(g)`$. Notice that $`𝒟_{n,w}^{(n)}`$ are invariant under right translations of $`g`$ by $`SU(k)`$ transformations, and under the $`U(1)`$ defined by $`T_{k^2+2k}`$ they have a definite charge $`n`$, up to the $`k`$-dependent normalization factor. Since they are not $`U(1)`$ invariant, they are really not functions on $`\mathrm{𝐂𝐏}^k`$, but sections of a line bundle on $`SU(k+1)/U(k)`$, the rank of the bundle being $`n`$. This is exactly what we should expect for quantization of $`\mathrm{𝐂𝐏}^k`$ since this space is given as $`SU(k+1)/U(k)`$.
### 1.4 Star products for fuzzy $`\mathrm{𝐂𝐏}^k`$
As for the flat case, we can construct a star product for functions on $`\mathrm{𝐂𝐏}^k`$ which captures the noncommutative algebra of functions . First we need to establish some notation. Let $`t_A`$ denote the generators of the Lie algebra of $`SU(k+1)`$, realized as $`(k+1\times k+1)`$-matrices. (The $`T`$’s given in equation (25) correspond to the generators of $`U(k)SU(k+1)`$ in the rank $`n`$ symmetric representation; they are the rank $`n`$ representatives of $`t_{k^2+2k}`$ and $`t_a`$. The remaining generators are of two types, $`t_i`$, $`i=1,2,\mathrm{},k`$, which are lowering operators and $`t_{+i}`$ which are raising operators.) Left and right translation operators on $`g`$ are defined by the equations
$$\widehat{L}_Ag=t_Ag,\widehat{R}_Ag=gt_A$$
(26)
If $`g`$ is parametrized by $`\phi ^A`$, some of which are he $`\xi `$’s, then we write
$$g^1dg=it_AE_B^Ad\phi ^B,dgg^1=it_A\stackrel{~}{E}_B^Ad\phi ^B$$
(27)
The operators $`\widehat{L}_A`$ and $`\widehat{R}_A`$ are then realized as differential operators
$$\widehat{L}_A=i(\stackrel{~}{E}^1)_A^B\frac{}{\phi ^B},\widehat{R}_A=i(E^1)_A^B\frac{}{\phi ^B}$$
(28)
The state $`|n,w`$, used in $`𝒟_{rw}^{(n)}(g)=n,r|\widehat{g}|n,w`$, is the lowest weight state, which means that we have the condition
$$\widehat{R}_i𝒟_{r,w}^{(n)}=0$$
(29)
This is essentially a holomorphicity condition. Notice that $`f(\xi )`$ are holomorphic in the $`\xi `$’s; the $`𝒟_{rw}^{(n)}`$ have an additional factor $`(1+\overline{\xi }\xi )^{n/2}`$, which can be interpreted as due to the nonzero connection in $`\widehat{R}_i`$, ultimately due to the nonzero curvature of the bundle. Equation (29) tells us that $`𝒟_{r,w}^{(n)}`$ are sections of a rank $`n`$ holomorphic line bundle.
We define the symbol corresponding to a matrix $`A_{ms}`$ as the function
$`A(g)`$ $`=`$ $`A(\xi ,\overline{\xi })={\displaystyle \underset{ms}{}}𝒟_{m,w}^{(n)}(g)A_{ms}𝒟_{s,w}^{(n)}(g)`$ (30)
$`=`$ $`w|\widehat{g}^T\widehat{A}\widehat{g}^{}|w`$
The symbol corresponding to the product of two matrices $`A`$ and $`B`$ can be simplified as follows.
$`(AB)`$ $`=`$ $`{\displaystyle \underset{r}{}}A_{mr}B_{rs}𝒟_{m,w}^{(n)}(g)𝒟_{s,w}^{(n)}(g)`$ (31)
$`=`$ $`{\displaystyle \underset{rr^{}p}{}}𝒟_{m,w}^{(n)}(g)A_{mr}𝒟_{r,p}^{(n)}(g)𝒟_{r^{},p}^{(n)}(g)B_{r^{}s}𝒟_{s,w}^{(n)}(g)`$
where we use the fact that $`g^{}g^T=1`$, which reads in the rank $`n`$ symmetric representation as $`\delta _{rr^{}}=_p𝒟_{r,p}^{(n)}(g)𝒟_{r^{},p}^{(n)}(g)`$. In the sum over $`p`$ on the right hand side of (31), the term with $`p=n`$ (corresponding to the lowest weight state $`|n,w`$) gives the product of the symbols for $`A`$ and $`B`$. The terms with $`p>n`$ may be written in terms of powers of the raising operators $`R_{+1}`$, $`R_{+2},\mathrm{}`$, $`R_{+k}`$, as
$$𝒟_{r,p}^{(n)}(g)=\left[\frac{(ns)!}{n!i_1!i_2!\mathrm{}i_k!}\right]^{\frac{1}{2}}\widehat{R}_{+1}^{i_1}\widehat{R}_{+2}^{i_2}\mathrm{}\widehat{R}_{+k}^{i_k}𝒟_{r,w}^{(n)}(g).$$
(32)
Here $`s=i_1+i_2+\mathrm{}+i_k`$ and the eigenvalue for the $`U(1)`$ generator $`T_{k^2+2k}`$ for the state $`|n,p`$ is $`(nk+sk+s)/\sqrt{2k(k+1)}`$.
We also get
$$\left[\widehat{R}_{+i}𝒟_{r^{},w}^{(n)}(g)\right]B_{r^{}s}𝒟_{s,w}^{(n)}(g)=\left[\widehat{R}_{+i}𝒟_{r^{},w}^{(n)}B_{r^{}s}𝒟_{s,w}^{(n)}(g)\right]=\widehat{R}_{+i}B(g).$$
(33)
where we used the fact that $`\widehat{R}_{+i}𝒟_{s,w}^{(n)}=0`$. Keeping in mind that $`\widehat{R}_+^{}=\widehat{R}_{}`$, equations (31 -33) combine to give
$`(AB)(g)`$ $`=`$ $`{\displaystyle \underset{s}{}}(1)^s\left[{\displaystyle \frac{(ns)!}{n!s!}}\right]{\displaystyle \underset{i_1+\mathrm{}+i_k=s}{\overset{n}{}}}{\displaystyle \frac{s!}{i_1!i_2!\mathrm{}i_k!}}`$ (34)
$`\times \widehat{R}_1^{i_1}\widehat{R}_2^{i_2}\mathrm{}\widehat{R}_k^{i_k}A(g)\widehat{R}_{+1}^{i_1}\widehat{R}_{+2}^{i_2}\mathrm{}\widehat{R}_{+k}^{i_k}B(g)`$
$``$ $`A(g)B(g)`$
This expression gives the star product for functions on $`\mathrm{𝐂𝐏}^k`$. As expected, the first term of the sum on the right hand side gives the ordinary product $`A(g)B(g)`$, successive terms involve derivatives and are down by powers of $`n`$, as $`n\mathrm{}`$. Since the dimension of the matrices is given by $`N=(n+k)!/n!k!`$, the large $`n`$ is what we need for the limit of the continuous manifold, and the star product, as written here, is suitable for extracting this limit for various quantities. For example, for the symbol corresponding to the commutator of $`A,B`$, we have
$`\left([A,B]\right)(g)`$ $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{k}{}}}(\widehat{R}_iA\widehat{R}_{+i}B\widehat{R}_iB\widehat{R}_{+i}A)+𝒪(1/n^2)`$ (35)
$`=`$ $`{\displaystyle \frac{i}{n}}\{A,B\}+𝒪(1/n^2)`$
The term involving the action of $`\widehat{R}`$’s on the functions can indeed be verified to be the Poisson bracket on $`\mathrm{𝐂𝐏}^k`$. Equation (35) is again the general correspondence of commutators and Poisson brackets, here realized for the specific case of $`\mathrm{𝐂𝐏}^k`$.
We also note that traces of matrices can be converted to phase space integrals. For a single matrix $`A`$, and for the product of two matrices $`A,B`$, we find
$`\mathrm{Tr}A`$ $`=`$ $`{\displaystyle \underset{m}{}}A_{mm}=N{\displaystyle 𝑑\mu (g)𝒟_{m,w}^{(n)}A_{mm^{}}𝒟_{m^{},w}^{(n)}}`$
$`=`$ $`N{\displaystyle 𝑑\mu (g)A(g)}`$
$`\mathrm{Tr}AB`$ $`=`$ $`N{\displaystyle 𝑑\mu (g)A(g)B(g)}`$ (36)
### 1.5 The large $`n`$-limit of matrices
Consider the symbol for the product $`\widehat{T}_B\widehat{A}`$, where $`\widehat{T}_B`$ are the generators of $`SU(k+1)`$, viewed as linear operators on the states. We can simplify it along the following lines.
$`(\widehat{T}_B\widehat{A})_{rs}`$ $`=`$ $`r|\widehat{g}^T\widehat{T}_B\widehat{A}\widehat{g}^{}|s`$ (37)
$`=`$ $`S_{BC}r|\widehat{T}_C\widehat{g}^T\widehat{A}\widehat{g}^{}|s`$
$`=`$ $`S_{Ba}(T_a)_{rp}p|\widehat{g}^T\widehat{A}\widehat{g}^{}|s+S_{B+i}r|\widehat{T}_i\widehat{g}^T\widehat{A}\widehat{g}^{}|s`$
$`+S_{Bk^2+2k}r|\widehat{T}_{k^2+2k}\widehat{g}^T\widehat{A}\widehat{g}^{}|s`$
$`=`$ $`_Br|\widehat{g}^T\widehat{A}\widehat{g}^{}|s`$
$`=`$ $`_BA(g)_{rs}`$
where we have used $`\widehat{g}^T\widehat{T}_B\widehat{g}^{}=S_{BC}\widehat{T}_C`$, $`S_{BC}=2\mathrm{T}\mathrm{r}(g^Tt_Bg^{}t_C)`$. (Here $`t_B,t_C`$ and the trace are in the fundamental representation of $`SU(k+1)`$.) We have also used the fact that the states $`|r`$, $`|s`$ are $`SU(k)`$-invariant. (They are both equal to $`|n,w`$, but we will make this identification only after one more step of simplification.) $`_B`$ is defined as
$$_B=\frac{nk}{\sqrt{2k(k+1)}}S_{Bk^2+2k}+S_{B+i}\widehat{\stackrel{~}{R}}_i$$
(38)
and $`\widehat{\stackrel{~}{R}}_i`$ is a differential operator defined by $`\widehat{\stackrel{~}{R}}_ig^T=T_ig^T`$; it can be written in terms of $`\widehat{R}_i`$ but the precise formula is not needed here.
By taking $`\widehat{A}`$ itself as a product of $`\widehat{T}`$’s, we can iterate this calculation and obtain the symbol for any product of $`\widehat{T}`$’s as
$$(\widehat{T}_A\widehat{T}_B\mathrm{}\widehat{T}_M)=_A_B\mathrm{}_M1.$$
(39)
where we have now set $`|n,r=|n,s=|n,w`$.
A function on fuzzy $`\mathrm{𝐂𝐏}^k`$ is an $`N\times N`$-matrix. It can be written as a linear combination of products of $`\widehat{T}`$’s, and by using the above formula, we can obtain its large $`n`$ limit. When $`n`$ becomes very large, the term that dominates in $`_A`$ is $`S_{Ak^2+2k}`$. We then see that for any matrix function we have the relation, $`F(\widehat{T}_A)F(S_{Ak^2+2k})`$.
We are now in position to define a set of “coordinates” $`X_A`$ by
$$X_A=\frac{1}{\sqrt{C_2(k+1,n)}}T_A$$
(40)
where $`T_A`$ is the matrix corresponding to $`\widehat{T}_A`$ and
$$C_2(k+1,n)=\frac{n^2k^2}{2k(k+1)}+\frac{nk}{2}$$
(41)
is the value of the quadratic Casimir for the symmetric rank $`n`$ representation. The coordinates $`X_A`$ are $`N\times N`$-matrices and can be taken as the coordinates of fuzzy $`\mathrm{𝐂𝐏}^k`$, embedded in $`𝐑^{k^2+2k}`$. In the large $`n`$ limit, we evidently have $`X_AS_{Ak^2+2k}=2\mathrm{T}\mathrm{r}(g^Tt_Ag^{}t_{k^2+2k})`$. From the definition, we can see that $`S_{Ak^2+2k}`$ obey algebraic constraints which can be verified to be the correct ones for describing $`\mathrm{𝐂𝐏}^k`$ as embedded in $`𝐑^{k^2+2k}`$ .
## 2 Noncommutative plane, fuzzy $`\mathrm{𝐂𝐏}^1`$, $`\mathrm{𝐂𝐏}^2`$, etc.
The noncommutative plane has already been described. The basic commutation rules are given by (2), with the indices taking only one value, $`1`$. The star product is given by (1.2). While the coherent state basis is very ideal for considering the commutative limit $`\theta 0`$, for many purposes, it is easy enough to deal with the representation of $`Z,\overline{Z}`$ as infinite-dimensional matrices. In fact, one can also use real coordinates and characterize them by the commutation rules
$$[X_i,X_j]=i\theta ϵ_{ij}$$
(42)
More generally, one may consider $`𝐑^{2k}`$, with the commutation rules
$$[X_i,X_j]=i\theta _{ij}$$
(43)
where the constant matrix $`\theta _{ij}`$ characterizes the noncommutativity.
Fuzzy $`\mathrm{𝐂𝐏}^1`$ is the same as the fuzzy two-sphere and has been studied for a long time . It can be treated as the special case $`k=1`$ of our analysis. The Hilbert space corresponds to representations of $`SU(2)`$, and they are given by the standard angular momentum theory. Representations are labeled by the maximal angular momentum $`j=\frac{n}{2}`$, with $`N=2j+1=n+1`$. The generators are the angular momentum matrices, and the coordinates of fuzzy $`S^2`$ are given by $`X_i=J_i/\sqrt{j(j+1)}`$, as in equation(40). These coordinate matrices obey the commutation rule
$$[X_i,X_j]=\frac{i}{\sqrt{j(j+1)}}ϵ_{ijk}X_k.$$
(44)
We get commuting coordinates only at large $`n`$.
If $`g`$ is an element of $`SU(2)`$ considered as a $`2\times 2`$-matrix, we can parametrize it, apart from an overall $`U(1)`$ factor and along the lines of (22), as
$$g=\frac{1}{\sqrt{(1+\overline{\xi }\xi )}}\left[\begin{array}{cc}1& \xi \\ \overline{\xi }& 1\end{array}\right].$$
(45)
The large $`n`$ limit of the coordinates is given by $`X_iS_{i3}(g)`$, which can be worked out as
$`S_{13}={\displaystyle \frac{\xi +\overline{\xi }}{(1+\xi \overline{\xi })}},`$ $`S_{23}=i{\displaystyle \frac{\xi \overline{\xi }}{(1+\xi \overline{\xi })}},S_{33}={\displaystyle \frac{\xi \overline{\xi }1}{\xi \overline{\xi }+1}}.`$ (46)
The quantities $`S_{i3}`$ obey the condition $`S_{i3}S_{i3}=1`$ corresponding to a unit two-sphere embedded in $`𝐑^3`$; $`\xi ,\overline{\xi }`$ are the local complex coordinates for the sphere. The matrix coordinates obey the condition $`X_iX_i=1`$. Thus we may regard them as giving the fuzzy two-sphere, which approximates to the continuous two-sphere as $`n\mathrm{}`$.
We can also study functions on fuzzy $`S^2`$, which are given as $`N\times N`$-matrices. At the matrix level, there are $`N^2=(n+1)^2`$ independent “functions”. A basis for them is given by $`\mathrm{𝟏}`$, $`X_i`$, $`X_{(i}X_{j)}`$, etc., where $`X_{(i}X_{j)}`$ denotes the symmetric part of the product $`X_iX_j`$ with all contractions of indices $`i,j`$ removed; i.e., $`X_{(i}X_{j)}=\frac{1}{2}(X_iX_j+X_jX_i)\frac{1}{3}\delta _{ij}\mathrm{𝟏}`$. Since we have finite-dimensional matrices, the last independent function corresponds to the symmetric $`n`$-fold product of $`X_i`$’s with all contractions removed.
On the smooth $`S^2`$, a basis for functions is given by the spherical harmonics, labeled by the integer $`l=0,1,2,\mathrm{}`$. They are given by the products of $`S_{i3}`$ with all contractions of indices removed. There are $`(2l+1)`$ such functions for each value of $`l`$. If we consider a truncated set of functions with a maximal value of $`l`$ equal to $`n`$, the number of functions is $`_0^n(2l+1)=(n+1)^2`$. Notice that this number coincides with the number of “functions” at the matrix level. There is one-to-one correspondence with the spherical harmonics, for $`l=0,1,2`$, etc., up to $`l=n`$. Further, by using the relation $`X_iS_{i3}`$, we can see that the matrix functions, $`\mathrm{𝟏}`$, $`X_i`$, $`X_{(i}X_{j)}`$, etc., in the large $`n`$ limit, approximate to the the spherical harmonics. The set of functions at the matrix level go over to the set of functions on the smooth $`S^2`$ as $`n\mathrm{}`$. Fuzzy $`S^2`$ may thus be viewed as a regularized version of the smooth $`S^2`$ where we impose a cut-off on the number of modes of a function; $`n`$ is the regulator or cut-off parameter.
Fuzzy $`\mathrm{𝐂𝐏}^2`$ is the case $`k=2`$ of our general anlysis. The coordinates are given by
$$X_A=\frac{3}{\sqrt{n(n+3)}}T_A$$
(47)
The large $`n`$ limit of the coordinates $`X_A`$ are $`S_{A8}=2\mathrm{T}\mathrm{r}(g^Tt_Ag^{}t_8)`$. In this limit, the coordinates obey the condition
$$\begin{array}{cc}\hfill X_AX_A=& 1\hfill \\ \hfill d_{ABC}X_BX_C=& \frac{1}{\sqrt{3}}X_C\hfill \end{array}$$
(48)
where $`d_{ABC}=2\mathrm{T}\mathrm{r}t_A(t_Bt_C+t_Ct_B)`$. These conditions are known to be the equations for $`\mathrm{𝐂𝐏}^2`$ as embedded in $`𝐑^8`$. Thus, our definition of fuzzy $`\mathrm{𝐂𝐏}^2`$ does approximate to the smooth $`\mathrm{𝐂𝐏}^2`$ in the large $`n`$ limit. Equations (48) can also be imposed at the level of matrices to get a purely matrix-level definition of fuzzy $`\mathrm{𝐂𝐏}^2`$ .
The dimension of the Hilbert space is given by $`N=\frac{1}{2}(n+1)(n+2)`$. Matrix functions are $`N\times N`$-matrices; a basis for them is given by products of the $`T`$’s with up to $`N1`$ factors. There are $`N^2`$ independent functions possible. On the smooth $`\mathrm{𝐂𝐏}^2`$, a basis of functions is given by products of the form $`\overline{u}_{\beta _1}\overline{u}_{\beta _2}\mathrm{}\overline{u}_{\beta _l}u^{\alpha _1}u^{\alpha _2}\mathrm{}u^{\alpha _l}`$, where $`u^\alpha =g_3^\alpha `$. The number of such functions, for a given value of $`l`$, is
$$\left[\frac{1}{2}(l+1)(l+2)\right]^2\left[\frac{1}{2}l(l+1)\right]^2=(l+1)^3$$
(49)
(All traces for these functions correspond to lower ones and can be removed from the counting at the level $`l`$.) If we consider a truncated set of functions, with values of $`l`$ going up to $`n`$, the number of independent functions will be
$$\underset{0}{\overset{n}{}}(l+1)^3=\frac{1}{4}(n+1)^2(n+2)^2=N^2.$$
(50)
It is thus possible to consider the fuzzy $`\mathrm{𝐂𝐏}^2`$ as a regularization of the smooth $`\mathrm{𝐂𝐏}^2`$ with a cut-off on the number of modes of a function. Since any matrix function can be written as a sum of products of $`\widehat{T}`$’s, the corresponding large $`n`$ limit has a sum of products of $`S_{A8}`$’s. The independent basis functions are thus given by representations of $`SU(3)`$ obtained from reducing symmetric products of the adjoint representation with itself. These are exactly what we expect based on the fact the smooth $`\mathrm{𝐂𝐏}^2`$ is given by the embedding conditions (48). The algebra of matrix functions for the fuzzy $`\mathrm{𝐂𝐏}^2`$, as we have constructed it, does go over to the algebra of functions on the smooth $`\mathrm{𝐂𝐏}^2`$.
Since the fuzzy spaces, the fuzzy $`\mathrm{𝐂𝐏}^k`$ in particular, can be regarded as a regularization of the smooth $`\mathrm{𝐂𝐏}^k`$ with a cut-off on the number of modes of a function, they can be used for regularization of field theories, in much the same way that lattice regularization of field theories is carried out. There are some interesting features or fuzzy regularization; for example, it may be possible to evade fermion doubling problem on the lattice .
## 3 Fields on fuzzy spaces, Schrödinger equation
A scalar field on a fuzzy space can be written as $`\mathrm{\Phi }(X)`$, indicating that it is a function of the coordinate matrices $`X_A`$. Thus $`\mathrm{\Phi }`$ is an $`N\times N`$-matrix. Further, equation (35) tells us that
$`[T_A,\mathrm{\Phi }]`$ $``$ $`{\displaystyle \frac{i}{n}}{\displaystyle \frac{nk}{\sqrt{2k(k+1)}}}\{S_{Ak^2+2k},\mathrm{\Phi }\}`$ (51)
$``$ $`iD_A\mathrm{\Phi }.`$
$`D_A`$, as defined by this equation, are the derivative operators on the space of interest. For example, for the fuzzy $`S^2`$, they are given by
$$\begin{array}{cc}\hfill D_1& =\frac{1}{2}(\overline{\xi }^2_{\overline{\xi }}+_\xi \xi ^2_\xi _{\overline{\xi }})\hfill \\ \hfill D_2& =\frac{i}{2}(\overline{\xi }^2_{\overline{\xi }}+_\xi +\xi ^2_\xi +_{\overline{\xi }})\hfill \\ \hfill D_3& =\overline{\xi }_{\overline{\xi }}\xi _\xi \hfill \end{array}$$
(52)
These obey the $`SU(2)`$ algebra, $`[D_A,D_B]=iϵ_{ABC}D_C`$. They are the translation operators on the two-sphere and correspond to the three isometry transformations.
These equations show that we can define the derivative, at the matrix level, as the commutator $`i[T_A,\mathrm{\Phi }]`$, which is the adjoint action of $`T_A`$ on $`\mathrm{\Phi }`$. The Laplacian on $`\mathrm{\Phi }`$ is then given by $`\mathrm{\Delta }\mathrm{\Phi }=[T_A,[T_A,\mathrm{\Phi }]]`$. The Euclidean action for a scalar field can be taken as
$$𝒮=\frac{1}{N}\mathrm{Tr}\left[\mathrm{\Phi }^{}[T_A,[T_A,\mathrm{\Phi }]]+V(\mathrm{\Phi })\right]$$
(53)
where $`V(\mathrm{\Phi })`$ is a potential term; it does not involve derivatives.
The identification of derivatives also lead naturally to gauge fields. We introduce a gauge field $`𝒜_A`$ by defining the covariant derivative as
$$i𝒟_A\mathrm{\Phi }=[T_A,\mathrm{\Phi }]+𝒜_A\mathrm{\Phi }$$
(54)
where $`𝒜_A`$ is a set of hermitian matrices. In the absence of the gauge field, we have the commutation rules $`[T_A,T_B]=if_{ABC}T_C`$. The field strength tensor $`_{AB}`$, which is the deviation from this algebra, is thus given by
$$i_{AB}=[T_A+𝒜_A,T_B+𝒜_B]if_{ABC}(T_C+𝒜_C).$$
(55)
The action for Yang-Mills theory on a fuzzy space is then given by
$$𝒮=\frac{1}{N}\mathrm{Tr}\left[\frac{1}{4}_{AB}_{AB}\right]$$
(56)
The quantum theory of these fields can be defined by the functional integral over actions such as (53) and (56). Perturbation theory, Feynman diagrams, etc., can be worked out. Our main focus will be on particle dynamics, so we will not do this here. However, some of the relevant literature can be traced from .
One can also write down the Schrödinger equation for particle quantum mechanics on a fuzzy space . The wave function $`\mathrm{\Psi }(X)`$ is matrix and its derivative is given by $`iD_A\mathrm{\Psi }=[T_A,\mathrm{\Psi }]`$. Coupling to an external potential may be taken to be of the form $`V(X)\mathrm{\Psi }`$. The Schrödinger equation is then given by
$$i\frac{\mathrm{\Psi }}{t}+D_A(D_A\mathrm{\Psi })V\mathrm{\Psi }=0$$
(57)
When it comes to gauge fields, there is a slight subtlety. The covariant derivative is of the form (54). To distinguish the action of the gauge field from the potential $`V`$, for the covariant derivative we may use the definition $`i𝒟_A\mathrm{\Psi }=[T_A,\mathrm{\Psi }]+\mathrm{\Psi }𝒜_A`$. (One could also change the action of the potential.) The Schródinger equation retains the usual form,
$$i𝒟_0\mathrm{\Psi }+\frac{1}{2m}𝒟_A(𝒟_A\mathrm{\Psi })V\mathrm{\Psi }=0$$
(58)
## 4 The Landau problem on $`𝐑_{NC}^2`$ and $`𝐒_F^2`$
As a simple example of the application of the ideas given above, we shall now work out the quantum mechanics of a charged particle in a magnetic field on the fuzzy two-plane . This is the fuzzy version of the classic Landau problem. We shall also include an oscillator potential to include the case of an ordinary potential as well. At the operator level, the inclusion of a background magnetic field is easily achieved by changing the commutation rules for the momenta. The modified algebra of observables is given by
$$\begin{array}{cc}\hfill [X_1,X_2]& =i\theta \hfill \\ \hfill [X_i,P_j]& =i\delta _{ij}\hfill \\ \hfill [P_1,P_2]& =iB\hfill \end{array}$$
(59)
where $`i,j=1,2`$, and $`B`$ is the magnetic field. The Hamiltonian may be taken as
$$H=\frac{1}{2}\left[P_1^2+P_2^2+\omega ^2(X_1^2+X_2^2)\right]$$
(60)
We have chosen the isotropic oscillator (with frequency $`\omega `$), and $`H`$ is invariant under rotations. The form of various operators can be slightly different from the usual ones because of the noncommutativity of the coordinates. The angular momentum is given by
$$L=\frac{1}{1\theta B}\left[X_1P_2X_2P_1+\frac{B}{2}(X_1^2+X_2^2)+\frac{\theta }{2}(P_1^2+P_2^2)\right]$$
(61)
$`L`$ commutes with $`H`$, as can be checked easily.
The strategy for solving this problem involves expressing $`X_i,P_i`$ in terms of a a usual canonical set, so that thereafter, it can be treated as an ordinary quantum mechanical system. This change of variables will be different for $`B<1/\theta `$ and for $`B>1/\theta `$. For $`B<1/\theta `$, we define a change of variables
$$\begin{array}{cc}\hfill X_1=l\alpha _1,& P_1=\frac{1}{l}\beta _1+q\alpha _2\hfill \\ \hfill X_2=l\beta _1,& P_2=\frac{1}{l}\alpha _1q\beta _2\hfill \end{array}$$
(62)
where $`l^2=\theta `$ and $`q^2=(1B\theta )/\theta `$. $`\alpha _i,\beta _i`$ form a standard set of canonical variables, with
$$\begin{array}{cc}\hfill [\alpha _i,\alpha _j]& =0\hfill \\ \hfill [\alpha _i,\beta _j]& =i\delta _{ij}\hfill \\ \hfill [\beta _i,\beta _j]& =0\hfill \end{array}$$
(63)
The Hamiltonian is now given by
$$H=\frac{1}{2}\left[\left(\omega ^2l^2+\frac{1}{l^2}\right)(\alpha _1^2+\beta _1^2)+q^2(\alpha _2^2+\beta _2^2)+\frac{2q}{l}(\alpha _1\beta _2+\alpha _2\beta _1)\right]$$
(64)
We can eliminate the mixing of the two sets of variables in the last term, and diagonalize $`H`$, by making a Bogoliubov transformation which will express $`\alpha _i,\beta _i`$ in terms of a canonical set $`q_i,p_i`$ as
$$\left(\begin{array}{c}\alpha _1\\ \alpha _2\\ \beta _1\\ \beta _2\end{array}\right)=\mathrm{cosh}\lambda \left(\begin{array}{c}q_1\\ q_2\\ p_1\\ p_2\end{array}\right)+\mathrm{sinh}\lambda \left(\begin{array}{c}p_2\\ p_1\\ q_2\\ q_1\end{array}\right)$$
(65)
The Hamiltonian can be diagonalized by the choice
$$\mathrm{tanh}2\lambda =\frac{2ql}{1+\omega ^2l^4+q^2l^2}$$
(66)
and is given by
$$H=\frac{1}{2}\left[\mathrm{\Omega }_+(p_1^2+q_1^2)+\mathrm{\Omega }_{}(p_2^2+q_2^2)\right]$$
(67)
with
$$\mathrm{\Omega }_\pm =\frac{1}{2}\sqrt{(\omega ^2\theta B)^2+4\omega ^2}\pm \frac{1}{2}(\omega ^2\theta +B)$$
(68)
From equation (67), we see that the problem is equivalent to that of two harmonic oscillators with frequencies $`\mathrm{\Omega }_+`$ and $`\mathrm{\Omega }_{}`$.
The case of $`B>1/\theta `$ can be treated in a similar way. We again have two oscillators, with the Hamiltonian (67), but with frequencies given as
$$\mathrm{\Omega }_\pm =\pm \frac{1}{2}\sqrt{(\omega ^2\theta B)^2+4\omega ^2}+\frac{1}{2}(\omega ^2\theta +B)$$
(69)
Notice that $`B\theta =1`$ is a special value, for both regions of $`B`$, with one of the frequencies becoming zero. The symplectic two-form which leads to the commutation rules (59) is given by
$$\mathrm{\Omega }=\frac{1}{1B\theta }\left(dP_1dX_1+dP_2dX_2+\theta dP_1dP_2+BdX_1dX_2\right)$$
(70)
The phase space volume is given by
$$d\mu =\frac{1}{|1B\theta |}d^2Xd^2P$$
(71)
A semiclassical estimate of the number of states is given by the volume divided by $`(2\pi )^2`$. The formula (71) shows that the density of states diverges at $`B\theta =1`$, again indicating that it is a special value.
The Landau problem on the fuzzy sphere can be formulated in a similar way. On the sphere, the translation operators are the angular momenta $`J_i`$ and the algebra of observables is given by
$$\begin{array}{cc}\hfill [X_i,X_j]& =0\hfill \\ \hfill [J_i,X_j]& =iϵ_{ijk}X_k\hfill \\ \hfill [J_i,J_j]& =iϵ_{ijk}J_k\hfill \end{array}$$
(72)
$`X^2`$ commutes with all operators and its value can be fixed to be $`a^2`$, where $`a`$ is the radius of the sphere. The other Casimir operator is $`XJ`$; its value is written as $`a(n/2)`$, where $`n`$ must be an integer and gives the strength of the magnetic field; it is the charge of the monopole at the center of the sphere (if we think of it as being embedded in $`𝐑^3`$).
For the fuzzy case, the coordinates themselves are noncommuting and are given, up to normalization, by $`SU(2)`$ operators $`R_i`$ as $`X_i=aR_i/\sqrt{C_2}`$. The algebra of observables becomes
$$\begin{array}{cc}\hfill [R_i,R_j]& =iϵ_{ijk}R_k\hfill \\ \hfill [J_i,R_j]& =iϵ_{ijk}R_k\hfill \\ \hfill [J_i,J_j]& =iϵ_{ijk}J_k\hfill \end{array}$$
(73)
The Casimir operators are now $`R^2`$ and $`RJ\frac{1}{2}J^2`$, the latter being related to the strength of the magnetic field. The algebra (73) can be realized by two independent $`SU(2)`$ algebras $`\{R_i\}`$ and $`\{K_i\}`$, with $`J_i=R_i+K_i`$. The two Casimirs are now $`R^2`$ and $`K^2`$, which we fix to the values $`r(r+1)`$ and $`k(k+1)`$, $`r`$, $`k`$ being positive half-integers. The difference $`kr=n/2`$. The limit of the smooth sphere is thus obtained by taking $`k,r\mathrm{}`$, with $`kr`$ fixed. As the generalization of $`P^2/2m`$, we take the Hamiltonian as
$$H=\frac{\gamma }{2a^2}J^2$$
(74)
where $`\gamma `$ is some constant. The spectrum of the Hamilonian is now easily calculated as
$$E=\frac{\gamma }{2a^2}j(j+1),j=\frac{|n|}{2},\frac{|n|}{2}+1,\mathrm{},j+k$$
(75)
From the commutation rule for the coordinates $`X_i=aR_i/\sqrt{r(r+1)}`$, we may identify the noncommutativity parameter as $`\theta a^2/r`$, for large $`r`$. The limit of this problem to the noncommutative plane can be obtained by taking $`r`$ large, but keeping $`\theta `$ fixed. Naturally, this will require a large radius for the sphere. The strength of the magnetic field in the plane is related to $`n`$ by $`(1B\theta )n=2Ba^2`$. For more details, see ; also the Landau problem on general noncommutative Riemann surfaces has been analyzed, see .
## 5 Lowest Landau level and fuzzy spaces
There is an interesting connection between the Landau problem on a smooth manifold $`M`$ and the construction of the fuzzy version of $`M`$; we shall explain this now.
The splitting of Landau levels is controlled by the magnetic field and, if the field is sufficiently strong, transitions between levels are suppressed and the dynamics is restricted to one level, say, the lowest. The observables are given as hermitian operators on this subspace of the Hilbert space corresponding to the lowest Landau level; they can be obtained by projecting the full operators to this subspace. The commutation rules can change due to this projection. The position coordinates, for example, when projected to the lowest Landau level (or any other level), are no longer mutually commuting. The dynamics restricted to the lowest Landau level is thus dynamics on a noncommutative space. In fact, the Hilbert (sub)space of the lowest Landau level can be taken as the Hilbert space $`_N`$ used to define the fuzzy version of $`M`$. Thus the solution of the Landau problem on smooth $`M`$ gives a construction of the fuzzy version of $`M`$.
We can see how this is realized explicitly by analyzing the two-sphere . Since $`S^2=SU(2)/U(1)`$, the wave functions can be obtained in terms of functions on the group $`SU(2)`$, i.e., in terms of the Wigner functions $`𝒟_{rs}^{(j)}(g)`$. We need two derivative operators which can be taken as two of the right translations of $`g`$, say, $`R_\pm =R_1\pm iR_2`$. With the correct dimensions, the covariant derivatives can be written as
$$D_\pm =i\frac{R_\pm }{a}$$
(76)
The $`SU(2)`$ commutation rule $`[R_+,R_{}]=2R_3`$ shows that the covariant derivatives do not commute and we may identify the value of $`R_3`$ as the field strength. In fact, comparing this commutation rule to $`[D_+,D_{}]=2B`$, we see that $`R_3`$ should be taken to be $`(n/2)`$, where $`n`$ is the monopole number, $`n=2Ba^2`$. Thus the wave functions on $`S^2`$ with the magnetic field background are of the form $`\mathrm{\Psi }_m𝒟_{m,\frac{n}{2}}^{(j)}(g)`$.
The one-particle Hamiltonian is given by
$$H=\frac{1}{4\mu }\left(D_+D_{}+D_{}D_+\right)=\frac{1}{2\mu a^2}\left(\underset{A=1}{\overset{3}{}}R_A^2R_3^2\right)$$
(77)
where $`\mu `$ is the particle mass. The eigenvalue $`\frac{n}{2}`$ must occur as one of the possible values for $`R_3`$, so that we can form $`𝒟_{m,\frac{n}{2}}^{(j)}(g)`$. This means that $`j`$ should be of the form $`j=\frac{|n|}{2}+q`$, $`q=0,1,..`$. Since $`R^2=j(j+1)`$, the energy eigenvalues are easily obtained as
$$E_q=\frac{1}{2\mu a^2}\left[\left(\frac{n}{2}+q\right)\left(\frac{n}{2}+q+1\right)\frac{n^2}{4}\right]=\frac{B}{2\mu }(2q+1)+\frac{q(q+1)}{2\mu a^2}$$
(78)
The integer $`q`$ is the Landau level index, $`q=0`$ being the lowest energy state or the ground state. The gap between levels increases as $`B`$ increases, and, in the limit of large magnetic fields, it is meaningful to restrict dynamics to one level, say the lowest, if the available excitation energies are small compared to $`B/2\mu `$. In this case, $`j=\frac{|n|}{2}`$, $`R_3=\frac{n}{2}`$, so that we have the lowest weight state for the right action of $`SU(2)`$, taking $`n`$ to be positive. The condition for the lowest Landau level is $`R_{}\mathrm{\Psi }=0`$.
The Hilbert space of the lowest Landau level is spanned by $`\mathrm{\Psi }_m𝒟_{m,\frac{n}{2}}^{(\frac{n}{2})}`$. Notice that this is exactly the Hilbert space for fuzzy $`S^2`$. Hence all observables for the lowest Landau level correspond to the observables of the fuzzy $`S^2`$.
This correspondence can be extended to the Landau problem on other spaces, say, $`\mathrm{𝐂𝐏}^k`$ with a $`U(1)`$ background field, for example. The background field specifies the choice of the representation of $`T_a`$ and $`T_{k^+2k}`$, the $`U(k)`$ subalgebra of $`SU(k+1)`$, in the Wigner $`𝒟`$-functions. For zero $`SU(k)`$ background field, $`T_a\mathrm{\Psi }=0`$ and the eigenvalue of $`T_{k^2+2k}`$ gives the magnetic field, which must obey appropriate quantization conditions. In fact, we get the equations (25) and the Hilbert subspace of the lowest Landau level is the same as the Hilbert space $`_N`$ used for the construction of fuzzy $`\mathrm{𝐂𝐏}^k`$ .
The lowest Landau level wave functions are holomorphic, except possibly for a common prefactor, which has to do with the inner product. This is also seen from (20). In fact, the condition (29), namely, $`R_i\mathrm{\Psi }=0`$, which selects the lowest level, are the holomorphicity conditions. The higher levels are not necessarily holomorphic. This will be useful later in writing the Yang-Mills amplitudes in terms of a Landau problem on $`\mathrm{𝐂𝐏}^1=S^2`$.
## 6 Twistors, supertwistors
### 6.1 The basic idea of twistors
The idea of twistors is due to Roger Penrose, many years ago, in 1967 . There are many related ways of thinking about twistors, but a simple approach is in terms of constructing solutions to massless wave equations or their Euclidean counterparts.
We start by considering the two-dimensional Laplace equation, which may be written in complex coordinates as
$$\overline{}f=0$$
(79)
where $`z=x_1+ix_2`$. The solution is then obvious, $`f(x)=h(z)+g(\overline{z})`$, where $`h(z)`$ is a holomorphic function of $`z`$ and $`g(\overline{z})`$ is antiholomorphic. For a given physical problem such as electrostatics or two-dimensional hydrodynamics, we then have to simply guess the holomorphic function with the required singularity structure. Further, the problem has conformal invariance and one can use the techniques of conforml mapping to simplify the problem.
We now ask the question: Can we do an analogous trick to find solutions of the four-dimensional problem, say, the Dirac or Laplace equations on $`S^4`$? Clearly this is not so simple as in two dimensions, there are some complications. First of all, $`S^4`$ does not admit a complex structure. Even if we consider $`𝐑^4`$, which is topologically equivalent to $`S^4`$ with a point removed, there is no natural choice of complex coordinates. We can, for example, combine the four coordinates into two complex ones as in
$$\left(\begin{array}{c}z_1\\ z_2\end{array}\right)=\left(\begin{array}{c}x_1+ix_2\\ x_3+ix_4\end{array}\right)$$
(80)
Equally well we could have considered
$$\left(\begin{array}{c}z_1^{}\\ z_2^{}\end{array}\right)=\left(\begin{array}{c}x_1+ix_3\\ x_2+ix_4\end{array}\right)$$
(81)
or, in fact, an infinity of other choices. Notice that any particular choice will destroy the overall $`O(4)`$-symmetry of the problem. We may now ask: How many inequivalent choices can be made, subject to, say, preserving $`x^2=\overline{z}_1z_1+\overline{z}_2z_2`$? Given one choice, as in (79), we can do an $`O(4)`$ rotation of $`x_\mu `$ which will generate other possible complex combinations with the same value of $`x^2`$. However, if we do a $`U(2)`$-transformation of $`(z_1,z_2)`$, this gives us a new combination of the $`z`$’s preserving holomorphicity. In particular, a holomorphic function of the $`z_i`$ will remain a holomorphic function after a $`U(2)`$ rotation. Thus the number of inequivalent choices of local complex structure is given by $`O(4)/U(2)=S^2=\mathrm{𝐂𝐏}^1`$. The idea now is to consider $`S^4`$ with the set of all possible local complex structures at each point, in other words, a $`\mathrm{𝐂𝐏}^1`$ bundle over $`S^4`$. This bundle is $`\mathrm{𝐂𝐏}^3`$. The case of $`𝐑^4`$ is similar to considering a neighborhood of $`S^4`$.
An explicit realization of this is as follows. We represent $`\mathrm{𝐂𝐏}^1`$ by a two-spinor $`U^A`$, $`A=1,2`$, with the identification $`u^a\lambda U^A`$, where $`\lambda 𝐂\{0\}`$. We now take a $`4`$-spinor with complex element $`Z^\alpha `$, $`\alpha =1,2,3,4`$, and write it as $`Z^\alpha =(W_{\dot{A}},U^A)`$, where $`U^A`$ describes $`\mathrm{𝐂𝐏}^1`$ as above. The relation between $`W_{\dot{A}}`$ and $`U^A`$ is taken as
$$W_{\dot{A}}=x_{\dot{A}A}U^A$$
(82)
$`x_{\dot{A}A}`$ defined by this equation may be taken as the local coordinates on $`S^4`$. We can even write this out as
$$x_{\dot{A}A}=\left(\begin{array}{cc}x_4+ix_3& x_2+ix_1\\ x_2+ix_1& x_4ix_3\end{array}\right)=x_\mu e^\mu $$
(83)
with $`e^i=\sigma ^i`$ are the Pauli matrices and $`e^4=\mathrm{𝟏}`$, so that $`x_\mu `$ are the usual coordinates. One can read the equation (82) in another way, namely, as combining the $`x`$’s into complex combinations $`W_1`$ and $`W_2`$, in a manner specified by the choice of $`U^A`$. Thus a point on $`\mathrm{𝐂𝐏}^1`$, namely, a choice of $`U^A`$, gives a specific combination of complex coordinates.
We have the identification $`Z^\alpha \lambda Z^\alpha `$, which follows from $`U^A\lambda U^A`$ and the definition of $`x_{\dot{A}A}`$ as in (82). This means that $`Z^\alpha `$ define $`\mathrm{𝐂𝐏}^3`$. Further the indices $`\dot{A},A`$ correspond to $`SU(2)`$ spinor indices, right and left, in the splitting $`O(4)SU_L(2)\times SU_R(2)`$. $`Z^\alpha `$ are called twistors.
Given the above-described structure, there is a way of constructing solutions to massless wave equations (or their Euclidean versions), in terms of holomorphic functions defined on a neighborhood of $`\mathrm{𝐂𝐏}^3`$. Evidently, preserving the $`O(4)`$ symmetry requires some sort of integration over all $`u`$’s, consistent with holomorphicity. There is a unique holomorphic differential we can make out of the $`u`$’s which is $`O(4)`$ invariant, namely, $`UdU=ϵ_{AB}U^AdU^B`$. We will now do a contour integration of holomorphic functions using this. Let $`f(Z)`$ be aholomorphic function of $`Z^\alpha `$ defined on some region in twistor space. We can then construct the contour integral
$$\stackrel{~}{f}^{A_1A_2\mathrm{}A_n}(x)=_CU𝑑UU^{A_1}U^{A_2}\mathrm{}U^{A_n}f(Z)$$
(84)
For this to make sense on a neighborhood of $`\mathrm{𝐂𝐏}^3`$, $`f(Z)`$ should have degree of homogeneity $`n2`$, so that the integrand is invariant under the scaling $`Z^\alpha \lambda Z^\alpha `$, $`U^A\lambda U^A`$, and thus projects down to a proper differential on $`\mathrm{𝐂𝐏}^3`$. The contour $`C`$ will be taken to enclose some of the poles of the function $`f(Z)`$. Since we write $`W_{\dot{A}}=x_{\dot{A}A}U^A`$, after integration, we are left with a function of the $`x`$’s; $`\stackrel{~}{f}`$ is a function of the $`S^4`$ or $`𝐑^4`$ coordinates; it is also a multispinor of $`SU_L(2)`$.
Consider now the action of the chiral Dirac operator on this, namely, $`ϵ_{CA_1}^{\dot{B}C}\stackrel{~}{f}^{A_1A_2\mathrm{}A_n}`$. Since $`x_\mu `$ appear in $`f(Z)`$ only via the combination $`x_{\dot{A}A}U^A`$, we can write
$`ϵ_{CA_1}^{\dot{B}C}\stackrel{~}{f}^{A_1A_2\mathrm{}A_n}`$ $`=`$ $`ϵ_{CA_1}{\displaystyle _C}U𝑑UU^{A_1}U^{A_2}\mathrm{}U^{A_n}^{\dot{B}C}f(Z)`$ (85)
$`=`$ $`ϵ_{CA_1}{\displaystyle _C}U𝑑UU^{A_1}U^{A_2}\mathrm{}U^{A_n}U^C{\displaystyle \frac{f(Z)}{W_{\dot{B}}}}`$
$`=`$ $`0`$
since $`ϵ_{CA_1}U^CU^{A_1}=0`$ by antisymmetry. Thus $`\stackrel{~}{f}^{A_1A_2\mathrm{}A_n}(x)`$ is a solution to the chiral Dirac equation in four dimensions.
In a similar way, one can define
$$\stackrel{~}{g}^{\dot{A}_1\dot{A}_2\mathrm{}\dot{A}_n}(x)=_CU𝑑U\frac{}{W_{\dot{A}_1}}\frac{}{W_{\dot{A}_1}}\mathrm{}\frac{}{W_{\dot{A}_1}}g(Z)$$
(86)
where $`g(Z)`$ has degree of homogeneity equal to $`n2`$. It is then easy to check that
$$ϵ_{\dot{B}\dot{A}_1}^{B\dot{B}}\stackrel{~}{g}^{\dot{A}_1\dot{A}_2\mathrm{}\dot{A}_n}=0$$
(87)
The two sets of functions, $`\stackrel{~}{f}^{A_1A_2\mathrm{}A_n}(x)`$ and $`\stackrel{~}{g}^{\dot{A}_1\dot{A}_2\mathrm{}\dot{A}_n}(x)`$, give a complete set of solutions to the chiral Dirac equation in four dimensions. This is essentially Penrose’s theorem, for this case. (The theorem is more general, applicable to other manifolds which admit twistor constructions.) The mapping between holomorphic functions in twistor space and massless fields in spacetime is known as the Penrose transform. (Strictly speaking, we are not concerned with holomorphic functions. They are holomorphic in some neighborhood in twistor space, and further, they are not really defined on $`\mathrm{𝐂𝐏}^3`$, since they have nontrivial degree of homogeneity. The proper mathematical characterization would be as sections of holomorphic sheaves of appropriate degree of homogeneity.)
### 6.2 An explicit example
As an explicit example of the Penrose transform, consider the holomorphic function
$$f(Z)=\frac{1}{aWbWcU}$$
(88)
where $`aW=a^{\dot{A}}x_{\dot{A}A}U^AU^1w^2U^2w^1`$, $`bW=b^{\dot{A}}x_{\dot{A}A}U^AU^1v^2U^2v^1`$. Defining $`z=U^2/U^1`$, we find, for the Penrose integral,
$`\psi ^A`$ $`=`$ $`{\displaystyle U𝑑U\frac{u^A}{aWbWcU}}`$ (89)
$`=`$ $`{\displaystyle 𝑑z\frac{U^A}{U^1}\frac{1}{(w^2zw^1)(v^2zv^1)c^2zc^1)}}`$
Taking the contour to enclose the pole at $`w^2/w^1`$, we find
$`\psi ^A`$ $`=`$ $`ϵ^{AB}{\displaystyle \frac{a^{\dot{A}}x_{\dot{A}B}}{x^2wc}}{\displaystyle \frac{1}{ab}}`$ (90)
$`=`$ $`ϵ^{AB}{\displaystyle \frac{a^{\dot{A}}x_{\dot{A}B}}{x^2(axc)}}{\displaystyle \frac{1}{ab}}`$
where $`axc=a^{\dot{A}}x_{\dot{A}A}c^A`$. (We take $`ab0`$.) One can check directly that this obeys the equation
$$_{\dot{A}A}\psi ^A=0$$
(91)
### 6.3 Conformal transformations
There is a natural action of conformal transformations on twistors. We can consider $`Z^\alpha `$ as a four-spinor of $`SU(4)`$, the latter acting as linear transformations on $`Z^\alpha `$, explicitly given by
$$Z^\alpha Z^\alpha =(gZ)^\alpha =g_\beta ^\alpha Z^\beta $$
(92)
where $`gSU(4)`$. The generators of infinitesimal $`SU(4)`$ transformations are thus given by
$$L_\beta ^\alpha =Z^\alpha \frac{}{Z^\beta }\frac{1}{4}\delta _\beta ^\alpha \left(Z^\gamma \frac{}{Z^\gamma }\right)$$
(93)
This may be split into different types of transformations as follows.
$`J_{AB}`$ $`=`$ $`U_A{\displaystyle \frac{}{U^B}}+U_B{\displaystyle \frac{}{U^A}}SU_L(2)`$
$`J_{\dot{A}\dot{B}}`$ $`=`$ $`U_{\dot{A}}{\displaystyle \frac{}{U^{\dot{B}}}}+U_{\dot{B}}{\displaystyle \frac{}{U^{\dot{A}}}}SU_R(2)`$
$`P^{A\dot{A}}`$ $`=`$ $`U^A{\displaystyle \frac{}{W_{\dot{A}}}}\mathrm{Translation}`$ (94)
$`K_{\dot{A}A}`$ $`=`$ $`W_{\dot{A}}{\displaystyle \frac{}{U^A}}\mathrm{Special}\mathrm{conformal}`$
$`\mathrm{transformation}`$
$`D`$ $`=`$ $`W_{\dot{A}}{\displaystyle \frac{}{W_{\dot{A}}}}U^A{\displaystyle \frac{}{U^A}}\mathrm{Dilatation}`$
where we have also indicated the interpretation of each type of generators. We see that the $`SU(4)`$ group is indeed the Euclidean conformal group; it is realized in a linear and homogeneous fashion on the twistor variables $`Z^\alpha `$. At the level of the purely holomorphic transformations, one can also choose the Minkowski signature, where upon the transformations given above become conformal transformations in Minkowski space, forming the group $`SU(2,2)`$.
### 6.4 Supertwistors
On can generalize the twistor space to an $`𝒩`$\- extended supertwistor space by adding fermionic or Grassman coordinates $`\xi _i`$, $`i=1,2,\mathrm{},𝒩`$; thus supertwistor space is parametrized by $`(Z^\alpha ,\xi _i)`$, with the identification $`Z^\alpha \lambda Z^\alpha `$, $`\xi _i\lambda \xi _i`$, where $`\lambda `$ is any nonzero complex number . $`\lambda `$ is bosonic, so only one of the bosonic dimensions is removed by this identification. Thus the supertwistor space is $`\mathrm{𝐂𝐏}^{3|𝒩}`$.
The case of $`𝒩=4`$ is special. In this case, one can form a top-rank holomorphic form on the supertwistor space; it is given by
$$\mathrm{\Omega }=\frac{1}{4!}ϵ_{\alpha \beta \gamma \delta }Z^\alpha dZ^\beta dZ^\gamma dZ^\delta d\xi _1d\xi _2d\xi _3d\xi _4$$
(95)
Notice that the bosonic part gets a factor of $`\lambda ^4`$ under the transformation $`Z^\alpha \lambda Z^\alpha `$, $`\xi _i\lambda \xi _i`$, while the fermionic part has a factor of $`\lambda ^4`$. $`\mathrm{\Omega }`$ is thus invariant under such scalings and becomes a differential form on the supermanifold $`\mathrm{𝐂𝐏}^{3|4}`$.
At this point, it is worth recalling the Calabi-Yau theorem .
> Theorem: For a given complex structure and Kähler class on a Kähler manifold, there exists a unique Ricci flat metric if and only if the first Chern class of the manifold vanishes or if and only if there is a globally defined top-rank holomorphic form on the manifold.
This is for an ordinary manifold. For the supersymmetric case, we will define a Calabi-Yau supermanifold as one which admits a globally defined top-rank holomorphic differential form . Whether such spaces admit a generalization of the Calabi-Yau theorem is not known. (For $`𝒩=1`$ spaces, a counterexample is known. However, super-Ricci flatness may follow from the vanishing of the first Chern class for $`𝒩2`$ .)
### 6.5 Lines in twistor space
Holomorphic lines in twistor space will turn out to be important for the construction of Yang-Mills amplitudes. First we will consider a holomorphic straight line, or a curve of degree one, in twistor space, giving the generalization to supertwistor space later. Since $`\mathrm{𝐂𝐏}^3`$ has three complex dimensions, we need two complex conditions to reduce to a line in twistor space. Thus we can specify a line in twistor space as the solution set of the equations
$$A_\alpha Z^\alpha =0,B_\alpha Z^\alpha =0$$
(96)
where $`A_\alpha ,B_\alpha `$ are constant twistors which specify the placement of the line in twistor space. These equations can be combined as
$$a_A^iU^A+b_{\dot{A}}^iW^{\dot{A}}=0$$
(97)
where $`A_\alpha =(a_A^1,b_{\dot{A}}^1),B_\alpha =(a_A^2,b_{\dot{A}}^2)`$. $`a,b`$ can be considered as $`(2\times 2)`$-matrices; $`deta`$ and $`detb`$ may both be nonzero, but both cannot be zero simultaneously, since equations (96) are then not sufficient to reduce to a line. We will take $`detb0`$ in the following. (The arguments presented will go through with appropriate relabelings if $`detb=0`$, but $`deta0`$.) In this case, $`b`$ is invertible and we can solve the equations (97) by
$`W_{\dot{A}}`$ $`=`$ $`(b^1a)_{\dot{A}A}U^A`$ (98)
$``$ $`x_{\dot{A}A}U^A`$
This shows that the condition (82) identifying the spacetime coordinates may be taken as defining a line in twistor space. In fact, here, $`x_{\dot{A}A}`$ specify the placement and orientation of the line in twistor space, in other words, they are the moduli of the line. We see that the moduli space of straight lines (degree-one curves) in twistor space is spacetime.
There is another way to write the equation (98). Recall that a line in a real space $`M`$ can be defined as a mapping of the interval $`[0,1]`$ into the space $`M`$, $`L:[0,1]M`$. We can do a similar construction for the complex case. We will define an abstract $`\mathrm{𝐂𝐏}^1`$ space by a two-spinor $`u^a`$ with the identification $`u^a\rho u^a`$ for any nonzero complex number $`\rho `$. Then we can regard a holomorphic line in twistor space as a map $`\mathrm{𝐂𝐏}^1\mathrm{𝐂𝐏}^3`$, realized explicitly as
$$U^A=(a^1)_a^Au^a,W_{\dot{A}}=(b^1)_{\dot{A}a}u^a$$
(99)
One can do $`SL(2,𝐂)`$ transformations on the coordinates $`u^a`$ of $`\mathrm{𝐂𝐏}^1`$; using this freedom, we can set $`a=1`$, or equivalently,
$$U^A=u^A,W_{\dot{A}}=x_{\dot{A}A}u^A$$
(100)
This is identical to (98).
The generalization of this to supertwistor space is now obvious. We will consider a map of $`\mathrm{𝐂𝐏}^1`$ to the supertwistor space $`\mathrm{𝐂𝐏}^{3|4}`$, given explicitly by
$`U^A=u^A,W_{\dot{A}}=x_{\dot{A}A}u^A`$
$`\xi ^\alpha =\theta _A^\alpha u^A`$ (101)
We have the fermionic moduli $`\theta _A^\alpha `$ in addition to the bosonic ones $`x_{\dot{A}A}`$.
The construction of curves of higher degree can be done along similar lines. A curve of degree $`d`$ is given by
$`Z^\alpha `$ $`=`$ $`{\displaystyle \underset{\{a\}}{}}a_{a_1a_2\mathrm{}a_d}^\alpha u^{a_1}u^{a_2}\mathrm{}u^{a_d}`$
$`\xi ^\alpha `$ $`=`$ $`{\displaystyle \underset{\{a\}}{}}\gamma _{a_1a_2\mathrm{}a_d}^\alpha u^{a_1}u^{a_2}\mathrm{}u^{a_d}`$ (102)
The coefficients $`a_{a_1a_2\mathrm{}a_d}^\alpha `$, $`\gamma _{a_1a_2\mathrm{}a_d}^\alpha `$ give the moduli of the curve. One can use $`SL(2,𝐂)`$ to set three of the coefficients to fixed values.
Given that each index $`a`$ takes values $`1,2`$, and the fact that the coefficients are symmetric in $`a_1,a_2,\mathrm{},a_n`$, we see that there are $`4(d+1)`$ bosonic and fermionic coefficients. The identification of $`u^a`$ and $`\rho u^a`$ and $`Z\lambda Z,\xi \lambda \xi `$ tells us that we can remove an overall scale degree of freedom. In other words, for the moduli, we have the identification,
$`a_{a_1a_2\mathrm{}a_d}^\alpha \lambda a_{a_1a_2\mathrm{}a_d}^\alpha `$
$`\gamma _{a_1a_2\mathrm{}a_d}^\alpha \lambda \gamma _{a_1a_2\mathrm{}a_d}^\alpha `$ (103)
for $`\lambda 𝐂\{0\}`$. Thus the moduli space of the curves may be taken as $`\mathrm{𝐂𝐏}^{4d+3|4d+4}`$. For the expressions of interest, as we shall see later, there is an overall $`SL(2,𝐂)`$ invariance, and hence three of the bosonic parameters can be fixed to arbitrarily chosen values.
## 7 Yang-Mills amplitudes and twistors
### 7.1 Why twistors are useful
In this section, we will start our discussion of the twistor approach to amplitudes in Yang-Mills theory , or multigluon scattering amplitudes, as they are often referred to.
We begin with the question of why the calculation of multigluon amplitudes are interesting. One of the motivations in seeking a twistor string theory was to obtain a weak coupling version of the standard duality between string theory on anti-de Sitter space and $`𝒩=4`$ supersymmetric Yang-Mills theory . However, developments in the subject over the last year or so have focused on, and yielded, many interesting results on the calculation of the scattering amplitudes themselves, so we shall concentrate on this aspect of twistors . Scattering amplitudes in any gauge theory are not only interesting in a general sense of helping to clarify a complicated interacting theory, but also there is a genuine need for them from a very practical point of view. This point can be illustrated by taking the experimental determination of the strong coupling constant as an example. We quote three values from the Particle Data Group based on three different processes.
$$\begin{array}{cccc}\hfill \alpha _s& =0.116\hfill & +0.003(expt.)\hfill & \pm 0.003(theory)\hfill \\ & & 0.005\hfill & \\ & & & \\ & =0.120\hfill & \pm 0.002(expt.)\hfill & \pm 0.004(theory)\hfill \\ & & & \\ & =0.1224\hfill & \pm 0.002(expt.)\hfill & \pm 0.005(theory)\hfill \end{array}$$
(104)
These values are for the momentum scale corresponding to the mas of the $`Z`$-boson, namely, $`\alpha _s(M_Z)`$; they are based on the Bjorken spin sum rule, jet rates in $`e`$-$`p`$ collisions and the photoproduction of two or more jets, respectively. Notice that the theoretical uncertainty is comparable to, or exceeds, the experimental errors. The major part of this comes from lack of theoretical calculations (to the order required) for the processes from which this value is extracted. Small as it may seem, this uncertainty can affect the hadronic background analysis at the Large Hadron Collider (currently being built at CERN), for instance. The relative signal strength for processes of interest, such as the search for the Higgs particle, can be improved if this uncertainty is reduced. This can also affect the estimate of the grand unification scale and theoretical issues related to it.
One could then ask the question: Since we know the basic vertices involved, and these are ultimately perturbative calculations, why not just do the calculations, to whatever order is required? Unfortunately, the direct calculation of the amplitudes is very, very difficult since there are large numbers, of the order of millions, of Feynman diagrams involved. (It is easy to see that the number of diagrams involved increases worse than factorially as the number of external lines increases.) Twistors provide a way to improve the situation.
A natural next question is then: What can twistors do, what has been accomplished so far? The progress so far may be summarized as follows.
1. It has been possible to write down a formula for all the tree level amplitudes in $`𝒩=4`$ supersymmetric Yang-Mills theory . Being at the tree level, this formula applies to the tree amplitudes of the nonsupersymmetric Yang-Mills theory as well. The formula reduces the calculation to the evaluation of the zeros of a number of polynomial equations and the evaluation of an integral. A certain analogy with instantons might be helpful in explaining the nature of this formula. To find the instanton field configurations, one must solve the self-duality conditions which are a set of coupled first-order differential equations. However, the ADHM procedure reduces this to an algebraic problem, namely, of solving a set of matrix equations, which lead to the construction of the appropriate holomorphic vector bundles. This algebraic problem is still difficult for large instanton numbers, nevertheless, an algebraicization has been achieved. In a similar way, the formula for the tree amplitudes replaces the evaluation of large numbers of Feynman diagrams, or an equivalent functional integral, by an ordinary integral whose evaluation requires the solution of some polynomial equations. This may still be difficult for cases with large numbers of negative helicity gluons; nevertheless, it is a dramatic simplification.
2. At the one-loop level, a similar formula has been obtained for all the so-called maximally helicity violating (MHV) amplitudes . A number of results for next-to-MHV amplitudes have been obtained .
3. A set of new diagrammatic rules based on the MHV vertices has been developed . Also, new types of recursion rules have been developed . These are promising new directions for perturbative analysis of a field theory.
4. Twistor inspired techniques have been used for some processes involving massive particles, particularly for the electroweak calculations .
### 7.2 The MHV amplitudes
We will begin with a discussion of the maximally helicity violating (MHV) amplitudes. These refer to $`n`$ gluon scattering amplitudes with $`n2`$ gluons of positive helicity and $`2`$ gluons of negative helicity. With $`n`$ gluons, $`n2`$ is the maximum number of positive helicity possible by conservation laws and such amplitudes are often refered to as the maximally helicity-violating (MHV) amplitudes.The analysis of these amplitudes will lead the way to the generalization for all amplitudes.
The gluons are massless and have momenta $`p_\mu `$ which obey the condition $`p^2=0`$; i.e., $`p_\mu `$ is a null vector. Using the identity and the Pauli matrices, we can write the four-vector $`p_\mu `$ as a $`2\times 2`$-matrix
$$p_{\dot{A}}^A=(\sigma ^\mu )_{\dot{A}}^Ap_\mu =\left(\begin{array}{cc}p_0+p_3& p_1ip_2\\ p_1+ip_2& p_0p_3\end{array}\right)$$
(105)
This matrix has a zero eigenvalue and, using this fact, we can see that it can be written as
$$p_{\dot{A}}^A=\pi ^A\overline{\pi }_{\dot{A}}$$
(106)
There is a phase ambiguity in the definition of $`\pi ,\overline{\pi }`$; $`\pi ^{}=e^{i\theta }\pi ,\overline{\pi }^{}=e^{i\theta }\overline{\pi }`$ give the same momentum vector $`p_\mu `$. Thus, physical results should be independent of this phase transformation. For a particular choice of this phase, an explicit realization of $`\pi ,\overline{\pi }`$ is given by
$$\pi =\frac{1}{\sqrt{p_0p_3}}\left(\begin{array}{c}p_1ip_2\\ p_0p_3\end{array}\right),\overline{\pi }=\frac{1}{\sqrt{p_0p_3}}\left(\begin{array}{c}p_1+ip_2\\ p_0p_3\end{array}\right)$$
(107)
The fact that $`p_\mu `$ is real gives a condition between $`\pi `$ and $`\overline{\pi }`$, which may be taken as $`\overline{\pi }_{\dot{A}}=(\pi ^A)^{}`$. We see that, for the momentum for each massless particle, we can associate a spinor momentum $`\pi `$.
There is a natural action of the Lorentz group $`SL(2,𝐂)`$ on the dotted and undotted indices, given by
$$\pi ^A\pi ^A=(g\pi )^A,\overline{\pi }_{\dot{A}}\overline{\pi }_{\dot{A}}^{}=(g^{}\overline{\pi })_{\dot{A}}$$
(108)
where $`gSL(2,𝐂)`$. The scalar product which preserves this symmetry is given by $`12=\pi _1\pi _2=ϵ_{AB}\pi _1^A\pi _2^B`$, $`[12]=ϵ^{\dot{A}\dot{B}}\overline{\pi }_{1\dot{A}}\overline{\pi }_{2\dot{B}}`$. At the level of vectors, this corresponds to the Minkowski product; i.e., $`\eta ^{\mu \nu }p_{1\mu }p_{2\nu }=p_1p_2=12[12]`$. Because of the $`ϵ`$-tensor, $`11=0`$ and the factorization (106) is consistent with $`p^2=0`$. (More generally, $`12=0`$ and $`[12]=0`$ if $`\pi _1`$ and $`\pi _2`$ are proportional to each other.) The scattering amplitudes can be simplified considerably when expressed in terms of these invariant spinor products. We will also define raising and lowering of the spinorial indices using the $`ϵ`$-tensor.
It is also useful to specify the polarization states of the gluons by helicity. The polarization vector $`ϵ_\mu `$ may then be written as
$$ϵ_\mu ϵ_{\dot{A}}^A=(\sigma ^\mu )_{\dot{A}}^Aϵ_\mu =\{\begin{array}{cc}\lambda ^A\overline{\pi }_{\dot{A}}/\pi \lambda & +1\mathrm{helicity}\\ & \\ \pi ^A\overline{\lambda }_{\dot{A}}/\overline{\pi }\overline{\lambda }& 1\mathrm{helicity}\end{array}$$
(109)
The spinor $`(\lambda ^A,\overline{\lambda }_{\dot{A}})`$ characterizes the choice of helicity.
We can now state the MHV amplitude for scattering of $`n`$ gluons, originally obtained by Parke and Taylor . They carried out the explicit calculation of Feynman diagrams, for small values of $`n`$, using some supersymmetry tricks for simplifications. Based on this, they guessed the general form of the amplitude; this guess was proved by Berends and Giele by using recursion rules for scattering amplitudes . The results are the following.
$`𝒜(1_+^{a_1},2_+^{a_2},3_+^{a_3},\mathrm{},n_+^{a_n})`$ $`=`$ $`0`$
$`𝒜(1_{}^{a_1},2_+^{a_2},3_+^{a_3},\mathrm{},n_+^{a_n})`$ $`=`$ $`0`$
$`𝒜(1_{}^{a_1},2_{}^{a_2},3_+^{a_3},\mathrm{},n_+^{a_n})`$ $`=`$ $`ig^{n2}(2\pi )^4\delta (p_1+p_2+\mathrm{}+p_n)`$
$`+\mathrm{noncyclic}\mathrm{permutations}`$
$`(1_{}^{a_1},2_{}^{a_2},3_+^{a_3},\mathrm{},n_+^{a_n})`$ $`=`$ $`12^4{\displaystyle \frac{\mathrm{Tr}(t^{a_1}t^{a_2}\mathrm{}t^{a_n})}{1223\mathrm{}n1nn1}}`$
The first nonvanishing amplitude with the maximum difference of helicities has two negative helicity gluons and $`n2`$ positive helicity gluons. This is what is usually called the MHV amplitude. In equation (7.2), $`g`$ is the gauge coupling constant. Notice that the amplitude $``$ is cyclically symmetric in all the particle labels except for the prefactor $`12^4`$. The latter refers to the momenta of the two negative helicity gluons. The summation over the noncyclic permutations makes the full amplitude symmetric in the gluon labels. We have taken all gluons as incoming. One can use the standard crossing symmetry to write down the corresponding amplitudes, with appropriate change of helicities, if some of the gluons are outgoing.
We will now carry out three steps of simplification of this result to bring out the twistor connection. $`\underset{¯}{Thefirststep:ThechiralDiracdeterminanton\mathrm{𝐂𝐏}^1}`$ Consider the functional determinant of the Dirac operator of a chiral fermion coupled to a gauge field $`A_{\overline{z}}`$ in two dimensions. By writing $`\mathrm{log}detD_{\overline{z}}=\mathrm{Tr}\mathrm{log}D_{\overline{z}}`$ and expanding the logarithm, we find
$`\mathrm{Tr}\mathrm{log}D_{\overline{z}}`$ $`=`$ $`\mathrm{Tr}\mathrm{log}(_{\overline{z}}+A_{\overline{z}})`$
$`=`$ $`\mathrm{Tr}\mathrm{log}\left(1+{\displaystyle \frac{1}{_{\overline{z}}}}A_{\overline{z}}\right)+\mathrm{constant}`$
$`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2x_1}{\pi }\frac{d^2x_2}{\pi }\mathrm{}\frac{(1)^{n+1}}{n}\frac{\mathrm{Tr}[A_{\overline{z}}(1)A_{\overline{z}}(2)\mathrm{}A_{\overline{z}}(n)]}{z_{12}z_{23}\mathrm{}z_{n1n}z_{n1}}}`$
where $`z_{12}=z_1z_2`$, etc. In writing this formula we have used the result
$$\left(\frac{1}{_{\overline{z}}}\right)_{12}=\frac{1}{\pi (z_1z_2)}$$
(112)
We can regard the $`z`$’s as local coordinates on $`\mathrm{𝐂𝐏}^1`$. Recall that $`\mathrm{𝐂𝐏}^1`$ is defined by two complex variables $`\alpha `$ and $`\beta `$, which may be regarded as a two-spinor $`u^a`$, $`u^1=\alpha ,u^2=\beta `$, with the identification $`u^a\rho u^a`$, $`\rho 𝐂\{0\}`$. On the coordinate patch with $`\alpha 0`$, we take $`z=\beta /\alpha `$ as the local coordinate. We can then write
$`z_1z_2`$ $`=`$ $`{\displaystyle \frac{\beta _1}{\alpha _1}}{\displaystyle \frac{\beta _2}{\alpha _2}}={\displaystyle \frac{\beta _1\alpha _2\beta _2\alpha _1}{\alpha _1\alpha _2}}`$ (113)
$`=`$ $`{\displaystyle \frac{ϵ_{ab}u_1^au_2^b}{\alpha _1\alpha _2}}={\displaystyle \frac{u_1u_2}{\alpha _1\alpha _2}}`$
Further, if we define $`\alpha ^2A\overline{z}=\overline{𝒜}`$, equation (LABEL:YM8) becomes
$$\mathrm{Tr}\mathrm{log}D_{\overline{z}}=\frac{1}{n}\frac{\mathrm{Tr}[\overline{𝒜}(1)\overline{𝒜}(2)\mathrm{}\overline{𝒜}(n)]}{(u_1u_2)(u_2u_3)\mathrm{}(u_nu_1)}$$
(114)
Notice that if $`u^a`$ is replaced by the spinor momentum $`\pi ^A`$, the denominator is exactly what appears in (7.2). The factor of $`1/n`$ gets cancelled out because (114) generates all permutations which gives $`n`$ times the sum over all noncyclic permutations. $`\underset{¯}{Thesecondstep:Thehelicityfactors}`$ The denominator for the MHV amplitude can be related to the chiral Dirac determinant as above. The factor $`12^4`$ can also be obtained if we introduce supersymmetry. We take up this second step of simplification now.
The transformation (108) shows that the Lorentz generator for the $`\pi `$’s is given by
$$J_{AB}=\frac{1}{2}\left(\pi _A\frac{}{\pi ^B}+\pi _B\frac{}{\pi ^A}\right)$$
(115)
where $`\pi _A=ϵ_{AB}\pi ^B`$. The spin operator is given by $`S_\mu ϵ_{\mu \nu \alpha \beta }J^{\nu \alpha }p^\beta `$, where $`J^{\mu \nu }`$ is the full Lorentz generator. This works out to $`S_{A\dot{A}}=J_B^A\pi ^B\overline{\pi }_{\dot{A}}=p_{\dot{A}}^As`$, identifying the helicity as
$$s=\frac{1}{2}\pi ^A\frac{}{\pi ^A}$$
(116)
Thus $`s`$ is, up to a minus sign, half the degree of homogeneity in the $`\pi `$’s. If we start with a positive helicity gluon, which would correspond to two negative powers of the corresponding spinor momentum, then we should expect an additional four factors of $`\pi `$ for a negative helicity gluon. Notice that there are two factors of spinor momenta in the denominator of the scattering amplitude (7.2) for each positive helicity gluon; for the two negative helicity gluons, because of the extra factor of $`12^4`$, the net result is two positive powers of $`\pi `$.
We now notice that if we have an anticommuting spinor $`\theta _A`$, $`d^2\theta \theta _A\theta _B=ϵ_{AB}`$, so that
$$d^2\theta (\pi \theta )(\pi ^{}\theta )=d^2\theta (\pi ^A\theta _A)(\pi ^B\theta _B)=\pi \pi ^{}$$
(117)
We see that an $`𝒩=4`$ theory is what we need to get four such factors, so as to get a term like $`12^4`$. Therefore we define an $`𝒩=4`$ superfield
$$\overline{A}^a(\pi ,\overline{\pi })=a_+^a+\xi ^\alpha a_\alpha ^a+\frac{1}{2}\xi ^\alpha \xi ^\beta a_{\alpha \beta }^a+\frac{1}{3!}\xi ^\alpha \xi ^\beta \xi ^\gamma ϵ_{\alpha \beta \gamma \delta }\overline{a}^{a\delta }+\xi ^1\xi ^2\xi ^3\xi ^4a_{}^a$$
(118)
where $`\xi ^\alpha =(\pi \theta )^\alpha =\pi ^A\theta _A^\alpha `$, $`\alpha =1,2,3,4`$. We can interpret $`a_+^a`$ as the classical version of the annihilation operator for a positive helicity gluon (whose gauge charge is specified by the Lie algebra index $`a`$), $`a_{}^a`$ as the annihilation operator for a negative helicity gluon; $`a_\alpha ^a,\overline{a}^{a\alpha }`$ correspond to four spin-$`\frac{1}{2}`$ particles and $`a_{\alpha \beta }^a`$ correspond to six spin-zero particles. This is exactly the particle content of $`𝒩=4`$ Yang-Mills theory.
We now choose the gauge potential in (114) to be given by
$$\overline{𝒜}=gt^a\overline{A}^a\mathrm{exp}(ipx)$$
(119)
Using this in the chiral Dirac determinant (114), we construct the expression
$$\mathrm{\Gamma }[a]=\frac{1}{g^2}d^8\theta d^4x\mathrm{Tr}\mathrm{log}D_{\overline{z}}]_{u^a\pi ^A}$$
(120)
It is then clear that the MHV amplitude can be written as
$`𝒜(1_{}^{a_1},2_{}^{a_2},3_+^{a_3},\mathrm{},n_+^{a_n})`$ $`=`$ $`i\left[{\displaystyle \frac{\delta }{\delta a_{}^{a_1}(p_1)}}{\displaystyle \frac{\delta }{\delta a_{}^{a_2}(p_2)}}{\displaystyle \frac{\delta }{\delta a_+^{a_3}(p_3)}}\mathrm{}{\displaystyle \frac{\delta }{\delta a_+^{a_n}(p_n)}}\mathrm{\Gamma }[a]\right]_{a=0}`$
An alternate representation involves introducing a supersymmetric version of the factor $`\mathrm{exp}(ipx)`$ for the $`𝒩=4`$ supermultiplet. Consider the function
$`\mathrm{exp}(i\eta \xi )=1+i\eta \xi +{\displaystyle \frac{1}{2!}}i\eta \xi i\eta \xi +{\displaystyle \frac{1}{3!}}i\eta \xi i\eta \xi i\eta \xi `$
$`+{\displaystyle \frac{1}{4!}}i\eta \xi i\eta \xi i\eta \xi i\eta \xi `$ (122)
where $`\eta \xi =\eta _\alpha \xi ^\alpha `$. Here $`\eta _\alpha `$ are four Grassman variables. We may think of them as characterizing the state of the external particles, specifically, their helicity. (The state is thus specified by the spinor momentum $`\pi `$ and $`\eta `$.) Since each $`\xi `$ carries one power of $`\pi `$, we see that we can associate the first term with a positive helicity gluon, the last with a negative helicity gluon, and the others with the superpartners of gluons, accordingly. We can then take
$$\overline{𝒜}=gt^a\varphi ^a\mathrm{exp}(ipx+i\eta \xi )$$
(123)
The scattering amplitudes are given by $`\mathrm{\Gamma }[a]`$ again, where, to get negative helicity for particles labeled $`1`$ and $`2`$ we take the coefficient of the factor $`\eta _{11}\eta _{21}\eta _{31}\eta _{41}\eta _{12}\eta _{22}\eta _{32}\eta _{42}`$, where the first subscript gives the component of $`\eta `$ and the second refers to the particle. We should also look at the term with $`n`$ factors of $`\varphi ^a`$ for the $`n`$-gluon amplitude. $`\underset{¯}{Thethirdstep:Reductiontoalineintwistorspace}`$ The results (LABEL:YM18-123) were known for a long time. The importance of supertwistor space was also recognized . (Some of the earlier developments, with connections to the self-dual Yang-Mills theory, etc., can be traced from .) Notice that with $`U^A=u^A`$, $`W=xu`$ and $`\xi `$, alongwith the condition $`\pi ^A=U^A`$, we are close to the usual variables of supertwistor space. More recently, Witten achieved enormous advances in this field by relating this formula to twistor string theory and curves in twistor space . To arrive at this generalization, first of all, we notice that the amplitude is holomorphic in the spinor momenta except for the exponential factor $`\mathrm{exp}(ipx)`$. We can rewrite this factor as follows.
$`\mathrm{exp}(ipx)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2}}\overline{\pi }^{\dot{A}}x_{\dot{A}A}\pi ^A\right)`$ (124)
$`=`$ $`\mathrm{exp}\left({\displaystyle \frac{i}{2}}\overline{\pi }^{\dot{A}}W_{\dot{A}}\right)]_{u^A=\pi ^A}`$
where $`W_{\dot{A}}=x_{\dot{A}A}\pi ^A`$. The strategy is now to regard $`W_{\dot{A}}`$ as a free variable, interpreting the condition $`W_{\dot{A}}=x_{\dot{A}A}u^A`$ as the restriction to a line in twistor space. We shall also use $`U^A=u^A`$, which is the other condition defining a line in twistor space. see (100). We can then write
$$𝑑\sigma \delta \left(\frac{\pi ^2}{\pi ^1}\frac{U^2}{U^1}\right)\mathrm{exp}\left(\frac{i}{2}\overline{\pi }^{\dot{A}}\pi ^1\frac{W_{\dot{A}}}{U^1}\right)=\mathrm{exp}(\frac{i}{2}\overline{\pi }^{\dot{A}}x_{\dot{A}A}\pi ^A)=\mathrm{exp}(ipx)$$
(125)
where $`\sigma =u^2/u^1`$, and we have used the restriction to the line $`W_{\dot{A}}=x_{\dot{A}A}u^A`$, $`U^A=u^A`$. The integration is along a line which contains the support of the $`\delta `$-function.
We can also treat $`\xi `$ as an independent variable, interpreting the condition $`\xi ^\alpha =\theta _A^\alpha u^A`$ as part of the line in supertwistor space, as in (101). The amplitude for $`n`$ particle scattering, with particle momenta labeled by $`\pi _i^A,\overline{\pi }_i^{\dot{A}}`$ and helicity factors $`\eta _{\alpha i}`$, can then be written as
$`𝒜`$ $`=`$ $`ig^{n2}{\displaystyle d^4xd^8\theta 𝑑\sigma _1\mathrm{}𝑑\sigma _n\frac{\mathrm{Tr}(t^{a_1}\mathrm{}t^{a_n})}{(\sigma _1\sigma _2)(\sigma _2\sigma _3)\mathrm{}(\sigma _n\sigma _1)}}`$ (126)
$`\times {\displaystyle \underset{i}{}}\delta ({\displaystyle \frac{\pi _i^2}{\pi _i^1}}{\displaystyle \frac{U^2(\sigma _i)}{U^1(\sigma _i)}})\mathrm{exp}({\displaystyle \frac{i}{2}}\overline{\pi }_i^{\dot{A}}\pi _i^1{\displaystyle \frac{W_{\dot{A}}(\sigma _i)}{U^1(\sigma _i)}}+i\pi _i^1\eta _{\alpha i}{\displaystyle \frac{\xi ^\alpha (\sigma _i)}{U^1(\sigma _i)}})`$
$`+\mathrm{noncyclic}\mathrm{permutations}`$
where the functions $`W_{\dot{A}},U^A,\xi ^\alpha `$ are given by
$`U^A=u^A,W_{\dot{A}}=x_{\dot{A}A}u^A`$
$`\xi ^\alpha =\theta _A^\alpha u^A`$ (127)
exactly as in (101). The variable $`\sigma `$ is given by $`\sigma =u^2/u^1`$. Notice that the overall factor of $`u^1`$ in $`W_{\dot{A}},U^A,\xi ^\alpha `$, cancels out in the formula (126).
### 7.3 Generalization to other helicities
The amplitude, in the form given in (126), shows a number of interesting properties. First of all, the amplitude is entirely holomorphic in the twistor variables $`Z^\alpha =(W_{\dot{A}},U^A)`$, $`\xi ^\alpha `$. It is also holomorphic in the variable $`\sigma `$ or $`u^a`$. (It is not holomorphic in $`\pi `$ since there is $`\overline{\pi }`$ in the exponentials, but this is immaterial for our arguments given below.) Secondly, the amplitude is invariant under the scalings $`Z^\alpha \lambda Z^\alpha `$, $`\xi ^\alpha \lambda \xi ^\alpha `$, so that it is a properly defined function on some neighborhood in the supertwistor space. Further, the amplitude has support only on a curve of degree one in supertwistor space given by (127). The moduli of this curve are given by $`x_{\dot{A}A}`$ and $`\theta _A^\alpha `$; there is integration over all these in the amplitude.
We may interpret this as follows. We consider a holomorphic map $`\mathrm{𝐂𝐏}^1\mathrm{𝐂𝐏}^{3|4}`$ which is of degree one. We pick $`n`$ points $`\sigma _1,\sigma _2,\mathrm{},\sigma _n`$ and then evaluate the integral in (126) over all $`\sigma `$’s and the moduli of the chosen curve.
The generalization of the formula suggested by Witten is to use curves of higher degree . In fact Witten argued, based on twistor string theory, that one should consider curves of degree $`d`$ and genus $`g`$, with
$$d=q1+l,gl$$
(128)
for $`l`$-loop Yang-Mills amplitudes with $`q`$ gluons of negative helicity. This generalization has been checked for various cases as mentioned before. For the tree amplitudes, the generalized formula reads
$`𝒜`$ $`=`$ $`ig^{n2}{\displaystyle 𝑑\mu 𝑑\sigma _1\mathrm{}𝑑\sigma _n\frac{\mathrm{Tr}(t^{a_1}\mathrm{}t^{a_n})}{(\sigma _1\sigma _2)(\sigma _2\sigma _3)\mathrm{}(\sigma _n\sigma _1)}}`$ (129)
$`\times {\displaystyle \underset{i}{}}\delta ({\displaystyle \frac{\pi _i^2}{\pi _i^1}}{\displaystyle \frac{U^2(\sigma _i)}{U^1(\sigma _i)}})\mathrm{exp}({\displaystyle \frac{i}{2}}\overline{\pi }_i^{\dot{A}}\pi _i^1{\displaystyle \frac{W_{\dot{A}}(\sigma _i)}{U^1(\sigma _i)}}+i\pi _i^1\eta _{\alpha i}{\displaystyle \frac{\xi ^\alpha (\sigma _i)}{U^1(\sigma _i)}})`$
$`+\mathrm{noncyclic}\mathrm{permutations}`$
where the curves of degree $`d`$ are
$`W_{\dot{A}}(\sigma )=(u^1)^d{\displaystyle \underset{0}{\overset{d}{}}}b_{\dot{A}k}\sigma ^k,U^A(\sigma )=(u^1)^d{\displaystyle \underset{0}{\overset{d}{}}}a_k^A\sigma ^k`$
$`\xi ^\alpha (\sigma )=(u^1)^d{\displaystyle \underset{0}{\overset{d}{}}}\gamma _k^\alpha \sigma ^k`$ (130)
This is exactly as in (102). The measure of integration for the moduli in (129) is given by
$$d\mu =\frac{d^{2d+2}ad^{2d+2}bd^{4d+4}\gamma }{vol[GL(2,𝐂)]}$$
(131)
The division by the volume of $`GL(2,𝐂)`$ arises as follows. There is an overall scale invariance for the integrand in (129), which means that we can remove one complex scale factor, corresponding to the moduli space being $`\mathrm{𝐂𝐏}^{4d+3|4d+4}`$. The integrand is also holomorphic in $`\sigma `$ and so has invariance under the $`SL(2,𝐂)`$ transformations $`u^au^a=(gu)^a`$ where $`g`$ is a $`(2\times 2)`$-matrix with unit determinant, or an element of $`SL(2,𝐂)`$. We must remove this factor to get an integral which does not diverge.
The actual evaluation of the integral can still be quite involved. One has to identify the zeros of the functions $`U^2(\sigma )/U^1(\sigma )`$ to integrate over the $`\delta `$-functions. This can be difficult to do explicitly for arbitrary values of the moduli. This is then followed by the integration over the moduli. Nevertheless, the formulae (129, 130) constitute a significant achievement. They reduce the problem of amplitude calculations in the gauge theory to an ordinary, multidimensional integral. As mentioned at the beginning of this section, this reduction has a status somewhat similar to what the ADHM construction has achieved for instantons.
Another important qualification about the formula (129) is that the integrals have to be defined by a continuation to real variables. The spacetime signature has to be chosen to be $`(++)`$ to be compatible with this. One has to carry out the analytic continuation after the integrals are done.
The justification for the generalization embodied in (129) comes from twistor string theory. We shall now briefly review this connection following Witten’s construction of twistor string theory ; there is an alternative string theory proposed by Berkovits which can also be used . Some of the structure of the latter will be used in section 9.
## 8 Twistor string theory
As mentioned in section 6, supertwistor space $`\mathrm{𝐂𝐏}^{3|4}`$ is a Calabi-Yau space. This allows the construction of a topological $`B`$-model with $`\mathrm{𝐂𝐏}^{3|4}`$ as the target space. In this theory, one considers open strings which end on $`D5`$-branes, with the condition $`\overline{\xi }=0`$. The gauge fields which characterize the dynamics of the ends of the open strings is then a potential $`\overline{𝒜}(Z,\overline{Z},\xi )`$ which can be checked to have the same content as the $`𝒩=4`$ gauge theory. An effective action for the topological sector can be written down; it is given by
$$=\frac{1}{2}_Y\mathrm{\Omega }\mathrm{Tr}(\overline{𝒜}\overline{}\overline{𝒜}+\frac{2}{3}\overline{𝒜}^3)$$
(132)
Here $`Y`$ is a submanifold of $`\mathrm{𝐂𝐏}^{3|4}`$ with $`\overline{\xi }=0`$. To linear order in the fields, the equations of motion for (132) correspond to $`\overline{}\overline{𝒜}=0`$. Thus, the fields are holomorphic in the $`Z`$’s and, via the Penrose transform, they correspond to massless fields in ordinary spacetime. Including the nonlinear terms, one can ask for an action in terms of the fields in spacetime which is equivalent to (132). This is given by
$$=\mathrm{Tr}\left[G^{AB}F_{AB}+\overline{\chi }^{A\alpha }D_{A\dot{A}}\chi _\alpha ^{\dot{A}}+\mathrm{}\right]$$
(133)
$`G^{AB}`$ is a self-dual field, $`F_{AB}`$ of the self-dual part of the usual field strength $`F_{\mu \nu }`$ of an ordinary gauge potential, $`\chi ,\overline{\chi }`$ are fermionic fields, etc. In terms of helicities $`G^{AB}`$ corresponds to $`1`$, the nonvanishing field strength $`F_{\dot{A}\dot{B}}`$ corresponds to $`+1`$ and so on. The action (133) cannot generate amplitudes with arbitrary number of negative helicity gluons; it is not the (super) Yang-Mills action either. The usual Yang-Mills term can be generated by the effect of a $`D1`$-instanton, which can lead to a term of the form $`\frac{1}{2}ϵG^2`$, where $`ϵ`$ is related to the action for the instanton. Integrating out the $`G`$-field, we then get a Yang-Mills term of the form $`F^2/4g^2`$ with a Yang-Mills coupling constant $`g^2ϵ`$. A term with $`q`$ factors of $`G`$, corresponding to $`q`$ particles of helicity $`1`$, will require $`q1`$ powers of $`ϵ`$. This corresponds to instanton number $`q1=d`$. Such $`D1`$-instantons are described by holomorphic curves of degree $`d`$. This is in agreement with the formula (128). This is the basic argument, schematically, why we should expect curves of degree $`d`$ to lead to tree-level amplitudes in the $`𝒩=4`$ Yang-Mills theory.
We have not discussed amplitudes at the loop level yet. A similar approach with curves of genus one has been verified at the one-loop level . However, the general formula in terms of higher genus curves has not been very useful for actual computations. This is partially due to the complexity of the formula. A more relevant reason is the emergence of new rules to calculate both tree level and loop level amplitudes using a sewing procedure with the basic MHV amplitudes as the vertices . The MHV amplitudes have to be continued off-shell for this reason. An off-shell extension has been proposed and used in .
The one-loop amplitudes which emerge naturally are the amplitudes for the $`𝒩=4`$ Yang-Mills theory. If the external (incoming and outgoing) particles are gluons, the superpartners can only occur in loops. As a result, at the tree-level, one can get the amplitudes in the pure Yang-Mills theory with no supersymmetry by restricting the external lines to be gluons. But at the one-loop level all superpartners can contribute. While the one-loop amplitudes for the $`𝒩=4`$ Yang-Mills theory are interesting in their own right, the corresponding amplitudes in the nonsupersymmetric theory are of even greater interest, since they pertain to processes which are experimentally accessible. It is not trivial to extract the nonsupersymmetric amplitudes from the $`𝒩=4`$ theory. One approach is to subtract out the contributions of the superpartners. An alternative is to build up the one-loop amplitude from the unitarity relation, using tree amplitudes. In principle, this can only yield the imaginary part of the one-loop amplitude. (One could attempt to construct the real part via dispersion relations. While this is usually ambiguous due to subtractions needed for the dispersion integrals, the $`𝒩=4`$ theory, which is finite, is special. There are relations among amplitudes which can be used for this theory.) If one makes an ansatz for some off-shell extensions of the tree amplitudes, one can obtain, via unitarity relations, some of the one-loop amplitudes in the nonsupersymmetric theory as well. The off-shell extensions can be checked for consistency in soft-gluon limits, etc., so they are fairly unambiguous. The results quoted in the beginning of section 7 emerged from such analyses. Notice that the state of the art here is a combination of rules emerging from the twistor approach and unitarity relations and a bit of guess work. The new set of recursion rules also has been very useful .
## 9 Landau levels and Yang-Mills amplitudes
### 9.1 The general formula for amplitudes
There is an interesting relationship between the amplitudes of the Yang-Mills theory and the Landau problem or the problem of quantum Hall effect. (This is also related to Berkovits’ twistor string theory .) To see how this connection arises, we start by rewriting the formula (126) in more compact form as follows. Define a one-particle wave function for the $`𝒩=4`$ supermultiplet by
$`\mathrm{\Phi }(\pi ,\overline{\pi },\eta )`$ $`=`$ $`\delta \left({\displaystyle \frac{\pi ^2}{\pi ^1}}{\displaystyle \frac{U^2(\sigma )}{U^1(\sigma )}}\right)\mathrm{exp}\left({\displaystyle \frac{i}{2}}\overline{\pi }^{\dot{A}}\pi ^1{\displaystyle \frac{W_{\dot{A}}(\sigma )}{U^1(\sigma )}}+i\pi ^1\eta _\alpha {\displaystyle \frac{\xi ^\alpha (\sigma )}{U^1(\sigma )}}\right)`$
$`=`$ $`\delta (\mathrm{\Pi }Z(\sigma )){\displaystyle \frac{Z(\sigma )A}{\mathrm{\Pi }A}}\mathrm{exp}\left({\displaystyle \frac{i}{2}}{\displaystyle \frac{\overline{\mathrm{\Pi }}Z(\sigma )\mathrm{\Pi }A}{Z(\sigma )A}}+i{\displaystyle \frac{\mathrm{\Pi }A}{Z(\sigma )A}}\eta \xi (\sigma )\right)`$
where we introduced the twistors,
$$\mathrm{\Pi }^\alpha =(0,\pi ^A)=(0,0,\pi ^1,\pi ^2),A_\alpha =(0,0,1,0)$$
(135)
which gives $`ZA=U^1`$, $`\mathrm{\Pi }A=\pi ^1`$. Notice again that $`\mathrm{\Phi }`$ is holomorphic in the twistor variables $`(Z,\xi )`$ and is invariant under the scaling $`Z^\alpha \lambda Z^\alpha `$, $`\xi ^\alpha \lambda \xi ^\alpha `$. It is also invariant under the scaling of the twistor $`A_\alpha `$. There is also an obvious $`SU(4)`$ or $`SU(2,2)`$ invariance, if we transform $`A_\alpha `$ as well. Thus the expression for $`\mathrm{\Phi }`$ can be used with a more general choice of $`A_\alpha `$ than the one given in (135). The twistor $`A_\alpha `$ plays the role of the reference momentum which has been used in many discussions of scattering amplitudes; ultimately, it drops out of the physical results due to conservation of momentum.
On the $`\mathrm{𝐂𝐏}^1`$ with the homogeneous coordinates $`u^a`$, we define the holomorphic differential
$$(udu)udu=ϵ_{ab}u^adu^b=(u^1)^2d\sigma $$
(136)
As a result we can write
$`{\displaystyle 𝑑\sigma _1𝑑\sigma _2\mathrm{}𝑑\sigma _n\frac{1}{(\sigma _1\sigma _2)(\sigma _2\sigma _3)\mathrm{}(\sigma _n\sigma _1)}}`$
$`={\displaystyle (udu)_1(udu)_2\mathrm{}(udu)_n\frac{1}{(u_1u_2)(u_2u_3)\mathrm{}(u_nu_1)}}`$ (137)
Using the formulae (LABEL:LL1, 137), we can write the amplitude (129) as
$`𝒜`$ $`=`$ $`{\displaystyle 𝑑\mu \underset{i}{}(udu)_i\mathrm{\Phi }(\pi _i,\overline{\pi }_i,\eta _i)\frac{\mathrm{Tr}(t^{a_1}\mathrm{}t^{a_n})}{(u_1u_2)(u_2u_3)\mathrm{}(u_nu_1)}}`$ (138)
$`+\mathrm{noncyclic}\mathrm{permutations}`$
Finally, the denominators arise from the chiral Dirac determinant, so we can also write this as
$$𝒜=d\mu \underset{i}{}(udu)_i\mathrm{\Phi }(\pi _i,\overline{\pi }_i,\eta _i)\left(\frac{\delta }{\delta \overline{𝒜}^{a_1}(u_1)}\mathrm{}\frac{\delta }{\delta \overline{𝒜}^{a_n}(u_n)}\right)\mathrm{Tr}\mathrm{log}D_{\overline{z}}]_{\overline{𝒜}=0}$$
(139)
This formula takes care of the permutations as well. The holomorphic curves of degree $`d`$ are as given in (102).
### 9.2 A field theory on $`\mathrm{𝐂𝐏}^1`$
We now consider a field theory on $`\mathrm{𝐂𝐏}^1`$ or the two-sphere. The action is given by
$$𝒮=𝑑\mu (\mathrm{𝐂𝐏}^1)\left[\overline{q}(\overline{}+\overline{𝒜})q+\overline{Y}(\overline{D}Q)\right]$$
(140)
Here $`q`$ and $`\overline{q}`$ are standard fermionic fields, so that the first term is the chiral Dirac action on $`\mathrm{𝐂𝐏}^1`$. These fields are analogous to what generates the extra current algebra in Berkovits’ paper . In the second term, $`Q`$ stands for the supertwistor variables $`(Z^\alpha ,\xi ^\alpha )`$. $`\overline{Y}`$ is another field with values in twistor variables again. Thus the second term corresponds to a two-dimensional action with target space $`𝐂^{4|4}`$. Further, $`\overline{D}=\overline{}+\overline{A}`$, where $`\overline{A}`$ is a $`GL(1,𝐂)`$ gauge field. This term has invariance under the scaling $`Q\lambda Q`$, $`\overline{Y}\lambda ^1\overline{Y}`$, with the gauge transformation $`\overline{A}\overline{A}\overline{}\mathrm{log}\lambda `$. Since there are equal number of fermions and bosons, the $`GL(1,𝐂)`$ transformation has no anomaly and the use of the chiral action for the second term is consistent with this symmetry at the quantum level.
We will consider the functional integral of $`e^𝒮`$. The integration over the fermions $`q,\overline{q}`$ leads to an effective action $`\mathrm{Tr}\mathrm{log}\overline{D}`$, and hence, the connected correlator with $`n`$ factors of $`\overline{𝒜}`$ is given by
$`M`$ $`=`$ $`{\displaystyle }[d\overline{Y}dZd\overline{A}]\mathrm{exp}({\displaystyle }\overline{Y}\overline{D}Z)\times `$ (141)
$`{\displaystyle \frac{d^2\sigma _1}{\pi }\mathrm{}\frac{d^2\sigma _n}{\pi }\frac{\mathrm{Tr}[\overline{𝒜}(1)\overline{𝒜}(2)\mathrm{}\overline{𝒜}(n)]}{(\sigma _1\sigma _2)(\sigma _2\sigma _3)\mathrm{}(\sigma _n\sigma _1)}}`$
We will now take $`\overline{𝒜}`$ to be of the form $`\overline{𝒜}=\overline{}\mathrm{\Phi }`$. We can use this in (141) and carry out the integration over $`d\overline{\sigma }`$’s. Partial integration can produce $`\delta `$-functions like $`\delta ^{(2)}(z_1z_2)`$. If we exclude coincidence of points, which will correspond to coincidence of external momenta in the context of multigluon scattering, then these $`\delta `$-functions have no support and the only contribution is from the boundary. The correlator (141) then becomes
$`M`$ $`=`$ $`{\displaystyle }[d\overline{Y}dZd\overline{A}]\mathrm{exp}({\displaystyle }\overline{Y}\overline{D}Z)\times `$ (142)
$`{\displaystyle \frac{d\sigma _1}{2\pi i}\mathrm{}\frac{d\sigma _n}{2\pi i}\frac{\mathrm{Tr}[\mathrm{\Phi }(1)\mathrm{\Phi }(2)\mathrm{}\mathrm{\Phi }(n)]}{(\sigma _1\sigma _2)(\sigma _2\sigma _3)\mathrm{}(\sigma _n\sigma _1)}}`$
Using the result (137), we can rewrite this as
$`M`$ $`=`$ $`{\displaystyle }[d\overline{Y}dZd\overline{A}]\mathrm{exp}({\displaystyle }\overline{Y}\overline{D}Z)\times `$
$`{\displaystyle \frac{(udu)_1}{2\pi i}\mathrm{}\frac{(udu)_n}{2\pi i}\frac{\mathrm{Tr}[\mathrm{\Phi }(1)\mathrm{\Phi }(2)\mathrm{}\mathrm{\Phi }(n)]}{(u_1u_2)(u_2u_3)\mathrm{}(u_nu_1)}}`$
We now turn to the integration over the gauge field $`\overline{A}`$. There are nontrivial $`U(1)`$ bundles over $`S^2\mathrm{𝐂𝐏}^1`$, corresponding to Dirac monopoles. The same result holds for $`GL(1,𝐂)`$. The space of gauge potentials then splits up into a set of disconnected pieces, one for each monopole number. In each sector, we can write $`\overline{A}=\overline{A}_d+\delta \overline{A}`$, where $`\overline{A}_d`$ is a fixed configuration of monopole number $`d`$ and $`\delta \overline{A}`$ is a fluctuation (of zero monopole number). In two dimensions, we can always write $`\delta \overline{A}=\overline{}\mathrm{\Theta }`$ for some complex function $`\mathrm{\Theta }`$ on $`\mathrm{𝐂𝐏}^1`$.The measure of integration splits up as $`[d\overline{A}]=[d\mathrm{\Theta }]det\overline{}`$. The determinant can contribute to the conformal anomaly, if we interpret this as the world-sheet formulation of a string theory calculation. For us, thinking of this as a two-dimensional field theory, this anomaly is not relevant. We can now eliminate $`\mathrm{\Theta }`$ by absorbing it into the definition of the fields. $`\mathrm{\Phi }`$ will be chosen to be $`GL(1,𝐂)`$ gauge-invariant, so this will not affect it. The integration over the fields $`\overline{A}`$ is thus reduced to a summation over the different monopole numbers, with a fixed representative background field $`\overline{A}_d`$ for each value of $`d`$.
For the integration over the twistor fields, we need a mode expansion. For the sector with monopole number $`d`$, we can expand the fields as
$`Z^\alpha `$ $`=`$ $`{\displaystyle \underset{\{a\}}{}}a_{a_1a_2\mathrm{}a_d}^\alpha u^{a_1}u^{a_2}\mathrm{}u^{a_d}+\mathrm{higher}\mathrm{Landau}\mathrm{levels}`$
$`\xi ^\alpha `$ $`=`$ $`{\displaystyle \underset{\{a\}}{}}\gamma _{a_1a_2\mathrm{}a_d}^\alpha u^{a_1}u^{a_2}\mathrm{}u^{a_d}+\mathrm{higher}\mathrm{Landau}\mathrm{levels}`$ (144)
with similar expansions for $`\overline{Y}`$. The first set of terms correspond to the lowest Landau level, the higher terms, which we have not displayed explicitly, correspond to higher Landau levels. The wave functions (or mode functions) for the lowest Landau level are holomorphic in the $`u`$’s, the higher levels involve $`\overline{u}`$’s as well. The functional integration is now over the coefficients $`a_{a_1a_2\mathrm{}a_d}^\alpha `$, $`\gamma _{a_1a_2\mathrm{}a_d}^\alpha `$, etc. Notice that the zero modes define a holomorphic curve of degree $`d`$ in supertwistor space. Our choice of $`\mathrm{\Phi }`$ will have no $`\overline{Y}`$, so, in integrating over the nonzero modes, one cannot have any propagators or loops generated by $`\overline{Y}Z`$ Wick contractions. As a result, the nonzero modes give only an overall normalization factor. In any correlator involving only $`Z`$’s and $`\xi `$’s, we can saturate the fields by the zero modes. Then, apart from constant factors, the correlator (142) becomes
$`M`$ $`=`$ $`{\displaystyle \underset{d}{}}C_d_d`$
$`_d`$ $`=`$ $`{\displaystyle 𝑑\mu (a,\gamma )\frac{(udu)_1}{2\pi i}\mathrm{}\frac{(udu)_n}{2\pi i}\frac{\mathrm{Tr}[\mathrm{\Phi }(1)\mathrm{\Phi }(2)\mathrm{}\mathrm{\Phi }(n)]}{(u_1u_2)(u_2u_3)\mathrm{}(u_nu_1)}}`$
We now choose the $`\mathrm{\Phi }`$’s as given in (LABEL:LL1), and the correlator (LABEL:LL11) becomes the multigluon amplitude given in (139), once we can show that the integration over the moduli space is given by the invariant measure (131). To see how this arises, consider an $`SL(2,𝐂)`$ transformation of the coefficients of the mode expansion given by
$`a_{a_1a_2\mathrm{}a_d}^\alpha `$ $`=`$ $`a_{b_1b_2\mathrm{}b_d}^\alpha g_{a_1}^{b_1}g_{a_2}^{b_2}\mathrm{}g_{a_n}^{b_n}`$
$`\gamma _{a_1a_2\mathrm{}a_d}^\alpha `$ $`=`$ $`\gamma _{a_1a_2\mathrm{}a_d}^\alpha g_{a_1}^{b_1}g_{a_2}^{b_2}\mathrm{}g_{a_n}^{b_n}`$ (146)
where $`gSL(2,𝐂)`$. If we use $`(a^{},\gamma ^{})`$, this is equivalent to using $`(a,\gamma )`$ and redefined $`u`$’s, $`u^a=g_b^au^b`$ in $`\mathrm{\Phi }`$. Since $`udu`$ and scalar products like $`(u_1u_2)`$ are invariant under such transformations, we can change the variables $`u^{}u`$; the integrand for the integration over the moduli $`(a,\gamma )`$ is thus $`SL(2,𝐂)`$-invariant. We must therefore consider the measure (131) to obtain well-defined correlators.
What we have shown is that there is a set of correlators in the two-dimensional problem defined by (140) which can be calculated exactly, being saturated by the lowest Landau levels, and which give the multigluon amplitudes. The YM amplitudes are thus obtained as a set of “holomorphic” correlators of the two-dimensional problem. The fermionic integration obviously leads to an expression which is similar to the Laughlin wave function for the $`\nu =1`$ quantum Hall state. If we consider just the value $`\alpha =1`$, the fermionic integration has the form
$$[d\gamma ]\underset{i}{}\gamma _{a_1a_2\mathrm{}a_n}^1u_i^{a_1}\mathrm{}u_i^{a_n}\underset{i<j}{}(u_iu_j)$$
(147)
(With four sets of such terms for $`\alpha =1`$ to $`\alpha =4`$, we get the fourth power of the right hand side.) In some sense, we can interpret the integration over the bosonic moduli as defining the bosonic version of the Laughlin wave function, since it is related to the fermionic one by the natural supersymmetry of the expression $`\mathrm{\Phi }`$ in (LABEL:LL1). Whether this Landau level point of view for the amplitudes can lead to new insights into the YM problem is yet to be seen.
This work was supported by the National Science Foundation grant number PHY-0244873 and by a PSC-CUNY grant.
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# Electromagnetically induced transparency for Λ - like systems with a structured continuum
## Abstract
Electric susceptibility of a laser-dressed atomic medium is calculated for a model $`\mathrm{\Lambda }`$-like system including two lower states and a continuum structured by a presence of an autoionizing state or a continuum with a laser-induced structure. Depending on the strength of a control field it is possible to obtain a significant reduction of the light velocity in a narrow frequency window in the conditions of a small absorption. A smooth transition is shown between the case of a flat continuum and that of a discrete state serving as the upper state of a $`\mathrm{\Lambda }`$ system.
It has been known for a long time that atoms in the $`\mathrm{\Lambda }`$ configuration irradiated by two laser fields may exhibit peculiar dynamical properties. In particular the population may be trapped in the so-called dark state, being a combination of the two lower states. It constitutes a basis for subtle coherent dynamical manipulations like an efficient population transfer (stimulated Raman adiabatic passage - STIRAP), robust against any broadening of the upper state b1 ; b2 . Effects of such a type are present also in the case in which the upper state of a $`\mathrm{\Lambda }`$ configuration has been replaced by a continuum b3 ; b4 ; b5 or, in strong laser fields, by a set of coupled continua b6 .
Dynamical effects in particular atoms are reflected in unusual light propagation effects in atomic media. One of the most important propagation effects in a medium of atoms in the $`\mathrm{\Lambda }`$ configuration is the electromagnetically induced transparency (EIT), which consists in making the medium transparent for a weak probe laser beam by irradiating it by a strong copropagating control beam b7 ; b8 . Instead of an absorption line there appears a transparency window and a normal dispersion. By switching the control field off when the pulse is inside the sample one can slow the light pulse down and finally stop or store it b9 ; b10 ; b11 . The photon energy of the signal is transferred to the control beam while the information about the pulse is written down in the form of a coherence (nondiagonal element of the density matrix) between the two lower states. Switching the signal on again results in reading the information out: the pulse is reconstructed in a coherent way.
A natural question arises whether similar propagation effects may occur in the case of a continuum serving as an upper state. Some answer has been given by van Enk et al. b12 ; b13 , who considered a general case of two beams propagating inside the medium. Their conclusion was that a necessary condition for EIT was that the asymmetry parameter in the continuum model was zero. In their case the losses due to photoionization are of second order with respect to the probe field. Thus, in the linear approximation with respect to the probe beam and neglecting propagation effects for the strong control beam, one still may apply the standard approach in which the probe propagation is discussed in terms of an atomic susceptibility. Below we give analytic expressions for the susceptibility in the case of the bound-continuum dipole matrix elements being modeled according to Fano autoionization theory b14 . By changing the Fano asymmetry parameter we pass from the case of a flat continuum to that of a discrete and broadened bound state. We examine the shape of the transparency window depending on the amplitude of the coupling field.
We consider a generalization of a $`\mathrm{\Lambda }`$ system in which the upper state is replaced by a continuum: an atomic system including two lower discrete states $`b`$ (initial) and $`c`$ and a continuum $`E`$ coupled with the state $`b`$ by a weak probe field $`(\frac{1}{2}ϵ_1(z,t)\mathrm{exp}[i(k_1z\omega _1t)]+c.c`$ and with the state $`c`$ with a relatively strong control field $`ϵ_2\mathrm{cos}\omega _2t`$. The continuum may have a Fano density of states, due to a configurational coupling with an autoionizing state or to a structure induced by another laser b15 . As usually in EIT, the control field, the propagation effects for which are neglected, dresses the atomic medium to create new conditions for the propagation of the probe pulse. The evolution of the atomic system is described by the von Neumann equation, which after transforming-off the rapidly oscillating terms and making the rotating-wave approximation is reduced in the first order perturbation with respect to the probe field to the set of the following equations for the density matrix $`\sigma `$
$`i\mathrm{}\dot{\sigma }_{Eb}=(EE_b\mathrm{}\omega _1)\sigma _{Eb}{\displaystyle \frac{1}{2}}d_{Eb}ϵ_1{\displaystyle \frac{1}{2}}d_{Ec}ϵ_2\sigma _{cb},`$
$`i\mathrm{}\dot{\sigma }_{cb}=(E_c+\mathrm{}\omega _2E_b\mathrm{}\omega _1i\gamma _{cb})\sigma _{cb}{\displaystyle \frac{1}{2}}ϵ_2^{}{\displaystyle d_{cE}\sigma _{Eb}𝑑E}.`$ (1)
In the above equation $`d`$ is the dipole moment and $`\gamma _{cb}`$ is a phenomenological relaxation rate for the coherence $`\sigma _{cb}`$. Stationary solutions of Eqs (1) can be found by first expressing $`\sigma _{cb}`$ in terms of $`\sigma _{Eb}`$ and then by solving the integral equation for the latter. The component of the polarization of the medium connected with the $`bE`$ coupling is
$$P^+(\omega _1)=Nd_{bE}\sigma _{Eb}𝑑E=ϵ_0\chi (\omega _1)ϵ_1(\omega _1),$$
(2)
with $`ϵ_0`$ being the vacuum electric permittivity, $`N`$ \- the atom density and with the medium susceptibility given by
$$\chi (\omega _1)=\frac{N}{2ϵ_0}\left(R_{bb}+\frac{1}{4}\frac{ϵ_2^2R_{bc}R_{cb}}{E_b+\mathrm{}\omega _1E_c\mathrm{}\omega _2+i\gamma _{cb}\frac{1}{4}ϵ_2^2R_{cc}}\right).$$
(3)
The function $`R_{ij}(\omega _1)`$, $`i,j=b,c`$, is given by
$$R_{ij}(\omega _1)=_0^{\mathrm{}}\frac{d_{iE}d_{Ej}}{E_b+\mathrm{}\omega _1E+i\eta }𝑑E,$$
(4)
where $`\eta 0^+`$ assures the correct behavior at the branch cut. The presence of an autoionizing state coupled with the continuum results in a modification of the density of the continuum states, which leads to the following Fano formula for the bound-structured continuum dipole matrix element
$$|d_{iE}|^2=B_i^2\frac{(EE_a+q\gamma )^2}{(EE_a)^2+\gamma ^2},$$
(5)
where $`E_a`$ and $`\gamma `$ characterize the position and width of the autoionizing resonance, $`q`$ is the Fano asymmetry parameter and $`B_i`$ is the bound-flat continuum matrix element. We assume that the matrix elements for the states $`b`$ and $`c`$ differ only by a constant factor $`B_{b,c}`$, i.e. the presence of the autoionizing state is the only reason for the continuum to exhibit any structure. For real systems the energy dependence of the matrix element is more complicated, in particular an additional form factor is necessary to account for a threshold behavior or a correct asymptotic properties. Taking such a form factor into account is not essential in an approach like ours, in which we concentrate on the frequency dependence of the susceptibility in a narrow band close to the resonance.
If the lasers are tuned sufficiently far from the ionization threshold and so is the localization of the autoionizing state we can extend the lower limit of the energy integral to $`\mathrm{}`$ and get from Eqs (4) and (5)
$$R_{ij}(\omega _1)=B_iB_j\pi \frac{\gamma (\mathrm{}\omega _1E_a)(q^21)+2q\gamma ^2i(\mathrm{}\omega _1E_a+q\gamma )^2}{(\mathrm{}\omega _1E_a)^2+\gamma ^2}$$
(6)
It is possible to evaluate the width of the structure in the susceptibility given by Eqs (3) and (6) if it is significantly smaller than an autoionizing width, i.e $`\gamma >>q\omega _1`$. The width for $`q>>1`$ is in atomic units of order of $`ϵ_2^2q^2`$. This result is to be compared with the width of the transparency window for EIT in a typical $`\mathrm{\Lambda }`$ system, i.e. one with a discrete broadened upper state. The width of the latter window is in atomic units of order of $`ϵ_2^2/\mathrm{\Gamma }`$, $`\mathrm{\Gamma }`$ being the relaxation rate for the nondiagonal element of the atomic coherence. One can see that for the same control fields the transparency window here is much more narrow. For a typical value of $`\mathrm{\Gamma }10^8`$ a.u. we need $`q10^4`$ to make the windows in the two cases comparable. The light absorption in the case of the continuum is by a few orders of magnitude weaker (again except for very large $`q`$) than in the case of a discrete upper level but due to the transparency window being narrow the dispersion can be very steep, so a significant light slowdown is possible.
As an illustration of the above formalism we show the atomic susceptibility calculated from Eqs (3) and (6). We have taken $`\gamma =10^9`$ a.u.$`6.6`$ MHz. The values of the field amplitude $`ϵ_2`$ ranged from $`10^9`$ to $`10^6`$ a.u., i.e. the width $`2\pi ϵ_2^2B_j^2/4`$ of the lower states, induced by the bound-continuum coupling for the flat continuum (without the Fano factor) would be of order of $`10^8\gamma `$ to $`10^2\gamma `$. The values of the parameters $`B_j`$ were taken arbitrarily $`B_1=2`$ a.u., $`B_2=3`$ a.u. The atomic density was $`N=0.67\times 10^{12}`$ cm<sup>-3</sup>. The asymmetry parameter $`q`$ was of order of 10-100. The position of the states was taken $`E_a=E_b/2=E_c=0.1`$ a.u.= 2.72 eV. The relaxation rate $`\gamma _{cb}`$, possibly due to the spontaneous emission from the level $`c`$ or to an incoherent ionization channel by the field $`ϵ_2`$ from the level $`b`$, was neglected.
In Fig. 1 we show the medium susceptibility for $`q`$=10 and $`ϵ_2=0`$. As expected, a normal dispersion is observed except close to the resonance, together with a single absorption line. Note that the maximum values of both the real and imaginary parts of $`\chi `$ are smaller than $`10^9`$. Fig. 2 shows the susceptibility after switching on a relatively strong control field $`ϵ_2=10^6`$ a.u. Note the transparency window of the width of order of $`\gamma `$ and the corresponding interval of a normal dispersion at a resonance frequency. Due to the Fano asymmetry the curves in the figure lack symmetry too. The group index $`n_g=1+\frac{\omega _1}{2}\frac{d}{d\omega _1}\mathrm{Re}\chi (\omega _1)`$, being the factor by which the probe signal velocity is reduced, is only about 1.1. Figs 3 and 4 show how the frequency window is narrowed for smaller control fields. The widths are of order of $`10^2\gamma `$ and $`10^4\gamma `$ respectively, while the corresponding group indices are of order of 10 and 1000. The spectral width of the laser beam needed in the latter case is however not yet accessible. The effects of a reduced absorption and a steep normal dispersion are made stronger if the Fano parameter $`q`$ is increased. In Fig. 5 we show the suspecptibility for $`q=100`$ and $`ϵ_2=10^9`$ a.u. ($`q`$ has been increased and $`ϵ_2`$ has been decreased by the factor of 10 compared with Fig. 4). Note that absolute values of $`\chi `$ are increased by two orders of magnitude compared with those of Fig. 4 while the transparency window is the same. This means that the group index is about $`10^5`$.
We have investigated the optical properties of an atomic medium in the $`\mathrm{\Lambda }`$ configuration in which the upper state is replaced by a continuum. Such a system exhibits properties similar to those of a typical $`\mathrm{\Lambda }`$ system, except that the orders of their magnitude are changed. Here the windows of a reduced absorption (which is in general much smaller for the systems examined here) and of a normal dispersion are much more narrow, if the control fields are not made stronger. The dispersion curve in the transparency window may still be steep, which gives rise to a significant light slowdown. By increasing the Fano asymmetry parameter by a few orders of magnitude we smoothly pass to the case of a typical $`\mathrm{\Lambda }`$ system.
###### Acknowledgements.
The work has been supported in part by the Committee for Scientific Research, grant No. 1 P03B 010 28. The subject belongs to the scientific program of the National Laboratory of AMO Physics in Toruń, Poland.
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# A note on 𝑞-Volkenborn Integration
## 1. Introduction
Let $`p`$ be a fixed odd prime, and let $`_p`$ denote the $`p`$-adic completion of the algebraic closure of $`_p.`$ For $`d`$ a fixed positive integer with $`(p,d)=1`$, let
$$X=X_d=\underset{}{\mathrm{lim}}_N/dp^N,X_1=_p,$$
$$\frac{X^{}=}{0<a<dp(a,p)=1a+dp_p,}$$
$$a+dp^N_p=\{xXxa(moddp^N)\},$$
where $`a`$ lies in $`0a<dp^N,`$ (cf. , ).
The $`p`$-adic absolute value in $`_p`$ is normalized so that $`|p|_p=\frac{1}{p}.`$ Let $`q`$ be variously considered as an indeterminate a complex number $`q`$, or a $`p`$-adic number $`q_p.`$ If $`q,`$ we always assume $`|q|<1.`$ If $`q_p,`$ we always assume $`|q1|_p<p^{\frac{1}{p1}}`$, so that $`q^x=\mathrm{exp}(x\mathrm{log}q)`$ for $`|x|_p1.`$ Throughout this paper, we use the following notation :
$$[x]_q=[x:q]=\frac{1q^x}{1q}.$$
We say that $`f`$ is uniformly differentiable function at a point $`a_p`$– and denote this property by $`fUD(_p)`$– if the difference quotients
$$F_f(x,y)=\frac{f(x)f(y)}{xy},$$
have a limit $`l=f^{}(a)`$ as $`(x,y)(a,a),`$ cf. . For $`fUD(_p),`$ let us start with the expression
$$\frac{1}{[p^N]_q}\underset{0j<p^N}{}q^if(j)=\underset{0j<p^N}{}f(j)\mu _q(j+p^N_p),\text{ cf. [2], }$$
representing $`q`$-analogue of Riemann sums for $`f`$.
The integral of $`f`$ on $`_p`$ will be defined as limit ($`n\mathrm{}`$) of these sums, when it exists. The $`q`$-Volkenborn integral of a function $`f\text{UD}(_p)`$ is defined by
$$__pf(x)𝑑\mu _q(x)=\underset{N\mathrm{}}{lim}\frac{1}{[p^N]_q}\underset{0j<p^N}{}f(j)q^j.$$
Note that if $`f_nf`$ in $`UD(Z_p)`$; then
$$__pf_n(x)𝑑\mu _q(x)__pf(x)𝑑\mu _q(x).$$
In , the $`q`$-Bernoulli polynomials are defined by
$$\beta _{m,q}(x)=__p[x+x_1]_q^m𝑑\mu _q(x_1).$$
In this paper, we investigate some properties of $`q`$-Volkenborn integrals on $`_p`$ in the meaning of fermionic and give some formulae which are related to $`q`$-numbers and polynomials in these integrals.
## §2. On the applications of $`q`$-Volkenborn integrals
For any positive integer $`N,`$ we set
$$\mu _q(a+dp^N_p)=\frac{q^a}{[dp^N]_q},$$
and this can be extended to distribution on $`X`$ (see ).
This distribution yields an integral for each non-negative integer $`m`$:
$`{\displaystyle __p}[a]_q^m𝑑\mu _q(a)`$ $`={\displaystyle _X}[a]_q^m𝑑\mu _q(a)`$
$`=\beta _{m,q}={\displaystyle \frac{1}{(1q)^m}}{\displaystyle \underset{i=0}{\overset{m}{}}}\left(\begin{array}{c}m\\ i\end{array}\right)(1)^i{\displaystyle \frac{i+1}{[i+1]}}.`$
The $`q`$-Bernoulli polynomials in variable $`x`$ in $`_p`$ with $`|x|_p1`$ are defined by
$$\beta _{n,q}(x)=__p[x+t]_q^n𝑑\mu _q(t).$$
These can be written as
$`\beta _{n,q}(x)`$ $`={\displaystyle \underset{i=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ i\end{array}\right)q^{ix}\beta _{i,q}[x]_q^{ni}`$
$`={\displaystyle \frac{1}{(1q)^m}}{\displaystyle \underset{i=0}{\overset{n}{}}}\left(\begin{array}{c}n\\ i\end{array}\right)(1)^iq^{xi}{\displaystyle \frac{i+1}{[i+1]_q}},(\text{cf.}[3]).`$
In the meaning of fermionic, we consider the $`q`$-numbers by using $`q`$-Volkenborn integrals as follows:
$$_{X_m}[x]_q^k𝑑\mu _q(x)=__p[x]_q^k𝑑\mu _q(x)=K_{k,q}\text{ for }k,m\text{ }.$$
Thus, we note that
$$K_{k,q}=[2]_q\left(\frac{1}{1q}\right)^k\underset{l=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{k}{l}\right)(1)^l\frac{1}{1+q^{l+1}},$$
$`1`$
where $`\left(\genfrac{}{}{0pt}{}{k}{i}\right)`$ is the binomial coefficient.
For $`x_p,`$ we define $`q`$-polynomials as follows:
$$__p[x+y]_q^k𝑑\mu _q(y)=K_{k,q}(x).$$
$`2`$
By (2), we easily see that
$$K_{k,q}(x)=\underset{n=0}{\overset{k}{}}\left(\genfrac{}{}{0pt}{}{k}{n}\right)[x]_q^{kn}q^{nx}K_{n,q}.$$
$`3`$
We now observe that
$$\frac{1}{[p^N]_q}\underset{n=0}{\overset{p^N1}{}}[x+y]_q^n(q)^y=[2]_q\left(\frac{1}{1q}\right)^n\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)(1)^kq^{xk}\frac{1}{1+q^{p^N}}\frac{1+q^{p^N(k+1)}}{1+q^{k+1}}.$$
$`4`$
Thus, we have
$$K_{n,q}(x)=__p[x+y]_q^n𝑑\mu _q(y)=[2]_q\left(\frac{1}{1q}\right)^n\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)(1)^kq^{xk}\frac{1}{1+q^{k+1}}.$$
$`5`$
By using the definition of Eq.(2), we will give the distribution of $`q`$-polynomials. From the definition of $`q`$-Volkenborn integral, we derive the below formula:
$$_{X_m}[x+y]_q^n𝑑\mu _q(y)=\frac{[m]_q^m}{[m]_q}\underset{a=0}{\overset{m1}{}}(1)^aq^a__p[\frac{a+x}{m}+y]_{q^m}^n𝑑\mu _{q^m}(y),\text{ if }m\text{ is odd}.$$
$`6`$
By (2) and (6), we obtain the below equation. Let $`m`$ be the odd integer. Then we have
$$K_{n,q}(x)=\frac{[m]_q^n}{[m]_q}\underset{a=0}{\overset{m1}{}}(1)^aq^aK_{n,q^m}(\frac{a+x}{m}).$$
$`7`$
Let $`\chi `$ be the primitive Dirichlet’s character with conductor $`f.`$ Then we also consider the below generalized $`q`$-Bernoulli numbers attached to $`\chi `$:
$$K_{n,\chi ,q}=_{X_f}\chi (x)[x]_q^n𝑑\mu _q(x).$$
$`8`$
By simple calculation in Eq.(8), we easily see that
$$K_{n,\chi ,q}=\frac{[f]_q^n}{[f]_q}\underset{a=0}{\overset{f1}{}}\chi (a)(1)^aq^a__p[\frac{a}{f}+x]_{q^f}^n𝑑\mu _{q^f}(x),\text{ if }f\text{ is odd integer.}$$
$`9`$
By (9) and (2), we obtain the following:
$$K_{n,\chi ,q}=\frac{[f]_q^n}{[f]_q}\underset{a=0}{\overset{f1}{}}\chi (a)(1)^aq^aK_{n,q^f}(\frac{a}{f}),$$
$`10`$
where $`\chi `$ is a primitive Dirichlet’s character with conductor odd integer $`f`$. It seems to be interesting to study the Kummer congruences for the generalized $`q`$-numbers attached to $`\chi `$ in the Eq.(10) because the $`q`$-Volkenborn measure $`d\mu _q(x)\frac{[2]_q}{2}(1)^xq^x`$.
## §3. Conclusion
In complex plane, we assume that $`q`$ with $`|q|<1.`$ It was well known that the Euler numbers are defined by
$$\frac{2}{e^t+1}=\underset{n=0}{\overset{\mathrm{}}{}}E_n\frac{t^n}{n!},\text{ for }|t|<2\pi \text{, cf. [1, 2, 3].}$$
$`11`$
Let $`F_q(t)=_{n=0}^{\mathrm{}}K_{n,q}\frac{t^n}{n!}`$. Then we have
$$F_q(t)=e^{\frac{t}{1q}}\underset{j=0}{\overset{\mathrm{}}{}}\frac{1+q}{1+q^{j+1}}(1)^j\left(\frac{1}{1q}\right)^j\frac{t^j}{j!}.$$
$`12`$
By (12), we can find the nice generating function of $`K_{n,q}`$ as follows:
$$F_q(t)=\underset{n=0}{\overset{\mathrm{}}{}}K_{n,q}\frac{t^n}{n!}=[2]_q\underset{n=0}{\overset{\mathrm{}}{}}(1)^nq^ne^{[n]_qt}.$$
$`13`$
Note that $`lim_{q1}F_q(t)=\frac{2}{e^t+1}.`$ By (11), we see that $`lim_{q1}K_{n,q}=E_n.`$ Hence, $`K_{n,q}`$ seems to be the nice $`q`$-analogue of $`E_n`$.
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# TKYNT-05-16, 2005/June Lee-Yang zero analysis for the study of QCD phase structure
## I Introduction
In the last few years, remarkable progress in exploring the QCD phase structure in the temperature $`(T)`$ and quark chemical potential $`(\mu _q)`$ plane have been made in numerical studies of lattice QCD. The phase transition line, separating hadron phase and quark-gluon plasma phase, was investigated from $`\mu _q=0`$ to finite $`\mu _q`$ , and the equation of state was also analyzed quantitatively at low density .
Among others, study of the endpoint of the first order phase transition line in the $`(T,\mu _q)`$ plane, whose existence is suggested by phenomenological studies , is particularly important both from the experimental and theoretical point of view. To locate the critical endpoint, Fodor and Katz investigated the positions of the Lee-Yang zeros (to be explained more below) in the complex $`\beta =6/g^2`$ plane using lattices with different spatial volumes, and examined the finite-volume scaling behavior of a Lee-Yang zero closest to the real axis. There are also studies in which the behavior of the critical endpoint as a function of the quark masses is examined by using the property that a critical endpoint exists at $`\mu _q=0`$ in the very small quark mass region for QCD with three flavors having degenerate quark masses. For extrapolating the result to the case with physical quark masses, an approach on the basis of the Taylor expansion in terms of $`\mu _q/T`$ and that of the imaginary chemical potential have been developed. Moreover, a study of phase-quenched finite density QCD, i.e. simulations with an isospin chemical potential, has been discussed in Ref. . The radius of convergence in the framework of the Taylor expansion of the grand canonical potential can establish a lower bound of the location of the critical endpoint . Also, the Glasgow method is an interesting approach for the study of QCD at non-zero baryon density.
In this paper, we focus our attention on the method of the Lee-Yang zero applied to finite density QCD. The Lee-Yang zero analysis is a popular method that is used to investigate the order of phase transitions. In order to study the existence of singularities in the thermodynamic limit (infinite volume limit), Lee and Yang proposed the following approach: Because there are no thermodynamic singularities as long as the volume is finite, the partition function $`𝒵`$ is always non-zero; it can develop zeros only in the infinite volume limit. However, if a real parameter of the model is extended into the complex parameter plane, a singularity, characterized by $`𝒵=0`$, can appear in the complex parameter plane even in a finite volume. These zeros are called the Lee-Yang zeros. Therefore, one can find the position of a singularity by exploring the position of the Lee-Yang zero in the complex parameter plane as a function of volume, and extrapolating the position of the Lee-Yang zero in the thermodynamic limit. For a system with a first order phase transition, the position of the nearest Lee-Yang zero approaches to the real axis in inverse proportion to the volume. On the other hand, the Lee-Yang zeros do not reach the real axis for crossover transition, i.e. rapid change without any thermodynamic singularities. For QCD at non-zero baryon density, we expect a rapid crossover transition in the low density regime which changes into a first order phase transition beyond a critical value of the density.
As we mentioned, Fodor and Katz have investigated the finite volume dependence of the Lee-Yang zeros in the complex $`\beta 6/g^2`$ plane for various values of the chemical potential. To carry out such analysis, the reweighting method is adopted, in which one performs simulations at $`\mu _q=0`$, and then corrects for the modified Boltzmann weight in the measurement of observables. In this case a famous problem arises for large $`\mu _q/T`$ and large volume, which is called “sign problem”. The sign problem is caused by complex phase fluctuations of the fermion determinant. In the region of small $`\mu _q/T`$, the phase fluctuations are not large and the sign problem is not serious. However, if the sign of the modification factor changes frequently during subsequent Monte-Carlo steps for large $`\mu _q/T`$, the statistical error becomes larger than the expectation values in general.
We find that the sign problem induces a serious problem in the finite volume scaling analysis of the Lee-Yang zero, which is used by Fodor and Katz. For any non-zero $`\mu _q`$ the normalized partition function calculated on the real axis with necessarily limited statistics will numerically always be consistent with zero once the volumes grow large. This means that the scaling behavior suggesting a crossover transition will not be obtained for the case with the sign problem, which is in contrast with the usual expectation.
Before discussing the case of non-zero density QCD, we study in Sec. II, as a simple example, the Lee-Yang zeros in the complex $`\beta `$ plane for SU(3) pure gauge theory by analyzing data from Monte-Carlo simulations. This model has a first order phase transition and simulations are much easier than for QCD at non-zero baryon density. Moreover, the pure gauge theory does, of course, not have a sign problem. Hence it is a good example to demonstrate how the Lee-Yang analysis works in the complex $`\beta `$ plane. In addition, it will become clear during this exercise that the complex phase fluctuations arising from the imaginary part of $`\beta `$ near Lee-Yang zeros are quite similar to those coming from the quark determinant where the sign problem exists for non-zero density QCD. The problem of the complex measure is reviewed in Sec. III. There we also comment on the reweighting method for the study at non-zero baryon density. In Sec. IV, we discuss a problem which arises when we apply the Lee-Yang zero analysis for non-zero baryon density QCD by using the reweighting technique, and consider possible other approaches in the framework of the Lee-Yang zero analysis for the investigation of the critical endpoint. Conclusions and discussions are given in Sec. V.
## II Lee-Yang zero for SU(3) pure gauge theory
### A General remarks
In this section, we apply the method of Lee-Yang zeros to the SU(3) pure gauge theory (quenched QCD).<sup>*</sup><sup>*</sup>* A pioneering study has been done for lattices with $`N_\tau =2`$ in Ref. . The phase transition of the SU(3) pure gauge theory is known to be of first order , which is expected from the corresponding Z(3) spin models. The pure gauge theory is controlled by only one parameter $`\beta =6/g^2`$ with the partition function,
$`𝒵={\displaystyle 𝒟Ue^{6\beta N_{\mathrm{site}}P}},`$ (1)
where $`P`$ is an averaged plaquette $`P=(_{x,\mu <\nu }W_{\mu \nu }^{1\times 1}(x))/(6N_{\mathrm{site}})`$, and $`W_{\mu \nu }^{1\times 1}`$ is $`1\times 1`$ Wilson loop operator for the lattice size $`N_{\mathrm{site}}=N_\sigma ^3\times N_\tau `$. Here, $`N_\sigma `$ and $`N_\tau `$ are spatial and temporal extension of the lattice. We extend the real parameter $`\beta `$ into the complex plane $`(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}})`$, and determine the position of Lee-Yang zeros, at which $`𝒵(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}})=0`$ is satisfied, by numerical simulations. We use standard Monte-Carlo techniques; configurations $`\{U_\mu \}`$ are generated with the probability of the Boltzmann weight. The expectation value of an operator $`𝒪[U_\mu ]`$, $`𝒪`$, is then calculated by taking an average over the configurations. We expect that the position of the Lee-Yang zero $`(\beta _{\mathrm{Re}}^0,\beta _{\mathrm{Im}}^0)`$ approaches the real $`\beta `$ axis in the infinite volume limit; with $`\beta _{\mathrm{Im}}^01/VN_\sigma ^3`$ for a first order phase transition.
In carrying out the above calculation two problems arise: One is that the Monte Carlo method is applicable only to the expectation values of physical quantities but not to the partition function itself. Another problem is that the measure is complex for a complex coupling $`\beta `$, and hence we cannot apply the Monte Carlo method directly, since the probabilities (Boltzmann weights) must be real and positive. To avoid these problems, we introduce the normalized partition function $`𝒵_{\mathrm{norm}}`$ together with the reweighting technique,
$`𝒵_{\mathrm{norm}}(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}})`$ $``$ $`\left|{\displaystyle \frac{𝒵(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}})}{𝒵(\beta _{\mathrm{Re}},0)}}\right|=\left|{\displaystyle \frac{𝒟Ue^{6(\beta _{\mathrm{Re}}+i\beta _{\mathrm{Im}})N_{\mathrm{site}}P}}{𝒟Ue^{6\beta _{\mathrm{Re}}N_{\mathrm{site}}P}}}\right|`$ (2)
$`=`$ $`\left|e^{6i\beta _{\mathrm{Im}}N_{\mathrm{site}}P}_{(\beta _{\mathrm{Re}},0)}\right|=\left|e^{6i\beta _{\mathrm{Im}}N_{\mathrm{site}}\mathrm{\Delta }P}_{(\beta _{\mathrm{Re}},0)}\right|.`$ (3)
Here $`\mathrm{\Delta }P=PP`$ and $`\left|\mathrm{exp}(6i\beta _{\mathrm{Im}}N_{\mathrm{site}}P)\right|=1`$. Because the denominator $`𝒵(\beta _{\mathrm{Re}},0)`$ is always finite for any finite volume, the position of $`𝒵(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}})=0`$ can be identified by analyzing $`𝒵_{\mathrm{norm}}`$. Although the partition function is not zero for $`\beta _{\mathrm{Im}}=0`$, it can be zero at some points in the $`(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}})`$ plane, when the complex phase factor in Eq.(3) changes sign frequently on the generated configurations. For the determination of the critical point in the original theory, i.e. on the real $`\beta `$ axis, the position of the nearest Lee-Yang zero should be investigated as a function of the volume $`V=N_\sigma ^3`$.
The mechanism that leads the occurrence of a Lee-Yang zero in Eq.(3) is quite similar to that which limits the applicability of the reweighting method for QCD with finite chemical potential . We will discuss this in more detail in Sec. III. At a point for which the width of the probability distribution of $`6\beta _{\mathrm{Im}}N_{\mathrm{site}}P`$ is smaller than $`O(\pi /2)`$, the sign of the complex phase does not change. Therefore, the standard deviation of the plaquette distribution is required to be larger than $`\pi /(12\beta _{\mathrm{Im}}N_{\mathrm{site}})`$ in the region where Lee-Yang zeros exist. Moreover, because the square of standard deviation is in proportion to the value of the plaquette susceptibility, the position of its maximum must agree with the position where $`Z_{\mathrm{norm}}`$ becomes minimal as function of $`\beta _{\mathrm{Re}}`$ for fixed $`\beta _{\mathrm{Im}}`$. Hence, the real part, $`\beta _{\mathrm{Re}}`$, of the position of the nearest Lee-Yang zero must be consistent with the peak position of the plaquette susceptibility. The method to find a critical point from the location of Lee-Yang zeros thus is essentially the same as the method which determines a critical point through the location of the peak position of the susceptibility and its finite volume scaling.
Here, it is instructive to introduce a probability distribution function for the plaquette, $`w(P)`$, which is defined by
$`w(P^{})={\displaystyle \frac{1}{𝒵}}{\displaystyle 𝒟U\delta (P^{}P)e^{6\beta _{\mathrm{Re}}N_{\mathrm{site}}P}},`$ (4)
where $`\delta (x)`$ is the delta function. Then, Eq.(3) can be rewritten as
$`𝒵_{\mathrm{norm}}(\beta )=\left|{\displaystyle dPe^{6i\beta _{\mathrm{Im}}N_{\mathrm{site}}\mathrm{\Delta }P}w(P)}\right|.`$ (5)
This means that the partition function $`𝒵_{\mathrm{norm}}`$ as a function of $`6\beta _{\mathrm{Im}}N_{\mathrm{site}}`$ is obtained through a Fourier transformation of $`w(P)`$.
Using this equation, the relation between the scaling behavior of the Lee-Yang zeros in the infinite volume limit and the distribution function of the plaquette becomes clearer. In Monte-Carlo simulations, configurations with probabilities proportional to their Boltzmann weight are generated by a computer, and we obtain a distribution function of the plaquette from the histogram of the plaquette. The histogram has usually a Gaussian shape at a normal, non-critical point, but it deviates from the Gaussian form near a critical point, and attains a double peak shape at a first order transition point, corresponding to the coexistence of two phases.
For the case of a non-singular point of $`\beta _{\mathrm{Re}}`$ or a crossover pseudo-critical point, where the distribution is expected to be a Gaussian function, the point of $`Z_{\mathrm{norm}}=0`$ does not exist except in the limit of $`\beta _{\mathrm{Im}}N_{\mathrm{site}}\mathrm{}`$ or $`\mathrm{}`$, because the function which is obtained through a Fourier transformation of a Gaussian function again is a Gaussian function. Of course, results of numerical simulations have statistical errors, hence $`𝒵_{\mathrm{norm}}`$ can become zero “within errors”, if the expectation value and the error become of the same order. However, in this case, the point at which $`Z_{\mathrm{norm}}=0`$ appears at random in terms of $`\beta _{\mathrm{Im}}\times N_{\mathrm{site}}`$. Therefore, the volume dependence of the position of the Lee-Yang zero $`(\beta _{\mathrm{Re}}^0,\beta _{\mathrm{Im}}^0)`$ does not necessary to be $`\beta _{\mathrm{Im}}^01/V(N_\sigma ^3)`$ for fixed $`N_\tau `$.
On the other hand, in the case of a first order phase transition, we expect that the plaquette histogram has two peaks having the same peak height at the transition point. Performing the Fourier transformation of such a double peaked function leads to a function which has zeros periodically. For example, a distribution function $`w(P)`$ having two Gaussian peaks at $`\mathrm{\Delta }P=\pm A`$ leads to a normalized partition function $`𝒵_{\mathrm{norm}}`$ which has zeros at
$`\beta _{\mathrm{Im}}^0={\displaystyle \frac{\pi (2n+1)}{12N_{\mathrm{site}}A}},(n=0,1,2,3,\mathrm{}).`$ (6)
This is mathematically the same as that for the interference experiment using a laser and a double-slit. The Lee-Yang zeros correspond to dark lines (destructive interference) and they appear periodically as given in Eq.(6). Moreover, for a first order phase transition the difference of plaquette values in cold and hot phases, $`2A`$, is related to the latent heat $`\mathrm{\Delta }\epsilon `$, i.e. the energy difference between the hot and cold phases,
$`{\displaystyle \frac{\mathrm{\Delta }\epsilon }{T^4}}12AN_\tau ^4a{\displaystyle \frac{d\beta }{da}},`$ (7)
where $`a`$ is the lattice spacing. Since $`\mathrm{\Delta }\epsilon `$ is non-zero, $`A`$ does not vanish in the infinite volume limit $`(N_\sigma ^3\mathrm{})`$. Therefore, we find that in the infinite volume limit the nearest Lee-Yang zero approaches the real $`\beta `$ axis like $`\beta _{\mathrm{Im}}^01/V`$, which is consistent with the general argument on the Lee-Yang zero for a first order phase transition. We also emphasize that the isolated Lee-Yang zeros appear periodically. The distances to these points from the real axis are $`1,3,5,\mathrm{}`$ in units of the distance to the first Lee-Yang zero. This is also an important property, which is not observed for a crossover transition.
In addition, the discussion given for the plaquette distribution function can be extended to the analysis of fourth order Binder cumulants,
$`B_4{\displaystyle \frac{\mathrm{\Delta }P^4}{\mathrm{\Delta }P^2^2}},`$ (8)
which is an alternative to the method of Lee-Yang zeros often used to identify the order of a phase transition. The value of the Binder cumulant at the critical point depends on the universality class. In the case of a first order phase transition, assuming the plaquette distribution is a double peaked function, the Binder cumulants are estimated as
$`B_4={\displaystyle \frac{dP\mathrm{\Delta }P^4w(P)}{\left(dP\mathrm{\Delta }P^2w(P)\right)^2}}{\displaystyle \frac{A^4}{(A^2)^2}}1,`$ (9)
where the distance between two peaks is $`2A`$ and is wider than the width of each peak. On the other hand, when the distribution function can be modeled by a Gaussian function for a crossover transition, the Binder cumulants are given by
$`B_4{\displaystyle \frac{\sqrt{x/\pi }𝑑P\mathrm{\Delta }P^4e^{x\mathrm{\Delta }P^2}}{\left(\sqrt{x/\pi }𝑑P\mathrm{\Delta }P^2e^{x\mathrm{\Delta }P^2}\right)^2}}=\left(\sqrt{{\displaystyle \frac{x}{\pi }}}{\displaystyle \frac{d^2\sqrt{\pi /x}}{dx^2}}\right)/\left(\sqrt{{\displaystyle \frac{x}{\pi }}}{\displaystyle \frac{d\sqrt{\pi /x}}{dx}}\right)^2=3.`$ (10)
In a region where a first order phase transition changes to a crossover, the Binder cumulant changes rapidly from one to three. We expect to find such a region for full QCD at high temperature and density. The value of the Binder cumulant at the endpoint of the first order transition line, which is of second order, is determined by the universality class. Hence, the plaquette distribution function plays an important role for both methods to identify the order of a phase transition.
### B Numerical results
We calculate the normalized partition function for SU(3) pure gauge theory to find Lee-Yang zeros in the complex $`\beta `$ plane, using data for plaquettes obtained by QCDPAX in Ref. . There are five data sets measured at the transition point for $`N_\tau =4`$ and $`6`$. The spatial lattice sizes are $`24^2\times 36\times 4,12^2\times 24\times 4,36^2\times 48\times 6,24^3\times 6,`$ and $`20^3\times 6`$. $`O(10^6)`$ configurations are available for the analysis of each data set. The reweighting technique is also used for the real $`\beta `$ direction to analyze the Lee-Yang zeros in the complex $`\beta `$ plane for a data set obtained at only one $`\beta `$ point $`(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}})=(\beta _0,0).`$ The normalized partition function is given by
$`𝒵_{\mathrm{norm}}(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}})=\left|{\displaystyle \frac{e^{6i\beta _{\mathrm{Im}}N_{\mathrm{site}}\mathrm{\Delta }P}e^{6\beta _{\mathrm{Re}}N_{\mathrm{site}}\mathrm{\Delta }P}_{(\beta _0,0)}}{e^{6\beta _{\mathrm{Re}}N_{\mathrm{site}}\mathrm{\Delta }P}_{(\beta _0,0)}}}\right|.`$ (11)
Figure 1 shows the contour plot of $`𝒵_{\mathrm{norm}}`$ for the $`24^2\times 36\times 4`$ lattice. The simulation point is $`\beta _0=3.6492`$. In this definition, $`𝒵_{\mathrm{norm}}`$ is normalized to be one on the real $`\beta `$ axis. Circles at $`(\beta _{\mathrm{Re}}^0,\beta _{\mathrm{Im}}^0)(5.6925,0.0021)`$ and $`(5.6931,0.0056)`$ are Lee-Yang zeros. Since the SU(3) pure gauge theory has a first order phase transition, Lee-Yang zeros appear periodically. For this data, two clear peaks are visible in the plaquette histogram . The distance between these two peaks is $`2A0.003`$. The positions of the Lee-Yang zeros are consistent with $`\beta _{\mathrm{Im}}^0\pi /(12N_{\mathrm{site}}A)0.002`$ and $`3\pi /(12N_{\mathrm{site}}A)0.006`$, as given in Eq.(6).
Above property is not seen so clearly for lattices having small $`N_\sigma `$ and large $`N_\tau `$. The position of the next-to-leading zero points of $`𝒵_{\mathrm{norm}}`$ appear at random for the other data sets relative to the nearest zero point. The positions of Lee-Yang zeros are shown in Table I. We could not obtain clearly isolated Lee-Yang zeros for the lattices $`24^3\times 6`$ and $`20^3\times 6`$. The second nearest Lee-Yang zero to the real axis could be measured only for the $`24^2\times 36\times 4`$ lattice. The result on the $`36^2\times 48\times 6`$ lattice $`(\beta _0=5.8936)`$ is also shown in Fig. 2. Only the nearest Lee-Yang zero is obtained clearly. The Lee-Yang zero becomes less clear as $`N_\tau `$ increases and $`N_\sigma `$ decreases, hence simulations on lattices having large $`N_\sigma /N_\tau `$ seem to be necessary for the study of Lee-Yang zeros.
The values for $`\beta _{\mathrm{Im}}^0V`$ on $`24^2\times 36\times 4`$ and $`12^2\times 24\times 4`$ lattices are $`43.9(5)`$ and $`42.0(6)`$, respectively. These are roughly constant and suggest the scaling behavior of $`\beta _{\mathrm{Im}}^01/V`$ for a first order phase transition. Also, in the previous study for lattices with $`N_\tau =2`$ , the $`1/V`$ scaling behavior has been confirmed for $`N_\sigma =6,8,10`$ and $`12`$. However, for a more precise quantitative investigation that takes into account the errors, the spatial lattice size $`12^2\times 24`$ may not be large enough to check the $`1/V`$ scaling for $`N_\tau =4`$, since the difference of $`\beta _{\mathrm{Im}}^0V`$ is larger than the statistical error. We should fit the data obtained on more than two lattices by a curved function of $`1/V`$, to confirm through a $`1/V`$ scaling analysis that the phase transition of the SU(3) pure gauge theory is first order. E.g. for the study of the SU(2) gauge-Higgs model , the following fitting functions have been used, $`\mathrm{Im}\kappa _0(V)=\kappa _0^c+CV^\nu `$ and $`\mathrm{Im}\kappa _0(V)=\kappa _0^c+CV^1+DV^2`$ for a complex parameter $`\kappa `$ in the model with fitting parameters $`\kappa _0^c,C,D,`$ and $`\nu `$.
## III Complex phase fluctuation and sign problem
As seen in the previous section, the investigation of Lee-Yang zeros in the complex $`\beta `$ plane seems to be useful to identify the order of phase transition. However, if we try to extend this analysis to full QCD at non-zero baryon density, a serious problem arises. This problem is closely related to the sign problem for finite density QCD, since the normalized partition function can be zero in the complex $`\beta `$ plane due to fluctuations of the complex phase related to $`\beta _{\mathrm{Im}}`$ and also due to the complex phase from the quark determinant that causes the sign problem. Before discussing the Lee-Yang zero analysis for finite density QCD, we would like to review the sign problem briefly.
The main difficulty for studies at finite baryon density is that the Boltzmann weight is complex if the chemical potential is non-zero. In this case the Monte-Carlo method is not applicable directly, since configurations cannot be generated with a complex probability. One approach to avoid this problem is the reweighting method. We perform simulations at $`\mu =0`$, and incorporate the remaining part of the correct Boltzmann weight for finite $`\mu `$ in the calculation of expectation values. Expectation values $`𝒪`$ at $`(\beta ,\mu )`$ are thus computed by a simulation at $`(\beta ,0)`$ using the following identity,
$`𝒪_{(\beta ,\mu )}={\displaystyle \frac{𝒪e^{N_\mathrm{f}(\mathrm{ln}detM(\mu )\mathrm{ln}detM(0))}_{(\beta ,0)}}{e^{N_\mathrm{f}(\mathrm{ln}detM(\mu )\mathrm{ln}detM(0))}_{(\beta ,0)}}},`$ (12)
where $`M`$ is the quark matrix and $`N_\mathrm{f}`$ is the number of flavors ($`N_\mathrm{f}/4`$ for staggered type quarks instead of $`N_\mathrm{f}`$); $`\mu `$ is a quark chemical potential in lattice units, i.e. $`\mu \mu _qa=\mu _q/(N_\tau T)`$, and $`\mu _q`$ is the quark chemical potential in physical units. This is the basic formula of the reweighting method. However, because $`\mathrm{ln}detM(\mu )`$ is complex, the calculations of the numerator and denominator in Eq.(12) becomes in practice increasingly more difficult for larger $`\mu `$. We define the phase of the quark determinant $`\theta `$ by $`(detM(\mu ))^{N_\mathrm{f}/4}|detM(\mu )|^{N_\mathrm{f}/4}e^{i\theta }`$ for staggered type quarks. If the typical value of $`\theta `$ becomes larger than $`\pi /2`$, the real part of $`e^{i\theta }`$ $`(=\mathrm{cos}\theta )`$ changes its sign frequently. Eventually both the numerator and denominator of Eq.(12) become smaller than their statistical errors and Eq.(12) can no longer be evaluated. We call it the “sign problem”.
Here, the denominator of Eq.(12), or simply $`\mathrm{cos}\theta `$, is a good indicator for the occurrence of the sign problem. If this indicator is zero within statistical errors, Eq.(12) cannot be computed. In the following we give an estimate for the value of the complex phase. Since the direct calculation of the quark determinant is difficult except for calculations on small lattices, we expand $`\mathrm{ln}detM(\mu )`$ in a Taylor series,
$`\mathrm{ln}detM(\mu )\mathrm{ln}detM(0)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left[{\displaystyle \frac{^n(\mathrm{ln}detM)}{\mu ^n}}\right]{\displaystyle \frac{\mu ^n}{n!}}.`$ (13)
Then, we can easily separate it into real and imaginary parts because the even derivatives of $`\mathrm{ln}detM(\mu )`$ are real and the odd derivatives are purely imaginary . The complex phase $`\theta `$ is given by
$`\theta ={\displaystyle \frac{N_\mathrm{f}}{4}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{Im}{\displaystyle \frac{^{2n1}(\mathrm{ln}detM)}{\mu ^{2n1}}}{\displaystyle \frac{\mu ^{2n1}}{(2n1)!}},`$ (14)
for staggered type quarks at small $`\mu `$. The Taylor expansion coefficients are rather easy to calculate by using the stochastic noise method. The comparison between the value of $`\theta `$ with this approximation and the exact value has been done in Ref. , and the reliability and the application range have been discussed.
We use data for the Taylor expansion coefficients obtained in Ref. . The data were generated by using Symanzik-improved gauge and p4-improved staggered fermion actions. Coefficients up to $`O(\mu ^5)`$ have been calculated. Figure 3 shows the indicator $`\mathrm{cos}\theta `$ measured at $`\beta =\{3,60,3.65,`$ and $`3.68\}`$, for $`ma=0.1`$, corresponding to $`T/T_c=0.90,1.00,`$ and $`1.07`$, respectively, on a $`16^3\times 4`$ lattice As mentioned in Ref. , in the calculation using the stochastic noise method, the error due to the finite number of noise vector $`(N_{\mathrm{noise}})`$ is large for the calculation of $`\mathrm{cos}\theta `$ with $`N_{\mathrm{noise}}=10`$. For the purpose of this study we increased the number of noise vector to $`N_{\mathrm{noise}}=100`$. We checked that the difference between the results with $`N_{\mathrm{noise}}=50`$ and $`N_{\mathrm{noise}}=100`$ is about $`10\%`$ for the calculation of the position at which $`\mathrm{cos}\theta =0.1`$. . We also estimate the values of $`\mu _q/TN_\tau \mu `$ at which $`\mathrm{cos}\theta =0.1,0.01`$ and $`0.0`$. Results are given in Table II. The situation now is quite similar to the non-singular case of the normalized partition function in the complex $`\beta `$ plane discussed in the previous section. This becomes even more apparent if we consider for simplicity only the first term in the expansion of $`\theta `$ which is proportional to $`\mu `$. Then $`\mathrm{cos}\theta `$, $`\mu `$ and $`\mathrm{Im}[(\mathrm{ln}detM)/\mu ]`$ correspond to the normalized partition function, $`\beta _{\mathrm{Im}}`$ and plaquette, respectively.
Because the distribution of the complex phase $`\theta `$ is almost of Gaussian shape, the indicator, $`\mathrm{cos}\theta `$, decreases exponentially as $`\mu `$ increases, and it may cross zero at a point where the expectation value becomes smaller than the statistical error. Therefore, the points of $`\mathrm{cos}\theta =0`$ appear accidentally and the results given in Table II are unstable. Moreover, $`(\mathrm{ln}detM)/\mu `$ becomes larger as the volume increases, hence the indicator for the sign problem vanishes in the infinite volume limit for any non-zero $`\mu `$, which means that the range of applicability for the reweighting method approaches $`\mu =0`$ in the infinite volume limit.
## IV Lee-Yang zero analysis for finite density QCD
The Lee-Yang zero analysis for finite density QCD has been performed by Fodor and Katz . They measured the normalized partition function $`𝒵_{\mathrm{norm}}`$ using the reweighting method, and determined the points where $`𝒵_{\mathrm{norm}}=0`$ as a function of spatial volume. The normalized partition function is defined by
$`𝒵_{\mathrm{norm}}(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}},\mu )`$ $``$ $`\left|{\displaystyle \frac{𝒵(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}},\mu )}{𝒵(\beta _{\mathrm{Re}},0,0)}}\right|`$ (15)
$`=`$ $`\left|e^{6i\beta _{\mathrm{Im}}N_{\mathrm{site}}\mathrm{\Delta }P}e^{i\theta }\left|e^{(N_\mathrm{f}/4)(\mathrm{ln}detM(\mu )\mathrm{ln}detM(0))}\right|_{(\beta _{\mathrm{Re}},0,0)}\right|`$ (16)
or
$`𝒵_{\mathrm{norm}}(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}},\mu )`$ $``$ $`\left|{\displaystyle \frac{𝒵(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}},\mu )}{𝒵(\beta _{\mathrm{Re}},0,\mu )}}\right|`$ (17)
$`=`$ $`\left|{\displaystyle \frac{e^{6i\beta _{\mathrm{Im}}N_{\mathrm{site}}\mathrm{\Delta }P}e^{i\theta }\left|e^{(N_\mathrm{f}/4)(\mathrm{ln}detM(\mu )\mathrm{ln}detM(0))}\right|_{(\beta _{\mathrm{Re}},0,0)}}{e^{i\theta }\left|e^{(N_\mathrm{f}/4)(\mathrm{ln}detM(\mu )\mathrm{ln}detM(0))}\right|_{(\beta _{\mathrm{Re}},0,0)}}}\right|`$ (18)
for staggered type quarks. $`\theta `$ is the complex phase of $`\mathrm{exp}[(N_\mathrm{f}/4)(\mathrm{ln}detM(\mu )\mathrm{ln}detM(0))]`$. Since the numerator of Eq.(18), which is the same as $`𝒵_{\mathrm{norm}}`$ in Eq.(16), is required to be zero at a zero point of $`𝒵_{\mathrm{norm}}`$ in Eq.(18), we consider Eq.(16) as an indicator for the Lee-Yang zero.
Here, we notice that for $`\beta _{\mathrm{Im}}=0`$ this normalized partition function is exactly the same as the indicator for the sign problem, i.e. the denominator of Eq.(12). As discussed in the previous section, in any practical simulation this indicator will be consistent with zero within errors for large values of $`\mu `$. Moreover, the region where the indicator is non-zero becomes narrower as the volume increases, and this region vanishes in the infinite volume limit. Hence, the Lee-Yang zeros always approach the $`\beta _{\mathrm{Im}}=0`$ axis in the infinite volume limit for any finite $`\mu `$. This means that the scaling behavior for crossover will not be obtained for the case with the sign problem. This is clearly different from the usual expectation for the QCD phase diagram in the $`(T,\mu )`$ plane. Most model calculations suggest that the transition is crossover in the low density region. This might be a problem of the definition of the normalized partition function, Eq.(16). The normalized partition function on the real $`\beta `$ axis is exactly one for $`\mu =0`$, but it vanishes for finite $`\mu `$ in the infinite volume limit. Therefore, it is very difficult to distinguish the first order transition and crossover by investigating the position of Lee-Yang zeros as a function of spatial volume. This is the most important difference between the definitions for the pure gauge theory in Sec. II and QCD at non-zero density.
The critical endpoint is shown to be located at $`\mu _B=3\mu _q=725(35)\mathrm{MeV}`$ in Ref. and $`360(40)\mathrm{MeV}`$ in Ref. , which is inconsistent with the above argument. This may be a problem of the fitting function. In Refs. the position of Lee-Yang zeros has been fitted by $`\beta _{\mathrm{Im}}^0=A(1/V)+\beta _{\mathrm{Im}}^{\mathrm{}}`$, where $`A`$ and $`\beta _{\mathrm{Im}}^{\mathrm{}}`$ are fitting parameters. The first order transition and crossover have then been distinguished by the value of $`\beta _{\mathrm{Im}}^{\mathrm{}}`$. As we discussed in Sec. II, this fitting function is too simple to fit the data of Lee-Yang zeros obtained on lattices as small as those used in Ref. , i.e. $`V12^3`$. In fact, if one assumes a curved extrapolation function, all data in Table 1 of Ref. seems to approach $`\beta _{\mathrm{Im}}=0`$ in the $`1/V0`$ limit.
In our argument, the statistical error of $`Z_{\mathrm{norm}}`$, which is controlled by the number of configurations in the Monte-Carlo simulation, plays an important role. If the statistical error of $`Z_{\mathrm{norm}}`$ becomes much smaller than the mean value of $`Z_{\mathrm{norm}}`$ by increasing the statistics for each simulation, the method in Ref. would be applicable. However, one cannot satisfy this condition in general simply because we are looking for the Lee-Yang zero which gives $`Z_{\mathrm{norm}}=0`$. Namely, statistical error cannot be smaller than 0. Moreover, if the error of $`Z_{\mathrm{norm}}`$ is sizeable, there appear fake Lee-Yang zeros which are located even closer to the real $`\beta `$ axis than the true zero in the region where the mean value of $`Z_{\mathrm{norm}}`$ is smaller than the error. Since we adopt the closest zero as the Lee-Yang zero in the actual scaling analysis, we may thus misidentify the true zero by the fake one.
The above point can be seen explicitly for the second Lee-Yang zero of the SU(3) gauge theory with $`N_\tau =6`$ shown in Fig. 2. Theoretically, we expect the second Lee-Yang zero exists around $`\beta _{\mathrm{Im}}0.013`$, i.e. three times larger than that of the first Lee-Yang zero as shown in Fig. 1. However we find several $`\beta `$ which give $`Z_{\mathrm{norm}}=0`$ in the region of $`\beta _{\mathrm{Im}}>0.006`$, and they distribute randomly. If we identify the second nearest point as the second Lee-Yang zero, the resulting $`\beta _{\mathrm{Im}}`$ is much smaller than the theoretical expectation. This problem is caused by the existence of the region in the complex $`\beta `$-plane where the statistical error of $`Z_{\mathrm{norm}}`$ is larger than the mean value as discussed above. As the statistics is increased for fixed $`V`$, such a region should become smaller and fake Lee-Yang zeros should disappear.
Now we discuss how the above situation changes by increasing the volume $`V`$. Fortunately, in the SU(3) gauge theory, the location of Lee-Yang zero can be determined better also as the volume increases as shown in Sec. II. In this case, the validity of the scaling assumption becomes better as the volume increases. On the other hand, for the case with the sign problem, the error normalized by the mean value grows exponentially as a function of volume. Then the size of the region having fake Lee-Yang zeros cannot be made smaller unless one has exponentially large statistics. This leads to the conclusion that the quality of the scaling analysis is not improved by increasing the volume, and any reliable information about Lee-Yang zeros in the infinite volume limit cannot be obtained. This is the reason why the serious problem of identifying the critical endpoint is intimately related to the sign problem.
This discussion suggests that the finite volume scaling analysis suffers serious damage through the unsolved sign problem, and it is very difficult to apply the criterion used by Fodor and Katz for the investigation of the critical endpoint in practice. However, the property of the second nearest Lee-Yang zero characteristic for a first order transition in Fig. 1, i.e. the fact that the distance to the second Lee-Yang zero from the real axis is three times larger than that of the first Lee-Yang zero, could be investigated on a finite lattice. This study is possible for small $`\mu `$ without taking the infinite volume limit. On the other hand, we do not expect any isolated Lee-Yang zero for a crossover transition, hence we may be able to determine the order of phase transition by investigating the distribution of Lee-Yang zeros in the complex $`\beta `$ plane. Although the measurements of the second Lee-Yang zero may require large lattice sizes and high statistics, as seen for the case of pure SU(3) gauge theory, it may be possible to find the region of the first order phase transition, if the critical endpoint exists in the low density region.
## V Conclusions
We commented on the Lee-Yang zero analysis for the study of the critical endpoint in the $`(T,\mu _q)`$ phase diagram. It is found that the Lee-Yang zero analysis at non-zero baryon density encounters a serious problem. The complex phases of the quark determinant and the complex $`\beta `$ are mixed at non-zero chemical potential. In this case, in practical simulations with limited statistics the normalized partition function can develop zeros even on the real $`\beta `$ axis for large $`\mu _q`$ in finite volumes. Moreover, in the infinite volume the normalized partition function is always zero except for $`\mu _q=0`$. This means that the nearest Lee-Yang zero always approaches the real $`\beta `$ axis in the infinite volume limit. The scaling behavior suggesting a crossover transition thus will not be obtained. This is clearly different from usual expectations for the QCD phase diagram. To avoid this problem, the sign problem must be removed by careful treatments increasing the number of configurations exponentially as the volume or $`\mu _q/T`$ increases, otherwise the finite volume scaling behavior for the position of Lee-Yang zeros, which has been analyzed by Fodor and Katz , does not provide an appropriate criterion for the investigation of the order of the phase transition.
To make the underlying problem more transparent, we applied the Lee-Yang zero analysis to the SU(3) pure gauge theory, which does not have a sign problem and for which the simulations are much easier. Lee-Yang zeros are found in the complex $`\beta `$ plane. They appear periodically as expected by the discussion using a plaquette distribution function for a first order phase transition. The positions of the first Lee-Yang zero on two lattices having different volume sizes are roughly consistent with the finite size scaling behavior for a first order phase transition, i.e. $`\beta _{\mathrm{Im}}^01/V`$. However, for quantitative analysis it is necessary to fit data from more than two different lattice sizes by a curved function to study the order of the phase transition. It is found, in this analysis, that complex phase fluctuations arising from the imaginary part of $`\beta `$ play an important role, and the mechanism that leads to the appearance of the Lee-Yang zeros is quite similar to the situation in QCD where the sign problem is present.
The property of a first order phase transition that isolated Lee-Yang zeros appear periodically at $`\beta _{\mathrm{Im}}C(2n+1)`$, where $`C`$ is the distance to the nearest Lee-Yang zero and $`n`$ is an integer, is free from the problems that arise in the infinite volume limit. Therefore, to investigate the pattern of the appearance of Lee-Yang zeros in the $`(\beta _{\mathrm{Re}},\beta _{\mathrm{Im}})`$ plane is important. For this calculation, the simulations by using high statistics data and large lattice size are indispensable. Further studies are clearly important to find the endpoint of the first order phase transition line in the $`(T,\mu _q)`$ plane.
Recently, a close relation between the strength of the sign problem and the position of the phase transition line for pion condensation in phase quenched QCD has been discussed in Ref. . There it has been found that the endpoints of the first order transition line determined in Ref. are located near the phase transition line of pion condensation. These results may relate to our discussion given here.
Moreover, the pathologies in the Glasgow method have been discussed in Ref. . Similar problems arise also in the Glasgow method. It would be interesting to consider the relation between the pathologies in the Glasgow method and those in the Lee-Yang zero analysis.
## Acknowledgments
I would like to thank T. Hatsuda and K. Kanaya for fruitful discussions and useful comments on this manuscript. I also thank C.R. Allton, M. Döring, S.J. Hands, O. Kaczmarek, F. Karsch, E. Laermann, and K. Redlich for discussions, comments and allowing me to use the data of the Bielefeld-Swansea collaboration partly published in Ref. . The data given by QCDPAX in Ref. are used for the analysis in Sec. II. This work is supported by BMBF under grant No.06BI102.
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